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# 1 Introduction
## 1 Introduction
As the materials sciences and technologies continue their rapid development, realistic possibilities are emerging of realizing so–called *metamaterials* with novel and hitherto unconsidered optical/electromagnetic properties. A prime example is provided by the recently discovered metamaterials which support planewave propagation with *negative phase velocity* (NPV), and thereby negative refraction. Until 2000, little attention had been paid to the phenomenon of negative refraction. Since 2000, there has been an explosion of interest in negative refraction , following experimental reports of a metamaterial which supports negative refraction in the microwave regime .
Naturally–occurring uniaxial crystals have been extensively studied ever since the earliest days of the optical sciences. However, the electromagnetic properties of uniaxial mediums have recently been revisited by theoreticians in consideration of the prospects for NPV propagation in such mediums . A closely related issue concerns uniaxial dielectric–magnetic mediums with indefinite constitutive dyadics .
The defining characteristic of a uniaxial dielectric medium is a distinguished axis of symmetry, known as the optic axis. Mathematically, the permittivity dyadic of a uniaxial dielectric medium may be expressed as
$$\underset{¯}{\underset{¯}{ϵ}}=ϵ\underset{¯}{\underset{¯}{I}}+\left(ϵ_xϵ\right)\underset{¯}{\overset{^}{x}}\underset{¯}{\overset{^}{x}},$$
(1)
where a coordinate system has been selected in which the direction of the optic axis coincides with the direction of the unit vector $`\underset{¯}{\overset{^}{x}}`$ lying along the $`x`$ axis, and $`\underset{¯}{\underset{¯}{I}}`$ denotes the 3$`\times `$3 identity dyadic. The real–valued parameter
$$\gamma =\{\begin{array}{ccc}\frac{ϵ_x}{ϵ}\hfill & \text{for}& \hfill ϵ_x,ϵ\\ & & \\ \frac{\text{Re}\left\{ϵ_x\right\}}{\text{Re}\left\{ϵ\right\}}\hfill & \text{for}& \hfill ϵ_x,ϵ\end{array}$$
(2)
may be usefully employed to characterize planewave propagation in the medium specified by (1). The upper expression is appropriate to nondissipative mediums whereas the lower expression is appropriate to dissipative mediums.
The electromagnetic/optical properties of uniaxial mediums with $`\gamma >0`$ — this category includes naturally–occurring uniaxial crystals — have long been established. Comprehensive descriptions can be found in standard works . Uniaxial mediums with $`\gamma <0`$ are much more exotic. Interest in these mediums stems from their potential applications in negatively refracting scenarios and in diffraction gratings , for example.
Planewave propagation in a uniaxial medium is characterized in terms of a dispersion relation which is quadratic in terms of the corresponding wavevector components. The dispersion relations for nondissipative mediums with $`\gamma >0`$ have an elliptical representation, whereas a hyperbolic representation is associated with $`\gamma <0`$. In this communication we investigate the planewave characteristics and conceptualization of uniaxial dielectric mediums with hyperbolic dispersion relations.
## 2 Planewave analysis
The propagation of plane waves with field phasors
$$\begin{array}{c}\underset{¯}{E}(\underset{¯}{r})=\underset{¯}{E}_0\mathrm{exp}\left(i\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{r}\right)\hfill \\ \underset{¯}{H}(\underset{¯}{r})=\underset{¯}{H}_0\mathrm{exp}\left(i\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{r}\right)\hfill \end{array}\}$$
(3)
in the uniaxial dielectric medium specified by the permittivity dyadic (1) is investigated. The permittivity parameters are generally complex–valued; i.e., $`ϵ,ϵ_x`$. The wavevector $`\underset{¯}{k}`$ is taken to be of the form
$$\underset{¯}{k}=\alpha \underset{¯}{\overset{^}{x}}+\beta \underset{¯}{\overset{^}{z}},$$
(4)
where $`\alpha `$, $`\beta `$ and $`\underset{¯}{\overset{^}{z}}`$ is the unit vector directed along the $`z`$ axis. This form of $`\underset{¯}{k}`$ is appropriate to planar boundary value problems and from the practical viewpoint of potential optical devices . We note that the plane waves (3) are generally nonuniform.
The source–free Maxwell curl postulates
$$\begin{array}{c}\times \underset{¯}{E}(\underset{¯}{r})=i\omega \underset{¯}{B}(\underset{¯}{r})\hfill \\ \times \underset{¯}{H}(\underset{¯}{r})=i\omega \underset{¯}{D}(\underset{¯}{r})\hfill \end{array}\}$$
(5)
yield the vector Helmholtz equation
$$\left[\left(\times \underset{¯}{\underset{¯}{I}}\right)\text{ }\text{}\text{ }\left(\times \underset{¯}{\underset{¯}{I}}\right)\mu _0\omega ^2\underset{¯}{\underset{¯}{ϵ}}\right]\text{ }\text{}\text{ }\underset{¯}{E}(\underset{¯}{r})=\underset{¯}{0},$$
(6)
with $`\mu _0`$ being the permeability of free space. Combining (3) with (6) yields the planewave dispersion relation
$$\left(\alpha ^2+\beta ^2ϵ\mu _0\omega ^2\right)\left(\alpha ^2ϵ_x+\beta ^2ϵϵ_xϵ\mu _0\omega ^2\right)=0.$$
(7)
In the following we consider the time–averaged Poynting vector
$$\underset{¯}{P}(\underset{¯}{r})=\frac{\mathrm{exp}\left(2\text{Im}\left\{\beta \right\}z\right)}{2\mu _0\omega }\text{Re}\left\{|\underset{¯}{E}_0|^2\underset{¯}{k}^{}\left(\underset{¯}{E}_0\text{ }\text{}\text{ }\underset{¯}{k}^{}\right)\underset{¯}{E}_0^{}\right\}.$$
(8)
Evanescent plane waves are characterized by $`\text{Im}\left\{\beta \right\}>0`$. The scenario characterized by $`\text{Im}\left\{\beta \right\}<0`$ is not physically plausible for passive mediums and is therefore not considered here.
### 2.1 Ordinary wave
The ordinary wavevector
$$\underset{¯}{k}_{or}=\alpha \underset{¯}{\overset{^}{x}}+\beta _{or}\underset{¯}{\overset{^}{z}},$$
(9)
arises from the dispersion relation (7) with components satisfying
$`\alpha ^2+\beta _{or}^2`$ $`=`$ $`\omega ^2ϵ\mu _0.`$ (10)
The vector Helmholtz equation (6) yields the eigenvector solution $`\underset{¯}{E}_0=E_y\underset{¯}{\overset{^}{y}}`$, directed parallel to the unit vector $`\underset{¯}{\overset{^}{y}}`$ lying along the $`y`$ axis, where the complex–valued magnitude $`E_y`$ is determined by the initial/boundary conditions. Consequently, the time–averaged Poynting vector reduces to
$$\underset{¯}{P}(\underset{¯}{r})=\frac{\mathrm{exp}\left(2\text{Im}\left\{\beta _{or}\right\}z\right)}{2\omega \mu _0}|E_y|^2\text{Re}\left\{\underset{¯}{k}_{or}^{}\right\}.$$
(11)
Since
$`\text{Re}\left\{\underset{¯}{k}_{or}\right\}\text{ }\text{}\text{ }\underset{¯}{P}(\underset{¯}{r})`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left(2\text{Im}\left\{\beta _{or}\right\}z\right)}{2\omega \mu _0}}|E_y|^2\left[\alpha ^2+\left(\text{Re}\left\{\beta _{or}\right\}\right)^2\right]\mathrm{\hspace{0.17em}0},`$ (12)
we say that ordinary plane waves have positive phase velocity (PPV) for all directions of propagation.
Let us focus attention on a nondissipative medium (i.e., $`ϵ,ϵ_x`$). From (10) we see that $`\text{Im}\left\{\beta _{or}\right\}0`$ for (i) $`ϵ>0`$ when $`\omega ^2ϵ\mu _0<\alpha ^2`$; and (ii) $`ϵ<0`$. Thus, nonevanescent ordinary plane waves propagate in a nondissipative medium only when $`ϵ>0`$ and $`\omega \sqrt{ϵ\mu _0}<\alpha <\omega \sqrt{ϵ\mu _0}`$. In geometric terms, the wavevector components have a circular representation in $`(\alpha ,\beta _{or})`$ space.
### 2.2 Extraordinary wave
The extraordinary wavevector
$$\underset{¯}{k}_{ex}=\alpha \underset{¯}{\overset{^}{x}}+\beta _{ex}\underset{¯}{\overset{^}{z}},$$
(13)
arises from the dispersion relation (7) with components satisfying
$`\alpha ^2ϵ_x+\beta _{ex}^2ϵ`$ $`=`$ $`\omega ^2ϵϵ_x\mu _0.`$ (14)
In the case where $`\beta _{ex}=0`$ the mathematical description of the extraordinary wave is isomorphic to that for the ordinary wave. Therefore, we exclude this possibility from our consideration in this section. The eigenvector
$$\underset{¯}{E}_0=\left(\underset{¯}{\overset{^}{x}}\frac{ϵ_x\alpha }{ϵ\beta _{ex}}\underset{¯}{\overset{^}{z}}\right)E_x,$$
(15)
arises as a solution to the vector Helmholtz equation (6); the complex–valued magnitude $`E_x`$ is determined by the initial/boundary conditions. The corresponding time–averaged Poynting vector is provided as
$`\underset{¯}{P}(\underset{¯}{r})`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left(2\text{Im}\left\{\beta _{ex}\right\}z\right)}{2\omega \mu _0}}\text{Re}\{\alpha \left(\right|{\displaystyle \frac{ϵ_x}{ϵ\beta _{ex}}}|^2\alpha ^2+{\displaystyle \frac{ϵ_x\beta _{ex}^{}}{ϵ\beta _{ex}}})\underset{¯}{\overset{^}{x}}`$ (16)
$`+(\beta _{ex}^{}+\alpha ^2{\displaystyle \frac{ϵ_x^{}}{ϵ^{}\beta _{ex}^{}}})\underset{¯}{\overset{^}{z}}\}|E_x|^2.`$
Hence, we find
$`\text{Re}\left\{\underset{¯}{k}_{ex}\right\}\text{ }\text{}\text{ }\underset{¯}{P}(\underset{¯}{r})`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left(2\text{Im}\left\{\beta _{ex}\right\}z\right)}{2\omega \mu _0}}[\left(\text{Re}\left\{\beta _{ex}\right\}\right)^2`$ (17)
$`+\alpha ^2\left(\alpha ^2\right|{\displaystyle \frac{ϵ_x}{ϵ\beta _{ex}}}|^2+\text{Re}\left\{{\displaystyle \frac{ϵ_x\beta _{ex}^{}}{ϵ\beta _{ex}}}\right\}+\text{Re}\left\{\beta _{ex}\right\}\text{Re}\left\{{\displaystyle \frac{ϵ_x^{}}{ϵ^{}\beta _{ex}^{}}}\right\})].`$
We analytically explore the nondissipative scenario for nonevanescent and evanescent planewave propagation in Sections 2.3 and 2.4, respectively, whereas both the dissipative and the nondissipative scenarios are treated graphically in Section 2.5.
### 2.3 Nonevanescent propagation
By (14), the inequality
$$\omega ^2ϵ_x\mu _0\alpha ^2\gamma >0$$
(18)
is satisfied for nonevanescent planewave propagation in a nondissipative medium, where $`\gamma `$ is defined in (2). Thus, $`\text{Im}\left\{\beta _{ex}\right\}=0`$. We explore the cases $`\gamma >0`$ and $`\gamma <0`$ separately.
* If $`\gamma >0`$ then we require $`\omega \sqrt{ϵ\mu _0}<\alpha <\omega \sqrt{ϵ\mu _0}`$ in order to comply with (18). This implies that $`ϵ>0`$ and $`ϵ_x>0`$. In geometric terms, the wavevector components have an elliptical representation in $`(\alpha ,\beta _{ex})`$ space.
* If $`\gamma <0`$ then the inequality (18) reduces to $`\omega ^2ϵ\mu _0<\alpha ^2`$. Therefore, we see that nonevanescent propagation arises for (a) $`\alpha >\omega \sqrt{ϵ\mu _0}`$ and $`\alpha <\omega \sqrt{ϵ\mu _0}`$ when $`ϵ>0`$; and (b) $`\mathrm{}<\alpha <\mathrm{}`$ when $`ϵ<0`$. In geometric terms, the wavevector components have a hyperbolic representation in $`(\alpha ,\beta _{ex})`$ space.
For $`\text{Im}\left\{\beta _{ex}\right\}=0`$ and $`ϵ_x,ϵ`$, we find that (17) reduces to
$`\text{Re}\left\{\underset{¯}{k}_{ex}\right\}\text{ }\text{}\text{ }\underset{¯}{P}(\underset{¯}{r})`$ $`=`$ $`{\displaystyle \frac{\omega ^3\mu _0\gamma ^2ϵ_x^2}{2\beta _{ex}^2}}.`$ (19)
Hence, nonevanescent plane waves have PPV regardless of the sign of $`\gamma `$ or $`ϵ_x`$.
### 2.4 Evanescent propagation
We turn to evanescent planewave propagation in a nondissipative medium as characterized by the inequality
$$\omega ^2ϵ_x\mu _0\alpha ^2\gamma <0.$$
(20)
Hence, we have $`\text{Re}\left\{\beta _{ex}\right\}=0`$. As in the previous subsection, we explore the cases $`\gamma >0`$ and $`\gamma <0`$ separately.
* If $`\gamma >0`$ then the situation mirrors that which we described earlier for hyperbolic nonevanescent propagation. That is, evanescent propagation arises for (a) $`\alpha >\omega \sqrt{ϵ\mu _0}`$ and $`\alpha <\omega \sqrt{ϵ\mu _0}`$ when $`ϵ>0`$; and (b) $`\mathrm{}<\alpha <\mathrm{}`$ when $`ϵ<0`$. In geometric terms, the wavevector components have a hyperbolic representation in $`(\alpha ,\text{Im}\left\{\beta _{ex}\right\})`$ space.
* If $`\gamma <0`$ then evanescent propagation arises provided that $`ϵ>0`$, $`ϵ_x<0`$ and $`\omega \sqrt{ϵ\mu _0}<\alpha <\omega \sqrt{ϵ\mu _0}`$. In geometric terms, the wavevector components have an elliptical representation in $`(\alpha ,\text{Im}\left\{\beta _{ex}\right\})`$ space.
For $`\text{Re}\left\{\beta _{ex}\right\}=0`$ and $`ϵ_x,ϵ`$, we find that (17) reduces to
$`\text{Re}\left\{\underset{¯}{k}_{ex}\right\}\text{ }\text{}\text{ }\underset{¯}{P}(\underset{¯}{r})`$ $`=`$ $`{\displaystyle \frac{\omega \alpha ^2ϵ_x\gamma }{2\left(\alpha ^2\gamma \omega ^2\mu _0ϵ_x\right)}}\mathrm{exp}\left(2\text{Im}\left\{\beta _{ex}\right\}z\right).`$ (21)
Hence, evanescent plane waves have PPV if (a) $`\gamma <0`$ or (b) $`\gamma >0`$ and $`ϵ_x>0`$. However, negative phase velocity (NPV) propagation arises if $`\gamma >0`$ and $`ϵ_x<0`$.
### 2.5 Illustrative examples
Let us illustrate the geometric aspect of the dispersion relations with some representative numerical examples.
First, suppose we consider the case $`\gamma >0`$ with $`ϵ=2ϵ_0`$ and $`ϵ_x=6ϵ_0`$, where $`ϵ_0`$ is the free–space permittivity. In Figure 1 the real and imaginary parts of $`\beta _{ex}`$ are plotted against $`\alpha `$. The elliptical nonevanescent nature of the dispersion relation is clear for $`\omega \sqrt{2ϵ_0\mu _0}<\alpha <\omega \sqrt{2ϵ_0\mu _0}`$, while the hyperbolic evanescent nature is apparent for $`\alpha <\omega \sqrt{2ϵ_0\mu _0}`$ and $`\alpha >\omega \sqrt{2ϵ_0\mu _0}`$. The elliptical/hyperbolic geometric interpretation breaks down when dissipative mediums are considered. However, the corresponding dispersion relations are geometrically reminiscent of their nondissipative counterparts. This can be observed in Figure 2 in which the graphs corresponding to Figure 1 are displayed for $`ϵ=\left(2+i0.5\right)ϵ_0`$ and $`ϵ_x=\left(6+i0.75\right)ϵ_0`$.
Second, we turn to the case $`\gamma <0`$ with $`ϵ=2ϵ_0`$ and $`ϵ_x=6ϵ_0`$. The real and imaginary parts of $`\beta _{ex}`$ are graphed against $`\alpha `$ in Figure 3. The graphs mirror those of Figure 1 but with nonevanescent and evanscent aspects interchanged; i.e., we observe hyperbolic nonevanescent characteristics for $`\alpha <\omega \sqrt{2ϵ_0\mu _0}`$ and $`\alpha >\omega \sqrt{2ϵ_0\mu _0}`$, and elliptical evanescent characteristics for $`\omega \sqrt{2ϵ_0\mu _0}<\alpha <\omega \sqrt{2ϵ_0\mu _0}`$. The corresponding graphs for $`ϵ=\left(2+i0.5\right)ϵ_0`$ and $`ϵ_x=\left(6+i0.75\right)ϵ_0`$ are presented in Figure 4. Notice that the shapes of the graphs in Figures 4 and 2 are similar but not identical.
## 3 Numerical conceptualization
Although uniaxial dielectric mediums with $`\gamma <0`$ do not occur in nature (to the best of the authors’ knowledge), they can be conceptualized as metamaterials by means of homogenization.
For example, let us consider the homogenization of a composite comprising two component materials phases, labelled as $`a`$ and $`b`$. Both component material phases are taken to be isotropic dielectric mediums: $`ϵ^a`$ and $`ϵ^b`$ denote the permittivity scalars of phases $`a`$ and $`b`$, respectively. The component material phases are envisioned as random distributions of identically–oriented, spheroidal particles. The spheroidal shape — which is taken to be the same for all spheroids in component material phase $`a`$ and $`b`$ — is parameterized via the shape dyadic $`\underset{¯}{\underset{¯}{U}}=\text{diag}(U_x,U,U)`$. That is, we take the spheroid’s principal axis to lie along the $`x`$ axis. The spheroid’s surface is prescribed by the vector
$$\underset{¯}{r}_s(\theta ,\varphi )=\eta \underset{¯}{\underset{¯}{U}}\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}(\theta ,\varphi ),$$
(22)
with $`\underset{¯}{\overset{^}{r}}`$ being the radial unit vector specified by the spherical polar coordinates $`\theta `$ and $`\varphi `$. The linear dimensions of the spheroid, as determined by the parameter $`\eta `$, are assumed to be small relative to the electromagnetic wavelength(s).
The permittivity dyadic of the resulting homogenized composite medium (HCM)
$$\underset{¯}{\underset{¯}{ϵ}}^{HCM}=\text{diag}(ϵ_x^{HCM},ϵ^{HCM},ϵ^{HCM}),$$
(23)
as estimated using the Bruggeman homogenization formalism, is provided implicitly via
$$f_a\underset{¯}{\underset{¯}{a}}^a+f_b\underset{¯}{\underset{¯}{a}}^b=\underset{¯}{\underset{¯}{0}},$$
(24)
where $`f_a`$ and $`f_b=1f_a`$ denote the respective volume fractions of the material component phases $`a`$ and $`b`$. The polarizability dyadics in (24) are defined as
$$\underset{¯}{\underset{¯}{a}}^{\mathrm{}}=\left(ϵ^{\mathrm{}}\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\text{ }\text{}\text{ }\left[\underset{¯}{\underset{¯}{I}}+i\omega \underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^{\mathrm{}}\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\right]^1,(\mathrm{}=a,b),$$
(25)
wherein the depolarization dyadic is given by the surface integral
$$\underset{¯}{\underset{¯}{D}}=\frac{1}{i\omega 4\pi }_0^{2\pi }𝑑\varphi _0^\pi 𝑑\theta \mathrm{sin}\theta \left(\frac{1}{\underset{¯}{\overset{^}{r}}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{HCM}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}}\right)\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}\underset{¯}{\overset{^}{r}}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1.$$
(26)
Closed–form expressions for the depolarization dyadic for uniaxial mediums are available in terms of hyperbolic functions . However, we note that these exact results are not valid for nondissipative mediums with $`\gamma <0`$, and numerical evaluation of $`\underset{¯}{\underset{¯}{D}}`$ has to be resorted to.
The Jacobi iteration scheme
$$\underset{¯}{\underset{¯}{ϵ}}^{HCM}\left[p\right]=𝒯\left\{\underset{¯}{\underset{¯}{ϵ}}^{HCM}\left[p1\right]\right\},\left(p=1,2,3,\mathrm{}\right),$$
(27)
where the operator $`𝒯`$ is defined via
$`𝒯\left\{\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right\}`$ $`=`$ $`\left\{f_aϵ^a\left[\underset{¯}{\underset{¯}{I}}+i\omega \underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^a\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\right]^1+f_bϵ^b\left[\underset{¯}{\underset{¯}{I}}+i\omega \underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^b\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\right]^1\right\}`$
$`\text{}\text{ }\left\{f_a\left[\underset{¯}{\underset{¯}{I}}+i\omega \underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^a\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\right]^1+f_b\left[\underset{¯}{\underset{¯}{I}}+i\omega \underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^b\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{HCM}\right)\right]^1\right\}^1,`$
may be employed to solve (24) for $`\underset{¯}{\underset{¯}{ϵ}}^{HCM}`$. Suitable initial values for the iterative scheme are provided by
$$\underset{¯}{\underset{¯}{ϵ}}^{HCM}\left[\mathrm{\hspace{0.17em}0}\right]=\left(f_aϵ^a+f_bϵ^b\right)\underset{¯}{\underset{¯}{I}}.$$
(29)
For further details on the Bruggeman homogenization formalism the reader is referred to and to references therein.
Let us consider the homogenization scenario wherein material component phase $`a`$ is taken to be iron at 670 nm free–space wavelength. Correspondingly, we take $`ϵ^a=\left(4.34+i20.5\right)ϵ_0`$. The material component phase $`b`$ is assumed to be free space; i.e., $`ϵ^b=ϵ_0`$. The shape of the component spheroids is specified by $`U_x/U=12`$. The Bruggeman estimates of the HCM permittivity parameters $`ϵ^{HCM}`$ and $`ϵ_x^{HCM}`$ are plotted as functions of volume fraction $`f_a`$ in Figure 5. At intermediate values of $`f_a`$ we see that $`\gamma <0`$ for a substantial range of $`f_a`$ values.
Extensive accounts of similar numerical homogenizations, based on the Bruggeman formalism and more general approaches, can be found elsewhere .
## 4 Concluding remarks
The dispersion relations for uniaxial dielectric mediums have been characterized with respect to the parameter $`\gamma `$ (2). For $`\gamma <0`$, the dispersion relations are hyperbolic for nondissipative mediums and hyperbolic–like for dissipative mediums. Similarly, the dispersion relations are elliptical for nondissipative mediums and elliptical–like for dissipative mediums with $`\gamma >0`$. Through the homogenization of isotropic component material phases based on spheroidal topology, we demonstrate that metamaterials with $`\gamma <0`$ may be straightforwardly conceptualized. Thus, a practical means of achieving the exotic electromagnetic properties associated with hyperbolic and hyperbolic–like uniaxial mediums is presented.
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# Transport coefficients for an inelastic gas around uniform shear flow: Linear stability analysis
## I Introduction
The understanding of granular systems still remains a topic of interest and controversy. Under rapid flow conditions, they can be modeled as a fluid of hard spheres dissipating part of their kinetic energy during collisions. In the simplest model, the grains are taken to be smooth so that the inelasticity of collisions is characterized through a constant coefficient of normal restitution $`\alpha 1`$. Energy dissipation has profound consequences on the behavior of these systems since they exhibit a rich phenomenology with many qualitative differences with respect to molecular systems. In particular, the absence of energy conservation yields subtle modifications of the conventional Navier-Stokes equations for states with small gradients of the hydrodynamic fields. The dependence of the corresponding transport coefficients on dissipation may be determined from the Boltzmann kinetic equation conveniently modified to account for inelastic binary collisions GS95 ; BDS97 . The idea is to extend the Chapman-Enskog method CC70 to the inelastic case by expanding the velocity distribution function around the local version of the homogeneous cooling state, namely, a homogeneous state whose dependence on time occurs only through the temperature. In the first order of the expansion, explicit expressions for the transport coefficients as functions of the coefficient of restitution have been obtained in the case of a single gas single as well as for granular mixtures mixture , showing good agreement the analytical results with those obtained from Monte Carlo simulations DSMC .
Although the Chapman-Enskog method can be in principle applied to get higher orders in the gradients (Burnett and super-Burnett corrections,$`\mathrm{}`$), it is extremely difficult to evaluate those terms especially for inelastic systems. In addition, questions about its convergence remains still open GS03 . This gives rise to the search for alternative approaches to characterize transport for strongly inhomogeneous situations (i.e., beyond the Navier-Stokes limit). One possibility is to expand in small gradients around a more relevant reference state than the (local) homogenous cooling state. For example, consider states near a shearing reference steady state such as the so-called uniform (simple) shear flow (USF) C90 . Such an application of the Chapman-Enskog method to a nonequilibrium state requires some care as recently discussed in Ref. L05 . The USF state is probably the simplest flow problem since the only nonzero hydrodynamic gradient is $`u_x/ya=\text{const}`$, where $`𝐮`$ is the flow velocity and $`a`$ is the constant shear rate. Due to its simplicity, this state has been widely used in the past both for elastic GS03 and inelastic gases C90 to shed light on the complexities associated with the nonlinear response of the system to the action of strong shearing. However, the nature of this state for granular systems is different from that of the elastic fluids since the source of energy due to the macroscopic imposed shear field drives the granular system into rapid flow and a steady state is achieved when the amount of energy supplied by shearing work is balanced by the lost one due to the inelastic collisions between the particles. As a consequence, in the steady state the reduced shear rate $`a^{}a/\sqrt{T}`$ (which is the relevant nonequilibrium parameter of the problem) is not an independent quantity but becomes a function of the coefficient of restitution $`\alpha `$. This means that the quasielastic limit ($`\alpha 1`$) naturally implies the limit of small shear rates ($`a^{}1`$) and vice versa. The study of the rheological properties of the USF state has received a great deal of attention in recent years in the case of monocomponent LSJC84 ; JR88 ; C89 ; HS92 ; SK94 ; LB94 ; GT96 ; SGN96 ; BRM97 ; CR98 ; MGSB99 and multicomponent systems CH02 ; MG02 ; G02 ; L04 ; AL03 ; MGAL05 .
The aim of this paper is to determine the heat and momentum fluxes of a gas of inelastic hard spheres under simple shear flow in the framework of the Boltzmann equation. The physical situation is such that the gas is in a state that deviates from the simple shear flow by small spatial gradients. The starting point of this study is a recent approximate solution of the Boltzmann equation which is based on Grad’s method MG02 ; G02 ; SGD04 . In spite of this approach, the relevant transport properties obtained from this solution compare quite well with Monte Carlo simulations even for strong dissipation BRM97 ; MG02 ; AS05 , showing again the reliability of Grad’s approximation to compute the lowest velocity moments of the velocity distribution function. Since the system is slightly perturbed form the USF, the Boltzmann equation is solved by applying the Chapman-Enskog method around the (local) shear flow state rather than the (local) homogeneous cooling state. This is the main feature of this expansion since the reference state is not restricted to small values of the shear rate. One important point is that, for general small deviations from the shear flow state, the zeroth-order distribution is not a stationary distribution since the collisional cooling cannot be compensated locally for viscous heating. This fact gives rise to new conceptual and practical difficulties not present in the previous analysis made for elastic gases to describe transport in thermostatted shear flow states LD97 . Due to the difficulties involved in this expansion, here general results will be restricted to particular perturbations for which steady state conditions apply. In the first order of the expansion, the generalized transport coefficients are given in terms of the solutions of linear integral equations. To get explicit expressions for these coefficients, one needs to know the fourth-degree moments of USF. This requires to consider higher-order terms in Grad’s approximation for the reference distribution function, which is quite an intricate problem. In order to overcome such difficulty, here I have used a convenient kinetic model BDS99 that preserves the essential properties of the inelastic Boltzmann equation but admits more practical analysis. The mathematical and physical basis for this model as a good representation of the Boltzmann equation has been discussed in Ref. BDS99 . In particular, it is worth noting that the results derived from this model coincides with those given from the Boltzmann equation at the level of the rheological properties BRM97 ; BDS99 . Furthermore, recent computer simulation results AS05 have also shown good agreement between the kinetic model and the Boltzmann equation for the fourth-degree moments, covering this agreement a wide range of values of dissipation (say, for instance, $`\alpha 0.5`$). This good agreement extends that previously demonstrated for Couette flow in dilute gases TTMGSD01 and for USF in dense systems MGSB99 and shows the reliability of the kinetic model to capture the main trends of the Boltzmann equation, especially those related to transport properties.
The knowledge of the above generalized transport coefficients allows one to determine the hydrodynamic modes from the associated linearized hydrodynamic equations. This is quite an interesting problem widely analyzed in the literature. As noted by the different molecular dynamics experiments carried out for the USF problem GT96 ; AL03 ; MD , it becomes apparent the development of inhomogeneities and formation of clusters as the flow progresses. Consequently, the USF state is unstable for long enough wavelength spatial perturbations. In order to understand this phenomenon, several stability analysis have been undertaken S92 ; B93 ; AN97 ; K00 ; K01 . Most of them are based on the Navier-Stokes equations S92 ; B93 ; AN97 and, therefore, they are limited to small velocity gradients, which for the USF problem means small dissipation. Another alternative has been to solve the Boltzmann equation by means of an expansion in a set of basis functions K00 ; K01 . The coefficients of this expansion are then determined by using also an expansion in powers of the parameter $`ϵ\sqrt{1\alpha ^2}`$, which is assumed to be small. All these analytical results have shown that the USF becomes unstable for certain kind of disturbances. My approach is different from previous works since the conditions for stability are obtained from a linear stability analysis involving the transport coefficients of the perturbed USF state instead of the usual Navier-Stokes coefficients. Furthermore, the analysis is not restricted to the low-dissipation limit since the reference state goes beyond this range of values of $`\alpha `$. Two different perturbations to the reference state have been considered here: (i) perturbations along the velocity gradient ($`y`$ direction) only and (ii) perturbations along the vorticity direction ($`z`$ direction) only. The results show that the USF is linearly stable in the first case while it becomes unstable in the second case. These results agree qualitatively with those previously derived S92 ; AN97 in the context of the Navier-Stokes description. On the other hand, at a quantitative level, the comparison carried out here shows significant differences between the Navier-Stokes description and the present results as the collisions become more inelastic. In addition, our results also confirm that the instability is confined to long wavelengths (small wave numbers) and so it can be avoided for small enough systems.
The plan of the paper is as follows. In Sec. II, the Boltzmann kinetic equation is introduced and a brief summary of relevant results concerning the USF problem is given. In Sec. III, the problem we are interested in is described and the set of generalized transport coefficients characterizing the transport around USF is defined. Explicit expressions for these coefficients are provided in Sec. IV by using a kinetic model of the Boltzmann equation. The details of the calculations are displayed along several Appendices. Section V is devoted to the linear stability analysis around the steady USF state and presents the form of the hydrodynamic modes. The paper is closed in Sec. VI with a discussion of the results obtained here.
## II Boltzmann kinetic equation and uniform shear flow
Let us consider a granular gas composed by smooth spheres of mass $`m`$ and diameter $`\sigma `$. The inelasticity of collisions among all pairs is accounted for by a constant coefficient of restitution $`0\alpha 1`$ that only affects the translational degrees of freedom of grains. In a kinetic theory description all the relevant information on the state of the system is given by the one-particle velocity distribution function $`f(𝐫,𝐯,t)`$. At low density the inelastic Boltzmann equation GS95 ; BDS97 gives the time evolution of $`f(𝐫,𝐯,t)`$. In the absence of an external force, it has the form
$$\left(\frac{}{t}+𝐯\right)f(𝐫,𝐯,t)=J[𝐯|f(t),f(t)],$$
(1)
where the Boltzmann collision operator is
$`J\left[𝐯_1|f,f\right]`$ $`=`$ $`\sigma ^2{\displaystyle 𝑑𝐯_2𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝝈}𝐠)(\widehat{𝝈}𝐠)}`$ (2)
$`\times \left[\alpha ^2f(𝐫,𝐯_1^{})f(𝐫,𝐯_2^{},t)f(𝐫,𝐯_1,t)f(𝐫,𝐯_2,t)\right].`$
Here, $`\widehat{𝝈}`$ is a unit vector along their line of centers, $`\mathrm{\Theta }`$ is the Heaviside step function, and $`𝐠=𝐯_1𝐯_2`$ is the relative velocity. The primes on the velocities denote the initial values $`\{𝐯_1^{},𝐯_2^{}\}`$ that lead to $`\{𝐯_1,𝐯_2\}`$ following a binary collision:
$$𝐯_1^{}=𝐯_1\frac{1}{2}\left(1+\alpha ^1\right)(\widehat{𝝈}𝐠)\widehat{𝝈},𝐯_2^{}=𝐯_2+\frac{1}{2}\left(1+\alpha ^1\right)(\widehat{𝝈}𝐠)\widehat{𝝈}$$
(3)
The first five velocity moments of $`f`$ define the number density
$$n(𝐫,t)=𝑑𝐯f(𝐫,𝐯,t),$$
(4)
the flow velocity
$$𝐮(𝐫,t)=\frac{1}{n(𝐫,t)}𝑑\mathrm{𝐯𝐯}f(𝐫,𝐯,t),$$
(5)
and the granular temperature
$$T(𝐫,t)=\frac{m}{3n(𝐫,t)}𝑑𝐯V^2(𝐫,t)f(𝐫,𝐯,t),$$
(6)
where $`𝐕(𝐫,t)𝐯𝐮(𝐫,t)`$ is the peculiar velocity. The macroscopic balance equations for density $`n`$, momentum $`m𝐮`$, and energy $`\frac{3}{2}nT`$ follow directly from Eq. (1) by multiplying with $`1`$, $`m𝐯`$, and $`\frac{1}{2}mv^2`$ and integrating over $`𝐯`$:
$$D_tn+n𝐮=0,$$
(7)
$$D_tu_i+(mn)^1_jP_{ij}=0,$$
(8)
$$D_tT+\frac{2}{3n}\left(𝐪+P_{ij}_ju_i\right)=\zeta T,$$
(9)
where $`D_t=_t+𝐮`$. The microscopic expressions for the pressure tensor $`𝖯`$, the heat flux $`𝐪`$, and the cooling rate $`\zeta `$ are given, respectively, by
$$𝖯(𝐫,t)=𝑑𝐯m\mathrm{𝐕𝐕}f(𝐫,𝐯,t),$$
(10)
$$𝐪(𝐫,t)=𝑑𝐯\frac{1}{2}mV^2𝐕f(𝐫,𝐯,t),$$
(11)
$$\zeta (𝐫,t)=\frac{1}{3n(𝐫,t)T(𝐫,t)}𝑑𝐯mV^2J[𝐫,𝐯|f(t)].$$
(12)
We assume that the gas is under uniform (or simple) shear flow (USF). This idealized macroscopic state is characterized by a constant density, a uniform temperature and a simple shear with the local velocity field given by
$$u_i=a_{ij}r_j,a_{ij}=a\delta _{ix}\delta _{jy},$$
(13)
where $`a`$ is the constant shear rate. This linear velocity profile assumes no boundary layer near the walls and is generated by the Lee-Edwards boundary conditions LE72 , which are simply periodic boundary conditions in the local Lagrangian frame moving with the flow velocity. For elastic gases, the temperature grows in time due to viscous heating and so a steady state is not possible unless an external (artificial) force is introduced GS03 . However, for inelastic gases, the temperature changes in time due to the competition between two (opposite) mechanisms: on the one hand, viscous (shear) heating and, on the other hand, energy dissipation in collisions. A steady state is achieved when both mechanisms cancel each other and the fluid autonomously seeks the temperature at which the above balance occurs. Under these conditions, in the steady state the balance equation (9) becomes
$$aP_{xy}=\frac{3}{2}\zeta p,$$
(14)
where $`p=nT`$ is the hydrostatic pressure. Note that for given values of the shear rate $`a`$ and the coefficient of restitution $`\alpha `$, the relation (14) gives the temperature $`T`$ in the steady state as a unique function of the density $`n`$.
The USF problem is perhaps the nonequilibrium state most widely studied in the past few years both for granular and conventional gases GS03 ; C90 . At a microscopic level, it becomes spatially homogeneous when the velocities of the particles are referred to the Lagrangian frame of reference co-moving with the flow velocity $`𝐮`$ DSBR86 . Therefore, the one-particle distribution function adopts the uniform form, $`f(𝐫,𝐯)f(𝐕)`$, and the Boltzmann equation (1) reads
$$aV_y\frac{}{V_x}f(𝐕)=J\left[𝐕|f,f\right].$$
(15)
This equation is invariant under the transformations
$$V_zV_z,(V_x,V_y)(V_x,V_y),(V_x,a)(V_x,a).$$
(16)
The elements of the pressure tensor provide information on the relevant transport properties of the USF problem. These elements can be obtained by multiplying the Boltzmann equation (15) by $`mV_iV_j`$ and integrating over $`𝐕`$. The result is
$`a_i\mathrm{}P_j\mathrm{}+a_j\mathrm{}P_i\mathrm{}`$ $`=`$ $`m{\displaystyle 𝑑𝐕V_iV_jJ[𝐕|f,f]}`$ (17)
$``$ $`\mathrm{\Lambda }_{ij}.`$
The exact expression of the collision integral $`\mathrm{\Lambda }_{ij}`$ is not known, even in the elastic case. However, a good estimate can be expected by using Grad’s approximation:
$$f(𝐕)f_0(𝐕)\left[1+\frac{m}{2T}\left(\frac{P_{ij}}{p}\delta _{ij}\right)V_iV_j\right],$$
(18)
where $`f_0(𝐕)`$
$$f_0(𝐕)=n(m/2\pi T)^{3/2}\mathrm{exp}(mV^2/2T)$$
(19)
is the local equilibrium distribution function. When Eq. (18) is substituted into the definition of $`\mathrm{\Lambda }_{ij}`$ and nonlinear terms in $`P_{ij}/nT\delta _{ij}`$ are neglected, one gets G02
$$\mathrm{\Lambda }_{ij}=\nu \left[\beta \left(P_{ij}p\delta _{ij}\right)+\zeta ^{}P_{ij}\right],$$
(20)
where
$$\nu (T)=\frac{16}{5}n\sigma ^2\sqrt{\frac{\pi T}{m}},$$
(21)
is an effective collision frequency,
$$\zeta ^{}=\frac{\zeta }{\nu }=\frac{5}{12}(1\alpha ^2),$$
(22)
is the dimensionless cooling rate evaluated in the local equilibrium approximation and
$$\beta =\frac{1+\alpha }{2}\left(1\frac{1\alpha }{3}\right).$$
(23)
The set of coupled equations for $`P_{ij}`$ can be now easily solved when one takes into account the approach (20). The expressions for the reduced elements $`P_{ij}^{}=P_{ij}/p`$ are
$$P_{xx}^{}=32P_{yy}^{},P_{yy}^{}=P_{zz}^{}=\frac{\beta }{\beta +\zeta ^{}},P_{xy}^{}=\frac{\beta }{(\beta +\zeta ^{})^2}a^{},$$
(24)
where the (reduced) shear rate $`a^{}=a/\nu `$ is given by
$$a^{}=\sqrt{\frac{3}{2}\frac{\zeta ^{}}{\beta }}\left(\beta +\zeta ^{}\right).$$
(25)
The expression (25) clearly indicates the intrinsic connection between the (reduced) velocity gradient and dissipation in the system. In fact, in the elastic limit ($`\alpha =1`$, which implies $`a^{}=0`$), the equilibrium results of the ordinary gas are recovered, i.e., $`P_{ij}^{}=\delta _{ij}`$. This means that $`\alpha `$ (or $`a^{}`$) can be considered as the relevant nonequilibrium parameter of the problem. The analytical results given by Eqs. (24) and (25) agree quite well SGD04 ; MG02 with Monte Carlo simulations of the Boltzmann equation MG02 ; AS05 , even for strong dissipation.
## III Small perturbations from the uniform shear flow: Transport coefficients
In general, the USF state can be disturbed by small spatial perturbations. The response of the system to these perturbations gives rise to additional contributions to the momentum and heat fluxes, which can be characterized by generalized transport coefficients. This section is devoted to the study of such small perturbations.
In order to analyze this problem we have to start from the Boltzmann equation with a general time and space dependence. Let $`𝐮_0=𝖺𝐫`$ be the flow velocity of the undisturbed USF state. Here, the only nonzero element of the tensor $`𝖺`$ is $`a_{ij}=a\delta _{ix}\delta _{jy}`$. In the disturbed state, however the true velocity $`𝐮`$ is in general different from $`𝐮_0`$ since $`𝐮=𝐮_0+\delta 𝐮`$, $`\delta 𝐮`$ being a small perturbation to $`𝐮_0`$. As a consequence, the true peculiar velocity is now $`𝐜𝐯𝐮=𝐕\delta 𝐮`$, where $`𝐕=𝐯𝐮_0`$. In the Lagrangian frame moving with $`𝐮_0`$, the Boltzmann equation can be written as
$$\frac{}{t}faV_y\frac{}{V_x}f+\left(𝐕+𝐮_0\right)f=J[𝐕|f,f],$$
(26)
where here the derivative $`f`$ is taken at constant $`𝐕`$. The corresponding macroscopic balance equations associated with this disturbed USF state follows from the general equations (7)–(9) when one takes into account that $`𝐮=𝐮_0+\delta 𝐮`$. The result is
$$_tn+𝐮_0n=(n\delta 𝐮),$$
(27)
$$_t\delta 𝐮+𝖺\delta 𝐮+(𝐮_0+\delta 𝐮)\delta 𝐮=(mn)^1𝖯,$$
(28)
$$\frac{3}{2}n_tT+\frac{3}{2}n(𝐮_0+\delta 𝐮)T+aP_{xy}+𝐪+𝖯:\delta 𝐮=\frac{3}{2}p\zeta ,$$
(29)
where the pressure tensor $`𝖯`$, the heat flux $`𝐪`$ and the cooling rate $`\zeta `$ are defined by Eqs. (10)–(12), respectively, with the replacement $`𝐕𝐜`$.
We assume now that the deviations from the USF state are small, which means that the spatial gradients of the hydrodynamic fields
$$A(𝐫,t)\{n(𝐫,t),T(𝐫,t),\delta 𝐮(𝐫,t)\}$$
(30)
are small. Under these conditions, a solution to the Boltzmann equation (26) can be obtained by means of a generalization of the conventional Chapman-Enskog method CC70 where the velocity distribution function is expanded about a local shear flow reference state in terms of the small spatial gradients of the hydrodynamic fields relative to those of USF. This type of Chapman-Enskog-like expansion has been considered in the case of elastic gases to get the set of shear-rate dependent transport coefficients GS03 ; LD97 in a thermostatted shear flow problem and it has also been recently considered L05 in the context of inelastic gases.
To construct the Chapman-Enskog expansion let us look for a normal solution of the form
$$f(𝐫,𝐕,t)f(A(𝐫,t),𝐕).$$
(31)
This special solution expresses the fact that the space dependence of the reference shear flow is completely absorbed in the relative velocity $`𝐕`$ and all other space and time dependence occurs entirely through a functional dependence on the fields $`A(𝐫,t)`$. The functional dependence can be made local by an expansion of the distribution function in powers of the hydrodynamic gradients:
$$f(𝐫,𝐕,t)=f^{(0)}(A(𝐫,t),𝐕)+f^{(1)}(A(𝐫,t),𝐕)+\mathrm{},$$
(32)
where the reference zeroth-order distribution function corresponds to the USF distribution function but taking into account the local dependence of the density and temperature and the change $`𝐕𝐕\delta 𝐮(𝐫,t)`$ \[see Eqs. (138) and (139) for the explicit form of $`f^{(0)}`$ in the steady state given by a kinetic model of the Boltzmann equation\]. The successive approximations $`f^{(k)}`$ are of order $`k`$ in the gradients of $`n`$, $`T`$, and $`\delta 𝐮`$ but retain all the orders in the shear rate $`a`$. This is the main feature of this expansion. In this paper, only the first order approximation will be considered. More details on this Chapman-Enskog-like type of expansion can be found in Ref. L05 .
The expansion (32) yields the corresponding expansion for the fluxes and the cooling rate when one substitutes (32) into their definitions (10)–(12):
$$𝖯=𝖯^{(0)}+𝖯^{(1)}+\mathrm{},𝐪=𝐪^{(0)}+𝐪^{(1)}+\mathrm{},\zeta =\zeta ^{(0)}+\zeta ^{(1)}+\mathrm{}.$$
(33)
Finally, as in the usual Chapman-Enskog method, the time derivative is also expanded as
$$_t=_t^{(0)}+_t^{(1)}+_t^{(2)}+\mathrm{},$$
(34)
where the action of each operator $`_t^{(k)}`$ is obtained from the hydrodynamic equations (27)–(29). These results provide the basis for generating the Chapman-Enskog solution to the inelastic Boltzmann equation (26).
### III.1 Zeroth-order approximation
Substituting the expansions (32) and (34) into Eq. (26), the kinetic equation for $`f^{(0)}`$ is given by
$$_t^{(0)}f^{(0)}aV_y\frac{}{V_x}f^{(0)}=J[𝐕|f^{(0)},f^{(0}].$$
(35)
To lowest order in the expansion the conservation laws give
$$_t^{(0)}n=0,_t^{(0)}T=\frac{2}{3n}aP_{xy}^{(0)}T\zeta ^{(0)},$$
(36)
$$_t^{(0)}\delta u_i+a_{ij}\delta u_j=0.$$
(37)
As said before, for given values of $`a`$ and $`\alpha `$, the steady state condition (14) establishes a mapping between the density and temperature so that every density corresponds to one and only one temperature. Since the density $`n(𝐫,t)`$ and temperature $`T(𝐫,t)`$ are specified separately in the local USF state, the viscous heating only partially compensates for the collisional cooling and so, $`_t^{(0)}T0`$. Consequently, the zeroth-order distribution $`f^{(0)}`$ depends on time through its dependence on the temperature. Because of the steady state condition (14) does not apply in general locally, the reduced shear rate $`a^{}=a/\nu (n,T)`$ depends on space and time so that, $`a^{}`$ and $`\alpha `$ must be considered as independent parameters for general infinitesimal perturbations around the USF state. Since $`f^{(0)}`$ is a normal solution, then
$`_t^{(0)}f^{(0)}`$ $`=`$ $`{\displaystyle \frac{f^{(0)}}{n}}_t^{(0)}n+{\displaystyle \frac{f^{(0)}}{T}}_t^{(0)}T+{\displaystyle \frac{f^{(0)}}{\delta u_i}}_t^{(0)}\delta u_i`$ (38)
$`=`$ $`\left({\displaystyle \frac{2}{3n}}aP_{xy}^{(0)}+T\zeta ^{(0)}\right){\displaystyle \frac{}{T}}f^{(0)}a_{ij}\delta u_j{\displaystyle \frac{}{\delta u_i}}f^{(0)}`$
$`=`$ $`\left({\displaystyle \frac{2}{3n}}aP_{xy}^{(0)}+T\zeta ^{(0)}\right){\displaystyle \frac{}{T}}f^{(0)}+a_{ij}\delta u_j{\displaystyle \frac{}{c_i}}f^{(0)},`$
where in the last step we have taken into account that $`f^{{}_{}{}^{(}0)}`$ depends on $`\delta 𝐮`$ only through the peculiar velocity $`𝐜`$. Substituting (38) into (35) yields the following kinetic equation for $`f^{(0)}`$:
$$\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)\frac{}{T}f^{(0)}ac_y\frac{}{c_x}f^{(0)}=J[𝐕|f^{(0)},f^{(0}].$$
(39)
The zeroth-order solution leads to $`𝐪^{(0)}=\mathrm{𝟎}`$. On the other hand, to solve Eq. (39) one needs to know the temperature dependence of the zeroth momentum flux $`P_{xy}^{(0)}`$. A closed set of equations for $`𝖯^{(0)}`$ is obtained when one considers Grad’s approximation (18):
$$\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)\frac{}{T}P_{ij}^{(0)}+a_i\mathrm{}P_j\mathrm{}^{(0)}+a_j\mathrm{}P_i\mathrm{}^{(0)}=\nu \left[\beta \left(P_{ij}^{(0)}p\delta _{ij}\right)+\zeta ^{}P_{ij}^{(0)}\right],$$
(40)
where
$$\zeta ^{}=\frac{\zeta ^{(0)}}{\nu }=\frac{5}{12}(1\alpha ^2).$$
(41)
The steady state solution of Eq. (40) is given by Eqs. (24) and (25). However, in general the equations (40) must be solved numerically to get the dependence of the zeroth-order pressure tensor $`P_{ij}^{(0)}(T)`$ on temperature. A detailed study on the unsteady hydrodynamic solution of Eqs. (40) has been carried out in Ref. SGD04 . In what follows, $`P_{ij}^{(0)}(T)`$ will be considered as a known function of $`T`$.
### III.2 First-order approximation
The analysis to first order in the gradients is worked out in Appendix A. Only the final results are presented in this Section. The distribution function $`f^{(1)}`$ is of the form
$$f^{(1)}=𝐗_nn+𝐗_TT+𝖷_u:\delta 𝐮,$$
(42)
where the vectors $`𝐗_n`$ and $`𝐗_T`$ and the tensor $`𝖷_u`$ are functions of the true peculiar velocity $`𝐜`$. They are the solutions of the following linear integral equations:
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+ac_y\frac{}{c_x}\right]X_{n,i}+\frac{T}{n}\left[\frac{2a}{3p}(1n_n)P_{xy}^{(0)}\zeta ^{(0)}\right]X_{T,i}=Y_{n,i},$$
(43)
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+\frac{2a}{3p}T(_TP_{xy}^{(0)})+\frac{3}{2}\zeta ^{(0)}+ac_y\frac{}{c_x}\right]X_{T,i}=Y_{T,i},$$
(44)
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+ac_y\frac{}{c_x}\right]X_{u,k\mathrm{}}a\delta _{ky}X_{u,x\mathrm{}}\zeta _{u,k\mathrm{}}T_Tf^{(0)}=Y_{u,k\mathrm{}},$$
(45)
where $`𝐘_n(𝐜)`$, $`𝐘_T(𝐜)`$, and $`𝖸_u(𝐜)`$ are defined by Eqs. (91)–(93), respectively, and $`\zeta _{u,k\mathrm{}}`$ is defined by Eq. (96). While the $`Y`$ functions are given in terms of the reference state distribution $`f^{(0)}`$, $`\zeta _{u,k\mathrm{}}`$ is a functional of the unknown $`X_{u,k\mathrm{}}`$. In addition, $``$ is the linearized Boltzmann collision operator around the reference state
$$X\left(J[f^{(0)},X]+J[X,f^{(0)}]\right).$$
(46)
A good estimate of $`\zeta _{u,k\mathrm{}}`$ can be obtained by expanding $`X_{u,k\mathrm{}}`$ in a complete set of polynomials (for instance, Sonine polynomials) and then truncating the series after the first few terms. In practice, the leading term in these expansions provides a very accurate result over a wide range of dissipation. This contribution has been obtained in Appendix B and is given by Eq. (113).
With the distribution function $`f^{(1)}`$ determined by (42), the first-order corrections to the fluxes are
$$P_{ij}^{(1)}=\eta _{ijk\mathrm{}}\frac{\delta u_k}{r_{\mathrm{}}},$$
(47)
$$q_i^{(1)}=\kappa _{ij}\frac{T}{r_j}\mu _{ij}\frac{n}{r_j},$$
(48)
where
$$\eta _{ijk\mathrm{}}=𝑑𝐜mc_ic_jX_{u,k\mathrm{}}(𝐜),$$
(49)
$$\kappa _{ij}=𝑑𝐜\frac{m}{2}c^2c_iX_{T,j}(𝐜),$$
(50)
$$\mu _{ij}=𝑑𝐜\frac{m}{2}c^2c_iX_{n,j}(𝐜).$$
(51)
Upon writing Eqs. (47)–(51) use has been made of the symmetry properties of $`X_{n,i}`$ $`X_{T,i}`$ and $`X_{u,ij}`$. In general, the set of generalized transport coefficients $`\eta _{ijk\mathrm{}}`$, $`\kappa _{ij}`$, and $`\mu _{ij}`$ are nonlinear functions of the coefficient of restitution $`\alpha `$ and the reduced shear rate $`a^{}`$. The anisotropy induced in the system by the shear flow gives rise to new transport coefficients, reflecting broken symmetry. The momentum flux is expressed in terms of a viscosity tensor $`\eta _{ijk\mathrm{}}(\alpha )`$ of rank 4 which is symmetric and traceless in $`ij`$ due to the properties of the pressure tensor $`P_{ij}^{(1)}`$. The heat flux is expressed in terms of a thermal conductivity tensor $`\kappa _{ij}(\alpha )`$ and a new tensor $`\mu _{ij}(\alpha )`$.
### III.3 Steady state conditions
As shown in the above subsections, the evaluation of the complete nonlinear dependence of the generalized transport coefficients on the shear rate and dissipation requires the analysis of the hydrodynamic behavior of the unsteady reference state. This involves the corresponding numerical integrations of the differential equations obeying the velocity moments of the zeroth-order solution. This is a quite intricate and long problem. However, given that here we are mainly interested in performing a linear stability analysis of the hydrodynamic equations with respect to the steady state, we want to evaluate the transport coefficients in this special case. As a consequence, $`_t^{(0)}T=0`$ and so the condition
$$a^{}P_{xy}^{}=\frac{3}{2}\zeta ^{}.$$
(52)
applies. In Eq. (52), it is understood that $`a^{}`$ and $`P_{xy}^{}=P_{xy}^{(0)}/p`$ are evaluated in the steady state, namely, they are given by Eqs. (24) and (25), respectively. A consequence of Eq. (52) is that the first term on the left hand side of the integral equations (43)–(45) vanishes. In addition, the dependence of the pressure tensor $`P_{ij}^{(0)}`$ on density and temperature occurs explicitly through $`p=nT`$ and through its dependence on $`a^{}`$. In this case, the derivatives $`_nP_{ij}^{(0)}`$ and $`_TP_{ij}^{(0)}`$ can be written more explicitly as
$$n_nP_{ij}^{(0)}=n_npP_{ij}^{}(a^{})=p\left(1a^{}\frac{}{a^{}}\right)P_{ij}^{}(a^{}),$$
(53)
$$T_TP_{ij}^{(0)}=T_TpP_{ij}^{}(a^{})=p\left(1\frac{1}{2}a^{}\frac{}{a^{}}\right)P_{ij}^{}(a^{}).$$
(54)
The dependence of $`P_{ij}^{}`$ on $`a^{}`$ near the steady state is determined in the Appendix C so that, all the terms appearing in the integral equations are explicitly known in the steady state. Under the above conditions, Eqs. (43)–(45) become
$$\left(ac_y\frac{}{c_x}+\right)X_{n,i}+\frac{2a}{3}\frac{T}{n}(P_{xy}^{}+a^{}_a^{}P_{xy}^{})X_{T,i}=Y_{n,i},$$
(55)
$$\left(ac_y\frac{}{c_x}\frac{1}{3}a\left(P_{xy}^{}a^{}_a^{}P_{xy}^{}\right)+\right)X_{T,i}=Y_{T,i},$$
(56)
$$\left(ac_y\frac{}{c_x}+\right)X_{u,k\mathrm{}}a\delta _{ky}X_{u,x\mathrm{}}\zeta _{u,k\mathrm{}}T_Tf^{(0)}=Y_{u,k\mathrm{}},$$
(57)
where it is understood again that in Eqs. (55)–(57) all the quantities are evaluated in the steady state. Henceforth, I will restrict my calculations to this particular case.
Given that in the steady state the coefficient of restitution and the reduced shear rate are coupled, the usual Navier-Stokes transport coefficients for ordinary gases are recovered for elastic collisions ($`a^{}=0`$). Thus, when $`\alpha 1`$ the coefficients become
$$\eta _{ijk\mathrm{}}\eta _0\left(\delta _{ik}\delta _j\mathrm{}+\delta _{jk}\delta _i\mathrm{}\frac{2}{3}\delta _{ij}\delta _k\mathrm{}\right),\kappa _{ij}\kappa _0\delta _{ij},\mu _{ij}0,$$
(58)
where $`\eta _0=p/\nu `$ and $`\kappa _0=15\eta _0/4m`$ are the shear viscosity and thermal conductivity coefficients given by the (elastic) Boltzmann equation.
## IV Results from a simple kinetic model
The explicit form of the generalized transport coefficients $`\mu _{ij}`$, $`\kappa _{ij}`$ and $`\eta _{ijk\mathrm{}}`$ requires to solve the integral equations (55)–((57). Apart from the mathematical difficulties embodied in the Boltzmann collision operator $``$, the fourth-degree velocity moments of the distribution $`f^{(0)}`$ are also needed to determine $`\mu _{ij}`$ and $`\kappa _{ij}`$ and they are not provided in principle by the Grad approximation. Nevertheless, an accurate estimate of these moments from the Boltzmann equation is a formidable task since it would require at least to include the fourth-degree moments in Grad’s solution. In this case, to overcome such difficulties it is useful to consider a model kinetic equation of the Boltzmann equation. As for elastic collisions, the idea is to replace the true Boltzmann collision operator with a simpler, more tractable operator that retains the most relevant physical properties of the Boltzmann operator. Here, I consider a kinetic model BDS99 based on the well-known Bhatnagar-Gross-Krook (BGK) GS03 for ordinary gases where the operator $`J[f,f]`$ is note
$$J[f,f]\beta \nu (ff_0)+\frac{\zeta }{2}\frac{}{𝐜}\left(𝐜f\right).$$
(59)
Here, $`\nu `$ and $`\beta `$ are given by Eqs. (21) and (23), respectively, $`f_0`$ is the local equilibrium distribution (19) and $`\zeta `$ is the cooling rate defined by Eq. (12). As said before, an estimate of $`\zeta `$ to first order in the gradients has been derived in Appendix B. In general, the quantity $`\beta `$ can be considered as an adjustable parameter to optimize the agreement with the Boltzmann equation. In this paper, $`\beta `$ has been chosen to reproduce the true Navier-Stokes shear viscosity coefficient of an inelastic gas of hard spheres single . A slightly different choice for $`\beta `$, namely $`\beta =(1+\alpha )/2`$, is considered in Ref. AS05 .
By taking moments with respect to $`1`$, $`𝐜`$ and $`c^2`$, the model kinetic equation (59) yields the same form of the macroscopic balance equations for mass, momentum, and energy, Eqs. (7)–(9), as those given from the Boltzmann equation. When $`\alpha =1`$, then $`\beta =1`$, $`\zeta =0`$ and so the kinetic model (59) reduces to the BGK equation whose utility to address complex states not accessible via the Boltzmann equation is well-established for elastic gases GS03 . In the case of granular gases, it is easy to show that the kinetic model leads to the same results for the pressure tensor in the USF problem as those given from Grad’s solution to the Boltzmann equation, Eqs. (24)–(25). This result, along with those of Refs. MGSB99 and TTMGSD01 , confirms the reliability of the kinetic model for granular media as well. A summary of the USF results derived from the kinetic model is provided in Appendix D. In particular, beyond rheological properties, recent computer simulations AS05 have confirmed the accuracy of the kinetic model to capture the dependence of the fourth-degree velocity moments (whose expressions are needed to get the coefficients $`\mu _{ij}`$ and $`\kappa _{ij}`$ on dissipation in the USF state. To illustrate it, in Fig. 1 we plot the fourth-degree moment
$$c^4=𝑑𝐜c^4f(𝐜)$$
(60)
relative to its local equilibrium value $`c^4_0=15nT^2/m^2`$. The symbols refer to the numerical results obtained from the DSMC method AS05 . It is quite apparent that the analytical results agree well with simulation data (the discrepancies between both results are smaller than 3%), showing again that the reliability of the kinetic model goes beyond the quasielastic limit.
Let us consider the perturbed USF problem in the context of the kinetic model. By using the model (59), the integral equations (55)–(57) still apply with the only replacement
$$X\nu \beta X\frac{\zeta ^{(0)}}{2}\frac{}{𝐜}\left(𝐜X\right),$$
(61)
in the case of $`X_{n,i}`$ and $`X_{T,i}`$ and
$$X_{ij}\nu \beta X_{ij}\frac{\zeta ^{(0)}}{2}\frac{}{𝐜}\left(𝐜X_{ij}\right)\frac{\zeta _{u,ij}}{2}\frac{}{𝐜}\left(𝐜f^{(0)}\right),$$
(62)
in the case of $`X_{u,ij}`$. In the above equations, $`\zeta ^{(0)}`$ is the zeroth-order approximation to $`\zeta `$ which is given by Eq. (41). With the changes (61) and (62) all the generalized transport coefficients can be easily evaluated from Eqs. (55)–(57). Details of these calculations are also given in Appendix B; a more complete listing can be obtained on request from the author.
The dependence of the generalized transport coefficients on the coefficient of restitution $`\alpha `$ is illustrated in Figs. 2, 3 and 4 for the (reduced) coefficients $`\mu _{ij}^{}`$, $`\kappa _{xy}^{}`$, $`\kappa _{yy}^{}`$, $`\eta _{xxyy}^{}`$, $`\eta _{yyyy}^{}`$, $`\eta _{zzyy}^{}`$, and $`\eta _{xyyy}^{}`$. Here, $`\mu _{ij}^{}=n\mu _{ij}/T\kappa _0`$, $`\kappa _{ij}^{}=\kappa _{ij}/\kappa _0`$ and $`\eta _{ijk\mathrm{}}^{}=\eta _{ijk\mathrm{}}/\eta _0`$, where $`\eta _0=p/\nu `$ and $`\kappa _0=5\eta _0/2m`$ are the elastic values of the shear viscosity and thermal conductivity coefficients given by the BGK kinetic model. In general, we observe that the influence of dissipation on the transport coefficients is quite significant.
With all the transport coefficients known, the new constitutive equations (47) and (48) are completed and the corresponding set of closed hydrodynamic equations (27)–(29) can be derived. They are given by
$$_tn+𝐮_0n+(n\delta 𝐮)=0,$$
(63)
$$_t\delta u_i+a_{ij}\delta u_j+\left(𝐮_0+\delta 𝐮\right)\delta u_i+\frac{1}{mn}\frac{}{r_j}\left(P_{ij}^{(0)}\eta _{ijk\mathrm{}}\frac{\delta u_k}{r_{\mathrm{}}}\right)=0,$$
(64)
$`{\displaystyle \frac{3}{2}}n_tT+{\displaystyle \frac{3}{2}}n(𝐮_0+\delta 𝐮)Ta\eta _{xyij}{\displaystyle \frac{\delta u_i}{r_j}}`$ (65)
$``$ $`{\displaystyle \frac{}{r_i}}\left(\mu _{ij}{\displaystyle \frac{n}{r_j}}+\kappa _{ij}{\displaystyle \frac{T}{r_j}}\right)+\left(P_{ij}^{(0)}\eta _{ijk\mathrm{}}{\displaystyle \frac{\delta u_k}{r_{\mathrm{}}}}\right){\displaystyle \frac{\delta u_i}{r_j}}+aP_{xy}^{(0)}`$
$`=`$ $`{\displaystyle \frac{3}{2}}nT\zeta {\displaystyle \frac{3}{2}}nT\zeta _{u,ij}{\displaystyle \frac{\delta u_i}{r_j}}`$
Note also that consistency would require to consider the term $`aP_{xy}^{(2)}`$ which is of second order in gradients and so, it should be retained. Given that this would require to determine the second order contributions to the fluxes, this term will be neglected in our study. An important feature of our linearized hydrodynamic equations is that they are not restricted to small values of the (reduced) shear rate or, equivalently, to small inelasticity. This allows us to go beyond the usual Navier-Stokes hydrodynamics. The hydrodynamic equations (63)–(65) are the starting point of the linear stability analysis of the USF of the next Section.
## V Linear stability analysis of the steady shear flow state
As said in the Introduction, computer simulations MD have clearly shown that the USF state is unstable with respect to long enough wavelength perturbations. These results have been also confirmed by different analytical results S92 ; B93 ; AN97 ; K00 , most of them based on the Navier-Stokes description that applies to first order in the shear rate. However, given that USF is inherently non-Newtonian SGD04 , the full nonlinear dependence of the transport coefficients on the shear rate is required to perform a consistent linear stability analysis of the nonlinear hydrodynamic equations (63)–(65) with respect to the USF state for small initial excitations. This analysis allows one to determine the hydrodynamic modes for states near USF as well the conditions for instabilities at long wavelengths. A growth of these modes signals the onset of instability, which is ultimately controlled by the dominance of nonlinear terms. Note also that while all the works have been mainly devoted to dense systems, much less attention has been paid to dilute gases.
Let us assume that the deviations $`\delta x_\mu (𝐫,t)=x_\mu (𝐫,t)x_{0\mu }(𝐫)`$ are small, where $`\delta x_\mu (𝐫,t)`$ denotes the deviation of $`\{n,𝐮,T\}`$ from their values in the USF state $`\{n_0,𝐮_0,T_0\}`$. The quantities in the USF verify
$$n_0=T_0=0,𝐮_0=𝖺𝐫,_tT_0=0.$$
(66)
Now, let us linearize Eqs. (63)–(65) with respect to
$$\{\delta x_\mu (𝐫,t)\}\{\delta n(𝐫,t),\delta T(𝐫,t),\delta 𝐮(𝐫,t)\}.$$
(67)
The resulting set of five linearized hydrodynamic equations follows from Eqs. (63)–(65):
$$_t\delta n+ay\frac{}{x}\delta n+n_0\delta 𝐮=0,$$
(68)
$`{\displaystyle \frac{3}{2}}n_0_t\delta T+ay{\displaystyle \frac{}{x}}\delta T+a\delta _{ix}\delta u_y+a\left[(_nP_{xy}^{(0)})\delta n+(_TP_{xy}^{(0)})\delta T\right]`$ (69)
$`+\left(P_k\mathrm{}^{(0)}a\eta _{xyk\mathrm{}}\right){\displaystyle \frac{\delta u_k}{r_{\mathrm{}}}}\mu _{ij}{\displaystyle \frac{^2\delta n}{r_ir_j}}\kappa _{ij}{\displaystyle \frac{^2\delta T}{r_ir_j}}`$
$`=`$ $`{\displaystyle \frac{3}{2}}\zeta _0n_0T_0\left(2{\displaystyle \frac{\delta n}{n_0}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{\delta T}{T_0}}\right){\displaystyle \frac{3}{2}}n_0T_0\zeta _{u,k\mathrm{}}{\displaystyle \frac{\delta u_k}{\mathrm{}}},`$
$$_t\delta u_k+ay\frac{}{x}\delta u_k+a\delta _{kx}\delta u_y+\frac{1}{mn_0}\left[(_nP_k\mathrm{}^{(0)})\frac{\delta n}{r_{\mathrm{}}}+(_TP_k\mathrm{}^{(0)})\frac{\delta T}{r_{\mathrm{}}}\eta _{k\mathrm{}ij}\frac{^2\delta u_i}{r_{\mathrm{}}r_j}\right]=0.$$
(70)
Here, it is understood that the pressure tensor $`P_{ij}^{(0)}`$ and its derivatives with respect to $`n`$ and $`T`$, the cooling rate $`\zeta _0`$ and the transport coefficients $`\eta _{ijk\mathrm{}}`$, $`\mu _{ij}`$, and $`\kappa _{ij}`$ are evaluated in the steady USF state.
To analyze the linearized hydrodynamic equations (68)–(70) it is convenient to transform to the local Lagrangian frame, $`r_i^{}=r_ita_{ij}r_j`$. The Lees-Edwards boundary conditions then become simple periodic boundary conditions in the variable $`𝐫^{}`$ LD97 . A Fourier representation is defined as
$$\delta \stackrel{~}{x}_\mu (𝐤,t)=𝑑𝐫^{}e^{i𝐤𝐫^{}}\delta x_\mu (𝐫,t)=𝑑𝐫e^{i𝐤(t)𝐫}\delta x_\mu (𝐫,t),$$
(71)
where in the second equality $`k_i(t)=k_j(\delta _{ij}ta_{ji})`$. Periodicity conditions requires that $`k_i=2n_i\pi /L_i`$, where $`n_i`$ are integers and $`L_i`$ are the linear dimensions of the system. In this Fourier representation, the resulting set of five linear equations defines the hydrodynamic modes, i.e., linear response excitations to small perturbations. If at least one of the modes grows in time, the reference USF state is linearly unstable. Given the mathematical difficulties involved in the general problem, for the sake of simplicity, here I consider two kind of perturbations: (i) perturbations along the velocity gradient direction only ($`k_x=k_z=0;k_y0`$) and (ii) perturbations in the vorticity direction only ($`k_x=k_y=0;k_z0`$). In both cases, the linearized hydrodynamic equations have time-independent coefficients.
### V.1 Perturbations in the velocity gradient direction ($`k_x=k_z=0;k_y0`$)
Let us consider first perturbations along the $`y`$ direction only. In this case, Eqs. (68)–(70) in this Fourier representation can be written in the matrix form
$$_\tau \delta \stackrel{~}{x}_\mu ^{}+F_{\mu \nu }\delta \stackrel{~}{x}_\nu ^{}=0,$$
(72)
where the dimensionless quantities $`\tau =\nu _0t`$ and $`\delta \stackrel{~}{x}_\mu ^{}\{\rho _k,\theta _k,𝐰_k\}`$, with
$$\rho _k=\frac{\delta \stackrel{~}{n}}{n_0},\theta _k=\frac{\delta \stackrel{~}{T}}{T_0},𝐰_k=\frac{\delta \stackrel{~}{𝐮}}{\sqrt{T_0/m}},$$
(73)
have been introduced. The matrix $`F_{\mu \nu }`$ is
$$F_{\mu \nu }=2C\delta _{\mu 2}\delta _{\nu 1}+C\delta _{\mu 2}\delta _{\nu 2}+a^{}\delta _{\mu 3}\delta _{\nu 4}ik^{}G_{\mu \nu }+k^2H_{\mu \nu },$$
(74)
where $`a^{}=a/\nu _0`$, $`\nu _0`$ is the collision frequency (21) of the reference state and $`k^{}=\mathrm{}_0k`$, $`\mathrm{}_0=\sqrt{T_0/m}/\nu _0`$ being of the order of the mean free path. In addition, we have introduced the coefficient
$$C(\alpha )=\frac{1}{3}a^{}\left(1+a^{}_a^{}\right)P_{xy}^{},$$
(75)
and the square matrices
$$𝖦=\left(\begin{array}{ccccc}0& 0& 0& 1& 0\\ 0& 0& \frac{2}{3}(P_{xy}^{}a^{}\eta _{xyxy}^{})+\zeta _{xy}^{}& \frac{2}{3}(P_{yy}^{}a^{}\eta _{xyyy}^{})+\zeta _{yy}^{}& 0\\ \left(1a^{}_a^{}\right)P_{xy}^{}& \left(1\frac{1}{2}a^{}_a^{}\right)P_{xy}^{}& 0& 0& 0\\ \left(1a^{}_a^{}\right)P_{yy}^{}& \left(1\frac{1}{2}a^{}_a^{}\right)P_{yy}^{}& 0& 0& 0\\ 0& 0& 0& 0& 0\end{array}\right),$$
(76)
$$𝖧=\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ \frac{5}{3}\mu _{yy}^{}& \frac{5}{3}\kappa _{yy}^{}& 0& 0& 0\\ 0& 0& \eta _{xyxy}^{}& \eta _{xyyy}^{}& 0\\ 0& 0& \eta _{yyxy}^{}& \eta _{yyyy}^{}& 0\\ 0& 0& 0& 0& \eta _{zyzy}^{}\end{array}\right),$$
(77)
have been also introduced. Here, $`P_{ij}^{}=P_{ij}^{(0)}/n_0T_0`$ and
$$\zeta _{ij}^{}=\frac{1}{48}(1\alpha ^2)\left(P_k\mathrm{}^{}\delta _k\mathrm{}\right)\eta _{k\mathrm{}ij}^{}.$$
(78)
The eigenvalues $`\lambda _\mu (k,\alpha )`$ of the matrix $`𝖥(k,\alpha )`$ determine the time evolution of $`\delta \stackrel{~}{x}_\mu ^{}(k,t)`$. In the case that the real parts of the eigenvalues $`\lambda _\mu (k,\alpha )`$ are positive, then the USF state will be linearly stable. Before considering the general case , it is convenient to consider some special limits. Thus, in the elastic limit ($`\alpha =1`$), the hydrodynamic modes of the Navier-Stokes equations (for the particular case considered here and in the context of the BGK model) are recovered RL77 , namely, two sound modes, a heat mode and a two-fold degenerate shear mode. To second order in $`k^{}`$ they are given by
$$\lambda _\mu (k,\alpha =1)\{i\sqrt{\frac{5}{3}}k^{}+k^2,i\sqrt{\frac{5}{3}}k^{}+k^2,k^2,k^2,k^2\},$$
(79)
and consequently, excitations around equilibrium are damped. It is also quite illustrative to get the modes by setting $`k=0`$, namely, consider small, homogenous deviations from the steady shear flow state. In this case, it is easy to see that $`\rho _k`$ and $`w_{y,k}`$ are constant and
$$w_{x,k}(\tau )=w_{x,k}(0)a\tau w_{y,k}(0),$$
(80)
$$\theta _k(\tau )=\theta _k(0)e^{C\tau }2\rho _k(0).$$
(81)
The mode associated with $`w_{x,k}`$ is unstable to an initial perturbation in $`w_{y,k}`$, leading to an unbounded linear change in time. However, stability is still possible at finite $`k`$ if this behavior is modulated by exponential damping factors. With respect to the temperature field, initial disturbances decay at $`\tau \mathrm{}`$ if the coefficient $`C(\alpha )>0`$. Figure 5 shows that the coefficient $`C`$ is positive for any value of $`\alpha `$ and so, this mode is stable with a finite decay constant.
The analysis for $`k0`$ requires to get the eigenvalues $`\lambda _\mu (k^{},\alpha )`$ with the full nonlinear dependence of $`k^{}`$. However, the structure of $`𝖥(k,\alpha )`$ shows that the perturbation $`\delta \stackrel{~}{x}_5^{}\delta \stackrel{~}{u}_z`$ is decoupled from the other four modes and hence can be obtained more easily. This is due to the choice of gradients along the $`y`$ direction only. The eigenvalue associated with this mode is positive and is simply given by
$$\lambda _5(k,\alpha )=\eta _{zyzy}^{}k^2,\eta _{zyzy}^{}=\frac{\beta }{(\beta +\zeta ^{})^2},$$
(82)
where $`\zeta ^{}`$ is defined by Eq. (22). The remaining modes correspond to $`\rho _k`$, $`\theta _k`$ and the components of the velocity field $`w_{x,k}`$ and $`w_{y,k}`$. They are the solutions of a quartic equation with coefficients that depend on $`k^{}`$ and $`\alpha `$. The results show that $`\text{Re}\lambda _\mu (k^{},\alpha )>0`$ for all the values of the coefficient of restitution $`\alpha `$ and consequently, the flow remains stable to this kind of perturbations. As an illustration, the dispersion relations for a gas with $`\alpha =0.8`$ are plotted in Fig. 6. It is apparent that all the real parts of the eigenvalues $`\lambda _\mu `$ are positive in the range of values of wavenumber $`k^{}`$ considered. Our conclusion agrees with previous stability analysis S92 ; AN97 based on the Navier-Stokes constitutive equations where it was found a minimum value of solid fraction (around 0.156) below which the USF is stable. Given that our system is a dilute gas (zero density), the present results confirm previous findings when one uses the improved transport coefficients.
### V.2 Perturbations in the vorticity direction ($`k_x=k_y=0;k_z0`$)
The variation of the hydrodynamic modes with wavenumber $`k=k_z`$ in the vorticity direction is considered next. This situation has not been widely studied in the literature since most of the studies have been focussed on 2-D flows due to the relative computational efficiency with which they can be analyzed. Here, for the sake of simplicity, I consider perturbations for which $`\delta u_x=\delta u_y=0`$ and so, the eigenvalues $`\lambda _\mu (k^{},\alpha )`$ obey a cubic equation. The analysis is similar to the one carried out in the previous section and so, details will be omitted. For a given value of $`\alpha `$, it can be seen that this dispersion relation has one real root and a complex conjugate pair of damping modes. The instability arises from the real root since this mode $`\lambda _\mu (k^{},\alpha )>0`$ if $`k^{}`$ is larger than a certain threshold value $`k_s^{}(\alpha )`$. This value can be obtained by solving $`\lambda _\mu (k^{},\alpha )=0`$. As a consequence, the USF state is linearly stable against excitations with a wavenumber $`k^{}>k_s^{}(\alpha )`$. The stability line $`k_s^{}(\alpha )`$ is plotted in Fig. 7 as a function of the coefficient of restitution. Above this line the modes are stable, while below this line they are unstable. For comparison, the corresponding stability line obtained from the approximations made in previous works S92 ; AN97 is also plotted. This line can be formally obtained from the results derived in this paper when one replaces the expressions of the coefficients $`\eta _{ijk\mathrm{}}`$, $`\kappa _{ij}`$, and $`\mu _{ij}`$ by their corresponding Navier-Stokes expressions single . It is apparent that the Navier-Stokes approximation captures the qualitative dependence of $`k_s^{}`$ on $`\alpha `$, although as expected quantitative discrepancies between both descriptions appear as the dissipation increases. Thus, for instance, for $`\alpha =0.8`$ the discrepancies between both approaches are about 22 $`\%`$ while for $`\alpha =0.5`$ the discrepancies are about 49$`\%`$. The prediction of a long-wavelength instability for the USF state has been observed in early molecular dynamics simulations MD and qualitatively agrees with the previous analytical results based on the Navier-Stokes equations S92 ; B93 ; AN97 ; K00 . At a quantitative level, the lack of numerical results from the Boltzmann equation prevent us to carry out a more detailed comparison to confirm the results derived from this kinetic model. We hope that the results offered here will stimulate the performance of such computer simulations.
## VI Summary and Discussion
The objective of this paper has been to study the transport properties of a granular gas of inelastic hard spheres for the special nonequilibrium states near the uniform (simple) shear flow (USF). Although the derivation of the Navier-Stokes equations (with explicit expressions for the transport coefficients appearing in them) from a microscopic description has been widely worked out in the past single ; mixture , the analysis of transport in a strongly shearing granular gas has received little attention due perhaps to its complexity and technical difficulties. Very recently, a generalized Chapman-Enskog method has been proposed to analyze transport around nonequilibrium states in granular gases L05 . In the case of the USF state, due to the anisotropy induced in the system by the presence of shear flow, tensorial quantities are required to describe the momentum and heat fluxes instead of the usual Navier-Stokes transport coefficients single ; mixture . In this paper we have been interested in a physical situation where weak spatial gradients of density, velocity and temperature coexist with a strong shear rate. Under these conditions, the corresponding generalized transport coefficients characterizing heat and momentum transport are nonlinear functions of both the (reduced) shear rate $`a^{}`$ and the coefficient of restitution $`\alpha `$. The determination of such transport coefficients has been the primary aim of this paper.
Due to the difficulties embodied in this problem, a low-density gas described by the inelastic Boltzmann equation has been considered. Although the exact solution to the Boltzmann equation in the (steady) USF is not known, a good estimate of the relevant transport properties can be obtained by means of Grad’s method MG02 ; G02 ; SGD04 . The reliability of this approximation has been recently assessed by comparison with Monte Carlo simulations of the Boltzmann equation MG02 ; AS05 . Assuming that the USF state is slightly perturbed, the Boltzmann equation has been solved by a Chapman-Enskog-like expansion where the shear flow state is used as the reference state rather than the local equilibrium or the (local) homogeneous cooling state. Due to the spatial dependence of the zeroth-order distribution $`f^{(0)}`$ (reference state), this distribution is not in general stationary and only in very special conditions has a simple relation with the (steady) USF distribution L05 . Here, since one the main goals has been to address a stability analysis of the USF state, for practical purposes my results have been specialized to the steady state, namely, when the hydrodynamic variables satisfy the balance condition (52). In this situation, the (reduced) shear rate $`a^{}`$ is coupled with the coefficient of restitution $`\alpha `$ \[see Eq. (25)\] so that the latter is the relevant parameter of the problem. In the first order of the expansion the momentum and heat fluxes are given by Eqs. (47) and (48), respectively, where the set of generalized transport coefficients $`\eta _{ijk\mathrm{}}`$, $`\mu _{ij}`$, and $`\kappa _{ij}`$ are given in terms of the solutions of the linear integral equations (55)–(57). As expected, there are many new transport coefficients in comparison to the case of states near equilibrium or cooling state. These coefficients provide all the information on the physical mechanisms involved in the transport of momentum and energy under shear flow.
Practical applications require to solve the integral equations (55)–(57), which is in general quite a complex problem. In addition, the fourth-degree velocity moments of USF (whose evaluation would require to consider higher-order terms in Grad’s solution (18) of the Boltzmann equation) are needed to determine the coefficients $`\kappa _{ij}`$ and $`\mu _{ij}`$. To overcome such mathematical difficulties, here a kinetic model of the Boltzmann equation BDS99 has been used. This kinetic model can be considered as an extension of the well-known BGK equation to inelastic gases. Although the kinetic model is only a crude representation of the Boltzmann equation, it does preserve the most important features for transport, such as the homogeneous cooling state and the macroscopic conservation laws. The model has a free parameter $`\beta `$ to be adjusted to fit a given property of the Boltzmann equation. Here, $`\beta `$ is given by Eq. (23) to get good quantitative agreement of the Navier-Stokes shear viscosity coefficient obtained from the Boltzmann equation. Furthermore, this choice yields the same results for rheological properties in the USF problem as those derived from the Boltzmann equation by means of Grad’s method. On the other hand, given that the model does not intend to mimic the behavior of the true distribution function beyond the thermal velocity region, discrepancies between the kinetic model and the Boltzmann equation are expected beyond the second-degree velocity moments (which quantify the elements of the pressure tensor). Nevertheless, a recent comparison with Monte Carlo simulations of the Boltzmann equation AS05 have shown the accuracy of the kinetic model predictions for the fourth-degree moments. As illustrated in Fig. 1, the semi-quantitative agreement between theory and simulation is not restricted to the quasielastic limit ($`\alpha 0.99`$) since it covers values of large dissipation ($`\alpha 0.5`$). The use of this kinetic model allows one to get the explicit dependence of the generalized transport coefficients on the coefficient of restitution. This dependence has been illustrated in some cases showing that in general the deviation of the transport coefficients from their corresponding elastic values is quite significant.
With these new expressions for the fluxes, a closed set of generalized hydrodynamic equations for states close to USF has been derived. A stability analysis of these linearized hydrodynamic equations with respect to the USF state have been also carried out to identify the conditions for stability in terms of dissipation. Two different kind of perturbations to the USF state has been analyzed: (i) perturbations along the velocity gradient only ($`k_y0`$) and (ii) perturbations along the vorticity direction only ($`k_z0`$). In the first case, previous results S92 ; AN97 have shown that the USF is stable for a dilute gas while the USF becomes unstable in the second case for all $`\alpha `$ L04 . These results agree with these findings and the USF is unstable for any finite value of dissipation at sufficiently long wave lengths when disturbances are generated in the orthogonal direction to the shear flow plane. On the other hand, as expected, quantitative discrepancies between our results and those given S92 ; AN97 from the Navier-Stokes approximation become significant as the dissipation increases. These differences have been illustrated in Fig. 7 for the stability line. Although the instability of the USF has been extensively studied for many authors by using a Navier-Stokes description S92 ; B93 ; AN97 as well as solutions of the Boltzmann equation in the quasielastic limit K00 ; K01 , I am not aware of any previous solution of the hydrodynamic equations where the generalized transport coefficients describing transport around USF were taken into account. The analytical results found in this paper allows a quantitative comparison with numerical solutions to the Boltzmann equation for finite dissipation. As happens for the USF problem for elastic LD97 ; LDMSL96 ; MSLDL98 and inelastic MG02 ; AS05 gases, one expects that the results reported here compare well with such simulations, confirming again the reliability of the kinetic theory results to characterize the onset and the first stages of evolution of the clustering instability. We hope to carry out these simulations in the next future.
On the other hand, the stability analysis performed here has only considered spatial variations along the $`y`$ and $`z`$ directions. More complex dynamics is expected in the general case of arbitrary direction for the spatial perturbation. This will be worked elsewhere along with comparison with direct Monte Carlo computer simulations of the Boltzmann equation. Another possible direction of study is the extension of the present approach to other physically interesting reference states, such as the nonlinear Couette flow. This is a more realistic shearing problem than the USF state since combined heat and momentum transport appears in the system. Given that an exact solution to the kinetic model used here is known for the Couette flow problem TTMGSD01 , the reference distribution for the Chapman-Enskog-like expansion is available.
###### Acknowledgements.
I am grateful to Dr. James Lutsko for pointing out on some conceptual errors that were present in a previous version of this manuscript. It is also a pleasure to thank to Dr. Andrés Santos for valuable discussions. Partial support of the Ministerio de Ciencia y Tecnología (Spain) through Grant No. FIS2004-01399 (partially financed by FEDER funds) is acknowledged.
## Appendix A Chapman-Enskog expansion
Inserting the expansions (32) and (34) into Eq. (26), one gets the kinetic equation for $`f^{(1)}`$,
$$\left(_t^{(0)}aV_y\frac{}{V_x}+\right)f^{(1)}=\left[_t^{(1)}+(𝐕+𝐮_0)\right]f^{(0)},$$
(83)
where $``$ is the linearized Boltzmann collision operator
$$X\left(J[f^{(0)},X]+J[X,f^{(0)}]\right).$$
(84)
The velocity dependence on the right side of Eq. (83) can be obtained from the macroscopic balance equations to first order in the gradients. They are given by
$$_t^{(1)}n+𝐮_0n=(n\delta 𝐮),$$
(85)
$$_t^{(1)}\delta 𝐮+(𝐮_0+\delta 𝐮)\delta 𝐮=\frac{1}{\rho }𝖯^{(0)},$$
(86)
$$\frac{3}{2}n_t^{(1)}T+\frac{3}{2}n(𝐮_0+\delta 𝐮)T+aP_{xy}^{(1)}+𝖯^{(0)}:\delta 𝐮=\frac{3}{2}p\zeta ^{(1)},$$
(87)
where $`\rho =mn`$ is the mass density,
$$P_{ij}^{(1)}=𝑑𝐜mc_ic_jf^{(1)}(𝐜),$$
(88)
and
$$\zeta ^{(1)}=\frac{1}{3p}𝑑𝐜mc^2f^{(1)}.$$
(89)
Use of Eqs. (85)–(87) in Eq. (83) yields
$$\left(_t^{(0)}aV_y\frac{}{V_x}+\right)f^{(1)}\zeta ^{(1)}T\frac{f^{(0)}}{T}=𝐘_nn+𝐘_TT+𝖸_u:\delta 𝐮,$$
(90)
where
$$Y_{n,i}=\frac{f^{(0)}}{n}c_i+\frac{1}{\rho }\frac{f^{(0)}}{\delta u_j}\frac{P_{ij}^{(0)}}{n},$$
(91)
$$Y_{T,i}=\frac{f^{(0)}}{T}c_i+\frac{1}{\rho }\frac{f^{(0)}}{\delta u_j}\frac{P_{ij}^{(0)}}{T},$$
(92)
$$Y_{u,ij}=n\frac{f^{(0)}}{n}\delta _{ij}\frac{f^{(0)}}{\delta u_i}c_j+\frac{2}{3n}\frac{f^{(0)}}{T}\left(P_{ij}^{(0)}a\eta _{xyij}\right).$$
(93)
According to Eqs. (91)–(92), $`Y_{u,ij}`$ has the same symmetry properties (16) as the distribution function $`f^{(0)}`$ while $`Y_{n,i}`$ and $`Y_{T,i}`$ are odd functions in the velocity $`𝐜`$.
The solution to Eq. (90) has the form
$$f^{(1)}=X_{n,i}(𝐜)_in+X_{T,i}(𝐜)_iT+X_{u,ji}(𝐜)_i\delta u_j.$$
(94)
Note that in Eq. (93) the coefficients $`\eta _{ijk\mathrm{}}`$ are defined through Eq. (49). Substitution of the solution (94) into the relation (89) allows one to write the cooling rate in the form
$$\zeta ^{(1)}=\zeta _{n,i}_in+\zeta _{T,i}_iT+\zeta _{u,ji}_i\delta u_j,$$
(95)
where
$$\left(\begin{array}{c}\zeta _{n,i}\\ \zeta _{T,i}\\ \zeta _{u,ij}\end{array}\right)=\frac{1}{3p}𝑑𝐜mc^2\left(\begin{array}{c}X_{n,i}\\ X_{T,i}\\ X_{u,ij}\end{array}\right).$$
(96)
However, given that $`X_{n,i}`$ and $`X_{T,i}`$ are odd functions in $`𝐜`$ \[see for instance, Eqs. (101) and (102) below\], the terms proportional to $`n`$ and $`T`$ vanish by symmetry, i.e.,
$$\zeta _{n,i}=\zeta _{T,i}=0.$$
(97)
Thus, the only nonzero contribution to $`\zeta ^{(1)}`$ comes from the term proportional to the tensor $`_i\delta u_j`$:
$$\zeta ^{(1)}=\zeta _{u,ji}_i\delta u_j.$$
(98)
An estimate of the tensor $`\zeta _{u,ij}`$ has been made in Appendix B by considering the leading terms in a Sonine polynomial expansion of the distribution $`f^{(1)}`$. Its expression is given by Eq. (113). As expected, $`\zeta _{u,ij}`$ vanishes in the elastic limit ($`\alpha =1`$).
The coefficients $`X_{n,i}`$, $`X_{T,i}`$, and $`X_{u,ij}`$ are functions of the peculiar velocity $`𝐜`$ and the hydrodynamic fields. In addition, there are contributions from the time derivative $`_t^{(0)}`$ acting on the temperature and velocity gradients given by
$`_t^{(0)}_iT`$ $`=`$ $`_i_t^{(0)}T`$ (99)
$`=`$ $`\left({\displaystyle \frac{2a}{3n^2}}(1n_n)P_{xy}^{(0)}{\displaystyle \frac{\zeta ^{(0)}T}{n}}\right)_in\left({\displaystyle \frac{2a}{3n}}_TP_{xy}^{(0)}+{\displaystyle \frac{3}{2}}\zeta ^{(0)}\right)_iT,`$
$$_t^{(0)}_i\delta u_j=_i_t^{(0)}\delta u_j=a_{jk}_i\delta u_k.$$
(100)
Substituting (98) into (90) and identifying coefficients of independent gradients gives the set of equations
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+ac_y\frac{}{c_x}\right]X_{n,i}+\frac{T}{n}\left[\frac{2a}{3p}(1n_n)P_{xy}^{(0)}\zeta ^{(0)}\right]X_{T,i}=Y_{n,i},$$
(101)
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+\frac{2a}{3p}T(_TP_{xy}^{(0)})+\frac{3}{2}\zeta ^{(0)}+ac_y\frac{}{c_x}\right]X_{T,i}=Y_{T,i},$$
(102)
$$\left[\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_T+ac_y\frac{}{c_x}\right]X_{u,k\mathrm{}}a\delta _{ky}X_{u,x\mathrm{}}\zeta _{u,k\mathrm{}}T_Tf^{(0)}=Y_{u,k\mathrm{}}.$$
(103)
Upon writing Eqs. (101)–(103), use has been made of the property
$`_t^{(0)}X`$ $`=`$ $`{\displaystyle \frac{X}{T}}_t^{(0)}T+{\displaystyle \frac{X}{\delta u_i}}_t^{(0)}\delta u_i`$ (104)
$`=`$ $`\left({\displaystyle \frac{2}{3n}}aP_{xy}^{(0)}+T\zeta ^{(0)}\right){\displaystyle \frac{X}{T}}+a_{ij}\delta u_j{\displaystyle \frac{X}{c_i}},`$
where in the last step we have taken into account that $`X`$ depends on $`\delta 𝐮`$ through $`𝐜=𝐕\delta 𝐮`$.
## Appendix B Evaluation of the cooling rate
In this Appendix the contribution $`\zeta _{u,ij}`$ to the cooling rate $`\zeta ^{(1)}`$ is evaluated by expanding $`X_{u,ij}`$ as series in Sonine polynomials and taking the lowest order truncation. The tensor $`\zeta _{u,ij}`$ is given by
$`\zeta _{u,ij}`$ $`=`$ $`{\displaystyle \frac{1}{3p}}{\displaystyle 𝑑𝐜_1mc_1^2X_{u,ij}}`$ (105)
$`=`$ $`{\displaystyle \frac{1}{3p}}{\displaystyle 𝑑𝐜_1mc_1^2\left\{J[𝐜_1|f^{(0)},X_{u,ij}]+J[𝐜_1|X_{u,ij},f^{(0)}]\right\}}.`$
A useful identity for an arbitrary function $`h(𝐜_1)`$ is
$$𝑑𝐜_1h(𝐜_1)J[𝐜_1|f,g]=\sigma ^2𝑑𝐜_1𝑑𝐜_2f(𝐜_1)g(𝐜_2)𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝝈}𝐠)(\widehat{𝝈}𝐠)\left[h(𝐜_1^{\prime \prime })h(𝐜_1)\right],$$
(106)
where $`𝐠=𝐜_1𝐜_2`$ and
$$𝐜_1^{\prime \prime }=𝐜_1\frac{1}{2}(1+\alpha )(\widehat{𝝈}𝐠)\widehat{𝝈}.$$
(107)
Using (106), Eq. (105) can be written as
$$\zeta _{u,ij}=\frac{m}{6p}\sigma ^2(1\alpha ^2)𝑑𝐜_1𝑑𝐜_2f^{(0)}(𝐜_1)X_{u,ij}(𝐜_2)𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝝈}𝐠)(\widehat{𝝈}𝐠)^3.$$
(108)
The integration over $`\widehat{𝝈}`$ in (108) yields
$$\zeta _{u,ij}=\frac{m}{12p}\pi \sigma ^2(1\alpha ^2)𝑑𝐜_1𝑑𝐜_2f^{(0)}(𝐜_1)X_{u,ij}(𝐜_2).$$
(109)
This equation is still exact. To perform the integrals over $`𝐜_1`$ and $`𝐜_2`$ one takes the Grad approximation (18) to $`f^{(0)}`$ and expands $`X_{u,ij}`$ in Sonine polynomials. In this case and according to the anisotropy of the USF problem, one takes the approximation
$$X_{u,k\mathrm{}}(𝐜)\frac{1}{2nT^2}D_{ij}\eta _{ijk\mathrm{}}f_0(c),$$
(110)
where
$$f_0(c)=n\left(\frac{m}{2\pi T}\right)^{3/2}\mathrm{exp}\left(\frac{mc^2}{2T}\right)$$
(111)
is the Maxwellian distribution and
$$D_{ij}(𝐜)=m\left(c_ic_j\frac{1}{3}c^2\delta _{ij}\right).$$
(112)
Next, change variables to the (dimensionless) relative velocity $`𝐠^{}=(𝐜_1𝐜_2)/v_0`$ and center of mass $`𝐆^{}=(𝐜_1+𝐜_2)/2v_0`$, where $`v_0=\sqrt{2T/m}`$ is the thermal velocity. A lengthy calculation leads to
$`\zeta _{u,ij}`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{v_0\sigma ^2}{\pi ^2T}}(1\alpha ^2){\displaystyle 𝑑𝐠^{}𝑑𝐆^{}g^3e^{2G^2}e^{g^2/2}}`$ (113)
$`\times \left[G_k^{}G_{\mathrm{}}^{}G_m^{}G_n^{}{\displaystyle \frac{1}{18}}g^2G^2\left(\delta _{km}\delta _\mathrm{}n+\delta _{kn}\delta _\mathrm{}m\right)+{\displaystyle \frac{1}{16}}g_k^{}g_{\mathrm{}}^{}g_m^{}g_n^{}\right]\left({\displaystyle \frac{P_{mn}}{nT}}\delta _{mn}\right)\eta _{k\mathrm{}ij}`$
$`=`$ $`{\displaystyle \frac{1}{15}}\sigma ^2\sqrt{{\displaystyle \frac{\pi }{mT}}}(1\alpha ^2)\left({\displaystyle \frac{P_k\mathrm{}}{nT}}\delta _k\mathrm{}\right)\eta _{k\mathrm{}ij}.`$
Of course, when $`\alpha =1`$, then $`\zeta _{u,ij}=0`$.
## Appendix C Behavior of the zeroth-order velocity moments near the steady state
This Appendix addresses the behavior of the velocity moments of the zeroth-order distribution $`f^{(0)}`$ near the steady state. Let us start with the elements of the pressure tensor $`P_{ij}^{(0)}`$. In the context of the Boltzmann equation and by using Grad’s approximation (18), they verify the equation
$$\left(\frac{2}{3n}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)\frac{}{T}P_{ij}^{(0)}+a_i\mathrm{}P_j\mathrm{}^{(0)}+a_j\mathrm{}P_i\mathrm{}^{(0)}=\nu \left[\beta \left(P_{ij}^{(0)}p\delta _{ij}\right)+\zeta ^{}P_{ij}^{(0)}\right].$$
(114)
Since we are interested in the hydrodynamic solution, the temperature derivative term can be written as
$$T_TP_{ij}^{(0)}=T_TpP_{ij}^{}=p\left(1\frac{1}{2}a^{}\frac{}{a^{}}\right)P_{ij}^{},$$
(115)
where $`P_{ij}^{}=P_{ij}^{(0)}/p`$. Upon deriving (115), use has been made of the fact that the dimensionless pressure tensor $`P_{ij}^{}`$ depends on $`T`$ only through its dependence on the reduced shear rate $`a^{}=a/\nu (n,T)`$. In dimensionless form, the set of equations (114) become
$$\left(\frac{2}{3}a^{}P_{xy}^{}+\zeta ^{}\right)\left(1\frac{1}{2}a^{}\frac{}{a^{}}\right)P_{ij}^{}+a_i\mathrm{}^{}P_j\mathrm{}^{}+a_j\mathrm{}^{}P_i\mathrm{}^{}=\left[\beta \left(P_{ij}^{}\delta _{ij}\right)+\zeta ^{}P_{ij}^{}\right],$$
(116)
where
$$\zeta ^{}=\frac{\zeta ^{(0)}}{\nu }=\frac{5}{12}(1\alpha ^2).$$
(117)
Let us consider the elements $`P_{xy}^{}`$ and $`P_{yy}^{}=P_{zz}^{}`$. From Eq. (114), one gets
$$\left(\frac{2}{3}a^{}P_{xy}^{}+\zeta ^{}\right)\left(1\frac{1}{2}a^{}\frac{}{a^{}}\right)P_{xy}^{}+a^{}P_{yy}^{}=\left(\beta +\zeta ^{}\right)P_{xy}^{},$$
(118)
$$\left(\frac{2}{3}a^{}P_{xy}^{}+\zeta ^{}\right)\left(1\frac{1}{2}a^{}\frac{}{a^{}}\right)P_{yy}^{}=\left(\beta +\zeta ^{}\right)P_{yy}^{}+\beta .$$
(119)
This set of equations have a singular point corresponding to the steady state solution, i.e., when $`a^{}(T)=a_s^{}`$ where $`a_s^{}(\alpha )`$ is the steady state value of $`a^{}`$ given by Eq. (25). Since we are interested in the solution of Eqs. (118) and (119) near the steady state, we assume that in this region $`P_{xy}^{}`$ and $`P_{yy}^{}`$ behave as
$$P_{xy}^{}=P_{xy,s}^{}+\left(\frac{P_{xy}^{}}{a^{}}\right)_s(a^{}a_s^{})+\mathrm{},$$
(120)
$$P_{yy}^{}=P_{yy,s}^{}+\left(\frac{P_{yy}^{}}{a^{}}\right)_s(a^{}a_s^{})+\mathrm{},$$
(121)
where the subscript $`s`$ means that the quantities are evaluated in the steady state. Substitution of (120) and (121) into Eqs. (118) and (119) allows one to determine the corresponding derivatives. The result is
$$\left(\frac{P_{yy}^{}}{a^{}}\right)_s=4P_{yy,s}^{}\frac{a_s^{}C+P_{xy,s}^{}}{2a_s^2C+6\beta +3\zeta ^{}},$$
(122)
where $`C\left(P_{xy}^{}/a^{}\right)_s`$ is the real root of the cubic equation
$$2a_s^4C^3+12a_s^2(\zeta ^{}+\beta )C^2+\frac{9}{2}(7\zeta ^2+14\zeta ^{}\beta +4\beta ^2)C+9\beta (\zeta ^{}+\beta )^2(2\beta ^22\zeta ^2\beta \zeta ^{}).$$
(123)
Equations (122) and (123) can be also obtained from a different way. Let us write the set of equations (118) and (119) as
$$\frac{P_{xy}^{}}{a^{}}=\frac{2P_{yy}^{}\frac{2}{a^{}}P_{xy}^{}\left(\beta \frac{2}{3}P_{xy}^{}a^{}\right)}{\zeta ^{}+\frac{2}{3}a^{}P_{xy}^{}},$$
(124)
$$\frac{P_{yy}^{}}{a^{}}=\frac{2\beta 2P_{yy}^{}\left(\beta \frac{2}{3}P_{xy}^{}a^{}\right)}{a^{}\left(\zeta ^{}+\frac{2}{3}a^{}P_{xy}^{}\right)}.$$
(125)
In the steady state limit ($`a^{}a_s^{}`$), the numerators and denominators of Eqs. (124) and (125) vanish. Evaluating the corresponding limit by means of l’Hopital’s rule, one reobtains the above results (122) and (123). This procedure can be used to get the behavior of the remaining velocity moments near the steady state.
The behavior of the fourth-degree velocity moments of the distribution $`f^{(0)}`$ near the steady state is also needed to determine the transport coefficients $`\mu _{ij}`$ and $`\kappa _{ij}`$ associated with the heat flux in the first-order solution. To evaluate this behavior we use the Boltzmann kinetic model (59). Let us introduce the velocity moments of the zeroth-order distribution
$$M_{k_1,k_2,k_3}^{(0)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}f^{(0)}(𝐜)$$
(126)
These moments verify the equation
$`\left({\displaystyle \frac{2}{3n}}aP_{xy}^{(0)}+T\zeta ^{(0)}\right)_TM_{k_1,k_2,k_3}^{(0)}`$ $`+`$ $`ak_1M_{k_11,k_2+1,k_3}^{(0)}=\nu \beta (M_{k_1,k_2,k_3}^{(0)}`$ (127)
$`N_{k_1,k_2,k_3})k{\displaystyle \frac{\zeta ^{(0)}}{2}}M_{k_1,k_2,k_3}^{(0)},`$
where $`k=k_1+k_2+k_3`$ and $`N_{k_1,k_2,k_3}`$ are the velocity moments of the Gaussian distribution $`f_0`$. As before, the derivative $`_TM_{k_1,k_2,k_3}^{(0)}`$ can be written as
$`T_TM_{k_1,k_2,k_3}^{(0)}`$ $`=`$ $`T_Tn\left({\displaystyle \frac{2T}{m}}\right)^{k/2}M_{k_1,k_2,k_3}^{}(a^{})`$ (128)
$`=`$ $`n\left({\displaystyle \frac{2T}{m}}\right)^{k/2}{\displaystyle \frac{1}{2}}\left(ka^{}_a^{}\right)M_{k_1,k_2,k_3}^{}(a^{}).`$
In dimensionless form, Eq. (127) become
$`\left({\displaystyle \frac{2}{3}}a^{}P_{xy}^{}+\zeta ^{}\right){\displaystyle \frac{1}{2}}(ka^{}_a^{})M_{k_1,k_2,k_3}^{}`$ $`+`$ $`k_1a^{}M_{k_11,k_2+1,k_3}^{}+(\beta +{\displaystyle \frac{k}{2}}\zeta ^{})M_{k_1,k_2,k_3}^{}`$ (129)
$`\beta N_{k_1,k_2,k_3}^{}=0,`$
where $`N_{k_1,k_2,k_3}^{}`$ are the reduced moments of the the Gaussian distribution given by
$$N_{k_1,k_2,k_3}^{}=\pi ^{3/2}\mathrm{\Gamma }\left(\frac{k_1+1}{2}\right)\left(\frac{k_2+1}{2}\right)\left(\frac{k_3+1}{2}\right)$$
(130)
if $`k_1`$, $`k_2`$, and $`k_3`$ are even, being zero otherwise. Equation (129) gives the expressions of the reduced moments $`M_{k_1,k_2,k_3,s}^{}`$ in the steady state. To get $`_a^{}M_{k_1,k_2,k_3}^{}`$ in the steady state, we differentiate with respect to $`a^{}`$ both sides of Eq. (129) and then takes the limit $`aa_s^{}`$. In general, it is easy to see that the problem becomes linear so that it can be easily solved. To illustrate the procedure, let us consider for simplicity the moment $`M_{040}^{}`$, which obeys the equation
$$\left(\frac{2}{3}a^{}P_{xy}^{}+\zeta ^{}\right)\left(2\frac{1}{2}a^{}_a^{}\right)M_{040}^{}+\left(\beta +2\zeta ^{}\right)M_{040}^{}\frac{3}{4}\beta =0.$$
(131)
From this equation, one gets the identity
$`\left({\displaystyle \frac{2}{3}}a^{}P_{xy}^{}+\zeta ^{}\right)_a^{}\left[\left(2{\displaystyle \frac{1}{2}}a^{}_a^{}\right)M_{040}^{}\right]`$ $``$ $`{\displaystyle \frac{2}{3}}\left[P_{xy}^{}+a^{}(_a^{}P_{xy}^{})\right]\left(2{\displaystyle \frac{1}{2}}a^{}_a^{}\right)M_{040}^{}`$ (132)
$`+(\beta +2\zeta ^{})_a^{}M_{040}^{}=0`$
In the steady state limit, Eq. (52) applies and the first term on the left hand side vanishes. In this case, one easily gets
$$\left(\frac{}{a^{}}M_{040}^{}\right)_s=\frac{4\chi _s}{a_s^{}\chi _s+2\beta +4\zeta ^{}}M_{040,s}^{},$$
(133)
where
$$\chi _s=\frac{2}{3}\left[P_{xy,s}^{}+a_s^{}\left(\frac{P_{xy}^{}}{a^{}}\right)_s\right]$$
(134)
is a known function and
$$M_{040,s}^{}=\frac{3}{4}\frac{\beta }{\beta +2\zeta ^{}}.$$
(135)
Proceeding in a similar way, all the derivatives of the form $`_a^{}M^{}`$ can be analytically computed in the steady state.
## Appendix D Kinetic model results in the steady state
In this Appendix, I display the results obtained from the model kinetic equation chosen here for the determination of the generalized transport coefficients. In the model, the Boltzmann collision operator is replaced by the term BDS99
$$J[f,f]\beta \nu (ff_0)+\frac{\zeta }{2}\frac{}{𝐜}\left(𝐜f\right),$$
(136)
where $`\nu `$ and $`\beta `$ are given by Eqs. (21) and (23), respectively, $`f_0`$ is the local equilibrium distribution (19) and $`\zeta `$ is the cooling rate (12).
### D.1 Steady state solution for the (unperturbed) USF
Let us consider first the steady state solution to the (unperturbed) USF problem. In this case, $`\delta 𝐮=\mathrm{𝟎}`$ and so $`𝐜=𝐕`$. The one-particle distribution function $`f(𝐕)`$ obeys the kinetic equation
$$aV_y\frac{}{V_x}f(𝐕)=\beta \nu (ff_0)+\frac{\zeta ^{(0)}}{2}\frac{}{𝐕}\left(𝐕f\right),$$
(137)
where here $`\zeta `$ has been approximated by its local equilibrium approximation $`\zeta ^{(0)}`$ given by Eq. (22). The main advantage of using a kinetic model instead of the Boltzmann equation is that the model lends itself to an exact solution GS03 ; AS05 . It can be written as
$$f(𝐕)=n\left(\frac{m}{2T}\right)^{3/2}f^{}(𝝃),𝝃=\sqrt{\frac{m}{2T}}𝐕,$$
(138)
where the reduced velocity distribution function $`f^{}`$ is a function of the coefficient of restitution $`\alpha `$ and the reduced peculiar velocity $`𝝃`$:
$$f^{}(𝝃)=\pi ^{3/2}_0^{\mathrm{}}𝑑se^{(1\frac{3}{2}\overline{\zeta })s}\mathrm{exp}\left[e^{\overline{\zeta }s}\left(𝝃+s\overline{𝖺}𝝃\right)^2\right].$$
(139)
Here, $`\overline{𝖺}=𝖺/(\nu \beta )`$ and $`\overline{\zeta }=\zeta ^{(0)}/(\nu \beta )`$. It has been recently shown that the distribution function (139) presents an excellent agreement with Monte Carlo simulations in the region of thermal velocities, even for strong dissipation AS05 .
The explicitly knowledge of the velocity distribution function allows one to compute all the velocity moments. We introduce the moments
$$M_{k_1,k_2,k_3}=𝑑𝐯V_x^{k_1}V_y^{k_2}V_z^{k_3}f(𝐕)$$
(140)
According to the symmetry of the USF distribution (139), the only nonvanishing moments correspond to even values of $`k_1+k_2`$ and $`k_3`$. In this case, after some algebra, one gets AS05
$$M_{k_1,k_2,k_3}=n\left(\frac{2T}{m}\right)^{k/2}M_{k_1,k_2,k_3}^{},$$
(141)
where the reduced moments $`M_{k_1,k_2,k_3}^{}`$ are given by
$`M_{k_1,k_2,k_3}^{}`$ $`=`$ $`\pi ^{3/2}{\displaystyle \underset{\stackrel{q=0}{q+k_1=\text{even}}}{\overset{k_1}{}}}(\overline{a})^q\left(1+{\displaystyle \frac{\overline{\zeta }}{2}}k\right)^{(1+q)}{\displaystyle \frac{k_1!}{(k_1q)!}}`$ (142)
$`\times \mathrm{\Gamma }\left({\displaystyle \frac{k_1q+1}{2}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{k_2+q+1}{2}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{k_3+1}{2}}\right),`$
with $`\overline{a}=a/(\nu \beta )=a^{}/\beta `$. It is easy to see that the expressions for the second degree-degree velocity moments (rheological properties) coincide with those given from the Boltzmann equation by using Grad’s approximation, Eqs. (24)–(25).
### D.2 Transport coefficients
Let us now evaluate the generalized transport coefficients $`\eta _{ijk\mathrm{}},\kappa _{ij}`$, and $`\mu _{ij}`$ in the steady state. They can be obtained from Eqs. (55)–(57) with the replacement given by Eqs. (61) and (62). With these changes, Eqs. (55)–(57) become
$$\left(ac_y\frac{}{c_x}+\nu \beta \frac{\zeta ^{(0)}}{2}\frac{}{𝐜}𝐜\right)X_{n,i}+\frac{2a}{3}\frac{T}{n}(P_{xy}^{}+a^{}_a^{}P_{xy}^{})X_{T,i}=Y_{n,i},$$
(143)
$$\left(ac_y\frac{}{c_x}\frac{1}{3}a\left(P_{xy}^{}a^{}_a^{}P_{xy}^{}\right)+\right)X_{T,i}=Y_{T,i},$$
(144)
$$\left(ac_y\frac{}{c_x}+\nu \beta \frac{\zeta ^{(0)}}{2}\frac{}{𝐜}𝐜\right)X_{u,j\mathrm{}}\frac{1}{2}\zeta _{u,j\mathrm{}}\left[\frac{}{𝐜}(𝐜f^{(0)})+2T\frac{}{T}f^{(0)}\right]a\delta _{jy}X_{u,x\mathrm{}}=Y_{u,j\mathrm{}}.$$
(145)
In order to get the transport coefficients $`\kappa _{ij}`$, $`\mu _{ij}`$, and $`\eta _{ijk\mathrm{}}`$, it is convenient to introduce the velocity moments
$$A_{k_1,k_2,k_3}^{(i)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}X_{n,i},$$
(146)
$$B_{k_1,k_2,k_3}^{(i)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}X_{T,i},$$
(147)
$$C_{k_1,k_2,k_3}^{(ij)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}X_{u,ij}.$$
(148)
The knowledge of these moments allows one to get all the transport coefficients of the perturbed USF problem. Now, we multiply Eqs. (143)–(145) by $`c_x^{k_1}c_y^{k_2}c_z^{k_3}`$ and integrate over velocity. The result is
$$ak_1A_{k_11,k_2+1,k_3}^{(i)}+\left(\nu \beta +\frac{1}{2}k\zeta ^{(0)}\right)A_{k_1,k_2,k_3}^{(i)}+\omega _nB_{k_1,k_2,k_3}^{(i)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{n,i},$$
(149)
$$ak_1B_{k_11,k_2+1,k_3}^{(i)}+\left(\nu \beta +\frac{1}{2}k\zeta ^{(0)}+\omega _T\right)B_{k_1,k_2,k_3}^{(i)}=𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{T,i},$$
(150)
$`ak_1C_{k_11,k_2+1,k_3}^{(j\mathrm{})}+\left(\nu \beta +{\displaystyle \frac{1}{2}}k\zeta ^{(0)}\right)`$ $`C_{k_1,k_2,k_3}^{(j\mathrm{})}+{\displaystyle \frac{1}{2}}\zeta _{u,j\mathrm{}}\left(k2T_T\right)M_{k_1,k_2,k_3}^{(0)}`$ (151)
$`a\delta _{jy}C_{k_1,k_2,k_3}^{(x\mathrm{})}={\displaystyle 𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{u,j\mathrm{}}}.`$
Here, $`M_{k_1,k_2,k_3}^{(0)}`$ are the moments of the zeroth-order distribution $`f^{(0)}`$ and we have introduced the quantities
$$\omega _n=\frac{2a}{3}\frac{T}{n}(P_{xy}^{}+a^{}_a^{}P_{xy}^{}),\omega _T=\frac{1}{3}a\left(P_{xy}^{}a^{}_a^{}P_{xy}^{}\right).$$
(152)
The right-hand side terms of Eqs. (149)–(151) can be easily evaluated with the result
$`𝒜_{k_1,k_2,k_3}^{(\mathrm{})}`$ $``$ $`{\displaystyle 𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{n,\mathrm{}}}`$ (153)
$`=`$ $`{\displaystyle \frac{}{n}}M_{k_1+\delta _\mathrm{}x,k_2+\delta _\mathrm{}y,k_3+\delta _\mathrm{}z}+{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{P_\mathrm{}j^{(0)}}{n}}`$
$`\times \left(\delta _{jx}k_1M_{k_11,k_2,k_3}+\delta _{jy}k_2M_{k_1,k_21,k_3}+\delta _{jz}k_3M_{k_1,k_2,k_31}\right)`$
$`=`$ $`\left({\displaystyle \frac{2T}{m}}\right)^{\frac{k+1}{2}}[(1a^{}_a^{})M_{k_1+\delta _\mathrm{}x,k_2+\delta _\mathrm{}y,k_3+\delta _\mathrm{}z}^{}{\displaystyle \frac{1}{2}}(1a^{}_a^{})P_\mathrm{}j^{}`$
$`\times (\delta _{jx}k_1M_{k_11,k_2,k_3}^{}+\delta _{jy}k_2M_{k_1,k_21,k_3}^{}+\delta _{jz}k_3M_{k_1,k_2,k_31}^{})],`$
$`_{k_1,k_2,k_3}^{(\mathrm{})}`$ $``$ $`{\displaystyle 𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{T,\mathrm{}}}`$ (154)
$`=`$ $`{\displaystyle \frac{}{T}}M_{k_1+\delta _\mathrm{}x,k_2+\delta _\mathrm{}y,k_3+\delta _\mathrm{}z}+{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{P_\mathrm{}j^{(0)}}{T}}`$
$`\times \left(\delta _{jx}k_1M_{k_11,k_2,k_3}+\delta _{jy}k_2M_{k_1,k_21,k_3}+\delta _{jz}k_3M_{k_1,k_2,k_31}\right)`$
$`=`$ $`n\left({\displaystyle \frac{2T}{m}}\right)^{\frac{k+1}{2}}[{\displaystyle \frac{1}{2T}}(k+1a^{}_a^{})M_{k_1+\delta _\mathrm{}x,k_2+\delta _\mathrm{}y,k_3+\delta _\mathrm{}z}^{}{\displaystyle \frac{1}{2T}}(1{\displaystyle \frac{1}{2}}a^{}_a^{})P_\mathrm{}j^{}`$
$`\times (\delta _{jx}k_1M_{k_11,k_2,k_3}^{}+\delta _{jy}k_2M_{k_1,k_21,k_3}^{}+\delta _{jz}k_3M_{k_1,k_2,k_31}^{})],`$
$`𝒞_{k_1,k_2,k_3}^{(j\mathrm{})}`$ $``$ $`{\displaystyle 𝑑𝐜c_x^{k_1}c_y^{k_2}c_z^{k_3}Y_{u,j\mathrm{}}}`$ (155)
$`=`$ $`\delta _j\mathrm{}\left(1n{\displaystyle \frac{}{n}}\right)M_{k_1,k_2,k_3}+{\displaystyle \frac{2}{3n}}\left(P_j\mathrm{}^{(0)}a\eta _{xyj\mathrm{}}\right){\displaystyle \frac{}{T}}M_{k_1,k_2,k_3}`$
$`M_{k_1,k_2,k_3}\left(\delta _{jx}\delta _\mathrm{}xk_1+\delta _{jy}\delta _\mathrm{}yk_2+\delta _{jz}\delta _\mathrm{}zk_3\right)`$
$`k_1\delta _{jx}\left(\delta _\mathrm{}yM_{k_11,k_2+1,k_3}+\delta _\mathrm{}zM_{k_11,k_2,k_3+1}\right)`$
$`k_2\delta _{jy}\left(\delta _\mathrm{}xM_{k_1+1,k_21,k_3}+\delta _\mathrm{}zM_{k_1,k_21,k_3+1}\right)`$
$`k_3\delta _{jz}\left(\delta _\mathrm{}xM_{k_1+1,k_2,k_31}+\delta _\mathrm{}yM_{k_1,k_2+1,k_31}\right)`$
$`=`$ $`n\left({\displaystyle \frac{2T}{m}}\right)^{k/2}[\delta _j\mathrm{}a^{}_a^{}M_{k_1,k_2,k_3}^{}`$
$`{\displaystyle \frac{1}{3nT}}\left(P_j\mathrm{}^{(0)}a\eta _{xyj\mathrm{}}\right)(ka^{}_a^{})M_{k_1,k_2,k_3}^{}`$
$`+M_{k_1,k_2,k_3}^{}\left(\delta _{jx}\delta _\mathrm{}xk_1+\delta _{jy}\delta _\mathrm{}yk_2+\delta _{jz}\delta _\mathrm{}zk_3\right)`$
$`+k_1\delta _{jx}\left(\delta _\mathrm{}yM_{k_11,k_2+1,k_3}^{}+\delta _\mathrm{}zM_{k_11,k_2,k_3+1}^{}\right)`$
$`+k_2\delta _{jy}\left(\delta _\mathrm{}xM_{k_1+1,k_21,k_3}^{}+\delta _\mathrm{}zM_{k_1,k_21,k_3+1}^{}\right)`$
$`+k_3\delta _{jz}(\delta _\mathrm{}xM_{k_1+1,k_2,k_31}^{}+\delta _\mathrm{}yM_{k_1,k_2+1,k_31}^{})].`$
Here, $`M_{k_1,k_2,k_3}^{}`$ are the reduced moments of the distribution $`f^{(0)}`$ defined by Eq. (141). In the steady state, $`M_{k_1,k_2,k_3}^{}`$ is given by Eq. (142) while the derivatives $`_a^{}M_{k_1,k_2,k_3}^{}`$ can be obtained by following the procedure described in Appendix C.
The solution to Eqs. (149)–(151) can be written as
$$A_{k_1,k_2,k_3}^{(i)}=(\nu \beta )^1\underset{q=0}{\overset{k_1}{}}(\overline{a})^q\left(1+\frac{k\overline{\zeta }}{2}\right)^{(1+q)}\frac{k_1!}{(k_1q)!}\left[𝒜_{k_1q,k_2+q,k_3}^{(i)}\omega _nB_{k_1q,k_2+q,k_3}^{(i)}\right],$$
(156)
$$B_{k_1,k_2,k_3}^{(i)}=(\nu \beta )^1\underset{q=0}{\overset{k_1}{}}(\overline{a})^q\left(1+\overline{\omega }_T+\frac{k\overline{\zeta }}{2}\right)^{(1+q)}\frac{k_1!}{(k_1q)!}_{k_1q,k_2+q,k_3}^{(i)},$$
(157)
$`C_{k_1,k_2,k_3}^{(j\mathrm{})}`$ $`=`$ $`(\nu \beta )^1{\displaystyle \underset{q=0}{\overset{k_1}{}}}(\overline{a})^q\left(1+{\displaystyle \frac{k\overline{\zeta }}{2}}\right)^{(1+q)}{\displaystyle \frac{k_1!}{(k_1q)!}}`$
$`\times \left[𝒞_{k_1q,k_2+q,k_3}^{(j\mathrm{})}+a\delta _{jy}C_{k_1q,k_2+q,k_3}^{(x\mathrm{})}{\displaystyle \frac{1}{2}}n\left({\displaystyle \frac{2T}{m}}\right)^{k/2}\zeta _{u,j\mathrm{}}a^{}_a^{}M_{k_1q,k_2+q,k_3}^{}\right],`$
where $`\overline{\omega }_T=\omega _T/(\nu \beta )`$. From Eqs. (156)–(D.2) one can get the expressions for the transport coefficients $`\kappa _{ij}`$, $`\mu _{ij}`$, and $`\eta _{ijk\mathrm{}}`$ in terms of $`\beta `$, $`\overline{\zeta }`$ and $`\overline{a}`$.
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# Photon-axion conversion as a mechanism for supernova dimming: Limits from CMB spectral distortion
## I Introduction
Supernovae of type Ia (SNe Ia) at redshifts $`0.3z1.7`$ appear fainter than expected from the luminosity-redshift relation in a decelerating Universe supnovR ; supnovP ; Riess:2004nr . On the other hand, the cosmic microwave background (CMB) anisotropy and large-scale structure observations suggest that the Universe is spatially flat, with a matter density of approximately 30% of the critical density Spergel:2003cb ; Tegmark:2003ud . The “concordance model” thus implies that the Universe must be accelerating today because it is dominated by a “dark energy” component (about 70% of the critical density) with an equation of state $`w=p/\rho 1`$.
The lack of a satisfactory fundamental explanation for this component has triggered wide-ranging theoretical investigations of more or less exotic scenarios Carroll:2003qq . Some years ago Csáki, Kaloper and Terning Csaki:2001yk (CKT I) suggested that the observed achromatic dimming of distant SNe Ia may be a consequence of the mixing of photons with very light and weakly coupled axion-like particles in the intergalactic magnetic fields. Though still requiring some non-standard fluid (e.g. with $`p/\rho 1/3`$) to fit the flatness of the universe, this model seemed capable to explain the SN dimming through a completely different mechanism without apparently affecting other cosmological observations.
Later it was recognized that the conclusions of CKT I can be significantly modified when the effects of the intergalactic plasma on the photon-axion oscillations are taken into account Deffayet:2001pc . Assuming an electron density $`n_en_{\mathrm{baryons}}=n_\gamma \eta 10^7\mathrm{cm}^3`$, the model is ruled out in most of the parameter space because of either an excessive photon conversion or a chromaticity of the dimming. Only fine-tuned parameters for the statistical properties of the extragalactic magnetic fields would still allow this explanation. On the other hand, Csáki, Kaloper and Terning Csaki:2001jk (CKT II) criticized the assumed value of $`n_e`$ as being far too large for most of the intergalactic space, invoking observational hints for a value at least one order of magnitude smaller. For $`n_e2.5\times 10^8`$ cm<sup>-3</sup>, the photon-axion mixing hypothesis works even better with the plasma, because the constraints from CMB anisotropies via photon-axion conversion can be relaxed.
If distant SNe Ia are dimmed by this mechanism, the same would apply to other sources. In particular, one would expect a dispersion in the observed quasar (QSO) spectra. An analysis based on the first data release of the Sloan Digital Sky Survey excludes a large part of the parameter space Ostman:2004eh , suggesting that only for $`n_e10^{10}`$ cm<sup>-3</sup> the axion mechanism is still able to explain a dimming by $`0.1`$ magnitudes or more. If the QSO spectra had an intrinsic dispersion at the 5% level would rule out axion dimming exceeding $``$ 0.05 mag. Future data will be sensitive to yet larger regions in parameter space, yet QSOs will never be sensitive to very low $`n_e`$.
A similar bound has been obtained by a possible violation of the reciprocity relation between the luminosity distance and the angular-diameter distance bakuI ; bakuII . However, this constraint is less robust than the QSO one because it is affected by possibly large systematic errors that are difficult to quantify uam .
The purpose of our paper is to further constrain the photon-axion conversion model by studying its effect on the CMB spectral shape. We will show that the low-$`n_e`$ region of parameters left open by the QSO limit is ruled out by our new limit, leaving little if any room for the axion hypothesis to mimic cosmic acceleration.
In Sec. II we discuss the formalism of photon-axion conversion and in Sec. III we summarize its effect on SN Ia dimming. In Sec. IV we describe the constraints coming from spectral CMB distortions and in Sec. V we combine our new limits with those from QSO dispersion. Finally, in Sec. VI we draw our conclusions and comment on the viability of the photon-axion conversion mechanism.
## II Photon-axion conversion
Axions and photons oscillate into each other in an external magnetic field sikivie ; Raffelt:1987im ; Anselm:1987vj ; Raffeltbook due to the interaction term
$$_{a\gamma }=\frac{1}{4}g_{a\gamma }F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }a=g_{a\gamma }𝐄𝐁a,$$
(1)
where $`F_{\mu \nu }`$ is the electromagnetic field tensor, $`\stackrel{~}{F}_{\mu \nu }=\frac{1}{2}ϵ_{\mu \nu \rho \sigma }F^{\rho \sigma }`$ is its dual, $`a`$ is the axion field, and $`g_{a\gamma }`$ is the axion-photon coupling (with dimension of inverse energy). We always use natural units with $`\mathrm{}=c=k_\mathrm{B}=1`$. For very relativistic axions, the equations of motion in the presence of an external magnetic field $`B`$ reduce to the linearized form Raffelt:1987im
$$\left(\omega i_z+\right)\left(\begin{array}{ccc}A_x& & \\ A_y& & \\ a& & \end{array}\right)=0,$$
(2)
where $`z`$ is the direction of propagation, $`A_x`$ and $`A_y`$ correspond to the two linear polarization states of the photon field, and $`\omega `$ is the photon or axion energy. The mixing matrix is
$$=\left(\begin{array}{ccc}\mathrm{\Delta }_{xx}& \mathrm{\Delta }_{xy}& \frac{1}{2}g_{a\gamma }B_x\\ \mathrm{\Delta }_{yx}& \mathrm{\Delta }_{yy}& \frac{1}{2}g_{a\gamma }B_y\\ \frac{1}{2}g_{a\gamma }B_x& \frac{1}{2}g_{a\gamma }B_y& \mathrm{\Delta }_a\end{array}\right),$$
(3)
where $`\mathrm{\Delta }_a=m_a^2/2\omega `$. The component of $`𝐁`$ parallel to the direction of motion does not induce photon-axion mixing. The quantities $`\mathrm{\Delta }_{ij}`$ with $`i,j=x,y`$ mix the photon polarization states. They are energy dependent and are determined both by the properties of the medium and the QED vacuum polarization effect. We ignore the latter, being sub-dominant for the problem at hand Deffayet:2001pc .
For a homogeneous magnetic field we may choose a coordinate system aligned with the field direction. The linear photon polarization state parallel to the transverse field direction $`𝐁_T`$ is denoted as $`A_{}`$ and the orthogonal one as $`A_{}`$. Equation (2) becomes then
$$\left(\omega i_z+\right)\left(\begin{array}{ccc}A_{}& & \\ A_{}& & \\ a& & \end{array}\right)=0,$$
(4)
with mixing matrix
$$=\left(\begin{array}{ccc}\mathrm{\Delta }_{}& \mathrm{\Delta }_\mathrm{R}& 0\\ \mathrm{\Delta }_R& \mathrm{\Delta }_{}& \mathrm{\Delta }_{a\gamma }\\ 0& \mathrm{\Delta }_{a\gamma }& \mathrm{\Delta }_a\end{array}\right).$$
(5)
Here, $`\mathrm{\Delta }_{}=\mathrm{\Delta }_{\mathrm{pl}}+\mathrm{\Delta }_{}^{\mathrm{CM}}`$, $`\mathrm{\Delta }_{}=\mathrm{\Delta }_{\mathrm{pl}}+\mathrm{\Delta }_{}^{\mathrm{CM}}`$, $`\mathrm{\Delta }_{\mathrm{pl}}=\omega _{\mathrm{pl}}^2/2\omega `$, $`\mathrm{\Delta }_{a\gamma }=g_{a\gamma }|𝐁_T|/2`$, and $`\omega _{\mathrm{pl}}^2=4\pi \alpha n_e/m_e`$ defines the plasma frequency, $`m_e`$ being the electron mass and $`\alpha `$ the fine-structure constant. The $`\mathrm{\Delta }_,^{\mathrm{CM}}`$ terms describe the Cotton-Mouton effect, i.e. the birefringence of fluids in the presence of a transverse magnetic field where $`|\mathrm{\Delta }_{}^{\mathrm{CM}}\mathrm{\Delta }_{}^{\mathrm{CM}}|B_T^2`$. These terms are of little importance for the following arguments and will thus be neglected. The Faraday rotation term $`\mathrm{\Delta }_\mathrm{R}`$, which depends on the energy and the longitudinal component $`B_z`$, couples the modes $`A_{}`$ and $`A_{}`$. While Faraday rotation is important when analyzing polarized sources of photons, it plays no role for the problem at hand.
With this simplification the $`A_{}`$ component decouples, and the propagation equations reduce to a 2-dimensional mixing problem with a purely transverse field $`𝐁=𝐁_T`$
$$\left(\omega i_z+_2\right)\left(\begin{array}{cc}A_{}& \\ a& \end{array}\right)=0,$$
(6)
with a 2-dimensional mixing matrix
$$_2=\left(\begin{array}{cc}\mathrm{\Delta }_{\mathrm{pl}}& \mathrm{\Delta }_{a\gamma }\\ \mathrm{\Delta }_{a\gamma }& \mathrm{\Delta }_a\end{array}\right).$$
(7)
The solution follows from diagonalization through the rotation angle
$$\vartheta =\frac{1}{2}\mathrm{arctan}\left(\frac{2\mathrm{\Delta }_{a\gamma }}{\mathrm{\Delta }_{\mathrm{pl}}\mathrm{\Delta }_a}\right).$$
(8)
In analogy to the neutrino case Kuo:1989qe , the probability for a photon emitted in the state $`A_{}`$ to convert into an axion after traveling a distance $`s`$ is
$`P_0(\gamma a)`$ $`=`$ $`\left|A_{}(0)|a(s)\right|^2`$ (9)
$`=`$ $`\mathrm{sin}^2\left(2\vartheta \right)\mathrm{sin}^2(\mathrm{\Delta }_{\mathrm{osc}}s/2)`$
$`=`$ $`\left(\mathrm{\Delta }_{a\gamma }s\right)^2{\displaystyle \frac{\mathrm{sin}^2(\mathrm{\Delta }_{\mathrm{osc}}s/2)}{(\mathrm{\Delta }_{\mathrm{osc}}s/2)^2}},`$
where the oscillation wavenumber is given by
$$\mathrm{\Delta }_{\mathrm{osc}}^2=(\mathrm{\Delta }_{\mathrm{pl}}\mathrm{\Delta }_a)^2+4\mathrm{\Delta }_{a\gamma }^2.$$
(10)
The conversion probability is energy-independent when $`2|\mathrm{\Delta }_{a\gamma }||\mathrm{\Delta }_{\mathrm{pl}}\mathrm{\Delta }_a|`$ or whenever the oscillatory term in Eq. (9) is small, i.e. $`\mathrm{\Delta }_{\mathrm{osc}}s/21`$, implying the limiting behavior $`P_0=\left(\mathrm{\Delta }_{a\gamma }s\right)^2`$.
The propagation over many random $`B`$-field domains is a truly 3-dimensional problem, because different photon polarization states play the role of $`A_{}`$ and $`A_{}`$ in different domains. This is enough to guarantee that the conversion probability over many domains is an incoherent average over magnetic field configurations and photon polarization states. The probability after travelling over a distance $`rs`$, where $`s`$ is the domain size, is Grossman:2002by
$$P_{\gamma a}(r)=\frac{1}{3}\left[1\mathrm{exp}\left(\frac{3P_0r}{2s}\right)\right],$$
(11)
with $`P_0`$ given by Eq. (9). As expected one finds that for $`r/s\mathrm{}`$ the conversion probability saturates, so that on average one third of all photons converts to axions.
## III Photon-axion conversion and Supernova dimming
To explore the effect of photon-axion conversion on SN dimming we recast the relevant physical quantities in terms of natural parameter values. The energy of optical photons is a few eV. The strength of widespread, all-pervading $`B`$-fields in the intergalactic medium must be less than a few $`10^9`$ G over coherence lengths $`s`$ crudely at the Mpc scale, according to the constraint coming from the Faraday effect of distant radio sources Kronberg:1993vk . Along a given line of sight, the number of such domains in our Hubble radius is about $`NH_0^1/s4\times 10^3`$ for $`s1`$ Mpc. The mean diffuse intergalactic plasma density is bounded by $`n_e2.7\times 10^7`$ cm<sup>-3</sup>, corresponding to the recent WMAP measurement of the baryon density Spergel:2003cb . Recent results from the CAST experiment Andriamonje:2004hi give a direct experimental bound on the axion-photon coupling of $`g_{a\gamma }1.16\times 10^{10}`$ GeV<sup>-1</sup>, comparable to the long-standing globular-cluster limit Raffeltbook . For ultra-light axions a stringent limit from the absence of $`\gamma `$-rays from SN 1987A gives $`g_{a\gamma }1\times 10^{11}`$ GeV<sup>-1</sup> Brockway:1996yr or even $`g_{a\gamma }3\times 10^{12}`$ GeV<sup>-1</sup> Grifols:1996id . Therefore, suitable numerical values of the mixing parameters are
$`{\displaystyle \frac{\mathrm{\Delta }_{a\gamma }}{\mathrm{Mpc}^1}}`$ $`=`$ $`0.15g_{10}B_{\mathrm{nG}},`$
$`{\displaystyle \frac{\mathrm{\Delta }_a}{\mathrm{Mpc}^1}}`$ $`=`$ $`7.7\times 10^{28}\left({\displaystyle \frac{m_a}{1\mathrm{eV}}}\right)^2\left({\displaystyle \frac{\omega }{1\mathrm{eV}}}\right)^1,`$
$`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{pl}}}{\mathrm{Mpc}^1}}`$ $`=`$ $`11.1\left({\displaystyle \frac{\omega }{1\mathrm{eV}}}\right)^1\left({\displaystyle \frac{n_e}{10^7\mathrm{cm}^3}}\right),`$ (12)
where we have introduced $`g_{10}=g_{a\gamma }/10^{10}`$ GeV<sup>-1</sup> and $`B_{\mathrm{nG}}`$ is the magnetic field strength in nano-Gauss.
The mixing angle defined in Eq. (8) is too small to yield a significant conversion effect for the allowed range of axion masses because $`|\mathrm{\Delta }_a||\mathrm{\Delta }_{a\gamma }|,|\mathrm{\Delta }_{\mathrm{pl}}|`$. Therefore, to ensure a sufficiently large mixing angle one has to require nearly massless pseudo-scalars, sometimes referred to as “arions” Anselm:1981aw ; Anselm:1982ip . Henceforth we will consider the pseudoscalars to be effectively massless, so that our remaining independent parameters are $`g_{10}B_{\mathrm{nG}}`$ and $`n_e`$. Note that $`m_a`$ only enters the equations via the term $`m_a^2\omega _{\mathrm{pl}}^2`$, so that for tiny but non-vanishing values of $`m_a`$, the electron density should be interpreted as $`n_{e,\mathrm{eff}}=|n_em_a^2m_e/(4\pi \alpha )|`$.
The distance relevant for SN Ia dimming is the luminosity distance $`d_L`$ at redshift $`z`$, defined by
$$d_L^2(z)=\frac{}{4\pi },$$
(13)
where $``$ is the absolute luminosity of the source and $``$ is the energy flux arriving at Earth supnovR ; supnovP . Usually the data are expressed in terms of magnitudes
$$m=M+5\mathrm{log}_{10}\left(\frac{d_L}{\mathrm{Mpc}}\right)+25,$$
(14)
where $`M`$ is the absolute magnitude, equal to the value that $`m`$ would have at $`d_L=10`$ pc. After a distance $`r`$, photon-axion conversion has reduced the number of photons emitted by the source and thus the flux $``$ to the fraction $`P_{\gamma \gamma }=1P_{\gamma a}`$. Therefore, the luminosity distance becomes
$$d_Ld_L/(P_{\gamma \gamma })^{1/2}$$
(15)
and the brightness
$$mm\frac{5}{2}\mathrm{log}_{10}(P_{\gamma \gamma }).$$
(16)
Distant SNe Ia would eventually saturate ($`P_{\gamma \gamma }=2/3`$), and hence they would appear $`(3/2)^{1/2}`$ times farther away than they really are. This corresponds to a maximum dimming of approximately 0.4 mag.
In Fig. 1 we show qualitatively the regions of $`n_e`$ and $`g_{10}B_{\mathrm{nG}}`$ relevant for SN dimming at cosmological distances. To this end we show iso-dimming contours obtained from Eq. (16) for a photon energy 4.0 eV and a magnetic domain size $`s=1`$ Mpc. For simplicity we neglect the redshift evolution of the intergalactic magnetic field $`B`$, domain size $`s`$, plasma density $`n_e`$ and photon frequency $`\omega `$. Our iso-dimming curves are intended to illustrate the regions where the photon-axion conversion could be relevant. In reality, the dimming should be a more complicated function since the intergalactic medium is expected to be very irregular: there could be voids of low $`n_e`$ density, but there will also be high density clumps, sheets and filaments and these will typically have higher $`B`$ fields as well. However, the simplifications used in this work are consistent with the ones adopted in CKT II model and do not alter our main results.
The iso-dimming contours are horizontal in the low-$`n_e`$ and low-$`g_{10}B_{\mathrm{nG}}`$ region. They are horizontal for any $`g_{10}B_{\mathrm{nG}}`$ when $`n_e`$ is sufficiently low. From the discussion in Sec. II we know that the single-domain probability $`P_0`$ of Eq. (9) is indeed energy independent when $`|\mathrm{\Delta }_{\mathrm{osc}}s|1`$, i.e. for $`|\mathrm{\Delta }_{\mathrm{pl}}|s/21`$ and $`|\mathrm{\Delta }_{a\gamma }|s1`$. When $`n_e\mathrm{few}10^8`$ cm<sup>-3</sup> and $`g_{10}B_{\mathrm{nG}}4`$, we do not expect an oscillatory behavior of the probability. This feature is nicely reproduced in our iso-dimming contours. From Fig. 1 we also deduce that a significant amount of dimming is possible only for $`g_{10}B_{\mathrm{nG}}4\times 10^2`$.
In CKT I, where the effect of $`n_e`$ was neglected, a value $`m_a10^{16}`$ eV was used. In terms of our variables, this corresponds to $`n_{e,\mathrm{eff}}6\times 10^{12}`$ cm<sup>-3</sup>. As noted in CKT II, when plasma effects are taken into account, any value $`n_e2.5\times 10^8`$ cm<sup>-3</sup> guarantees the required achromaticity of the dimming below the 3% level between the B and V bands. The choice $`B_{\mathrm{nG}}`$ of a few and $`g_{10}0.1`$ in CKT I and II falls in the region where the observed SN dimming could be explained while being marginally compatible with the bounds on $`B`$ and $`g_{10}`$.
## IV CMB Constraints
If $`\gamma a`$ conversion over cosmological distances is responsible for the SN Ia dimming, the same phenomenon should also leave an imprint in the CMB. We note that a similar argument was previously considered for photon$``$graviton conversion Chen:1994ch . Qualitatively, in the energy-dependent region of $`P_{\gamma a}`$ one expects a rather small effect due to the low energy of CMB photons ($`\omega 10^4`$ eV). However, when accounting for the incoherent integration over many domains crossed by the photon, appreciable spectral distortions may arise in view of the accuracy of the CMB data (at the level of one part in $`10^4`$$`10^5`$). For the same reason, in the energy-independent region, at much lower values of $`n_e`$ than for the SNe Ia, the constraints on $`g_{10}B_{\mathrm{nG}}`$ are expected to be quite severe. The depletion of CMB photons in the patchy magnetic sky and its effect on the CMB anisotropy pattern have been previously considered in Csaki:2001yk . However, more stringent limits come from the distortion of the overall blackbody spectrum.
To this end we use the COBE/FIRAS data for the experimentally measured spectrum, corrected for foregrounds Fixsen:1996nj . Note that the new calibration of FIRAS Mather:1998gm is within the old errors and would not change any of our conclusions. The $`N=43`$ data points $`\mathrm{\Phi }_i^{\mathrm{exp}}`$ at different energies $`\omega _i`$ are obtained by summing the best-fit blackbody spectrum to the residuals reported in Ref. Fixsen:1996nj . The experimental errors $`\sigma _i`$ and the correlation indices $`\rho _{ij}`$ between different energies are also available. In the presence of photon-axion conversion, the original intensity of the “theoretical blackbody” at temperature $`T`$
$$\mathrm{\Phi }^0(\omega ,T)=\frac{\omega ^3}{2\pi ^2}\left[\mathrm{exp}(\omega /T)1\right]^1$$
(17)
would convert to a deformed spectrum that is given by $`\mathrm{\Phi }(\omega ,T)=\mathrm{\Phi }^0(\omega ,T)P_{\gamma \gamma }(\omega )`$. We then build the reduced chi-squared function
$$\chi _\nu ^2(T,\lambda )=\frac{1}{N1}\underset{i,j=1}{\overset{N}{}}\mathrm{\Delta }\mathrm{\Phi }_i(\sigma ^2)_{ij}^1\mathrm{\Delta }\mathrm{\Phi }_j,$$
(18)
where
$$\mathrm{\Delta }\mathrm{\Phi }_i=\mathrm{\Phi }_i^{\mathrm{exp}}\mathrm{\Phi }^0(\omega _i,T)P_{\gamma \gamma }(\omega _i,\lambda )$$
(19)
is the $`i`$-th residual, and
$$\sigma _{ij}^2=\rho _{ij}\sigma _i\sigma _j$$
(20)
is the covariance matrix. We minimize this function with respect to $`T`$ <sup>1</sup><sup>1</sup>1 In principle, one should marginalize also over the galactic foreground spectrum Fixsen:1996nj . However, we neglect it since it is a subleading effect with respect to the deformation induced on the CMB blackbody by the photon-axion conversion. for each point in the parameter space $`\lambda =(n_e,g_{10}B_{\mathrm{nG}})`$, i.e. $`T`$ is an empirical parameter determined by the $`\chi _\nu ^2`$ minimization for each $`\lambda `$ rather than being fixed at the standard value $`T_0=2.725\pm 0.002`$ K Mather:1998gm .
In Fig. 2 we show our exclusion contour in the plane of $`n_e`$ and $`g_{10}B_{\mathrm{nG}}`$. The region above the continuous curve is the excluded region at 95% C.L., i.e. in this region the chance probability to get larger values of $`\chi _\nu ^2`$ is lower than 5%. We also show the corresponding 99% C.L. contour which is very close to the 95% contour so that another regression method and/or exclusion criterion would not change the results very much. Within a factor of a few, the same contours also hold if one varies the domain size $`s`$ within a factor 10.
Comparing our exclusion plot with the iso-dimming curves of Fig. 1 we conclude that the entire region $`n_e10^9`$ cm<sup>-3</sup> is excluded for SN dimming.
A few comments are in order. Intergalactic magnetic fields probably are a relatively recent phenomenon in the cosmic history, arising only at redshifts of a few. As a first approximation we have then considered the photon-axion conversion as happening on present ($`z=0`$) CMB photons. Since $`P_{\gamma \gamma }`$ is an increasing function of the photon energy $`\omega `$, our approach leads to conservative limits. Moreover, we assumed no correlation between $`n_e`$ and the intergalactic magnetic field strength. It is however physically expected that the fields are positively correlated with the plasma density so that relatively high values of $`g_{10}B_{\mathrm{nG}}`$ should be more likely when $`n_e`$ is larger. Our constraints in the region of $`n_e10^{10}`$ cm<sup>-3</sup> are thus probably tighter than what naively appears.
## V QSO Constraints
Our limits are nicely complementary to the ones obtained from the effects of photon-axion conversion on quasar colors and spectra Ostman:2004eh . In Fig. 3 we superimpose our CMB exclusion contours with the schematic region excluded by quasars <sup>2</sup><sup>2</sup>2We use the exclusion regions of astro-ph/0410501v1. In the published version Ostman:2004eh , corresponding to astro-ph/0410501v2, the iso-dimming curves were erroneously changed. The difference is that in version 1 the angle $`\alpha `$ in Eq. (3) of Ref. Ostman:2004eh that characterizes the random magnetic field direction was correctly taken in the interval 0–360 whereas in version 2 it was taken in the interval 0–90 (private communication by the authors).. The region to the right of the dot-dashed line is excluded by requiring achromaticity of SN Ia dimming Csaki:2001jk . The region inside the dashed lines is excluded by the dispersion in QSO spectra. Moreover, assuming an intrinsic dispersion of 5% in these spectra, the excluded region could be enlarged up to the dotted lines. Our CMB argument excludes the region above the solid curve at 95% C.L.
A cautionary remark is in order when combining the two constraints. As we have discussed in the previous section, our CMB limits on photon-axion conversion are model independent. Conversely, the limits placed by the QSO spectra are possibly subjected to loop holes, since they are based on a full correlation between the intergalactic electron density and the magnetic field strength, which is reasonable but not well established observationally.
## VI Conclusions
We have examined the conversion of CMB photons into very low-mass axions in the presence of intergalactic magnetic fields. The resulting CMB spectral deformation excludes a previously allowed parameter region corresponding to very low densities of the intergalactic medium. Our new limits are complementary to the ones derived from QSO dispersion which place serious constraints on the axion-photon conversion mechanism. As a result, it appears that this mechanism can hardly play a leading role for the apparent SN Ia dimming.
The axion-photon conversion hypothesis has also been advocated to explain trans-GZK cutoff events in Ultra High Energy Cosmic Rays (UHECRs) Csaki:2003ef . In principle, UHECR photons, produced in cosmological sources far away, could drastically reduce energy losses while propagating in the intergalactic medium as axions. Some of these particles would eventually convert back to photons within a few GZK radii, thus justifying the observations of extremely high energy events as well as their isotropy. While one can not rule out the possibility that some UHE “photon-like” events at energies $`E4\times 10^{19}`$ eV might be due to this mechanism, our bounds imply that it can play only a subdominant role. Moreover, photons anyway are disfavored as candidates for the majority of the UHECRs.
In summary, the CMB constraints together with previous limits suggest that the fascinating mechanism of photon-axion conversion in the intergalactic magnetic fields does not play an important role for either the phenomenon of SN Ia dimming or for UHECR propagation. A definitive verdict would probably require a common analysis of SN Ia dimming, QSO spectra, and the Faraday effect of distant radio sources, based on mutually consistent assumptions about the intergalactic matter density and its distribution, the intergalactic $`B`$-field strength and its distribution and correlation with the electron density, and the redshift evolution of these quantities. Our results show that the low-$`n_e`$ escape route from the QSO limits is definitely closed.
###### Acknowledgements.
We thank S. Hannestad, M. Kachelrieß and T. Rashba for useful discussions and L. Ostman and E. Mörtsell for clarifying some issues about their quasar bounds Ostman:2004eh . A.M. thanks G.L. Fogli for interesting discussions, and E. Lisi and D. Montanino for carefully reading the manuscript. We thank C. Csáki, N. Kaloper, and J. Terning for critical comments on an earlier version of this manuscript. We acknowledge partial support by the Deutsche Forschungsgemeinschaft under Grant No. SFB-375 and by the European Union under the Ilias project, contract No. RII3-CT-2004-506222. The work of A.M. is supported in part by the Italian “Istituto Nazionale di Fisica Nucleare” (INFN) and by the “Ministero dell’Istruzione, Università e Ricerca” (MIUR) through the “Astroparticle Physics” research project.
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# Optical and Near-Infrared spectroscopy of Nova V1494 Aquilae 1999 # 2
## 1 Introduction
Nova V1494 Aquilae 1999 # 2 was discovered by Pereira on Dec 1.875, 1999 at a visual magnitude of 6.0 (Pereira, 1999). Spectroscopic observations showed emission lines of the Balmer series, O i, Mg ii, Fe ii, all having P-Cyg profiles (Fujii 1999 ; Ayani 1999 ; Moro et al. 1999). This confirmed the object to be an “Fe-II” nova in its early stage. The nova reached a maximum brightness of 4.0 on December 3.4, 1999, followed by a rapid decline with characteristic timescales of $`t_2=6.6\pm 0.5`$ days and $`t_3=16\pm 0.5`$ days, making this a fast nova (Kiss & Thomson, 2000). The light curve showed oscillations during the transition phase, between mid-January to April 2000, which was followed by a smooth decline to 15 mag (Kiyota et al., 2004). Early spectral evolution in the optical was reported by Kiss & Thomson (2000) and Anupama, Sahu & Mayya (2001). Spectral evolution in the optical indicates likely continued mass ejection for over 195 days (Ijima & Esenoglu 2003 ; Eyres et al. 2005). The interstellar extinction to the nova has been estimated to be E(B-V) = 0.6 from the equivalent widths of the interstellar absorption components of Na i D1 and D2 (Ijima & Esenoglu, 2003). Using this result, the authors have estimated the distance to the nova as 1.6 kpc. We shall adopt all the above parameters in our discussion and calculations (see Sec. 4.5).
In this paper we present multi-epoch optical and near-infrared spectra of nova V1494 Aquilae obtained over a period of 18 months since outburst at various phases of its evolution.
## 2 Observations
### 2.1 Optical
Optical CCD spectra were obtained from the Vainu Bappu Observatory on several nights in December 1999 using the OMR spectrograph at the cassegrain focus of the 2.3m Vainu Bappu Telescope (VBT). The spectra were obtained at a resolution of 11 Å. FeAr and FeNe spectra were used for wavelength calibration. 58 Aql was used as the spectrophotometric standard star (Hamuy et al., 1994). Spectra were also obtained from the Guillermo Haro Astrophysical Observatory (GHAO) during December 4-7 using the B&C spectrograph on the 2.12m telescope, at resolutions of 10 Å and 2 Å. The nova was observed during the nebular phase on two occasions in 2001 (April 29 and May 9) from VBT. All spectra were bias subtracted and flat-field corrected in the standard manner, and the one dimensional spectra extracted using the optimal extraction method. The spectra were wavelength calibrated and corrected for instrumental response in the standard manner. All data were analysed using various tasks within IRAF.
The spectra are not corrected for telluric absorption. These features are marked in Fig. 1. The absolute flux level is accurate to 10%. On a relative scale, the emission line fluxes are generally accurate to 10%. However, weaker lines have errors upto 20% while strong lines like H$`\beta `$ are accurate to 5%.
### 2.2 Near-Infrared
Near-infrared spectra covering the early decline and transition phases were obtained with PRLNIC, an imager-spectrometer based on a 256 $`\times `$ 256 HgCdTe NICMOS3 array detector, at the cassegrain focus of the 1.2m telescope of Mt. Abu IR Observatory (MAIRO). To remove the background, spatially offset spectra were obtained by nodding the star along the slit oriented in the north-south direction. One or more spectra were obtained at each position. Nearby standard stars were observed in the same manner. HR 7953, $`\beta `$ Tau, $`\delta `$ Aql, HR 5407, $`\lambda `$ Ser and $`\alpha `$ Leo were used as standard stars on different nights. Wavelength calibration was done using the OH sky lines. Known stellar absorption features were removed from the standard star spectra and the nova spectrum was then divided by this. The result was multiplied by the spectrum of a blackbody of temperature commensurate with the spectral type of the standard star. This process removes the effects of atmospheric absorption and intrumental response in the wavelength regions common to the nova and standard star spectra; the nova spectra have slightly different spectral coverage at different epochs. The data were reduced using standard procedures available within IRAF. The fluxes are on a relative scale and accurate to around 10%.
Table 1 gives the details of both optical and near-infrared observations.
## 3 Evolution of the spectra
Nova spectra undergo different phases of evolution depending on the physical conditions of the outflowing gas and temperature of the central remnant. A nova is initially seen as an optically-thick fireball, with numerous absorption lines and hardly any emission lines. As the ejecta expands outwards and becomes optically thin with time numerous emission lines begin to appear. A high-ionisation coronal phase is reached in some novae. Eventually, the nova returns to its pre-outburst state.
### 3.1 Fireball phase
Our first spectrum was obtained on 4 December 1999, at near-maximum. This spectrum (Fig.1 – top) displays the fireball phase of the nova explosion. Many members of the Hydrogen Balmer series, right upto Balmer 21, can be seen along with several lines of O i. Signature of the outflowing optically-thick wind can be seen as P-Cyg profiles on all these lines. Numerous lines of Fe ii, N i , Mg i and C i are also present, although the P-Cyg profiles of some of them are not apparent due to blending with adjacent lines. The absorption components of some lines like the Ca ii IR triplet are much stronger than the emission. The absorption velocites lie in the range of 1000-1500 km s<sup>-1</sup>. Thus, the spectrum is that of optically thick gas flowing outwards from the central source. There possibly are some warm, dense packets of gas in the outflowing material.
### 3.2 Early decline and transition phase
In the early decline phase of novae, the P-Cygni profiles disappear and the lines acquire a more rounded emission peak This transition can be seen in our spectra obtained on 6 December 1999 (Fig. 1– bottom). Polarisation observations by Kawabata et al. (2001) during this period revealed a drastic change in the position angle of the intrinsic polarisation (from 65°to 140°). They have interpreted this in terms of an optical depth effect or a geometric change in the nova wind. A combination of both these effects are likely to be responsible for this change.
Subsequent optical spectra displayed in Fig. 2 show typical characteristics of Fe-II novae – numerous Fe ii and N ii lines – along with prominent lines of H$`\alpha `$, H$`\beta `$ and O i. The emission lines show P-Cyg profiles on the first few days and become more rounded subsequently. An optically thin wind gives a rounded emission peak, while an optically thin shell produces a more rectangular line profile (Williams 1992 and references therein). The profile of the H$`\alpha `$ line shown in Fig. 3 clearly demonstrates the prominence of the wind and shell components at different times. Initially, it shows a strong P-Cyg absorption and is representative of an optically-thick wind. Later it is characteristic of optically thin wind and eventually shows the typical saddle shape of the polar blob - equatorial ring geometry of the ejected shell (see Sec. 4.1). The hardening of the spectrum with time can be seen in the weakening of the Fe ii and N ii lines, and the continuing strength of the O i 8446 Å line. Some lines like He i (4471 Å , 5876 Å ) and \[O i\] (6300, 6363 Å ) are clearly seen in emission only on some days. The FWHM of the H$`\beta `$ line lies in the range of 2280 – 2875 km s<sup>-1</sup> in this period, with no apparent secular change.
Near-infrared spectra obtained on various days are shown in the figures 4, 5 and 6. In each of these panels, the spectra have been offset from each other and scaled for illustrative purposes. Lines of hydrogen (Paschen and Brackett series) and oxygen are prominent in the 5 December 1999 spectra. The $`J`$ band has several lines due to C i and N i. The overall spectrum is typical of an Fe-II nova in the very early decline stage. None of the lines show significant P-Cygni profiles, sugesting that the nova ejecta had become optically thin in the near-infrared as early as day 2. The increasing level of ionisation of the ejecta is reflected in the gradual decrease in the line strengths of Fe ii, C i and N i, and the appearance of various lines of He i. The strong line near 1.08 $`\mu `$m seen on 5 December 1999 is likely to be Fe ii, as is corroborated by its reduced strength two days later. The He i lines at 1.252, 1.700 and 2.058 $`\mu `$m are seen starting from respectively 28, 70 and 70 days since maximum. The He i lines are very strong in March 2000. The FWHMs of the hydrogen and helium lines lie in the range 2100 – 2800 km s<sup>-1</sup>. Again, no temporal trend can be seen.
### 3.3 Late stage nebular spectra
V1494 Aql showed fairly strong nebular lines and weak coronal lines at around 65 and 80 days respectively (Ijima & Esenoglu, 2003), during the transition phase. Infrared coronal lines seen in July 2000 (day 226), along with lines of H i, He i and He ii (Venturini et al., 2000) show that the nova ejecta had zones of low or no ionisation coeval with the highly ionised zone. Optical spectra (Arkhipova et al. 2002 ; Ijima & Esenoglu 2003) show that by September 2000 (day $``$ 280) the coronal lines had strenghtened considerably while the low ionisation lines, such as He i, had weakened significantly. Our spectra (see Figs. 7 and 8) show no He i lines, indicating that the helium has been completely ionised by day $``$ 510. Similarly, the absence of \[O i\] 6300 Å line implies that the zone where neutral oxygen was present earlier also has been ionised. Both spectra show many lines of highly ionised iron.
CHANDRA observations show that the nova had become a super-soft X-ray source by August 2000 (Starrfield et al., 2000). This is the signature of the energy emitted by hydrostatic hydrogen burning (of the unejected matter) on the surface of the white dwarf. The April 2001 spectrum (see figure 7), shows that the lines of highly ionised iron have sharp, single-peaked profiles. The nebular lines are double peaked and show similar structure as that of the Balmer lines. This suggests that the coronal lines could be arising in a region different (possibly closer to the hot white dwarf) from that of the nebular (and other) lines. This conforms to the model of a nova shell presented by Saizar & Ferland (1994) wherein the ejecta consist of warm (T $`\mathrm{\hspace{0.17em}10}^4`$ K), dense clouds embedded in a hot (T $`\mathrm{\hspace{0.17em}10}^6`$ K), tenuous gas. The two phases are photoionised by the hard white dwarf continuum and the warm phase is further photoionised by the free-free continuum generated in the hot gas. Thus, the nebular spectra require some kind of non-uniform (in temperature and density) shell structure in order to explain the observed range of ionisation.
In section 4.5 below we have derived values for the physical conditions of the ejecta in this stage.
### 3.4 Spectral evolutionary sequence
The Tololo classification system for novae has been evolved to define the temporal evolution of a nova spectrum (Williams et al. 1991, Williams et al 1994). Every nova can be assigned an evolutionary sequence according to the various phases and sub-classes for each phase that are observed in the optical spectra. In the early decline phase (early-December 1999), permitted lines of Fe II were the strongest non-Balmer lines ; hence, the spectral classification of the nova during this phase is P<sub>fe</sub>. Also, the O i 8446 Å line remained stronger than H$`\beta `$ while the nova was still in the permitted line phase. Hence a classification of P$`{}_{}{}^{\mathrm{o}}{}_{\mathrm{fe}}{}^{}`$ can be assigned. If \[Fe x\] 6375 Å emission is clearly present and stronger than \[Fe vii\] 6087 Å, the nova spectrum is considered to be in the coronal phase, regardless of any other line strengths. This is the case during April 2001. The strongest non-Balmer line in this spectrum is \[O iii\] 5007 Å, and therefore the nova is in the C<sub>o</sub> phase. Thus, the evolutionary sequence for V1494 Aql is P<sub>fe</sub>P$`{}_{}{}^{\mathrm{o}}{}_{\mathrm{fe}}{}^{}`$C<sub>o</sub>.
## 4 Discussion
### 4.1 Ejecta geometry
Nova shells are not uniform, but show considerable structure because of the clumpy nature of the ejecta (for example, Anupama & Prabhu 1993). A double-peaked, saddle-shaped emission profile is the signature of emission from a shell of material, which has an equatorial-ring, polar-cone/blob morphology (Hutchings, 1972). Substructures within the shell manifest themselves as multiply-peaked spectral lines (Gill & O’Brien, 1999). Nova V1494 Aql showed triangular lines in the initial phases (see figures 9 and 10). They broadened and acquired a more rectangular, multi-peaked profile with time. The velocities deduced from emission peaks range from 500 to 1200 kms<sup>-1</sup>, although the O i lines show higher velocity components of upto 2500 km s<sup>-1</sup>. Spectropolarimetric observations of this nova have shown that an asymmetric geometry was present even prior to maximum brightness (Kawabata et al., 2001). At around 10 days after maximum light, rapidly variable components of polarisation were observed, and their contribution increased with time. This is due to clumping in ejection near the nova. A radio image obtained using MERLIN on day 136 shows that the ejecta continued to have this clumpy structure (Eyres et al., 2005).
The line profiles (Figures 9, 10) seen in the later decline phases are similar to those seen in nova shell models having shell inclination angles between 60-90 $`\mathrm{°}`$ (Gill & O’Brien, 1999). Kiss & Thomson (2000) infer that the shell is seen edge-on whereas Eyres et al. (2005) favour a low inclination angle. It is to be noted that similar profiles may be produced by different parameters of inclination angles, ellipticities and positions of rings (Gill & O’Brien, 1999). In the absence of detailed modelling, we can only say that the nova spectra reflect the structure (equatorial ring - polar ring / cap) of the shell. As seen in the insets in figures 7 and 8, the nebular line profiles show changes in the red-blue asymmetry during the two epochs, reflecting the fact that shell-shaping is continuing even at this late stage.
### 4.2 Ly$`\beta `$ fluorescence
A striking feature of the spectra is the presence of strong lines of oxygen. The 1.128 $`\mu `$m line is much stronger than the 1.316 $`\mu `$m line on all days. Also, the 8446 Å line is stronger than would be expected from recombination alone, which should produce F<sub>8446</sub>/F$`{}_{H\alpha }{}^{}`$ 10<sup>-4</sup> (Rudy, Rossano & Puetter, 1989). The 1.128 $`\mu `$m line is predominantly produced from fluorescent excitation of O i by Ly$`\beta `$ (Bowen fluorescence). This is also the dominant mechanism for the 8446 Å line. The main implication of strong O i lines generated by Ly$`\beta `$ fluorescence is the existence of a region or regions where a sizeable fraction of the O and H is neutral but where a high Ly$`\beta `$ flux density is also present. Such regions are dense, warm, and quite optically thick in the Lyman lines.
### 4.3 Unidentified lines
Several novae have shown unidentified lines at 0.8926, 1.1110, 1.1900, 1.5545 and 2.0996 $`\mu `$ m (for example, V4633 Sgr (Lynch et al., 2001), V723 Cas (Rudy et al., 2002), CI Aql (Lynch et al., 2004)). They first appear about the same time as the He ii lines and generally persist till the coronal phase. Thus, they are of medium to high excitation and require relatively low electron densities for appearance (Lynch et al., 2004). Near-infrared spectra of V1494 Aql display at least two of these lines – 1.1110 and 1.5545 $`\mu `$ m . The Brackett 16 line appears unusually strong in the 31 December 1999 spectrum and it could be blended with the unidentified line at 1.5545 $`\mu `$ m. The shoulder on Pa$`\gamma `$ line on 11 Feb 2000 could be the 1.1110 $`\mu `$ m line. There are a few other unidentified lines which are seen only in some novae. For example, the lines at 1.6983 $`\mu `$ m (V2487 Oph – Lynch et al. 2000) and 1.770 $`\mu `$ m (V838 Her – Harrison & Stringfellow 1994). The 1.6983 $`\mu `$ m line is seen in V1494 Aql on 31 December 1999 as a blend with the Brackett 11 line, which appears unusually strong. The unidentified line at 1.770 $`\mu `$ m is seen in the 1 March 2000 spectrum.
### 4.4 Ejecta ionisation
The ejecta of V1494 Aql displayed only low excitation lines in the first month following outburst. He i emission lines were not seen in our December 1999 spectra (only the He i 1.252 $`\mu `$ m line is seen in the 31 December 1999 spectrum), indicating that the level of ionisation was less than 25 eV. The appearance and rapid rise of the He i lines in February and March 2000 (see figures 5 and 6) suggests that the ejecta was evolving to a higher excitation state in this period. By July 2000, the ejecta ionisation had reached more than 329 eV (Venturini et al., 2000). However, lines of lower ionisation were also seen, suggesting that there were still some denser, cooler zones within the clumpy ejecta. About 17 months after the outburst, such zones hardly existed and the ejecta were almost completely ionised as can be seen in figure 7. Increased emission line strengths in May 2001 (see figure 8) suggest a further hardening of the nova spectrum.
### 4.5 Physical parameters of the ejecta
Emission line fluxes corrected for extinction enable us to determine physical conditions in the ejecta. We have used the nebular stage spectra because they are ideal for abundance estimates.
Since the standard nebular lines used for density estimate are not available in our data, we use the H$`\beta `$ luminosity to estimate the electron number density, N<sub>e</sub>. The volume of the line emitting region is estimated assuming the shell to be spherical, with a filling factor of 0.01 (Ijima & Esenoglu (2003) obtained a value of 0.016 in June 2000) and uniformly expanding with a velocity of 2500 km s<sup>-1</sup>. The electron temperature, T<sub>e</sub>, in this region is assumed to be 1.5 $`\times `$ 10<sup>4</sup> K. Using these simplistic assumptions, the N<sub>e</sub> is found to be 1.1 ($`\pm `$ 0.06) $`\times `$ 10<sup>5</sup> cm<sup>-3</sup>. Only the uncertainty in flux is considered for calculating the error. As mentioned in sec. 4.1, the shell is aspherical and there could exist several velocity components. From the derived N<sub>e</sub> we estimate the mass of hydrogen in the ejecta, M<sub>H</sub>, to be 6 $`\times `$ 10<sup>-6</sup> M. The total ejecta mass would be higher than this value.
Assuming that all the helium is ionised, the helium abundance by number is estimated using the He ii 4686/H$`\beta `$ ratio. The hydrogen and He ii emissivities are from Hummer & Storey (1987). The helium abundance of V1494 Aql is found to be 0.24, and is similar to that observed in other novae such as V1425 Aql (Kamath et al., 1997).
We have used the nebular package (Shaw & Dufour, 1995) within IRAF to calculate T<sub>e</sub> in the zone of nebular lines and ionic abundances. T<sub>e</sub> as determined using \[O iii\] line fluxes is 1.0 ($`\pm `$ 0.02) $`\times `$ 10<sup>5</sup> K. The ionic abundances with respect to H<sup>+</sup> are calculated assming that both the emission lines – nebular and H$`\beta `$ – arise in regions with the same T<sub>e</sub> and N<sub>e</sub>. This assumption is not strictly true (for example, the O<sup>+</sup> ion requires T<sub>e</sub> $``$ 2$`\times `$ 10<sup>4</sup> K) but provides a first-order estimate of the abundances in the ejecta. The abundances of O<sup>+</sup>, O<sup>2+</sup> and S<sup>2+</sup> are shown in Table 4. As mentioned in Sections 3.3, 4.1 and 4.4, there is a density as well as temperature stratification in the ejecta. In the absence of detailed modelling, the values in Table 4 are considered to be representative of the conditions in the ejected material.
The early spectra show numerous lines of neutral oxygen, carbon and nitrogen. Lines such as O i 9264 Å and N i 9042 Å are produced solely by recombination. This indicates an abundance enhancement of these elements. Also, since the mechanism of Ly$`\beta `$ fluorescence is very strong and persistent for a long time, it indicates the possible presence of a region of fairly high oxygen abundance. Therefore, the CNO abundances seem to be enhanced. This is broadly in line with present nova theories (see Hernanz (2004) for a recent review).
## 5 Remarks
We have presented optical and near-infrared spectra of the fast nova V1494 Aquilae 1999 # 2 in the early decline, transition and nebular phases, covering 18 months since outburst. Based on our data and observations reported in literature, the following picture emerges.
Nova V1494 Aql was a fast nova which ejected matter asymmetrically at velocities of 1000-2500 km s<sup>-1</sup>. The ejected matter was optically thick initially, but partial thinness set in soon. The ejecta displayed low ionisation levels during the first month after outburst. Higher ionisation lines were evident after about day 65. At about 17 months after outburst, the ejecta were largely ionised and showed strong coronal lines. The clumpy nature of the ejecta was evident in polarisation observations, spectral line profiles, nebular lines and the MERLIN radio image. The nebular spectra present evidence of temperature and density stratification within the ejecta. The calculated elemental and ionic abundances in the ejecta are similar to those found in other novae.
## Acknowledgments
We thank all staff members of the respective observatories for help during the observations. Research work at Physical Research Laboratory is funded by the Department of Space, Government of India. IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc. , under cooperative agreement with the National Science Foundation.
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# Scaling behavior of fragment shapes
## Abstract
We present an experimental and theoretical study of the shape of fragments generated by explosive and impact loading of closed shells. Based on high speed imaging, we have determined the fragmentation mechanism of shells. Experiments have shown that the fragments vary from completely isotropic to highly anisotropic elongated shapes, depending on the microscopic cracking mechanism of the shell-material. Anisotropic fragments proved to have self-affine character described by a scaling exponent. The distribution of fragment shapes exhibits a power law decay. The robustness of the scaling laws is illustrated by a stochastic hierarchical model of fragmentation. Our results provide a possible improvement of the representation of fragment shapes in models of space debris.
Spacecrafts and satellites, during their mission and service time, are exposed to the danger of impact with pieces of space debris, which is a growing population of rocket bodies, non-functioning spacecrafts, rocket fuel ejecta and pieces of fragmented material accumulated during 40 years of space exploration debris\_nasa\_1997 . In order to minimize the potential hazard, objects of size larger than 10 cm are continuously tracked in space and their orbits are taken into account for space activities. Fragmentation events like on-orbit explosions of fuel containers of upper rocket stages and secondary breakups of fragments due to mutual collisions, are the main source of the proliferation of space debris, creating a large number of small fragments which cannot be tracked. For safety reasons it is essential to work out models of fragmentation, i.e. the breaking of objects into smaller pieces, which are able to predict the consequences of on-orbit explosions and impacts nasa\_model\_2001 . The NASA breakup model EVOLVE05 nasa\_model\_2001 , also implemented by other space agencies, represents the fragments in terms of their characteristic length $`L_c`$, surface-to-mass ratio $`A/m`$ and velocity $`\stackrel{}{v}`$. Other quantities like the fragment mass $`m`$ are determined from scaling relations. Model calculations are performed in a phenomenological way, i.e. based on experiments and on-orbit observations of breakup events, probability distributions of the above quantities are prescribed. Monte Carlo simulations are carried out taking into account the specific initial conditions of the event studied nasa\_model\_2001 . The orbits of fragments are determined by their velocity, however, the lifetime of the orbits is limited by the atmospheric drag which mainly depends on the shape of the fragments. The probability of impact of debris pieces with a spacecraft and the resulting damage can be calculated from their velocity, mass, size, and shape. The precision and predictive power of model calculations strongly rely on the quality of the input distributions and the validity of scaling relations used.
General studies on fragmentation phenomena mostly focused on the understanding of the mass distribution $`F(m)`$ of fragments. For bulk solids, power law fragment mass distributions $`F(m)m^\tau `$ have been obtained under widely varying conditions with universal exponents $`\tau `$ depending mainly on the spatial dimension $`d`$ turcotte\_1986 ; glassplate\_kadono\_1997 ; kun\_transition\_1999 ; astrom\_dynfrag\_2004 . Recently, we have pointed out that the fragmentation of shell-like objects, like fuel containers or rocket bodies relevant for space debris, forms an independent universality class shell\_prl ; shell\_exp\_imp , which is also supported by other studies linna .
In this paper we present a study of the shape of fragments generated by explosive and impact loading of closed shells. High speed imaging investigation of the explosion process showed that the shape of fragments is determined by the underlying cracking mechanism of shells which strongly depends on microscopic material properties: fragments of isotropic shape are obtained for materials where the branching-merging mechanism of cracks governs the breakup, while the formation of long straight cracks results in fragments of a high degree of anisotropy. Fragments of anisotropic shape proved to have self-affine character described by a scaling exponent. The distribution of fragment shapes has a power law decay with a material dependent exponent. We illustrate the robustness of the scaling laws of the shape of shell fragments by a hierarchical stochastic model. Our results suggest a possible improvement of the description of fragment shapes in phenomenological breakup models of space debris by clearly separating bulk and shell fragmentation and by taking into account the effect of the cracking mechanism of different shell materials on fragment shapes.
In order to understand how the breakup mechanism of shells determines the shape of fragments, we have carried out explosion experiments of closed shells made of brittle materials with different microscopic properties such as hen and quail egg-shells and hollow glass spheres. The egg-shells were used as a cheap and easy-to-handle brittle bio-ceramics with a highly disordered microstructure egg\_icf11 . A Photron APX ultima high speed camera with frame rate 15000/s and spatial resolution $`256\times 256`$ pixels was used to follow the time evolution of the explosion process, which also enabled us to study the dynamics of crack formation on the surface, and for the first time in the literature provided direct access to the mechanism of fragmentation. The shells were filled with a stoichiometric hydrogen-oxygen mixture which was electrically ignited approximately in the center of the shell. The analysis of the explosion of 20 egg-shells showed that the breakup process starts with the nucleation of a few cracks at the flatter end of the egg (see Fig. 1). Since the energy stored in the expanded shell at the instant of crack nucleation is high compared to the energy released by the free fracture surface, the cracks propagate at a rather high speed. The instability of the propagating crack results in sequential splitting of the crack tip at almost regular distances, triggered by the heterogeneities of the material (Fig. 1$`a`$) marder\_physrep\_1999 .
The propagating sub-branches accelerate due to the expansion of the shell and can undergo further splittings giving rise to a hierarchical tree-like crack pattern (Fig. 1$`b`$). Fragments are formed along the main cracks by the merging of adjacent side branches at almost right angles (Fig. 1$`c`$). At most $`34`$ splittings can be observed along a main branch, however, merging typically occurs at the first two levels of the hierarchy. We note that the area of fragments generated at this stage does not show large variations, it is practically determined by the inherent length scale of crack tip splitting, furthermore, the shape of fragments is more or less isotropic. The branching-merging process is initially governed by the in-plane deformation of the shell, however, as the expansion increases the out-of-plane deformation dominates giving rise to further cracks mainly perpendicular to the former ones (Fig. 1$`d`$). This cracking proceeds again in a sequential manner, typically breaking the fragments into two pieces until a stable configuration is reached.
The cracking mechanism discussed above should be generic to materials with a strongly disordered microstructure astrom\_dynfrag\_2004 . However, for shells made of materials like glass, the breaking mechanism can be significantly different glassplate\_kadono\_1997 ; glass\_crack\_nature . Experiments on exploding hollow glass spheres have shown that the breakup starts at one hot spot with random position on the surface (Fig. 2$`a`$) from which long straight cracks radiate without any apparent branching (Fig. 2$`a,b`$). This cracking mechanism results in a large number of long thin fragments having also a relatively large curvature which makes them unstable against bending. During the expansion of the sphere these primary fragments undergo a sequential breakup process due to the out-of-plane bending deformation (Fig. 2$`c,d`$). Further details on the dynamics of cracking of shells will be provided in Ref. wittel\_unpub . Egg-shells and hollow glass spheres were also fragmented by impact with a hard wall, which produced the same type of fragments shell\_exp\_imp .
In the final state of the breakup process, the fragments were carefully collected and digitized with a scanner for further evaluation. It can be observed in Figs. 1, 2 and in the inset of Fig. 3 that for the different types of materials considered, the fragments are always compact two-dimensional objects with little surface roughness, however, their overall shape can vary from completely isotropic (egg-shell) to highly anisotropic (glass) depending on the cracking mechanisms. The mass $`m`$ and surface $`A`$ of fragments is defined as the number of pixels $`N`$ and the contour length of the spots in the digital image, respectively. We characterize the linear extension of fragments by their radius of gyration as $`R_g^2=(1/N)_{ij=1}^N\left(\stackrel{}{r}_i\stackrel{}{r}_j\right)^2`$, where the sum goes over the $`N`$ pixels $`\stackrel{}{r}_i`$ of the fragments. In order to reveal how the shape of fragments varies with their size, in Fig. 3 the average fragment mass $`m`$ is presented as a function of $`R_g`$ for different materials from impact and explosion experiments. It is important to note that in all cases power law functional forms $`mR_g^\alpha `$ are obtained with a high quality, however, the exponent $`\alpha `$ depends on the structure of the crack pattern. Since the egg-shell pieces have regular isotropic shape, their mass increases with the square of $`R_g`$ and hence $`\alpha =2\pm 0.05`$ was fitted. The large glass fragments are characterized by a significantly lower value of the exponent $`\alpha =1.5\pm 0.08`$, while for small glass pieces one observes a crossover to isotropic shape with $`\alpha =2\pm 0.08`$. The value $`\alpha <2`$ implies that the fragments have self-affine character, i.e. the larger they are, the more elongated they get. Note that similar anisotropy and self-affinity of fragment shapes was not observed in $`d`$ dimensional bulk fragmentation ($`d=2,3`$) astrom\_dynfrag\_2004 ; glassplate\_kadono\_1997 ; turcotte\_1986 ; kun\_transition\_1999 .
To quantify this behavior, let us consider that the shell fragments have a rectangular shape with side lengths $`a`$ and $`b`$, hence, the surface, mass, and radius of gyration can be obtained as $`A=2(a+b)`$, $`m=ab`$, and $`R_g=\sqrt{a^2+b^2}/(2\sqrt{3})`$, respectively. The fragment mass can be expressed in terms of the aspect ratio $`r=a/b`$ and $`R_g`$ as $`mR_g^2/\left(r+1/r\right)`$, where even for moderately elongated fragments the approximation $`mR_g^2/r`$ is valid. Consequently, for fragments with exponents $`\alpha <2`$, the aspect ratio must increase as a power of $`R_g`$ so that $`rR_g^\delta `$. Hence, $`mR_g^{2\delta }`$ follows, and $`\alpha =2\delta `$, where $`\delta 1/2`$ was obtained in the experiments. It is interesting to note that the value $`\delta =1/2`$ has been found in a broad class of systems producing self-affine structures, for instance, for the scaling of the width with the length of the arms of noise reduced DLA clusters, for clusters of directed percolation or for the Hölder exponent of one-dimensional random walks dla .
Apart from the shape of individual fragments, it is also important to know the probability of occurrence of a specific fragment shape in the final state of a breakup process. The NASA breakup model characterizes the shape of fragments by the surface-to-mass ratio $`A/m`$, the distribution of which is fitted by a linear combination of Gaussian distributions nasa\_model\_2001 . The functional form of the corresponding distributions $`g(A/m)`$ of our shell pieces in the inset of Fig. 4 again shows a strong dependence on the cracking mechanism. For isotropic fragments a reasonable fit could be obtained with Gaussians in agreement with the NASA model nasa\_model\_2001 , however, for anisotropic fragments $`g`$ increases monotonically. The small sized isotropic and the large very elongated anisotropic fragments both have large $`A/m`$ value which prevents clear shape identification. To obtain a better characterization of fragment shapes, we introduce a dimensionless shape parameter $`S`$ defined as $`S=\frac{A}{m}R_g`$, multiplying the surface-to-mass ratio $`A/m`$ by the radius of gyration $`R_g`$. Assuming rectangular objects, the shape parameter $`S`$ takes the form $`S=(a+b)\sqrt{a^2+b^2}/\left(\sqrt{3}ab\right)`$. For fragments of isotropic shape $`ab`$, it follows that $`S1.63`$, which is indicated by the vertical dashed line in Fig. 4. If the fragments are elongated $`ab`$, the shape parameter $`Sa/b`$ coincides with the aspect ratio $`r`$ characterizing the degree of anisotropy. Corresponding to the cracking mechanisms, the distributions $`f(S)`$ of different materials and fragmentation modes (explosion and impact) form two groups in Fig. 4. Fragments of a low degree of anisotropy, irrespective of their size, contribute to the maximum of $`f(S)`$ in the vicinity of $`S1.63`$. Since egg-shell fragments are mostly isotropic at all sizes, the distribution $`f(S)`$ decreases rapidly over a narrow interval of $`S`$. The remarkable feature of the distribution $`f`$ is that it follows a power law decay $`f(S)S^\beta `$, where the exponent $`\beta =6.8\pm 0.3`$ is obtained for isotropic fragments over a limited scaling range, while $`\beta =3.5\pm 0.2`$ follows when the cracking mechanism favors the formation of anisotropic fragments.
Discrete element models of shell fragmentation usually consider highly disordered brittle materials and provide isotropic fragment shapes, but they have difficulties to capture microscopic mechanisms resulting in long straight cracks shell\_prl ; shell\_exp\_imp . We propose a simple stochastic binary breakup model in the spirit of Refs. turcotte\_1986 ; gonzalo\_1995 to better understand the experimental findings. The model focuses on the binary breakup of fragments formed by the primary cracking mechanism of the shell. Representing the fragments by rectangles with a continuous mass distribution at a fixed aspect ratio $`r=a/b`$, the effect of the material dependent primary cracking mechanism on the shape of fragments can be taken into account by setting $`r1`$ for fragments of isotropic shape (egg), and $`r1`$ for the highly anisotropic needle-like pieces (glass). Based on Figs. 1 and 2, these fragments are then assumed to undergo a sequential binary breakup process, where at each step of the hierarchy they break into two pieces of equal mass with a probability $`p1`$. Note that the fragments have $`1p`$ chance to keep their actual size. To capture the effect of out-of-plain deformations, we choose a side of a rectangle to break with a probability proportional to its length. Computer simulations of the model were performed starting from a continuous distribution of fragment sizes varying the initial aspect ratio $`r`$ to model different materials, while $`p`$ was fixed. The hierarchical process was followed up to $`n=30`$ generations resulting in $`10^7`$ fragments in the final state, where the mass $`m`$, the radius of gyration $`R_g`$ and the shape parameter $`S`$ of fragments were determined. It can be observed in Fig. 5 that similar to the experiments, for all values of $`r`$ the average fragment mass exhibits a power law dependence on the radius of gyration $`mR_g^\alpha `$.
Starting the process with isotropic shapes $`r1`$, the fragments remain isotropic at all levels of the hierarchy implying $`\alpha =2\pm 0.05`$ as it was observed for the egg pieces. Modelling glass fragments by a high initial anisotropy $`r1`$, the crack mostly occurs along the same side of the rectangle lowering $`r`$, however, when the two sides become comparable $`r`$ fluctuates about one. Consequently, large fragments are characterized by an exponent $`\alpha `$ significantly lower than 2, while for small pieces a crossover occurs to isotropic shape with $`\alpha =2`$. In the experiments we estimated the initial aspect ratio of glass fragments from Fig. 2 to fall in the range of $`15r35`$. Simulations with these aspect ratios $`r`$ proved to provide values of the exponent $`\alpha `$ in the vicinity of $`1.5`$. In Fig. 5 the results are presented for $`r=20`$ where $`\alpha =1.5\pm 0.06`$ was obtained . When the fragments initially have an isotropic shape, the hierarchical process gives rise to a rapidly decreasing distribution of the shape parameter $`f(S)`$ over a narrow range as it was observed for egg pieces. Starting the simulation with elongated fragments ($`r=20`$ in Fig. 5) the distribution $`f`$ shows a power law decay $`f(S)S^\beta `$ with an exponent $`\beta =3.8\pm 0.3`$ very close to the experimental value of glass fragments. It is very important to notice that these results for the exponent $`\alpha `$ and for the distribution of fragment shapes are practically independent on $`p`$.
Summarizing, based on high speed imaging techniques we have determined the fragmentation mechanism of closed shells: after the material-dependent primary cracking mechanism governed by the in-plane deformation of the shell, a hierarchical secondary breakup process sets in due to out-of-plane deformations. Contrary to bulk systems, the shape of shell fragments shows large variations from completely isotropic to highly anisotropic fragments depending on the primary cracking mechanism. We pointed out that the anisotropic fragments have a self-affine character with a scaling exponent $`\delta =1/2`$. To give a quantitative characterization of fragment shapes we proposed a shape parameter the distribution of which was found to exhibit a power law decay. A hierarchical stochastic breakup model provided quantitative agreement with the experimental findings, which demonstrates the robustness of the scaling laws of fragment shapes in shell fragmentation. The results imply that the characterization of fragment shapes in breakup models of space debris production should be improved by a clear distinction of bulk and shell fragmentation and by using scaling laws with exponents depending on the cracking mechanism of the material.
The authors are grateful to H. Klinkrad and C. Wiedeman of ESA for valuable discussions. This work was supported by the project SFB381. F. K. was supported by OTKA T049209, M041537 and by the Gy. Békési Foundation of HAS.
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# Measurement of ϕ₃
## 1 Introduction
Determinations of the Cabibbo-Kobayashi-Maskawa (CKM) $`^\mathrm{?}`$ matrix elements provide important checks on the consistency of the Standard Model and ways to search for new physics. Various methods using $`CP`$ violation in $`BDK`$ decays have been proposed $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ to measure the Unitarity Triangle angle $`\varphi _3`$. These methods are based on two key observations: neutral $`D^0`$ and $`\overline{D^0}`$ mesons can decay to a common final state, and the decay $`B^+D^{()}K^+`$ can produce neutral $`D`$ mesons of both flavors via $`\overline{b}\overline{c}u\overline{s}`$ (Fig. 1a) and $`\overline{b}\overline{u}c\overline{s}`$ (Fig. 1b) transitions, with a relative phase $`\theta _+`$ between the two interfering amplitudes that is the sum, $`\delta +\varphi _3`$, of strong and weak interaction phases. For the charge conjugate mode, the relative phase is $`\theta _{}=\delta \varphi _3`$.
The results are based on a data sample containing 275 million $`B\overline{B}`$ pairs, collected with the Belle detector $`^\mathrm{?}`$ at the KEKB asymmetric energy $`e^+e^{}`$ collider operating at the $`\mathrm{{\rm Y}}(4S)`$ resonance.
## 2 $`B^\pm D_{CP}K^\pm `$ and $`D_{CP}^{}K^\pm `$
Recent theoretical studies on $`B`$ meson dynamics have demonstrated a method to access $`\varphi _3`$ using the process $`B^{}D^{()0}K^{()}`$ $`^\mathrm{?}`$. When a $`D^0`$ is reconstructed as a $`CP`$ eigenstate, the $`bc`$ and $`bu`$ processes interfere. This interference leads to direct $`CP`$ violation. To extract $`\varphi _3`$ and assuming no $`D^0\overline{D}^0`$ mixing, some necessary observables sensitive to $`CP`$ violation are:
$`𝒜_{1,2}`$ $``$ $`{\displaystyle \frac{(B^{}D_{1,2}K^{})(B^+D_{1,2}K^+)}{(B^{}D_{1,2}K^{})+(B^+D_{1,2}K^+)}}={\displaystyle \frac{2r\mathrm{sin}\delta ^{}\mathrm{sin}\varphi _3}{1+r^2+2r\mathrm{cos}\delta ^{}\mathrm{cos}\varphi _3}}`$
$`_{1,2}`$ $``$ $`{\displaystyle \frac{R^{D_{1,2}}}{R^{D^0}}}=1+r^2+2r\mathrm{cos}\delta ^{}\mathrm{cos}\varphi _3,\delta ^{}=\{\begin{array}{cc}\delta \hfill & \text{for }D_1\hfill \\ \delta +\pi \hfill & \text{for }D_2\hfill \end{array},`$
where the ratios $`R^{D_{1,2}}`$ and $`R^{D^0}`$ are defined as
$`R^{D_{1,2}}`$ $`=`$ $`{\displaystyle \frac{(B^{}D_{1,2}K^{})+(B^+D_{1,2}K^+)}{(B^{}D_{1,2}\pi ^{})+(B^+D_{1,2}\pi ^+)}}`$
$`R^{D^0}`$ $`=`$ $`{\displaystyle \frac{(B^{}D^0K^{})+(B^+\overline{D}^0K^+)}{(B^{}D^0\pi ^{})+(B^+\overline{D}^0\pi ^+)}}`$
where $`D_1`$ and $`D_2`$ are $`CP`$-even and $`CP`$-odd eigenstates respectively. The asymmetries $`𝒜_1`$ and $`𝒜_2`$ have opposite signs. The ratio $`r`$ is defined as $`r=|A(B^{}\overline{D}^0K^{})/A(B^{}D^0K^{})|`$ and is the ratio of the two tree diagrams shown in Fig. 1 where $`\delta `$ is their strong-phase difference. The size of the ratio $`r`$ governs the magnitude of the maximum possible CP asymmetry; this ratio is suppressed by both CKM $`(0.45)`$ and color $`(0.40)`$ factors. The asymmetries and double ratios can be calculated for $`D^{}`$ in a similar manner. The analysis is described in detail elsewhere $`^\mathrm{?}`$.
Fig. 2 shows the $`\mathrm{\Delta }E`$ distributions for $`B^\pm DK^\pm `$ events. Table 1 summarizes the yields from $`\mathrm{\Delta }E`$ fit and the corresponding asymmetries with statistical errors. In the control samples, no large deviation from 0 is seen. The modes of interest are $`D_1K`$ and $`D_2K`$ where the $`B^+`$ and $`B^{}`$ events are used to calculate asymmetries and double ratios.
The final asymmetries for $`𝒜_1`$ and $`𝒜_2`$ are found do be
| $`𝒜_1`$ | = | 0.07 $`\pm `$ 0.14 (stat) $`\pm `$ 0.06 (sys) |
| --- | --- | --- |
| $`𝒜_2`$ | = | -0.11 $`\pm `$ 0.14 (stat) $`\pm `$ 0.05 (sys) |
agreeing with theoretical expectations where they should have opposite signs. The double ratios:
| $`R_1`$ | = | 0.98 $`\pm `$ 0.18 (stat) $`\pm `$ 0.10 (sys) |
| --- | --- | --- |
| $`R_2`$ | = | 1.29 $`\pm `$ 0.16 (stat) $`\pm `$ 0.08 (sys) |
Fig. 3 shows the $`\mathrm{\Delta }E`$ distributions for $`B^\pm D^{}K^\pm `$ events. Table 2 contains the yields of the distributions with statistical errors and asymmetries. The statistical significance of $`D_1^{}K`$ and $`D_2^{}K`$ signals are 5.6 and 4.5 respectively.
Asymmetries were found to be:
| $`𝒜_1`$ | = | -0.27 $`\pm `$ 0.25 (stat) $`\pm `$ 0.04 (sys) |
| --- | --- | --- |
| $`𝒜_2`$ | = | 0.26 $`\pm `$ 0.26 (stat) $`\pm `$ 0.03 (sys) |
where the systematic errors were calculated in a similar way to the $`Dh`$ case. Double ratios found are:
| $`R_1`$ | = | 1.43 $`\pm `$ 0.28 (stat) $`\pm `$ 0.06 (sys) |
| --- | --- | --- |
| $`R_2`$ | = | 0.94 $`\pm `$ 0.28 (stat) $`\pm `$ 0.06 (sys) |
In summary, the partial rate asymmetries $`𝒜_{1,2}`$ are measured for the decays $`B^\pm D_{CP}^{()}K^\pm `$ and are consistent with zero. A first observation is seen for $`D_1^{}K`$ and $`D_2^{}K`$.
## 3 Measurement of $`\varphi _3`$ with Dalitz Plot Analysis of $`𝑩^\mathbf{\pm }\mathbf{}𝑫^{\mathbf{(}\mathbf{}\mathbf{)}}𝑲^\mathbf{\pm }`$ Decay
Recently, three body final states common to $`D^0`$ and $`\overline{D^0}`$, such as $`K_S\pi ^+\pi ^{}`$ $`^\mathrm{?}`$, were suggested as promising modes for the extraction of $`\varphi _3`$. This method is based on two key observations: neutral $`D^0`$ and $`\overline{D}^0`$ mesons can decay to a common final state such as $`K_s\pi ^+\pi ^{}`$, and the decay $`B^+D^{()}K^+`$ can produce neutral $`D`$ mesons of both flavors via $`\overline{b}\overline{c}u\overline{s}`$ and $`\overline{b}\overline{u}c\overline{s}`$ transitions, where the relative phase $`\theta _+`$ between the two interfering amplitudes is the sum, $`\delta +\varphi _3`$, of strong and weak interaction phases. In the charge conjugate mode, the relative phase $`\theta _{}=\delta \varphi _3`$, so both phases can be extracted from the measurements of such $`B`$ decays and their charge conjugate modes. The phase measurement is based on the analysis of Dalitz distribution of the three body final state of the $`D^0`$ meson. The analysis is described in detail elsewhere $`^\mathrm{?}`$.
The Dalitz plots of $`D`$ decaying to $`K_s\pi ^+\pi ^{}`$, which contain information about CP violation in $`B`$ decays, are fitted for $`B^{}`$ and $`B^+`$ data sets. A combined unbinned maximum likelihood fit to the $`B^+`$ and $`B^{}`$ samples with $`r`$, $`\varphi _3`$ and $`\delta `$ as free parameters yields the following values: $`r=0.25\pm 0.07`$, $`\varphi _3=64^{}\pm 15^{}`$, $`\delta =157^{}\pm 16^{}`$ for the $`B^\pm \stackrel{~}{D}K^\pm `$ sample and $`r=0.25\pm 0.12`$, $`\varphi _3=75^{}\pm 25^{}`$, $`\delta =321^{}\pm 25^{}`$ for the $`B^\pm \stackrel{~}{D}^{}K^\pm `$ sample. The errors quoted here are obtained from the likelihood fit. These errors are a good representation of the statistical uncertainties for a Gaussian likelihood distribution, however in our case the distributions are highly non-Gaussian. In addition, the errors for the strong and weak phases depend on the values of the amplitude ratio $`r`$ (e.g. for $`r=0`$ there is no sensitivity to the phases). A more reliable estimate of the statistical uncertainties is obtained using a large number of MC pseudo-experiments as discussed below.
We use a frequentist technique to evaluate the statistical significance of the measurements. To obtain the probability density function (PDF) of the fitted parameters as a function of the true parameters, which is needed for this method, we employ a “toy” MC technique that uses a simplified MC simulation of the experiment which incorporates the same efficiencies, resolution and backgrounds as used in the data fit. This MC is used to generate several hundred experiments for a given set of $`r`$, $`\theta _+`$ and $`\theta _{}`$ values. For each simulated experiment, Dalitz plot distributions are generated with equal numbers of events as in the data, 137 and 139 events for $`B^{}`$ and $`B^+`$ decays, correspondingly, for $`B^\pm \stackrel{~}{D}K^\pm `$ mode and 34 and 35 events for $`B^{}`$ and $`B^+`$ for $`B^\pm \stackrel{~}{D}^{}K^\pm `$ mode. The simulated Dalitz plot distributions are subjected to the same fitting procedure that is applied to the data. This is repeated for different values of $`r`$, producing distributions of the fitted parameters that are used to produce a functional form of the PDFs of the reconstructed values for any set of input parameters.
The confidence regions for the pairs of parameters $`(\varphi _3,\delta )`$ and $`(\varphi _3,r)`$ are shown in Fig. 4 ($`B^\pm \stackrel{~}{D}K^\pm `$ mode) and Fig. 5 ($`B^\pm \stackrel{~}{D}^{}K^\pm `$ mode). They are the projections of the corresponding confidence regions in the three-dimensional parameter space. We show the 20%, 74% and 97% confidence level regions, which correspond to one, two, and three standard deviations for a three-dimensional Gaussian distribution.
For the final results, we use the central values that are obtained by maximizing the PDF and the statistical errors corresponding to the 20% confidence region (one standard deviation). Of the two possible solutions ($`\varphi _3`$, $`\delta `$ and $`\varphi _3+180^{}`$, $`\delta +180^{}`$) we choose the one with $`0<\varphi _3<180^{}`$. The final results are
$$r=0.21\pm 0.08\pm 0.03\pm 0.04,\varphi _3=64^{}\pm 19^{}\pm 13^{}\pm 11^{},\delta =157^{}\pm 19^{}\pm 11^{}\pm 21^{}$$
(2)
for the $`B^\pm \stackrel{~}{D}K^\pm `$ mode and
$$r=0.12_{0.11}^{+0.16}\pm 0.02\pm 0.04,\varphi _3=75^{}\pm 57^{}\pm 11^{}\pm 11^{},\delta =321^{}\pm 57^{}\pm 11^{}\pm 21^{}$$
(3)
for the $`B^\pm \stackrel{~}{D}^{}K^\pm `$ mode. The first, second, and third errors are statistical, systematic, and model dependent errors.
The significance of $`CP`$ violation is 94% for the $`B^\pm \stackrel{~}{D}K^\pm `$ sample and 38% for $`B^\pm \stackrel{~}{D}^{}K^\pm `$.
The two events samples, $`B^\pm DK^\pm `$ and $`B^\pm D^{}K^\pm `$, are combined in order to obtain a more accurate measurement of $`\varphi _3`$. The $`\varphi _3`$ result from the combined analysis is
$$\varphi _3=68^{}{}_{15^{}}{}^{+14^{}}\pm 13^{}\pm 11^{},$$
(4)
where the first error is statistical, the second is experimental systematics, and the third is model uncertainty. The two standard deviation interval including the systematic and model uncertainties is $`22^{}<\varphi _3<113^{}`$. The statistical significance of $`CP`$ violation for the combined measurement is 98%.
## 4 Study of the Suppressed Decays $`B^{}[K^+\pi ^{}]_DK^{}`$ and $`B^{}[K^+\pi ^{}]_D\pi ^{}`$
As noted by Atwood, Dunietz and Soni (ADS) $`^\mathrm{?}`$, $`CP`$ violation effects are enhanced if the final state is chosen so that the interfering amplitudes have comparable magnitudes; the archetype uses $`B^{}[K^+\pi ^{}]_DK^{}`$, where $`[K^+\pi ^{}]_D`$ indicates that the $`K^+\pi ^{}`$ pair originates from a neutral $`D`$ meson. The analysis is described in detail elsewhere $`^\mathrm{?}`$.
The ratio of branching fractions is defined as
$$R_{Dh}\frac{(B^{}D_{\mathrm{sup}}h^{})}{(B^{}D_{\mathrm{fav}}h^{})}=\frac{N_{D_{\mathrm{sup}}h^{}}/ϵ_{D_{\mathrm{sup}}h^{}}}{N_{D_{\mathrm{fav}}h^{}}/ϵ_{D_{\mathrm{fav}}h^{}}},$$
where $`N_{D_{\mathrm{sup}}h}`$ ($`N_{D_{\mathrm{fav}}h}`$) and $`ϵ_{D_{\mathrm{sup}}h^{}}`$ ($`ϵ_{D_{\mathrm{fav}}h^{}}`$) are the number of signal events and the reconstruction efficiency for the decay $`B^{}D_{\mathrm{sup}}h^{}`$ ($`B^{}D_{\mathrm{fav}}h^{}`$), and are given in Table 3.
The ratios $`R_{Dh}`$ are calculated to be
$`R_{DK}`$ $`=`$ $`(2.3_{1.4}^{+1.6}(\mathrm{stat})\pm 0.1(\mathrm{syst}))\times 10^2,`$
$`R_{D\pi }`$ $`=`$ $`(3.5_{0.9}^{+1.0}(\mathrm{stat})\pm 0.2(\mathrm{syst}))\times 10^3.`$
Since the signal for $`B^{}D_{\mathrm{sup}}K^{}`$ is not significant, we set an upper limit at the $`90\%`$ confidence level (C.L.) of $`R_{DK}<4.4\times 10^2`$.
The product branching fractions for $`B^{}D_{\mathrm{sup}}h^{}`$ are determined as
$$(B^{}D_{\mathrm{sup}}h^{})=(B^{}D_{\mathrm{fav}}h^{})\times R_{Dh},$$
and are given in Table 3. A third uncertainty arises due to the error in the branching fraction of $`B^{}D_{\mathrm{fav}}h^{}`$, which is taken from $`^\mathrm{?}`$. The uncertainties are statistics-dominated. For the $`B^{}D_{\mathrm{sup}}K^{}`$ branching fraction, we set an upper limit at the $`90\%`$ C.L. of $`(B^{}D_{\mathrm{sup}}K^{})<6.3\times 10^7`$. For $`B^{}D_{\mathrm{sup}}\pi ^{}`$, our measured branching fraction is consistent with expectation neglecting the contribution from $`B^{}\overline{D}^0\pi ^{}`$.
The ratio $`R_{DK}`$ is related to $`\varphi _3`$ by
$`R_{DK}=r_B^2+r_D^2+2r_Br_D\mathrm{cos}\varphi _3\mathrm{cos}\delta ,`$
where $`^\mathrm{?}`$
$`r_B`$ $``$ $`\left|{\displaystyle \frac{A(B^{}\overline{D}^0K^{})}{A(B^{}D^0K^{})}}\right|,\delta \delta _B+\delta _D,`$
$`r_D`$ $``$ $`\left|{\displaystyle \frac{A(D^0K^+\pi ^{})}{A(D^0K^{}\pi ^+)}}\right|=0.060\pm 0.003,`$
and $`\delta _B`$ and $`\delta _D`$ are the strong phase differences between the two $`B`$ and $`D`$ decay amplitudes, respectively. Using the above result, we obtain a limit on $`r_B`$. The least restrictive limit is obtained allowing $`\pm 1\sigma `$ variation on $`r_D`$ and assuming maximal interference ($`\varphi _3=0^{},\delta =180^{}`$ or $`\varphi _3=180^{},\delta =0^{}`$) and is found to be $`r_B<0.27`$.
We search for partial rate asymmetries $`𝒜_{Dh}`$ in $`B^{}D_{\mathrm{sup}}h^{}`$ decay, fitting the $`B^{}`$ and $`B^+`$ yields separately for each mode, where $`𝒜_{Dh}`$ is determined as
$$𝒜_{Dh}\frac{(B^{}D_{\mathrm{sup}}h^{})(B^+D_{\mathrm{sup}}h^+)}{(B^{}D_{\mathrm{sup}}h^{})+(B^+D_{\mathrm{sup}}h^+)}.$$
The peaking background for $`B^{}D_{\mathrm{sup}}K^{}`$ is subtracted assuming no $`CP`$ asymmetry. The fit results are shown in Fig. 7 and Table 4. We find
$`𝒜_{DK}`$ $`=`$ $`0.88_{0.62}^{+0.77}(\mathrm{stat})\pm 0.06(\mathrm{syst}),`$
$`𝒜_{D\pi }`$ $`=`$ $`0.30_{0.25}^{+0.29}(\mathrm{stat})\pm 0.06(\mathrm{syst}).`$
In summary, we observe $`B^{}D_{\mathrm{sup}}\pi ^{}`$ for the first time, with a significance of $`6.4\sigma `$. The size of the signal is consistent with expectation based on measured branching fractions $`^\mathrm{?}`$. The significance for $`B^{}D_{\mathrm{sup}}K^{}`$ is $`2.3\sigma `$ and we set an upper limit on the ratio of $`B`$ decay amplitudes $`r_B<0.27`$ at $`90\%`$ confidence level.
## Acknowledgments
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the NII for valuable computing and Super-SINET network support. We acknowledge support from MEXT and JSPS (Japan); ARC and DEST (Australia); NSFC (contract No. 10175071, China); DST (India); the BK21 program of MOEHRD and the CHEP SRC program of KOSEF (Korea); KBN (contract No. 2P03B 01324, Poland); MIST (Russia); MHEST (Slovenia); SNSF (Switzerland); NSC and MOE (Taiwan); and DOE (USA).
## References
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# Group Gradings on Simple Lie Algebras of Type “A”
## 1 Introduction
In this paper we describe all gradings by finite abelian groups on the matrix Lie algebras $`\mathrm{sl}(n)`$ over an algebraically closed field of characteristic zero. We introduce two types of gradings of $`\mathrm{sl}(n)`$, type I, induced from the gradings of the full matrix algebras $`M_n`$ , described in , and type II, obtained by a simple procedure from so called involution gradings of the full matrix algebras $`M_n`$, described in , and show that they exhaust all possible gradings of $`\mathrm{sl}(n)`$. A survey of known results about the gradings of simple Lie algebras can be found in \[10, Section 3.3\], including the results of V. Kac classifying cyclic gradings of all simple Lie algebras. Many examples of grading on simple Lie algebras can be found in the papers of J. Patera and coauthors (see, for example, ).
## 2 Some notation and simple facts
Let $`F`$ be an arbitrary field, $`A`$ a not necessarily associative algebra over an $`F`$ and $`G`$ a group. We say that $`A`$ is a $`G`$-graded algebra, if there is a vector space sum decomposition
$$A=\underset{gG}{}A_g,$$
(1)
such that
$$A_gA_hA_{gh}\text{ for all }g,hG.$$
(2)
A subspace $`VA`$ is called graded (or homogeneous) if $`V=_{gG}(VA_g)`$. An element $`aR`$ is called homogeneous of degree $`g`$ if $`aA_g`$. We also write $`\mathrm{deg}a=g`$. The support of the $`G`$-grading is a subset
$$\mathrm{Supp}A=\{gG|A_g0\}.$$
If $`N`$ is a normal subgroup of $`G`$ then $`A`$ naturally acquires a $`G/N`$-grading if one sets
$$A_x=\underset{gx}{}A_g\text{ for any }xG/N.$$
(3)
If $`A`$ is an associative algebra then (1) is called a Lie grading if instead of (2) one has
$$[A_g,A_h]A_{gh}\text{ for all }g,hG.$$
(4)
Suppose now that $`F`$ is of characteristic different from $`2`$. If $`A`$ is an associative algebra with involution $``$ and, in addition to (2), one has
$$(A_g)^{}A_g\text{ for all }gG.$$
(5)
then we say that (1) is an involution preserving grading or simply an involution grading. In this case, given a graded subspace $`BA`$ we set
$$B^{(+)}=\{bB|b^{}=b\},\text{ the set of symmetric elements of }B$$
(6)
and
$$B^{()}=\{bB|b^{}=b\},\text{ the set of skew-symmetric elements of }B.$$
(7)
If $`B`$ is an associative subalgebra of $`A`$ then $`B^{()}`$ is a Lie subalgebra of $`A`$, that is, with respect to $`[x,y]=xyyx`$ while $`B^{(+)}`$ is a Jordan subalgebra of $`A`$, that is, with respect to $`xy=xy+yx`$. We always have $`B=B^{()}B^{(+)}`$.
In this paper we study group gradings on simple Lie algebras. The following simple remark formally shows why in this case we may restrict ourselves to the case of abelian groups. We give a proof for completeness.
###### Lemma 1
Let $`L=_{gG}L_g`$ be a simple Lie algebra over an arbitrary field $`F`$, graded by a (possibly, noncommutative) group $`G`$. Then the support $`\mathrm{Supp}L`$ generates in $`G`$ an abelian subgroup.
Proof. First we note that $`gh=hg`$ for $`g,hG`$ as soon as $`[L_g,L_h]0`$ since $`[L_g,L_h]=[L_h,L_g]L_{gh}L_{hg}`$. We are going to generalize this property to the case of more than two factors. Namely, we will prove that inequality $`[L_{g_1},\mathrm{},L_{g_k}]0`$ implies that $`g_1,\mathrm{},g_k`$ pairwise commute.
Suppose $`k3`$ and $`[x_1,\mathrm{},x_k]0`$ for some $`x_1L_{g_1},\mathrm{},x_kL_{g_k}`$. Then $`y=[x_1,\mathrm{},x_{k1}]0`$ and by induction all $`g_1,\mathrm{},g_{k1}`$ commute. Also $`g_k`$ commutes with the product $`g_1\mathrm{}g_{k1}`$. Clearly, at least one of two products $`a=[x_1,\mathrm{},x_{k2},x_k]`$, $`b=[[x_1,\mathrm{},x_{k2}],[x_{k1},x_k]]`$ is non-zero. If $`a0`$ then $`g_k`$ commutes with all $`g_1,\mathrm{},g_{k2}`$ and with $`g_1\mathrm{}g_{k1}`$. Hence $`g_k`$ commutes with $`g_{k1}`$. Similarly, if $`b0`$ then $`g_{k1}g_k=g_kg_{k1}`$ commutes with all $`g_1,\mathrm{},g_{k2}`$. Hence $`g_k`$ commutes with all $`g_1,\mathrm{},g_{k1}`$.
Now we consider arbitrary $`g,h\mathrm{Supp}L`$. Since $`L`$ is simple and $`L_g0`$, there exist $`g_1,\mathrm{},g_k\mathrm{Supp}L`$ such that $`[L_g,L_{g_1},\mathrm{},L_{g_k}]0`$ and $`gg_1\mathrm{}g_k=h`$. Then all $`g,g_1,\mathrm{},g_k`$ commute and hence $`gh=hg`$. $`\mathrm{}`$
## 3 Two types of Lie gradings on associative algebras
If $`R`$ is an associative algebra over an arbitrary field $`F`$, graded by an abelian group $`G`$, $`L`$ a Lie subalgebra in $`R`$ with respect to the bracket operation $`[x,y]=xyyx`$ and $`L`$ is a graded subspace of $`R`$ then $`L`$ becomes a $`G`$-graded algebra with $`L_g=LR_g`$. The inclusion $`[L_g,L_h]L_{gh}`$ easily follows. Thus some gradings of a Lie algebra can be induced from the gradings of an associative algebra. In certain cases all gradings of important Lie algebras can be obtained in this way.
For example, this is the case when $`R=M_n`$, the matrix algebra of order $`n`$ over an algebraically closed field $`F`$ of characteristic 0 and $`L`$ is a Lie subalgebra of all matrices which are skew-symmetric under a symplectic involution (simple Lie algebra of the type $`C_k`$, $`n=2k`$, $`k3`$) or a Lie subalgebra of all matrices which are skew-symmetric under a transpose involution (simple Lie algebra of the type $`B_k`$, $`n=2k+1`$, $`k2`$ or simple Lie algebra of the type $`D_k`$, $`n=2k`$, $`k>4`$). This made it possible to give a complete description of all abelian gradings on simple Lie algebras of the types just mention, in an earlier paper .
Below we briefly recall the results of , where the full description of abelian group gradings on the full matrix algebra has been given.
### 3.1 Abelian Gradings on Matrix Algebras
A grading $`R=_{gG}R_g`$ on the matrix algebra $`R=M_n(F)`$ is called elementary if there exists an $`n`$-tuple $`(g_1,\mathrm{},g_n)G^n`$ such that the matrix units $`E_{ij},1i,jn`$ are homogeneous and $`E_{ij}R_gg=g_i^1g_j.`$
A grading is called fine if $`dimR_g=1`$ for any $`g\mathrm{Supp}R`$. A particular case of fine gradings is the so-called $`\epsilon `$-grading where $`\epsilon `$ is $`n^{\mathrm{th}}`$ primitive root of $`1`$. Let $`G=a_n\times b_n`$ be the direct product of two cyclic groups of order $`n`$ and
$$X_a=\left(\begin{array}{cccc}\epsilon ^{n1}& 0& \mathrm{}& 0\\ 0& \epsilon ^{n2}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1\end{array}\right),X_b=\left(\begin{array}{cccc}0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1\\ 1& 0& \mathrm{}& 0\end{array}\right).$$
(8)
Then
$$X_aX_bX_a^1=\epsilon X_b,X_a^n=X_b^n=I$$
(9)
and all $`X_a^iX_b^j,1i,jn`$, are linearly independent. Clearly, the elements $`X_a^iX_b^j,i,j=1,\mathrm{},n`$, form a basis of $`R`$ and all the products of these basis elements are uniquely defined by (9).
Now for any $`gG,g=a^ib^j`$, we set $`X_g=X_a^iX_b^j`$ and denote by $`R_g`$ a one-dimensional subspace
$$R_g=X_a^iX_b^j.$$
(10)
Then from (9) it follows that $`R=_{gG}R_g`$ is a $`G`$-grading on $`M_n(F)`$ which is called an $`\epsilon `$-grading.
Now let $`R=M_n(F)`$ be the full matrix algebra over $`F`$ graded by an abelian group $`G`$. The folowing result has been proved in \[4, Section 4, Theorems 5, 6\] and \[2, Subsection 2.2, Theorem 6\].
###### Theorem 1
Let $`F`$ be an algebraically closed field of characteristic zero. Then as a $`G`$-graded algebra $`R`$ is isomorphic to the tensor product
$$R^{(0)}R^{(1)}\mathrm{}R^{(k)}$$
where $`R^{(0)}=M_{n_0}(F)`$ has an elementary $`G`$-grading, $`\mathrm{Supp}R^{(0)}=S`$ is a finite subset of $`G`$, $`R^{(i)}=M_{n_i}(F)`$ has the $`\epsilon _i`$ grading, $`\epsilon _i`$ being a primitive $`n_i^{\mathrm{th}}`$ root of $`1`$, $`\mathrm{Supp}R^{(i)}=H_i_{n_i}\times _{n_i},i=1,\mathrm{},k`$. Also $`H=H_1\mathrm{}H_kH_1\times \mathrm{}\times H_k`$ and $`SH=\{e\}`$ in $`G`$.
Since $`L=\mathrm{sl}(n)=[R,R]`$, $`L`$ is a graded subspace of $`R=M_n`$, and the following is true.
###### Corollary 1
(Type I Gradings) Given any grading $`R=_{gG}R_g`$ by a finite abelian group $`G`$, setting $`L_g=R_gL`$ makes $`L`$ into a $`G`$-graded Lie algebra.
If $`R=M_n`$ and $`L=\mathrm{sl}(n)M_n`$ then not all gradings of $`L`$ are induced from $`R`$. For example, if $`n>2`$, $`G=_2`$, $`L_0`$ is the set of all skew-symmetric matrices under the ordinary transpose involution $`X^tX`$, $`L_1`$ the set of all symmetric matrices of trace zero. This is not induced from any $`_2`$-grading of $`M_n`$ since, according to the general theory of \[4, Section 4, Theorems 5, 6\] any such grading is elementary, that is, there are two natural numbers $`k,l`$ such that $`dimR_0=k^2+l^2`$ and $`dimR_1=2kl`$. In this case $`dimL_0=k^2+l^21=\frac{n(n1)}{2}`$, and $`dimL_1=2kl=\frac{n(n+1)}{2}1`$. Since $`k^2+l^22kl`$ we easily derive $`n2`$.
Quite a general result, as we will see in the future, dealing with the gradings of an associative algebra $`R`$ in the presence of involutions, is the following.
###### Theorem 2
Let $`R`$ be an associative algebra over a field $`F`$ of characteristic different from $`2`$, $`G`$ be a finite abelian group, $`h`$ an element of order 2 in $`G`$. If $``$ is an involution on $`R`$ and
$$R=\underset{gG}{}\stackrel{~}{R}_g$$
is an involution $`G`$-grading then a Lie grading by $`G`$ on $`R`$ can be given by
$$R_g=\stackrel{~}{R}_g^{()}\stackrel{~}{R}_{gh}^{(+)}.$$
Here, as in Section 2, $`\stackrel{~}{R}_g^{(\pm )}`$ is the set of symmetric (skew-symmetric) elements in $`\stackrel{~}{R}_g`$ with respect to $``$. $`\mathrm{}`$
The proof of this theorem is a direct verification of the relations $`[R_g,R_g^{}]R_{gg^{}}`$, for all $`g,g^{}G`$. If we consider the restriction of the grading of Theorem 2 on $`R=M_n`$ to the matrices of trace zero then we obtain the following result.
###### Corollary 2
(Type II Gradings) Let $`R=M_n`$ and suppose $`R`$ satisfies the hypotheses of Theorem 2. If also $`\mathrm{char}Fn`$ then a grading of the Lie algebra $`L=\mathrm{sl}(n)`$ will be obtained if we set
$$L_g=\{\begin{array}{cc}\stackrel{~}{R}_g^{()}\stackrel{~}{R}_{gh}^{(+)}\hfill & \text{ if }gh\hfill \\ \stackrel{~}{R}_h^{()}(\stackrel{~}{R}_e^{(+)}L)\hfill & \text{ otherwise}.\hfill \end{array}$$
Note that we need the restriction $`\mathrm{char}Fn`$ to make sure that $`R=FI[R,R]`$.
All involution gradings of $`R=M_n`$ have been described in \[5, Sections 6, 7, 8\].
### 3.2 Abelian Involution Gradings on Matrix Algebras
Throughout this subsection we assume the base field $`F`$ being algebraically closed of characteristic zero.
We first recall that any involution $``$ of $`R=M_n`$ can always be written as
$$X^{}=\mathrm{\Phi }^1(^tX)\mathrm{\Phi }$$
(11)
where $`\mathrm{\Phi }`$ is a nondegenerate matrix which is either symmetric or skew-symmetric and $`X^tX`$ is the ordinary transpose map. In the case where $`\mathrm{\Phi }`$ is symmetric we call $``$ a transpose involution. If $`\mathrm{\Phi }`$ is skew-symmetric $``$ is called a symplectic involution. Before we formulate the theorem describing involution gradings on $`M_n`$ we need three (slightly modified) lemmas from . The first one deals with certain fine involution gradings while the last two with elementary involution gradings. If $`R`$ has an involution $``$ then by $`R^{(\pm )}`$ we denote the space of symmetric (respectively skew-symmetric) matrices in $`R`$ under $``$.
###### Lemma 2
Let $`R=M_2(F)`$ be a $`2\times 2`$ matrix algebra endowed with an involution $`:RR`$ corresponding to a symmetric or skew-symmetric non-degenerate bilinear form with the matrix $`\mathrm{\Phi }`$. Then the $`(1)`$-grading of $`M_2`$ is an involution grading if and only if one of the following holds:
* $`\mathrm{\Phi }`$ is skew-symmetric,
$$\mathrm{\Phi }=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),R^{()}=\left\{\left(\begin{array}{cc}a& b\\ c& a\end{array}\right)\right\},R^{(+)}=\left\{\left(\begin{array}{cc}a& 0\\ 0& a\end{array}\right)\right\};$$
and
$$\left(\begin{array}{cc}x& y\\ z& t\end{array}\right)^{}=\left(\begin{array}{cc}t& y\\ z& x\end{array}\right);$$
* $`\mathrm{\Phi }`$ is symmetric,
$$\mathrm{\Phi }=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),R^{()}=\left\{\left(\begin{array}{cc}a& 0\\ 0& a\end{array}\right)\right\},R^{(+)}=\left\{\left(\begin{array}{cc}a& b\\ c& a\end{array}\right)\right\};$$
and
$$\left(\begin{array}{cc}x& y\\ z& t\end{array}\right)^{}=\left(\begin{array}{cc}t& y\\ z& x\end{array}\right);$$
* $`\mathrm{\Phi }`$ is symmetric,
$$\mathrm{\Phi }=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),R^{()}=\left\{\left(\begin{array}{cc}0& b\\ b& 0\end{array}\right)\right\},R^{(+)}=\left\{\left(\begin{array}{cc}a& b\\ b& c\end{array}\right)\right\};$$
and
$$\left(\begin{array}{cc}x& y\\ z& t\end{array}\right)^{}=\left(\begin{array}{cc}x& z\\ y& t\end{array}\right);$$
* $`\mathrm{\Phi }`$ is symmetric,
$$\mathrm{\Phi }=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),R^{()}=\left\{\left(\begin{array}{cc}0& b\\ b& 0\end{array}\right)\right\},R^{(+)}=\left\{\left(\begin{array}{cc}a& b\\ b& c\end{array}\right)\right\};$$
and
$$\left(\begin{array}{cc}x& y\\ z& t\end{array}\right)^{}=\left(\begin{array}{cc}x& z\\ y& t\end{array}\right).$$
The next lemma handles the case of an elementary grading compatible with an involution defined by a symmetric non-degenerate bilinear form.
###### Lemma 3
Let $`R=M_n(F)`$, $`n`$ a natural number, be a matrix algebra with involution * defined by a symmetric non-degenerate bilinear form. Let $`G`$ be an abelian group and let $`R`$ be equipped with an elementary involution $`G`$-grading defined by an $`n`$-tuple $`(g_1,\mathrm{},g_n)`$. Then $`g_1^2=\mathrm{}=g_m^2=g_{m+1}g_{m+l+1}=\mathrm{}=g_{m+l}g_{m+2l}`$ for some $`0l\frac{n}{2}`$ and $`m+2l=n`$. The involution * acts as $`X^{}=(\mathrm{\Phi }^1)^tX\mathrm{\Phi }`$ where
$$\mathrm{\Phi }=\left(\begin{array}{ccc}I_m& 0& 0\\ 0& 0& I_l\\ 0& I_l& 0\end{array}\right),$$
where $`I_s`$ is the $`s\times s`$ identity matrix. Moreover, $`R^{()}`$ consists of all matrices of the type
$$\left(\begin{array}{ccc}P& S& T\\ ^tT& A& B\\ ^tS& C& ^tA\end{array}\right),$$
(12)
where $`{}_{}{}^{t}P=P,{}_{}{}^{t}B=B,{}_{}{}^{t}C=C`$ and
$$PM_m(F),A,B,C,DM_l(F),S,TM_{m\times l}(F)$$
while $`R^{(+)}`$ consists of all matrices of the type
$$\left(\begin{array}{ccc}P& S& T\\ {}_{}{}^{t}T& A& B\\ {}_{}{}^{t}S& C& {}_{}{}^{t}A\end{array}\right),$$
(13)
where $`{}_{}{}^{t}P=P,{}_{}{}^{t}B=B,{}_{}{}^{t}C=C`$ and
$$PM_m(F),A,B,C,DM_l(F),S,TM_{m\times l}(F).$$
The last lemma deals with the case of an elementary grading compatible with an involution defined by a skew-symmetric non-degenerate bilinear form.
###### Lemma 4
Let $`R=M_n(F)`$, $`n=2k`$, be the matrix algebra with involution $``$ defined by a skew-symmetric non-degenerate bilinear form. Let $`G`$ be an abelian group and let $`R`$ be equipped with an elementary involution $`G`$-grading defined by an $`n`$-tuple $`(g_1,\mathrm{},g_n)`$. Then $`g_1g_{k+1}=\mathrm{}=g_kg_{2k}`$, the involution $``$ acts as $`X^{}=(\mathrm{\Phi }^1)^tX\mathrm{\Phi }`$ where
$$\mathrm{\Phi }=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right),$$
$`I`$ is the $`k\times k`$ identity matrix, $`R^{()}`$ consists of all matrices of the type
$$\left(\begin{array}{cc}A& B\\ C& ^tA\end{array}\right),A,B,C,M_k(F),{}_{}{}^{t}B=B,^tC=C$$
(14)
while $`R^{(+)}`$ consists of all matrices of the type
$$\left(\begin{array}{cc}A& B\\ C& {}_{}{}^{t}A\end{array}\right),A,B,C,M_k(F),{}_{}{}^{t}B=B,^tC=C.$$
(15)
We can now explicitly describe gradings on a matrix algebra with involution. Recall that in the case of an algebraically closed field any involution $``$ on a matrix algebra is, up to isomorphism, either the transpose or the symplectic involution.
###### Theorem 3
Let $`R=M_n(F)=_{gG}R_g`$ be a matrix algebra over an algebraically closed field of characteristic zero graded by the group $`G`$ and $`\mathrm{Supp}R`$ generates $`G`$. Suppose that $`:RR`$ is a graded involution. Then $`G`$ is abelian, $`n=2^km`$ and $`R`$ as a $`G`$-graded algebra with involution is isomorphic to the tensor product $`R^{(0)}R^{(1)}\mathrm{}R^{(k)}`$ where
* $`R^{(0)},\mathrm{},R^{(k)}`$ are graded subalgebras stable under the involution $``$;
* $`R^{(0)}=M_m(F)`$ is as in Lemma 4 if $``$ is symplectic on $`R^{(0)}`$ or as in Lemma 3 if $``$ is transpose on $`R_0`$;
* $`R^{(1)}\mathrm{}R^{(k)}`$ is a $`T=T_1\times \mathrm{}\times T_k`$-graded algebra and any $`R^{(i)},1ik`$, is $`T_i_2\times _2`$-graded algebra as in Lemma 2.
* A graded basis of $`R`$ is formed by the elements $`YX_{t_1}\mathrm{}X_{t_k}`$, where $`Y`$ is an element of a graded basis of $`R^{(0)}`$ and the elements $`X_{t_i}`$ are of the type (8), with $`n=2`$, $`t_iT_i`$. The involution on these elements is given canonically by
$$(YX_{t_1}\mathrm{}X_{t_k})^{}=Y^{}X_{t_1}^{}\mathrm{}X_{t_k}^{}=\mathrm{sgn}(t)(Y^{}X_{t_1}\mathrm{}X_{t_k}),$$
where $`YR^{(0)}`$, $`X_{t_i}`$ are the elements of the basis of the canonical $`(1)`$-grading of $`M_2`$, $`i=1,\mathrm{},k`$, $`t=t_1\mathrm{}t_kT`$, $`\mathrm{sgn}(t)=\pm 1`$, depending on the cases in Lemma 2. $`\mathrm{}`$
As we have seen, not every grading of $`L=\mathrm{sl}(n)R=M_n`$ is induced by a grading of $`R`$. The best way to see why this happens in the “A” case but not in the case of “B”,“ C”, “D” is to consider the connection between the gradings by an abelian group $`G`$ and the action by automorphism of the dual group $`\widehat{G}`$.
## 4 Dual group action
Let $`F`$ be a field. Denote by $`\widehat{G}`$ the dual group for $`G`$. Thus the elements of $`\widehat{G}`$ are all irreducible characters $`\chi :GF^{}`$, where $`F^{}`$ is the multiplicative group of the field $`F`$.
If $`\mathrm{\Lambda }`$ is a subset of $`\widehat{G}`$ then we denote by $`\mathrm{\Lambda }^{}`$ the subgroup of $`G`$ given by
$$\mathrm{\Lambda }^{}=\{gG|\lambda (g)=1\text{ for all }\lambda \mathrm{\Lambda }\}.$$
(16)
Similarly one defines $`S^{}`$, a subgroup in $`\widehat{G}`$, for any subset $`SG`$.
If $`G`$ is a a finite abelian group and $`\widehat{G}`$ is its group of characters over a field $`F`$ containing a primitive root of degree $`n=|G|`$ then $`\widehat{G}G`$. More precisely, if
$$G=a_1_{n_1}\times \mathrm{}\times a_j_{n_j}\times \mathrm{}\times a_k_{n_k}$$
then
$$\widehat{G}=\xi _1_{n_1}\times \mathrm{}\times \xi _i_{n_i}\times \mathrm{}\times \xi _k_{n_k}$$
where $`\xi _j(a_j)`$ is a primitive root of $`1`$ of degree $`n_j`$ and $`\xi _j(a_i)=1`$ if $`ji`$.
Obviously, $`\widehat{G}`$ separates the elements of $`G`$ in the sense that for any $`g_1g_2`$ there exists a $`\chi \widehat{G}`$ such that $`\chi (g_1)\chi (g_2)`$. Then there is a natural isomorphism between $`G`$ and $`\widehat{\widehat{G}}`$ given by $`g\psi _g`$ where $`\psi _g(\chi )=\chi (g)`$.
Notice that if $`\mathrm{\Lambda }`$ is a subgroup of $`\widehat{G}`$ and $`H=\mathrm{\Lambda }^{}`$ then there is a natural isomorphism $`\alpha :\mathrm{\Lambda }\widehat{G/H}`$ given by $`\alpha (\lambda )(gH)=\lambda (g)`$ where $`\lambda \mathrm{\Lambda }`$. Its inverse $`\beta :\widehat{G/H}\mathrm{\Lambda }`$ is given by $`\beta (\pi )(g)=\pi (gH)`$. This, in particular, shows that $`|\mathrm{\Lambda }||\mathrm{\Lambda }^{}|=|G|`$.
We briefly recall the relation between the $`G`$-gradings and $`\widehat{G}`$-actions in the case where $`G`$ is finite (see e.g. ) and $`F`$ contains a primitive root of degree $`n=|G|`$. Let $`R`$ be an arbitrary, i.e. not necessarily associative, algebra graded by a finite abelian group $`G`$, $`R=_{gG}R_g`$. Any element $`a`$ in $`R`$ can be uniquely decomposed as the sum of homogeneous components, $`a=_{gG}a_g,a_gR_g`$. Given $`\chi \widehat{G}`$ we can define
$$\chi a=\underset{gG}{}\chi (g)a_g.$$
(17)
It is easy to observe that (17) defines a $`\widehat{G}`$-action on $`R`$ by automorphisms and a subspace $`VR`$ is a graded subspace if and only if $`V`$ is invariant under this action, i.e. $`\widehat{G}V=V`$. In particular, $`a`$ is a homogeneous iff $`a`$ is a common eigenvector for all $`\chi \widehat{G}`$. The elements of the identity component $`R_e`$, $`e`$ the identity element of $`G`$, are precisely the fixed points of the above action.
We will also use the following relation between actions by the subgroups of $`\widehat{G}`$ and gradings by the factor-groups of $`G`$. Notice that thanks to the isomorphism $`\alpha :\mathrm{\Lambda }\widehat{G/H}`$ we have the action of $`\mathrm{\Lambda }`$ on $`G/H`$ given by $`\lambda (gH)=\lambda (g)`$ and on $`R`$ by $`\lambda (_{xG/H}r_x)=_{xG/H}\lambda (x)r_x`$.
###### Lemma 5
Let $`H`$ and $`\mathrm{\Lambda }`$ be subgroups of $`G`$ and $`\widehat{G}`$, respectively, such that $`H=\mathrm{\Lambda }^{}`$. Then $`aR`$ is homogeneous in the natural $`G/H`$-grading if and only if $`a`$ is a common eigenvector for all $`\chi \mathrm{\Lambda }`$. Similarly, $`VR`$ is a $`G/H`$-graded subspace if and only if $`\mathrm{\Lambda }VV`$.
Proof. In view of what was said about the connection between gradings and actions, both claim will follow if we prove that the isomorphisms $`\alpha :\mathrm{\Lambda }\widehat{G/H}`$ and $`\beta :\widehat{G/H}\mathrm{\Lambda }`$, are compatible with the restriction of the action of $`\widehat{G}`$ to $`\mathrm{\Lambda }`$ and the natural action of $`\widehat{G/H}`$ on the $`G/H`$-graded algebra $`R`$. It is sufficient to check this compatibility for one of them, say, $`\beta `$. For instance, every element is the sum of those ones which are homogeneous in the $`G/H`$-grading. If $`r`$ has degree $`gH`$ then $`r=_{hH}r_{gh}`$ and for $`\pi \widehat{G/H}`$ we should have:
$`\pi r`$ $`=`$ $`\pi (gH)r=\beta (\pi )(g)r={\displaystyle \underset{hH}{}}\beta (\pi )(g)r_{gh}`$
$`=`$ $`{\displaystyle \underset{hH}{}}\beta (\pi )(gh)r_{gh}={\displaystyle \underset{hH}{}}\beta (\pi )r_{gh}=\beta (\pi )r.`$
It follows that indeed the action of $`\mathrm{\Lambda }`$ is equivalent to the action of the whole of $`\widehat{G/H}`$, which is what we needed. $`\mathrm{}`$
Now the difference between the “A” case and those of “B”,“ C”, “D” in the case of an algebraically closed field $`F`$ of characteristic zero is that each algebra $`L`$ from the latter list, with the exception of some small rank algebras, has a canonical embedding in a matrix algebra $`R`$ in such a way that any automorphism of $`L`$ is induced by an automorphism of $`R`$. Our example above shows that this cannot be done in the case of simple Lie algebras of type “A”. Still the following result from the classical Lie Theory tells us the following.
###### Theorem 4
(\[8, Chapter IX, Theorem 5\]) The automorphism group of any Lie algebra $`L=\mathrm{sl}(n)M_n`$ over an algebraically closed field $`F`$ of characteristic zero is generated by the inner automorphisms $`XT^1XT`$, $`T`$ a nondegenerate matrix in $`M_n`$ and an outer automorphism of order 2, $`X^tX`$, where $`{}_{}{}^{t}X`$ is the transpose of $`X`$. In the case $`n=2`$ this latter mapping of $`L`$ is also induced by an inner automorphism of $`M_2`$.
Now suppose we are given a grading of $`L=\mathrm{sl}(n)`$ by a finite abelian group $`G`$. We consider the dual group $`\widehat{G}`$. Then we have the action of $`\widehat{G}`$ by the automorphisms of $`L`$. We will call a $`G`$-grading $`L=_{gG}L_g`$ on $`L=\mathrm{sl}(n)`$ inner if $`\widehat{G}`$ acts on $`L`$ by inner automorphisms. Otherwise we call this grading outer.
Suppose first that the grading is inner, that is, for each $`\chi \widehat{G}`$ there is a nondegenerate matrix $`T_\chi M_n`$ such that
$$\chi X=T_\chi ^1XT_\chi ,\text{ for any }XL.$$
(18)
It is then obvious that the same formula (18) defines an action of $`\widehat{G}`$ on $`R`$, thus a $`G`$-grading of $`R=M_n`$. This will be a unique $`G`$-grading of $`R=M_n`$ that induces the given grading of $`L=\mathrm{sl}(n)`$, and it is given by $`R_e=FIL_e`$ and $`R_g=L_g`$ for $`ge`$. The description of such gradings of $`\mathrm{sl}(n)`$ is identical to that of $`M_n`$ given in Theorem 1. For convenience, we formulate this result below in the language of Lie algebras.
###### Theorem 5
Let $`F`$ be an algebraically closed field of characteristic zero, $`G`$ a finite abelian group. If $`L=\mathrm{sl}(n)`$ and $`L=_{gG}L_g`$ is an inner $`G`$-grading then there is a grading $`R=_{gG}R_g`$ of $`R=M_n`$ described in Theorem 1 such that $`L_g=R_gL`$ for any $`gG`$. $`\mathrm{}`$
## 5 Outer gradings. General Results
Main results and arguments of the remaining sections of this paper heavily depend on , where it is assumed that the base field $`F`$ is algebraically closed of characteristic zero. Therefore, although some intermediate results and arguments could be proved under milder restrictions, we still prefer from now on to work under this assumption.
So we have to consider the case of an outer grading, namely where there is an element of $`\widehat{G}`$ which acts on $`L`$ as an outer automorphism. In this case there is a subgroup $`\mathrm{\Lambda }\widehat{G}`$ of index 2, which acts by inner automorphisms on $`L`$. An element, say $`\phi \widehat{G}\mathrm{\Lambda }`$ is the composition of an inner automorphism $`X\mathrm{\Phi }^1X\mathrm{\Phi }`$ and the canonical automorphism $`X^tX`$. Thus the action of $`\phi `$ is given by
$$\phi X=\mathrm{\Phi }^1(^tX)\mathrm{\Phi }.$$
(19)
This formula defines also an action of $`\phi `$ on the associative algebra $`R=M_n`$ by a Lie automorphism such that $`\phi I=I`$. Moreover, we now can extend the action of $`\widehat{G}`$ on $`R`$ by Lie automorphisms if we set $`\chi I=I`$ for $`\chi \mathrm{\Lambda }`$ and $`\chi I=I`$ otherwise.
Notice that in (19) the matrix $`\mathrm{\Phi }`$ is defined up to a scalar multiple. A simple but useful remark about the form of $`\mathrm{\Phi }`$ is the following.
###### Lemma 6
If one applies an inner automorphism to $`M_n`$ induced by a matrix $`C`$ then $`\mathrm{\Phi }`$ in (19) is changed to $`{}_{}{}^{t}C\mathrm{\Phi }C`$. In other words, $`\mathrm{\Phi }`$ behaves as the matrix of a bilinear form.
Proof. Indeed, we must have $`\phi \stackrel{~}{X}=\stackrel{~}{\phi X}`$ where tildes mean the matrices modified as a result of the inner automorphism in question. If we denote by $`\mathrm{\Psi }`$ the modified $`\mathrm{\Phi }`$ then we will have
$$\mathrm{\Psi }^1{}_{}{}^{t}(C^1XC)\mathrm{\Psi }=C^1\mathrm{\Phi }^1{}_{}{}^{t}X\mathrm{\Phi }C\text{ or }(^tC^1\mathrm{\Psi })^1{}_{}{}^{t}X(^tC^1\mathrm{\Psi })=(\mathrm{\Phi }C)^1{}_{}{}^{t}X(\mathrm{\Phi }C).$$
It follows that $`{}_{}{}^{t}(C^1)\mathrm{\Psi }=\alpha \mathrm{\Phi }C`$, for some coefficient $`\alpha `$ and so we may take $`\mathrm{\Psi }=^tC\mathrm{\Phi }C`$, as claimed. $`\mathrm{}`$
###### Proposition 1
One can choose the character $`\phi `$, the elements $`a,hG`$, the subgroups $`\mathrm{\Lambda }_1\widehat{G}`$ and $`KG`$ in such a way that the following hold:
1. The order of $`\phi `$ is a 2-power: $`o(\phi )=2m=2^k`$ for a natural $`k1`$;
2. $`\widehat{G}=\phi \times \mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }=\phi ^2\times \mathrm{\Lambda }_1`$;
3. $`G=a\times K`$ so that $`\phi (K)=1`$, $`\phi (a)=\rho `$, $`\rho `$ a $`(2m)^{\mathrm{th}}`$ primitive root of $`1`$, $`\mathrm{\Lambda }^{}=h`$ where $`h=a^m`$, $`\mathrm{\Lambda }_1^{}=a`$, $`\phi (h)=1`$.
Proof. This is an easy exercise, which we accomplish here for completeness. Let us denote $`\widehat{G}`$ by $`\mathrm{\Gamma }`$. We can decompose $`\mathrm{\Gamma }`$ as $`\mathrm{\Gamma }=\mathrm{\Gamma }_2\times \mathrm{\Gamma }_2^{}`$ where $`\mathrm{\Gamma }_2`$ is the Sylow $`2`$-subgroup and $`\mathrm{\Gamma }_2^{}`$ is the product of all other Sylow $`p`$-subgroups. Since $`\mathrm{\Gamma }^2\mathrm{\Lambda }`$ it easily follows that $`\mathrm{\Gamma }_2^{}\mathrm{\Lambda }`$. Now we can decompose $`\mathrm{\Gamma }_2`$ as the direct product of cyclic subgroups
$$\mathrm{\Gamma }_2=\chi _1_{k_1}\times \mathrm{}\times \chi _i_{k_i}\times \mathrm{}\times \chi _m_{k_m}\text{ with }k_1|\mathrm{}|k_i|\mathrm{}|k_m$$
(20)
where $`i`$ is the smallest index such that $`\chi _i\mathrm{\Lambda }`$. If we replace in (20) each $`\chi _j\mathrm{\Lambda }`$ by $`\chi _i^1\chi _j`$ then the decomposition remains but only one direct summand will be outside $`\mathrm{\Lambda }`$. Now we can set $`\phi =\chi _i`$, and choose $`\mathrm{\Lambda }_1`$ to be the direct product of the remaining cyclic factors of the transformed $`\mathrm{\Gamma }_2`$, and $`\mathrm{\Gamma }_2^{}`$. This then proves (i) and (ii).
For the proof of (iii) the easiest way is to recall that $`\widehat{\widehat{G}}G`$. Then if we have a direct cyclic decomposition
$$\widehat{G}=\xi _1_{n_1}\times \mathrm{}\times \xi _i_{n_i}\times \mathrm{}\times \xi _k_{n_k}$$
of $`\widehat{G}`$ it causes a similar decomposition
$$G\widehat{\widehat{G}}=a_1_{n_1}\times \mathrm{}\times a_j_{n_j}\times \mathrm{}\times a_k_{n_k}$$
of $`G\widehat{\widehat{G}}`$ such that $`\xi _j(a_j)`$ is a primitive root of $`1`$ of degree $`n_j`$ and $`\xi _j(a_i)=1`$ if $`ji`$. Since we can assume $`\phi `$ being one of $`\xi _i`$, our claim (iii) is an immediate consequence of these two decompositions. $`\mathrm{}`$
Let us denote by $`H`$ the subgroup of order 2 generated by $`h`$, as in the preceding Proposition 1. Then $`\mathrm{\Lambda }\widehat{G/H}`$, with a well-defined isomorphism $`\alpha :\mathrm{\Lambda }\widehat{G/H}`$ given by $`\alpha (\lambda )(\overline{g})=\lambda (g)`$. Here $`\overline{g}=gH`$, an element of $`G/H`$. Therefore it is easy to check that we have a $`G/H`$-grading on $`L=\mathrm{sl}(n)`$ defined by $`L_{\overline{g}}=L_gL_{gh}`$. As noted above, we have an extension of the action of $`\widehat{G}`$ on $`R=M_n`$ by Lie automorphisms if we set $`\chi I=I`$ for $`\chi \mathrm{\Lambda }`$ and $`\chi I=I`$ otherwise. This is equivalent to saying that the outer grading of $`L`$ extends to a Lie grading of $`R`$ if we set $`R_h=FIL_h`$ and $`R_g=L_g`$ if $`gh`$. For the factor-grading we then have $`R_{\overline{e}}=FIL_{\overline{e}}`$ and $`R_{\overline{g}}=L_{\overline{g}}`$ if $`\overline{g}\overline{e}`$.
###### Lemma 7
Let $`\mathrm{sl}(n)=L=_{gG}L_g`$ be an outer grading on $`L`$ and $`G,H,\mathrm{\Lambda },\phi `$ be as in Proposition 1. Then the $`G/H`$-grading on $`L`$ is inner and induced from a $`G/H`$-grading on $`R=M_n`$. Moreover, if $`R=R^{(0)}R^{(1)}\mathrm{}R^{(k)}`$ as in Theorem 1, with respect to $`G/H`$-grading, then all $`LR^{(i)}`$, $`i=0,1,\mathrm{},k`$, and $`L(R^{(1)}\mathrm{}R^{(k)})`$ are $`G`$-graded subalgebras of $`L`$.
Proof. Since $`\mathrm{\Lambda }\widehat{G/H}`$ acts on $`L`$ by inner automorphisms it follows that the $`G/H`$-grading of $`L`$ is inner hence induced from a unique $`G/H`$-grading of $`R`$. It follows that the $`G/H`$-Lie-grading of $`R`$ described just before this lemma is actually an associative grading. If we decompose $`R`$ as the tensor product $`R=R^{(0)}R^{(1)}\mathrm{}R^{(k)}`$ following Theorem 1 then all $`R^{(i)}`$ are stable under $`\mathrm{\Lambda }`$-action. Since $`\mathrm{\Lambda }`$ and $`\phi `$ commute, and each $`R_{\overline{g}}`$ is a weight subspace under the action of $`\mathrm{\Lambda }`$, it follows that for each $`\overline{g}G/H`$ we have $`\phi R_{\overline{g}}R_{\overline{g}}`$. Moreover,
$$R^{(0)}=\underset{\overline{g}\mathrm{Supp}R^{(0)}}{}R_{\overline{g}}\text{ hence }\phi R^{(0)}=R^{(0)}.$$
Since $`\mathrm{Supp}R^{(i)}\mathrm{Supp}R^{(j)}=\{\overline{e}\}`$ for any $`ij`$ and the centralizer $`C_R(R^{(0)})`$ of $`R^{(0)}`$ in $`R`$ equals $`R^{(1)}\mathrm{}R^{(k)}`$ it follows that if $`R_{\overline{g}}^{(i)}0`$, $`i1`$, and $`\overline{g}\overline{e}`$ then $`R_{\overline{g}}^{(i)}=R_{\overline{g}}C_R(R^{(0)})`$. It is immediate then that $`\phi R_{\overline{g}}^{(i)}=R_{\overline{g}}^{(i)}`$ for all $`\overline{g}\overline{e}`$ and $`1ik`$. Now let us consider $`R_{\overline{e}}^{(i)}`$. If $`i=1,\mathrm{},k`$ then $`R_{\overline{e}}^{(i)}=\mathrm{Span}\{I\}`$. Since $`\phi I=I`$ we have $`\phi R_{\overline{e}}^{(i)}=R_{\overline{e}}^{(i)}`$ for any $`i=1,\mathrm{},k`$. Finally, from the intersection property of the supports it follows that
$$R_{\overline{e}}=R_{\overline{e}}^{(0)}R_{\overline{e}}^{(1)}\mathrm{}R_{\overline{e}}^{(k)}=R_{\overline{e}}^{(0)}.$$
As before, we then have $`\phi R_{\overline{e}}^{(0)}=R_{\overline{e}}^{(0)}`$. Now if $`S=R^{(1)}\mathrm{}R^{(k)}`$ then it follows from $`\phi (xy)=\phi (y)\phi (x)`$ that $`\phi S=S`$. As a result, all $`R^{(i)}`$ and $`R^{(1)}\mathrm{}R^{(k)}`$ are stable under $`\widehat{G}`$.
Finally, since the action of $`\widehat{G}`$ leaves $`L`$ invariant we have that each $`LR^{(i)}`$, $`i=0,1,\mathrm{},k`$, and $`L(R^{(1)}\mathrm{}R^{(k)})`$ are $`G`$-graded subalgebras of $`L`$, as required. $`\mathrm{}`$
Since the action of $`\mathrm{\Lambda }`$ on $`L`$ is inner with each $`\lambda \mathrm{\Lambda }`$ one can associate a non-degenerate matrix $`T_\lambda M_n`$ such that
$$\lambda X=T_\lambda ^1XT_\lambda ,\text{ for any }XL.$$
(21)
We can also say that $`X`$ is homogeneous of degree $`\overline{g}`$ if and only if
$$T_\lambda ^1XT_\lambda =\lambda (g)X.$$
(22)
Now, as before, we can extend the action of $`\mathrm{\Lambda }`$ on $`L`$ to $`R=M_n`$. It is given by the same formula (21). Therefore, we have a $`G/H`$-grading of $`M_n`$. Actually, as mentioned before, for any $`\overline{g}\overline{e}`$ we have that $`R_{\overline{g}}=L_{\overline{g}}`$ while $`R_{\overline{e}}=FIL_{\overline{e}}`$.
If we know the $`\overline{G}=G/H`$-grading of $`L`$ and an automorphism $`\phi `$ satisfying $`\phi L_{\overline{g}}=L_{\overline{g}}`$ then we can recover the $`G`$-grading using the following procedure. We recall that $`L_{\overline{g}}=L_gL_{gh}`$.
###### Proposition 2
Using our previous notation, one can find $`L_a`$, for any $`a\overline{g}`$, in the form as follows:
$$L_a=\{X+\phi (a)^1(\phi X)|XL_{\overline{g}}\}.$$
(23)
Proof. Indeed, let $`M_a`$ denote the right side of (23). Since any element in $`M_a`$ is still in $`L_{\overline{g}}`$ the action of any $`\lambda \mathrm{\Lambda }`$ amounts to the scalar multiplication by $`\lambda (g)=\lambda (a)`$, for any $`a\overline{g}`$. Now let us check the action of $`\phi `$.
$$\phi (X+\phi (a)^1(\phi X))=\phi X+\phi (a)^1(\phi ^2X)).$$
Since $`\phi ^2\mathrm{\Lambda }`$ we have that $`\phi ^2X=\phi ^2(\overline{g})X=\phi ^2(a)X=(\phi (a))^2X`$. Plugging this value in (23) we obtain
$$\phi (X+\phi (a)^1(\phi X))=\phi (a)(\phi (a)^1(\phi X)+X).$$
Since $`\phi `$ and $`\mathrm{\Lambda }`$ generate the whole of $`\widehat{G}`$ it follows that $`X+\phi (a)^1(\phi X)`$ is an eigenvector for the action of any $`\chi \widehat{G}`$ with eigenvalue $`\chi (a)`$, proving that the degree of this element in the $`G`$-grading is $`a`$. So $`M_aL_a`$. Notice also that
$$M_{ah}=\{X+\phi (ah)^1(\phi X)|XL_{\overline{g}}\}=\{X\phi (a)^1(\phi X)|XL_{\overline{g}}\}.$$
It is immediate then that $`M_a+M_{ah}=L_{\overline{g}}`$. Since, being graded components of the $`G`$-grading of $`L`$, the subspaces $`L_a`$ and $`L_{ah}`$ have trivial intersection we must conclude that $`M_a=L_a`$ and $`M_{ah}=L_{ah}`$, so that our claim is correct. $`\mathrm{}`$
Now we know that the action of $`\phi `$ on $`L`$ is given by (19) and that both $`\phi `$ and $`\mathrm{\Lambda }`$ belong to the same abelian group $`\widehat{G}`$. This causes a number of relations between the matrix $`\mathrm{\Phi }`$, as in (19) and the matrices $`T_\lambda `$, as in (21), for each $`\lambda \mathrm{\Lambda }`$.
We have
$$\lambda (\phi X)=\lambda (\mathrm{\Phi }^1{}_{}{}^{t}X\mathrm{\Phi })=T_\lambda ^1\mathrm{\Phi }^1{}_{}{}^{t}X\mathrm{\Phi }T_\lambda =(\mathrm{\Phi }T_\lambda )^1{}_{}{}^{t}X(\mathrm{\Phi }T_\lambda ).$$
and also
$$\phi (\lambda X)=\phi (T_\lambda ^1XT_\lambda )=\mathrm{\Phi }^1{}_{}{}^{t}(T_\lambda ^1XT_\lambda )\mathrm{\Phi }=(^tT^1\mathrm{\Phi })^1{}_{}{}^{t}X(^tT^1\mathrm{\Phi }).$$
Since $`X`$ is any matrix with trace zero, it follows that the conjugating matrices on the right sides of both equations differ only by a scalar $`\beta `$. We then have
$$\mathrm{\Phi }T_\lambda =\beta ^tT_\lambda ^1\mathrm{\Phi }\text{ or }^tT_\lambda \mathrm{\Phi }T_\lambda =\beta \mathrm{\Phi }.$$
(24)
Now let us extend the action of $`\phi `$ to $`M_n`$, using the same formula (19). It is obvious that $`\phi `$ and $`\mathrm{\Lambda }`$ remain commutative. It follows then that for any homogeneous component $`L_{\overline{g}}`$, respectively, $`R_{\overline{g}}`$ we have $`\phi L_{\overline{g}}=L_{\overline{g}}`$, respectively, $`\phi R_{\overline{g}}=R_{\overline{g}}`$. Let us notice that the action of $`\phi `$ on $`M_n`$ satisfies the following relation:
$$\phi (XY)=\phi (Y)\phi (X)\text{ for all }X,YM_n.$$
(25)
It is clear, conversely, that any mapping $`\phi :M_nM_n`$ satisfying (25) must be of the form (19). Indeed, if we combine such $`\phi `$ with $`X^tX`$ then we obtain an automorphism of $`M_n`$, hence a conjugation by an appropriate non-degenerate matrix $`\mathrm{\Phi }`$, as needed.
If $`\phi `$ is of order $`2`$ then $`\phi `$ becomes an involution and so the factor-grading on $`M_n`$ by $`G/H`$ becomes an involution grading. Such gradings have been completely described in Subsection 3.2.
We cannot immediately apply these results to outer gradings of $`L=\mathrm{sl}(n)`$ since our outer automorphism $`\phi `$ need not be an element of order $`2`$. But there is a way to replace, in certain situations, outer automorphisms of arbitrary order by those of order two. Notice that in the statement of the theorem that follows we model on the form of $`G`$ and $`\widehat{G}`$ as it is given in Proposition 1.
###### Theorem 6
Let $`L=\mathrm{sl}(n)R=M_n`$ be graded by a group $`G`$ of the form $`G=a_{2m}\times K`$. Suppose the dual group has the form of $`\widehat{G}=\phi _{2m}\times \mathrm{\Lambda }_1`$ where $`\phi (a)=\rho `$, a primitive root of degree $`2m`$, $`\phi (K)=1`$, $`\mathrm{\Lambda }_1(a)=1`$, $`\mathrm{\Lambda }_1\widehat{K}`$. Let the action of $`\phi `$ on $`L`$ be outer, while that of $`\mathrm{\Lambda }=\phi ^2\times \mathrm{\Lambda }_1`$ inner. Suppose $`\psi `$ is an inner automorphism of $`L`$ such that the action of $`\phi ^2`$ coincides with $`\psi ^2`$ and $`\psi `$ commutes with the action of $`\widehat{G}`$. Then if we set $`h=a^m`$, there is an involution $``$ of $`R`$ and an involution preserving $`G`$-grading
$$R=\underset{gG}{}\stackrel{~}{R}_g$$
of $`R`$ such that
$$L_g=\stackrel{~}{R}_g^{()}(\stackrel{~}{R}_{gh}^{(+)}L).$$
(26)
for all $`gG`$.
Proof. Note that by Corollary 2 the above formulas define a $`G`$-grading of $`L`$.
To prove the converse, we consider a new abelian group $`\stackrel{~}{G}=c_2\times G`$ and its dual $`\widehat{\stackrel{~}{G}}=\omega _2\times \widehat{G}`$. The action of $`\widehat{\stackrel{~}{G}}`$ on $`\stackrel{~}{G}`$ naturally extends that of $`\widehat{G}`$ on $`G`$ by setting $`\omega (c)=1`$.
Now the action of $`\widehat{G}`$ on $`L`$ naturally extends to that of $`\widehat{\stackrel{~}{G}}`$ if we set $`\omega X=(\phi \psi ^1)(X)`$. As a consequence, $`L`$ acquires a $`\stackrel{~}{G}`$ grading given by
$$L_{(c^i,g)}=\{XL|(\omega ^j,\chi )X=(1)^{ij}\chi (g)X\}.$$
(27)
Here
$$(\omega ,\chi )X=\omega (\chi X)=\chi (\omega X)\text{ and }i,j=0,1,\chi \widehat{G}.$$
It is obvious that then if we know the $`\stackrel{~}{G}`$-grading of $`L`$ then the $`G`$-grading can be recovered by setting
$$L_g=L_{(e,g)}L_{(c,g)}.$$
To find the components of the $`\stackrel{~}{G}`$-grading of $`L`$, we proceed as follows. We rewrite $`\widehat{\stackrel{~}{G}}`$ in the form
$$\widehat{\stackrel{~}{G}}=\omega _2\times \phi \omega ^1_{2m}\times \mathrm{\Lambda }_1.$$
In this case, the subgroup $`P=\phi \omega ^1_{2m}\times \mathrm{\Lambda }_1`$ acts on $`L`$ by inner automorphisms. As previously, this action naturally extends to an associative action on $`R`$. Under this extended action $`\omega `$ is an involution $``$ of an associative algebra $`R`$ and so we write $`X^{}=\omega X`$. The components of the $`\stackrel{~}{G}`$-grading of $`L`$ are the symmetric and skew-symmetric components of the components of the factor-grading by $`\stackrel{~}{G}/P^{}`$, under this involution (see Lemma 5).
Now
$$P^{}=\phi \omega ^1^{}(\mathrm{\Lambda }_1)^{}.$$
Since $`(\mathrm{\Lambda }_1)^{}=c_2\times a_{2m}`$ and $`\omega ^1=\omega `$, we need only the annihilator of $`\phi \omega `$ on $`c_2\times a_{2m}`$, which is, obviously,
$$Q=(c,a^m)_2=(c,h)_2.$$
So we have to determine the $`\stackrel{~}{G}/Q`$-grading of $`L`$ and then partition it into the skew-symmetric and symmetric components. Obviously, $`\stackrel{~}{G}/QG`$, so that actually we have an inner $`G`$-grading.
We consider the components of the $`\stackrel{~}{G}/Q`$-grading of $`L`$ by $`L_{(c^i,g)Q}`$. Then
$$L_{(c^i,g)Q}=L_{(c^i,g)}L_{(c^{i+1},gh)}\text{ where }i=0,1(\text{mod}\mathrm{\hspace{0.17em}2}),gG.$$
Now we can apply Proposition 2 to conclude that for any $`XL_{(e,g)Q}`$ the element
$$X+\omega ((e,g))^1(\omega X)$$
will be of degree $`(e,g)`$ while
$$X+\omega ((c,gh))^1(\omega X)$$
of degree $`(c,gh)`$. Therefore, the elements of the form $`X+\omega X`$, that is the skew-symmetric elements of $`=\omega `$ are of degree $`(e,g)`$ while the elements of the form $`X\omega X`$, that is the symmetric elements are of degree $`(c,gh)`$.
Now $`L_g=L_{(e,g)}L_{(c,g)}`$. By the above, $`L_{(e,g)}`$ is the set of skew-symmetric elements in $`L_{(e,g)Q}`$ while $`L_{(c,g)}`$ consists of the symmetric elements in $`L_{(c,g)Q}`$, which is the same as $`L_{(e,gh)Q}`$.
As we mentioned above, $`\stackrel{~}{G}/QG`$ allows us to make the above inner grading of $`L`$ by $`\stackrel{~}{G}/Q`$ into an inner grading by $`G`$. We may set $`\stackrel{~}{L}_g=L_{(e,g)Q}=L_{(c,gh)Q}`$. Because this grading of $`L`$ is inner it is the restriction of the $`G`$-grading of $`R`$ defined by $`\stackrel{~}{R}_e=FI\stackrel{~}{L}_e`$ and $`\stackrel{~}{R}_g=\stackrel{~}{L}_g`$ for $`ge`$. Our involution $`\omega `$ of $`R`$ is then compatible with this new grading. If we denote by $`\stackrel{~}{R}_g^{(\pm )}`$ the set of $`\omega `$-symmetric (resp., skew-symmetric) elements of $`\stackrel{~}{R}_g`$, and recall that skew-symmetric elements of an involution always have trace 0, we will have
$$L_g=\stackrel{~}{R}_g^{()}(\stackrel{~}{R}_{gh}^{(+)}L).$$
$`\mathrm{}`$
###### Corollary 3
Suppose the outer automorphism $`\phi `$ as above is of order $`2`$. Then $`\phi `$ is an involution $``$ and there is an inner involution compatible $`G`$-grading
$$R=\underset{gG}{}\stackrel{~}{R}_g$$
of $`R`$ such that
$$L_g=\stackrel{~}{R}_g^{()}\text{ if }gK\text{ and }L_g=\stackrel{~}{R}_{gh}^{(+)}L\text{ if }gK.$$
(28)
Proof. In this case we have $`\psi =\mathrm{id}`$, $`\omega `$ acts as $`\phi `$, and $`P=\phi \omega \times \mathrm{\Lambda }_1`$. Now $`\stackrel{~}{L}_g=L_{(e,g)Q}=L_{(e,g)}L_{(c,gh)}`$. By definition (27), $`L_{(e,g)}=\{XL_g|\omega ^jX=X\}`$, for any $`j=0,1`$. Now $`\omega `$ acts as $`\phi `$ and so $`L_{(e,g)}=\{XL_g|X=X\}=0`$ for $`gK`$. Also, $`L_{(e,g)}=L_g`$ for $`gK`$. Similarly, $`L_{(c,g)}=\{XL_g|\omega ^jX=(1)^jX\}`$. Now $`L_{(c,g)}=\{XL_g|X=X\}=0`$ for $`gK`$. Also, $`L_{(c,g)}=L_g`$ for $`gK`$. By the proof of our theorem, $`L_{(e,g)}`$ is the set of skew-symmetric elements in $`\stackrel{~}{L}_g`$ while $`L_{(c,g)}`$ is the set of symmetric elements in $`\stackrel{~}{L}_{gh}=\stackrel{~}{R}_{gh}\stackrel{~}{L}`$ and our lemma follows. $`\mathrm{}`$
To study outer gradings on $`L=\mathrm{sl}(n)`$ we are going to apply Lemma 7. According to this lemma, the matrix algebra $`R=M_n`$ possesses a grading by the group $`\overline{G}=G/H`$. One can decompose $`R`$ as the tensor product $`R=PQ`$ of $`\overline{G}`$-graded subalgebras $`P=R^{(0)}M_p`$ and $`Q=R^{(1)}\mathrm{}R^{(k)}M_q`$, the $`\overline{G}`$-grading on $`P`$ being elementary and the $`\overline{G}`$-grading on $`Q`$ being fine. Moreover, both $`P`$ and $`Q`$ are invariant under $`\phi `$. The intersections $`LP`$ and $`LQ`$ are thus $`G`$-graded Lie algebras $`\mathrm{sl}(p)P=M_p`$ and $`\mathrm{sl}(q)Q=M_q`$, respectively, and the $`G/H`$-grading on $`P`$ is inner and elementary while $`\overline{G}`$-grading on $`Q`$ is inner and fine. Since the action of $`\phi `$ satisfies $`\phi (XY)=\phi (Y)\phi (X)`$ on $`R`$, it follows that on each $`P`$ and $`Q`$ we can write $`\phi `$ in the form (19) for appropriate $`\mathrm{\Phi }`$.
These remarks allow us to restrict ourselves to the cases where the original $`G`$-grading is either such that the induced $`G/H`$-grading is elementary or such that the induced $`G/H`$-grading is fine. When we get information about these two cases we will have to return to the general case and consider the tensor products.
## 6 Fine gradings
We start with a lemma close to \[5, Lemma 4\].
###### Lemma 8
Let $`R=M_n(F)=_{tT}R_t`$ be the $`n\times n`$-matrix algebra with an $`\epsilon `$-grading, $`T=a_n\times b_n`$. Let also $`\phi :RR`$ be a mapping on $`R`$ defined by $`\phi X=\mathrm{\Phi }^{1t}X\mathrm{\Phi }`$. If $`\phi R_t=R_t`$ for all $`tT`$ then $`n=2`$ and $`\phi `$ acts on $`\mathrm{sl}(2)`$ as the conjugation by one of the matrices $`X_a`$, $`X_b`$ or $`X_{ab}`$ (see (8)).
Proof. First we consider the $`\phi `$-action on $`X_a`$. Since $`R_a`$ is stable under $`\phi `$,
$$\mathrm{\Phi }^{1t}X_a\mathrm{\Phi }=\mathrm{\Phi }^1X_a\mathrm{\Phi }=\alpha X_a$$
for some scalar $`\alpha 0`$. Then
$$X_a^1\mathrm{\Phi }X_a=\beta \mathrm{\Phi }$$
(29)
with some $`\beta =\alpha ^1`$. Since $`X_a^n=I`$, we obtain $`\beta ^n=1`$, so that $`\beta =\epsilon ^j`$ for some $`0jn1`$.
Denote by $`P`$ the linear span of $`I,X_a,\mathrm{},X_a^{n1}`$. Then $`R=PX_bP\mathrm{}X_b^{n1}P`$ as a vector space and the conjugation by $`X_a`$ acts on $`X_b^iP`$ as the multiplication by $`\epsilon ^i`$. In particular, all eigenvectors with eigenvalue $`\epsilon ^j`$ are in $`X_b^jP`$. It follows that $`\mathrm{\Phi }X_b^jP`$, that is, $`\mathrm{\Phi }=X_b^jQ`$ for some $`QP`$.
Now we consider the action of $`\phi `$ on $`X_b`$:
$$\phi X_b=\mathrm{\Phi }^{1t}X_b\mathrm{\Phi }=\mathrm{\Phi }^1X_b^1\mathrm{\Phi }=\gamma X_b,$$
that is, $`X_b\mathrm{\Phi }X_b=\mu \mathrm{\Phi }`$ with $`\mu =\gamma ^10`$. If we write $`Q=\alpha _iX_a^i`$ then
$$X_b\mathrm{\Phi }X_b=X_b^j\underset{i}{}\alpha _iX_bX_a^iX_b=X_b^j\underset{i}{}\alpha _i^{}X_a^iX_b^2=\mu \mathrm{\Phi }=\mu X_b^j\underset{i}{}\alpha _iX_a^i,$$
(30)
In this case $`X_b^j_i\alpha _i^{}X_a^iX_b^2=\mu X_b^j_j\alpha _iX_a^i`$ where the scalars $`\alpha _j^{}`$ can be explicitly computed using (9). Since the degrees in $`X_a,X_b`$ define the degrees in the $`T`$-grading, we can see that (30) immediately implies $`X_b^2=I`$, i.e. $`n=2`$.
As we have shown before, (29) implies $`\mathrm{\Phi }=X_b^jQ`$ with $`Q=\alpha _0I+\alpha _1X_a`$. Since $`n=2`$, the argument following (29) applies if we change places of $`a`$ and $`b`$ so that $`\mathrm{\Phi }=X_a^k(\beta _0I+\beta _1X_b)`$. Comparing these two expressions we obtain that $`\mathrm{\Phi }`$ must be one of $`I`$, $`X_a`$, $`X_b`$, or $`X_{ab}`$, up to a scalar multiple. Finally note that $`X_{ab}^1YX_{ab}=^tY`$ for any traceless $`2\times 2`$ matrix $`Y`$ and then $`\mathrm{\Phi }^{1t}Y\mathrm{\Phi }`$ is the conjugation by one of the matrices $`X_{ab}I`$, $`X_{ab}X_a`$, $`X_{ab}X_b`$ or $`X_{ab}^2`$ which are scalar multiples of $`X_{ab},X_a,X_b`$ and $`I`$, respectively. $`\mathrm{}`$
A quick corollary is as follows.
###### Corollary 4
Let $`R=R^{(0)}R^{(1)}\mathrm{}R^{(k)}`$ be the decomposition of $`R=M_n(F)`$ corresponding to the factor-grading of $`R`$ by $`\overline{G}=G/H`$ where $`R^{(0)}=M_{n_0}(F)`$ has an elementary $`\overline{G}`$-grading while the $`\overline{G}`$-grading on $`R^{(1)}\mathrm{}R^{(k)}`$ is fine, as in Lemma 7. Then $`R^{(1)}\mathrm{}R^{(k)}M_2(F)`$ have fine $`(1)`$-grading. Also, the restriction of the outer automorphism $`\phi \widehat{G}\mathrm{\Lambda }`$ to any factor $`R^{(i)},1ik`$, coincides with the action of some inner automorphism. The order of the restriction of $`\phi `$ to $`R^{(1)}\mathrm{}R^{(k)}`$ equals $`2`$.
Proof. Only the claim about the order requires some justification. To do this, we have to apply Lemma 8 and the formula (25). $`\mathrm{}`$
To formulate the next theorem we recall that if a grading by a group $`G`$ on a matrix algebra $`R=M_n`$ is fine then it follows from the results of Subsection 3.1 that the support of the grading is a subgroup, say $`T`$, and a canonical graded basis can be chosen of nondegenerate matrices $`X_t`$, $`tT`$, such that $`X_tX_t^{}=\alpha (t,t^{})X_{tt^{}}`$ for a $`2`$-cocycle $`\alpha :T\times TK^{}`$ . If there is an involution $``$ on $`R`$, which respects the grading, then there is a function $`\beta :T\{\pm 1\}`$ such that $`X_t^{}=\beta (t)X_t`$.
###### Theorem 7
Let $`L=\mathrm{sl}(n)=_{gG}L_gR=M_n`$ be given an outer grading by a finite abelian group $`G`$ such that the respective $`\overline{G}`$-grading of $`R`$ is fine. Then $`n=2^k`$ and there is a fine involution grading on $`R=_{gG}\stackrel{~}{R}_g`$ with a subgroup $`T`$ as its support and an element $`h`$ of order $`2`$ in $`G`$ such that $`R`$, as a $`G`$-graded algebra with involution is isomorphic to the tensor product $`R=R^{(1)}\mathrm{}R^{(k)}`$ of graded involution stable subalgebras $`R^{(i)}`$ each of which is a matrix algebra of order $`2`$ of one of the four types in Lemma 2, with support $`T_i_2\times _2`$. The support of $`R`$ is $`T=T_1\times \mathrm{}\times T_k`$. The basis of $`R`$ is formed by the Kronecker products $`X_t=X_{t_1}\mathrm{}X_{t_k}`$ where $`t_iT_i`$. Further,
$$X_t^{}=X_{t_1}^{}\mathrm{}X_{t_k}^{}=\beta (t)X_{t_1}\mathrm{}X_{t_k}\text{ and }\beta (t)=\beta (t_1)\mathrm{}\beta (t_k).$$
Finally, the original $`G`$-grading of $`L`$ can be recovered as follows (we mention only nonzero components).
*Case 1*: $`hT`$.
1. $`L_t=\mathrm{Span}\{X_t\}`$ for $`tT`$ such that $`\beta (t)=1`$;
2. $`L_g=\mathrm{Span}\{X_t\}`$ for $`gG\{h\}`$ such that $`g=th`$, $`tT`$ such that and $`\beta (t)=1`$.
*Case 2*: $`hT`$.
1. $`L_t=\mathrm{Span}\{X_t,X_{th}\}`$ for $`tT\{h\}`$ such that $`\beta (t)=1`$ and $`\beta (th)=1`$;
2. $`L_t=\mathrm{Span}\{X_t\}`$ for $`tT`$ such that $`\beta (t)=1`$ and $`\beta (th)=1`$;
3. $`L_t=\mathrm{Span}\{X_{th}\}`$ for $`tT\{h\}`$ such that $`\beta (t)=1`$ and $`\beta (th)=1`$
Proof. By the argument preceding Lemma 7 and the definition of the factor-grading,
$$R_{\overline{e}}=\mathrm{Span}\{I\}L_{\overline{e}}=\mathrm{Span}\{I\}L_eL_h.$$
By our hypotheses, $`dimR_{\overline{e}}=1`$ and so $`L_e=L_h=\{0\}`$. Now we adopt the notation and the argument of Theorem 6. By Corollary 4 any character $`\phi G\mathrm{\Lambda }`$ acts as an automorphism of order $`2`$ and so Theorem 6 and Corollary 3 state that then $`R`$ possesses an inner grading $`R=_{gG}\stackrel{~}{R}_g`$ stable under the involution $`=\phi `$. Following the proof of Theorem 6, we write $`G=a_{2m}K`$, its dual $`\widehat{G}=\phi _{2m}\times \mathrm{\Lambda }_1`$, $`\stackrel{~}{G}=c_2\times G`$ and its dual $`\widehat{\stackrel{~}{G}}=\omega _2\times \widehat{G}`$. One chooses $`P=\phi \omega ^1_{2m}\times \mathrm{\Lambda }_1`$ and $`Q=P^{}=(c,a^m)_2=(c,h)_2.`$ Then $`\stackrel{~}{R}_g`$ is defined by $`\stackrel{~}{R}_g=R_{(e,g)Q}`$. Since $`\phi `$ commutes with $`\widehat{\stackrel{~}{G}}`$ this grading is stable under the involution. To check that this grading of $`R`$ is fine we only need to check $`dim\stackrel{~}{R}_e=dimR_{(e,e)Q}=1`$ (\[4, Section 4\]). Now
$$\mathrm{Span}\{I\}R_{(e,e)Q}=R_{(e,e)}R_{(c,h)}R_eR_h=L_eL_h\mathrm{Span}\{I\}=\mathrm{Span}\{I\},$$
proving that indeed $`dim\stackrel{~}{R}_e=1`$. After this we can invoke Theorem 3 to obtain the structure of $`R=_{gG}\stackrel{~}{R_g}`$, as it is claimed in our theorem. To obtain the components of the original grading of $`L`$, we recall that by Theorem 6, $`L_g=\stackrel{~}{R}_g^{()}\stackrel{~}{R}_{gh}^{(+)}L`$. Also, $`X_tL`$ for $`te`$ and $`X_tR^{()}`$ if $`\beta (t)=1`$ while $`X_tR^{(+)}`$ if $`\beta (t)=1`$. Now the computation leading to the explicit form of the components of the original grading becomes obvious. $`\mathrm{}`$
## 7 Elementary Gradings
In this section we continue our study of outer gradings on $`L=\mathrm{sl}(n)`$. We use the notation and results of Proposition 1. Therefore, $`\widehat{G}=\phi \times \mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }=\phi ^2\times \mathrm{\Lambda }_1`$, $`G=a\times K`$ so that $`\phi (K)=1`$, $`\phi (a)=\rho `$, $`\rho `$ the $`(2m)^{\mathrm{th}}`$ primitive root of $`1`$, $`\mathrm{\Lambda }^{}=h`$ where $`h=a^m`$, $`\mathrm{\Lambda }_1^{}=a`$, $`\phi (h)=1`$. The Lie automorphism $`\phi `$ is outer. By Lemma 7 the $`G/H`$ grading of $`L`$ is inner. We set $`\overline{G}=G/H`$. Let now the $`\overline{G}`$-grading on $`R=M_n`$ be elementary. Suppose it is given by an $`n`$-tuple $`\tau =(\overline{g}_1,\mathrm{},\overline{g}_n)`$. Then we know from \[5, Section 7\] that the action of $`\widehat{\overline{G}}\mathrm{\Lambda }`$ is given as follows. If $`\lambda \mathrm{\Lambda }`$ and $`XR`$ then
$$\lambda X=T_\lambda ^1XT_\lambda \text{ where }T_\lambda =\mathrm{diag}\{\lambda (\overline{g}_1),\mathrm{},\lambda (\overline{g}_n)\}.$$
(31)
If we use (24), then in order for $`\phi `$ to respect the elementary $`\overline{G}`$-grading in question, we must have
$$T_\lambda \mathrm{\Phi }T_\lambda =\beta \mathrm{\Phi }\text{ for all }\lambda \mathrm{\Lambda },\beta \text{ a nonzero scalar, depending on }\lambda .$$
(32)
Now we closely follow the argument in to determine the best possible form of $`\mathrm{\Phi }`$, up to an inner automorphism of $`M_n`$, making sure that the $`\overline{G}`$-grading is still elementary. If we identify $`M_n`$ with $`\mathrm{End}V`$ where $`V`$ is an $`n`$-dimensional vector space with basis $`\{e_1,\mathrm{},e_n\}`$ then an elementary grading of $`M_n`$ corresponding to (31) is induced by a $`\overline{G}`$-grading of $`V`$ where the elements of the basis are homogeneous of degrees forming $`\tau `$. If we change the places of the elements of the basis then we may assume, without loss of generality, that
$$\tau =(\underset{p_1}{\underset{}{a_1,\mathrm{},a_1}},\underset{p_2}{\underset{}{a_2,\mathrm{},a_2}},\mathrm{},\underset{p_q}{\underset{}{a_q,\mathrm{},a_q}})\text{ or, shorter, }\tau =(a_1^{(p_1)},\mathrm{},a_q^{(p_q)})$$
where $`\{a_1,\mathrm{},a_q\}=\{\overline{g}_1,\mathrm{},\overline{g}_n\}`$, $`a_ia_j`$ for $`ij`$, $`p_1+\mathrm{}+p_q=n`$.
We split $`\mathrm{\Phi }`$ into $`q^2`$ blocks by drawing horizontal and vertical lines according to the partition $`n=p_1+\mathrm{}+p_q`$. Let $`\mathrm{\Phi }_{ij}`$ denote the block of dimension $`p_i\times p_j`$ in the position $`(i,j)`$. If we use (32) then we have $`\lambda (a_i)\mathrm{\Phi }_{ij}\lambda (a_j)=\beta \mathrm{\Phi }_{ij}`$ for all $`1i,jq`$. Since $`\lambda `$ is a character, if $`\mathrm{\Phi }_{ij}0`$ then $`\lambda (a_ia_j)=\beta `$. Now if $`\mathrm{\Phi }_{ij}0`$ and $`\mathrm{\Phi }_{i^{}j}0`$ then $`\lambda (a_ia_j)=\lambda (a_i^{}a_j)`$ and so $`\lambda (a_i)=\lambda (a_i^{})`$ for any $`\lambda \mathrm{\Lambda }`$. Because $`\widehat{\overline{G}}\mathrm{\Lambda }`$ we must have $`a_i=a_i^{}`$. By our hypothesis then $`i=i^{}`$. Recalling that $`\mathrm{\Phi }`$ is non-degenerate we determine that in each row of blocks there is one block $`\mathrm{\Phi }_{ij}`$ different from zero. The same is true for the columns of blocks of $`\mathrm{\Phi }`$.
Now it follows from the commutativity $`a_ia_j=a_ja_i`$ that if in the $`i^{\mathrm{th}}`$ row we have $`\mathrm{\Phi }_{ij}0`$ then in the $`j^{\mathrm{th}}`$ row we must have $`\mathrm{\Phi }_{ji}0`$. Our main property used for this and other claims is that the value of $`a_ia_j`$ is constant for all cases where $`\mathrm{\Phi }_{ij}0`$. Let us denote by $`x_0`$ the element of the group $`\overline{G}`$ to which all such products are equal. It is obvious now that by rearranging the basis of $`V`$ and changing notation for the components of $`\tau `$ we may assume that
$$\tau =(x_1^{(k_1)},\mathrm{},x_s^{(k_s)},y_1^{(m_1)},z_1^{(m_1)},\mathrm{},y_t^{(m_t)},z_t^{(m_t)}).$$
Also, there exists an element $`x_0\overline{G}`$ such that
$$x_1^2=\mathrm{}=x_s^2=y_1z_1=\mathrm{}y_tz_t=x_0.$$
In the same basis we must have
$$\mathrm{\Phi }=\mathrm{diag}\{\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_s,\left[\begin{array}{cc}0& \mathrm{\Psi }_1\\ \mathrm{\Psi }_1^{}& 0\end{array}\right],\mathrm{},\left[\begin{array}{cc}0& \mathrm{\Psi }_t\\ \mathrm{\Psi }_t^{}& 0\end{array}\right]\}.$$
(33)
It should be noted that thanks to the non-degeneracy of $`\mathrm{\Phi }`$ we must have all ingredient matrices square and $`\mathrm{\Psi }_i`$ of the same order as $`\mathrm{\Psi }_i^{}`$.
Now we have to recall $`\phi ^2\mathrm{\Lambda }`$. Let us set $`\lambda _0=\phi ^2`$. Then $`\phi ^2X=T_{\lambda _0}^1XT_{\lambda _0}`$, for any $`XL`$. Then we should have
$$\phi ^2X=\mathrm{\Phi }^1{}_{}{}^{t}(\mathrm{\Phi }^1{}_{}{}^{t}X\mathrm{\Phi })\mathrm{\Phi }=(^t\mathrm{\Phi }^1\mathrm{\Phi })^1X(^t\mathrm{\Phi }^1\mathrm{\Phi }).$$
It follows that
$$^t\mathrm{\Phi }^1\mathrm{\Phi }=\alpha T_{\lambda _0},\text{ that is, }\mathrm{\Phi }=\alpha ^t\mathrm{\Phi }T_{\lambda _0}.$$
(34)
Considering the explicit form (33) for $`\mathrm{\Phi }`$, we find that the following are true:
$$\mathrm{\Phi }_i=\alpha \lambda _0(x_i)^t\mathrm{\Phi }_i\text{ for all }1is.$$
(35)
and
$$\left[\begin{array}{cc}0& \mathrm{\Psi }_j\\ \mathrm{\Psi }_j^{}& 0\end{array}\right]=\left[\begin{array}{cc}0& {}_{}{}^{t}\mathrm{\Psi }_{j}^{}\\ {}_{}{}^{t}\mathrm{\Psi }_{j}^{}& 0\end{array}\right]\left[\begin{array}{cc}\alpha \lambda _0(y_j)& 0\\ 0& \alpha \lambda _0(z_j)\end{array}\right].$$
(36)
It follows from (35) that $`\alpha ^2\lambda _0(x_i)^2=\alpha ^2\lambda _0(x_0)=1`$. Therefore each $`\mathrm{\Phi }_i`$ is either symmetric or skew-symmetric. We also have
$$\mathrm{\Psi }_j=\alpha \lambda _0(z_j)^t\mathrm{\Psi }_j^{}=\alpha ^2\lambda _0(z_j)\lambda _0(y_j)\mathrm{\Psi }_j.$$
(37)
So again $`\alpha ^2\lambda _0(z_j)\lambda _0(y_j)=\alpha ^2\lambda _0(x_0)=1`$.
Now let us conjugate our matrix algebra by a nondegenerate matrix of a shape similar to $`\mathrm{\Phi }`$:
$$P=\mathrm{diag}\{A_1,\mathrm{},A_s,\left[\begin{array}{cc}B_1& 0\\ 0& C_1\end{array}\right],\mathrm{},\left[\begin{array}{cc}B_t& 0\\ 0& C_t\end{array}\right]\}$$
In this case the matrices $`T_\lambda `$ will not change while the blocks of $`\mathrm{\Phi }`$ will change as follows. By Lemma 6 each $`\mathrm{\Phi }_i`$ will be replaced by $`{}_{}{}^{t}A_{i}^{}\mathrm{\Phi }_iA_i`$. This allows us to assume, without loss of generality, that, depending on whether $`\alpha \lambda _0(x_i))`$ is $`1`$ or $`1`$, each $`\mathrm{\Phi }_i`$ is either an identity matrix $`I_{k_i}`$ of appropriate order $`k_i`$ or a matrix $`S_{l_p}`$, $`S_{l_p}=\left[\begin{array}{cc}0& I_{l_p}\\ I_{l_p}& 0\end{array}\right]`$ of even order $`k_p=2l_p`$.
Now the block $`\left[\begin{array}{cc}0& \mathrm{\Psi }_j\\ \mathrm{\Psi }_j^{}& 0\end{array}\right]`$ will be replaced by $`\left[\begin{array}{cc}0& {}_{}{}^{t}B_{j}^{}\mathrm{\Psi }_jC_j\\ {}_{}{}^{t}(^tB_j\mathrm{\Psi }_j^{}C_j)& 0\end{array}\right]`$. Considering (37), this allows us to assume, without loss of generality, that the whole block in question can be replaced by $`\left[\begin{array}{cc}0& I_{m_u}\\ \alpha \lambda _0(y_u)I_{m_u}& 0\end{array}\right]`$ where $`I_{m_u}`$ is the identity matrix of an appropriate order $`m_u`$.
In the new basis $`\mathrm{\Phi }`$ will look like the following:
$`\mathrm{\Phi }`$ $`=`$ $`\mathrm{diag}\{I_{k_1},\mathrm{},I_{k_r},\left[\begin{array}{cc}0& I_{l_{r+1}}\\ I_{l_{r+1}}& 0\end{array}\right],\mathrm{},\left[\begin{array}{cc}0& I_{l_s}\\ I_{l_s}& 0\end{array}\right],`$ (47)
$`\left[\begin{array}{cc}0& I_{m_1}\\ \alpha \lambda _0(y_1)I_{m_1}& 0\end{array}\right],\mathrm{},\left[\begin{array}{cc}0& I_{m_t}\\ \alpha \lambda _0(y_t)I_{m_t}& 0\end{array}\right]\},`$
At this point we are ready to construct an inner automorphism $`\psi `$ of $`L`$ such that $`\psi `$ commutes with the action of the whole of $`\widehat{G}`$ and also the action of $`\phi ^2=\lambda _0`$ coincides with $`\psi ^2`$, which will allow us to apply Theorem 6.
We will look for $`\psi `$ in the form of an inner automorphism given by a diagonal matrix $`T_\psi `$ with respect to a basis of $`V`$ which results after all the above transformations. Notice that thanks to (34) and (47), in which we set $`\gamma _u=\alpha \lambda _0(y_u)`$ for $`u=1,\mathrm{},t`$, we can write
$`T_{\lambda _0}`$ $`=`$ $`\alpha ^1{}_{}{}^{t}\mathrm{\Phi }_{}^{1}\mathrm{\Phi }=\alpha ^1\mathrm{diag}\{I_{k_1},\mathrm{},I_{k_s},I_{l_{r+1}},I_{l_{r+1}}\mathrm{},I_{l_s},I_{l_s},`$ (48)
$`\gamma _1I_{m_1},\gamma _1^1I_{m_1},\mathrm{},\gamma _tI_{m_t},\gamma _t^1I_{m_t}\}.`$
This suggests that the matrix we want to find has the form of
$`T_\psi `$ $`=`$ $`\mathrm{diag}\{\epsilon I_{k_1},\mathrm{},\epsilon I_{k_s},\pi I_{l_{r+1}},\rho I_{l_{r+1}}\mathrm{},\pi I_{l_s},\rho I_{l_s},`$ (49)
$`\mu _1I_{m_1},\nu _1I_{m_1},\mathrm{},\mu _tI_{m_t},\nu _tI_{m_t}\}.`$
Now $`T_\psi ^2=\xi T_{\lambda _0}`$ for some scalar $`\xi `$. This gives us the following relations
$`\epsilon ^2`$ $`=`$ $`\xi \alpha ^1,`$
$`\pi ^2`$ $`=`$ $`\xi \alpha ^1,\rho ^2=\xi \alpha ^1,`$
$`\mu _u^2`$ $`=`$ $`\xi \gamma _u,\nu _u^2=\xi \gamma _u^1,u=1,\mathrm{},t.`$ (50)
Now in order for $`\psi `$ to commute with $`\phi `$ we must have an equation of the form (32) satisfied for $`T_\psi `$, that is,
$$T_\psi \mathrm{\Phi }T_\psi =\delta \mathrm{\Phi },$$
(51)
for an appropriate parameter $`\delta `$. If we use (51) then, considering all nonzero entries of $`\mathrm{\Phi }`$, we arrive at the following relations
$`\epsilon ^2`$ $`=`$ $`\delta ,`$
$`\pi \rho `$ $`=`$ $`\delta ,`$
$`\mu _u\nu _u`$ $`=`$ $`\delta .`$ (52)
In the case where at least one of the diagonal blocks of $`\mathrm{\Phi }`$ is nonzero, we can modify our $`n`$-tuple $`\tau `$ by dividing all entries by $`x_1`$. Then we will have $`x_0=\overline{e}`$ and $`\alpha ^2=1`$. Resolving (50) and (52) in this case we obtain $`\xi \alpha ^1=\delta `$ and then $`\pi \rho =\xi \alpha ^1`$ and $`\mu _u\nu _u=\xi \alpha ^1`$, for all $`u`$. It follows that, in this case, one can choose $`\epsilon `$ any square root of $`\xi \alpha ^1`$, $`\pi `$ any square root of $`\xi \alpha ^1`$ while $`\rho =\pi `$. Then $`\pi \rho =\xi \alpha ^1`$, as required. We can also take as $`\mu _u`$ any square root of $`\xi \gamma _u`$ and $`\nu _u=\frac{\xi \alpha }{\mu _u}`$. Then $`\nu _u^2=\xi \gamma _u^1`$, as needed. At the same time, $`\mu _u\nu _u=\xi \alpha =\xi \alpha ^1`$ because we have assumed $`\alpha ^2=1`$.
In the case where we have no nonzero diagonal blocks in $`\mathrm{\Phi }`$ the only equations we have to resolve are
$$\mu _u^2=\xi \gamma _u,\nu _u^2=\xi \gamma _u^1\text{ and }\mu _u\nu _u=\delta .$$
In this case we take $`\mu _u`$ any square root of $`\xi \gamma _u`$ and $`\nu _u=\frac{\xi \alpha }{\mu _u}`$. Then the value of $`\mu _u\nu _u`$ is a constant $`\xi \alpha `$ so that we can set $`\delta `$ equal to this value.
Now $`T_\psi `$ satisfies (50) hence $`\psi `$ commutes with $`\phi `$. Since $`T_\psi `$ is diagonal the conjugation by $`T_\psi `$ commutes with the conjugation by all $`T_\lambda `$, $`\lambda \mathrm{\Lambda }`$.
Thus the existence of $`\psi `$ with the properties desired has been proved. Once we found $`\psi `$ we can use Theorem 6 and make the following conclusion.
###### Theorem 8
Let $`L=\mathrm{sl}(n)`$ be given an outer grading by a finite abelian group $`G`$ such that the respective $`\overline{G}`$-grading is elementary. Then there is an involution $``$ on $`R=M_n`$, an element $`h`$ of order $`2`$ in $`G`$, and an elementary involution $`G`$-grading $`R=_{gG}\stackrel{~}{R}_g`$ such that
$$L_g=\{\begin{array}{cc}\stackrel{~}{R}_g^{()}\stackrel{~}{R}_{gh}^{(+)}\hfill & \text{ if }gh\hfill \\ \stackrel{~}{R}_h^{()}(\stackrel{~}{R}_e^{(+)}L)\hfill & \text{ otherwise}.\hfill \end{array}$$
Here $`\stackrel{~}{R}_g^{(\pm )}`$ is the set of symmetric (skew-symmetric) elements in $`\stackrel{~}{R}_g`$ with respect to the involution $``$.
Proof. All claims have been proved except that the grading $`R=_{gG}\stackrel{~}{R}_g`$ is elementary. This, however, easily follows because the action of $`\psi `$ is given as the conjugation by a diagonal matrix. So every matrix unit $`E_{ij}`$ is an eigenvector of $`\psi `$. Since the same is true for the action of $`\mathrm{\Lambda }_1`$, every matrix unit $`E_{ij}`$ is graded. It is well-known in this case that the grading must be elementary. $`\mathrm{}`$
## 8 Mixed Gradings
Our approach to handling the outer gradings on $`L=\mathrm{sl}(n)`$, is to apply Theorem 6. Therefore, given an outer automorphism of $`L=\mathrm{sl}(n)`$, we have to find an inner automorphism $`\psi `$ which commutes with the action of $`\widehat{G}`$ and such that the action of $`\phi ^2`$ coincides with $`\psi ^2`$.
At this point we have $`R=R^{(0)}R^{(1)}\mathrm{}R^{(k)}`$ where the $`\overline{G}`$-grading on $`R^{(0)}`$ is elementary and that on $`R^{(1)}\mathrm{}R^{(k)}`$ is fine. Recall also that $`G`$-grading on Lie algebra $`L`$ induces a $`G`$-grading on $`R`$ as on a Lie algebra. The subspaces $`R^{(0)}`$ and $`R^{(1)}\mathrm{}R^{(k)}`$ are also $`G`$-graded and the $`G`$-grading on them has been described in Sections 6 and 7. Let us choose $`G`$-graded bases in these two subspaces, such that every element is either traceless or is the identity matrix. Then we have a $`G`$-graded basis on $`R`$ consisting of $`II`$ and some traceless matrices of the form $`uv`$ where at least one of $`u,v`$ has trace zero. Then these latter matrices will form a $`G`$-graded basis of $`L=\mathrm{sl}(n)`$. To prove this claim we have to apply the generators of $`\widehat{G}`$ to these matrices. Assume $`g_1=\mathrm{deg}_Gu`$, $`g_2=\mathrm{deg}_Gv`$. If none of $`u`$, $`v`$ is $`I`$, and $`\lambda \mathrm{\Lambda }`$ we will have
$$\lambda (uv)=(\lambda u)(\lambda v)=\lambda (g_1)\lambda (g_2)(uv)=\lambda (g_1g_2)(uv)=\lambda (g_1g_2h)(uv).$$
If we apply $`\phi `$ then using (25) we will get
$$\phi (uv)=(\phi u)(\phi v)=\phi (g_1)\phi (g_2)(uv)=\phi (g_1g_2h)(uv).$$
It follows that $`\mathrm{deg}(uv)=g_1g_2h`$. As for the elements of the form $`uI`$ and $`Iv`$ then they retain their degrees as the elements of $`R^{(0)}`$ or $`R^{(1)}\mathrm{}R^{(k)}`$, that is, $`\mathrm{deg}(uI)=g_1`$, $`\mathrm{deg}(Iv)=g_2`$.
Using this notation, we define a mapping $`\psi :RR`$ by the formula
$$\psi (uv)=(\psi u)v.$$
(53)
Here $`\psi u`$ is defined as an inner automorphism, the result of our argument in Section 7 leading to Theorem 8. For that $`\psi `$ we had $`\psi ^2=\phi ^2`$ on $`R^{(0)}`$. We also remember that according to Corollary 4, $`\phi ^2=\mathrm{id}`$ when restricted to $`R^{(1)}\mathrm{}R^{(k)}`$. So it is immediate $`\psi ^2=\phi ^2`$ for the mapping defined by (53). Clearly, $`\psi `$ is inner, given by the matrix $`T_\psi I`$ where $`T_\psi `$ has been found in Section 7. It is obvious that $`\psi `$ commutes with $`\widehat{G}`$. Thus our $`G`$-grading of $`L=\mathrm{sl}(n)`$ can be recovered by Theorem 2 from a $`G`$-grading, which respects an involution of $`R=M_n`$. All such gradings have been completely described in Theorem 3.
Our final results will then look as follows.
###### Theorem 9
Let $`F`$ be an algebraically closed field of characteristic zero. Any grading $`L=_{gG}L_g`$ of $`L=\mathrm{sl}(n)`$ by a finite abelian group $`G`$ on $`R=M_n`$ is conjugate by an inner automorphism of $`R`$ to one of the following types.
1. The restriction to $`L`$ of any associative $`G`$-grading of $`M_n`$;
2. Given any involution $`G`$-grading $`R=_{gG}\stackrel{~}{R}_g`$ and an element $`h`$ of order $`2`$ in $`G`$, the grading defined by
$$L_g=\{\begin{array}{cc}\stackrel{~}{R}_g^{()}\stackrel{~}{R}_{gh}^{(+)}\hfill & \text{ if }gh\hfill \\ \stackrel{~}{R}_h^{()}(\stackrel{~}{R}_e^{(+)}L)\hfill & \text{ otherwise}.\hfill \end{array}$$
If we want a more explicit form, we have to notice that in both “inner” and “outer” types of the gradings we have to start with a $`G`$-grading $`R=AB`$ where $`AM_p`$ is a $`G`$-graded subalgebra with an elementary $`G`$-grading and $`BM_q`$ a $`G`$-graded subalgebra with a fine grading. There is a subgroup $`TG`$, $`T_{n_1}^2\times \mathrm{}\times _{n_k}^2`$ with $`T\mathrm{Supp}M_p=\{e\}`$, which supports $`M_q`$. Thus a basis of $`B`$ can be chosen in the form of $`\{X_t|tT\}`$ where each $`X_t`$ is the Kronecker product $`X_{t_1}\mathrm{}X_{t_k}`$ where $`t_s=a_s^{i_s}b_s^{j_s}_{n_s}^2`$ and $`X_{t_s}=X_{a_s}^{i_s}X_{b_s}^{j_s}`$, as in (8), $`s=1,\mathrm{},k`$. Also, there is a $`p`$-tuple $`\tau =(g_1,\mathrm{},g_p)`$ of elements of $`G`$, which defines the elementary grading of $`A`$ by $`\mathrm{deg}E_{ij}=g_i^1g_j`$.
In the case of Type I gradings, to obtain any grading of $`L`$ we have to choose a basis of $`R`$ in the form $`\{E_{ij}X_t|\mathrm{\hspace{0.17em}1}i,jp,tT\}`$ and set
$`L_g`$ $`=`$ $`\mathrm{Span}\{\{E_{ij}X_t|g_i^1g_jt=g,\mathrm{\hspace{0.17em}1}ijn\}`$
$`{\displaystyle }\{(E_{11}E_{ii})X_g|\mathrm{\hspace{0.17em}1}<in\}\}.`$
In the case of Type II gradings, the $`p`$-tuple $`\tau `$ has to be chosen as prescribed in Lemmas 3 and 4. The same lemmas define an involution $``$ on $`A`$. Also, $`n_1=\mathrm{}=n_k=2`$, and an involution $``$ is defined on $`B`$ by $`X_t^{}=(\mathrm{sgn}t)X_t`$ where $`\mathrm{sgn}t=\pm 1`$, as defined in Theorem 3. Now an involution $``$ is defined on $`R`$ by $`(YX_t)^{}=Y^{}X_t^{}`$ so that $`YX_t`$ is symmetric if either $`Y`$ is symmetric and $`\mathrm{sgn}t=1`$ or $`Y`$ is skew-symmetric and $`\mathrm{sgn}t=1`$. Similarly for skew-symmetric elements. It is obvious from the disjoint property for the supports of $`A`$ and $`B`$ that such symmetric (skew-symmetric) elements span $`R^{(\pm )}`$. Now we can define
$`L_g`$ $`=`$ $`\mathrm{Span}\{\{YX_t|(\mathrm{deg}Y)t=g,YX_t\text{ skew-symmetric}\}`$
$`{\displaystyle }\{YX_t|(\mathrm{deg}Y)t=gh,YX_t\text{ symmetric}\}\}`$
if $`gh`$ and
$`L_h`$ $`=`$ $`\mathrm{Span}\{\{YX_t|(\mathrm{deg}Y)t=h,YX_t\text{ skew-symmetric}\}`$
$`{\displaystyle }\{YX_t|(\mathrm{deg}Y)t=e,YX_t\text{ symmetric,}`$
$`\text{ and }\mathrm{Tr}Y=0\text{ if }t=e\}\}.`$
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# On 𝑚-dimensional toric codes
## 1. Introduction
In , J. Hansen introduced the notion of a toric code. Let $`P^m`$ be an integral convex polytope (the convex hull of some set of integer lattice points). We suggest as a good general reference for the geometry of polytopes. Suppose that $`P^m`$ is properly contained in the rectangular box $`[0,q2]^m`$ (which we denote $`\mathrm{}_{q1})`$, for some prime power $`q`$. Then a toric code is obtained by evaluating linear combinations of the monomials with exponent vector in $`P^m`$ at some subset (usually all) of the points of $`(𝔽_q^{})^m`$. We formalize this in the following definition.
###### Definition 1.1.
Let $`𝔽_q`$ be a finite field with primitive element $`\alpha `$. For $`f^m`$ with $`0f_iq2`$ for all $`i`$, let $`p_f=(\alpha ^{f_1},\mathrm{},\alpha ^{f_m})`$ in $`(𝔽_q^{})^m`$. For each $`e=(e_1,\mathrm{},e_m)P^m`$, let $`x^e`$ be the corresponding monomial and write
$$(p_f)^e=(\alpha ^{f_1})^{e_1}\mathrm{}(\alpha ^{f_m})^{e_m}.$$
The toric code $`C_P(𝔽_q)`$ over the field $`𝔽_q`$ associated to $`P`$ is the linear code of block length $`n=(q1)^m`$ with generator matrix
$$G=((p_f)^e),$$
where the rows are indexed by the $`eP^m`$, and the columns are indexed by the $`p_f(𝔽_q^{})^m`$. In other words, letting $`L=\text{Span}\{x^e:eP^m\}`$, we define the evaluation mapping
$`\text{ev}:L`$ $``$ $`𝔽_q^{(q1)^m}`$
$`g`$ $``$ $`(g(p_f):f(𝔽_q^{})^m)`$
Then $`C_P=\text{ev}(L)`$. If the field is clear from the context, we will often omit it in the notation and simply write $`C_P`$. The matrix $`G`$ will be called the standard generator matrix for the toric code.
Because of the close connection between integral polytopes and the theory of toric varieties, Hansen and others have proposed techniques from algebraic geometry such as intersection theory on algebraic surfaces and higher dimensional varieties to study toric codes and their parameters. The articles , , , and have used this approach.
In this article, for the most part, we use a more elementary viewpoint, based on the observation that the square submatrices of the standard generator matrix of a toric code are examples of a multivariate generalization of the familiar univariate Vandermonde matrices. These multivariate Vandermonde matrices have been studied by a number of different techniques in the context of multivariate polynomial interpolation. The literature there is truly vast because of the many applications of interpolation in numerical analysis and other parts of applied mathematics. We direct the reader to the bibliography of . To our knowledge, these techniques have not been used before in coding theory in this form.
Our contributions are as follows. In §2, we begin with a lower bound on the minimum distance of toric codes based on Vandermonde determinants (Proposition 2.1). We use this to study the minimum distances of toric codes from rectangular polytopes and simplices (see Theorems 2.4, 2.5, 2.9). We do this by identifying special configurations of points in $`(𝔽_q^{})^m`$ for which the Vandermonde determinant is nonzero. In the context of polynomial interpolation, such sets are called poised sets for the interpolation problem using linear combinations of monomials corresponding to the lattice points in some polytope. (They are the sets for which the interpolation problem has a unique solution for all function values assigned at those points.) Our methods here are suggested by the interpolation-theoretic computations in .
All of the 2-dimensional examples from and can be handled with our methods. But in fact our theorems are more general since they apply for all $`m2`$, not only the case $`m=2`$.
In §3 we prove the general statement that lattice equivalent polytopes $`P_1,P_2`$ yield monomially equivalent toric codes $`C_{P_1},C_{P_2}`$. We apply this result to consider a classification of $`m=2`$ toric codes up to monomial equivalence for small $`k`$.
Finally, some comments about the utility of toric codes are probably in order. In the case $`m=1`$, a toric code is just a Reed-Solomon code since $`P`$ is a line segment in $`[0,q2]`$ with integer endpoints. Higher dimensional toric codes are in a sense a natural extension of Reed-Solomon codes and have many similar properties. For instance, it is easy to see that they are all $`m`$-dimensional cyclic codes (see ). So one might hope that toric codes exist having similarly good parameters. And indeed, contains a number of examples showing that toric codes can have very good parameters, equaling or bettering the best known minimum distance for a given $`n,k`$ in . Not all toric codes perform this well, however. Our main results on toric codes from rectangular polytopes and simplices show, in fact, that their minimum distances are often quite small for their dimensions. It is an interesting problem, we believe, to determine criteria concerning which polytopes yield good toric codes.
We will use this notation in the following sections. Suppose that $`P\mathrm{}_{q1}^m`$ is an integral convex polytope. We will write $`\mathrm{\#}(P)`$ for the number of integer lattice points in $`P`$ (that is, $`\mathrm{\#}(P)=|P^m|`$). We will write
$$P^m=\{e(i):i=1,\mathrm{},\mathrm{\#}(P)\}$$
for the set of those integer lattice points. For any set $`S^n`$, $`\text{conv}(S)`$ denotes the convex hull of $`S`$.
## 2. Minimum distances via Vandermonde matrices
We begin by describing the Vandermonde matrices involved here. Using the notation introduced in §1, let $`P`$ be an integral convex polytope, and suppose $`P^m=\{e(i):i=1,\mathrm{},\mathrm{\#}(P)\}`$, listed in some particular order. Let $`S=\{p_j:j=1,\mathrm{},\mathrm{\#}(P)\}`$ be any set of $`\mathrm{\#}(P)`$ points in $`(𝔽_q^{})^m`$, also ordered. Then $`V(P;S)`$, the Vandermonde matrix associated to $`P`$ and $`S`$, is the $`\mathrm{\#}(P)\times \mathrm{\#}(P)`$ matrix
$$V(P;S)=\left(p_j^{e(i)}\right),$$
where we use the standard multi-index notation $`p_j^{e(i)}`$ to indicate the value of the monomial $`x^{e(i)}`$ at the point $`p_j`$. For example, if $`P=\text{conv}\{(0,0),(2,0),(0,2)\}`$ in $`^2`$, and $`S=\{(x_j,y_j)\}`$ is any set of 6 points in $`(𝔽_q^{})^2`$, for one particular choice of ordering of the lattice points in $`P`$, we have
$`(1)`$
$$V(P;S)=\left(\begin{array}{cccccc}1& 1& 1& 1& 1& 1\\ x_1& x_2& x_3& x_4& x_5& x_6\\ y_1& y_2& y_3& y_4& y_5& y_6\\ x_1^2& x_2^2& x_3^2& x_4^2& x_5^2& x_6^2\\ x_1y_1& x_2y_2& x_3y_3& x_4y_4& x_5y_5& x_6y_6\\ y_1^2& y_2^2& y_3^2& y_4^2& y_5^2& y_6^2\end{array}\right)$$
Since we assume $`P\mathrm{}_{q1}`$, the monomials $`x^{e(i)}`$ define linearly independent functions on $`(𝔽_q^{})^m`$ and we can also view the $`V(P;S)`$ as square submatrices of the standard generator matrix for the toric code $`C_P=C_P(𝔽_q)`$. Our first observation is that Vandermonde determinants may be used to bound the minimum distance of a toric code.
###### Proposition 2.1.
Let $`P^m`$ be an integral convex polytope. Let $`d`$ be a positive integer and assume that in every set $`T(𝔽_q^{})^m`$ with $`|T|=(q1)^md+1`$ there exists some $`ST`$ with $`|S|=\mathrm{\#}(P)`$ such that $`detV(P;S)0`$. Then the minimum distance satisfies $`d(C_P)d`$.
###### Proof.
All codewords of $`C_P`$ are linear combinations of the rows of the $`\mathrm{\#}(P)\times (q1)^m`$ standard generator matrix. If a codeword has zeroes in the locations corresponding to the subset $`T(𝔽_q^{})^m`$, then by looking at the entries corresponding to $`ST`$ we get a system of $`\mathrm{\#}(P)`$ homogeneous linear equations in the coefficients of the linear combination, whose matrix is $`V(P;S)`$. By hypothesis, this matrix is nonsingular, so all the coefficients in the linear combination must be zero. Since this is true for all $`T`$, all nonzero codewords of $`C_P`$ have at most $`(q1)^md`$ zero entries, which implies $`d(C_P)d`$. ∎
We will apply this proposition by identifying particular configurations of points $`S`$ for which $`detV(P;S)0`$, based on the particular monomials appearing in $`P`$. From examples such as (1) above, it should be relatively clear that giving characterizations of $`S`$ such that $`detV(P;S)=0`$ or $`detV(P;S)0`$ is difficult in general. Indeed, even in the context of multivariate polynomial interpolation, only partial results in special cases are well understood.
We will begin with the case where $`P`$ is a rectangular polytope in $`^m`$ (special cases are rectangles in the plane and rectangular solids in $`^3`$). For these polytopes, the extension from the case $`m=2`$ to general $`m2`$ is almost immediate, so for notational simplicity, we will treat only the case $`m=2`$ in detail.
Let $`P_{k,\mathrm{}}`$ be the rectangle $`P_{k,\mathrm{}}=\text{conv}\{(0,0),(k,0),(0,\mathrm{}),(k,\mathrm{})\}`$. Note that $`\mathrm{\#}(P_{k,\mathrm{}})=(k+1)(\mathrm{}+1)`$. We will call any set $`S`$ of $`(k+1)(\mathrm{}+1)`$ points in $`(𝔽_q^{})^2`$ consisting of $`(\mathrm{}+1)`$ distinct points on each of $`(k+1)`$ distinct vertical lines $`x=a_i`$ a $`(k+1)\times (\mathrm{}+1)`$ configuration.
When we construct $`V(P_{k,\mathrm{}};S)`$ for a $`(k+1)\times (\mathrm{}+1)`$ configuration we get a matrix as in the following proposition.
###### Proposition 2.2.
Suppose $`A=(c_{ij})`$ is an $`a\times a`$ matrix and $`B_1,B_2,\mathrm{},B_a`$ are $`b\times b`$ matrices. Let M be the $`ab\times ab`$ block matrix:
$$𝐌=\left(\begin{array}{cccc}c_{11}B_1& c_{12}B_2& \mathrm{}& c_{1a}B_a\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ c_{a1}B_1& c_{a2}B_2& \mathrm{}& c_{aa}B_a\end{array}\right)$$
Then $`det(𝐌)=\pm det(A)^bdet(B_1)det(B_2)\mathrm{}det(B_a)`$.
The matrix $`𝐌`$ is similar to a tensor product matrix, but the $`B_i`$ may be different matrices, so this construction is somewhat more general.
###### Proof.
If $`det(A)=0`$ or $`det(B_i)=0`$ for some $`i`$, then $`det(𝐌)=0`$ as well. So, we assume that $`det(A)0`$ and $`det(B_i)0`$ for all $`i`$. In order to find the determinant of $`𝐌`$, we may transform $`𝐌`$ into a block upper triangular matrix using blockwise row operations, obtaining a matrix $`𝐌^{}`$ in the following form:
$$𝐌^{}=\left(\begin{array}{ccccc}c_{11}^{}B_1& & & & \\ 0& c_{22}^{}B_2& & & \\ \mathrm{}& 0& \mathrm{}& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ 0& 0& \mathrm{}& 0& c_{aa}^{}B_a\end{array}\right),$$
in which the $`c_{ii}^{}`$’s are the same as the entries obtained by the corresponding row operations applied to the matrix $`A`$. Now we have a block-upper triangular matrix. This implies that the determinant of the matrix is the product of the determinants of the diagonal block entries. Thus,
$`det(𝐌^{})`$ $`=`$ $`(c_{11}^{})^bdet(B_1)(c_{22}^{})^bdet(B_2)\mathrm{}(c_{aa}^{})^bdet(B_a)`$
$`=`$ $`(c_{11}^{}\mathrm{}c_{aa}^{})^bdet(B_1)det(B_2)\mathrm{}det(B_a)`$
$`=`$ $`det(A^{})^bdet(B_1)det(B_2)\mathrm{}det(B_a)`$
We know that $`det(A^{})=\pm det(A)`$ (some row interchanges might have been necessary in the reduction to upper triangular form). So, we can substitute to yield:
$$det(𝐌)=\pm det(A)^bdet(B_1)det(B_2)\mathrm{}det(B_a),$$
which is what we wanted to show. ∎
When $`S`$ is a $`(k+1)\times (\mathrm{}+1)`$ configuration, consisting of $`\mathrm{}+1`$ distinct points on each of $`k+1`$ distinct lines $`x=a_u`$, $`u=1,\mathrm{},k+1`$, then $`V(P_{k,\mathrm{}};S)`$ has the form given in the proposition, where $`A`$ is the ordinary $`(k+1)\times (k+1)`$ univariate Vandermonde matrix of the $`x=a_u`$, and $`B_u`$ is the $`(\mathrm{}+1)\times (\mathrm{}+1)`$ univariate Vandermonde matrix of the $`y`$-coordinates of the points on $`x=a_u`$. Hence we obtain the following consequence.
###### Corollary 2.3.
Let $`S`$ be a $`(k+1)\times (\mathrm{}+1)`$ configuration in $`(𝔽_q^{})^2`$. Then
$$detV(P_{k,\mathrm{}};S)0.$$
###### Proof.
This follows from the factorization of the univariate Vandermonde determinants as products of differences of the $`x`$\- and $`y`$-values at the points in $`S`$. ∎
We are now ready to prove the first major result of this section.
###### Theorem 2.4.
Let $`k,\mathrm{}<q1`$ so that $`P_{k,\mathrm{}}\mathrm{}_{q1}^2`$. Then the minimum distance of the two-dimensional toric code $`C_{P_{k,\mathrm{}}}`$ is
$$d(C_{P_{k,\mathrm{}}})=(q1)^2(k+\mathrm{})(q1)+k\mathrm{}=((q1)k)((q1)\mathrm{}).$$
###### Proof.
We write $`d=d(C_{P_{k,\mathrm{}}})`$. In order to show equality, we will show that both $`d(q1)^2(k+\mathrm{})(q1)+k\mathrm{}`$ and $`d(q1)^2(k+\mathrm{})(q1)+k\mathrm{}`$. We start with the former. In $`L=\text{Span}\{1,x,x^2,\mathrm{},x^k,y,xy,x^2y,\mathrm{},x^ky,\mathrm{},x^ky^{\mathrm{}}\}`$ consider a polynomial $`p(x,y)=q_1(x)q_2(y)`$ where $`q_1(x)`$ and $`q_2(y)`$ factor completely:
$$p(x,y)=(xa_1)(xa_2)\mathrm{}(xa_k)(yb_1)(yb_2)\mathrm{}(yb_{\mathrm{}}),$$
for some distinct $`a_u`$ and $`b_v`$ in $`𝔽_q^{}`$ This means that the codeword $`\text{ev}(p)`$ has zeros in the positions corresponding to the points along $`k`$ distinct vertical lines and $`\mathrm{}`$ distinct horizontal lines in $`(𝔽_q^{})^2`$. We see that $`\text{ev}(p)`$ has $`(k+\mathrm{})(q1)k\mathrm{}`$ zeroes. Accordingly, we have $`d(q1)^2(k+\mathrm{})(q1)+k\mathrm{}`$.
Now we show $`d(q1)^2(k+\mathrm{})(q1)+k\mathrm{}`$. By Corollary 2.3 and Proposition 2.1, it suffices to show that every set $`T`$ of size $`(k+\mathrm{})(q1)k\mathrm{}+1`$ in $`(𝔽_q^{})^2`$ contains a $`(k+1)\times (\mathrm{}+1)`$ configuration $`S`$. Since $`q1>k`$, we see that $`(k+\mathrm{})(q1)k\mathrm{}+1>\mathrm{}(q1)`$. So, by the pigeonhole principle, some vertical line $`x=a_1`$ contains $`\mathrm{}+1`$ distinct points from $`T`$. There are at most $`(q1)(\mathrm{}+1)=q\mathrm{}2`$ other points on that line. Therefore, there are at least $`(k+\mathrm{}1)(q1)k\mathrm{}+1`$ other points in $`T`$ not on $`x=a_1`$. To repeat this argument to find the rest of the configuration $`S`$, we must show that for all $`1jk`$, $`(k+\mathrm{}j)(q1)k\mathrm{}+1>\mathrm{}(q1j)`$. To do this, we subtract $`\mathrm{}(q1j)`$ and perform some arithmetic:
$$(k+\mathrm{}j)(q1)k\mathrm{}+1\mathrm{}(q1j)=(kj)((q1)\mathrm{})+1>0.$$
So, after we have found the first $`j`$ sets of $`\mathrm{}+1`$ points on vertical lines in $`T`$, there are still enough additional points left in $`T`$ so that we can find the next set of $`\mathrm{}+1`$ points. After $`k+1`$ steps, we have a complete $`(k+1)\times (\mathrm{}+1)`$ configuration $`ST`$. ∎
The method used in the proof of Theorem 2.4 extends without difficulty to toric codes constructed from $`k_1\times k_2\times \mathrm{}\times k_m`$ rectangular polytopes $`P_{k_1,\mathrm{},k_m}^m`$ for all $`m2`$. The result is as follows.
###### Theorem 2.5.
Let $`k_1,\mathrm{},k_m`$ be small enough so that $`P_{k_1,\mathrm{},k_m}\mathrm{}_{q1}^m`$. Then the minimum distance of the $`m`$-dimensional toric code $`C_{P_{k_1,\mathrm{},k_m}}`$ is
$$d(C_{P_{k_1,\mathrm{},k_m}})=\underset{i=1}{\overset{m}{}}((q1)k_i).$$
We next turn to the toric codes $`C_{P_{\mathrm{}}(m)}`$ for $`P_{\mathrm{}}(m)`$ an $`m`$-dimensional simplex of the form
$$P_{\mathrm{}}(m)=\text{conv}\{\mathrm{𝟎},\mathrm{}𝐞_1,\mathrm{},\mathrm{}𝐞_m\},$$
where the $`𝐞_i`$ are the standard basis vectors in $`^m`$. The monomials corresponding to the $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)`$ integer lattice points in $`P_{\mathrm{}}(m)`$ are all of the monomials in $`m`$ variables of total degree $`\mathrm{}`$. Naturally enough, the corresponding Vandermonde matrices arise in the study of multivariate polynomial interpolation using polynomials of bounded total degree. The next recursive definition gives the special configurations of points where we will be able to compute Vandermonde determinants for the $`P_{\mathrm{}}(m)`$.
###### Definition 2.6.
If $`m=1`$, an $`\mathrm{}`$th order simplicial configuration is any collection of $`\left(\genfrac{}{}{0pt}{}{1+\mathrm{}}{\mathrm{}}\right)`$ distinct points in $`𝔽_q^{}`$. For $`m2`$, we will say that a collection $`S`$ of $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)`$ points in $`(𝔽_q^{})^m`$ is an $`m`$-dimensional $`\mathrm{}`$th order simplicial configuration if the following conditions hold:
1. For some $`i`$, $`1im`$, there are hyperplanes $`x_i=a_1,x_i=a_2,\mathrm{},x_i=a_{\mathrm{}+1}`$ such that for each $`1j\mathrm{}+1`$, $`S`$ contains exactly $`\left(\genfrac{}{}{0pt}{}{m1+j1}{j1}\right)`$ points with $`x_i=a_j`$. (Note that
$`(2)`$
$$\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)=\underset{j=1}{\overset{\mathrm{}+1}{}}\left(\genfrac{}{}{0pt}{}{m1+j1}{j1}\right)$$
by a standard binomial coefficient identity.)
2. For each $`j`$, $`1j\mathrm{}+1`$, the points in $`x_i=a_j`$ form an $`(m1)`$-dimensional simplicial configuration of order $`j1`$.
We call these special configurations of points simplicial configurations because they mimic, to an extent, the arrangement of the integer lattice points in the corresponding simplex $`P_{\mathrm{}}(m)`$. For instance, Figure 1 shows a $`2`$-dimensional simplicial configuration of order 2 in $`(𝔽_8^{})^2`$. It consists of six points. (We write $`\alpha `$ for a primitive element in the field $`𝔽_8`$.) Note that there is $`1=\left(\genfrac{}{}{0pt}{}{1+0}{0}\right)`$ point in $`S`$ on the line $`x_1=a_1=\alpha ^4`$, $`2=\left(\genfrac{}{}{0pt}{}{1+1}{1}\right)`$ points on the line $`x_1=a_2=\alpha ^3`$, and $`3=\left(\genfrac{}{}{0pt}{}{1+2}{2}\right)`$ on the line $`x_1=a_3=\alpha `$.
Figure 1. A 2-dimensional simplicial configuration $`S`$ of order 2.
In order to state our next result, a sort of recurrence relation for the Vandermonde determinants $`detV(P_{\mathrm{}}(m);S)`$ where $`S`$ is a an $`m`$-dimensional simplicial configuration, we introduce some notation.
Let $`S`$ be an $`m`$-dimensional $`\mathrm{}`$th order simplicial configuration consisting of $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)`$ points, in hyperplanes $`x_m=a_1,\mathrm{},x_m=a_{\mathrm{}+1}`$. Write $`S=S^{}S^{\prime \prime }`$ where $`S^{}`$ is the union of the points in $`x_i=a_1,\mathrm{},a_{\mathrm{}}`$, and $`S^{\prime \prime }`$ is the set of points in $`x_i=a_{\mathrm{}+1}`$. Also, let $`\pi :𝔽_q^m𝔽_q^{m1}`$ be the projection on the first $`m1`$ coordinates. By the definition, it follows that both $`S^{}`$ and $`\pi (S^{\prime \prime })`$ are themselves simplicial configurations, with $`S^{}`$ of dimension $`m`$ and order $`\mathrm{}1`$, and $`\pi (S^{\prime \prime })`$ of dimension $`m1`$ and order $`\mathrm{}`$.
###### Theorem 2.7.
Let $`P_{\mathrm{}}(m)`$ be as above and let $`S`$ be an $`\mathrm{}`$th order simplicial configuration of $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)`$ points as in the paragraph above. Then writing $`p=(p_1,\mathrm{},p_m)`$ for points $`p(𝔽_q^{})^m`$,
$$detV(P_{\mathrm{}}(m);S)=\pm \underset{pS^{}}{}(p_ma_{\mathrm{}+1})detV(P_\mathrm{}1(m);S^{})detV(P_{\mathrm{}}(m1);\pi (S^{\prime \prime })).$$
Before we give the proof of this theorem, we will do two things. First, we give an example to illustrate what the theorem is saying. The idea for this computation comes from , where corresponding sets of points in $`^m`$ are identified as poised sets for interpolation by polynomials of degree bounded bounded by $`\mathrm{}`$. Consider all polynomials of degree $`2`$ in three variables and the Vandermonde matrix $`V(P_2(3);S)`$. For notational simplicity, write points in a 3-dimensional simplicial configuration $`S(𝔽_q^{})^3`$ of order 2 as $`(x_i,y_i,z_i)`$, for $`i=1,\mathrm{},10=\left(\genfrac{}{}{0pt}{}{3+2}{2}\right)`$. Here $`S^{}`$ consists of the first four points in $`S`$, and $`S^{\prime \prime }`$ consists of the other six points. Under the hypothesis that $`S`$ is a simplicial configuration, we have $`z_5=z_6=\mathrm{}=z_{10}=c`$ for some $`c=a_3`$. Noting this, but ignoring other equalities between the coordinates, we see $`V(P_2(3);S)=`$
$$\left(\begin{array}{cccccccccc}1& 1& 1& 1& 1& 1& 1& 1& 1& 1\\ x_1& x_2& x_3& x_4& x_5& x_6& x_7& x_8& x_9& x_{10}\\ y_1& y_2& y_3& y_4& y_5& y_6& y_7& y_8& y_9& y_{10}\\ z_1& z_2& z_3& z_4& c& c& c& c& c& c\\ x_1^2& x_2^2& x_3^2& x_4^2& x_5^2& x_6^2& x_7^2& x_8^2& x_9^2& x_{10}^2\\ x_1y_1& x_2y_2& x_3y_3& x_4y_4& x_5y_5& x_6y_6& x_7y_7& x_8y_8& x_9y_9& x_{10}y_{10}\\ y_1^2& y_2^2& y_3^2& y_4^2& y_5^2& y_6^2& y_7^2& y_8^2& y_9^2& y_{10}^2\\ x_1z_1& x_2z_2& x_3z_3& x_4z_4& x_5c& x_6c& x_7c& x_8c& x_9c& x_{10}c\\ y_1z_1& y_2z_2& y_3z_3& y_4z_4& y_5c& y_6c& y_7c& y_8c& y_9c& y_{10}c\\ z_1^2& z_2^2& z_3^2& z_4^2& c^2& c^2& c^2& c^2& c^2& c^2\end{array}\right).$$
To evaluate the determinant of this matrix, we perform row operations to introduce zeroes. First subtract $`c`$ times row 4 from row 10, then $`c`$ times row 3 from row 9, $`c`$ times row 2 from row 8, and finally $`c`$ times row 1 from row 4. After rearranging rows, we find a matrix with a block of zeroes:
$`(3)`$
$$\left(\begin{array}{ccccccccc}z_1c& \mathrm{}& z_4c& 0& 0& 0& 0& 0& 0\\ x_1(z_1c)& \mathrm{}& x_4(z_4c)& 0& 0& 0& 0& 0& 0\\ y_1(z_1c)& \mathrm{}& y_4(z_4c)& 0& 0& 0& 0& 0& 0\\ z_1(z_1c)& \mathrm{}& z_4(z_4c)& 0& 0& 0& 0& 0& 0\\ 1& \mathrm{}& 1& 1& 1& 1& 1& 1& 1\\ x_1& \mathrm{}& x_4& x_5& x_6& x_7& x_8& x_9& x_{10}\\ y_1& \mathrm{}& y_4& y_5& y_6& y_7& y_8& y_9& y_{10}\\ x_1^2& \mathrm{}& x_4^2& x_5^2& x_6^2& x_7^2& x_8^2& x_9^2& x_{10}^2\\ x_1y_1& \mathrm{}& x_4y_4& x_5y_5& x_6y_6& x_7y_7& x_8y_8& x_9y_9& x_{10}y_{10}\\ y_1^2& \mathrm{}& y_4^2& y_5^2& y_6^2& y_7^2& y_8^2& y_9^2& y_{10}^2\end{array}\right).$$
Up to a sign, the determinant of $`V(P_2(3);S)`$ is therefore equal to the product of
$$det\left(\begin{array}{cccc}z_1c& z_2c& z_3c& z_4c\\ x_1(z_1c)& x_2(z_2c)& x_2(z_3c)& x_4(z_4c)\\ y_1(z_1c)& y_2(z_2c)& y_3(z_3c)& y_4(z_4c)\\ z_1(z_1c)& z_2(z_2c)& z_3(z_3c)& z_4(z_4c)\end{array}\right),$$
which equals
$$(z_1c)(z_2c)(z_3c)(z_4c)detV(P_1(3),S^{}),$$
and the determinant of the lower right $`6\times 6`$ block in (3), which is $`V(P_2(2),\pi (S^{\prime \prime }))`$. Hence
$$detV(P_2(3);S)=\pm \underset{j=1}{\overset{4}{}}(z_jc)detV(P_1(3),S^{})detV(P_2(2),\pi (S^{\prime \prime })),$$
as in the statement of the theorem.
Second, we note the following immediate consequence of the theorem.
###### Corollary 2.8.
Let $`P_{\mathrm{}}(m)`$ be as above and let $`S`$ be an $`\mathrm{}`$th order simplicial configuration of $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)`$ points. Then $`detV(P_{\mathrm{}}(m);S)0`$.
###### Proof.
This is seen easily by a double induction on $`m`$ and $`\mathrm{}`$. In the base cases $`m=1`$, $`\mathrm{}`$ arbitrary, we have an ordinary univariate Vandermonde determinant, which is nonzero by the definition of a simplicial configuration. For the induction step, the recurrence given in the theorem then establishes the corollary since all the factors are nonzero under the hypothesis that $`S`$ is a simplicial configuration. ∎
We will now give a proof of Theorem 2.7, following the methods from used in the example above.
###### Proof.
Since $`S`$ is simplicial, the $`\left(\genfrac{}{}{0pt}{}{m+\mathrm{}1}{\mathrm{}1}\right)`$ points in $`S^{\prime \prime }S`$ all have the same $`x_m`$-coordinate, say $`x_m=a_{\mathrm{}+1}=c`$. For each pair of monomials $`x^e`$ and $`x^ex_m`$ corresponding to points in $`P_{\mathrm{}}(m)`$, we perform a row operation on the Vandermonde matrix, subtracting $`c`$ times the row of $`V(P_{\mathrm{}}(m),S)`$ for $`x^e`$ from the row for $`x^ex_m`$ in decreasing order by the degree in $`x_m`$. After rearranging rows to put all the zeroes created by these operations in the upper right block, the lower right block in the columns corresponding to $`S^{\prime \prime }`$ is the matrix $`V(P_{\mathrm{}}(m1),\pi (S^{\prime \prime }))`$. In the upper left block (in the columns corresponding to $`S^{}`$), all the entries in a column are divisible by one of the $`(p_mc)`$ for $`pS^{}`$. Factoring out those factors from each column, the matrix that is left is $`V(P_\mathrm{}1(m),S^{})`$, and the statement of the theorem follows. ∎
We note that if we apply the recurrence relation repeatedly, another corollary of Theorem 2.7 is a closed formula for $`detV(P_{\mathrm{}}(m);S)`$ for $`S`$ a simplicial configuration in terms of univariate Vandermonde determinants. We will not need this formula, so we leave the derivation of its exact form as an exercise.
We will now use Corollary 2.8 to establish the minimum distances of the toric codes $`C_{P_{\mathrm{}}(m)}`$.
###### Theorem 2.9.
Let $`\mathrm{}<q1`$, and let $`P_{\mathrm{}}(m)`$ be the simplex in $`^m`$ defined above. Then the minimum distance of the toric code $`C_{P_{\mathrm{}}(m)}`$ is given by
$$d(C_{P_{\mathrm{}}(m)})=(q1)^m\mathrm{}(q1)^{m1}.$$
###### Proof.
Let $`d=d(C_{P_{\mathrm{}}(m)})`$. As in the proof of Theorem 2.4 above, we show both inequalities $`d(q1)^m\mathrm{}(q1)^{m1}`$ and $`d(q1)^m\mathrm{}(q1)^{m1}`$ hold. The first of these follows as in Theorem 2.4, since the completely reducible polynomials
$$p(x_m)=(x_ma_1)\mathrm{}(x_ma_{\mathrm{}})$$
for $`a_i`$ distinct in $`𝔽_q^{}`$ are contained in $`L=\text{Span}\{x^e:eP_{\mathrm{}}(m)\}`$. Such a polynomial has zeroes at all the $`𝔽_q`$-rational points with nonzero coordinates on the union of the hyperplanes $`x_m=a_i`$. There are $`\mathrm{}(q1)^{m1}`$ such points, so $`d(q1)^m\mathrm{}(q1)^{m1}`$ as claimed.
To establish the reverse inequality, by Corollary 2.8 and Proposition 2.1, it suffices to show that every set $`T`$ of $`\mathrm{}(q1)^{m1}+1`$ points contains an $`m`$-dimensional simplicial configuration $`S`$ of order $`\mathrm{}`$. By the pigeonhole principle, $`T`$ contains some set of $`\mathrm{}(q1)^{m2}+1`$ points with the same $`x_m`$-coordinate $`x_m=a_{\mathrm{}+1}`$. By an easy induction on $`m`$, it follows that
$$\mathrm{}(q1)^{m2}+1>\left(\genfrac{}{}{0pt}{}{m1+\mathrm{}}{\mathrm{}}\right)$$
(the first term on the right of (2)) for all $`m2`$ (and $`q1>\mathrm{}`$, of course) There are at least $`(\mathrm{}1)(q1)^{m1}+1`$ points in $`T`$ that do not lie on $`x_m=a_{\mathrm{}+1}`$. Hence we can apply the same pigeonhole principle argument repeatedly (or argue by induction on $`\mathrm{}`$) to see that $`T`$ contains
$$\underset{j=1}{\overset{\mathrm{}+1}{}}\left(\genfrac{}{}{0pt}{}{m1+j1}{j1}\right)=\left(\genfrac{}{}{0pt}{}{m+\mathrm{}}{\mathrm{}}\right)$$
points making up an $`m`$-dimensional $`\mathrm{}`$th order simplicial configuration. ∎
We also have the following consequence for more general simplices. Let $`\mathrm{}_i1`$ for all $`i`$ and define
$$P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}=\text{conv}\{\mathrm{𝟎},\mathrm{}_1𝐞_1,\mathrm{},\mathrm{}_m𝐞_m\},$$
where again the $`𝐞_i`$ are the standard basis vectors in $`^m`$.
###### Corollary 2.10.
If $`P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}\mathrm{}_{q1}^m`$, and $`\mathrm{}=\mathrm{max}_i\mathrm{}_i`$, then
$$d(C_{P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}})=(q1)^m\mathrm{}(q1)^{m1}.$$
###### Proof.
By definition, $`P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}P_{\mathrm{}}(m)`$. Hence $`C_{P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}}`$ is a subcode of $`C_{P_{\mathrm{}}(m)}`$ and $`d(C_{P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}})d(C_{P_{\mathrm{}}(m)})=(q1)^m\mathrm{}(q1)^{m1}`$. But if $`\mathrm{}=\mathrm{}_{i_0}`$, then $`C_{P_{\mathrm{}_1,\mathrm{},\mathrm{}_m}}`$ also contains codewords of weight exactly $`(q1)^m\mathrm{}(q1)^{m1}`$ obtained from evaluation of completely reducible polynomials
$$p(x_i)=(x_{i_0}a_1)\mathrm{}(x_{i_0}a_{\mathrm{}})$$
for some $`i`$ and distinct $`a_j𝔽_q^{}`$. This establishes the corollary. ∎
The same sort of reasoning applies to any toric code from a polytope $`PP_{\mathrm{}}(m)`$ that contains one complete edge $`\text{conv}\{\mathrm{𝟎},\mathrm{}𝐞_i\}`$ or $`\text{conv}\{\mathrm{}𝐞_i,\mathrm{}𝐞_j\}`$ of the simplex $`PP_{\mathrm{}}(m)`$ and gives the same minimum distance for $`C_P`$.
## 3. Classification of toric codes
In this section we will begin by stating and proving a theorem guaranteeing that two toric codes have the same parameters. We begin by introducing some terminology.
###### Definition 3.1.
Let $`C_1`$ and $`C_2`$ be two codes of block length $`n`$ and dimension $`k`$ over $`𝔽_q`$. Let $`G_1`$ be a generator matrix for $`C_1`$. Then $`C_1`$ and $`C_2`$ are said to be monomially equivalent if there is an invertible $`n\times n`$ diagonal matrix $`\mathrm{\Delta }`$ and an $`n\times n`$ permutation matrix $`\mathrm{\Pi }`$ such that
$$G_2=G_1\mathrm{\Delta }\mathrm{\Pi }$$
is a generator matrix for $`C_2`$.
It is easy to see that monomial equivalence is actually an equivalence relation on codes since a product $`\mathrm{\Pi }\mathrm{\Delta }`$ equals $`\mathrm{\Delta }^{}\mathrm{\Pi }`$ for another invertible diagonal matrix $`\mathrm{\Delta }^{}`$. It is also a direct consequence of the definition that monomially equivalent codes $`C_1`$ and $`C_2`$ have the same dimension and the same minimum distance (indeed, the same full weight enumerator).
Next we turn to a natural notion of equivalence for polytopes. Recall that an affine transformation of $`^m`$ is a mapping of the form $`T(𝐱)=M𝐱+\lambda `$, where $`\lambda `$ is a fixed vector and $`M`$ is an $`m\times m`$ matrix. The affine mappings $`T`$ where $`M,\lambda `$ have integer entries and $`M\text{GL}(m,)`$ (so $`det(M)=\pm 1`$) are precisely the bijective affine mappings from the integer lattice $`^m`$ to itself.
###### Definition 3.2.
We will say that two integral convex polytopes $`P_1`$ and $`P_2`$ in $`^m`$ are are lattice equivalent if there exists an invertible integer affine transformation $`T`$ as above such that $`T(P_1)=P_2`$.
This brings us to our next theorem which relates the two concepts we have just defined.
###### Theorem 3.3.
If two polytopes $`P_1`$ and $`P_2`$ are lattice equivalent, then the toric codes $`C_{P_1}`$ and $`C_{P_2}`$ are monomially equivalent.
###### Proof.
Suppose we have two lattice equivalent polytopes $`P_1`$ and $`P_2`$. Both $`P_1`$ and $`P_2`$ contain integer lattice points corresponding to monomials of the form $`x^e`$ where $`e^m`$. By our hypothesis on $`P_1`$ and $`P_2`$, there exists an invertible integer transformation
$$T(𝐱)=M(𝐱)+\lambda $$
such that $`T(P_1)=P_2`$ and $`M`$ is an element of $`GL(m,)`$ so $`det(M)=\pm 1`$. Hence $`\mathrm{\#}(P_1)=\mathrm{\#}(P_2)`$. Let $`P_1^m=\{e(i):i=1,\mathrm{},\mathrm{\#}(P_1)\}`$. So, $`C_{P_1}`$ is spanned by $`\text{ev}(x^{e(i)})`$ for $`1in`$, and similarly $`C_{P_2}`$ is spanned by $`\text{ev}(x^{T(e(i))})`$. Write $`\alpha `$ for a primitive element in $`𝔽_q^{}`$. Let $`e(i)P_1^m`$ and define $`\alpha ^f=(\alpha ^{f_1},\mathrm{},\alpha ^{f_m})(𝔽_q^{})^m`$. The component of $`\text{ev}(x^{e(i)})C_{P_1}`$ corresponding to $`\alpha ^f`$ is $`\alpha ^{e(i),f}`$, where $`e(i),f`$ is the usual dot product. The corresponding entry in the codeword $`\text{ev}(x^{T(e(i))})`$ in $`C_{P_2}`$ is is $`\alpha ^{T(e(i)),f}`$. This can be rewritten as
$$\alpha ^{Me(i)+\lambda ,f}=\alpha ^{Me(i),f}\alpha ^{\lambda ,f}$$
The second term of the product is not dependent on $`e(i)`$. These nonzero scalars are the diagonal entries in the matrix $`\mathrm{\Delta }`$ as in the definition of monomially equivalent codes. By a standard property of dot products,
$$\alpha ^{Me(i),f}=\alpha ^{e(i),M^tf}.$$
The transposed matrix $`M^t`$ also defines a bijective mapping from $`^m`$ to $`^m`$ since $`det(M^t)=det(M)=\pm 1`$. Now we must show that $`M^t`$ induces a permutation of $`(𝔽_q^{})^m`$. Suppose $`M^tfM^tg(modq1)`$. Since $`det(M^t)=\pm 10`$, we know that $`M^t`$ is invertible and $`(M^t)^1`$ is also an integer matrix. So, we can multiply by $`(M^t)^1`$ on the left. Hence, $`fg(modq1)`$ and $`M^t`$ defines a permutation of the points $`\alpha ^f`$, as desired. Note that $`M^t`$ permutes all of the codewords in the same way. This gives the permutation matrix $`\mathrm{\Pi }`$. Hence $`C_{P_1}`$ is monomially equivalent to $`C_{P_2}`$. ∎
In the remainder of this section, we will show how this result leads to a complete classification for toric codes with $`m=2`$ and $`k5`$. The classification could also be continued, of course, using a census of lattice equivalence classes of lattice polytopes with given $`\mathrm{\#}(P)`$.
###### Proposition 3.4.
Every toric surface code $`C_P`$ with $`k=2`$ is monomially equivalent to the toric code $`C_{P_2}`$ for $`P_2=\text{conv}\{(0,0),(1,0)\}`$.
###### Proof.
Let $`e(1),e(2)^2`$ be the integer lattice points in $`P`$. We can use a translation to map $`e(1)`$ to $`(0,0)`$. Then let $`e(2)=(a,b)^2`$. By convexity we have that gcd$`(a,b)=1`$, since otherwise there would be additional integer lattice points on the line from $`e(1)`$ to $`e(2)`$ and $`k=\mathrm{\#}(P)`$ would be greater than 2. Since $`\mathrm{gcd}(a,b)=1`$, there exist integers $`r,s`$ such that $`ra+sb=1`$, and this implies that there exists an invertible integer matrix $`M=\left(\begin{array}{cc}r& s\\ b& a\end{array}\right)`$ such that $`M\left(\begin{array}{c}a\\ b\end{array}\right)=\left(\begin{array}{c}1\\ 0\end{array}\right).`$ Hence there is an affine equivalence between $`P`$ and $`P_2`$. By Theorem 3.1, this completes the proof. ∎
Next, we wish to find a “nice” lattice polygon in each possible lattice equivalence class with $`\mathrm{\#}(P)=3,4,5`$. One way is to add additional points to $`P_2`$. Using Pick’s Theorem: $`A(P)=\mathrm{\#}(P)+\frac{1}{2}(P)1,`$ (where $`(P)`$ is the number of lattice points in the boundary of $`P`$) then eliminating lattice equivalent polygons, we obtain the following.
###### Theorem 3.5.
Every toric surface code with $`3k5`$ is monomially equivalent to one constructed from one of the following polygons.
Figure 2. Polygons yielding toric codes with $`k=3`$.
Figure 3. Polygons yielding toric codes with $`k=4`$.
Figure 4. Polygons yielding toric codes with $`k=5`$.
The final step in our classification is to show that no two of the toric surface codes constructed from these polygons can be monomially equivalent, hence they lie in distinct monomial equivalence classes. We do this by applying results from §1 to show that the minimum distances (or in some cases, other parts of the complete weight enumerators) are distinct.
###### Theorem 3.6.
Let $`q>5`$. No two of the toric codes $`C_P(𝔽_q)`$ constructed from the polygons in Theorem 3.5 are monomially equivalent.
###### Proof.
If the dimensions are different, the toric codes are certainly not monomially equivalent. Hence we only need to consider each $`k`$ separately.
For the two codes with $`k=3`$, Theorem 2.4 (which applies when $`\mathrm{}=0`$ also) shows $`d(C_{P_3^{(1)}})=(q1)^22(q1)`$. On the other hand, Theorem 2.9 gives $`d(C_{P_3^{(2)}})=(q1)^2(q1)`$. Hence these two codes are not equivalent.
For the codes with $`k=4`$, Theorem 2.4 shows $`d(C_{P_4^{(1)}})=(q1)^23(q1)`$. Corollary 2.10 shows $`d(C_{P_4^{(2)}})=(q1)^22(q1)`$. Theorem 2.4 applies to $`C_{P_4^{(3)}}`$ also, and shows $`d(C_{P_4^{(3)}})=((q1)1)^2=(q1)^2(2q3)`$. Finally, we must analyze $`d(C_{P_4^{(4)}})`$. Write $`C_{P_4^{(4)}}(𝔽_q)=C(𝔽_q)`$. In this case, some more advanced tools are needed. If we translate this polygon by $`(1,1)`$ to place it in $`\mathrm{}_{q1}`$, then we are evaluating polynomials in $`\text{Span}\{1,xy,x^2y,xy^2\}`$ to get the codewords of the corresponding (monomially equivalent) code. Any linear combination of these monomials in which the coefficient of $`x^2y`$ or $`xy^2`$ is nonzero defines an absolutely irreducible curve of degree 3, whose closure in $`^2`$ has arithmetic genus $`1`$ by Theorem 4.2 of . Hence by the general version of the Hasse-Weil bound from , there can be at most $`1+q+2\sqrt{q}`$ $`𝔽_q`$-rational points on the corresponding affine curve. This means that the minimum distance of $`C`$ is at least $`(q1)^2(1+q+2\sqrt{q})`$. On the other hand, the other $`k=4`$ examples have minimum distance no larger than $`d(C_{P_4^{(3)}})=(q1)^2(2q3)`$. It is easy to see from the quadratic formula that $`(1+q+2\sqrt{q})<2q3`$ for all $`q>11`$. Hence $`d(C)`$ is strictly larger than any of the others for $`q>11`$. For the remaining small values of $`q`$ we check directly that $`d(C)`$ is different from the others using the Magma code described in . The results are:
| $`q`$ | $`d(C(𝔽_q))`$ | $`(q1)^2(2q3)`$ |
| --- | --- | --- |
| 7 | 27 | 25 |
| 8 | 40 | 36 |
| 9 | 52 | 49 |
| 11 | 85 | 81 |
Table 1. $`d(C_{P_4^{(4)}}(𝔽_q))`$
These are also different from any of the other $`k=4`$ codes over those fields.
For the $`k=5`$ codes, Theorem 2.4 shows $`d(C_{P_5^{(1)}})=(q1)^24(q1)`$. Corollary 2.10 shows $`d(C_{P_5^{(2)}})=(q1)^23(q1)`$. $`C_{P_5^{(3)}}`$ is a subcode of the code $`C_{P_2^{(2)}}`$ from Theorem 2.9, and contains codewords of the same minimum weight as the supercode. So $`d(C_{P_5^{(3)}})=(q1)^22(q1)`$. The other three $`k=5`$ codes also have $`d(C_{P_5^{(i)}})=(q1)^22(q1)`$, which can be seen, for example, using Minkowski-decomposable subpolytopes (the sets of three collinear points) as in . In $`C_{P_5^{(4)}}`$, for example, we have codewords $`\text{ev}(b(xya_1)(xya_2))`$, where $`a_1,a_2,b𝔽_q^{}`$ and $`a_1a_2`$, which have $`2(q1)`$ zeroes in $`(𝔽_q^{})^2`$.
To show that the four codes with $`d=(q1)^22(q1)`$ are not equivalent, we need to look at finer invariants. For instance, $`C_{P_5^{(5)}}`$ can be distinguished from the other three by the number of words of minimum weight. In $`P_5^{(5)}`$, there are two different sets of three collinear lattice points while in the others, there is only one. This means that there will be more words of the minimum weight in $`C_{P_5^{(5)}}`$ than in $`C_{P_5^{(i)}}`$ for $`i=3,4,6`$. $`C_{P_5^{(5)}}`$ has at least $`2\left(\genfrac{}{}{0pt}{}{q1}{2}\right)(q1)`$ such words because there are two distinct families of reducible polynomials: $`b(xa_1)(xa_2)`$ with $`b,a_i𝔽_q^{}`$ and $`a_1a_2`$ and $`b(ya_1)(y^1a_2)`$ $`b,a_i𝔽_q^{}`$ and $`a_1a_{2}^{}{}_{}{}^{1}`$. On the other hand, $`C_{P_5^{(i)}}`$ for $`i=3,4,6`$ have (at least) $`\left(\genfrac{}{}{0pt}{}{q1}{2}\right)(q1)`$ such words. There are more for some small $`q`$, but never as many as $`2\left(\genfrac{}{}{0pt}{}{q1}{2}\right)(q1)`$. See the weight enumerators for $`C_{P_5^{(6)}}`$ over $`𝔽_{11}`$ and $`𝔽_{16}`$ in Table 2 below. For sufficiently large $`q`$, we claim in fact that there are exactly $`\left(\genfrac{}{}{0pt}{}{q1}{2}\right)(q1)`$ such words. This follows from the general Hasse-Weil bound from . For instance, for $`C_{P_5^{(6)}}`$, if $`q`$ is sufficiently large, then we claim all words in $`C_{P_5^{(6)}}`$ of weight $`(q1)^22(q1)`$ come from evaluations $`\text{ev}(b(ya_1)(ya_2))`$. Any other such word could come only from evaluating a linear combination of $`\{1,x,y,y^2,x^1y^1\}`$ in which $`y^2,x,x^1y^1`$ all appear with nonzero coefficients (since otherwise we are in a case previously covered). Any such curve is absolutely irreducible, of arithmetic genus 2 (because of the 2 interior lattice points in this case, see Theorem 4.2 of ). A simple argument shows that $`1+q+4\sqrt{q}<2q2`$ for all $`q23`$. For smaller values of $`q`$, we verify directly that the weight enumerators of $`C_{P_5^{(5)}}(𝔽_q)`$ do not match the weight enumerators of the other codes using the Magma code from . See Table 2.
To distinguish $`C_{P_5^{(3)}}`$ and $`C_{P_5^{(4)}}`$, we use the codewords of weight $`(q1)^2(2q3)`$ (one more than the minimum weight). Both of these codes contain such words coming from evaluation of the polynomials corresponding to the $`1\times 1`$ squares contained in these polygons (copies of $`P_4^{(3)}`$). Any such square yields $`(q1)^3`$ words of this weight since the polynomials in question have the form $`c(xa)(yb)`$ and $`a,b,c𝔽_q^{}`$ are arbitrary. However $`C_{P_5^{(4)}}`$ has precisely $`(q1)^3`$ words of weight $`(q1)^2(2q3)`$, while $`C_{P_5^{(3)}}`$ has more of them, $`3(q1)^3`$ to be specific. This can be seen by considering the reducible polynomials $`d(xa)(ybxc)`$ that evaluate to give codewords in $`C_{P_5^{(3)}}`$. We get codewords of weight $`(q1)^2(2q3)`$ if $`b=0`$, or if $`c=0`$, or if $`b,c0`$ and $`a=c/b`$.
Finally, to distinguish $`C_{P_5^{(6)}}`$ from the other three codes with $`d=(q1)^22(q1)`$, we must argue as in the last case of the $`k=4`$ codes. If $`q`$ is sufficiently large, then we claim $`C_{P_5^{(6)}}`$ contains no words at all of weight $`(q1)^2(2q3)`$. By Corollary 2.10 and the previous cases, we see that any such word could come only from a linear combination of $`\{1,x,y,y^2,x^1y^1\}`$ in which $`y^2,x,x^1y^1`$ all appear with nonzero coefficients. As before, any such curve is absolutely irreducible, of arithmetic genus 2. A simple argument shows that $`1+q+4\sqrt{q}<2q3`$ for all $`q>23`$. Hence the Hasse-Weil bound from shows that there are no words of this weight for large $`q`$. For smaller values of $`q`$, we again verify directly that the weight enumerators of $`C_{P_5^{(6)}}(𝔽_q)`$ do not match the weight enumerators of the other codes. See Table 2 below.
The following table gives the first three nonzero terms in the weight enumerators:
$$W_C(x)=\underset{i=0}{\overset{(q1)^2}{}}A_ix^i,$$
where $`A_i=|\{wC:\text{wt}(w)=i\}`$, for the $`k=5`$ toric codes with $`d=(q1)^22(q1)`$. These were all computed using the Magma code from .
Over $`𝔽_7`$:
$$\begin{array}{cc}P_5^{(3)}& 1+90x^{24}+648x^{25}+\mathrm{}\\ P_5^{(4)}& 1+90x^{24}+216x^{25}+\mathrm{}\\ P_5^{(5)}& 1+180x^{24}+324x^{26}+\mathrm{}\\ P_5^{(6)}& 1+90x^{24}+432x^{26}+\mathrm{}\end{array}$$
Over $`𝔽_8`$:
$$\begin{array}{cc}P_5^{(3)}& 1+147x^{35}+1029x^{36}+\mathrm{}\\ P_5^{(4)}& 1+147x^{35}+343x^{36}+\mathrm{}\\ P_5^{(5)}& 1+294x^{35}+343x^{37}+\mathrm{}\\ P_5^{(6)}& 1+147x^{35}+1029x^{37}+\mathrm{}\end{array}$$
Over $`𝔽_9`$:
$$\begin{array}{cc}P_5^{(3)}& 1+224x^{48}+1536x^{49}+\mathrm{}\\ P_5^{(4)}& 1+224x^{48}+512x^{49}+\mathrm{}\\ P_5^{(5)}& 1+448x^{48}+512x^{51}+\mathrm{}\\ P_5^{(6)}& 1+224x^{48}+512x^{50}+\mathrm{}\end{array}$$
Over $`𝔽_{11}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+450x^{80}+3000x^{81}+\mathrm{}\\ P_5^{(4)}& 1+450x^{80}+1000x^{81}+\mathrm{}\\ P_5^{(5)}& 1+900x^{80}+1500x^{84}+\mathrm{}\\ P_5^{(6)}& 1+650x^{80}+1000x^{82}+\mathrm{}\end{array}$$
Over $`𝔽_{13}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+792x^{120}+5184x^{121}+\mathrm{}\\ P_5^{(4)}& 1+792x^{120}+1728x^{121}+\mathrm{}\\ P_5^{(5)}& 1+1584x^{120}+7776x^{126}+\mathrm{}\\ P_5^{(6)}& 1+792x^{120}+1728x^{125}+\mathrm{}\end{array}$$
Over $`𝔽_{16}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+1575x^{195}+10125x^{196}+\mathrm{}\\ P_5^{(4)}& 1+1575x^{195}+3375x^{196}+\mathrm{}\\ P_5^{(5)}& 1+3150x^{195}+13500x^{203}+\mathrm{}\\ P_5^{(6)}& 1+2250x^{195}+13500x^{203}+\mathrm{}\end{array}$$
Over $`𝔽_{17}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+1920x^{224}+12288x^{225}+\mathrm{}\\ P_5^{(4)}& 1+1920x^{224}+4096x^{225}+\mathrm{}\\ P_5^{(5)}& 1+3840x^{224}+5120x^{232}+\mathrm{}\\ P_5^{(6)}& 1+1920x^{224}+4096x^{230}+\mathrm{}\end{array}$$
Over $`𝔽_{19}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+2754x^{288}+17496x^{289}+\mathrm{}\\ P_5^{(4)}& 1+2754x^{288}+5832x^{289}+\mathrm{}\\ P_5^{(5)}& 1+5508x^{288}+32076x^{298}+\mathrm{}\\ P_5^{(6)}& 1+2754x^{288}+5832x^{294}+\mathrm{}\end{array}$$
Over $`𝔽_{23}`$:
$$\begin{array}{cc}P_5^{(3)}& 1+5082x^{440}+31944x^{441}+\mathrm{}\\ P_5^{(4)}& 1+5082x^{440}+10648x^{441}+\mathrm{}\\ P_5^{(5)}& 1+10164x^{440}+154396x^{454}+\mathrm{}\\ P_5^{(6)}& 1+5082x^{440}+21296x^{450}+\mathrm{}\end{array}$$
Table 2.
Hence the enumerators never coincide for these four codes, even in exceptional cases for small $`q`$. ∎
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# Critical dynamics of phase transition driven by dichotomous Markov noise
## I Introduction
Over the last decade, the dynamics of ferromagnetic systems below their critical temperatures in a periodically oscillating magnetic field have been studied both theoretically tome ; lo ; acharyya ; sides ; rikvold ; korn1 ; fujisaka1 ; yasui ; tutu and experimentally jiang . The systems exhibit two qualitatively different behaviors referred to as symmetry-restoring oscillation (SRO) and symmetry-breaking oscillation (SBO), depending on the frequency $`\mathrm{\Omega }`$ and the amplitude $`h`$ of the applied magnetic field. It has been established that there exists a sharp transition line between SRO and SBO on the ($`\mathrm{\Omega }`$, $`h`$) plane, which is called the dynamical phase transition (DPT). The DPT was first observed numerically in the deterministic mean-field system for a ferromagnet in a periodically oscillating field tome , and has subsequently been studied in numerous Monte Carlo simulations of the kinetic Ising system below critical temperature lo ; acharyya ; sides ; rikvold ; korn1 . It has also been observed experimentally in an ultra-thin Co film on Cu(100) jiang .
Recently, we investigated the DPT by introducing the model equation $`\dot{s}(t)=(T_cT)ss^3+h\mathrm{cos}\mathrm{\Omega }t`$ fujisaka1 . This equation is a simplified model for the Ising spin system at the temperature $`T`$ below its critical value $`T_c`$ in an external periodic magnetic field. By appropriately scaling the magnetization $`s`$, time $`t`$, and the applied field, this equation is written as
$$\dot{s}(t)=ss^3+h\mathrm{cos}\mathrm{\Omega }t.$$
(1)
The SBO and SRO are observed in Eq. (1) and the transition line between them on the ($`\mathrm{\Omega }`$, $`h`$) plane is determined analytically fujisaka1 ; tutu .
It is quite interesting to ask whether DPT is observed under another kind of applied field, especially random field with bounded amplitude. The fundamental aim of the present paper is to study the dynamics of $`s(t)`$ with a dichotomous Markov noise (DMN) $`F(t)`$ instead of periodically oscillating external field $`h\mathrm{cos}\mathrm{\Omega }t`$ (see, e.g., kampen ).
The equation of motion
$$\dot{s}=f(s)+F(t)$$
(2)
with a nonlinear function $`f(s)`$ and the DMN $`F(t)`$ has been extensively studied by many authors kitahara ; bena ; heureux ; hanggi . It is well known that the master equation for the system can be derived, and then transition phenomena of stationary probability densities concerning the intensity of $`F(t)`$, for example, are studied, which are referred to as the noise-induced phase transition kitahara ; horsthemke . The asymptotic drift velocity $`\dot{s}`$ in the case of $`f(s)`$ being periodic functions are also discussed as a specific dynamic property bena . Furthermore, the mean first-passage time (MFPT) and transition rates are investigated as another important dynamic property when $`f(s)`$ is the force associated with the bistable potential given by Eq. (2heureux ; hanggi . For a review of works on DMN system, see Bena bena2 .
The fundamental aim of the present paper is to propose a phenomenological approach to the critical dynamics near the transition point between the symmetry-restoring motion (SRM) and the symmetry-breaking motion (SBM) observed in Eq. (2) with $`f(s)=ss^3`$ (Sec. II) and to investigate the distribution of passage times through channels, switching times between two different motions, and the power spectrum of the order parameter evolution. The present paper is constructed as follows. In Sec. II, we discuss the dynamics of symmetry-breaking motion and symmetry-restoring motion of the model equation (2). In Sec. III, the jumping process of the magnetization through channels, which are defined in the text, is investigated and the MFPT is obtained. In Sec. IV, a phenomenological approach simplifying the dynamics of passing through a channel in the SRM phase is introduced and three statistical characteristics are analytically developed. The results are compared with numerical simulations. Concluding remarks are given in Sec. V.
## II Model equation and symmetry-breaking transition
### II.1 Model equation and noise-induced phase transition
We consider the equation of motion driven by the external field $`F(t)`$,
$$\frac{ds(t)}{dt}=f(s)+F(t),(f(s)=ss^3)$$
(3)
where $`F(t)`$ is a symmetric DMN with taking the values $`\pm H_0`$. Here the probability $`p(\tau )`$ that $`F(t)`$ continues to take the identical value $`+H_0`$ or $`H_0`$ longer than time $`\tau `$ is given by
$$p(\tau )=e^{\tau /\tau _f}.$$
(4)
This implies that the correlation time of $`F(t)`$ is equal to $`\tau _f/2`$. Throughout this paper, numerical integrations of Eq. (3) are carried out by using the Euler difference scheme with the time increment $`\mathrm{\Delta }t=1/100`$.
Without DMN, $`s(t)`$ eventually approaches either of the stationary fixed points $`\pm 1`$, one of which is achieved according to the initial condition $`s(0)`$ as shown in Fig. 1.
In the presence of DMN, if $`H_0<H_c`$, $`H_c`$ being defined by
$$H_c2(1/3)^{3/2}=0.3849\mathrm{},$$
(5)
then $`f(s)+H_0=0(f(s)H_0=0)`$ has three real roots $`s_{j+}(s_j)`$, ($`j=1`$, 2, and 3). Each value of $`s_{j\pm }`$ is graphically shown in Fig. 1(a). On the other hand, if $`H_0>H_c`$, then $`f(s)+H_0=0`$ ($`f(s)H_0=0`$) has only one real root $`s_+`$ ($`s_{}`$) given by
$$s_\pm =\left[\frac{1}{2}\left(\pm H_0+\sqrt{H_0^2H_c^2}\right)\right]^{1/3}+\left[\frac{1}{2}\left(\pm H_0\sqrt{H_0^2H_c^2}\right)\right]^{1/3},$$
(6)
which are indicated in Fig. 1(b).
Next let us consider the dynamics described by Eq. (3) for $`H_0<H_c`$ and for $`H_0>H_c`$, and discuss similarity and difference between the dynamics in the periodically oscillating field case and those in the present DMN case. A part of our results belongs to the context of the noise-induced phase transition and MFPT in Refs. kitahara ; horsthemke ; bena ; heureux ; hanggi . In the case of $`H_0<H_c`$, three motions numerically integrated are shown in Figs. 2(a) and (b). Two motions confined in the ranges $`s_1<s(t)<s_{1+}`$ and $`s_3<s(t)<s_{3+}`$ are both stable. The long time average $`s(t)`$ of each motion does not vanish, and the motion is called SBM in relation to DPT in the oscillating external field case. On the other hand, the motion $`s_u(t)`$ confined in the range $`s_{2+}<s_u(t)<s_2`$ is unstable. The long time average of $`s_u(t)`$ vanishes, and in this sense the motion is called SRM. It should be noted that this unstable SRM is located between two stable SBM, which has a similar characteristic to SBO of DPT fujisaka1 .
The motion of $`s(t)`$ for $`H_0>H_c`$ is shown in Figs. 2(c) and (d). One observes that there exists a stable SRM confined in the range $`s_{}<s(t)<s_+`$. For SRM, the time average of $`s(t)`$ vanishes, i.e., $`s(t)=0`$. The comparison between Figs. 2(b) and (d) suggests that the SRM for $`H_0>H_c`$ is generated via the “attractor merging crisis” ott of the two SBM’s and one unstable SRM, i.e., the two SBM’s and one unstable SRM disappear and then one stable SRM takes place at $`H_0=H_c`$. This situation is similar to that in the DPT case. However, in contrast to the DPT case, as will be shown in Sec. II.2, the transition line on the ($`\tau _f^1`$, $`H_0`$) plane is independent of the correlation time $`\tau _f`$ of $`F(t)`$ and the average $`s(t)`$ depends discontinuously on $`H_0`$.
### II.2 Stationary distribution functions and phase diagram
In this subsection, we discuss the stationary distribution functions for SBM and SRM. To this aim, we first consider a slightly general nonlinear Langevin equation of motion driven by DMN,
$$\dot{x}(t)=f(x)+g(x)F(t),$$
(7)
where $`f(x)`$ and $`g(x)`$ are generally nonlinear functions of $`x`$ and $`F(t)`$ is DMN markov . The temporal evolution of the distribution function $`P(x,F,t)`$ that $`x(t)`$ and $`F(t)`$ respectively take the values $`x`$ and $`F(=\pm H_0)`$ is determined by kitahara ; horsthemke
$`{\displaystyle \frac{}{t}}P(x,t)`$ $`=`$ $`{\displaystyle \frac{}{x}}\left[f(x)P(x,t)+H_0g(x)q(x,t)\right],`$
$`{\displaystyle \frac{}{t}}q(x,t)`$ $`=`$ $`{\displaystyle \frac{2}{\tau _f}}q(x,t){\displaystyle \frac{}{x}}\left[f(x)q(x,t)+H_0g(x)P(x,t)\right],`$ (8)
where we put $`P(x,t)P(x,+H_0,t)+P(x,H_0,t)`$ and $`q(x,t)P(x,+H_0,t)P(x,H_0,t)`$. The stationary distribution $`P^{st}(x)P(x,\mathrm{})`$ is solved to yield
$$P^{st}(x)=N\frac{g(x)}{H_0^2g(x)^2f(x)^2}\mathrm{exp}\left\{\frac{1}{\tau _f}^x𝑑x^{}\left[\frac{1}{f(x^{})H_0g(x^{})}+\frac{1}{f(x^{})+H_0g(x^{})}\right]\right\},$$
(9)
provided that each of the equations
$$\dot{x}=f(x)+H_0g(x),\dot{x}=f(x)H_0g(x)$$
(10)
has at least one stable fixed point, where $`N`$ is the normalization constant.
By substituting $`f(x)=xx^3`$ and $`g(x)=1`$, (Eq. (3)), into Eq. (9), the stationary distribution function $`P_{SBM}^{st}(s)`$ for SBM ($`H_0<H_c`$) for $`s_3<s<s_{3+}`$ or $`s_1<s<s_{1+}`$ is written as
$`P_{SBM}^{st}(s)`$ $``$ $`|s^2s_{1+}^2|^{\beta _{1+}}|s^2s_1^2|^{\beta _1}|s^2s_{2+}^2|^{\beta _{2+}},`$ (11)
$`\beta _{j\pm }`$ $`=`$ $`1\tau _f^1|(s_{j\pm }s_{k\pm })(s_{j\pm }s_{l\pm })|,`$ (12)
where $`(j,k,l)=(1,2,3)`$, (2,3,1), and (3,1,2).
On the other hand, the stationary distribution function $`P_{SRM}^{st}(s)`$ for the SRM ($`H_0>H_c`$) for $`s_{}<s<s_+`$ is obtained as
$`P_{SRM}^{st}(s)`$ $``$ $`|s^2s_+^2|^{\frac{\tau _f^1}{3s_+^21}1}\left[(s^2+s_+^21)^2s_+^2s^2\right]^{\frac{\tau _f^1}{3s_+^21}1}`$ (13)
$`\times `$ $`\mathrm{exp}\{{\displaystyle \frac{\tau _f^1s_+}{(s_+^21/3)\sqrt{3s_+^24}}}[\mathrm{arctan}\left({\displaystyle \frac{2ss_+}{\sqrt{3s_+^24}}}\right)`$
$`\mathrm{arctan}\left({\displaystyle \frac{2s+s_+}{\sqrt{3s_+^24}}}\right)]\}.`$
The analytic solutions (11) and (13) are numerically confirmed in Fig. 3.
As $`H_0`$ is increased, the form of the stationary distribution function changes drastically from the forms in Eq. (11) to Eq. (13) at $`H_0=H_c`$. This phenomenon which is induced by the disappearance of two pairs of stable and unstable fixed points bena is an example of the noise-induced phase transitions horsthemke .
It turns out that the transition line between SRM and SBM on the ($`\tau _f^1`$, $`H_0`$) plane is given by $`H_0=H_c`$. The phase diagram is given in Fig. 4. Furthermore, the long time average of $`s(t)`$, $`s(t)`$, depends discontinuously on $`H_0`$ at $`H_0=H_c`$ as shown in Fig. 5. These behaviors are quite different from those of the DPT case driven by periodically oscillating field, $`F(t)=h\mathrm{cos}(\mathrm{\Omega }t)`$ acharyya ; fujisaka1 . The transition point $`h_c`$ for a fixed $`\mathrm{\Omega }`$ between SRO and SBO depends on the frequency $`\mathrm{\Omega }`$, and $`s(t)`$ is a continuous function of $`h`$.
## III MFPT through the channels
We hereafter discuss the dynamics for $`H_0`$ slightly above $`H_c`$. Let us first consider the behavior obeying the equations
$$\dot{s}=ss^3+ϵH_0,(ϵ=+\mathrm{or})$$
(14)
for $`H_0>H_c`$, i.e., $`F(t)`$ is fixed to be either $`+H_0`$ or $`H_0`$. Equation (14) for $`H_0>H_c`$ is integrated to yield
$`t`$ $`=`$ $`{\displaystyle \frac{1}{2(3s_ϵ^21)}}\mathrm{ln}{\displaystyle \frac{(ss_ϵ)^2}{s^2+s_ϵs+s_ϵ^21}}{\displaystyle \frac{s_0^2+s_ϵs_0+s_ϵ^21}{(s_0s_ϵ)^2}}`$ (15)
$`+{\displaystyle \frac{6s_ϵ}{2(3s_ϵ^21)\sqrt{3s_ϵ^24}}}\left[\mathrm{arctan}\left({\displaystyle \frac{2s+s_ϵ}{\sqrt{3s_ϵ^24}}}\right)\mathrm{arctan}\left({\displaystyle \frac{2s_0+s_ϵ}{\sqrt{3s_ϵ^24}}}\right)\right],`$
where $`s_0=s(0)`$ and $`s_ϵ`$ has been defined in Eq. (6).
Figure 6 displays three orbits given by Eq. (15) with $`s_0=s_{}`$ and $`ϵ=+`$, which shows that $`s(t)`$ approaches $`s_+`$ in the limit $`t\mathrm{}`$. One observes that $`s(t)`$ stays for a long time in the vicinity of $`s=1/\sqrt{3}`$ for $`H_0`$ slightly above $`H_c`$. The small region including the position $`s=1/\sqrt{3}`$ is called the ‘channel’. From the symmetry of the system, there also exists the channel near $`s=1/\sqrt{3}`$ for $`F(t)=H_0`$, as shown in Fig. 1(b). Let us express the positions $`s_{ch}`$ of the channels as
$$s_{ch}=\{\begin{array}{cc}1/\sqrt{3},& \mathrm{if}F(t)=+H_0\\ +1/\sqrt{3},& \mathrm{if}F(t)=H_0\end{array}.$$
(16)
The characteristic time $`\tau _{ch}`$ is then defined as the time span that the state point $`s(t)`$ passes through one of the channels for a constant $`F(t)`$, either $`+H_0`$ or $`H_0`$. $`\tau _{ch}`$ can be estimated by integrating Eq. (14) around $`ss_{ch}`$ as follows. First, consider the case $`F(t)=H_0`$. By setting $`u(t)=s(t)s_{ch}`$ and assuming $`|u|s_{ch}`$, Eq. (14) is approximated as
$$\dot{u}=3s_{ch}u^2(H_0H_c).$$
(17)
This can be integrated to give
$$u(t)=\sqrt{\frac{H_0H_c}{3s_{ch}}}\mathrm{tan}\left[\sqrt{3s_{ch}(H_0H_c)}t\right]$$
(18)
with the initial condition $`u(0)=0`$. $`\tau _{ch}`$ is estimated by the condition $`u(\tau _{ch})=\mathrm{}`$ and thus
$$\tau _{ch}=\frac{C}{(H_0H_c)^{1/2}},C=\frac{\pi }{2\sqrt{3s_{ch}}}.$$
(19)
Let us next consider the process that the state point $`s(t)`$ passes through the channels under DMN. Figure 7 shows temporal evolutions of $`s(t)`$ numerically obtained for $`H_0=0.388`$ and 0.385. One finds that the time of passing through channels increases as $`H_0`$ approaches $`H_c`$. The MFPT $`\overline{\tau }`$ through channels was calculated in Refs. heureux ; hanggi by analyzing the master equation. In the present section, we will derive MFPT in terms of the time scales $`\tau _f`$ and $`\tau _{ch}`$ from a phenomenological viewpoint without use of the analysis made in Refs. heureux ; hanggi .
The condition for passing through a channel is that $`F(t)`$ continues to take the identical value either $`+H_0`$ or $`H_0`$ for time longer than $`\tau _{ch}`$. For $`H_0`$ satisfying $`\tau _f>\tau _{ch}`$, we obtain $`p(\tau _{ch})=e^{\tau _{ch}/\tau _f}1`$, which implies that $`F(t)`$ almost always satisfies the condition for passing through the channel. Therefore, $`\overline{\tau }`$ in the case of $`\tau _f>\tau _{ch}`$ is nearly equivalent to $`\tau _{ch}`$, i.e.,
$$\overline{\tau }\frac{C}{(H_0H_c)^{1/2}}.$$
(20)
In the case of $`\tau _f\tau _{ch}`$, on the other hand, Eq. (4) gives $`p(\tau _{ch})1`$. This fact implies that the probability that $`F(t)`$ continues to take the identical value for time longer than $`\tau _{ch}`$ is quite small and hence that $`\overline{\tau }`$ is much longer than $`\tau _{ch}`$ because it needs a long time to satisfy the condition for the state point to pass through the channel. $`\overline{\tau }`$ in the case of $`\tau _f\tau _{ch}`$ is explicitly determined as follows. For a long $`\overline{\tau }`$, let us divide $`\overline{\tau }`$ into subintervals each of which has the time span $`\tau _f`$. The divided individual time series are approximately independent of each other. Therefore, $`\tau _f/\overline{\tau }`$ is the probability that the state point passes through a channel once because $`\overline{\tau }`$ is MFPT through the channel. On the other hand, $`p(\tau _{ch})`$ is identical to the probability for $`s(t)`$ to pass through the channel once by definition of the probability. Therefore we get the relation $`p(\tau _{ch})\tau _f/\overline{\tau }`$, which leads to
$$\overline{\tau }^1\tau _f^1e^{\tau _{ch}/\tau _f}=\tau _f^1\mathrm{exp}\left[\frac{C}{\tau _f(H_0H_c)^{1/2}}\right]$$
(21)
with the constant $`C`$ defined in Eq. (19). This expression agrees with the result obtained in Refs. hanggi ; heureux .
Equation (21) reveals that MFPT through the channel depends on $`H_0H_c`$ in a stretched exponential form for $`\tau _f\tau _{ch}`$, and is quite different from the asymptotic form (20). The above dependence of $`\overline{\tau }`$ on $`H_0H_c`$ is confirmed in Fig. 8.
## IV Phenomenological Analysis
In order to discuss statistical characteristics of the dynamics passing through the channels for $`\tau _f\tau _{ch}`$, we here develop a phenomenological approach. The behaviors of $`s(t)`$ for which we attempt to model are first summarized. The initial condition of $`s(t)`$ is set to be in the vicinity of $`s_+`$. If a time interval of $`F(t)`$ satisfying the condition $`F(t)=H_0`$ becomes longer than $`\tau _{ch}`$ for the first time, then $`s(t)`$ passes through $`s_{ch}`$ and approaches $`s_{}`$ in the time interval. See Fig. 9. The event in which $`s(t)`$ jumps from $`s_+`$ to $`s_{}`$ occurs only in this case. It should be noted that the jumps from $`s(t)>0`$ ($`s(t)<0`$) to $`s(t)<0`$ ($`s(t)>0`$) are approximately independent of subsequent jumps.
Let us discretize the time $`t`$ in the form $`t=k\mathrm{\Delta }t`$, ($`k=1`$, 2, 3, $`\mathrm{}`$) as a simple approach to develop the phenomenological analysis according to the process noted above, where $`\mathrm{\Delta }t`$ is a certain small time step. Then, $`\tau _{ch}`$ is discretized as $`\tau _{ch}n_{ch}\mathrm{\Delta }t`$ with the corresponding integer $`n_{ch}`$. $`F(t)`$ is assumed to keep the same value for the interval $`\mathrm{\Delta }t`$, which is denoted as $`F_k=F(k\mathrm{\Delta }t)`$. The conditional probability $`p`$ that $`F_{j+1}`$ takes the same value as $`F_j`$ is given by
$$p=e^{\mathrm{\Delta }t/\tau _f},$$
(22)
and the probability $`q`$ that $`F_{k+1}`$ is different from $`F_k`$ is therefore given by
$$q=1p.$$
(23)
The system is analyzed phenomenologically as follows:
* We introduce the variable $`s_k`$ at a discretized time $`k\mathrm{\Delta }t`$ which takes two values $`\pm 1`$.
* $`s_k`$ and $`F_k`$ are initially set to $`s_0=+1`$ and $`F_0=+H_0`$, respectively.
* $`s_k`$ jumps from $`+1`$ ($`1`$) to $`1`$ ($`+1`$) only if $`F_k`$ continues to take the identical value $`H_0`$ ($`+H_0`$) for a time interval longer than $`n_{ch}\mathrm{\Delta }t`$.
* $`s_k`$ does not jump from $`+1`$ ($`1`$) to $`1`$ ($`+1`$) even though $`F_k`$ continues to take $`+H_0`$ ($`H_0`$) for any time interval longer than $`n_{ch}\mathrm{\Delta }t`$.
### IV.1 MFPT $`\overline{\tau }`$ through the channel
We first derive the exact expression for the MFPT $`\overline{\tau }`$ through the channel with the phenomenological approach. In considering the time series having $`F_k`$, $`\overline{\tau }`$ is evaluated as
$$\overline{\tau }=\underset{l}{}l\mathrm{\Delta }t\underset{0kl}{}g_{k,l}^{(n_{ch})}q^kp^{lk}|_{q=1p}.$$
(24)
Here $`g_{k,l}^{(n_{ch})}`$ is the number of the time sequences $`\{F_j\}`$ for $`0j<l`$ satisfying that $`F_j`$ changed its value $`k`$ times in each $`\{F_j\}`$ and $`s_j`$ jumps from $`+1`$ to $`1`$ for the first time at $`t=l\mathrm{\Delta }t`$.
Equation (24) is, furthermore, rewritten as
$$\overline{\tau }=\widehat{T}Q_{n_{ch}}(q,p)|_{q=1p}$$
(25)
with the differential operator $`\widehat{T}`$ and the quantity $`Q_{n_{ch}}(q,p)`$ defined by
$$\widehat{T}\mathrm{\Delta }t\left(q\frac{}{q}+p\frac{}{p}\right)$$
(26)
and
$$Q_{n_{ch}}(q,p)\underset{l}{}\underset{0kl}{}g_{k,l}^{(n_{ch})}q^kp^{lk}.$$
(27)
One should note that the $`q`$\- and $`p`$-dependences in $`Q_{n_{ch}}`$ are crucial and that $`q`$ and $`p`$ are considered to be independent in Eq. (27).
The explicit form of $`Q_{n_{ch}}(q,p)`$ is then determined so as to satisfy the following conditions:
* In considering any length of time series giving $`F_k`$, there exists a time interval of length $`n_{ch}`$ in the last of the time series, where all the $`F_k`$ take the same value $`H_0`$, i.e., the condition that $`s_k`$ jumps from $`+1`$ to $`1`$ is satisfied.
* The condition for $`s_k`$ to jump from $`+1`$ to $`1`$ is not satisfied before the last time interval.
One should note that the equality $`Q_n(1p,p)=1`$ holds for any $`n`$, because the time interval described above always exists somewhere in a long time series. Particularly, for $`n=n_{ch}`$, $`Q_{n_{ch}}(1p,p)`$ is obviously equal to the probability that $`s_j`$ changes its sign, which must be unity for $`H_0>H_c`$.
As shown in Appendix A, the explicit form of $`Q_n(q,p)`$ is given by
$$Q_n(q,p)=\frac{(1p)qp^{n1}}{(1p)^2q^2(1p^{n1})},$$
(28)
where the condition $`Q_n(1p,p)=1`$ is easily confined. Applying the operator (26) to the explicit form (28) with $`n=n_{ch}`$ yields the relation
$`\overline{\tau }`$ $`=`$ $`\widehat{T}Q_{n_{ch}}(q,p)|_{q=1p}=\mathrm{\Delta }t{\displaystyle \frac{2p^{n_{ch}1}}{(1p)p^{n_{ch}1}}}`$ (29)
$`=`$ $`\mathrm{\Delta }t{\displaystyle \frac{2e^{\mathrm{\Delta }t/\tau _f}e^{\tau _{ch}/\tau _f}}{(1e^{\mathrm{\Delta }t/\tau _f})e^{\mathrm{\Delta }t/\tau _f}e^{\tau _{ch}/\tau _f}}},`$
where the last equality is obtained by using Eqs. (22) and (23) with the relation $`\tau _{ch}=n_{ch}\mathrm{\Delta }t`$. The exact expression of $`\overline{\tau }`$ is finally given by
$$\overline{\tau }=\tau _f\left(2e^{\tau _{ch}/\tau _f}1\right)$$
(30)
in the limit of $`\mathrm{\Delta }t0`$ by keeping $`\tau _{ch}`$ constant. Equation (30) qualitatively agrees for $`\tau _{ch}/\tau _f1`$ with the result (21).
### IV.2 Distribution function $`P(\tau )`$ for the passage time $`\tau `$
The distribution function $`P(\tau )`$ for the passage time $`\tau `$ through the channel $`s_{ch}`$ is determined by solving the equation
$$P(\tau )=\delta (\tau \widehat{T})Q_{n_{ch}}(q,p)|_{q=1p},$$
(31)
where $`\delta (x)`$ is the delta function. The Laplace transform $`[P](z)`$ should be calculated in order to solve Eq. (31). By using the series expansion of $`Q_{n_{ch}}(q,p)`$ given by Eq. (27), the Laplace transform of $`P(\tau )`$ is obtained as
$`[P](z)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}e^{\tau z}P(\tau )𝑑\tau =e^{z\widehat{T}}Q_{n_{ch}}(q,p)|_{q=1p}`$ (32)
$`=`$ $`{\displaystyle \underset{l}{}}{\displaystyle \underset{0kl}{}}g_{k,l}^{(n_{ch})}(e^{z\mathrm{\Delta }t}q)^k(e^{z\mathrm{\Delta }t}p)^{lk}|_{q=1p}.`$
Equation (32) implies that $`[P](z)`$ can be obtained by replacing $`q`$ and $`p`$ by $`e^{z\mathrm{\Delta }t}q`$ and $`e^{z\mathrm{\Delta }t}p`$ in $`Q_{n_{ch}}(q,p)`$, respectively, i.e.,
$$[P(\tau )](z)=Q_{n_{ch}}(e^{z\mathrm{\Delta }t}q,e^{z\mathrm{\Delta }t}p)|_{q=1p}.$$
(33)
Substituting the explicit form (28) for $`n=n_{ch}`$ into Eq. (33) yields the equation
$$[P(\tau )](z)=\frac{(\tau _fz+1)e^{(z+\tau _f^1)\tau _{ch}}}{\tau _f^2z^2+2\tau _fz+e^{(z+\tau _f^1)\tau _{ch}}}$$
(34)
in the limit of $`\mathrm{\Delta }t0`$ by keeping $`\tau _{ch}`$ constant. By applying the inverse Laplace transform to Eq. (34), the distribution function $`P(\tau )`$ is analytically evaluated in the series expansion as
$`P(\tau )`$ $`=`$ $`\tau _f^1e^{\tau /\tau _f}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\theta (t_{k+1}){\displaystyle \frac{(x)^k}{k!}}{\displaystyle \frac{d^k}{dx^k}}\mathrm{cosh}\sqrt{x}|_{x=(t_{k+1})^2}`$ (35)
$`=`$ $`\tau _f^1e^{\tau /\tau _f}[\theta (t_1)\mathrm{cosh}(t_1)\theta (t_2){\displaystyle \frac{t_2\mathrm{sinh}(t_2)}{2}}+\theta (t_3){\displaystyle \frac{t_3{}_{}{}^{2}\mathrm{cosh}(t_3)t_3\mathrm{sinh}(t_3)}{8}}`$
$`\theta (t_4){\displaystyle \frac{t_4{}_{}{}^{3}\mathrm{sinh}(t_4)3t_4{}_{}{}^{2}\mathrm{cosh}(t_4)+3t_4\mathrm{sinh}(t_4)}{48}}+\mathrm{}],`$
where $`t_k(\tau )(\tau k\tau _{ch})/\tau _f`$ and $`\theta (t)`$ is the Heaviside function defined by
$$\theta (t)=\{\begin{array}{ccc}1& \text{for}& t0\\ 0& \text{for}& t<0\end{array}.$$
(36)
For details of the derivation of Eq. (35), see Appendix B.
Let us suppose to truncate the expansion (35) at $`k=k_c`$ for an arbitrary $`k_c`$. It should be noted that Eq. (35) gives the exact distribution for $`0<\tau <k_c\tau _{ch}`$ even though the truncation is executed, since all the terms individually include $`\theta (t_k)`$ and so the terms for $`k>k_c`$ do not contribute to $`P(\tau )`$ for $`\tau <k_c\tau _{ch}`$. The analytical result (35) is compared with the numerically evaluated distribution in Fig. 10. One observes that the phenomenological analysis quantitatively explains the statistical property of passing through the channels. The characteristics obtained from the figure are summarized as follows.
* There exists a region where $`P(\tau )=0`$ for $`\tau <\tau _{ch}`$, which presents the minimal time of passing through the channels.
* $`P(\tau )`$ decreases exponentially for $`\tau \tau _{ch}`$, $`P(\tau )e^{\alpha \tau }`$ with a constant $`\alpha `$.
* The rate $`\alpha `$ increases as $`\tau _f`$ is increased. This tendency is consistent with the fact that the probability of passing through channels increases as $`\tau _f`$ is increased since DMN will often continue to take an identical value longer than $`\tau _{ch}`$.
On the other hand, the expansion (35) disagrees with the correct value in an exponential way for $`\tau >k_c\tau _{ch}`$. Let us try to obtain the asymptotic solution of $`P(\tau )`$ for $`\tau \tau _{ch}`$. Equation (62) is approximated as
$$e^{z(\tau \tau _{ch})}\frac{1}{1+(\overline{\tau }\tau _{ch})z}\text{for }|z|\tau _{ch}^1,$$
(37)
where $`\overline{\tau }`$ is MFPT given in Eq. (30). The inverse Laplace transform of Eq. (37) is straightforwardly calculated to give
$$P(\tau )\frac{1}{\overline{\tau }\tau _{ch}}\mathrm{exp}\left(\frac{\tau \tau _{ch}}{\overline{\tau }\tau _{ch}}\right),$$
(38)
which reveals that $`P(\tau )`$ decreases exponentially with the damping rate $`\alpha =(\overline{\tau }\tau _{ch})^1`$ for $`\tau \tau _{ch}`$.
### IV.3 Fourier spectrum determined by the phenomenological analysis
We derive the Fourier spectrum of a time series $`s(t)`$ by the phenomenological analysis to focus on the dynamical characteristics in the SRM phase. The Fourier spectrum $`I_x(\omega )`$ is defined by
$$I_x(\omega )=\underset{T\mathrm{}}{lim}\frac{1}{T}\left|_0^Tx(t)e^{i\omega t}𝑑t\right|^2,$$
(39)
i.e., the ensemble average of the Fourier transform of a time series $`x(t)`$.
Let us first consider $`s_0(t)\mathrm{sgn}[s(t)]`$. Then the time series $`s_0(t)`$ is expressed as
$$s_0(t)=(1)^{n1},\text{for }t_{n1}t<t_n$$
(40)
with $`n1`$, where $`t_n`$ denotes the $`n`$th time to cross zero for $`s(t)`$. Hereafter, $`t_0`$ is set to be zero without loss of generality. By identifying that $`\tau _nt_nt_{n1}`$ is independently distributed according to Eq. (35), one obtains the Fourier spectrum of $`s_0(t)`$ by the phenomenological analysis shown in Appendix C, in the form
$$I_{s_0}(\omega )=\frac{4}{\overline{\tau }\omega ^2}\mathrm{}\left(\frac{1e^{i\omega \tau _n}}{1+e^{i\omega \tau _n}}\right)=\frac{4}{\overline{\tau }\omega ^2}\frac{1|e^{i\omega \tau _n}|^2}{|1+e^{i\omega \tau _n}|^2}$$
(41)
where $`\mathrm{}(X)`$ represents the real part of $`X`$, and $`lim_N\mathrm{}\frac{t_N}{N}=\tau _n=\overline{\tau }`$ is used.
Substituting the explicit form of $`e^{i\omega \tau _n}`$ given in Eq. (34) with $`z=i\omega `$ into Eq. (41) yields
$$I_{s_0}(\omega )=\left(\frac{4\tau _f}{\overline{\tau }\omega }\right)\frac{\omega ^3\tau _f^3+(4e^{2\tau _{ch}/\tau _f})\omega \tau _f2e^{\tau _{ch}/\tau _f}(\omega \tau _f\mathrm{cos}\omega \tau _{ch}+2\mathrm{sin}\omega \tau _{ch})}{(4+\omega ^2\tau _f^2)(\omega ^2\tau _f^22\omega \tau _fe^{\tau _{ch}/\tau _f}\mathrm{sin}\omega \tau _{ch}+e^{2\tau _{ch}/\tau _f})}.$$
(42)
The above result is confirmed by comparing with the numerically evaluated Fourier spectrum for the normalized time series $`s_0(t)`$ in Fig. 11.
Let us finally modify the phenomenological analysis which is compatible with the numerically evaluated spectrum of the original time series $`s(t)`$ without normalization. Instead of Eq. (40), let us define
$$\stackrel{~}{s}(t)=(1)^{n1}[1a(tt_{n1})]\text{for }t_{n1}t<t_n$$
(43)
with $`n1`$, where $`a(\mathrm{\Delta }t)`$ incorporates the wave form of the time series passing through the channel and is assumed to be $`a(\mathrm{\Delta }t)=0`$ for $`\mathrm{\Delta }t>\tau _{ch}`$. Note that by setting $`a(\mathrm{\Delta }t)=0`$ also for $`\mathrm{\Delta }t\tau _{ch}`$ the result of original phenomenological analysis is recovered. As shown in appendix C, the Fourier spectrum $`I_{\stackrel{~}{s}}(\omega )`$ for $`\stackrel{~}{s}(t)`$ as a modification to $`I_{s_0}(\omega )`$ is obtained in the form
$$I_{\stackrel{~}{s}}(\omega )=I_{s_0}(\omega )\frac{1+|\widehat{a}(\omega )|^2+2\mathrm{}[\widehat{a}(\omega )]}{4},$$
(44)
where
$$\widehat{a}(\omega )1i\omega _0^{\tau _{ch}}a(t)e^{i\omega t}𝑑t.$$
(45)
By approximating as $`a(\mathrm{\Delta }t)=1+|s_{ch}|`$ for $`0<\mathrm{\Delta }t<\tau _{ch}`$, Eq. (44) reduces to
$$I_{\stackrel{~}{s}}(\omega )=I_{s_0}(\omega )\left(\frac{1+s_{ch}^2}{2}+\frac{1s_{ch}^2}{2}\mathrm{cos}\omega \tau _{ch}\right).$$
(46)
Equation (46) is confirmed by comparing with the numerically evaluated Fourier spectrum in Fig. 12.
## V Concluding remarks
In this paper, we used the model equation (3) under the dichotomous Markov noise (DMN) $`F(t)`$ with a finite correlation time in order to investigate the dynamics of the magnetization $`s(t)`$ of the ferromagnet system driven by the magnetic field applied in one direction, where its strength is constant and only the direction temporally changes. It was found that the dynamics of $`s(t)`$ show two kinds of motion, i.e., the symmetry-restoring motion (SRM) and the symmetry-breaking motion (SBM), which are respectively observed when $`H_0`$ is above and below the critical value $`H_c`$. The transition line between SRM and SBM was determined only by the strength of the applied DMN and is independent of the correlation time $`\tau _f`$ of DMN. By observing the distribution functions of $`s(t)`$ for SRM and SBM, the ensemble average of $`s(t)`$ discontinuously changes at $`H_0=H_c`$. These results are quite different from those in the system driven by a periodically oscillating field fujisaka1 .
We then discussed the mean first passage time (MFPT) slightly above $`H_c`$ and found that it depends on $`H_0H_c`$ and $`\tau _f`$ as
$$\overline{\tau }\tau _f\mathrm{exp}\left[\frac{C}{\tau _f(H_0H_c)^{1/2}}\right].$$
(47)
This anomalous characteristics has a form similar to that of the average duration between neighboring phase slips of the phase difference in the phase synchronization observed in coupled chaotic systems rosenblum . Furthermore, a phenomenological approach was proposed to analytically discuss the statistical characteristics for $`H_0>H_c`$. By obtaining the probability $`Q_n(q,p)`$ of $`s(t)`$ passing through the channel, the MFPT $`\overline{\tau }`$, the distribution function $`P(\tau )`$ of the passage time $`\tau `$ and the Fourier spectrum $`I(\omega )`$ of the time series were obtained by using the phenomenological analysis. The statistics obtained by the phenomenological analysis were found to be not only in qualitative but also in quantitative agreement with the numerically obtained results.
In closing the paper, it is worth noting that the effect of DMN on nonlinear dynamical systems is generally quite different from that of Gaussian noise and DMN produces a new dynamical response of the systems. It is highly desired to examine the statistical characteristics obtained in this paper in laboratory experiments, and also in other numerical simulations, e.g., Monte Carlo simulations.
###### Acknowledgements.
The authors thank N. Tsukamoto, N. Fujiwara, and H. Hata for valuable comments. This study was partially supported by Grant-in-Aid for Scientific Research (C) of the Ministry of Education, Culture, Sports, Science, and Technology, and the 21st Century COE Program “Center of Excellence for Research and Education on Complex Functional Mechanical Systems” at Kyoto University.
## Appendix A Explicit form of $`Q_n(q,p)`$
First of all, notice that each sample path of $`F_k`$ ($`k=0,1,2,\mathrm{}`$) corresponds to a symbol sequence of $`\{+,\}`$. Let $`w`$ be a symbol sequence of $`\{+,\}`$, which is referred to as a string. For given sets $`A`$ and $`B`$ of strings, let $`AB`$ be the set of all strings expressed as $`ab`$ for $`aA`$ and $`bB`$, where $`ab`$ denotes the concatenated string of $`a`$ followed by $`b`$. In the following, the set composed of only one string $`w`$ will be simply expressed as $`w`$. Furthermore, let $`(A)_n`$ be the set of all strings expressed as
$$a_1a_2\mathrm{}a_k$$
(48)
with $`a_iA`$ and $`0kn`$, where $`k=0`$ means the zero length string, including the case $`a_i=a_j`$ for $`ij`$. Then, with this notation, every possible sequence starting with $`+`$ and terminating with successive $``$’s of length $`n`$ $`(n2)`$, which firstly appears in that string, can be expressed as
$$S_n(+(+)_{\mathrm{}}()_{n2})_{\mathrm{}}+(+)_{\mathrm{}}\stackrel{n\text{ symbols}}{\stackrel{}{\mathrm{}}}.$$
(49)
For a given set $`S`$ of strings, let us define its “probability” as a function of $`p`$ and $`q`$ by
$$P(q,p;S)\underset{wS}{}P_w(p,q)$$
(50)
where $`P_w(p,q)p^{k(w)}q^{l(w)}`$ denotes the probability for a string $`w`$, and $`k(w)`$ and $`l(w)`$ denote the numbers of pairs of identical symbols ($``$, $`++`$) and different symbols ($`+`$, $`+`$) appearing in $`w`$, respectively. For given strings $`w_1`$ and $`w_2`$, obviously, the identity
$$P_{w_1w_2}(q,p)=P_{w_1\sigma }(q,p)P_{w_2}(q,p)$$
(51)
holds with $`\sigma `$ being the first symbol in $`w_2`$, and we call $`w_1\sigma `$ and $`w_2`$ a decomposition of the string $`ww_1w_2`$. By noting that
$$Q_n(q,p)=P(q,p;S_n)$$
(52)
and considering decompositions of each element in $`S_n`$, $`Q_n(q,p)`$ can be expressed as
$$Q_n(q,p)=P(q,p;S_n^{})P(q,p;R)p^{n1},$$
(53)
where
$$S_n^{}(+(+)_{\mathrm{}}()_{n2})_{\mathrm{}}+$$
(54)
and
$$R+(+)_{\mathrm{}}.$$
(55)
Since each element of $`S_n^{}`$ can be uniquely decomposed into a multiple of elements in
$$S_n^{\prime \prime }+(+)_{\mathrm{}}()_{n2}+$$
(56)
and inversely every multiple of elements in $`S_n^{\prime \prime }`$ uniquely corresponds to an element in $`S_n^{}`$ as its decomposition, we obtain
$`P(q,p;S_n^{})`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left[P(q,p;S_n^{\prime \prime })\right]^j`$ (57)
$`=`$ $`{\displaystyle \frac{1}{1P(q,p;S_n^{\prime \prime })}}.`$ (58)
Finally, $`P(q,p;R)`$ and $`P(q,p;S_n^{\prime \prime })`$ are calculated as
$$P(q,p;R)=q\underset{j=0}{\overset{\mathrm{}}{}}p^j=\frac{q}{1p}$$
(59)
and
$$P(q,p;S_n^{\prime \prime })=\underset{j=0}{\overset{\mathrm{}}{}}p^jq\underset{i=0}{\overset{n2}{}}p^iq=q^2\frac{1p^{n1}}{(1p)^2},$$
(60)
which yield Eq. (28) together with (53) and (58).
## Appendix B Derivation of the probability distribution function $`P(\tau )`$
Let us evaluate the inverse Laplace transform of Eq. (34) in order to derive the explicit form of $`P(\tau )`$. For simplicity, the time is rescaled as $`\tau _f=1`$. Then Eq. (34) is rewritten as
$$e^{z\tau }=\frac{(1+z)e^{(1+z)\tau _{ch}}}{(1+z)^2(1e^{(1+z)\tau _{ch}})},$$
(61)
which reads
$$[P(\tau +\tau _{ch})]=e^{z(\tau \tau _{ch})}=\frac{e^{\tau _{ch}}}{1+z\frac{1e^{(1+z)\tau _{ch}}}{1+z}}=e^{\tau _{ch}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\left(1e^{(1+z)\tau _{ch}}\right)^n}{(1+z)^{2n+1}}.$$
(62)
Repeatedly applying the formula
$$^1\left[e^{(1+z)\tau _{ch}}\widehat{f}(z)\right]=e^{\tau _{ch}}f(\tau \tau _{ch})=e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }}f(\tau )=e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }}^1[\widehat{f}(z)]$$
(63)
with $`f(\tau )^1[\widehat{f}(z)]`$, one obtains
$`^1\left[{\displaystyle \frac{\left(1e^{(1+z)\tau _{ch}}\right)^n}{(1+z)^{2n+1}}}\right]`$ $`=`$ $`(1e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }})^n^1\left[{\displaystyle \frac{1}{(1+z)^{2n+1}}}\right]`$ (64)
$`=`$ $`(1e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }})^n{\displaystyle \frac{\tau ^{2n}}{(2n)!}}e^\tau \theta (\tau ),`$
where $`\theta (\tau )`$ denotes the Heaviside function Eq. (36), and the formula
$$^1\left[\frac{1}{(1+z)^m}\right]=\frac{\tau ^{m1}}{(m1)!}e^\tau \theta (\tau )$$
(65)
for any positive integer $`m`$ was used. Thus, the inverse Laplace transform of Eq. (62) reads
$`P(\tau +\tau _{ch})`$ $`=`$ $`e^{\tau _{ch}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(1e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }}\right)^n{\displaystyle \frac{\tau ^{2n}}{(2n)!}}e^\tau \theta (\tau )`$ (66)
$`=`$ $`e^{(\tau +\tau _{ch})}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(1e^{\tau _{ch}\frac{d}{d\tau }}\right)^n{\displaystyle \frac{\tau ^{2n}}{(2n)!}}\theta (\tau ),`$
where the identity
$$e^{\tau _{ch}}e^{\tau _{ch}\frac{d}{d\tau }}e^\tau =e^\tau e^{\tau _{ch}\frac{d}{d\tau }}$$
(67)
was applied. By noting the identity
$$\left(1e^{\tau _{ch}\frac{d}{d\tau }}\right)^n=\underset{k=0}{\overset{n}{}}\frac{(e^{\tau _{ch}\frac{d}{d\tau }})^k}{k!}\frac{d^kx^n}{dx^k}|_{x=1}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{e^{k\tau _{ch}\frac{d}{d\tau }}(1)^k}{k!}\left(\frac{d}{dx}\right)^kx^n|_{x=1}$$
(68)
Eq. (66) is further simplified as
$`P(\tau +\tau _{ch})`$ $`=`$ $`e^{(\tau +\tau _{ch})}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}e^{k\tau _{ch}\frac{d}{d\tau }}{\displaystyle \frac{(1)^k}{k!}}\left({\displaystyle \frac{d}{dx}}\right)^kx^n{\displaystyle \frac{\tau ^{2n}}{(2n)!}}\theta (\tau )|_{x=1}`$ (69)
$`=`$ $`e^{(\tau +\tau _{ch})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}e^{k\tau _{ch}\frac{d}{d\tau }}{\displaystyle \frac{(1)^k}{k!}}\left({\displaystyle \frac{d}{dx}}\right)^k\mathrm{cosh}(\sqrt{x}\tau )\theta (\tau )|_{x=1}`$
$`=`$ $`e^{(\tau +\tau _{ch})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}e^{k\tau _{ch}\frac{d}{d\tau }}\theta (\tau ){\displaystyle \frac{(\tau ^2)^k}{k!}}\left({\displaystyle \frac{d}{dx}}\right)^k\mathrm{cosh}\sqrt{x}|_{x=\tau ^2}`$
$`=`$ $`e^{(\tau +\tau _{ch})}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\theta (t_k){\displaystyle \frac{(x)^k}{k!}}{\displaystyle \frac{d^k}{dx^k}}\mathrm{cosh}\sqrt{x}|_{x=t_k^2},`$
where $`t_k(\tau )\tau k\tau _{ch}`$ and the formula $`_{n=0}^{\mathrm{}}\frac{x^n}{(2n)!}=\mathrm{cosh}\sqrt{x}`$ was used. After the replacement $`\tau \tau \tau _{ch}`$ and the rescaling of time as $`\tau \tau /\tau _f`$ and $`\tau _{ch}\tau _{ch}/\tau _f`$, Eq. (69) reduces to Eq. (35).
## Appendix C Derivation of the Fourier Spectrum $`I(\omega )`$
In this appendix, we derive the Fourier spectrum of the time series of the magnetization $`s(t)`$ for the phenomenological analysis. Let $`\tau _1,\tau _2,\tau _3,\mathrm{}`$ be a sequence of mutually independent random variables having the probability density $`P(\tau )`$ obeying the condition $`P(\tau )=0`$ for $`\tau <\tau _{ch}`$, i.e., $`\tau _k\tau _{ch}`$. We introduce the variable
$$\stackrel{~}{s}(t)=(1)^n[1a(tt_{n1})]\text{ for }t_{n1}t<t_n,$$
(70)
where $`t_n_{k=1}^N\tau _k`$ and $`a(\mathrm{\Delta }t)`$ is a function satisfying $`a(\mathrm{\Delta }t)=0`$ for $`\mathrm{\Delta }t>\tau _{ch}`$. We assume that the time series of the magnetization $`s(t)`$ is approximately expressed by $`\stackrel{~}{s}(t)`$ with an appropriate form of $`a(t)`$. The Fourier transform of $`\stackrel{~}{s}(t)`$ follows
$`{\displaystyle _0^{t_N}}\stackrel{~}{s}(t)e^{i\omega t}𝑑t`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}(1)^{n1}{\displaystyle _{t_{n1}}^{t_n}}[1a(tt_{n1})]e^{i\omega t}𝑑t`$ (71)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}(1)^{n1}e^{i\omega t_{n1}}\left[{\displaystyle _0^{\tau _n}}e^{i\omega t}𝑑t{\displaystyle _0^{\tau _{ch}}}a(t)e^{i\omega t}𝑑t\right]`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}(1)^ne^{i\omega t_{n1}}{\displaystyle \frac{e^{i\omega \tau _n}\widehat{a}(\omega )}{i\omega }},`$
where
$$\widehat{a}(\omega )1i\omega _0^{\tau _{ch}}a(t)e^{i\omega t}𝑑t.$$
(72)
Considering the absolute square of Eq. (71) and then taking the ensemble average, we obtain
$`\omega ^2\left|{\displaystyle _0^{t_N}}\stackrel{~}{s}(t)e^{i\omega t}𝑑t\right|^2={\displaystyle \underset{n=1}{\overset{N}{}}}\left|e^{i\omega \tau _n}\widehat{a}(\omega )\right|^2`$ (73)
$`+2\mathrm{}\left[{\displaystyle \underset{1m<nN}{}}(1)^{nm}e^{i\omega _{k=m+1}^{n1}\tau _k}\left(e^{i\omega \tau _n}\widehat{a}(\omega )\right)\left(1e^{i\omega \tau _m}\widehat{a}^{}(\omega )\right)\right]`$
and
$`{\displaystyle \frac{\omega ^2}{N}}\left|{\displaystyle _0^{t_N}}\stackrel{~}{s}(t)e^{i\omega t}𝑑t\right|^2=1+|\widehat{a}(\omega )|^22\mathrm{}\left[e^{i\omega \tau _n}\widehat{a}^{}(\omega )\right]`$ (74)
$`+2N^1\mathrm{}\left[{\displaystyle \underset{1m<nN}{}}(1)^{nm}e^{i\omega \tau _k}^{nm1}\left(e^{i\omega \tau _k}\widehat{a}(\omega )\right)\left(1e^{i\omega \tau _k}\widehat{a}^{}(\omega )\right)\right].`$
Since $`|e^{i\omega \tau _n}|<1`$, in the limit of $`N\mathrm{}`$, the last term in Eq. (74) reads
$`2\mathrm{}\left[{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1)^ke^{i\omega \tau _n}^k\left(e^{i\omega \tau _n}\widehat{a}(\omega )\right)\left(1e^{i\omega \tau _n}\widehat{a}^{}(\omega )\right)\right]`$ (75)
$`=2\mathrm{}\left[{\displaystyle \frac{e^{i\omega \tau _n}\left(1+|\widehat{a}(\omega )|^2\right)\widehat{a}(\omega )e^{i\omega \tau _n}^2\widehat{a}^{}(\omega )}{1+e^{i\omega \tau _n}}}\right],`$
which leads to
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\left|_0^{t_N}\stackrel{~}{s}(t)e^{i\omega t}𝑑t\right|^2=\omega ^2\mathrm{}\left[\frac{1e^{i\omega \tau _n}}{1+e^{i\omega \tau _n}}\right]\left(1+|\widehat{a}(\omega )|^2+2\mathrm{}[\widehat{a}(\omega )]\right).$$
(76)
Noting that $`lim_N\mathrm{}\frac{t_N}{N}=\tau _n=\overline{\tau }`$, we obtain
$$I_{\stackrel{~}{s}}(\omega )=I_{s_0}(\omega )\frac{1+|\widehat{a}(\omega )|^2+2\mathrm{}[\widehat{a}(\omega )]}{4},$$
(77)
where
$$I_{s_0}(\omega )\frac{4}{\omega ^2\overline{\tau }}\mathrm{}\left[\frac{1e^{i\omega \tau _n}}{1+e^{i\omega \tau _n}}\right]$$
(78)
corresponds to the case that $`a(\mathrm{\Delta }t)=0`$ for all $`\mathrm{\Delta }t0`$ and thus $`s_0(t)=(1)^n`$ for $`t_{n1}t<t_n`$.
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# Insulating spin liquid in the lightly doped two-dimensional Hubbard model
## Abstract
We calculate the charge compressibility and uniform spin susceptibility for the two-dimensional (2D) Hubbard model slightly away from half-filling within a two-loop renormalization group scheme. We find numerically that both those quantities flow to zero as we increase the initial interaction strength from weak to intermediate couplings. This result implies gap openings in both charge and spin excitation spectra for the latter interaction regime. When this occurs, the ground state of the lightly doped 2D Hubbard model may be interpreted as an insulating spin liquid as opposed to a Mott insulating state.
After two decades of intensive research on the high-Tc superconductors, physicists are still puzzled by some of their very unusual electronic properties Millis . The prominent example is given by the cuprates. At zero doping, despite the fact that their highest occupied band is half-filled, they are charge insulators, and display antiferromagnetic long-range order. For this reason, they are said to be Mott insulators. As soon as one starts doping those compounds with holes, the long-range magnetic order becomes rapidly suppressed, and there are experimental evidences of an emergence of a spin gap in their corresponding excitation spectra Timusk . A charge gap is also observed by ARPES experiments in such lightly-doped systems Shen . Moreover, at finite temperatures, they turn themselves into poor conductors with electronic properties differing considerably from the predictions of Landau’s Fermi liquid theory. This scenario configures the so-called pseudogap regime. Although this phase continues to be not well understood, it is widely acknowledged to play a fundamental role in the underlying microscopic mechanism of such high-Tc superconductors. Indeed, upon some further doping, those poor metals become superconducting with $`d`$-wave symmetry up to relatively high temperatures around the optimal doping level.
From the theoretical viewpoint, it is widely accepted that the appropriate model for describing such systems is the two-dimensional (2D) Hubbard model (HM), since it is known to have a Mott insulating phase at half-filling, and is expected to become a $`d`$-wave superconductor at larger doping Zanchi . However, its intermediate doping regime, which could provide some insight to understand the physical nature of the pseudogap state, still remains elusive to this date. In this Letter we intend to address this question using renormalization group (RG) techniques in order to infer about the ground state of such model for electron densities slightly away from the half-filling limit.
Our considerations here will be based on a complete two-loop RG calculation of the uniform charge (CS) and spin (SS) susceptibilities of the 2D HM, taking into account simultaneously both the renormalization of the couplings, and the self-energy effects. (The CS is also called the charge compressibility of the system.) To best of our knowledge, it is the first time that such a full two-loop RG calculation is performed for the 2D lightly-doped HM, since previous estimates of the uniform susceptibilities followed random phase approximation (RPA) schemes.
In momentum space, the 2D Hubbard Hamiltonian on a square lattice is given by
$$H=\underset{𝐤,\sigma }{}\xi _𝐤\psi _{𝐤\sigma }^{}\psi _{𝐤\sigma }+\left(\frac{U}{N_{sites}}\right)\underset{𝐩,𝐤,𝐪}{}\psi _{𝐩+𝐤𝐪}^{}\psi _𝐪^{}\psi _𝐤\psi _𝐩,$$
(1)
where the energy dispersion is simply $`\xi _𝐤=2t\left[\mathrm{cos}(k_xa)+\mathrm{cos}(k_ya)\right]\mu `$, and $`\psi _{𝐤\sigma }^{}`$ and $`\psi _{𝐤\sigma }`$ are the usual creation and annihilation operators of electrons with momentum $`𝐤`$ and spin projection $`\sigma =,`$. Besides, $`\mu `$ stands for the chemical potential, whereas $`a`$ is the square lattice spacing. Another important parameter here is the width of the noninteracting band, which is given by $`W=8t`$. This model describes a system with many electrons interacting mutually via a local repulsive interaction $`U`$, and with a total number $`N_{sites}`$ of lattice sites. The electron band filling of the system is controlled by the ratio $`\mu /t`$. When $`\mu /t=0`$ the system is exactly at half-filling. As we start doping it with holes, $`\mu /t`$ takes slightly negative values.
Our starting point is a 2D nearly flat Fermi surface (FS) with no van Hove singularities (see Fig. 1). It correctly describes the 2D HM slightly away from the half-filling case. A similar FS has already been used by other groups to investigate the leading instabilities within either a parquet or, equivalently, a one-loop RG approach Yakovenko . All those investigations find diverging susceptibilities at finite energies (or finite temperatures) with the dominant instability being always the spin density wave (SDW). Their interpretation is that this implies a spontaneous symmetry breaking in the system, and the onset of a long-range ordered antiferromagnetic state.
In this Letter we argue that is not necessarily true since this result may also be an indicative of the limitations of the one-loop RG scheme. In low-dimensional systems, large quantum fluctuations are expected to suppress long-range order. The more those effects are taken into account, the more likely those long-range ordered states are transformed into short-range magnetically ordered phases. This also becomes clear as a result of our work. Taking into account quantum fluctuation effects up to two-loop order in our RG scheme, we are able to show that, for moderate coupling regimes, both charge compressibility and uniform spin susceptibility are strongly suppressed and flow unequivocally to zero. This behavior implies that there are gaps for both charge and spin excitations, and no trace of long-range symmetry breaking order in those cases. Such a state with a fully gapped charge and spin spectra is usually denominated an insulating spin liquid (ISL). The ISL is an example of a short-range resonant valence bond state, which was first proposed for a S=1/2 Heisenberg model Anderson , and clearly revealed in even-leg Hubbard ladders by both RG and bosonization approaches Schulz ; Balents ; Dagotto .
In order to implement a full RG calculation of the uniform susceptibilities, it is essential to consider at least two-loop order contributions. This is due to the fact that, at one-loop level, there is not a single infrared (IR) divergent diagram in the calculation of the so-called uniform response functions. In contrast, there are several of those diagrams in two loops (the nonparquet diagrams), and, as a result, one can reliably begin to derive appropriate RG flow equations for those quantities at that order of perturbation theory.
To keep a closer contact with well-known works in one-dimensional systems Solyom , we divide the FS into four different regions (two sets of solid and dashed line patches). Here we restrict the momenta at the FS to the flat parts only. The interaction processes connecting parallel patches of the FS are always logarithmically IR divergent due to quantum fluctuations. In contrast, those connecting perpendicular patches always remain finite, and do not contribute to the RG flow equations in our approach. For convenience, we restrict ourselves to one-electron states labeled by the momenta $`p_{}=k_x`$ and $`p_{}=k_y`$ associated with one of the two sets of perpendicular patches. The momenta parallel to the FS are restricted to the interval $`\mathrm{\Delta }p_{}\mathrm{\Delta }`$, with $`2\mathrm{\Delta }`$ being essentially the size of the flat patches. The energy dispersion of the single-particle states is given by $`\epsilon _a\left(𝐩\right)=v_F\left(\left|p_{}\right|k_F\right)`$, and depends only on the momenta perpendicular to the FS. The label $`a=\pm `$ refers to the flat sectors at $`p_{}=\pm k_F`$, respectively. In addition, we take $`k_F\lambda \left|p_{}\right|k_F+\lambda `$, where $`\lambda `$ is a fixed ultraviolet (UV) microscopic momentum cut-off.
We now write down the Lagrangian of the 2D HM as
$`L=`$ $`{\displaystyle \underset{𝐩,\sigma ,a=\pm }{}}\psi _{(a)\sigma }^{}\left(𝐩\right)\left[i_tv_F\left(\left|p_{}\right|k_F\right)\right]\psi _{(a)\sigma }\left(𝐩\right)`$
$`{\displaystyle \frac{1}{V}}{\displaystyle \underset{𝐩,𝐪,𝐤}{}}{\displaystyle \underset{\alpha ,\beta ,\delta ,\gamma }{}}\left[g_2\delta _{\alpha \delta }\delta _{\beta \gamma }g_1\delta _{\alpha \gamma }\delta _{\beta \delta }\right]`$
$`\times \psi _{(+)\delta }^{}\left(𝐩+𝐪𝐤\right)\psi _{()\gamma }^{}\left(𝐤\right)\psi _{()\beta }\left(𝐪\right)\psi _{(+)\alpha }\left(𝐩\right),`$
where $`\psi _{(\pm )}^{}`$ and $`\psi _{(\pm )}`$ are now fermionic fields associated to electrons located at the $`\pm `$ patches. The summation over momenta must be appropriately understood as $`_𝐩=V/(2\pi )^2d^2𝐩`$ in the thermodynamic limit. We linearized the energy dispersion about the lightly-doped FS, and the interaction term was parametrized in a manifestly SU(2) invariant form. Here we follow the well-known g-ology notation, with $`g_1`$ and $`g_2`$ standing for backscattering and forward scattering couplings, respectively. Since this should represent the 2D HM, these couplings must be initially defined as $`g_1=g_2=(V/4N_{sites})U`$. In addition, we do not include Umklapp processes since we are not at half-filling, and our FS does not intersect the so-called Umklapp surface at any point (see Fig. 1(a)). Consequently, our result is different from another evidence of ISL behavior reported in the literature in the context of the $`tt^{}`$ 2D HM Furukawa . In their case, the FS intersects the Umklapp surface even away from half-filling, and, for this reason, Umklapp processes become an essential ingredient for the correct description of the system at the lightly-doped regime. Moreover, their estimates of the uniform susceptibilities follow the already mentioned RPA-like scheme.
Since the HM is a microscopic model, all terms in the Lagrangian are defined at a scale of a few lattice spacings in real space (i.e at the UV cutoff scale in momentum space). These parameters are often inaccessible to every day experiments, since the latter probe only the low-energy and the long-wavelength dynamics of a given system. In RG theory, these unobserved quantities are known as the bare parameters. In fact, if one attempts to construct naive perturbative calculations with such parameters, one will obtain IR divergent Feynman diagrams in the computation of several quantities such as, e.g. the backscattering and the forward scattering four-point vertices, and the single-particle Green’s function. These divergences mean that the perturbation theory setup is not appropriately formulated. To solve this problem, one can redefine the perturbation scheme in a such a way as to circumvent these infinite results in the calculation of observable quantities. This is the strategy of the so-called renormalized perturbation theory. Thus, one rewrites all unobserved bare parameters in terms of the associated observable renormalized (or physical) quantities. The difference between them is given by a counterterm, whose fundamental role is to cancel, by construction, the related IR divergences in all orders of perturbation theory. If this program is successfully accomplished, then the theory is said to be properly renormalized.
As we explained in our previous paper Hermann , when one deals with this FS problem, the counterterms needed to renormalize the theory turn out to be continuous functions of the three momenta parallel to the FS, rather than being simply infinite constants. In addition, we computed the RG flow equations for the coupling functions and the quasiparticle weight up to two-loop order, and showed that they were in fact coupled integro-differential equations. We solved those equations self-consistently and, as a result, we found out for an intermediate coupling regime a possibly new physical regime, which was characterized by a strongly suppressed quasiparticle weight. Here we continue to explore this interacting regime, and our present calculation shows that the resulting quantum state may in fact be interpreted as an ISL.
To obtain the uniform susceptibilities of this system, we must first calculate the linear response due to an infinitesimal external field perturbation that couples to both charge and spin number operators. We do this by adding to the Lagrangian the new term
$$h_{external}\underset{𝐩,a=\pm }{}𝒯_B^{\alpha \alpha }(𝐩)\psi _{(a)\alpha }^B\left(𝐩\right)\psi _{(a)\alpha }^B\left(𝐩\right),$$
(3)
where $`B`$ stands for the bare quantities. This will generate an additional vertex (the one-particle irreducible function $`\mathrm{\Gamma }^{(2,1)}(𝐩,𝐪0)`$), which will in turn be afflicted by new IR divergences (see the nonparquet diagrams in Fig. 2). As a result, we must rewrite the bare quantity $`𝒯_B^{\alpha \alpha }`$ in terms of its renormalized counterpart (henceforth called $`𝒯_R^{\alpha \alpha }`$), and an appropriate counterterm $`\mathrm{\Delta }𝒯_R^{\alpha \alpha }`$, i.e.
$$𝒯_B^{\alpha \alpha }(p_{})=Z^1(p_{})\left[𝒯_R^{\alpha \alpha }(p_{})+\mathrm{\Delta }𝒯_R^{\alpha \alpha }(p_{})\right].$$
(4)
The quasiparticle weight $`Z`$ factor comes from the renormalization of the fermionic fields, which must be also taken into account in order to include the feedback of the self-energy effects into the RG flow equations. The Z function is calculated explicitly in Ref. Hermann . As was mentioned before, $`\mathrm{\Delta }𝒯_R^{\alpha \alpha }`$ must cancel exactly the divergences generated by the nonparquet diagrams. However, there are still several ways of choosing that counterterm. To solve this ambiguity, we must make a prescription establishing precisely that the $`𝒯_R^{\alpha \alpha }`$ is the experimentally observable response, i.e., $`\mathrm{\Gamma }^{(2,1)}(p_{},p_0=\omega ,p_{}=k_F;\text{q}0)=i𝒯_R^{\alpha \alpha }(p_{},\omega )`$, where $`\omega `$ is the RG energy scale parameter that denotes the proximity of the renormalized theory to the FS. In this way, to flow towards the FS we let $`\omega 0`$.
We are now ready to define the two different types of uniform response functions, which arise from a symmetrization with respect to the spin projection, namely
$`𝒯_{R,CS}(p_{},\omega )`$ $`=`$ $`𝒯_R^{}(p_{},\omega )+𝒯_R^{}(p_{},\omega ),`$ (5)
$`𝒯_{R,SS}(p_{},\omega )`$ $`=`$ $`𝒯_R^{}(p_{},\omega )𝒯_R^{}(p_{},\omega ),`$ (6)
where $`𝒯_{R,CS}`$ and $`𝒯_{R,SS}`$ are the response functions associated with the charge compressibility and spin susceptibility, respectively. In order to compute the RG flow equations for these response functions, one needs to recall that the bare parameters are independent of the renormalization scale $`\omega `$. Thus, using the RG condition $`\omega d𝒯_B^{\alpha \alpha }/d\omega =0`$, we obtain
$`\omega {\displaystyle \frac{d}{d\omega }}𝒯_{R,CS}(p_{})`$ $`=`$ $`\omega {\displaystyle \frac{d}{d\omega }}\mathrm{\Delta }𝒯_{R,CS}(p_{})+\gamma (p_{})𝒯_{R,CS}(p_{}),`$
$`\omega {\displaystyle \frac{d}{d\omega }}𝒯_{R,SS}(p_{})`$ $`=`$ $`\omega {\displaystyle \frac{d}{d\omega }}\mathrm{\Delta }𝒯_{R,SS}(p_{})+\gamma (p_{})𝒯_{R,SS}(p_{}),`$
where the anomalous dimension is given by $`\gamma (p_{})=\omega d\mathrm{ln}Z(p_{})/d\omega `$. Despite their apparent simplicity, it is impossible to solve these RG equations only by analytical means. To find their solutions, we have again to resort to numerics. Here we use the fourth-order Runge-Kutta numerical method. We discretize the FS continuum replacing the interval $`\mathrm{\Delta }p_{}\mathrm{\Delta }`$ by a discrete set of 33 points. We use $`\omega =\mathrm{\Omega }\mathrm{exp}(l)`$, where $`\mathrm{\Omega }=2v_F\lambda `$ with $`l`$ being our RG step. We also choose $`\mathrm{\Omega }=v_F\mathrm{\Delta }<W`$. In view of our choice of points for the FS, we are only allowed to go up to $`l2.8`$ in the RG flow to avoid the distance to the FS being smaller than the shortest distance between points in our discrete set.
Once the response functions are obtained, we can calculate the flow of the charge compressibility and the uniform spin susceptibility of the system. Using again our diagrammatic convention, it follows from Fig. 3 that they are given by
$$\chi _{CS(SS)}^R(\omega )=\frac{1}{4\pi ^2v_F}_\mathrm{\Delta }^\mathrm{\Delta }𝑑p_{}\left[𝒯_{R,CS(SS)}(p_{},\omega )\right]^2.$$
(9)
To evaluate these quantities we follow the same numerical procedure described above. Our results are displayed in Fig. 4. We note in this plot that both charge and spin susceptibilities flow at the same rate as we approach the FS even though their corresponding response functions have different flow equations. In addition, for initial couplings in which the quasiparticle weight flows to zero, the uniform susceptibilities become strongly suppressed in the low-energy limit. While the former asserts that there are no fermionic quasiparticle excitations present in the system, the latter is related to the complete absence of low-lying bosonic charge and spin collective excitations. Since this resulting state has only gapful excitations, it cannot be related to any broken symmetry state and, as a consequence, should possess only short-range ordering. This quantum state has, therefore, strong similarities with that predicted long ago by Anderson Anderson , which is commonly referred to as an ISL. In our present work, it becomes evident that such a state is produced by disordering effects induced by strong quantum fluctuations, and these are approximately taken into account by our two-loop RG scheme. Finally, we call attention to the fact that an insulating behavior in the lightly doped 2D HM was also reported recently in the literature Phillips . Our present result is clearly in agreement with their results.
In summary, we examined the flow of both charge compressibility and uniform spin susceptibility in the lightly doped 2D HM as a function of the initial interaction strength within a two-loop RG approach. For moderate interaction regimes, both quantities flow to zero as we approach the initial FS of the system. This is a strong indicative that there are gaps in both charge and spin excitation spectra of the lightly doped 2D HM. Hence the quantum state associated with that regime may be viewed as an ISL as opposed to a Mott insulator. This result may be of relevance for the cuprate high-Tc superconductors in view of the fact that the 2D HM in the intermediate coupling regime is widely believed to be appropriate to describe such compounds in all doping ranges.
This work was supported by the Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq).
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# Role of Pressure in Quasi-Spherical Gravitational Collapse
## I Introduction
In cosmology, solutions to Einstein’s field equations are obtained by imposing symmetries on space-time. Usually, spatial homogeneity is one of the reasonable assumptions (in an average sense) while for cosmological phenomena over galactic scale or in smaller scale, inhomogeneous solutions are useful. Szekeres in 1975 gave a class of inhomogeneous solutions representing irrotational dust. The space-time represented by these solutions has no killing vectors and it has invariant family of spherical hypersurfaces. Hence this space-time is referred as quasi-spherical space-time.
An extensive study of gravitational collapse \[3-8\] has been carried out of Tolman-Bondi-Lemaître (TBL) spherically symmetric space-times containing irrotational dust to support or disprove the cosmic censorship conjecture (CCC). A general conclusion from these studies is that a central curvature singularity forms but its local or global visibility depends on the initial data. Also over the past decades, the role of pressure within a collapsing cloud has been studied \[9-13\] but the actual role of pressure in determining the end state of a continual collapse is not yet clear.
On the other hand, there is very little progress in studying non-spherical collapse. Basically, the difficulty is the ambiguity of horizon formation in non-spherical geometries and the influence of gravitational radiation. Though there is hoop conjecture by Thorne to characterize the formation of horizon but only few works \[15-19\] have been done to confirm or refute the conjecture. Recently, an extensive study of irrotational dust collapse has been done in quasi-spherical Szekeres space-time both for four and higher dimensions. In this paper, we have done an extensive analysis of Szekeres model with matter containing pressure and studied the collapsing procedure to examine the role of pressure in characterizing the final singularity. The paper is organized as follows: The Szekeres model has been described in section II. The perfect fluid solution with asymptotic behaviour has been presented in section III. Section IV deals with the tangential stress only while section V deals with gravitational collapse with non-isotropic pressure. Finally, the paper ends with a brief discussion in section VI.
## II The Szekeres’ Model
The metric ansatz for the four dimensional Szekeres’ space-time is of the form
$$ds^2=dt^2e^{2\alpha }dr^2e^{2\beta }(dx^2+dy^2)$$
(1)
where $`\alpha `$ and $`\beta `$ are functions of all space-time variables i.e.,
$$\alpha =\alpha (t,r,x,y),\beta =\beta (t,r,x,y).$$
Considering both radial and transverse stresses the energy momentum tensor has the following structure
$$T_\mu ^\nu =\text{diag}(\rho ,p_r,p__T,p__T)$$
(2)
Hence the full set of Einstein equations are
$$2\dot{\alpha }\dot{\beta }+\dot{\beta }^2e^{2\beta }(\alpha _x^2+\alpha _y^2+\alpha _{xx}+\alpha _{yy}+\beta _{xx}+\beta _{yy})+e^{2\alpha }(2\alpha ^{}\beta ^{}3\beta ^22\beta ^{\prime \prime })=\rho $$
(3)
$$3\dot{\beta }^2+2\ddot{\beta }e^{2\alpha }\beta ^2e^{2\beta }(\beta _{xx}+\beta _{yy})=p_r$$
(4)
$$\ddot{\alpha }+\dot{\alpha }^2+\ddot{\beta }+\dot{\beta }^2+\dot{\alpha }\dot{\beta }+e^{2\alpha }(\alpha ^{}\beta ^{}\beta ^2\beta ^{\prime \prime })e^{2\beta }(\alpha _y^2+\alpha _{yy}\alpha _y\beta _y+\alpha _x\beta _x)=p__T$$
(5)
$$\ddot{\alpha }+\dot{\alpha }^2+\ddot{\beta }+\dot{\beta }^2+\dot{\alpha }\dot{\beta }+e^{2\alpha }(\alpha ^{}\beta ^{}\beta ^2\beta ^{\prime \prime })e^{2\beta }(\alpha _x^2+\alpha _{xx}\alpha _x\beta _x+\alpha _y\beta _y)=p__T$$
(6)
$$\alpha _y(\alpha _x+\beta _x)+\alpha _x\beta _y\alpha _{xy}=0$$
(7)
$$\dot{\alpha }\beta ^{}\dot{\beta }\beta ^{}\dot{\beta ^{}}=0$$
(8)
$$\dot{\alpha }\alpha _x+\dot{\beta }\alpha _x\dot{\alpha _x}\dot{\beta _x}=0$$
(9)
$$\dot{\alpha }\alpha _y+\dot{\beta }\alpha _y\dot{\alpha _y}\dot{\beta _y}=0$$
(10)
$$\alpha _x\beta ^{}\beta _x^{}=0$$
(11)
$$\alpha _y\beta ^{}\beta _y^{}=0$$
(12)
where dot, dash and subscript stands for partial differentiation with respect to $`t`$, $`r`$ and the corresponding variables respectively (e.g., $`\beta _x=\frac{\beta }{x}`$).
Combining the time derivatives of equations (11) and (12) with the differentiation of (8) with respect to $`x`$ and $`y`$ separately we have the integrability condition
$$\beta ^{}\dot{\beta }_x=0=\beta ^{}\dot{\beta }_y$$
(13)
Hence we may have three possibilities namely,
$$\begin{array}{cc}(a)\beta ^{}0,\dot{\beta _x}=0=\dot{\beta _y},& \\ & \\ (b)\beta ^{}=0,\dot{\beta _x}=0=\dot{\beta _y}& \\ & \\ (c)\beta ^{}=0,\dot{\beta _x}0,\dot{\beta _y}0& \end{array}$$
(14)
Following Szekeres’ the field equations are not solvable for the third case so we shall consider only the first two cases.
The energy conservation equation namely $`T^{\mu \nu };_\nu =0`$ gives
$$\dot{\rho }+\rho (\dot{\alpha }+2\dot{\beta })+(\dot{\alpha }p_r+2\dot{\beta }p__T)=0$$
(15)
$$\frac{}{x}(p__Te^\alpha )=p_r\alpha _xe^\alpha $$
(16)
$$\frac{}{y}(p__Te^\alpha )=p_r\alpha _ye^\alpha $$
(17)
$$\frac{}{r}(p_re^{2\beta })=p__T\frac{}{r}(e^{2\beta })$$
(18)
Now for the first choice namely $`\beta ^{}0,\dot{\beta }_x=0=\dot{\beta }_y`$ we have from the field equations the explicit form of the metric coefficients are as follows:
$$e^\beta =R(t,r)e^{\nu (r,x,y)}$$
(19)
$$e^\alpha =R^{}+R\nu ^{}$$
(20)
where $`R`$ and $`\nu `$ satisfy the following differential equations
$$2R\ddot{R}+\dot{R}^2+p_rR^2=f(r),(f(r)=\text{arbitrary separation function})$$
(21)
and
$$e^{2\nu }(\nu _{xx}+\nu _{yy})=f(r)1$$
(22)
Here we have assumed $`p_r=p_r(r,t)`$. Equation (21) has the first integral
$$\dot{R}^2=f(r)+\frac{G(r)}{R}\frac{1}{R}p_rR^2𝑑R$$
(23)
while one of the possible solutions of equation (22) can be taken as
$$e^\nu =A(r)(x^2+y^2)+B_1(r)x+B_2(r)y+D(r)$$
(24)
where the arbitrary functions $`A(r),B_1(r),B_2(r)`$ and $`D(r)`$ are related as
$$B_1^2+B_2^24AD=f(r)1$$
and $`G(r)`$ is an arbitrary function.
For the choice ($`b`$) the metric coefficients are of the form
$$e^\beta =R(t)e^{\nu (x,y)}$$
(25)
$$e^\alpha =R(t)\eta (r,x,y)+\mu (t,r)$$
(26)
Then as before from the field equation (4) we have similar differential equation in $`R`$ and $`\nu `$ as
$$2R\ddot{R}+\dot{R}^2+p_rR^2=K,(K\text{is a constant})$$
(27)
$$e^{2\nu }(\nu _{xx}+\nu _{yy})=K$$
(28)
with $`K`$ an arbitrary constant. Also for $`\nu `$ we choose as in case (a)
$$e^\nu =P(x^2+y^2)+Q_1x+Q_2y+S$$
(29)
where $`P,Q_1,Q_2`$ and $`S`$ are arbitrary constants restricted by the relation
$$Q_1^2+Q_2^24PS=K$$
Now to determine the function $`\eta `$ we have from the field equation (7)
$$\frac{^2(e^\nu \eta )}{xy}=0$$
(30)
and then from the field equations (5) and (6) we have a possible solution
$$e^\nu \eta =u(r)(x^2+y^2)+v_1(r)x+v_2(r)y+w(r)$$
(31)
where $`u(r),v_1(r),v_2(r)`$ and $`w(r)`$ are arbitrary functions of $`r`$ alone.
Also the differential equation in $`\mu `$ is
$$R\ddot{\mu }+\dot{R}\dot{\mu }+\mu (\ddot{R}+p__TR)+(p__Tp_r)R^2\eta =h(r)$$
(32)
with
$$h(r)=2(uS+wP)(v_1Q_1+v_2Q_2).$$
Now for explicit solutions, we shall consider the following cases in the next sections:
(i) the perfect fluid model (i.e., $`p_r=p__T`$)
(ii) the tangential stress only (i.e., $`p_r=0,p__T0`$)
(iii) the general case (i.e., $`p_r0,p__T0`$).
In fact, cosmological solutions are obtained (in sections III and IV) for the first two cases respectively while for the third case collapsing behaviour has been studied and the role of pressure has been examined.
## III The perfect fluid model
In this case due to energy conservation equations (16)-(18) the isotropic pressure is function of $`t`$ only i.e., $`p=p(t)`$ ($`p_r=p__T=p`$). As there is no restriction on the energy density so $`\rho `$ is in general a function of all the 4 variables i.e., $`\rho =\rho (t,r,x,y)`$ and hence no equation of state is imposed.
Now for explicit solution according to Szafron and Szafron and Wainwright we choose
$$p(t)=p_ct^s$$
(33)
($`p_c`$ and $`s`$ are positive constants) and we have the general solution for $`R`$ as
$$R^{\frac{3}{2}}=\{\begin{array}{ccc}\sqrt{t}\left\{C_1J_\xi [\frac{2\sqrt{\lambda }}{|s2|}t^{\frac{s2}{2}}]+C_2Y_\xi [\frac{2\sqrt{\lambda }}{|s2|}t^{\frac{s2}{2}}]\right\}\hfill & & \\ & & \\ \sqrt{t}\left\{C_1J_\xi [\frac{2\sqrt{\lambda }}{|s2|}t^{\frac{s2}{2}}]+C_2J_\xi [\frac{2\sqrt{\lambda }}{|s2|}t^{\frac{s2}{2}}]\right\}\hfill & & \\ & & \\ C_1t^{q_1}+C_2t^{1q_1}\hfill & & \end{array}$$
(34)
according as $`\xi `$ is an integer, non-integer and $`s=2`$. Here $`C_1`$ and $`C_2`$ are arbitrary functions of $`r`$ and we have chosen
$$\xi =\frac{1}{s2},\lambda =\frac{3p_c}{4},q_1=\frac{1}{2}(1+\sqrt{13p_c}).$$
It is to be noted that to derive the above solution we have chosen $`f(r)=0`$. Further, if we consider the dust model (i.e., $`p=0`$) then the above solution simplifies to $`Rt^{2/3}`$ which is the form of the scale factor in the usual Friedman model.
Now, the physical and kinematical parameters have the following expressions
$$\rho =\frac{G^{}(r)+3G\nu ^{}}{R^2(R^{}+R\nu ^{})}\frac{p_c}{t^s}$$
(35)
$$\theta =\frac{R\dot{R}^{}+3R\dot{R}\nu ^{}+2\dot{R}R^{}}{R(R^{}+R\nu ^{})}$$
(36)
$$\sigma ^2=\frac{1}{12}\left[\frac{R(Rf^{}2R^{}f)+(RG^{}(r)3R^{}G)}{\dot{R}R^2(R^{}+R\nu ^{})}\right]^2$$
(37)
The above solution is for the choice ($`a`$) (see eq.(13)). For the choice ($`b`$) (i.e., eq.(14)) the explicit form for $`R`$ is same as in equation (34) except here $`C_1`$ and $`C_2`$ are arbitrary constants and $`K=0`$. But we note that the differential equation (30) is not solvable for the above explicit solutions for $`R`$. The physical and kinematical parameters have the expressions as
$$\rho =\frac{2(\dot{R}^2K)}{R^2}\frac{2(\ddot{\mu }+\eta \ddot{R})}{\mu +\eta R}\frac{p_c}{t^s}$$
(38)
$$\theta =\frac{R\dot{\mu }+2\mu \dot{R}+3R\dot{R}\eta }{R(\mu +\eta R)}$$
(39)
$$\sigma ^2=\frac{1}{3}\left[\frac{R\dot{\mu }\dot{R}\mu }{R(\mu +\eta R)}\right]^2$$
(40)
$$\mathrm{𝐀𝐬𝐲𝐦𝐩𝐭𝐨𝐭𝐢𝐜}\mathrm{𝐁𝐞𝐡𝐚𝐯𝐢𝐨𝐮𝐫}$$
We shall now discuss the asymptotic behaviour of the above solutions. The co-ordinates vary over the range : $`t_0<t<\mathrm{};\mathrm{}<r<\mathrm{};\mathrm{}<x,y<\mathrm{}`$. For both the choices ($`a`$) and ($`b`$) as $`p0,p0`$ so we must have $`\frac{1}{2}<q_1<1`$. So for $`s=2`$ for large $`t`$,
$$\begin{array}{ccc}Rt^{\frac{2q_1}{3}}\hfill & & \\ & & \\ \rho t^2\hfill & & \\ & & \\ pt^2\hfill & & \\ & & \\ \theta t^1\hfill & & \\ & & \\ \sigma ^2t^2\hfill & & \end{array}$$
(41)
Hence as $`t\mathrm{}`$, ($`p,\rho `$) fall off faster compare to ($`\theta ,\sigma `$), while the scale factor $`R`$ gradually increases with time. So the model approaches isotropy along fluid world line as $`t\mathrm{}`$.
## IV Model with tangential stresses only
For this model we have from the conservation equation (18) $`\beta ^{}=0`$ i.e. choice (b) is only possible here. Also from the other conservation equations namely equation (16) and (17) we have
$$p__T=A(r,t)e^\alpha $$
(42)
where A(r,t) is an arbitrary function of $`r`$ and $`t`$. But for the consistency of the differential eq.(30) $`p__T`$ (as stated earlier) must be a function of $`r`$ and $`t`$ and hence $`\alpha `$ should be independent of $`x`$ and $`y`$. As a consequence, in the solution (26) for $`\alpha `$, $`\eta `$ must be independent of $`x`$ and $`y`$. Thus for the solution (32) for $`\eta `$ we should have
$$u(r)=P\eta _0(r),v_1(r)=Q_1\eta _0(r),v_2(r)=Q_2\eta _0(r)\text{and}w(r)=S\eta _0(r)$$
Hence we have $`\eta =\eta _0(r)`$, an arbitrary function of r alone and $`h(r)=K\eta _0(r)`$.
Now, the differential equation for $`R`$ has the simple form
$$\dot{R}^2=a_1+\frac{a_2}{R},(a_1,a_2\text{are constants})$$
(43)
which has a parametric solution of the form
$`\begin{array}{c}R=\frac{a_2}{2(a_1)}(1\text{cos}\varphi )\\ \\ t_ct=\frac{a_2}{2(a_1)^{3/2}}(\varphi \text{sin}\varphi )\end{array}\}\text{for}a_1<0(0<\varphi <2\pi )`$ (48)
(49)
$`\begin{array}{c}R=\frac{a_2}{2a_1}(\text{cosh}\varphi 1)\\ \\ t_ct=\frac{a_2}{2a_1^{3/2}}(\text{sinh}\varphi \varphi )\end{array}\}\text{for}a_1>0(\varphi >0)`$ (53)
$`\begin{array}{c}\text{and}\\ \\ R=\left(\frac{9a_2}{4}\right)^{1/3}(t_ct)^{2/3}\text{for}a_1=0\end{array}`$ (57)
Here $`t_c`$ is an integration constant that corresponds to the time of arrival of each shell to the central singularity.
Choosing the power law solution (i.e. $`a_1=0`$) for $`R`$ (i.e., if $`RT^{2/3}`$ then equation (27) implies that $`K=0`$) and assuming $`p__T=p_{_{0T}}/T^2,(T=t_ct)`$ (i.e., a function of $`T`$ alone), it is possible to have a solution for $`\mu `$ (from eq.(32)) as
$$\mu (r,t)=\{\begin{array}{c}C_1(r)T^{n_1}+C_2(r)T^{n_2},p_{_{0T}}<1/4\\ \\ C_1(r)T^{1/6}\text{cos}(k\text{ln}T)+C_2(r)T^{1/6}\text{sin}(k\text{ln}T),p_{_{0T}}>1/4\\ \\ C_1(r)T^{1/6}+C_2(r)T^{1/6}\text{ln}T,p_{_{0T}}=1/4\end{array}$$
(58)
where $`C_1`$ and $`C_2`$ are arbitrary functions of $`r`$, $`n_1=\frac{1}{6}+\frac{1}{2}\sqrt{14p_{_{0T}}}`$, $`n_2=\frac{1}{6}\frac{1}{2}\sqrt{14p_{_{0T}}}`$, $`k=\frac{1}{2}\sqrt{4p_{_{0T}}1}`$ and $`R_0=\left(\frac{9a_2}{4}\right)^{1/3}`$ (Here we have set $`\eta =0`$, which has no effect on the metric).
Further, the physical and kinematical parameters have the expressions
$$\rho =\frac{2(\dot{R}^2K)}{R^2}\frac{2(\ddot{\mu }+\eta \ddot{R})}{\mu +\eta R}2p__T$$
(59)
$$\theta =\frac{R\dot{\mu }+2\mu \dot{R}+3\eta R\dot{R}}{R(\mu +\eta R)}$$
(60)
$$\sigma ^2=\frac{1}{3}\left[\frac{\dot{\mu }R\mu \dot{R}}{R(\mu +\eta R)}\right]^2$$
(61)
It is to be noted that the solution for R does not depend on $`p__T`$ so R has same expression for dust model. But for the solution of $`\mu `$ we have only the power law form $`T^{2/3}`$ (or $`T^{1/3}`$) when $`R`$ has Friedmann like behaviour (i.e. $`RT^{2/3}`$). The difference comes in the matter density. For dust model $`\rho `$ is a function of all the four co-ordinate variables while in the presence of tangential stress $`\rho `$ is a function of $`t`$ and $`r`$ only. Finally, the asymptotic behaviour for both the model will be very similar.
## V Role of pressure in gravitational collapse
In the general case when both radial and tangential pressures are non-zero and distinct then from the Einstein equations they can be obtained in compact form as
$$\begin{array}{c}\rho =\frac{F^{}}{\zeta ^2\zeta ^{}}\\ \\ p_r=\frac{\dot{F}}{\zeta ^2\dot{\zeta }}\\ \\ p__T=p_r+\frac{\zeta p_r^{}}{2\zeta ^{}}\end{array}$$
(62)
where $`F(r,t)=Re^{3\nu }(\dot{R}^2f(r))`$ and $`\zeta =e^\beta `$.
Since $`p_r`$ is regular initially at the centre and blows up at the singularity so we can choose $`p_r`$ to be of the form:
$$p_r=\frac{g(r)}{R^n}$$
(63)
where $`g(r)`$ is an arbitrary function such that $`g(r)r^n`$ near $`r=0`$ to make initial matter density non-zero at the centre $`r=0`$ and $`n`$ is any constant. As a consequence, the expressions for matter density and tangential stress become
$$\rho =\frac{H^{}+3H\nu ^{}}{R^2(R^{}+R\nu ^{})}$$
(64)
$$p__T=\frac{g(r)}{R^n}\left[1\frac{nR^{}}{2(R^{}+R\nu ^{})}\right]+\frac{g^{}(r)}{2R^{n1}(R^{}+R\nu ^{})},(n3)$$
(65)
where $`H(R,t)=C(r)\frac{g(r)}{(3n)}R^{3n}`$ and $`C(r)`$ is an arbitrary integration function. Also the radial velocity of collapsing shells at a distance $`r`$ from the centre is given by
$$\dot{R}^2=f(r)+\frac{H(R,t)}{R}$$
(66)
Now if we choose $`R=r`$ initially then at the beginning of the collapse the density and the tangential stress have the initial values
$$\rho _i(r,x,y)=\rho _i(r,t_i,x,y)=\frac{c^{}+3c\nu ^{}}{r^2(1+r\nu ^{})}$$
(67)
$$p_{_{T_i}}=p__T(t=t_i)=\frac{g(r)}{r^n}\left[1\frac{n}{2(1+r\nu ^{})}\right]+\frac{g^{}(r)}{2r^{n1}(1+r\nu ^{})}$$
(68)
where $`c(r)=H(r,t_i)=C(r)\frac{g(r)}{3n}r^{3n}`$.
Here it is to be noted that for regular initial data $`C(r)`$ and $`g(r)`$ to be $`C^{\mathrm{}}`$ functions and hence we have the following series expansions
$$\begin{array}{c}g(r)=_{j=0}^{\mathrm{}}g_jr^{n+j}\\ \\ C(r)=_{j=0}^{\mathrm{}}C_jr^{3+j}\\ \\ \rho _i(r)=_{j=0}^{\mathrm{}}\rho _jr^j\\ \\ \nu ^{}=_{j=1}^{\mathrm{}}\nu _jr^j\\ \\ p_{_{T_i}}=_{j=0}^{\mathrm{}}p_jr^j\end{array}$$
(69)
where $`\nu __11`$.
In these series expansions the coefficients $`g_j`$’s and $`C_j`$’s are constants while $`\rho _j`$’s, $`\nu _j`$’s and $`p_j`$’s are functions of $`x`$ and $`y`$. These coefficients are related among themselves through the relations (54) and (55) as follows:
$$\begin{array}{c}p_0=g_0,p_1=g_1\left\{1+\frac{1}{2(1+\nu __1)}\right\},p_2=g_2\left\{1+\frac{1}{(1+\nu __1)}\right\}\frac{g_1\nu _0}{2(1+\nu __1)^2},\mathrm{}.\mathrm{}.\\ \\ \rho _0=3c_0,\rho _1=\frac{4+3\nu __1}{1+\nu __1}c_1,\rho _2=\frac{5+3\nu __1}{1+\nu __1}c_2\frac{\nu _0c_1}{(1+\nu __1)^2},\mathrm{}.\mathrm{}.\end{array}$$
(70)
or
$$\begin{array}{c}p_0=g_0+\frac{g_1}{2\nu _0},p_1=g_1\left(1\frac{\nu _1}{2\nu _0^2}\right)+\frac{g_2}{\nu _0},p_2=g_2\left(1\frac{\nu _1}{\nu _0^2}\right)+\frac{(\nu _1^2\nu _0\nu _2)}{2\nu _0^3}g_1+\frac{3g_3}{2\nu _0},\mathrm{}.\mathrm{}.\\ \\ \rho _0=3c_0+\frac{c_1}{\nu _0},\rho _1=\frac{2c_2}{\nu _0}+c_1\left(3\frac{\nu _1}{\nu _0^2}\right),\rho _2=\frac{3c_3}{\nu _0}+c_2\left(3\frac{2\nu _1}{\nu _0^2}\right)+c_1\frac{(\nu _1^2\nu _0\nu _2)}{\nu _0^3},\mathrm{}.\mathrm{}.\end{array}$$
(71)
according as $`\nu __1>1`$ or $`\nu __1=1`$ and $`c_i=C_i\frac{g_i}{3n},i=0,1,2,\mathrm{}`$.
The hypersurface $`t=t_s(r)`$ describing the shell focusing singularity is characterized by
$$R(t_s(r),r)=0$$
(72)
As the differential equation in $`R`$ (i.e., eq.(53)) is not solvable so we shall consider only the marginally bound case (i.e., $`f(r)=0`$). Hence in this case, the singularity hypersurface can be written in explicit form as
$$t_s(r)t_i=\frac{2r^{3/2}}{3\sqrt{C(r)}}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},b,b+1,\frac{g(r)r^{3n}}{C(r)(3n)}]$$
(73)
where $`{}_{2}{}^{}F_{1}^{}`$ is the usual hypergeometric function with $`b=\frac{3}{2(3n)}`$ .
### V.1 Formation of Trapped Surfaces
The event horizon of observers at infinity plays an important role in the nature of the singularity. As formation of event horizon depends greatly on the computation of null geodesics whose computation are almost impracticable for the present space-time geometry, so we consider closely related concept of a trapped surface (namely a compact space-like 2-surface whose normals on both sides are future pointing converging null geodesic families). In fact, if the 2-surface $`S_{r,t}`$ ($`r=`$constant, $`t=`$constant) is a trapped surface then it and its entire future development lie behind the event horizon provided the density falls off fast enough at infinity. So mathematically, if $`K^\mu `$ denotes the tangent vector field to the null geodesics which is normal to $`S_{r,t}`$ then we have
$$K_\mu K^\mu =0,K_{;\nu }^\mu K^\nu =0.$$
Now the convergence or divergence of the null geodesics is characterized by the sign of the invariant $`K_{;\mu }^\mu `$ evaluated on the surface $`S_{r,t}=0`$ (in fact, $`K_{;\mu }^\mu <0`$ indicates convergence while $`K_{;\mu }^\mu >0`$ stands for divergence). It can be shown that the inward geodesics converges initially and throughout the collapsing process while the outward geodesics diverges initially but becomes convergent after a time $`t_{ah}(r)`$ (time of formation of apparent horizon) given by
$$\dot{R}(t_{ah}(r),r)=\sqrt{1+f(r)}$$
Now using equations (23) and (50) we have
$$g(r)R^{3n}(t_{ah}(r),r)(n3)R(t_{ah}(r),r)+(n3)C(r)=0$$
(74)
or using the explicit solution for $`R`$ (from eq.(53))
$$t_{ah}(r)t_i=\frac{2r^{3/2}}{3\sqrt{C(r)}}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},b,b+1,\frac{g(r)r^{3n}}{C(r)(3n)}]\frac{2R^{3/2}(t_{ah},r)}{3\sqrt{C(r)}}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},b,b+1,\frac{g(r)R^{3n}(t_{ah},r)}{C(r)(3n)}]$$
(75)
From the equations (60) and (62) we see that the shell focusing singularity that appears at $`r>0`$ is in the future of the apparent horizon.
As we are interested in central shell focusing singularity (at $`r=0`$), so its time of occurrence is given by
$$\begin{array}{c}t_0=t_s(r)\\ limr0\\ \\ =t_i+\frac{2}{3\sqrt{C_0}}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},b,b+1,z],\left(z=\frac{g_0}{C_0(3n)}\right)\end{array}$$
(76)
where in evaluating the limit we have used the series form of $`g(r)`$ and $`C(r)`$ (from eq. (56)). Now if we restrict $`n<3`$ then we have a comparative expression between $`t_{ah}(r)`$ and $`t_0`$ as
$`t_{ah}(r)t_0=\left[C_0^{3/2}C_1{}_{2}{}^{}F_{1}^{}[{\displaystyle \frac{1}{2}},b,b+1,z]+{\displaystyle \frac{(C_0g_1C_1g_0)}{(3n)(92n)}}C_0^{5/2}{}_{2}{}^{}F_{1}^{}[{\displaystyle \frac{3}{2}},b+1,b+2,z]\right]r`$
$$+O(r^2)\frac{C_0^{3n}g_0}{(3n)(92n)}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},b,b+1,C_0^{3n}z]r^{92n}+\mathrm{}.\mathrm{}.,(n<3)$$
(77)
Note that here $`t_0`$ is the time of formation of singularity at $`r=0`$ while $`t_{ah}(r)`$ is the epoch at which a trapped surface is formed at a distance $`r`$. Thus if trapped surface is formed at a later instant than $`t_0`$ then it is possible that any light signal from the singularity can reach an observer. As the geometry of the present model is a class of spherical space-time having different centres, so the condition for formation of NS (or BH) will be same as TBL model. Therefore, $`t_{ah}(r)>t_0`$ is the necessary condition for formation of naked singularity, while to form black hole the sufficient condition is $`t_{ah}(r)t_0`$. It should be mentioned that this criterion for naked singularity is purely local.
Due to complicated form of equation (64) it is very difficult to make a comparative study between $`t_{ah}`$ and $`t_0`$. Hence for simplicity we choose $`n=3/2`$. Then the difference between $`t_{ah}`$ and $`t_0`$ has the form
$$t_{ah}(r)t_0=\frac{2\left(C_0g_1C_1g_0g_1\sqrt{C_0}\sqrt{C_0\frac{2}{3}g_0}\right)}{3g_0\sqrt{C_0}\sqrt{C_0\frac{2}{3}g_0}\left(\sqrt{C_0}+\sqrt{C_0\frac{2}{3}g_0}\right)}r+O(r^2)$$
(78)
Hence in the present problem it is possible to have local naked singularity or a black hole form under the conditions shown in the following table (see Table I):
TABLE-I
| Choice of the parameters | Naked Singularity | Black hole |
| --- | --- | --- |
| (i) $`g_1>0,C_1<0`$ | Always possible | Not possible |
| (ii) $`g_1<0,C_1>0`$ | Not possible | Always possible |
| (iii) $`g_1>0,C_1>0`$ | $`\frac{g_1}{C_1}>\frac{3}{2}\left(1+\sqrt{1\frac{2}{3}k_0}\right)`$ | $`\frac{g_1}{C_1}<\frac{3}{2}\left(1+\sqrt{1\frac{2}{3}k_0}\right)`$ |
| (iV) $`g_1<0,C_1<0`$ | $`|\frac{g_1}{C_1}|<\frac{3}{2}\left(1+\sqrt{1\frac{2}{3}k_0}\right)`$ | $`|\frac{g_1}{C_1}|>\frac{3}{2}\left(1+\sqrt{1\frac{2}{3}k_0}\right)`$ |
Here we note that for initial density gradient to be negative at the centre (i.e., $`\rho _1<0`$) we must have $`(C_1\frac{2g_1}{3})<0`$ (for $`\nu __1>1`$). In the first case (i.e., $`g_1>0,C_1<0`$) we have negative definite $`\rho _1`$ and there is always naked singularity as in the dust model. Similarly in the second case (i.e., $`g_1<0,C_1>0`$), $`\rho _1`$ is positive definite and we always get black hole same as dust model. For the third and fourth cases (when $`g_1`$ and $`C_1`$ are of same sign) both naked singularity (NS) and black hole (BH) are possible for the restrictions given in the table I. When both $`g_1`$ and $`C_1`$ are positive (third case) or negative (fourth case) then for formation of NS $`\rho _1`$ is negative but for BH case as there is no lower limit (or upper limit) of $`\frac{g_1}{C_1}`$ (or $`|\frac{g_1}{C_1}|`$) so $`\rho _1>0`$ or $`\rho _1<0`$ are possible. Further for $`g_1=C_1=0`$ we have $`\rho _1=0`$ then we have similar behaviour for the parameters ($`g_2,C_2`$). A diagrammatic representation of $`t_{ah}t_0`$ for variation of $`k_0(=g_0/C_0)`$ and $`k_1(=g_1/C_1)`$ has been shown in figures 1 and 2 for positive and negative $`C_1`$ respectively. In both the figures the vertical positive region corresponds to NS while the negative region stands for BH solution. Finally, we see that if the initial density or pressure has opposite behaviour (i.e., one increases and other decreases and vise-versa) near the centre $`r=0`$ then we have similar character of the singularity as in dust model i.e., pressure has no significant effect on the singularity formation. On the other hand, if initial density and pressure increase or decrease simultaneously near the centre then even for negative density gradient at the centre it is possible to have a BH formation at $`r=0`$, which is a distinct result in compare to dust model. Therefore, we may conclude that pressure tries to resist formation of NS.
### V.2 Study of Geodesics
For simplicity of calculation here we shall consider as before the marginally bound case ($`f(r)=0`$) with $`n=3/2`$. Then $`R(t,r)`$ has the explicit solution (choosing the initial time $`t_i=0`$) which can be written in convenient form as
$$t(r)=\frac{2}{g(r)}\left[\sqrt{C(r)\frac{2}{3}g(r)R^{3/2}}\sqrt{C(r)\frac{2}{3}g(r)r^{3/2}}\right]$$
(79)
To examine whether the singularity at $`t=t_0,r=0`$ is naked or not, we investigate whether there exist one or more outgoing null geodesics which terminate in the past at the central singularity. In particular, we shall concentrate to radial null geodesics only.
First we assume that it is possible to have one or more such geodesics and we choose the equation of the outgoing radial null geodesic (ORNG) which passes through the central singularity in the past as (near $`r=0`$)
$$t_{ORNG}=t_0+ar^\xi $$
(80)
to leading order in the ($`t,r`$) plane with $`a>0,\xi >0`$.
Now the expression for the singularity time (characterized by $`R(t_s(r),r)=0`$) from (66) is
$$t_s(r)=\frac{2}{g(r)}\left[\sqrt{C(r)}\sqrt{C(r)\frac{2}{3}g(r)r^{3/2}}\right]$$
(81)
and hence the time for central singularity is
$$t_0=\frac{2}{g_0}\left(\sqrt{C_0}\sqrt{C_0\frac{2}{3}g_0}\right)$$
(82)
Here we choose for $`C(r)`$ and $`g(r)`$ as
$$\begin{array}{c}C(r)=C_0r^3+C_kr^{k+3}\\ \\ g(r)=g_0r^{3/2}+g__lr^{l+3/2}\end{array}$$
(83)
where $`C_0,g_0`$ are constants and $`C_k(<0)`$ and $`g__l(<0)`$ are the first non-vanishing term beyond $`C_0`$ and $`g_0`$ respectively. As a consequence the expression for $`t_s(r)`$ becomes
$$t_s(r)=t_0+\frac{C_k}{g_0}\left(\frac{1}{\sqrt{C_0}}\frac{1}{\sqrt{C_0\frac{2}{3}g_0}}\right)r^k+\frac{2g__l}{g_0}\left(\frac{1}{3\sqrt{C_0\frac{2}{3}g_0}}\frac{\sqrt{C_0}\sqrt{C_0\frac{2}{3}g_0}}{g_0}\right)r^l+\mathrm{}.$$
(84)
we shall now study the following possibilities:
$$(i)k<l,(ii)k>l$$
Case I : $`k<l`$
Here for $`t_s(r)`$ we write
$$t_s(r)=t_0\frac{C_k}{g_0}\left(\frac{1}{\sqrt{C_0\frac{2}{3}g_0}}\frac{1}{\sqrt{C_0}}\right)r^k,(C_k<0)$$
(85)
Now comparing with the geodesic equation (67) we get the relations
$$(a)\xi >k\text{or}(b)\xi =k\text{and}a<\frac{C_k}{g_0}\left(\frac{1}{\sqrt{C_0\frac{2}{3}g_0}}\frac{1}{\sqrt{C_0}}\right)$$
(86)
When $`\xi >k`$ then near $`r=0`$ the solution for $`R`$ simplifies to
$$R=r\left[1\frac{3}{8}g_0t^2\frac{3}{2}t\left(\sqrt{C_0\frac{2}{3}g_0}+\frac{C_kr^k}{2\sqrt{C_0\frac{2}{3}g_0}}\right)\right]^{2/3}$$
(87)
Further for the given metric an ORNG should satisfy
$$\frac{dt}{dr}=R^{}+R\nu ^{}$$
(88)
To examine the feasibility of the null geodesic starting from the singularity, we combine equations (67) and (74) in equation (75) and we get (upto leading order in $`r`$)
$$a\xi r^{\xi 1}=\left(1+\nu __1+\frac{2k}{3}\right)\left[\frac{3C_kt_0}{4\sqrt{C_0\frac{2}{3}g_0}}\right]^{2/3}r^{\frac{2k}{3}},(\nu __10)$$
(89)
which implies
$$\xi =1+\frac{2k}{3}\text{and}a=\frac{1}{\xi }\left(1+\nu __1+\frac{2k}{3}\right)\left[\frac{3C_kt_0}{4\sqrt{C_0\frac{2}{3}g_0}}\right]^{2/3}$$
(90)
As $`\xi >k`$, so from (77) $`k<3`$. Since $`k`$ is an integer, we could have
$$\begin{array}{c}k=1,\xi =\frac{5}{3}\\ \\ \text{or}\\ \\ k=2,\xi =\frac{7}{3}\end{array}$$
(91)
On the other hand for $`\xi =k`$, as before we get $`k=3`$ and
$`a={\displaystyle \frac{1}{3}}[{\displaystyle \frac{3}{4}}(ag_0t_0+2a\sqrt{C_0{\displaystyle \frac{2}{3}}g_0}+{\displaystyle \frac{C_kt_0}{\sqrt{C_0\frac{2}{3}g_0}}})]^{1/3}\times `$
$$\left[\frac{3}{4}\left\{(1+\nu __1)\left(ag_0t_0+2a\sqrt{C_0\frac{2}{3}g_0}\right)+\frac{\left(2+\nu __1+\frac{4k}{3}\right)C_kt_0}{\sqrt{C_0\frac{2}{3}g_0}}\right\}\right]$$
(92)
Case II: $`k>l`$
In this case
$$t_s(r)=t_0\frac{2g__l}{g_0^2}\left(\sqrt{C_0}+\frac{g_03C_0}{3\sqrt{C_0\frac{2}{3}g_0}}\right)r^l$$
(93)
Now matching the geodesic equation as above we get
$$(a)\xi >l\text{or}(b)\xi =l\text{and}a<\frac{2g__l}{g_0^2}\left(\sqrt{C_0}+\frac{g_03C_0}{3\sqrt{C_0\frac{2}{3}g_0}}\right)$$
(94)
Then for $`\xi >l`$ we get
$$\begin{array}{c}l=1,\xi =\frac{5}{3}\text{or}l=2,\xi =\frac{7}{3};\\ \\ a=\frac{1}{\xi }\left(1+\nu __1+\frac{2l}{3}\right)\left[\frac{3g__lt_0}{8}\left(t_0\frac{4}{3\sqrt{C_0\frac{2}{3}g_0}}\right)\right]^{2/3}\end{array}$$
(95)
and for $`\xi =l`$
$`l=3,a={\displaystyle \frac{1}{3}}[{\displaystyle \frac{3}{8}}(g_0t_0^2+2at_0g_0){\displaystyle \frac{3}{2}}a\sqrt{C_0{\displaystyle \frac{2}{3}}g_0}+{\displaystyle \frac{g__lt_0}{2\sqrt{C_0\frac{2}{3}g_0}}}]^{1/3}\times `$
$$\left[\frac{3}{8}\left\{\left(1+\nu __1+\frac{2l}{3}\right)g_0t_0^2+2at_0(1+\nu __1)g_0\right\}\frac{3}{2}a(1+\nu __1)\sqrt{C_0\frac{2}{3}g_0}+\frac{\left(1+\nu __1+\frac{2l}{3}\right)g__lt_0}{2\sqrt{C_0\frac{2}{3}g_0}}\right]$$
(96)
We note that the expressions for ‘$`a`$’ is very complicated both in equations (79) and (83). So no definite conclusion is possible on the role of pressure in determining in the final state of collapse by ORNG.
## VI Discussions and Concluding Remarks
An extensive analysis of the four dimensional Szekeres model has been done for the matter containing pressure. When matter is in the form of perfect fluid then the isotropic pressure turns out to be a function of time only while the matter density is a function of all the four space-time variables. In this case, assuming a polynomial form for pressure, cosmological solutions have been obtained and their asymptotic behaviour have been studied. Both in quasi-spherical and quasi-cylindrical model the solution approaches isotropy along fluid world line as $`t\mathrm{}`$.
Secondly, for the matter with tangential stress only, solutions are possible for quasi-cylindrical model. Here both the tangential stress and the matter density turns out to be a function of $`t`$ and $`r`$ only. The scale factor $`R`$ has parametric solution as for dust model and does not depend on the tangential stress. However, choosing the parameter $`C_1=0`$, $`R`$ has a power law solution and it is possible to have a complete solution if we assume the tangential stress proportional to $`t^2`$.
Lastly, gravitational collapse has been studied in details for anisotropic pressure (i.e., both radial and tangential pressures are non-zero and distinct) in quasi-spherical model. Here we have to assume the radial pressure as a function of $`r`$ and $`t`$ of the form (see eq. (50)) $`p_r=g(r)/R^n`$. Also to solve the differential equation in $`R`$ (see eq. (53)) we consider only the marginally bound case (i.e., $`f=0`$) only. Then equation (64) shows a comparative study between the time of formation of trapped surface and the time of formation of central singularity. To simplified further, we choose $`n=3/2`$ and detailed analysis has been done using equation (65). Table I shows all possibilities for the parameters involved in the expression. If the initial density gradient at the centre is positive definite (or negative definite) then as in dust case we have definitely a black hole (or naked singularity) as the final state of collapse. But when $`\rho _1`$ has no definite sign (as in third and fourth cases) then for black hole solution it is possible to have negative density gradient at the centre initially. In fact, near the singularity if the initial density and pressure has identical behaviour (i.e., increase or decrease simultaneously) then even with negative density gradient (initially at the centre) we can have black hole as the end state but if the initial density and pressure has opposite behaviour (i.e., one decrease while other increases and vice versa) then we have identical character as in dust case. So we conclude that pressure tries to resist the formation of naked singularity. Finally we have studied the geodesics to examine whether it is possible to have any future directed non space-like geodesic terminating in the past at the singularity. For simplicity, we have considered only radial null geodesic and it is found that the end state of collapse is characterized by the coefficients of the series expansion of initial density and pressure (radial). Due to complicated expressions we can not definitely characterize the role of pressure. Therefore, in the context of local visibility, we say that pressure tries to cover the singularity.
Acknowledgement:
The authors are thankful to IUCAA for worm hospitality where the major part of the work has been done. One of the authors (U.D) is thankful to CSIR (Govt. of India) for awarding a Senior Research Fellowship.
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$`[21]`$ U. Debnath, S. Chakraborty and J. D. Barrow, Gen. Rel. Grav. 36 231 (2004); U. Debnath and S. Chakraborty, JCAP 05 001 (2004).
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$`[23]`$ S. Chakraborty and U. Debnath, Int. J. Mod. Phys. D 13 1085 (2004)).
$`[24]`$ Note that equation (22) has many possible solutions of which one is given in equation (24). Another solution can be taken as $`e^\nu =A_1(r)\text{sin}(\lambda x)+A_2(r)\text{cos}(\lambda x)+B_1(r)\text{sinh}(\lambda y)+B_2(r)\text{cosh}(\lambda y)`$ with $`\lambda ^2(A_1^2+A_2^2+B_1^2B_2^2)=f(r)1`$.
$`[25]`$ D. A. Szafron, J. Math. Phys. 18 1673 (1977).
$`[26]`$ D. A. Szafron and J. Wainwright J. Math. Phys. 18 1668 (1977).
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# Direct image of logarithmic complexes and infinitesimal invariants of cycles
(Date: Aug. 28, 2005, v.4)
## Abstract.
We show that the direct image of the filtered logarithmic de Rham complex is a direct sum of filtered logarithmic complexes with coefficients in variations of Hodge structures, using a generalization of the decomposition theorem of Beilinson, Bernstein and Deligne to the case of filtered $`D`$-modules. The advantage of using the logarithmic complexes is that we have the strictness of the Hodge filtration by Deligne after taking the cohomology group in the projective case. As a corollary, we get the total infinitesimal invariant of a (higher) cycle in a direct sum of the cohomology of filtered logarithmic complexes with coefficients, and this is essentially equivalent to the cohomology class of the cycle.
Introduction
Let $`X`$, $`S`$ be complex manifolds or smooth algebraic varieties over a field of characteristic zero. Let $`f:XS`$ be a projective morphism, and $`D`$ be a divisor on $`S`$ such that $`f`$ is smooth over $`SD`$. We have a filtered locally free $`𝒪`$-module $`(V^i,F)`$ on $`SD`$ underlying a variation of Hodge structure whose fiber $`V_s^i`$ at $`sSD`$ is the cohomology of the fiber $`H^i(X_s,)`$. If $`D`$ is a divisor with normal crossings on $`S`$, let $`\stackrel{~}{V}^i`$ denote the Deligne extension of $`V^i`$ such that the the eigenvalues of the residue of the connection are contained in $`[0,1)`$. The Hodge filtration $`F`$ is naturally extended to $`\stackrel{~}{V}^i`$ by . We have the logarithmic de Rham complex
$$\text{DR}_{\mathrm{log}}(\stackrel{~}{V}^i)=\mathrm{\Omega }_S^{}(\mathrm{log}D)_𝒪\stackrel{~}{V}^i,$$
which has the Hodge filtration $`F^p`$ defined by $`\mathrm{\Omega }_S^j(\mathrm{log}D)_𝒪F^{pj}\stackrel{~}{V}^i`$. In general, $`V^i`$ can be extended to a regular holonomic $`𝒟_S`$-module $`M^i`$ on which a local defining equation of $`D`$ acts bijectively. By , $`M^i`$ and hence the de Rham complex $`\text{DR}(M^i)`$ have the Hodge filtration $`F`$. If $`Y:=f^{}D`$ is a divisor with normal crossings on $`X,`$ then $`\mathrm{\Omega }_X^{}(\mathrm{log}Y)`$ has the Hodge filtration $`F`$ defined by the truncation $`\sigma `$ (see ) as usual, i.e. $`F^p\mathrm{\Omega }_X^{}(\mathrm{log}Y)=\mathrm{\Omega }_X^p(\mathrm{log}Y)`$.
Theorem 1. Assume $`Y=f^{}D`$ is a divisor with normal crossings. There is an increasing split filtration $`L`$ on the filtered complex $`𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ such that we have noncanonical and canonical isomorphisms in the filtered derived category:
$`𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ $`_i(\text{DR}(M^i),F)[i],`$
$`\text{Gr}_i^L𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ $`=(\text{DR}(M^i),F)[i].`$
If $`D`$ is a divisor with normal crossings, we have also
$`𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ $`_i(\text{DR}_{\mathrm{log}}(\stackrel{~}{V}^i),F)[i],`$
$`\text{Gr}_i^L𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ $`=(\text{DR}_{\mathrm{log}}(\stackrel{~}{V}^i),F)[i].`$
This follows from the decomposition theorem (see ) extended to the case of the direct image of $`(𝒪_X,F)`$ as a filtered $`𝒟`$-module, see . Note that Hodge modules do not appear in the last statement if $`D`$ is a divisor with normal crossings. The assertion becomes more complicated in the non logarithmic case, see Remark (i) in (2.5). A splitting of the filtration $`L`$ is given by choosing the first noncanonical isomorphism in the filtered decomposition theorem, see (1.4.2). A canonical choice of the splitting is given by choosing an relatively ample class, see .
Let $`\text{CH}^p(XY,n)`$ be Bloch’s higher Chow group, see . In the analytic case, we assume for simplicity that $`f:(X,Y)(S,D)`$ is the base change of a projective morphism of smooth complex algebraic varieties $`f^{}:(X^{},Y^{})(S^{},D^{})`$ by an open embedding of complex manifolds $`SS_{\mathrm{a}n}^{}`$, and an element of $`\text{CH}^p(XY,n)`$ is the restriction of an element of $`\text{CH}^p(X^{}Y^{},n)`$ to $`XY`$. If $`n=0`$, we may assume that it is the restriction of an analytic cycle of codimension $`p`$ on $`X`$. From Theorem 1, we can deduce
Corollary 1. With the above notation and assumption, let $`\xi \text{CH}^p(XY,n)`$. Then, choosing a splitting of the filtration $`L`$ in Theorem 1 (or more precisely, choosing the first noncanonical isomorphism in the filtered decomposition theorem (1.4.2)), we have the total infinitesimal invariant
$`\delta _{S,D}(\xi )=(\delta _{S,D}^i(\xi ))_{i0}𝐇^i(S,F^p\text{DR}(M^{2pni})),`$
$`(\text{resp.}`$ $`\overline{\delta }_{S,D}(\xi )=(\overline{\delta }_{S,D}^i(\xi ))_{i0}𝐇^i(S,\text{Gr}_F^p\text{DR}(M^{2pni})),)`$
where $`\delta _{S,D}^i(\xi )`$ (resp. $`\overline{\delta }_{S,D}^i(\xi ))`$ is independent of the choice of a splitting if the $`\delta _{S,D}^j(\xi )`$ (resp. $`\overline{\delta }_{S,D}^j(\xi ))`$ vanish for $`j<i`$. In the case $`D`$ is a divisor with normal crossings, the assertion holds with $`\text{DR}(M^{2pni})`$ replaced by $`\text{DR}_{\mathrm{log}}(\stackrel{~}{V}^{2pni})`$.
This shows that the infinitesimal invariants in , , , , , can be defined naturally in the cohomology of filtered logarithmic complexes with coefficients in variations of Hodge structures if $`D`$ is a divisor with normal crossings, see (2.4) for the compatibility with . Note that if $`S`$ is Stein or affine, then $`𝐇^i(S,F^p\text{DR}_{\mathrm{log}}(\stackrel{~}{V}^q))`$ is the $`i`$-th cohomology group of the complex whose $`j`$-th component is $`\mathrm{\Gamma }(S,\mathrm{\Omega }_S^j(\mathrm{log}D)_𝒪F^{pj}\stackrel{~}{V}^q)`$. If $`D`$ is empty, then an inductive definition of $`\delta _{S,D}^i(\xi )`$, $`\overline{\delta }_{S,D}^i(\xi )`$ was given by Shuji Saito using the filtered Leray spectral sequence together with the $`E_2`$-degeneration argument in . He also showed that the infinitesimal invariants depend only on the cohomology class of the cycle. If $`S`$ is projective, then it follows from that the total infinitesimal invariant $`(\delta _{S,D}^i(\xi ))`$ is equivalent to the cycle class of $`\xi `$ in $`H_{\text{DR}}^{2pn}(XY)`$ by the strictness of the Hodge filtration, and the filtration $`L`$ comes from the Leray filtration on the cohomology of $`XY`$, see Remark (iii) in (2.5).
Corollary 1 is useful to study the behavior of the infinitesimal invariants near the boundary of the variety. If $`D`$ is empty, let $`\delta _S^i(\xi )`$ denote $`\delta _{S,D}^i(\xi )`$. We can define $`\delta _{\text{DR},S}^i(\xi )`$ as in by omitting $`F^p`$ before DR in Corollary 1 where $`D=\mathrm{}`$.
Corollary 2. Assume $`S`$ is projective. Let $`U=SD`$. Then for each $`i0`$, $`\delta _{S,D}^i(\xi )`$, $`\overline{\delta }_{S,D}^i(\xi )`$, $`\delta _U^i(\xi )`$ and $`\delta _{\text{DR},U}^i(\xi )`$ are equivalent to each other, i.e. one of them vanishes if and only if the others do.
Indeed, $`(\delta _{\text{DR},U}^i(\xi ))`$ is determined by $`(\delta _U^i(\xi ))`$, and $`(\delta _U^i(\xi ))`$ by $`(\delta _{S,D}^i(\xi ))`$. Moreover, $`(\delta _{S,D}^i(\xi ))`$ is equivalent to $`(\delta _{\text{DR},U}^i(\xi ))`$ by the strictness of the Hodge filtration applied to $`(X,Y)`$ together with Theorem 1, see (2.3). For the relation with $`\overline{\delta }_{S,D}^i(\xi )`$, see (2.1). Note that the equivalence between $`\delta _U^i(\xi )`$ and $`\delta _{\text{DR},U}^i(\xi )`$ in the case of algebraic cycles (i.e. $`n=0)`$ was first found by J.D. Lewis and Shuji Saito in (assuming a conjecture of Brylinski and Zucker and the Hodge conjecture and using an $`L^2`$-argument). The above arguments seem to be closely related with their question, see also Remark (i) in (2.5) below.
As another corollary of Theorem 1 we have
Corollary 3. Assume $`f`$ induces an isomorphism over $`SD`$, and $`Y=f^{}D`$ is a divisor with normal crossings on $`X`$. Then
$$R^if_{}\mathrm{\Omega }_X^p(\mathrm{log}Y)=0\text{if}i+p>dimX.$$
This follows immediately from Theorem 1 since $`M^i=0`$ for $`i0`$. Corollary 3 is an analogue of the vanishing theorem of Kodaira-Nakano. However, this does not hold for a non logarithmic complex (e.g. if $`f`$ is a blow-up with a point center). This corollary was inspired by a question of A. Dimca.
I would like to thank Dimca, Lewis and Shuji Saito for good questions and useful suggestions.
In Section 1, we prove Theorem 1 after reviewing some basic facts on filtered differential complexes. In Section 2 we explain the application of Theorem 1 to the infinitesimal invariants of (higher) cycles. In Section 3 we give some examples using Lefschetz pencils.
1. Direct image of logarithmic complexes
1.1. Filtered differential complexes. Let $`X`$ be a complex manifold or a smooth algebraic variety over a field of characteristic zero. Let $`D^bF(𝒟_X)`$ (resp. $`D^bF(𝒟_X)^r)`$ be the bounded derived category of filtered left (resp. right) $`𝒟_X`$-modules. Let $`D^bF(𝒪_X,\text{Diff})`$ be the bounded derived category of filtered differential complexes $`(L,F)`$ where $`F`$ is exhaustive and locally bounded below (i.e. $`F_p=0`$ for $`p0`$ locally on $`X`$), see , 2.2. We have an equivalence of categories
$`(\mathrm{1.1.1})`$
$$\text{DR}^1:D^bF(𝒪_X,\text{Diff})D^bF(𝒟_X)^r,$$
whose quasi-inverse is given by the de Rham functor $`\text{DR}^r`$ for right $`𝒟`$-modules, see (1.2) below. Recall that, for a filtered $`𝒪_X`$-module $`(L,F)`$, the associated filtered right $`𝒟`$-module $`\text{DR}^1(L,F)`$ is defined by
$`(\mathrm{1.1.2})`$
$$\text{DR}^1(L,F)=(L,F)_𝒪(𝒟,F),$$
and the morphisms $`(L,F)(L^{},F)`$ in $`MF(𝒪_X,\text{Diff})`$ correspond bijectively to the morphisms of filtered $`𝒟`$-modules $`\text{DR}^1(L,F)\text{DR}^1(L^{},F)`$. More precisely, the condition on $`(L,F)(L^{},F)`$ is that the composition
$$F_pLLL^{}L^{}/F_qL^{}$$
is a differential operator of order $`pq1`$. The proof of (1.1.1) can be reduced to the canonical filtered quasi-isomorphism for a filtered right $`𝒟`$-module $`(M,F)`$
$$\text{DR}^1\text{DR}^r(M,F)(M,F),$$
which follows from a calculation of a Koszul complex.
Note that the direct image $`f_{}`$ of filtered differential complexes is defined by the sheaf-theoretic direct image $`𝐑f_{}`$, and this direct image is compatible with the direct image $`f_{}`$ of filtered $`𝒟`$-modules via (1.1.1), see , 2.3. So we get
$`(\mathrm{1.1.3})`$
$$𝐑f_{}=\text{DR}^rf_{}\text{DR}^1:D^bF(𝒪_X,\text{Diff})D^bF(𝒪_S,\text{Diff}),$$
where we use $`\text{DR}^r`$ for right $`𝒟`$-modules (otherwise there is a shift of complex).
1.2. De Rham complex. The de Rham complex $`\text{DR}^r(M,F)`$ of a filtered right $`𝒟`$-module $`(M,F)`$ is defined by
$`(\mathrm{1.2.1})`$
$$(\text{DR}^r(M,F))^i=^i\mathrm{\Theta }_X_𝒪(M,F[i])\text{for}i0.$$
Here $`(F[i])_p=F_{p+i}`$ in a compatible way with $`(F[i])^p=F^{pi}`$ and $`F_p=F^p`$. Recall that the filtered right $`𝒟`$-module associated with a filtered left $`𝒟`$-module $`(M,F)`$ is defined by
$`(\mathrm{1.2.2})`$
$$(M,F)^r:=(\mathrm{\Omega }_X^{dimX},F)_𝒪(M,F),$$
where $`\text{Gr}_p^F\mathrm{\Omega }_X^{dimX}=0`$ for $`pdimX`$. This induces an equivalence of categories between the left and right $`𝒟`$-modules. The usual de Rham complex $`\text{DR}(M,F)`$ for a left $`𝒟`$-module is defined by
$`(\mathrm{1.2.3})`$
$$(\text{DR}(M,F))^i=\mathrm{\Omega }_X^i_𝒪(M,F[i])\text{for}i0,$$
and this is compatible with (1.2.1) via (1.2.2) up to a shift of complex, i.e.
$`(\mathrm{1.2.4})`$
$$\text{DR}(M,F)=\text{DR}^r(M,F)^r[dimX].$$
1.3. Logarithmic complex. Let $`X`$ be as in (1.1), and $`Y`$ be a divisor with normal crossings on $`X`$. Let $`(V,F)`$ be a filtered locally free $`𝒪`$-module underlying a polarizable variation of Hodge structure on $`XY`$. Let $`(\stackrel{~}{V},F)`$ be the Deligne extension of $`(V,F)`$ to $`X`$ such that the eigenvalues of the residue of the connection are contained in $`[0,1)`$. Then we have the filtered logarithmic de Rham complex $`\text{DR}_{\mathrm{log}}(\stackrel{~}{V},F)`$ such that $`F^p`$ of its $`i`$-th component is
$$\mathrm{\Omega }_X^i(\mathrm{log}Y)F^{pi}\stackrel{~}{V}.$$
If $`(M,F)=(𝒪_X,F)`$ with $`\text{Gr}_p^F𝒪_X=0`$ for $`p0`$, then
$$\text{DR}_{\mathrm{log}}(𝒪_X,F)=(\mathrm{\Omega }_X^{}(\mathrm{log}X),F).$$
Let $`\stackrel{~}{V}(Y)`$ be the localization of $`\stackrel{~}{V}`$ by a local defining equation of $`Y`$. This is a regular holonomic left $`𝒟_X`$-module underlying a mixed Hodge module, and has the Hodge filtration $`F`$ which is generated by the Hodge filtration $`F`$ on $`\stackrel{~}{V}`$, i.e.
$$F_p\stackrel{~}{V}(Y)=_\nu ^\nu F^{p+|\nu |}\stackrel{~}{V},$$
where $`F_p=F^p`$ and $`^\nu =_i_i^{\nu _i}`$ with $`_i=/x_i`$. Here $`(x_1,\mathrm{},x_n)`$ is a local coordinate system such that $`Y`$ is contained in $`\{x_1\mathrm{}x_n=0\}`$. By , 3.11, we have a filtered quasi-isomorphism
$`(\mathrm{1.3.1})`$
$$\text{DR}_{\mathrm{log}}(\stackrel{~}{V},F)\stackrel{}{}\text{DR}(\stackrel{~}{V}(Y),F).$$
This generalizes the filtered quasi-isomorphism in
$`(\mathrm{1.3.2})`$
$$(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)\stackrel{}{}\text{DR}(𝒪_X(Y),F).$$
Note that the direct image of the filtered $`𝒟_X`$-module $`(\stackrel{~}{V}(Y),F)`$ by $`Xpt`$ in the case $`X`$ projective (or proper algebraic) is given by the cohomology group of the de Rham complex $`\text{DR}(\stackrel{~}{V}(Y),F)`$ (up to a shift of complex) by definition, and the Hodge filtration $`F`$ on the direct image is strict by the theory of Hodge modules. So we get
$`(\mathrm{1.3.3})`$ $`F^p𝐇^i(XY,\text{DR}(V))`$ $`:=F^p𝐇^i(X,\text{DR}(\stackrel{~}{V}(Y)))`$
$`=𝐇^i(X,F^p\text{DR}(\stackrel{~}{V}(Y)))`$ $`=𝐇^i(X,F^p\text{DR}_{\mathrm{log}}(\stackrel{~}{V})).`$
1.4. Decomposition theorem. Let $`f:XS`$ be a projective morphism of complex manifolds or smooth algebraic varieties over a field of characteristic zero. Then the decomposition theorem of Beilinson, Bernstein and Deligne is extended to the case of Hodge modules (, ), and we have noncanonical and canonical isomorphisms
$`(\mathrm{1.4.1})`$ $`f_{}(𝒪_X,F)_j^jf_{}(𝒪_X,F)[j]`$ $`\text{in}D^bF(𝒟_S),`$
$`^jf_{}(𝒪_X,F)=_{ZS}(M_Z^j,F)`$ $`\text{in}MF(𝒟_S),`$
where $`Z`$ are irreducible closed analytic or algebraic subsets of $`S`$, and $`(M_Z^j,F)`$ are filtered $`𝒟_S`$-modules underlying a pure Hodge module of weight $`j+dimX`$ and with strict support $`Z`$, i.e. $`M_Z^j`$ has no nontrivial sub nor quotient module whose support is strictly smaller than $`Z`$. (Here $`MF(𝒟_S)`$ denotes the category of filtered left $`𝒟_S`$-modules.) Indeed, the second canonical isomorphism follows from the strict support decomposition which is part of the definition of pure Hodge modules, see , 5.1.6. The first noncanonical isomorphism follows from the strictness of the Hodge filtration and the relative hard Lefschetz theorem for the direct image (see , 5.3.1) using the $`E_2`$-degeneration argument in together with the equivalence of categories $`D^bF(𝒟_S)D^bG(_S)`$. Here $`_S=_iF_i𝒟_S`$ and $`D^bG(_S)`$ is the derived category of bounded complexes of graded left $`_S`$-modules $`M_{}^{}`$ such that $`M_i^j=0`$ for $`i0`$ or $`|j|0`$, see , 2.1.12. We need a derived category associated to some abelian category in order to apply the argument in (see also ). In the algebraic case, we can also apply to the derived category of mixed Hodge modules on $`S`$ and it is also possible to use , 4.5.4 to show the first noncanonical isomorphism.
If $`f`$ is smooth over the complement of a divisor $`DS`$ and $`Y:=f^{}D`$ is a divisor with normal crossings, then the filtered direct image $`f_{}(𝒪_X(Y),F)`$ is strict (see , 2.15), and we have noncanonical and canonical isomorphisms
$`(\mathrm{1.4.2})`$ $`f_{}(𝒪_X(Y),F)_j^jf_{}(𝒪_X(Y),F)[j]`$ $`\text{in}D^bF(𝒟_S),`$
$`^jf_{}(𝒪_X(Y),F)=(M_S^j(D),F)`$ $`\text{in}MF(𝒟_S).`$
Here $`(M_S^j(D),F)`$ is the ‘localization’ of $`(M_S^j,F)`$ along $`D`$ which is the direct image of $`(M_S^j,F)|_U`$ by the open embedding $`U:=SDS`$ in the category of filtered $`𝒟`$-modules underlying mixed Hodge modules. (By the Riemann-Hilbert correspondence, this gives the direct image in the category of complexes with constructible cohomology because $`D`$ is a divisor.) The Hodge filtration $`F`$ on the direct image is determined by using the $`V`$-filtration of Kashiwara and Malgrange, and $`(M_S^j(D),F)`$ is the unique extension of $`(M_S^j,F)|_U`$ which underlies a mixed Hodge module on $`S`$ and whose underlying $`𝒟_S`$-module is the direct image in the category of regular holonomic $`𝒟_S`$-modules, see , 2.11. So the second canonical isomorphism follows because the left-hand side satisfies these conditions. (Note that $`(M_Z^j,F)`$ for $`ZS`$ vanishes by the localization, because $`ZD`$ if $`(M_Z^j,F)0`$ .) The first noncanonical isomorphism follows from the strictness of the Hodge filtration and the relative hard Lefschetz theorem by the same argument as above.
1.5. Proof of Theorem 1. Let $`r=dimXdimS`$. By (1.1.3), (1.3.2) and (1.4.2), we have isomorphisms
$`(\mathrm{1.5.1})`$ $`𝐑f_{}(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$ $`=\text{DR}^rf_{}\text{DR}^1(\mathrm{\Omega }_X^{}(\mathrm{log}Y),F)`$
$`=\text{DR}^rf_{}(𝒪_X(Y),F)[dimX]`$
$`_i\text{DR}(M_S^i(D),F)[ri],`$
where the shift of complex by $`r`$ follows from the difference of the de Rham complex for left and right $`𝒟`$-modules. Furthermore, letting $`L`$ be the filtration induced by $`\tau `$ on the complex of filtered $`𝒟_S`$-modules $`f_{}(𝒪_X(Y),F)[r]`$, we have a canonical isomorphism
$`(\mathrm{1.5.2})`$
$$\text{Gr}_i^Lf_{}(𝒪_X(Y),F)[r]=(M_S^{ir}(D),F)[i],$$
and the first assertion follows by setting $`M^i=M_S^{ir}(D)`$. The second assertion follows from the first by (1.3.1). This competes the proof of Theorem 1.
2. Infinitesimal invariants of cycles
2.1. Cycle classes. Let $`X`$ be a complex manifold, and $`𝒞^,`$ denote the double complex of vector spaces of currents on $`X`$. The associated single complex is denoted by $`𝒞^{}`$. Let $`F`$ be the Hodge filtration by the first index of $`𝒞^,`$ (using the truncation $`\sigma `$ in ). Let $`\xi `$ be an analytic cycle of codimension $`p`$ on $`X`$. Then it is well known that $`\xi `$ defines a closed current in $`F^p𝒞^{2p}`$ by integrating the restrictions of $`C^{\mathrm{}}`$ forms with compact supports on $`X`$ to the smooth part of the support of $`\xi `$ (and using a triangulization or a resolution of singularities of the cycle). So we have a cycle class of $`\xi `$ in $`H^{2p}(X,F^p\mathrm{\Omega }_X^{})`$.
Assume $`X`$ is a smooth algebraic variety over a field $`k`$ of characteristic zero. Then the last assertion still holds (where $`\mathrm{\Omega }_X^{}`$ means $`\mathrm{\Omega }_{X/k}^{}`$), see . Moreover, for the higher Chow groups, we have the cycle map (see , , , , )
$$cl:\text{CH}^p(X,n)F^pH_{\text{DR}}^{2pn}(X),$$
where the Hodge filtration $`F`$ is defined by using a smooth compactification of $`X`$ whose complement is a divisor with normal crossings, see . This cycle map is essentially equivalent to the cycle map to $`\text{Gr}_F^pH_{\text{DR}}^{2pn}(X)`$ because we can reduce to the case $`k=`$ where we have the cycle map
$$cl:\text{CH}^p(X,n)\text{Hom}_{\text{MHS}}(,H^{2pn}(X,)(p)),$$
and morphisms of mixed Hodge structures are strictly compatible with the Hodge filtration $`F`$.
2.2. Proof of Corollary 1. By (2.1) $`\xi `$ has the cycle class in
$$H^{2pn}(X,F^p\mathrm{\Omega }_X^{}(\mathrm{log}Y)).$$
By theorem 1, this gives the total infinitesimal invariant
$$\delta _{S,D}(\xi )=(\delta _{S,D}^{2pni}(\xi ))_i𝐇^{2pni}(S,F^p\text{DR}(M^i)),$$
and similarly for $`\overline{\delta }_{S,D}(\xi )`$. So the assertion follows.
2.3. Proof of Corollary 2. Choosing the first noncanonical isomorphism in the filtered decomposition theorem (1.4.2), we get canonical morphisms compatible with the direct sum decompositions
$`_{i0}𝐇^i(S,F^p\text{DR}(M^{qi}))`$ $`_{i0}𝐇^i(SD,F^p\text{DR}(M^{qi}))`$
$`_{i0}𝐇^i(SD,\text{DR}(M^{qi})),`$
and these are identified with the canonical morphisms
$`𝐇^q(X,F^p\mathrm{\Omega }_X^{}(\mathrm{log}Y))`$ $`𝐇^q(XY,F^p\mathrm{\Omega }_{XY}^{})`$
$`𝐇^q(XY,\mathrm{\Omega }_{XY}^{}).`$
By Deligne , the composition of the last two morphisms is injective because of the strictness of the Hodge filtration, see also (1.3). So we get the equivalence of $`\delta _{S,D}^i(\xi )`$, $`\delta _U^i(\xi )`$, $`\delta _{\text{DR},U}^i(\xi )`$. The equivalence with $`\overline{\delta }_{S,D}^i(\xi )`$ follows from (2.1).
2.4. Compatibility with the definition in . When $`D`$ is empty, the infinitesimal invariants are defined in by using the extension groups of filtered $`𝒟`$-modules together with the forgetful functor from the category of mixed Hodge modules to that of filtered $`𝒟`$-modules. Its compatibility with the definition in this paper follows from the equivalence of categories (1.1.1) and the compatibility of the direct image functors (1.1.3).
Note that for $`(L,F)D^bF(𝒪_X,\text{Diff})`$ in the notation of (1.1), we have a canonical isomorphism
$`(\mathrm{2.4.1})`$
$$\text{Ext}^i((\mathrm{\Omega }_X^{},F),(L,F))=𝐇^i(X,F_0L),$$
where the extension group is taken in $`D^bF(𝒪_X,\text{Diff})`$. Indeed, the left-hand side is canonically isomorphic to
$`\text{Ext}^i(\text{DR}^1(\mathrm{\Omega }_X^{},F),\text{DR}^1(L,F))`$
$`=𝐇^i(X,F_0om_𝒟(\text{DR}(𝒟_X,F),\text{DR}^1(L,F))),`$
$`=𝐇^i(X,F_0\text{DR}^r\text{DR}^1(L)),`$
and the last group is isomorphic to the right-hand side of (2.4.1) which is independent of a representative of $`(L,F)`$. If $`X`$ is projective, then this assertion follows also from the adjoint relation for filtered $`𝒟`$-modules.
If $`X`$ is smooth projective and $`Y`$ is a divisor with normal crossings, then the cycle class can be defined in
$`\text{Ext}^{2p}((\mathrm{\Omega }_X^{},F),\mathrm{\Omega }_X^{}(\mathrm{log}Y),F[p]))`$ $`=𝐇^{2p}(X,F^p\mathrm{\Omega }_X^{}(\mathrm{log}Y))`$
$`=F^p𝐇^{2p}(X,\mathrm{\Omega }_X^{}(\mathrm{log}Y)).`$
2.5. Remarks. (i) If we use (1.4.1) instead of (1.4.2) we get an analogue of Theorem 1 for non logarithmic complexes. However, the assertion becomes more complicated, and we get noncanonical and canonical isomorphisms
$`(\mathrm{2.5.1})`$ $`𝐑f_{}(\mathrm{\Omega }_X^{},F)`$ $`_{i,ZS}(\text{DR}(M_Z^{ir}),F)[i].`$
$`\text{Gr}_i^L𝐑f_{}(\mathrm{\Omega }_X^{},F)`$ $`=_{ZS}(\text{DR}(M_Z^{ir}),F)[i].`$
This implies an analogue of Corollary 1. If $`D`$ is a divisor with normal crossings, we have a filtered quasi-isomorphism for $`Z=S`$
$`(\mathrm{2.5.2})`$
$$(\text{DR}_{\mathrm{log}}(\stackrel{~}{M}_S^{ir}),F)\stackrel{}{}(\text{DR}(M_S^{ir}),F),$$
where $`\text{DR}_{\mathrm{log}}(\stackrel{~}{M}_S^{ir})`$ is the intersection of $`\text{DR}(M_S^{ir})`$ with $`\text{DR}_{\mathrm{log}}(\stackrel{~}{V}_S^i)`$. This seems to be related with a question of Lewis and Shuji Saito, see also .
(ii) If $`dimS=1`$, we can inductively define the infinitesimal invariants in Corollary 1 by an argument similar to using .
(iii) Assume $`S`$ is projective and $`D`$ is a divisor with normal crossings. Then the Leray filtration for $`XSpt`$ is given by the truncation $`\tau `$ on the complex of filtered $`𝒟_S`$-modules $`f_{}(𝒪_X(Y),F)`$, and gives the Leray filtration on the cohomology of $`XY`$ (induced by the truncation $`\tau `$ as in ). Indeed, the graded pieces $`^jf_{}(𝒪_X(Y),F)`$ of the filtration $`\tau `$ on $`S`$ coincide with $`(\stackrel{~}{V}^{j+r}(D),F)`$, and give the open direct images by $`US`$ of the graded pieces $`(V^{j+r},F)`$ of the filtration $`\tau `$ on $`U`$ as filtered $`𝒟`$-modules underlying mixed Hodge modules. Note that the morphism $`US`$ is open affine so that the direct image preserves regular holonomic $`𝒟`$-modules.
3. Examples
3.1. Lefschetz pencils. Let $`Y`$ be a smooth irreducible projective variety of dimension $`n`$ embedded in a projective space $`𝒫`$ over $``$. We assume that $`Y𝒫`$ and $`Y`$ is not contained in a hyperplane of $`𝒫`$ so that the hyperplane sections of $`Y`$ are parametrized by the dual projective spaces $`𝒫^{}`$. Let $`D𝒫^{}`$ denote the discriminant. This is the image of a projective bundle over $`Y`$ (consisting of hyperplanes tangent to $`Y`$), and hence $`D`$ is irreducible. At a smooth point of $`D`$, the corresponding hyperplane section of $`Y`$ has only one ordinary double point. We assume that the associated vanishing cycle is not zero in the cohomology of general hyperplane section $`X`$. This is equivalent to the non surjectivity of $`H^{n1}(Y)H^{n1}(X)`$.
A Lefschetz pencil of $`Y`$ is a line $`^1`$ in $`𝒫`$ intersecting the discriminant $`D`$ at smooth points of $`D`$ (corresponding to hyperplane sections having only one ordinary double point). We have a projective morphism $`\pi :\stackrel{~}{Y}^1`$ such that $`\stackrel{~}{Y}_t:=\pi ^1(t)`$ is the hyperplane section corresponding $`t^1𝒫`$ and $`\stackrel{~}{Y}`$ is the blow-up of $`Y`$ along a smooth closed subvariety $`Z`$ of codimension $`2`$ which is the intersection of $`\stackrel{~}{Y}_t`$ for any (or two of) $`t^1`$.
A Lefschetz pencil of hypersurface sections of degree $`d`$ is defined by replacing the embedding of $`Y`$ using $`𝒪_Y(d)`$ so that a hyperplane section corresponds to a hypersurface section of degree $`d`$. Here $`𝒪_Y(d)`$ for an integer $`d`$ denote the invertible sheaf induced by that on $`𝒫`$ as usual.
3.2. Hypersurfaces containing a subvariety. Let $`Y,𝒫`$ be as in (3.2). Let $`E`$ be a closed subvariety (which is not necessarily irreducible nor reduced). Let
$$E_{\{i\}}=\{xE:dimT_xE=i\}.$$
Let $`_E`$ be the ideal sheaf of $`E`$ in $`Y`$. Let $`\delta `$ be a positive integer such that $`_E(\delta )`$ is generated by global sections. By , (or ) we have the following
(3.2.1) If $`dimY>\mathrm{max}\{dimE_{\{i\}}+i\}`$ and $`d\delta `$, then there is a smooth hypersurface section of degree $`d`$ containing $`E`$.
We have furthermore
(3.2.2) If $`dimY>\mathrm{max}\{dimE_{\{i\}}+i\}+1`$ and $`d\delta +1`$, then there is a Lefschetz pencil of hypersurface sections of degree $`d`$ containing $`E`$.
Indeed, we have a pencil such that $`\stackrel{~}{Y}_t`$ has at most isolated singularities, because $`\stackrel{~}{Y}_t`$ is smooth near the center $`Z`$ which is the intersection of generic two hypersurfaces sections containing $`E`$, and hence is smooth, see , (or ). Note that a local equation of $`\stackrel{~}{Y}_t`$ near $`Z`$ is given by $`ftg`$ if $`t`$ is identified with an appropriate affine coordinate of $`^1`$ where $`f,g`$ are global sections of $`_E(d)`$ corresponding to smooth hypersurface sections.
To get only ordinary double points, note first that the parameter space of the hypersurfaces containing $`E`$ is a linear subspace of $`𝒫^{}`$. So it is enough to show that this linear subspace contains a point of the discriminant $`D`$ corresponding to an ordinary double point. Thus we have to show that an isolated singularity can be deformed to ordinary double points by replacing the corresponding section $`h\mathrm{\Gamma }(Y,_E(d))`$ with $`h+_it_ig_i`$ where $`g_i\mathrm{\Gamma }(Y,_E(d))`$ and the $`t_i`$ are general with sufficiently small absolute values. Since $`d\delta +1`$, we see that $`\mathrm{\Gamma }(Y,_E(d))`$ generates the $`1`$-jets at each point of the complement of $`E`$. So the assertion follows from the fact that for a function with an isolated singularity $`f`$, the singularities of $`\{f+_it_ix_i=0\}`$ are ordinary double points if $`t_1,\mathrm{},t_n`$ are general, where $`x_1,\mathrm{},x_n`$ are local coordinates. (Note that $`f`$ has an ordinary double point if and only if the morphism defined by $`(f/x_1,\mathrm{},f/x_n)`$ is locally biholomorphic at this point.)
3.3. Construction. For $`Y,𝒫`$ be as in (3.1), let $`i_{Y,𝒫}:Y𝒫`$ denote the inclusion. Assume
$`(\mathrm{3.3.1})`$
$$i_{Y,𝒫}^{}:H^j(𝒫)H^j(Y)\text{is surjective for any}jdimY,$$
where cohomology has coefficients in any field of characteristic zero. This condition is satisfied if $`Y`$ is a complete intersection.
Let $`E_1`$, $`E_2`$ be $`m`$-dimensional irreducible closed subvarieties of $`Y`$ such that
$$E_1E_2=\mathrm{},\mathrm{deg}E_1=\mathrm{deg}E_2.$$
Here $`dimY=n=2m+s+1`$ with $`m0`$, $`s1`$. Let $`E=E_1E_2`$. With the notation of (3.2), assume
$`(\mathrm{3.3.2})`$
$$d>\delta ,dimY>\mathrm{max}\{dimE_{\{i\}}+i\}+s,$$
$`(\mathrm{3.3.3})`$
$$i_{X^{(j)},Y}^{}:H^{nj}(Y)H^{nj}(X^{(j)})\text{is not surjective for}js,$$
where $`X^{(j)}`$ is a general complete intersection of multi degree $`(d,\mathrm{},d)`$ and of codimension $`j`$ in $`Y`$. (This is equivalent to the condition that the vanishing cycles for a hypersurface $`X^{(j)}`$ of $`X^{(j1)}`$ are nonzero.)
Let $`X`$ be a general hypersurface of degree $`d`$ in $`Y`$ containing $`E`$, see (3.2.1). Let $`L`$ denote the intersection of $`X`$ with a general linear subspace of codimension $`m+s`$ in the projective space. Then $`[E_a](a=1,2)`$ and $`c[LX]`$ have the same cohomology class in $`H^{2m+2s}(X)`$ for some $`c`$, because $`dimH^{2m+2s}(X)=1`$ by the weak and hard Lefschetz theorems together with (3.3.1). Let
$$\xi _a=[E_a]c[LX]\text{CH}^{m+s}(X)_{}(a=1,2).$$
These are homologous to zero. It may be expected that one of them is non torsion, generalizing an assertion in . More precisely, let $`S`$ be a smooth affine rational variety defined over a finitely generated subfield $`k`$ of $``$ and parametrizing the smooth hypersurfaces of degree $`d`$ containing $`E`$ as above so that there is the universal family $`𝒳S`$ defined over $`k`$ (see , ). Assume $`X`$ corresponds to a geometric generic point of $`S`$ with respect to $`k`$, i.e. $`X`$ is the geometric generic fiber for some embedding $`k(S)`$. Let
$$\xi _{a,𝒳}=[E_a\times _kS]c[L]_𝒳\text{CH}^{m+s}(𝒳)_{},$$
where $`[L]_𝒳`$ is the pull-back of $`[L]`$ by $`𝒳Y`$. Since the local system $`\{H^{2m+2sj}(𝒳_s)\}`$ on $`S`$ is constant for $`j<s`$ and $`S`$ is smooth affine rational, we see that $`\delta _S^j(\xi _{a,𝒳})=0`$ for $`j<s`$. Then it may be expected that $`\delta _S^s(\xi _{a,𝒳})0`$ for one of $`a`$, where $`S`$ can be replaced by any non empty open subvariety. We can show this for $`s=1`$ as follows. (For $`s>1`$, it may be necessary to assume further conditions on $`d`$, etc.)
3.4. Case $`s=1`$. Consider a Lefschetz pencil $`\pi :\stackrel{~}{Y}^1`$ such that $`\stackrel{~}{Y}_t:=\pi ^1(t)`$ for $`t^1`$ is a hypersurface of degree $`d`$ in $`Y`$ containing $`E`$. Here $`\stackrel{~}{Y}`$ is the blow-up of $`Y`$ along a smooth closed subvariety $`Z`$, and $`Z`$ is the intersection of $`\stackrel{~}{Y}_t`$ for any $`t^1`$. Note that $`\stackrel{~}{Y}_t`$ has an ordinary double point for $`t\mathrm{\Lambda }^1`$, where $`\mathrm{\Lambda }`$ denotes the discriminant, see (3.2.2).
Since $`Z`$ has codimension $`2`$ in $`Y`$, we have the isomorphism
$`(\mathrm{3.4.1})`$
$$H^n(\stackrel{~}{Y})=H^n(Y)H^{n2}(Z),$$
so that the cycle class of $`[E_a\times ^1]c[L]_{\stackrel{~}{Y}}\text{CH}^{m+1}(\stackrel{~}{Y})_{}`$ in $`H^n(\stackrel{~}{Y})`$ is identified with the difference of the cycle class $`cl_Z(E_a)H^{n2}(Z)`$ and the cycle class of $`L`$ in $`H^n(Y)`$. Indeed, the injection $`H^{n2}(Z)H^n(\stackrel{~}{Y})`$ in the above direct sum decomposition is defined by using the projection $`Z\times ^1Z`$ and the closed embedding $`Z\times ^1\stackrel{~}{Y}`$, and the injection $`H^n(Y)H^n(\stackrel{~}{Y})`$ is the pull-back by $`\stackrel{~}{Y}Y`$, see .
By assumption, one of the $`cl_Z(E_a)`$ is not contained in the non primitive part, i.e. not a multiple of the cohomology class of the intersection of general hyperplane sections. Indeed, if both are contained in the non primitive part, then $`cl_Z(E_1)=cl_Z(E_2)`$ and this implies the vanishing of the self intersection number $`E_aE_a`$ in $`Z`$.
We will show that the cycle class of $`[E_a\times ^1]c[L]_{\stackrel{~}{Y}}`$ does not vanish in the cohomology of $`\pi ^1(U)`$ for any non empty open subvariety of $`^1`$, in other words, it does not belong to the image of $`_{t\mathrm{\Lambda }^{}}H_{\stackrel{~}{Y}_t}^n(\stackrel{~}{Y})`$ where $`\mathrm{\Lambda }^{}`$ is any finite subset of $`^1`$ containing $`\mathrm{\Lambda }`$. (Note that the condition for the Lefschetz pencil is generic, and for any proper closed subvariety of the parameter space, there is a Lefschetz pencil whose corresponding line is not contained in this subvariety.)
Thus the assertion is reduced to that $`dimH_{\stackrel{~}{Y}_t}^n(\stackrel{~}{Y})`$ is independent of $`t^1`$ because this implies that the image of $`H_{\stackrel{~}{Y}_t}^n(\stackrel{~}{Y})H^n(\stackrel{~}{Y})`$ is independent of $`t`$. (Note that the Gysin morphism $`H^{n2}(\stackrel{~}{Y}_t)H^n(\stackrel{~}{Y})`$ for a general $`t`$ can be identified with the direct sum of the Gysin morphism $`H^{n2}(\stackrel{~}{Y}_t)H^n(Y)`$ and the restriction morphism $`H^{n2}(\stackrel{~}{Y}_t)H^{n2}(Z)`$ up to a sign, and the image of the last morphism is the non primitive part by the weak Lefschetz theorem.) By duality, this is equivalent to that $`R^n\pi _{}_{\stackrel{~}{Y}}`$ is a local system on $`^1`$. Then it follows from the assumption that the vanishing cycles are nonzero, see (3.3.3).
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# A note on the projective varieties of almost general type
## 1 Introduction
Throughout the paper every variety is defined over the field of complex numbers $`𝐂`$. We follow the notation and terminology in .
###### Definition 1.1.
A $`𝐐`$-Cartier divisor $`D`$ on a projective variety $`X`$ is almost numerically positive (almost nup, for short), if there exists a union $`F`$ of at most countably many prime divisors on $`X`$ such that $`(D,C)>0`$ for every curve $`CF`$ (i.e. if $`(D,C)>0`$ for every very general curve $`C`$). We say that $`D`$ is quasi-numerically positive (quasi-nup, for short), if $`D`$ is nef and almost nup.
###### Definition 1.2.
An algebraic variety $`X`$ is of almost general type, if there exists a projective variety $`M`$ with only terminal singularities such that the canonical divisor $`K_M`$ is almost nup and that $`M`$ is birationally equivalent to $`X`$.
Obviously, an algebraic variety $`X`$ is of almost general type, if it is of general type (i.e. if the geometric Kodaira dimension $`\kappa _{geom}(X)`$ equals the dimension $`dimX`$). However, the converse statement contains the conjecture that every quasi-nup canonical divisor on a minimal variety should be semi-ample, which is the essential part of the famous abundance conjecture (cf. Ambro ).
###### Definition 1.3.
The property $`P(x)`$ holds for every very general point $`x`$ on an algebraic variety $`X`$, if $`P(x)`$ holds for every point $`xXE`$ where $`E`$ is some fixed union of at most countably many prime divisors on $`X`$.
###### Definition 1.4.
For an ambient space $`X`$, a subset $`L`$ is covered by subsets $`D_i`$ ($`iI`$) if $`L_{iI}D_i`$.
###### Definition 1.5.
For an algebraic variety $`X`$, we denote by $`RatX`$ the rational function field of $`X`$.
In this article, we look into the complex algebraic varieties of almost general type and show that the class of varieties of this type coincides with the class of varieties that are neither uniruled nor covered by any family of varieties being birationally equivalent to minimal varieties with numerically trivial canonical divisors, under the minimal model conjecture (see Theorem 3.5). In the category of projective varieties with only terminal singularities, this class is proved to be the class of varieties with almost nup canonical divisors, under the minimal model conjecture (see Theorem 2.3). In contrast with the method in the previous paper , we proceed not depending on the (log) abundance conjecture, which claims the semi-ampleness of nef (log) canonical divisors, and the deformation invariance theorem of Kodaira dimension due to Tsuji and Siu . The key is Tsuji’s theory of numerically trivial fibrations.
We briefly explain the background. The minimal model conjecture states that every projective variety with only $`𝐐`$-factorial terminal singularities would reach a variety with nef canonical divisor (i.e. a minimal variety) or with some Mori fiber space structure, after a finite sequence of divisorial contractions and flips. We emphasize that this conjecture does not contain the abundance conjecture that the (nef) canonical divisor of every minimal variety should be semi-ample. Compared with the mechanism of the minimal model theory, the known proofs of the abundance theorems for surfaces and threefolds are case-by-case and very complicated. So it is meaningful to get geometric results, not depending on the abundance conjecture, but only on the minimal model conjecture.
###### Remark 1.6.
We state the current status of the minimal model conjecture for the convenience of the reader. The general theory reduces this conjecture to the existence and termination of flips. The existence is now a theorem due to Birkar, Cascini, Hacon and McKernan in all dimensions. The termination became a theorem (Shokurov for threefolds and Kawamata-Matsuda-Matsuki for fourfolds) in dimension $`4`$.
## 2 Stability under deformation
In this section, we show that the property “being of almost general type” is stable under deformation.
###### Proposition 2.1.
Let $`f:MS`$ be a surjective morphism between smooth projective varieties such that the field extension $`RatM/RatS`$ is algebraically closed. Then one of the following holds:
(i) The canonical divisor $`K_F`$ is not almost nup for every general fiber $`F`$ of $`f`$.
(ii) The canonical divisor $`K_F`$ is almost nup for every very general fiber $`F`$ of $`f`$.
###### Proof.
Let $`M\times Hilb(M)`$ be the universal family parametrized by the Hilbert scheme $`Hilb(M)`$. Consider the projection morphisms $`p_1:M`$ and $`p_2:Hilb(M)`$.
Let $`H`$ be a hyperplane section of $`S`$. The morphism $`p_2`$ is flat. Thus, for every curve $`C`$ on $`M`$, there exists some irreducible component $`V`$ of $``$ such that $`Vp_2^1([C])`$ and that $`dimVdimp_2(V)=1`$, where $`[C]`$ is the point representing the subscheme $`C`$ of $`M`$. Now define the set $`I:=\{V|V`$ is an irreducible component of $``$, such that $`Vp_2^1([C])`$ for some curve $`C`$ with $`(K_M,C)0`$ and $`(f^{}H,C)=0`$ and that $`dimVdimp_2(V)=1\}`$. Note that $`(f^{}H,C)=0`$ if and only if $`f(C)`$ is a point.
We divide the situation into two cases, according to the bigness of the set $`\{p_1(V)|VI\}`$ in $`M`$.
First we treat the case where $`p_1(V)=M`$ for some $`VI`$.
We apply the normalization, the Stein factorization and the flattening to the morphism $`p_2|_V:Vp_2(V)`$ and obtain the following commutative diagram among projective varieties
$$\begin{array}{ccc}W& \stackrel{\mu }{}& T\\ \rho & & \nu & & \\ V& \underset{p_2|_V}{}& p_2(V)\end{array}$$
(2.1)
with the properties:
(1) the morphism $`\rho `$ is birational
(2) the morphism $`\nu `$ is generically finite
(3) the extension $`RatW/RatT`$ is algebraically closed
(4) the morphism $`\mu `$ is flat
(5) the variety $`T`$ is smooth.
For every curve $`C^{}(\nu \mu )^1([C])=\rho ^1(p_2^1([C]))`$, the relation $`((p_1\rho )^{}K_M,C^{})0`$ and $`((fp_1\rho )^{}H,C^{})=0`$ holds. Thus for every fiber $`F`$ of $`\mu `$, we have the relation $`((p_1\rho )^{}K_M,F)0`$ and $`((fp_1\rho )^{}H,F)=0`$, because all fibers of the flat morphism $`\mu `$ represent the same homology class in $`H_2(W,𝐙)`$.
Let $`W^{}`$ be the image of the morphism $`(p_1\rho ,\mu ):WM\times T`$. We note that $`W^{}`$ ($`M\times T`$) is birationally equivalent to $`W`$, because the image of the morphism $`(id.,\nu )(p_1\rho ,\mu ):WM\times TM\times p_2(V)`$ is $`V`$. Consider the projection morphisms $`p_1^{}:W^{}M`$ and $`p_2^{}:W^{}T`$. Then $`p_1^{}(W^{})=M`$.
$$\begin{array}{c}W\\ & & \\ W^{}& \stackrel{\text{embedding}}{}& M\times T& \stackrel{\text{2nd proj.}}{}& T\\ & & \text{1st proj.}& & \\ & & M\\ & & f& & \\ & & S\end{array}$$
(2.2)
We have some open dense subset $`U`$ of $`T`$ such that every fiber $`G`$ of $`p_2^{}`$ over $`U`$ is irreducible. By applying the projection formula to the birational morphism $`WW^{}`$, the relation that $`((p_1^{})^{}K_M,G)0`$ and $`((fp_1^{})^{}H,G)=0`$ follows. Note that $`G`$ can be identified with a curve on $`M`$ that is $`f`$-exceptional and whose intersection number with $`K_M`$ is not positive. The property of constructible sets implies that $`p_1^{}((p_2^{})^1(U))`$ and $`fp_1^{}((p_2^{})^1(U))`$ contain open dense subsets of $`M`$ and $`S`$, respectively.
Consequently the statement (i) holds.
Next we treat the case where $`p_1(V)M`$ for every $`VI`$. The countability of components of the Hilbert scheme $`Hilb(M)`$ implies that the statement (ii) holds. ∎
We cite Tsuji’s existence theorem of numerically trivial fibrations.
###### Proposition 2.2 (Tsuji , cf. ).
Let $`X`$ be a normal projective variety and $`L`$ a nef divisor on $`X`$. Then there exist projective varieties $`Y`$ and $`Z`$ and morphisms $`\mu :YX`$ and $`\nu :YZ`$ with the following properties:
(i) the morphism $`\mu `$ is birational
(ii) the morphism $`\nu `$ is surjective
(iii) the extension $`RatY/RatZ`$ is algebraically closed
(iv) for some two open subvarieties $`U`$ and $`V`$ of $`X`$ and $`Z`$ respectively, $`\mu ^1(U)=\nu ^1(V)`$ and $`\mu |_{\mu ^1(U)}`$ is isomorphic
(v) the divisor $`\mu ^{}L|_F`$ is numerically trivial for every very general fiber $`F`$ of $`\nu `$
(vi) for every very general point $`xX`$ there does not exist a closed subvariety $`Sx`$ such that $`L|_S`$ is numerically trivial and that $`dimS>dimYdimZ`$.
The following means that the property for a canonical divisor to be almost nup is stable under birational transformation.
###### Theorem 2.3.
Assume that the minimal model conjecture holds in dimension $`n`$. Let $`X`$ be a projective variety with only terminal singularities of dimension $`n`$. Then the canonical divisor $`K_X`$ is almost nup if and only if $`X`$ is of almost general type.
###### Proof.
The “only if” part. Trivial from the definition of “being of almost general type”.
The “if” part. From Miyaoka-Mori (), $`X`$ is not uniruled. Thus a crepant $`𝐐`$-factorialization of $`X`$ has a minimal model $`Z`$. Here we have a common resolution $`\mu :YX`$ and $`\nu :YZ`$ such that $`\mu ^{}K_X=\nu ^{}K_Z+E`$ where $`E`$ is a $`\nu `$-exceptional effective $`𝐐`$-divisor. Suppose that $`K_Z`$ is not quasi-nup. Then for some open subvariety $`U`$ of $`Z`$ there exists a proper dominating morphism $`\rho `$ from $`U`$ to a lower-dimensional variety $`V`$ such that the extension $`RatU/RatV`$ is algebraically closed and that $`K_Z|_F`$ is numerically trivial for every very general fiber $`F`$ of $`\rho `$, from Proposition 2.2 (Tsuji). For every desingularization $`\alpha :MZ`$, the divisor $`(K_M\alpha ^{}K_Z)|_{\alpha ^1(F)}`$ is $`\alpha |_{\alpha ^1(F)}`$-exceptional. Thus $`(K_M\alpha ^{}K_Z)\alpha ^{}(H)^{dimF1}|_{\alpha ^1(F)}=0`$ where $`H`$ is a hyperplane section of $`Z`$. So $`K_M\alpha ^{}(H)^{dimF1}|_{\alpha ^1(F)}=0`$. Therefore $`K_M`$ is not almost nup and thus $`X`$ is not of almost general type. This is a contradiction! Consequently $`K_Z`$ is quasi-nup and thus $`\mu ^{}K_X`$ is almost nup. As a result, $`K_X`$ is almost nup. ∎
Here we have the main result of this section.
###### Theorem 2.4.
Assume that the minimal model conjecture holds in dimension $`n`$. Let $`f:MS`$ be a surjective morphism between (possibly singular) projective varieties with relative dimension $`n`$ such that the extension $`RatM/RatS`$ is algebraically closed. Then one of the following holds:
(i) Every general fiber $`F`$ of $`f`$ is not of almost general type.
(ii) Every very general fiber $`F`$ of $`f`$ is of almost general type.
###### Proof.
Proposition 2.1 and Theorem 2.3 imply the assertion. ∎
## 3 Some kind of hyperbolicity
In this section, we show that the varieties of almost general type are characterized by some kind of hyperbolicity.
###### Theorem 3.1.
Assume that the minimal model conjecture holds in dimension $`<n`$. Let $`X`$ be a projective variety with only terminal singularities of dimension $`n`$. If $`K_X`$ is almost nup, then the locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type $`\}`$ is covered by at most countably many prime divisors on $`X`$.
###### Proof.
The proof proceeds along the same line as in the paper , by using Theorem 2.4. But, for the readers’ convenience, we do not make the presentation rough.
Assuming that $`K_X`$ is almost nup and that however the locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type $`\}`$ cannot be covered by at most countably many prime divisors on $`X`$, we derive a contradiction.
Let $`X\times Hilb(X)`$ be the universal family parametrized by the Hilbert scheme $`Hilb(X)`$. By the countability of the components of $`Hilb(X)`$, we have an irreducible component $`V`$ of $``$ with surjective projection morphisms $`f:VX`$ and $`g:VT(Hilb(X))`$ from $`V`$ to projective varieties $`X`$ and $`T`$ respectively, such that $`f(g^1(t))X`$ for every $`tT`$ and that the locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type and $`D=f(g^1(t))`$ for some $`tT\}`$ cannot be covered by at most countably many prime divisors on $`X`$. Let $`\nu :V_{norm}V`$ be the normalization. We consider the Stein factorization of $`g\nu `$ into the finite morphism $`g_1:ST`$ from a projective normal variety $`S`$ and the morphism $`g_2:V_{norm}S`$ with an algebraically closed extension $`RatV/RatS`$. Put $`V^{}:=`$ \[the image of the morphism $`(f\nu ,g_2):V_{norm}X\times S`$\].
$$\begin{array}{ccccccc}V_{norm}& & V^{}& \stackrel{\text{embedding}}{}& X\times S& & S\\ & & & & & & g_1& & \\ & & V& \stackrel{\text{embedding}}{}& X\times T& & T\\ & & & & & & \\ & & & & X\end{array}$$
(3.1)
Note that every fiber of the morphism $`g:VT`$ consists of a finite number of fibers of the projection morphism from $`V^{}`$ to $`S`$. Thus we may replace $`(V,T)`$ by $`(V^{},S)`$ and assume that the extension $`RatV/RatT`$ is algebraically closed.
According to Theorem 2.4 (i) and (ii), we divide the situation into two cases.
First consider the case where $`g^1(t)`$ is of almost general type for very general $`tT`$. Then there exists a subvariety $`T_0T`$ such that the locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type and $`D=f(g^1(t))`$ for some $`tT_0\}`$ cannot be covered by at most countably many prime divisors on $`X`$. Thus we can replace $`(V,T)`$ by $`(V_1,T_1)`$, where $`V_1`$ and $`T_1`$ are projective varieties such that $`V_1`$ is some suitable irreducible component of $`g^1(T_0)`$ and $`T_1=g(V_1)`$ . Because $`dimV_1<dimV`$, by repeating this process of replacement, we can reduce the assertion to the next case.
Now consider the case where $`g^1(t)`$ is not of almost general type for general $`tT`$. If $`dimV>dimX`$, then $`dimf^1(x)1`$ for all $`xX`$, thus $`dimg(f^1(x))1`$ (this means that $`g(f^1(x))H\mathrm{}`$ for every hyperplane section $`H`$ of $`T`$). Therefore, by repeating the process of cutting $`T`$ by general hyperplanes, we can reduce the assertion to the subcase where $`dimV=dimX`$. Take birational morphisms $`\alpha :X^{}X`$ and $`\beta :V^{}V`$ from non-singular projective varieties with a generically finite morphism $`\gamma :V^{}X^{}`$ such that $`\alpha \gamma =f\beta `$.
$$\begin{array}{ccc}V& \stackrel{\beta }{}& V^{}\\ f& & \gamma & & \\ X& \underset{\alpha }{}& X^{}\end{array}$$
(3.2)
Because $`K_X`$ is almost nup and $`K_X^{}\alpha ^{}K_X`$ is effective, $`K_X^{}`$ is almost nup. Here we have $`K_V^{}=\gamma ^{}K_X^{}+R_\gamma `$, where $`R_\gamma `$ is the ramification divisor (which is effective) for $`\gamma `$. Thus $`K_V^{}`$ becomes almost nup. For a very general fiber $`F^{}=(g\beta )^1(t)`$ for $`g\beta `$ (i.e. the point $`t`$ does not belong to some fixed union of at most countably many prime divisors on $`T`$), also $`K_F^{}`$ is almost nup. Thus every very general fiber $`F`$ of $`g`$ is of almost general type. This is a contradiction! ∎
As a corollary we have
###### Theorem 3.2.
Let $`X`$ be a projective variety with only terminal singularities of dimension $`5`$. If $`K_X`$ is almost nup, then the locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type $`\}`$ is covered by at most countably many prime divisors on $`X`$.
###### Proof.
Theorem 3.1 implies the assertion by virtue of the existence theorem (Shokurov , cf. Hacon and McKernan ) and the termination theorem (Kawamata-Matsuda-Matsuki ) of four-dimensional flips. ∎
We cite the following proposition and proof for the convenience of the reader.
###### Proposition 3.3 ().
Let $`X`$ be a projective variety with only terminal singularities. If $`K_X`$ is almost nup, then $`X`$ is not birationally equivalent to any minimal variety with numerically trivial canonical divisor.
###### Proof ().
We give a sketchy proof. Let $`Z`$ be a minimal variety with numerically trivial canonical divisor $`K_Z`$. For any resolution $`\nu :YZ`$, the divisor $`K_Y\nu ^{}K_Z`$ is effective and $`\nu `$-exceptional. Thus $`K_Y(\nu ^{}H)^{dimY1}=0`$ for any ample divisor $`H`$ on $`Z`$. Consequently $`K_Y`$ is not almost nup. ∎
Now we consider the converse statement of Theorem 3.1 and Proposition 3.3.
###### Theorem 3.4.
Assume that the minimal model conjecture holds in dimension $`n`$. Let $`X`$ be a projective variety with only terminal singularities of dimension $`n`$. If $`X`$ is neither uniruled nor birationally equivalent to any variety fibered by minimal varieties with numerically trivial canonical divisors, then $`K_X`$ is almost nup.
###### Proof.
A crepant $`𝐐`$-factorialization of $`X`$ has a minimal model $`Z`$. Suppose that $`K_X`$ is not almost nup. From the inequality $`K_XK_Z`$, between bi-divisors (), $`K_Z`$ is not quasi-nup. Thus from Proposition 2.2 (Tsuji), $`Z`$ is birationally equivalent to some variety fibered by minimal varieties with numerically trivial canonical divisors. This is a contradiction! Consequently $`K_X`$ is almost nup. ∎
At last we establish the following
###### Theorem 3.5 (Main Theorem).
Assume that the minimal model conjecture holds in dimension $`n`$. Let $`X`$ be a projective variety with only terminal singularities of dimension $`n`$. Then the five conditions below are equivalent to each other:
(1) $`X`$ is of almost general type.
(2) $`K_X`$ is almost nup.
(3) The locus $`\{D;D`$ is a closed subvariety ($`X`$) not of almost general type $`\}`$ is covered by at most countably many prime divisors on $`X`$ and the variety $`X`$ is not birationally equivalent to any minimal variety with numerically trivial canonical divisor and is not a rational curve.
(4) $`X`$ is not uniruled and can not be covered by any family of varieties being birationally equivalent to minimal varieties with numerically trivial canonical divisors.
(5) $`X`$ is neither uniruled nor birationally equivalent to any variety fibered by minimal varieties with numerically trivial canonical divisors.
###### Proof.
(1) and (2) are equivalent (Theorem 2.3). (2) implies (3) by Theorem 3.1 and Proposition 3.3. (3) implies (4), because every minimal variety with numerically trivial canonical divisor is not of almost general type from Theorem 2.3. Obviously (4) implies (5). (5) implies (2) by Theorem 3.4. ∎
###### Remark 3.6 (cf. ).
The assumption of Main Theorem 3.5 is satisfied for $`n4`$, because of the existence (Shokurov , cf. Hacon and McKernan ) and the termination (Kawamata-Matsuda-Matsuki ) of four-dimensional flips.
## Appendix
Here we note the theory of pseudo-effective divisors due to Boucksom-Demailly-Paun-Peternell (). Let $`L`$ be a Cartier divisor on a smooth projective variety $`X`$. The divisor $`L`$ is pseudo-effective if and only if $`(L,C)0`$ for every curve $`C`$ through a very general point $`x`$ on $`X`$ (\[4, Theorem 0.2\]). For the pseudo-effective divisor $`L`$, we set the number $`p(L):=dimX`$ (the notation at \[4, Definition 8.3\]) if and only if $`L`$ is almost nup.
If $`X`$ is of dimension $`4`$, the canonical divisor $`K_X`$ is pseudo-effective and $`p(K_X)=4`$, then every proper subvariety $`SX`$ through a very general point $`x`$ on $`X`$ is of general type. (The paper overlaps this result \[4, Proposition 9.12\]. The presentation of the proof in the latter paper is casual.)
Faculty of Education, Gifu Shotoku Gakuen University
Yanaizu-cho, Gifu City, Gifu 501-6194, Japan
fukuda@ha.shotoku.ac.jp
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# Radiative Transfer and Acceleration in Magnetocentrifugal Winds
## 1. Introduction
A variety of observational signatures point to the importance of outflowing gas within many types of Active Galactic Nuclei (AGNs). Blueshifted absorption features (in Broad Absorption Line Quasars, or BALQSOs; see, e.g., Weymann et al., 1991) are seen in approximately 15% (Reichard et al., 2003) of radio-quiet quasars, with velocities up to 0.1c. In addition, radio-loud quasars display relativistic, collimated outflows. More recently, both UV and X-ray absorbing gas have been observed in approximately 60% of Seyfert 1 galaxies (Crenshaw et al., 1999). Observational estimates hint that the mass outflow rate is nearly equal to the mass inflow rate (for a review of mass outflow in AGNs, see Crenshaw, Kraemer, & George, 2003).
The development of models to explain the mass outflow rates, geometry, and general kinematics of these winds has proven difficult, but progress on a number of possibilities has been encouraging. From very early on, researchers examined both radiative wind models (e.g., Drew & Boksenberg, 1984; Vitello & Shlosman, 1988) and, to explain radio jets, hydromagnetic models (e.g., Blandford & Payne, 1982, hereafter BP82). Later models of radiatively-driven winds were able to explain the BALQSO outflow velocities and population fraction (Murray et al., 1995, hereafter, MCGV95) as well as the single-peaked emission lines (Murray & Chiang, 1997). The line-driven models were also demonstrated in hydrodynamical simulations (Proga, Stone & Drew, 1998; Proga, Stone, & Kallman, 2000; Pereyra et al., 2004), which elucidated the density structure, geometry, and kinematics of the flow in two dimensions. In addition, models with combinations of continuum- and line-driving by X-rays were presented by Chelouche & Netzer (2001). One of the concerns with line-driving in AGNs has been the possible over-ionization of the gas by X-rays in the central continuum (leaving the wind with too few lines to intercept flux in the lines): the need for “shielding gas” to prevent this overionization has been a persistent concern (e.g., MCGV95, Chelouche & Netzer, 2003b; Proga & Kallman, 2004).
In addition, hydromagnetic winds have also been developed to gain insight into the observations; some of these models have included radiative driving. In these magnetocentrifugal winds, gas is loaded onto and tied to large-scale magnetic field lines; those field lines then fling the matter centrifugally away from the disk, like beads along a wire. Magnetocentrifugal acceleration is commonly used to explain large-scale collimated outflows in radio galaxies (BP82; for reviews of magnetocentrifugal driving, see Spruit, 1996; Königl & Pudritz, 2000; Ferreira, 2003). In the context of the Broad Emission Line Region (BLR), magnetocentrifugal outflows made up of clouds were first called upon to explain the single-peaked broad emission lines, outflow velocities, and densities (Emmering et al., 1992) . The “torus” (obscuring gas that plays a central role in the Unification model, yielding a dependence of observed properties with inclination angle; see Antonucci, 1993) has been explained as a dusty, continuous magnetic wind (Königl & Kartje, 1994, hereafter KK94) where radiative acceleration on dust affects the wind geometry. Hydromagnetic disk winds with radiation pressure were also developed by de Kool & Begelman (1995) as an alternative explanation for the population fraction of BALQSOs and to understand cloud confinement within the outflow. In addition, another magnetocentrifugal wind model was developed to explain single peaked emission lines arising from a distribution of clouds (Bottorff et al., 1997), as well as the dynamics of the warm absorber in NGC 5548 (Bottorff, Korista, & Shlosman, 2000). Finally, some effects of magnetic fields (not including magnetocentrifugal driving as in BP82) have also been considered in two-dimensional hydrodynamic simulations (Proga, 2000, 2003).
At the present time, both radiatively driven winds and magnetocentrifugally driven winds (with radiative acceleration added in some models) offer compelling but competing pictures of the key physics in the cores of AGNs. No clear observational evidence yet discriminates the dominant physics of wind-launching in AGNs. For radiatively-driven winds, three main lines of evidence hint that radiative driving should be important in models of wind dynamics. First, Laor & Brandt (2002) find a correlation between UV luminosity and the observed outflow velocity, which agrees with a basic prediction of line-driven wind models (Proga, 1999). In addition, the “ghost of Ly-$`\alpha `$(Arav et al., 1995; Chelouche & Netzer, 2003b), as well as the realization that the radiative momentum removed from the continuum in blueshifted absorption lines in BALQSOs is a significant fraction of the total radiative momentum (MCGV95) are both also important clues that radiative acceleration is an important component to any self-consistent model of AGN winds. There are still, however, concerns about how the “shielding” of the wind works (Chelouche & Netzer, 2003b; Proga & Kallman, 2004). In addition, in the case of stellar disk winds, observations of two nova-like variables seem to show that their winds are not dominated by radiative driving (Hartley et al., 2002). (This conclusion, however, rests on the prediction that line equivalent widths are direct measures of mass outflow rate, which may not be the case; see Pereyra et al., 2004)
On the other hand, the leading model for collimated radio jets in AGNs (BP82) already calls upon large scale, dynamically important magnetic fields, as do the hydromagnetic wind models (mentioned above, e.g. Königl & Kartje, 1994; Kartje, 1995; Bottorff et al., 1997) that have also had success in explaining AGN observations. In addition, such a hydromagnetic wind would have no difficulty with overionization, and so might naturally serve as a “shield” for radiative acceleration further from the central source. There are, however, currently no models which address the interplay between magnetocentrifugal driving and radiative acceleration; if such wind models could be constrained, we may be able to observationally distinguish the physics of wind launching in the cores of AGNs, and gain insight into the role that outflows play both in accretion and in feedback of those winds in the galaxy and surrounding matter.
This paper develops a detailed photoionization and dynamical model for magnetocentrifugal winds in AGNs. The model is designed to explore the radiative transfer within such magnetocentrifugal winds, but also to help understand how radiative driving impacts the kinematic structure of such winds. Constructing such a model builds the foundation for later work to determine absorption and emission line profiles in order to compare with observations and check for the presence of magnetocentrifugal winds within AGNs. An overview of this model is presented in §2; the model is then defined in detail in §3. An examination of the structure of a particular “fiducial” magnetocentrifugal, radiatively-accelerated wind is described in §4 and then the dependences of radiative acceleration on some initial parameters are shown in §5. Conclusions and directions for future work are summarized in §6.
## 2. Model Overview
Before examining the model in detail, it is instructive to present a summary of the basic design. The semianalytic model developed here includes magnetic acceleration and radiative acceleration of a continuous wind launched from an accretion disk. A detailed treatment of radiative transfer is included by using Version 96.00 of the photoionization simulation program Cloudy, last described by Ferland et al. (1998). These elements are introduced in an approximate schematic of the wind model shown in Figure 1, depicting a portion of the accretion disk and outflowing wind. In this figure, radiation (entering from the left side of the schematic) first encounters a purely magnetocentrifugally accelerated wind, which will be referred to as the “shield”, as it absorbs radiation from the central continuum. The shield is introduced as a separate component in order to cleanly differentiate the effect of shielding from radiative acceleration; radiative driving of the shield is therefore not considered in this work. Beyond that shield is an optically thin, radiatively and magnetically accelerated wind streamline (which we will henceforth refer to as the “wind”). In this portion of the model, both magnetocentrifugal and radiative forces help accelerate the flow off of the accretion disk; the magnetic fieldlines are shown by the black lines bordering the outflow. The included radiative acceleration is calculated by first simulating the photoionization within both the shield and the wind along radial paths such as the thick, black lines in the figure.
The separation of the wind into two components is, of course, artificial. In reality, radiative acceleration would gradually increase in importance for portions of the wind that are increasingly shielded. However, splitting the outflow into these two components allows a first-order, qualitative solution that can be used to gain some understanding of how magnetic and radiative forces might interact, and how a magnetic wind may be able to act as a radiative “shield” to allow more efficient radiative acceleration. While artificial, this method of using two wind components has already been used successfully to examine winds with magnetocentrifugal and radiative driving on dust (e.g., KK94, Kartje, 1995).
Figure 2 presents a schematic flow chart of how model calculations proceed. The wind starts as a self-similar magnetohydrodynamic model that yields the pure magnetocentrifugal wind solution (covered in §3.1). Next, simulations of the photoionization balance of that wind streamline (§3.2) are run, and the resultant ionization balance and transmitted continuum are used to calculate the radiative acceleration behind the shield (§3.3). Next, the radiative acceleration is input (as a function of polar angle, $`\theta `$) back into the self-similar magnetohydrodynamic model, modifying the structure of the wind streamline, while leaving the shield unaffected. This process is then repeated, simulating the photoionization of that modified wind and recalculating the radiative acceleration terms. We typically iterate five to eight times to converge to a final equilibrium solution.
With the basic model now summarized, it may be instructive to compare and contrast it to the recent wind model examined in Proga (2003), where the combination of magnetic and radiative forces in disk winds is also investigated. Proga (2003) concentrates on numerical simulations of time-dependent winds with line driving and magnetic forces. These numerical simulations allow large-scale models of outflows that are valuable in understanding global wind structures in many different astrophysical contexts. In contrast, the radiatively-driven components of this model are more localized, since radiative acceleration is applied to the shielded streamline only. This model setup has been chosen for its flexibility in radiative acceleration modeling; using Cloudy for such radiative acceleration models enables computations not only of radiative line driving but also of continuum driving, and allows us the freedom to include dust and easily vary the incident spectrum, for example. In addition, the magnetic winds produced in Proga (2003) are not magnetocentrifugal outflows (as in BP82), as are the semianalytic winds presented here. Further, these new models are steady-state, not time-dependent. Finally, semi-analytic, steady-state models allow an exploration of general behaviors through many parameter variations; large-scale numerical simulations can usually vary only a few parameters. In summary, these models cover different facets of the disk-wind problem, yielding different perspectives on a complicated system.
## 3. The Two-Phase Hydromagnetic and Radiative Wind Model
In this section, we describe in detail the model’s components, and derive key equations. The magnetocentrifugal wind model is introduced first (§3.1), followed by the photoionization simulations (§3.2), radiative acceleration calculation (§3.3), and finally a discussion of the wind model equation of motion (§3.4).
### 3.1. Magnetocentrifugal Self-Similar Wind Solution
To derive the equations governing the continuous magnetocentrifugal wind, we start with the equations of a stationary, axisymmetric magnetohydrodynamic (MHD) flow in cylindrical coordinates, make the assumption of self-similarity in the spherical radial coordinate, and utilize the continuity equation, conservation of angular momentum along the flow, and both the radial and vertical momentum equations, very much as in BP82 and KK94 (see Appendix A). Thermal effects are neglected in the wind, therefore effectively assuming that the wind starts out supersonic (this assumption is checked later on, see §3.5). In deriving the wind equations, we use the same simplifications as BP82, except for the added complication that energy is not conserved in the radiatively-accelerated system due to the constant input of radiative energy into the wind. In the original formulation of BP82, conservation of energy supplied an additional constraint which allowed a simplification of the equations of motion to two first-order differential equations. Because energy is not conserved in this flow, the equivalent of three first-order differential equations must be integrated, solving for three parameters simultaneously instead of two as in BP82. The detailed setup and derivation of this set of equations of motion are given in Appendix A.
The integration of the momentum equations starts by specifying the following initial parameters: the mass loading of the wind (the ratio of mass flux to magnetic flux in the magnetocentrifugal wind, $`\kappa \frac{4\pi \rho v_p}{B_p}`$, where $`\rho `$ is the mass density of the wind, $`v_p`$ is the poloidal velocity of the wind ($`v_p(v_r^2+v_z^2)^{1/2}`$), and $`B_p`$ is the poloidal magnetic field strength), the specific angular momentum of gas and field in the wind, and the power-law exponent ($`b`$) that describes the change in density with spherical radius: $`\rho R^b`$. Also input, as parameters, are the mass of the central black hole, $`M_{}`$, the wind’s launch radius on the disk, $`r_0`$, and the density at the base of the wind at the launch radius, $`n_0`$. The program employs a “shooting” algorithm (using the SLATEC routine DNSQ; Powell, 1970) to integrate from the singular point (the Alfvén point) to the disk, solving for the height of the singular point above the disk ($`\chi _A`$) and the slope of the streamline at both the disk and the Alfvén point ($`\xi _0^{}`$ and $`\xi _A^{}`$) by matching the integration results to boundary conditions on the disk. After solving for the position of the Alfvén point, the equations of motion are integrated from the disk to a user-specified height beyond the Alfvén point; along the streamline, the run of velocity, density, and magnetic field are calculated.
This code has been tested (without radiative acceleration) against the solution given in BP82 and have duplicated their results to within 8%. This is close to the previously reported 4% variance in the recalculation of BP82 reported in Safier (1993). The difference in these new results comes not only from using higher precision calculations compared to BP82 (we use the same precision as Safier, 1993), but in addition, a more complex set of equations is being solved.
### 3.2. Photoionization Simulations of the Wind
Next, Version 96.00 of the photoionization code Cloudy (Ferland et al., 1998) is used to simulate the absorption of the magnetocentrifugal shield and wind as well as the ionization state at the wind streamline. Photoionization simulations of the shield and wind are used to calculate the radiative acceleration and allow considerable flexibility in gas parameters (such as gas abundances, dust, central continuum, etc). This flexibility is gained at the cost of simulating the photoionization state of the wind as if it were a static medium, as Cloudy assumes; this is only true of the recombination timescale for the gas is much shorter than the transit timescale of gas in the region simulated ($`\tau _{\mathrm{recomb}}<\tau _{\mathrm{flow}}`$). We have however, verified, a posteriori, that $`\tau _{\mathrm{recomb}}<\tau _{\mathrm{flow}}`$ for all of the radiatively accelerated ions at the base of the wind where radiative acceleration is important, using the recombination rate approximations in Arnaud & Raymond (1992) and Verner & Ferland (1996) to calculate the recombination times. Even at high latitude, $`\tau _{\mathrm{recomb}}<\tau _{\mathrm{flow}}`$ for all of the significant ions. At the highest latitudes, highly ionized O, which has the longest recombination time among the significantly radiatively accelerated ions, has a recombination timescale that is still a factor of two less than the transit time. Meanwhile, highly ionized Fe ions, which dominate the low level of radiative acceleration at high latitudes, have recombination timescales an order of magnitude less than the transit timescale. Using Cloudy is therefore a reasonable approximation for the radiative equilibrium within our winds (especially since, with our code, adiabatic and advective effects are added, see §3.2.1). As Cloudy enables a flexible, self-consistent calculation of radiative acceleration, we accept this approximation to enable these calculations.
Of basic importance to the photoionization simulations is the illuminating spectral energy distribution (SED). For the purposes of these simulations, an SED adapted from Risaliti & Elvis (2004) is used (an example of the SED is shown in Figure 6). This SED is input into Cloudy via the generic “AGN” continuum with $`T_{\mathrm{blackbody}}=1.5\times 10^5`$ K, $`\alpha _{\mathrm{ox}}=1.43`$ (Elvis, Risaliti, & Zamorani, 2002), $`\alpha _{\mathrm{UV}}=0.44`$, and $`\alpha _X=0.9`$ ($`\alpha _{\mathrm{ox}}`$ defines a single power-law that would describe the continuum between $`2500`$ Å and 2 keV, $`\alpha _{\mathrm{UV}}`$ is the slope of the low-energy component of the Big Blue Bump, and $`\alpha _\mathrm{X}`$ is the X-ray power-law exponent; our value of $`\alpha _{\mathrm{ox}}`$ is taken from the middle of the range 0.8 to 1.0 given in Risaliti & Elvis (2004)).
Spectral signatures and radiative acceleration also depend, of course, on the column density in the wind. Observations yield only rough constraints for this, so these columns as left as free parameters; the effect of varying these columns will be investigated in this paper. The columns throughout the shield and wind are set by the columns at the base of the shield and wind, denoted $`N_{\mathrm{H},0}`$ for the hydrogen column at the base of the wind. As the wind rises above the disk and accelerates, that column density ($`N_\mathrm{H}`$) drops as a function of height due to mass conservation (an example of this is shown later in Fig. 7). Investigating the shielding ability of such a dynamic shielding column is of central interest to this paper, and will be addressed in §5.1.
#### 3.2.1 Photoionization of the Continuous Wind
As depicted in Figure 1, Cloudy simulates the photoionization of the wind along radial sight lines through the shield and through the wind, ending at the site of radiative acceleration on the wind streamline. The continuum incident on the wind streamline is also calculated. The photoionization state and continuum at the end of the Cloudy calculations are recorded, and then used to compute the radiative acceleration of that gas. Finally, the acceleration is tabulated and applied as a function of $`\theta `$ along the wind streamline by inputing the angle-dependent radiative acceleration into the gravitational term (this is covered in more detail in Appendix A, see eqn. A3).
Whereas Cloudy is designed to simulate the photoionization balance of gases as in the shield and wind, it cannot easily incorporate adiabatic and advection effects: Cloudy has no knowledge of the particular velocity profile of the wind in the overarching model, or the temperature difference between successive photoionization models as the wind climbs above the disk. Therefore, this model calculates both advective heating and adiabatic cooling in the wind, and adds those terms manually into Cloudy’s simulations. Both terms largely cancel in the wind, and have only a negligible effect on outflow dynamics, but they are included in all of the models for completeness.
### 3.3. Radiative Acceleration Calculations
The model then incorporates the above-mentioned results for the ionization structure and radiation field to calculate the radiative forces felt by the wind. There are two different kinds of radiative acceleration to consider: continuum acceleration (including radiative acceleration on dust) and line acceleration. It is convenient to express the radiative acceleration in terms of $`\mathrm{\Gamma }(\theta )`$,
$`\mathrm{\Gamma }(\theta )`$ $``$ $`{\displaystyle \frac{a_{\mathrm{radiative}}(\theta )}{g}},`$ (1)
where $`a_{\mathrm{radiative}}`$ is the acceleration due to radiation, and $`g`$ is the local gravitational acceleration.
#### 3.3.1 Line and Continuum Acceleration
In general, for continuum and line acceleration, the radiative acceleration is given by
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{\frac{n_e\sigma _TF}{\rho c}(M_{\mathrm{cont}}+M_{\mathrm{lines}})}{\frac{GM_{}}{r^2+z^2}}},`$ (2)
where $`F`$ is the local flux (the flux transmitted through both the shield and wind column), $`n_e`$ is the electron density, $`\rho `$ is the gas density, $`c`$ is the speed of light, $`G`$ is the gravitational constant, $`M_{}`$ is the mass of the central black hole, $`r`$ and $`z`$ are cylindrical coordinates centered on the black hole, and $`M_{\mathrm{cont}}`$ & $`M_{\mathrm{lines}}`$ are the “force multipliers” that relate how much the radiative forces on the gas (on the line and continuum opacity, respectively) exceed the radiative forces on electrons alone. They are given below in terms of the continuum opacity, $`\chi _\nu `$, and the line opacity, $`\chi _l`$, for the continuum and lines, respectively:
$`M_{\mathrm{cont}}`$ $`=`$ $`{\displaystyle \frac{1}{n_e\sigma _TF}}{\displaystyle \chi _\nu F_\nu 𝑑\nu },`$ (3)
$`M_{\mathrm{lines}}`$ $`=`$ $`{\displaystyle \frac{1}{F}}{\displaystyle \underset{l}{}}F_l\mathrm{\Delta }\nu _l{\displaystyle \frac{1e^{\eta _lt}}{t}},`$ (4)
with
$$\eta _l\frac{\chi _l}{\sigma _Tn_e}t\frac{\sigma _Tn_ev_{\mathrm{th}}}{dv_R/dR},$$
(5)
where $`\nu `$ is the photon frequency, $`F_l`$ is the local (transmitted) flux in the line at the frequency of line number $`l`$, $`v_{\mathrm{th}}`$ is the sound speed in the gas, $`\mathrm{\Delta }\nu _l=\nu v_{\mathrm{th}}/c`$ is the thermal line width, and $`\eta _l`$ compares the opacity of the line (for a given ionization state of the gas) to the electron opacity, representing all of the atomic physics in the radiative acceleration calculation. The last remaining variable, $`t`$, is often called the “effective electron optical depth” and encodes the dynamical information of the wind in the radiative acceleration calculation. This dynamical information is important because in an accelerating medium, one must also account for the Doppler shift of the atomic line absorption energy in the accelerating gas relative to the emitted line photon’s energy: beyond the Sobolev length, $`v_{\mathrm{th}}/(dv_R/dR)`$, included in $`t`$, a line photon will be Doppler-shifted out of the thermal width of the absorption line and can escape the gas (see Sobolev, 1958; Castor, Abbott, & Klein, 1975; Mihalas & Weibel-Mihalas, 1999).
The above-mentioned force multipliers are calculated using the resonance line data of Verner et al. (1996) with solar abundances. Since $`M_{\mathrm{cont}}`$, the continuum multiplier, depends only on the ionization state, it is tabulated solely as a function of height in the wind. In contrast to $`M_{\mathrm{cont}}`$, the line multiplier ($`M_{\mathrm{line}}`$) is tabulated for a range of values of the parameter $`t`$. Later, when calculating the equation of motion for the wind, the local velocity gradient is used to compute the actual value of $`t`$, which is then used to linearly interpolate the table of $`M_{\mathrm{line}}`$ and then evaluate the radiative acceleration.
The force multiplier computation has been tested against Arav et al. (1994), who also calculated radiative acceleration from photoionization simulations. Figure 3 compares these new calculations results against their fits (noting that there is a typo in their eq. \[2.9\]; Z.-Y. Li, personal communication), where the force multipliers as a function of the ionization parameter $`U`$ is presented ($`U`$ is the ratio of hydrogen-ionizing photon density to hydrogen number density $`n`$, given by $`UQ/4\pi nR^2c`$, where $`Q`$ is the number of incident hydrogen-ionizing photons per second, and $`R`$ is the distance from the continuum source). Overall, good agreement is found, especially considering that Arav et al. (1994) point out that their fit deviates from their calculations at low values of $`U`$. The increase in the newly-calculated continuum force multiplier over Arav et al. (1994) is most likely due to the different continuum opacity database included in Cloudy 96 compared to the code (MAPPINGS) that was used in Arav et al. (1994). The multiplier values and trends with ionization parameter are still clearly very similar, however.
#### 3.3.2 Non-Sobolev Effects
Simply using the Sobolev length, $`v_{\mathrm{th}}/(dv_R/dR)`$, in Equation 4 can be misleading. In early simulations, we found that this Sobolev length could, in regions where the gas is slowly accelerating, be much larger than the physical size of the shield and wind combined. This is clearly not physical, so a simple non-Sobolev method was employed to calculate the size of the combined shield and wind column.
First, for the wind, the length of the absorbing column is simply limited to the minimum of the wind’s Sobolev length and its true (physical) length.
Second, for the shield, the length of the absorbing column is given by the minimum of the shield’s Sobolev length<sup>1</sup><sup>1</sup>1The Sobolev length calculation for the shield does take into account the offset in velocity between the wind and shield, which is important since the wind is radiatively accelerated in addition to the magnetic acceleration that the wind and shield share., its true physical length, and the average length of the column that absorbs photons at the wavelength of the dominant accelerating atomic lines. This last length-scale requires more explanation. To calculate this length, the top 20 line force multipliers (for individual atomic lines) are found for each polar angle $`\theta `$, and for each of those lines, the column of each particular ion in the shield is read from the Cloudy simulations. An average shielding column for all of the high-opacity lines is then calculated. To test this approximation, we have used lists of the top 10, 20, and 50 transitions in the wind and used them to calculate the limiting column in the shield. Changing the number of lines included does not significantly change the final wind solutions, especially since this effect is only important at the extreme base of the wind. Without considering all of these constraints on the size of the absorbing column, the Sobolev length can significantly overestimate the optical depth in the shield (even overestimating the physical length of the columns, for small accelerations), which results in the line driving acceleration dropping below the continuum acceleration. Other non-Sobolev effects (such as line blanketing within the shielding gas or the wind) are not considered in this model; for a consideration of these effects, see Chelouche & Netzer (2003a).
With a length-scale found from the sum of the wind and shield lengths found above, that length-scale is then substituted for the Sobolev length in the calculation of $`M_{\mathrm{lines}}`$ in Eq. 4. Including this physical length of the wind and shield will introduce dependences on the sizes of the wind and shield columns, which will be examined in §5.
### 3.4. Integrating the Euler Equation for the Wind
Given the results from the ionization and radiative acceleration calculations, the next step is to solve for the effect of radiative acceleration on the wind, taking the magnetocentrifugal wind model already computed and augmenting its acceleration with radiative forces. To do this, the full equation of motion (the Euler equation) is integrated along the streamline of the self-similar wind solution, recording $`\mathrm{\Gamma }(\theta )`$, the radiative acceleration. $`\mathrm{\Gamma }(\theta )`$ is then input into the self-similar model in the subsequent iteration (see Appendix A; specifically, eqn. A3).
In its simplest form, Euler’s Equation is given by
$`\rho \left({\displaystyle \frac{𝐯}{t}}+(𝐯)𝐯\right)={\displaystyle 𝐅_i},`$ (6)
where the $`𝐅_i`$ represent the various forces included.
As we are interested in steady-state winds, $`\frac{𝐯}{t}`$ in Eq. 6 is set to zero. As already mentioned, gravitational, radiation, and Lorentz forces are included, while thermal terms are not. This yields the expression below, with the magnetic force term is split into pressure and tension components.
$`\rho (𝐯)𝐯`$ $`=`$ $`[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM\rho }{R^2}}\widehat{R}{\displaystyle \frac{1}{8\pi }}B^2+{\displaystyle \frac{1}{4\pi }}(𝐁)𝐁`$ (7)
For these calculations, it is more intuitive to integrate the equation of motion along the flow already given by the magnetocentrifugal wind solution. Therefore, taking the dot product of Euler’s equation with $`\widehat{s}`$, which is defined as the direction along the flow, and expanding and simplifying the left-hand side of the equation, one finds:
$`(𝐯𝐯)\widehat{s}`$ $`=`$ $`v_p{\displaystyle \frac{v_p}{s}}{\displaystyle \frac{v_\varphi ^2}{r}}\mathrm{sin}\theta _F,`$ (8)
where
$`\theta _F`$ $``$ $`\mathrm{tan}^1\left({\displaystyle \frac{dr}{dz}}\right).`$ (9)
In the same way, if $`\theta `$ is defined via
$`\theta \mathrm{tan}^1\left({\displaystyle \frac{r}{z}}\right),`$ (10)
the gravitational term can be written as
$`[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM}{R^2}}\widehat{R}\widehat{s}`$ $`=`$ $`[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM}{r^2+z^2}}\mathrm{cos}(\theta \theta _F).`$ (11)
Next, we take the dot product of $`\widehat{s}`$ with the magnetic terms to find
$`\left[{\displaystyle \frac{1}{8\pi \rho }}B^2+{\displaystyle \frac{1}{4\pi \rho }}(𝐁)𝐁\right]\widehat{s}`$ $`=`$ $`{\displaystyle \frac{B_\varphi }{4\pi \rho r}}{\displaystyle \frac{(rB_\varphi )}{s}}.`$ (12)
Combining all of those terms, the full Euler Equation is obtained:
$`v_p{\displaystyle \frac{v_p}{s}}{\displaystyle \frac{v_\varphi ^2}{r}}\mathrm{sin}\theta _F`$ $`=`$ $`[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM}{(r^2+z^2)}}\mathrm{cos}(\theta \theta _F)`$ (13)
$`{\displaystyle \frac{B_\varphi }{4\pi \rho r}}{\displaystyle \frac{(rB_\varphi )}{s}}.`$
This equation is still dependent on $`v_\varphi `$, however, which can be eliminated by appealing to the induction equation and $`E_\varphi =0`$ for such axisymmetric systems (see, e.g., Königl & Pudritz, 2000), to yield a relation between $`v_p`$ and $`v_\varphi `$:
$`v_\varphi `$ $`=`$ $`{\displaystyle \frac{v_pB_\varphi }{B_p}}+\mathrm{\Omega }r.`$ (14)
Substituting this expression into the Euler Equation yields:
$`v_p{\displaystyle \frac{v_p}{s}}\left({\displaystyle \frac{v_pB_\varphi }{B_p}}+\mathrm{\Omega }r\right)^2{\displaystyle \frac{\mathrm{sin}\theta _F}{r}}=`$ (15)
$`[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM}{(r^2+z^2)}}\mathrm{cos}(\theta \theta _F){\displaystyle \frac{B_\varphi }{4\pi \rho r}}{\displaystyle \frac{(rB_\varphi )}{s}}.`$
To evaluate the effective optical depth $`t`$, an expression is required for $`dv_R/dR`$, the spherical radial gradient of the spherical radial velocity. Since the code integrates quantities only parallel to the flow, approximations to the perpendicular velocity gradients must be used. Assuming that the derivatives of $`\theta `$ and $`\theta _F`$ with distance along the streamline are small (verified a posteriori to be true):
$`{\displaystyle \frac{dv_R}{dR}}`$ $`=`$ $`\widehat{R}v_R,`$ (16)
$`=`$ $`\widehat{R}[v_p\mathrm{cos}(\theta \theta _F)],`$ (17)
$``$ $`\mathrm{cos}(\theta \theta _F)\left({\displaystyle \frac{dv_p}{dr}}\mathrm{sin}\theta +{\displaystyle \frac{dv_p}{dz}}\mathrm{cos}\theta \right),`$ (18)
$``$ $`\mathrm{cos}(\theta \theta _F){\displaystyle \frac{dv_p}{ds}}(\mathrm{sin}\theta _F\mathrm{sin}\theta +\mathrm{cos}\theta _F\mathrm{cos}\theta ),`$ (19)
$`=`$ $`\mathrm{cos}^2(\theta \theta _F){\displaystyle \frac{dv_p}{ds}}.`$ (20)
This integration procedure has been tested with radiative acceleration turned off, where it reproduces the original self-similar velocity profile to within one part in $`10^5`$. With the radiative acceleration turned on, the entire code has repeatedly converged within approximately eight iterations to an equilibrium magnetic wind structure (see Fig. 11). These tests show that consistent solutions are found and that radiation pressure does indeed affect the shielded component of the wind.
#### 3.4.1 Critical Points
As with any steady-state wind dynamics problem, one must search for and consistently pass all critical points (e.g., Vlahakis et al., 2000). Critical points mark the location in the wind where the flow speed is equal to the speed of information propagation in the wind, and mark locations in the solutions where solutions branches, or roots, meet; in the case of radiatively accelerated winds, the location of the critical point is set by the information propagation speed of radiative-acoustic wave, or Abbott speed (see, e.g., Abbott, 1980; Mihalas & Weibel-Mihalas, 1999). In integrating the equation of motion for the radiatively-accelerated wind, this code searches for critical points by looking for multiple roots in the solution to the equation of motion. However, no critical points due to radiative acceleration are present in any of the wind solutions we have found (the magnetocentrifugal wind does, however, always pass through its own Alfvén critical point).
To check this result, we have duplicated the work of Feldmeier & Shlosman (1999, hereafter, FS99), verifying that for simple wind geometries and without magnetocentrifugal acceleration, the integration code does indeed encounter a radiative critical point as predicted and found by FS99. In particular, for the field geometries and forces used in FS99, we have found identical solutions to both their analytical and numerical calculations.
We then gradually add, to the FS99 model, new components that are present in our new calculations. When the centrifugal acceleration and the enforced corotation near the base of the wind are introduced into the framework of FS99, the centrifugal acceleration overwhelms the acceleration of the second root that was present in FS99. Therefore, only a single root is found, and with only a single root, a critical point cannot be present in our solutions. Another way to check this is to examine the limit of winds launched at very large angles to the accretion disk. Indeed, for those large angles, the radiative acceleration begins to dominate the centrifugal acceleration, and the critical point reappears. Therefore, for the geometry of these magnetocentrifugal winds, where the angle of the outflow to the accretion disk surface is less than 60, no radiative critical point is expected within these solutions: a radiative critical point will not be present when centrifugal acceleration is dominant.
### 3.5. Model Assumptions and Limitations
This model includes simplifying assumptions about the outflow in order to make these calculations possible. In this subsection, the assumptions and limitations of the model are summarized, as are the reasons for allowing those assumptions.
First and foremost among these assumptions is self-similarity. While enabling a relatively quick and flexible model that can be used to survey a wide variety of outflows, this assumption does impose constraints on the dynamics of the model. However, the assumption of self-similarity is essential to the magnetocentrifugal wind solution, as it simplifies the complicated MHD equations. Simultaneously, the radial self-similarity accommodates, very easily, the radial geometry for the photoionization simulations.
The above-mentioned Cloudy photoionization simulations assume a static medium, which is an approximation as well. As we show in §3.2, however, this approximation is valid for the these wind models.
Since the accretion disk is a boundary condition in these models, these winds are assumed to be loaded with matter from the disk. Accurate models of the accretion disk structure are beyond the scope of this paper, so it is assumed that the full mass outflow rate of the the wind is indeed input onto the magnetic fieldlines at the accretion disk surface. In addition, the matter that is loaded onto those fieldlines is assumed to flow supersonically, i.e., the gas has already passed the sonic point in the flow. This is done to simplify the magnetocentrifugal wind equations and retain the basic model as outlined in BP82; as such, the same asymptotic expansions near the disk (from BP82) are utilized here. This treatment of the wind has been checked in several different ways. First, the magnetic pressure in the wind is indeed greater than the thermal pressure throughout the entire wind. Also, the wind’s final velocities are much greater than the sound speed at the base of the wind, showing that thermal effects are negligible in determining the final wind velocities. Finally, the lowest speeds found in the wind model are of order 20% of the sound speed, and such low Mach numbers are found only very near the disk, at the base of the wind. Given the above evidence of the dominance of magnetic fields and radiative acceleration, the “cold-wind” approximation is valid for these calculations.
Also, in the calculation of the radiative acceleration of the wind, an approximation to the velocity gradient along spherical rays is required. This approximation is necessary because the Euler integration for the wind yields velocity gradients only along the flow (the poloidal velocity gradients), and thus the other components must be approximated geometrically (see §3.4).
The approximations and limitations outlined above do constrain the use of this wind model, but in making these compromises, a versatile tool can be developed to study the chosen geometries and forces.
## 4. Radiation Transfer Within a Fiducial Magnetocentrifugal Wind
For definitiveness, this model is first employed to examine the radiative transfer within one fiducial irradiated magnetocentrifugal wind. For the purpose of this paper, ‘fiducial’ is defined to indicate the parameters listed in Table 1. These parameters are not meant to represent a proposed model for any one particular AGN, but to define a starting point from which to examine the structure of these outflows as well as the dependences of the outflows on the model parameters. For instance, the shielding column of $`N_{\mathrm{H},\mathrm{shield},0}=10^{23}`$ cm<sup>-2</sup> (where $`N_{\mathrm{H},\mathrm{shield},0}`$ represents the value of $`N_{\mathrm{H},\mathrm{shield}}`$ at the base of the wind, i.e. just above the accretion disk surface) is chosen because it displays an amount of radiative acceleration between the extremes of the smaller and larger columns that will be tested. Similarly, the radiative wind column that is defined is again between the extremes of nearly optically thick $`N_{\mathrm{H},\mathrm{rad},0}=10^{23}`$ cm<sup>-2</sup> and very optically thin $`N_{\mathrm{H},\mathrm{rad},0}=10^{19}`$ cm<sup>-2</sup>. Studying such a model first will help bring into focus important issues concerning the interplay of dynamics and photoionization, and represents a foundation from which one can explore the parameter dependencies of the model.
This fiducial model was therefore run with the parameters given in Table 1, and after eight iterations, converged to the final wind structure. The results for the fiducial model are shown in Figures 4 though 11. In these figures, an overview of the equilibrium state of this model is presented. The results displayed in these plots are discussed in detail below.
First, Figure 4 shows the height of the poloidal streamline as a function of radius in units of the launching radius. In addition, to illustrate the small difference in geometry between the final and initial wind models, Figure 5 gives the fractional change in height as a function of radius. Both of these figures show that the wind still maintains a collimated state, achieving a height of $`z/r_0100`$ (where $`r_0`$ is the launching radius) at a cylindrical radius of only $`r/r_030`$. So, in this fiducial model, despite the input from radiative acceleration, the wind maintains this streamline with only small changes in the structure of the wind throughout all iterations (see Fig. 5). Thus, for the case of the fiducial model with $`L/L_{\mathrm{Edd}}=0.01`$, this added acceleration does not significantly affect the structure of the magnetocentrifugal outflow. These models do show changes in the velocity structure of the wind near the disk surface (as will be shown in Fig. 11), but the poloidal wind structure does not change significantly: on the scale of Figure 4, the streamlines of the initial, purely magnetocentrifugal streamline would lie on top of the streamline shown.
At logarithmically-spaced co-latitudinal angles along the streamline, Cloudy photoionization simulations are run to determine the photoionization state of the gas, as well as the radiative transfer through the shield and wind. Changes in the continuum transmitted through the shield are shown in Figure 6; the various plots show the simulated continuum at various heights in the shield corresponding to the indicated columns. As the shielding column decreases as a function of height above the disk, the shield transmits progressively more and more of the ionizing radiation. This plot also displays how rapidly the column drops as a function of height above the disk: $`N_{\mathrm{H},\mathrm{shield},0}=10^{23}`$ cm<sup>-2</sup> occurs at $`\theta =89.9^{}`$, $`N_{\mathrm{H},\mathrm{shield}}=10^{22.5}`$ cm<sup>-2</sup> at $`\theta =89.8^{}`$, $`N_{\mathrm{H},\mathrm{shield}}=10^{22}`$ cm<sup>-2</sup> at $`\theta =89.4^{}`$, and $`N_{\mathrm{H},\mathrm{shield}}`$ drops to $`10^{21}`$ cm<sup>-2</sup> at $`\theta =85.1^{}`$.
The flux transmitted through the shield then illuminates the radiatively accelerated wind; results from the photoionization simulations for the wind are presented in Figure 7, where the streamline, velocity, density, ionization parameter and temperature in the wind are plotted. In Figure 7a, the height of a wind streamline as a function of distance along the streamline (labeled $`s`$, given in units of the initial radius, $`r_0`$) is shown; this plot simply recasts the structure of the flowline shown in Figure 4 in terms of $`s`$ for comparison with the remaining plots.
In Figure 7b, velocities along the streamline are plotted, showing not only the rapid acceleration in the wind, but also comparing the components of the wind’s velocities. All velocities are given in units of the Keplerian velocity at the base of the wind, $`v_{\mathrm{k},0}`$ ($`v_{\mathrm{k},0}=6.65\times 10^3`$ km s<sup>-1</sup> for the parameter values in Table 1). This plot shows that the vertical velocity is dominant in these winds at large distances (again showing the wind is somewhat collimated), with the radial and azimuthal velocities becoming approximately equal far from the launching radius of the outflow. Near the very base of the disk, the radial velocity quickly dominates both the azimuthal and vertical velocities (qualitatively similar to the velocity structure calculated for radiatively-dominated flows as in MCGV95). Most importantly, though, we note the extraordinarily rapid acceleration of the gas from the disk; such acceleration is a hallmark of both magnetocentrifugal as well as radiatively-dominated winds, which usually accelerate to their terminal velocities in a distance on the order of their launching radius.
Due to mass conservation, the extremely rapid acceleration of the wind causes a sharp drop in both the number density and the column density with height above the disk: both the number density and column density immediately drop by three orders of magnitude as the wind rises above the disk. This is displayed in Figure 7c. This overall drop in density is extremely important for the ionization state of the wind, not only for the observational ramifications (i.e., what ions are present in various parts of the wind) but for the acceleration of the wind as well, as will be shown very shortly. Figures 8 and 9 show in more detail how the radiative acceleration leads to a substantial change (by approximately a factor of two) in both velocity and density with height near the disk surface. The changes in velocity for the pure magnetocentrifugal wind as compared to the magnetocentrifugal and radiatively-accelerated wind is shown in Figure 8. The difference in the density profile for the pure magnetocentrifugal wind as compared to the magnetocentrifugal and radiatively-accelerated wind is shown in Figure 9. (Note that all densities are normalized to their value at the base of the wind.)
The corresponding ionization state of the wind and temperature are shown in Figure 7d. Most striking is the dramatic rise in the ionization parameter as the wind rises above the disk, which is simply due to the drop in density and in shielding already mentioned. The ionization parameter is of prime interest, as the radiative acceleration in resonant lines is dependent on the number of atomic lines in the gas; the rapid ionization of the gas prompts questions about how efficient line-driving will be within this magnetocentrifugal wind. In addition, can a wind with such a dramatic drop in column density form an effective “shield”? And for these models, how does the radiative acceleration then compare to magnetic acceleration?
These questions are addressed in Figure 10. From the Cloudy simulations summarized in Figure 7, both the bound-free and bound-bound radiative acceleration are calculated; the resultant acceleration (compared to the local gravity) for the fiducial model is shown in Figure 10. This model shows that both line-driving and continuum-driving have important roles to play, with the line-driving dominating the continuum driving in the high-density, low-ionization part of the wind, and continuum driving dominating line-driving at larger distances, when the density is much lower. It is important to note that for the parameters of the fiducial model with $`L/L_{\mathrm{Edd}}=0.01`$, magnetic acceleration is still much greater than either line or continuum acceleration. Line-driving is greater than continuum-driving in only part of the outflow (although this can change with the density at the base of the wind, as will be shown in §5.2), as the ionization parameter is low enough only at the base of the wind for significant numbers of atomic lines to exist. In addition, radiative driving is not immediately important at the extreme base of the outflow, because of the low fluxes that penetrate the columns there. Meanwhile, the acceleration due to continuum-driving stays close to the Eddington ratio, at $`0.01`$. This value is reasonable, as most of the continuum acceleration comes from electron scattering, so that $`\mathrm{\Gamma }0.01`$ would be expected. Minor increases above that value very near the disk surface are due to bound-free transitions in the portion of the wind closer to the disk, where $`\mathrm{\Gamma }_{\mathrm{Continuum}}`$ rises to $`0.02`$.
As has already been shown, at this low Eddington ratio the structure of the wind does not change significantly. However, the velocity at the base of the outflow is affected. The change in velocity due to radiative acceleration is shown in Figure 11. This figure shows both the poloidal velocity as a function of distance along the flowline, and the variation in that velocity with iterations of the model, therefore showing the convergence in the model. Figure 11 shows that line-driving near the disk does significantly accelerate the wind, but magnetic driving determines the terminal velocity at larger radii. This figure also displays how the code converges; the relatively slow convergence during the first four iterations is the result of the program slowly increasing the radiative acceleration to the computed value (increased slowly in order to avoid severely overestimating the radiative acceleration in the lines and causing sudden deceleration in later iterations). In the later iterations, the calculation converges to a final velocity profile using the full radiative acceleration. This profile shows the affect of line-driving near the base of the wind, where the velocity increases above that of the initial magnetocentrifugal wind. However, as already mentioned, the magnetocentrifugal wind (in this model, where $`L/L_{\mathrm{Edd}}=0.01`$) still determines the velocity at large distances. At those distances, the gas is too ionized to be appreciably accelerated by line driving.
It is important to note that considering non-Sobolev effects leads to large changes in $`M_{\mathrm{lines}}`$, the line force multiplier, found in the above calculations: “capping” the Sobolev absorption length-scale by the actual absorbing column length leads to smaller optical depths and larger line accelerations. This is illustrated in Figure 12, where the fiducial model has been calculated with the Sobolev approximation as well as our non-Sobolev treatment (see Section 3.3.2). The difference in force multiplier is due to the strict Sobolev treatment overestimating the column for the low-acceleration gas near the base of the wind. This relatively straightforward modification is very important to correctly estimate the optical depth, as can be seen in Figure 12.
Overall, in this section, the fiducial model has shown how the wind velocity, number density, column density, and radiative acceleration all interact to determine the final state of a magnetocentrifugal wind. These components have not previously been self-consistently combined in a magnetocentrifugal model, and so yield a new look at the state of these winds. In addition, the importance of the number densities and column densities to the final result are most apparent, and clearly merit further investigation, which will be addressed in §5.
## 5. Dependence of Wind Structure on Model Parameters
Having analyzed the fiducial model in detail, and observed how that model’s properties change with height, how sensitive are the trends in §4 to those fiducial parameters? This is an important question, and one of the key attributes of this self-similar model is that it allows some flexibility in the selection of initial parameters. In this section, we test for variations in the wind by examining how the wind changes as parameters are modified.
### 5.1. Variations with Shielding Column
One of the most difficult issues for radiative driving in AGNs is over-ionization of the wind. As shown in §4, as the wind accelerates and its density decreases, the magnetocentrifugal outflow can easily become too ionized to be efficiently accelerated to escape velocity solely by atomic lines. This is the problem of the “shielding gas” that was mentioned in §1: for pure radiative line-driving, some shielding gas is required to intercept the X-ray ionizing radiation so that the remaining UV resonant line photons can be absorbed by the wind and radiatively accelerate it to the escape velocity.
Some important papers have already been dedicated to examining the concept of shielding gas, such as Chelouche & Netzer (2003b, considering very detailed photoionization simulations of gas shields with constant column density) and Proga & Kallman (2004, where multidimensional hydrodynamics simulations with approximate radiative effects are considered). In contrast, the models presented here include a shield where the column density varies with height in a shielding wind, and where detailed photoionization simulations can be employed.
The MHD wind model presented here already launches a wind magnetocentrifugally, so it is immune to concerns of overionization. Is it therefore possible for a magnetocentrifugally-driven wind, with its commensurate drop in column density with height above the disk, to act as a shield, allowing for more efficient radiative acceleration beyond it? We can test this question by simply varying the shielding column in the fiducial model and checking the radiative acceleration seen by a wind launched behind the shield.
As shown in Figure 13, as the shielding column is increased from $`N_{\mathrm{H},\mathrm{shield},0}10^{21}`$ to $`10^{24}\mathrm{cm}^2`$, line-driving in the wind increases from $`\mathrm{\Gamma }_{\mathrm{lines}}0.05`$ to $`\mathrm{\Gamma }_{\mathrm{lines}}0.12`$. (Recall that for this model, continuum driving is $`\mathrm{\Gamma }_{\mathrm{continuum}}0.01`$.) This shows that, for large shielding columns ($`N_{\mathrm{H},\mathrm{shield},0}10^{24}\mathrm{cm}^2`$), line driving can be up to an order of magnitude more effective than continuum-driving at the base of the wind. This increase in acceleration is due to the absorption of the ionizing radiation by the shield, which allows a lower ionization state in the wind, and more line-driving due to more atomic lines. It is also apparent that, as the shielding is increased, the resultant lower total flux at the base of the wind means that the onset of significant radiative acceleration is delayed: this accounts for the offset of maximum radiative acceleration from the disk surface as the shielding column is increased.
Further, in Figure 14, the ratio of radiative acceleration to magnetic acceleration along a streamline is shown. Since the MHD effects are still supplying most of the acceleration (with an acceleration roughly equal to and opposite that of gravity), the ratio of radiative to magnetic acceleration looks much like that in Figure 13. Magnetic effects dominate in these wind models, even when large columns of shielding are included.
We have therefore shown that a magnetocentrifugal outflow can act as a shield and increase the efficiency of line-driving in the wind. However, it can also be seen that line-driving is important in these models only at the base of the wind. This arises not only from the drop in the shield’s column density with height above the disk, but the drop in the wind’s density as well (and the commensurate rise of the ionization parameter).
### 5.2. Variations with Initial Density
Owing to the increase of line-driving with decreasing ionization parameter, higher accelerations would also be expected at higher densities. Thus, we investigate the effect of changes in the initial density in the wind in Figures 15 and 16.
Displaying the effect of a range of initial densities, Figure 15 shows that line-driving is only effective in these magnetocentrifugal winds at relatively high densities. Since continuum driving is approximately constant (and relatively independent of density) at $`a/g0.01`$ for $`L/L_{\mathrm{Edd}}=0.01`$, any line-driving below that level is insignificant for these winds. Thus, for initial densities $`n_0<10^9\mathrm{cm}^3`$, line driving falls below the level of continuum driving and ceases to be important. For the highest density tested, $`n_0=10^{11}\mathrm{cm}^3`$, line driving dominates continuum driving for all locations in the wind. (The variations in each acceleration curve as a function of $`s`$ shown on this plot are chiefly due to variations in ionization parameter: line-driving is high near the disk due to shielding and the relatively high density, and rises towards the end of the streamline due to the dropping flux levels at large distances.)
Since $`n_0`$ clearly has a great impact on the radiative acceleration, how do such changes affect observables, such as the velocity? Figure 16 shows how the variation in radiative acceleration affects the poloidal velocity ($`v_\mathrm{p}`$) of the outflow. As the initial density and radiative acceleration increase, the wind’s velocity shows substantial variations from the pure magnetocentrifugal model (which dominates the $`n_0=10^7\mathrm{cm}^3`$ and $`n_0=10^8\mathrm{cm}^3`$ models). In the case of $`n_0=10^{11}\mathrm{cm}^3`$, where the greatest difference in $`v_\mathrm{p}`$ is seen, the velocity increases by a factor of $`2`$ to 3 close to the disk. Beyond the region close to the disk ($`s/r_0>1`$), however, magnetocentrifugal driving still dominates the final velocities for these winds. But the velocity differences near the disk may be observationally important, especially if acceleration near the base of the wind is the source of single-peaked emission lines (as in Murray & Chiang, 1997). If true, such emission lines may be critical in testing the differences between wind models.
### 5.3. Variations with Radiative Column
Having already tested the more obvious parameters of the initial density and shield column density, we now turn to one of the most crucial parameters for the efficiency of line-driving: the optical depth in the lines.
Under normal circumstances, where the gas velocity is sufficiently low, or where the gas column is very low, the optical depths in the lines are simply governed by the ionic columns themselves. For large accelerations or large columns, however, the optical depths are dominated by the Sobolev length, which is defined as the distance over which the relative velocity between atoms is equal to the thermal width, so that a photon emitted by one atom could be absorbed by another within that Sobolev length (see §3.3.1). With the Eddington ratio in this magnetocentrifugal wind model ($`L/L_{\mathrm{Edd}}=0.01`$), both regimes can be important. Depending on the initial parameters prescribed, the wind column can be small enough such that the column alone determines the opacity in the wind (instead of the velocity) and therefore the amount of observed acceleration, so that the Sobolev length is not important. This is demonstrated in Figure 17, where the variation in line-driving with the radiatively accelerated wind column is presented (the shielding column is held constant). For the larger columns, the optical depth in the lines increases and the radiative acceleration decreases. For the smallest columns, very large radiative acceleration is predicted due to the low opacity in the lines. This is critical for these models, for at significantly high column, the wind would see no significant line driving (this is true for magnetocentrifugal wind columns with $`N_{\mathrm{H},\mathrm{rad},0}10^{22}\mathrm{cm}^2`$).
### 5.4. Modifying the Eddington Ratio
Of key importance to applications to AGN is understanding the acceleration of outflows as a function of the Eddington ratio, $`L/L_{\mathrm{Edd}}`$. To investigate the impact of varying Eddington ratios in our model, we present Figure 18, which displays the radiative acceleration (both the combined line and continuum acceleration in panel *a* as well as the line driving in panel *b*) in the fiducial model for three different Eddington ratios: $`L/L_{\mathrm{Edd}}=0.001,0.01,`$ and $`0.1`$.
The largest variation in Figure 18a is the continuum driving increasing linearly with the Eddington ratio. As expected, the continuum acceleration, relative to gravity, is roughly equal to the Eddington ratio. The line driving, on the other hand, can be seen in both the deviations from the approximately constant continuum acceleration in Figure 18a and in Figure 18b.
We begin examining this figure by concentrating on the first three models in Figure 18, which have $`N_{\mathrm{H},\mathrm{shield},0}=10^{23}`$ cm<sup>-2</sup>. In these models in Figure 18a, the increase in acceleration due to line-driving, relative to the continuum-driving, decreases as the Eddington ratio increases. This can also be seen in Figure 18b, where the line-driving peaks near the disk ($`s/r_00.01`$) for $`L/L_{\mathrm{Edd}}=0.01`$, but decreases for Eddington ratios an order of magnitude larger (where the gas is overionized) and an order of magnitude smaller (where the radiation field doesn’t have the momentum to accelerate the wind as strongly as at $`L/L_{\mathrm{Edd}}=0.01`$).
Now we turn to the fourth model in Figure 18, where we keep $`L/L_{\mathrm{Edd}}=0.1`$ but increase the shielding level to $`N_{\mathrm{H},\mathrm{shield},0}=10^{24}`$cm<sup>-2</sup>. This model shows the importance of shielding gas for $`L/L_{\mathrm{Edd}}=0.1`$. The increase in $`L/L_{\mathrm{Edd}}`$ and the increase in shielding relative to the fiducial model allow for increased line radiative acceleration that jumps almost two orders of magnitude in strength near the base of the wind.
## 6. Results
If magnetocentrifugal winds power outflows in AGNs, they must certainly be affected by the intense radiation field they experience; in turn, such winds will also influence the efficiency of radiative acceleration. This paper has explored the radiative transfer through these magnetocentrifugal winds, and how they both are affected by and affect radiative acceleration of outflows from AGNs. The model has been used to explore the detailed dynamics and ionization of a fiducial magnetocentrifugal disk wind, showing the inter-relation between shielding column, initial number density, outflow velocity, Eddington ratio, and acceleration.
As a result of this study, these models have shown:
1. A magnetocentrifugal outflow, acting as a “shield”, can improve the efficiency of line-driving by factors of approximately two to three $`N_{\mathrm{H},\mathrm{shield},0}=10^{23}`$cm<sup>-2</sup>5.1) and by up to almost two orders of magnitude for $`N_{\mathrm{H},\mathrm{shield},0}=10^{24}`$cm<sup>-2</sup>5.4). A magnetocentrifugal wind has the advantage that it can be accelerated without regard to the ionization state, whereas radiatively-driven winds must have a low ionization parameter in order for a critical abundance of atomic lines to be present. Therefore, magnetocentrifugal winds could play an important role in acting as a radiation shield and allow large radiative accelerations. It may also be possible that pressure differences (MCGV95) or disk photons (Proga & Kallman, 2004) may help “lift” the shield; neither of those effects are considered here. Later work with this model will include the effect of disk-emitted photons.
2. The efficiency of line-driving is strongly dependent on the density at the base of the wind (§5.2). This is due to the very critical dependence of line acceleration on the ionization parameter. The lower the ionization state, the more lines exist to aid in the momentum transfer from outward-streaming photons. The density at the base of the disk is therefore crucial to setting to line-driving within these magnetocentrifugal models.
3. Small columns ($`N_{\mathrm{H},\mathrm{rad},0}10^{21}\mathrm{cm}^2`$) within magnetocentrifugal winds can be significantly accelerated by line-driving (§5.3). This point demonstrates the importance of “non-Sobolev” effects; i.e., that at low columns, the optical depths in the lines drop below the opacity given by the Sobolev length, and at such low columns, the radiative acceleration can be underestimated by the simple Sobolev approximation.
In addition, by examining the fiducial model and the above cases where model parameters are varied, these solutions have displayed the importance of considering the detailed interaction between the dynamics and photoionization in AGN outflows. Calculations of the ionization parameter along the flow and in the variation of shielding and optical depth along the flow are central issues to modeling these winds. The issues outlined above are a few of the dependences that arise from such modeling, and may indeed (by the variation in acceleration and therefore velocity along the streamlines) lead to tests to observationally determine the physics of wind launching in AGNs.
Thus, while this model has been developed to help address the above questions about shielding and the affects of radiative driving on magnetocentrifugal winds, the solutions available are not limited to the above, examined cases. Future papers will study further the variation of radiative acceleration on model parameters such as the SED, atomic line lists used to calculate the acceleration, initial densities, Eddington ratios, and other parameters of the model. The model can also be used to explore the absorption and emission features from such a wind, as well as, for instance, the possible role of “clouds” within a continuous wind (Everett et al., 2002). In addition, this model is in no way constrained to only study AGNs. The same basic physical framework could also be employed to study winds from accretion disks surrounding young stellar objects or cataclysmic variables.
## 7. Acknowledgments
I gratefully acknowledge my advisor, Arieh Königl, for many useful conversations throughout the course of this project. In addition, many thanks to Gary Ferland and his collaborators for developing Cloudy, making it freely available, and supporting it. The referee’s comments were very helpful, and those comments helped improve the paper. Thanks also to David Ballantyne, Pat Hall, Lewis Hobbs, John Kartje, Ruben Krasnopolsky, Bob Rosner, Nektarios Vlahakis, and Don York for valuable comments. This work would not have been possible without the generous support of NASA’s ATP program, in this case via grant NAG5-9063, and the support of the Natural Sciences and Engineering Research Council of Canada. This research has made use of NASA’s Astrophysics Data System.
## Appendix A Appendix A: Derivation of the Self-Similar Centrifugal Wind Equations
In this appendix, a rederivation of the system of self-similar wind equations for the magnetocentrifugal wind is presented. The equations utilized in this calculation advance upon those presented in BP82 and KK94: the wind not only has an arbitrary density power-law index, $`b`$, as in KK94, but energy conservation is not required. Since the radiation field continually inputs energy into the outflow, this is an important modification that was not fully considered in the derivation presented in KK94.
First, a stationary, axisymmetric, ideal, cold MHD flow in cylindrical coordinates $`(r,\varphi ,z)`$ is assumed. The equations are based on both the radial and vertical momentum equations:
$`v_r{\displaystyle \frac{v_r}{r}}+v_z{\displaystyle \frac{v_r}{z}}{\displaystyle \frac{v_\varphi ^2}{r}}`$ $`=`$ $`\rho {\displaystyle \frac{\mathrm{\Phi }}{r}}{\displaystyle \frac{B_z}{4\pi }}\left({\displaystyle \frac{B_z}{r}}{\displaystyle \frac{B_r}{z}}\right){\displaystyle \frac{B_\varphi }{4\pi r}}{\displaystyle \frac{(rB_\varphi )}{r}}`$ (A1)
$`\rho (𝐯)v_z`$ $`=`$ $`\rho {\displaystyle \frac{\mathrm{\Phi }}{z}}{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{B^2}{z}}+{\displaystyle \frac{1}{4\pi }}(𝐁)B_z,`$ (A2)
where $`𝐯`$ is the fluid velocity and $`𝐁`$ is the magnetic field. The thermal term is neglected in the limit that thermal affects are much less important than magnetocentrifugal and radiative-driving effects. $`\mathrm{\Phi }`$ is the effective gravitational potential, defined as
$`\mathrm{\Phi }=[1\mathrm{\Gamma }(\theta )]{\displaystyle \frac{GM_{}}{(r^2+z^2)^{1/2}}},`$ (A3)
where $`M_{}`$ is the mass of the central black hole, and $`\mathrm{\Gamma }(\theta )`$ gives the local radiative-to-gravitational radial acceleration (see eq. ).
These equations are solved by first stipulating mass conservation,
$`(\rho 𝐯)=0,`$ (A4)
and then relating the flow velocity to the magnetic field via
$`𝐯(𝐫)={\displaystyle \frac{k𝐁(𝐫)}{4\pi \rho (𝐫)}}+\omega (𝐫)\times 𝐫,`$ (A5)
(e.g., Chandrasekhar, 1956; Mestel, 1961), where $`k/4\pi `$ is the ratio of mass flux to magnetic flux, and $`\omega (𝐫)`$ and $`\rho (𝐫)`$ are the field angular velocity and gas mass density of the flow, respectively. Both $`\omega `$ and $`k`$ are constant along magnetic fieldlines. In addition, while the specific energy is not constant, the total specific angular momentum
$`l=rv_\varphi {\displaystyle \frac{rB_\varphi }{k}}`$ (A6)
is conserved.
Self-similarity is then imposed on this system by specifying
$`𝐫`$ $`=`$ $`[r_0\xi (\chi ),\varphi ,r_0\chi ],`$ (A7)
$`𝐯`$ $`=`$ $`[\xi ^{}(\chi )f(\chi ),g(\chi ),f(\chi )]v_{k,0},`$ (A8)
where $`v_{\mathrm{k},0}`$ is the Keplerian speed at the base of the outflow, $`v_{\mathrm{k},0}=(GM_{}/r_0)^{1/2}`$, and the prime indicates differentiation with respect to $`\chi `$. At the same time, the above constants are re-expressed in dimensionless form:
$`\lambda `$ $``$ $`{\displaystyle \frac{l}{(GM_{}r_0)^{1/2}}},`$ (A9)
$`\kappa `$ $``$ $`{\displaystyle \frac{k(1+\xi _{}^{}{}_{0}{}^{2})^{1/2}}{B_{p,0}}}v_{k,0},`$ (A10)
where $`B_{p,0}`$ is the poloidal magnetic field strength at the base of the wind.
As in KK94, a general power-law scaling of the density and magnetic field along the disk’s surface is defined:
$`\rho _0`$ $``$ $`r_0^b,`$ (A11)
$`B_0`$ $``$ $`r_0^{(b+1)/2}.`$ (A12)
With this self-similar specification, the radial and vertical momentum equations become, after some simplification:
$`{\displaystyle \frac{f\xi ^{}m^{}}{\kappa \xi J}}{\displaystyle \frac{f^2\xi ^{}}{\xi J}}+\xi ^{\prime \prime }f^2{\displaystyle \frac{(\lambda m\xi ^2)^2}{\xi ^3(m1)^2}}`$ $`=`$ $`\xi [1\mathrm{\Gamma }(\theta )]S^3{\displaystyle \frac{f}{\kappa \xi J^2}}({\displaystyle \frac{(1+\xi ^2)(b+1)}{2}}+`$ (A13)
$`{\displaystyle \frac{(\chi +\xi \xi ^{})\xi ^{}}{\xi }}{\displaystyle \frac{\xi ^{\prime \prime }}{JS^2}}){\displaystyle \frac{\kappa f}{\xi }}{\displaystyle \frac{(\lambda \xi ^2)}{(m1)}}`$
$`[{\displaystyle \frac{(\lambda \xi ^2)}{(m1)}}{\displaystyle \frac{(b+1)}{2}}+`$
$`\chi ({\displaystyle \frac{2\xi \xi ^{}}{(m1)}}+{\displaystyle \frac{(\lambda \xi ^2)m^{}}{(m1)^2}})],`$
$`{\displaystyle \frac{f}{\kappa \xi J}}(m^{}f\kappa \xi ^{}J+f\kappa \xi \chi \xi ^{\prime \prime })`$ $`=`$ $`[1\mathrm{\Gamma }(\theta )]\chi S^3+{\displaystyle \frac{f\xi ^{}}{\kappa \xi J^2}}({\displaystyle \frac{(1+\xi ^2)(b+1)}{2}}+`$ (A14)
$`{\displaystyle \frac{(\chi +\xi \xi ^{})\xi ^{}}{\xi }}{\displaystyle \frac{\xi ^{\prime \prime }(\chi ^2+\xi ^2)}{J}})`$
$`\xi ^{}\kappa f(\lambda \xi ^2)\left({\displaystyle \frac{(b+1)(\lambda \xi ^2)2(\lambda +\xi ^2)}{2\xi (m1)^2}}\right)+`$
$`{\displaystyle \frac{(\lambda \xi ^2)^2m^{}\kappa f}{(m1)^3}},`$
where
$`m`$ $``$ $`{\displaystyle \frac{4\pi \rho v_p^2}{B_p^2}}=\kappa \xi fJ=\mathrm{square}\mathrm{of}\mathrm{poloidal}\mathrm{Alfv}\stackrel{´}{\mathrm{e}}\mathrm{n}\mathrm{Mach}\mathrm{number},`$ (A15)
$`\kappa `$ $``$ $`{\displaystyle \frac{k(1\xi _{}^{}{}_{0}{}^{2})^{\frac{1}{2}}v_{k,0}}{B_{p,0}}}=\mathrm{dimensionless}\mathrm{ratio}\mathrm{of}\mathrm{mass}\mathrm{flux}\mathrm{to}\mathrm{magnetic}\mathrm{flux},`$ (A16)
$`\lambda `$ $``$ $`{\displaystyle \frac{l}{(GMr_0)^{\frac{1}{2}}}}=\mathrm{normalized}\mathrm{angular}\mathrm{momentum},`$ (A17)
$`J`$ $``$ $`\xi \chi \xi ^{},`$ (A18)
$`S`$ $``$ $`1/\sqrt{\xi ^2+\chi ^2}.`$ (A19)
The two equations (A13) and (A14) define the differential equations for $`m^{}`$ and $`\xi ^{\prime \prime }`$, which are, respectively, the spatial gradient in the poloidal Alfvén mach number (gradient with respect to height, $`\chi `$) and the (cylindrical) radial velocity gradient (again with respect to $`\chi `$).
One can see from close inspection of the above equations that many of the terms have a denominator of $`(m1)`$, showing that when the gas crosses the Alfvén point (where $`m=1`$), the equations become singular. Rewriting and solving the $`m^{}`$ equation for the value of $`m^{}`$ at the Alfvén singular point:
$`m_A^{}`$ $`=`$ $`2\xi J[8\chi \kappa ^2\lambda m^{}\xi ^{}J^3+4(1+b)\kappa ^2\lambda \xi ^2J^2(\chi +\xi \xi ^{})+m^2(\chi +\xi \xi ^{})(2\kappa ^2\lambda ^2S+`$ (A21)
$`(1+b)+2\kappa ^2\lambda ^34\chi \kappa ^2\lambda ^{\frac{3}{2}}(\lambda S)\xi ^{}+((1+b)+2\chi ^2\kappa ^2\lambda (\lambda S))\xi ^2+2\kappa ^2\lambda \mathrm{\Gamma }(\theta )J^2)]/`$
$`\left[4\xi J\left({\displaystyle \frac{4\kappa ^2\lambda \xi ^2J^2}{S^2}}+m^2(\chi +\xi \xi ^{})^2\right)\right].`$
This constraint is used to start the integral at the Alfvén point with the value of $`m_A^{}`$ given by equation (A21).
As covered in the main text, these equations are solved using a “shooting algorithm,” integrating from both the Alfvén point and the disk surface towards an intermediate point. Matching the integrals of three first-order equations (given by the first-order equation for $`m^{}`$ and the second-order equation for $`\xi ^{\prime \prime }`$ in eqs. A13 and A14) at the common point allows us to solve for the three free parameters in the system: $`\xi _0^{}`$, $`\xi _A^{}`$, and $`\chi _A`$.
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# 1 Introduction
## 1 Introduction
One of the most pressing problems in string theory is the issue of moduli stabilization. Lately, important progress has been accomplished by taking into account the freedom of switching on (quantized) RR and NS fluxes in the compact closed string background. This road has been particularly explored in the context of type IIB theory, where RR/NS fluxes create a superpotential that depends on the complex structure fields and the axi-dilaton and allows to fix these fields dynamically . In order to further determine the Kähler moduli, non-perturbative effects have also been put to work . For other proposals for fixing Kähler moduli, see . In simple IIB toroidal orientifolds in general the moduli are fixed in regions in which the compact volume is of order the string scale and/or the dilaton is of order one, so that the validity of an effective 4-dimensional supergravity action is open to question. One of the important reasons why this is the case is that the values of fluxes are strongly constrained by RR tadpole cancellation conditions. The situation is ameliorated in compactifications on IIB Calabi-Yau orientifolds since the flux contribution to tadpoles is typically large and one can generate vacua with all moduli and the dilaton in regions of parameter space where the effective supergravity approximation may be trusted .
In comparison, less effort has been devoted to the similar moduli-fixing problem in the case of type IIA compactifications. The fact that in type IIA there are fluxes with both even and odd rank suggests that both complex structure and Kähler moduli fields may be determined simultaneously without resorting to non-perturbative effects. This has been anticipated by several authors . Indeed, in the type IIA case, RR/NS backgrounds give rise to superpotentials depending both on Kähler and complex structure moduli, but with no terms mixing both kinds of moduli. Furthermore, in simple toroidal settings one can also include metric fluxes and generate superpotential terms coupling both kinds of moduli . The so-called metric fluxes can arise partially from T-duality of NS fluxes . More generally, turning on constant metric fluxes corresponds to Scherk-Schwarz reductions that can be understood as compactifications on twisted tori .
The purpose of this paper is twofold. We will first present a detailed study of minima of the moduli potential induced by RR, NS and metric fluxes in the simple $`\mathrm{T}^6/(\mathrm{\Omega }(1)^{F_L}I_3)`$ type IIA orientifold. We concentrate on the potential for the dilaton and the diagonal Kähler and complex structure moduli, which may be also viewed as the only untwisted moduli of a related $`_2\times _2`$ orientifold. We argue though that the results found for $`\mathrm{T}^6/(\mathrm{\Omega }(1)^{F_L}I_3)`$ ignoring off-diagonal moduli still constitute extrema of the potentials which are stable in relevant cases. We find four classes of (non-singular) vacua which correspond to $`𝒩`$=1 supersymmetric models in Minkowski space, no-scale, AdS with $`𝒩`$=1 supersymmetry and non-supersymmetric AdS models. In the Minkowski cases only a few of the moduli may be determined. On the other hand, the AdS vacua look particularly interesting since all moduli are stabilized (except for a combination of axion-like fields, we come back to this point below).
The structure of vacua in both Minkowski and AdS space depends very much on the existence or not of metric fluxes which lead to some remarkable new features. In particular, in $`𝒩`$=1 supersymmetric AdS vacua without metric fluxes, NS and RR fluxes always contribute to RR tadpoles like D6-branes do. The RR tadpole cancellation conditions restrict some of the flux parameters but some others (particularly the RR 4-form and 2-form fluxes) remain unconstrained. Due to this fact, one can easily find minima with all closed string moduli stabilized in regions with large volume and small dilaton, so that the effective supergravity action should be a good approximation. The (negative) cosmological constant may be made arbitrarily small for sufficiently large fluxes. $`𝒩`$=1 supersymmetric AdS minima in generic IIA orientifolds with NS and RR fluxes were recently analyzed in . In this case we obtain analogous results. We also find examples of AdS vacua with broken supersymmetry and all moduli stabilized.
In the $`𝒩`$=1 AdS vacua with RR/NS backgrounds and metric fluxes turned on a particular new property appears. The flux contribution to the RR tadpoles may be positive, negative or zero. This is due to the fact that the RR-charge $`Q_{RR}`$ has the schematic structure $`Q_{RR}(m\overline{H}_3+\omega \overline{F}_2)`$, where $`m`$ is the $`0`$-form of massive IIA supergravity and $`\omega `$ represents the metric flux parameters. The signs of the different fluxes are not arbitrary since they are correlated to the signs of the real parts of the moduli fixed at the minimum. The fact that we can add fluxes determining moduli but not contributing to RR tadpoles is important since this means that we have a rigid ‘corset’, namely the concrete AdS $`𝒩`$=1 background, which can be added to any RR tadpole-free configuration of D6-branes to stabilize all moduli. On the other hand, considering fluxes contributing like O6-planes to RR tadpoles is interesting since we can dispose of orientifold planes in certain cases. We find, at least for the massive $`m0`$ case, that the AdS supersymmetric minima may be made to reside at points with large compact volume and small dilaton so that corrections are under control.
The observation that one can have string backgrounds leading to vanishing or negative RR charges in AdS is not new, see e.g., and is also related to the fact that in the presence of metric fluxes one is really dealing with non-Calabi-Yau manifolds with peculiar topology. In our case we have a twisted torus with a half-flat structure. Then, a D6-brane wrapping a certain 3-cycle in the original torus with RR charges cancelled by some background including metric fluxes may be alternatively understood as a homologically trivial brane in the twisted torus which is however stable because it wraps a generalized calibrated 3-cycle .
The second main topic in this paper is the inclusion of D6-branes in models with fluxes and the construction of some semi-realistic examples. It turns out that adding branes gives rise to some new interesting features beyond the obvious one of their contribution to RR tadpoles. Stacks of D6-branes wrapping 3-cycles contain in general $`U(1)`$ fields which couple to RR fields, the imaginary parts of the complex structure moduli and the axi-dilaton. In particular, some of the $`U(1)`$’s get Stückelberg masses by combining with these RR fields. Now, if some of such RR fields get masses from fluxes some inconsistency is expected, in particular, the flux induced superpotential would violate gauge invariance. Therefore, we need that the linear combinations of RR fields combining with $`U(1)`$’s and those getting masses from fluxes should be orthogonal. We find that this is guaranteed as long as the Freed-Witten anomaly induced on the world-volume of D6-branes by the fluxes cancels. This turns out to be an important constraint on the Minkowski minima. In the case of AdS $`𝒩`$=1 supersymmetric minima, with the real parts of all moduli determined, we find the interesting result that the cancellation of the FW anomaly automatically forces the branes to preserve supersymmetry, i.e. to wrap special Lagrangian (slag) cycles. Stating it the other way around, any D6-brane wrapping slag cycles will automatically be free of the FW anomaly in this background.
It is important to see how far one can go in stabilizing all moduli in models with possible phenomenological relevance. We present examples of configurations of D6-branes wrapping 3-cycles on the torus and intersecting at angles with chiral MSSM-like spectra and fixed moduli. The models contain three generations of quarks and leptons and one Higgs set, with the gauge group of the SM extended by one or two extra $`U(1)`$’s and some extra heavy vector-like $`SU(2)_L`$ doublets and singlets. Some of the examples live in Minkowski space (either $`𝒩`$=1 supersymmetric or no-scale), in which case only a few moduli are fixed. On the other hand, we present the first semi-realistic $`𝒩`$=1 supersymmetric model in AdS with all closed string moduli stabilized. This model requires the presence of both metric and NS/RR fluxes so that a ‘wrong sign’ contribution to tadpoles, mimicking orientifold planes, is obtained. The minima are stabilized in a perturbative regime and all physical quantities like gauge and Yukawa couplings are given in terms of otherwise undetermined NS and RR fluxes.
The outline of the paper is as follows. In the next section we introduce the basic tools needed to describe the $`\mathrm{T}^6/(\mathrm{\Omega }(1)^{F_L}I_3)`$ type IIA orientifold with fluxes. In section 3 we display the structure of the flux-induced superpotential and discuss the contribution of different fluxes to RR charges. The systematic analysis of the different vacua of the flux-induced potential is presented in section 4. Examples of Minkowski and AdS vacua with/without supersymmetry and with/without metric fluxes are reported. In section 5 we examine the constraints coming from the Freed-Witten anomaly and their connection to the open string $`U(1)`$’s. Specific semi-realistic intersecting D6-brane models with MSSM-like spectrum are discussed in section 6. This includes the AdS $`𝒩`$=1 supersymmetric example with all closed string moduli stabilized. We present some comments and conclusions in section 7. Some related results are collected in two appendices. In appendix A we study the $`SU(3)`$ structure of the twisted torus and discuss $`𝒩`$=1 vacua in terms of torsion classes. In appendix B we present some non-supersymmetric D6-brane configurations with moduli stabilized in AdS.
## 2 Basic Features
The aim of this section is to present the concepts needed to describe the low-energy effective action of type IIA orientifolds in the presence of background fluxes. We first introduce the moduli fields and then exhibit the superpotential induced by NS and RR fluxes. We next define the so-called metric fluxes and recall how they can partially arise from T-duality of NS fluxes in type IIB.
### 2.1 IIA orientifolds: moduli and NS/RR fluxes
This section contains a brief review of the structure of IIA Calabi-Yau orientifolds. The treatment follows which the reader can consult for more details. We limit to a discussion of moduli fields, NS/RR fluxes, and flux-induced superpotentials. Our main purpose is to apply the general results to a simple IIA toroidal orientifold.
Compactification of type IIA strings on a Calabi-Yau 3-fold $`Y`$ gives a $`D`$=4, $`𝒩`$=2 theory with $`h_{11}`$ vector multiplets and $`(1+h_{12})`$ hypermultiplets . Turning on fluxes for the NS and RR field strengths generates a potential for the scalars in these multiplets . To obtain a $`𝒩`$=1 theory one can implement an orientifold projection by $`\mathrm{\Omega }_P(1)^{F_L}\sigma `$, where $`\mathrm{\Omega }_P`$ is the world-sheet parity operator, $`(1)^{F_L}`$ is the space-time fermionic number for left-movers and $`\sigma `$ is an order two involution of $`Y`$. The action of $`\sigma `$ on the Kähler form and the holomorphic 3-form is $`\sigma (J)=J`$ and $`\sigma (\mathrm{\Omega })=e^{2i\theta }\mathrm{\Omega }^{}`$. We take $`\theta =0`$ and $`\sigma (z^i)=\overline{z}^i`$, where $`z^i`$ are local complex coordinates. This implies O6-planes whose tadpoles can be cancelled by adding D6-branes or flux, as we will see.
The closed (1,1) forms split into $`h_{11}^+`$ and $`h_{11}^{}`$, according to whether they are even or odd under $`\sigma `$. There is an equal number $`(1+h_{12})`$ of even and odd 3-forms. Then, the resulting matter content from the closed string sector consists of $`h_{11}^+`$ vector multiplets, $`h_{11}^{}`$ chiral multiplets corresponding to Kähler moduli, and $`(1+h_{12})`$ chiral multiplets corresponding to the dilaton and the complex structure moduli . The scalar components of the Kähler moduli, denoted $`T_A`$, are defined in terms of the complexified Kähler form as
$$J_c=B+iJ=i\underset{A=1}{\overset{h_{11}^{}}{}}T_A\omega _A,$$
(2.1)
where $`\omega _A`$ are the $`\sigma `$-odd (1,1) closed forms. The complex structure moduli, denoted $`N_L`$, $`L=0,\mathrm{},h_{12}`$, can be extracted from
$$iN_L=_Y\mathrm{\Omega }_c\beta _L;\mathrm{\Omega }_c=C_3+i\mathrm{Re}(C\mathrm{\Omega }),$$
(2.2)
where $`C_3=_L\xi _L\alpha _L`$ provides the axions. Here $`\alpha _L`$ and $`\beta _L`$ are respectively the $`\sigma `$-even and $`\sigma `$-odd 3-forms. The field $`C`$ is in turn specified by
$$C=e^{\varphi _4}e^{K_{cs}/2};K_{cs}=\mathrm{log}[\frac{i}{8}_Y\mathrm{\Omega }\mathrm{\Omega }^{}],$$
(2.3)
where $`\varphi _4`$ is the T-duality invariant four-dimensional dilaton given by $`e^{\varphi _4}=e^{\varphi _{10}}/\sqrt{\mathrm{vol}Y}`$.
To be more explicit let us now consider the example of a factorized 6-torus $`_{j=1}^3\mathrm{T}_j^2`$. The area of each sub-torus is $`(2\pi )^2A_j`$, where $`A_j=R_x^jR_y^j`$. The $`A_j/\alpha ^{}`$ are thus the real part, $`t_j`$, of three Kähler moduli $`T_j`$. For each sub-torus we take a square lattice, consistent with the orientifold projection. The complex structure parameter of each $`\mathrm{T}_j^2`$ is then $`\tau _j=R_y^j/R_x^j`$. It is known, see e.g. , that in this setup the IIA $`D`$=4 fields $`S`$ and $`U_i`$, corresponding to the dilaton and complex structure moduli, have real parts
$$\mathrm{Re}Ss=\frac{e^{\varphi _4}}{\sqrt{\tau _1\tau _2\tau _3}};\mathrm{Re}U_iu_i=e^{\varphi _4}\sqrt{\frac{\tau _j\tau _k}{\tau _i}};ijk,$$
(2.4)
where $`e^{\varphi _4}=e^\varphi /\sqrt{t_1t_2t_3}`$. We will next obtain these results from the general analysis of . As usual, the holomorphic 3-form can be written as
$$\mathrm{\Omega }=(dx^1+i\tau _1dy^1)(dx^2+i\tau _2dy^2)(dx^3+i\tau _3dy^3),$$
(2.5)
where $`y^i=x^{i+3}`$. The orientifold involution acts as $`\sigma (x^i)=x^i`$ and $`\sigma (y^i)=y^i`$. The even and odd 3-forms with one leg on each sub-torus are
$`\alpha _0`$ $`=`$ $`dx^1dx^2dx^3;\beta _0=dy^1dy^2dy^3,`$
$`\alpha _1`$ $`=`$ $`dx^1dy^2dy^3;\beta _1=dy^1dx^2dx^3,`$ (2.6)
$`\alpha _2`$ $`=`$ $`dy^1dx^2dy^3;\beta _2=dx^1dy^2dx^3,`$
$`\alpha _3`$ $`=`$ $`dy^1dy^2dx^3;\beta _3=dx^1dx^2dy^3.`$
Our normalization is such that $`_{\mathrm{T}^6}\alpha _I\beta _J=\delta _{IJ}`$. Substituting in (2.3) we find $`C=\mathrm{Re}S`$. From (2.2) we then obtain the corresponding moduli $`N_0=S`$ and $`N_i=U_i`$.
We have only considered 3-forms with one leg on each $`\mathrm{T}_j^2`$ because these are the directions in which we are going to switch on fluxes. If the orientifold has an extra $`_2\times _2`$ symmetry these are in fact the only invariant forms.
The next step is to turn on background fluxes. The NS $`H_3`$ is odd under the orientifold action, thus the general flux allowed is
$$\overline{H}_3=\underset{L=0}{\overset{h_{12}}{}}h_L\beta _L.$$
(2.7)
For the RR forms, $`F_0`$ and $`F_4`$ are even while $`F_2`$ and $`F_6`$ are odd under the orientifold projection. We can then have the general expansions
$`\overline{F}_0=m`$ ; $`\overline{F}_6=e_0d\mathrm{vol}_6`$
$`\overline{F}_2={\displaystyle \underset{A=1}{\overset{h_{11}^{}}{}}}q_A\omega _A`$ ; $`\overline{F}_4={\displaystyle \underset{A=1}{\overset{h_{11}^{}}{}}}e_A\stackrel{~}{\omega }_A,`$ (2.8)
where $`\stackrel{~}{\omega }_A`$ are the $`h_{22}^+=h_{11}^{}`$ $`\sigma `$-even (2,2) forms. There are also quantization conditions
$$\frac{\mathrm{}^3\mu _1}{2\pi }_{\mathrm{\Pi }_3}\overline{H}_3;\frac{\mathrm{}^p\mu _{p2}}{2\pi }_{\mathrm{\Pi }_p}\overline{F}_p,$$
(2.9)
for any $`p`$-cycle $`\mathrm{\Pi }_p`$ in $`Y`$. Here $`\mathrm{}=2\pi \sqrt{\alpha ^{}}`$, and $`\mu _p=1/(2\pi )^p\alpha ^{(p+1)/2}`$ . We normally take the various forms, e.g. $`\beta _L`$, to belong to an integer basis so that in units of $`2\pi /\mu _{p2}\mathrm{}^p=1/\mathrm{}`$ the various coefficients, such as $`h_L`$, are integers. Actually, to avoid subtleties with exotic orientifold planes we take the coefficients to be even. Notice that by including the factors of $`\mathrm{}`$ explicitly, all forms have dimensions $`(\mathrm{length})^1`$. With these conventions the moduli fields are all dimensionless.
The RR fluxes generate a superpotential for the Kähler moduli that can be written as
$$W_K=_Ye^{J_c}\overline{F}_{RR},$$
(2.10)
where $`\overline{F}_{RR}`$ represents a formal sum of the even RR fluxes. This result can be obtained applying mirror symmetry to the type IIB superpotential . It can also be derived performing the explicit Kaluza-Klein reduction which also allows to determine the superpotential for the complex structure moduli due to NS flux, namely
$$W_Q=_Y\mathrm{\Omega }_c\overline{H}_3=i\underset{L=0}{\overset{h_{12}}{}}h_LN_L,$$
(2.11)
where we have used (2.2). The Kähler potential for both kinds of moduli are given by
$$K_K=\mathrm{log}[\frac{4}{3}_YJJJ];K_Q=\mathrm{log}e^{4\varphi _4}.$$
(2.12)
There are corrections to $`K_K`$ due to world-sheet instantons and to $`K_Q`$ due to D2 instantons . We will see later that one can locate the minima of the potential in regions with large volume and small dilaton in which these corrections should be in principle under control.
It is instructive to apply these results to the $`_{j=1}^3\mathrm{T}_j^2`$ example. The $`\beta _L`$ are given in (2.6) whereas
$$\omega _i=dx^idy^i;\stackrel{~}{\omega }_i=dx^jdy^jdx^kdy^k;ijk.$$
(2.13)
Notice that $`_{\mathrm{T}^6}\omega _i\stackrel{~}{\omega }_j=\delta _{ij}`$. It is straightforward to substitute the flux expansions in (2.10) and (2.11) to obtain
$`W_K`$ $`=`$ $`e_0+i{\displaystyle \underset{i=1}{\overset{3}{}}}e_iT_iq_1T_2T_3q_2T_1T_3q_3T_1T_2+imT_1T_2T_3,`$
$`W_Q`$ $`=`$ $`ih_0Si{\displaystyle \underset{i=1}{\overset{3}{}}}h_iU_i.`$ (2.14)
These superpotentials have been recently discussed in . Finally, the Kähler potential takes the usual expression
$$K=\mathrm{log}(S+S^{})\underset{i=1}{\overset{3}{}}\mathrm{log}(U_i+U_i^{})\underset{i=1}{\overset{3}{}}\mathrm{log}(T_i+T_i^{}).$$
(2.15)
### 2.2 Metric fluxes and twisted tori
In the next section we will see how superpotential terms mixing Kähler and complex structure moduli, including the dilaton, can be generated by switching on so-called metric fluxes. Such backgrounds appear naturally in the context of Scherk-Schwarz reductions . In turn these can be shown (see e.g. ) to be equivalent to compactification on a twisted torus defined by
$$d\eta ^P=\frac{1}{2}\omega _{MN}^P\eta ^M\eta ^N,$$
(2.16)
where $`\omega _{MN}^P`$ are constant coefficients, antisymmetric in the lower indices. The structure constants $`\omega _{MN}^P`$ are the metric fluxes we are interested in. The $`\eta ^P`$ are the tangent 1-forms and can depend linearly on the internal coordinates $`x^M`$, concretely
$$\eta ^M=N_N^M(x)dx^N;dx^N=N_M^N(x)\eta ^M.$$
(2.17)
One can define isometry generators as
$$Z_M=N_M^N\frac{}{x^N}.$$
(2.18)
The metric fluxes are actually the structure constants of the Lie algebra generated by the $`Z_M`$, i.e.
$$[Z_M,Z_N]=\omega _{MN}^PZ_P.$$
(2.19)
Either from the Jacobi identity of the algebra or from the Bianchi identity of (2.16) one finds that the metric fluxes must satisfy
$$\omega _{[MN}^P\omega _{R]P}^S=0.$$
(2.20)
It can further be shown that $`\omega _{PN}^P=0`$ .
We can derive a helpful result for the exterior derivative of a 2-form $`X=X_{MN}\eta ^M\eta ^N`$ using (2.16). For coefficients independent of the $`x^N`$ we readily find
$$(dX)_{LMN}=\omega _{[LM}^PX_{N]}{}_{P}{}^{}.$$
(2.21)
Similar formulas can be obtained for higher forms.
We will focus on the case in which only metric fluxes of type $`\omega _{ab}^i`$, $`\omega _{jk}^i`$, $`\omega _{ib}^a`$, $`i=1,2,3`$, $`a=4,5,6`$, are allowed. This can be implemented by imposing a symmetry of (2.16) under the orientifold involution $`\eta ^i\eta ^i`$, $`\eta ^a\eta ^a`$. As in the case of RR and NS fluxes discussed previously, we are only going to switch on metric fluxes along factorized directions. These correspond to the structure constants which are invariant under a $`_2\times _2`$ symmetry whose generators transform $`(\eta ^1,\eta ^2,\eta ^3,\eta ^4,\eta ^5,\eta ^6)`$ into $`(\eta ^1,\eta ^2,\eta ^3,\eta ^4,\eta ^5,\eta ^6)`$ and $`(\eta ^1,\eta ^2,\eta ^3,\eta ^4,\eta ^5,\eta ^6)`$. In the end there are twelve metric fluxes left. To write down the relations that follow from (2.20) we introduce the notation
$$\left(\begin{array}{c}a_1\\ a_2\\ a_3\end{array}\right)=\left(\begin{array}{c}\omega _{56}^1\\ \omega _{64}^2\\ \omega _{45}^3\end{array}\right);\left(\begin{array}{ccc}b_{11}& b_{12}& b_{13}\\ b_{21}& b_{22}& b_{23}\\ b_{31}& b_{32}& b_{33}\end{array}\right)=\left(\begin{array}{ccc}\omega _{23}^1& \omega _{53}^4& \omega _{26}^4\\ \omega _{34}^5& \omega _{31}^2& \omega _{61}^5\\ \omega _{42}^6& \omega _{15}^6& \omega _{12}^3\end{array}\right).$$
(2.22)
The Jacobi identities imply the twelve constraints
$`b_{ij}a_j+b_{jj}a_i`$ $`=`$ $`0;ij`$
$`b_{ik}b_{kj}+b_{kk}b_{ij}`$ $`=`$ $`0;ijk.`$ (2.23)
There are some obvious solutions of these constraints. For instance, (1): $`b_{ij}=0`$, $`i,j`$; (2): $`a_i=0`$, $`b_{ij}=b_i\delta _{ij}`$; (3): $`a_i=a`$, $`b_{ij}=b`$, $`ij`$, $`b_{ii}=b`$.
It is also enlightening to see how the twisted torus structure arises by T-dualizing a string background including constant NS-NS 3-form flux . To simplify the discussion, as internal space we take the flat torus $`\mathrm{T}^6`$, i.e. we neglect the warp factors needed to have a solution of the equations of motion. Concretely, we start from a type IIB background
$`ds^2`$ $`=`$ $`(dx^1)^2+\mathrm{}+(dx^6)^2,`$
$`_3`$ $`=`$ $`a_1dx^1dx^5dx^6a_2dx^4dx^2dx^6a_3dx^4dx^5dx^3.`$ (2.24)
We want to perform T-dualities in $`x^1,x^2,x^3`$. For the magnetic field we choose a gauge such that $`_2`$ does not depend on the dualized coordinates. We take
$$_2=a_1x^6dx^1dx^5a_2x^4dx^2dx^6a_3x^5dx^3dx^4.$$
(2.25)
Using standard results (see e.g. appendix A in ) gives the transformed metric
$`ds^2`$ $`=`$ $`(dx^1+a_1x^6dx^5)^2+(dx^2+a_2x^4dx^6)^2+(dx^3+a_3x^5dx^4)^2`$ (2.26)
$`+`$ $`(dx^4)^2+(dx^5)^2+(dx^6)^2.`$
Moreover, $`_3^{}=0`$. All of the NS-NS flux is traded for metric flux. From the new metric we read off the following tangent 1-forms
$`\eta ^1`$ $`=`$ $`dx^1+a_1x^6dx^5;\eta ^4=dx^4,`$
$`\eta ^2`$ $`=`$ $`dx^2+a_2x^4dx^6;\eta ^5=dx^5,`$ (2.27)
$`\eta ^3`$ $`=`$ $`dx^3+a_3x^5dx^4;\eta ^6=dx^6.`$
Taking the exterior derivatives we easily identify the structure constants: $`\omega _{56}^1=a_1`$, $`\omega _{64}^2=a_2`$, and $`\omega _{45}^3=a_3`$.
Finally, an important point is that metric fluxes are also quantized. For the $`a_i`$ we have just seen that they are obtained from T-duality of NS fluxes. In general this is needed for consistency of the twisted torus structure .
## 3 IIA superpotential and RR tadpoles due to general fluxes
It is instructive to check how a number of terms in the IIA superpotential (including some induced by metric fluxes) may be obtained applying T-duality transformations to the known type IIB results. Our starting point is the type IIB orientifold $`\mathrm{T}^6/\mathrm{\Omega }(1)^{F_L}I_6`$, where $`I_6`$ reflects the six internal coordinates $`x^M`$. There are 64 O3-planes whose charge can be cancelled by adding D3-branes and/or flux. To go to type IIA we will implement mirror symmetry which is the same as T-duality in $`x^1,x^2,x^3`$. In the type IIA picture there are then O6-planes, wrapping the $`x^i`$, $`i=1,2,3`$, and one can add intersecting D6-branes.
We consider a factorized geometry in which $`\mathrm{\Omega }`$ is given in (2.5). In order to generate a superpotential for the $`\tau _i`$ and the axion-dilaton we turn on NS $`\overline{}_3`$ and RR $`\overline{}_3`$ 3-form fluxes that are conveniently expanded in the basis (2.6). For $`\overline{}_3`$ we take the most general combination
$$\overline{}_3=m\alpha _0e_0\beta _0+\underset{i=1}{\overset{3}{}}(e_i\alpha _iq_i\beta _i).$$
(3.1)
Mirror symmetry transforms $`\overline{}_3`$ into RR fluxes $`(\overline{F}_0,\overline{F}_2,\overline{F}_4,\overline{F}_6)`$ in type IIA. For $`\overline{}_3`$ we instead restrict to
$$\overline{}_3=h_0\beta _0\underset{i=1}{\overset{3}{}}a_i\alpha _i.$$
(3.2)
Under the three T-dualities only $`h_0`$ remains as NS flux, i.e. $`\overline{}_3\overline{H}_3=h_0\beta _0`$. As discussed before, the $`a_i`$ become instead metric fluxes, in fact $`\omega _{56}^1=a_1`$, $`\omega _{64}^2=a_2`$ and $`\omega _{45}^3=a_3`$. We do not turn on $`\overline{}_3\alpha _0`$ because then $`_2`$ would depend on the $`x^i`$. We do not consider $`\overline{}_3\beta _i`$ fluxes either because they lead to more complicated metrics .
The type IIB superpotential induced by the fluxes is given by
$$𝒲=(\overline{}_3\tau \overline{}_3)\mathrm{\Omega }.$$
(3.3)
Here $`\tau =C_0+ie^\varphi `$, where $`C_0`$ is the RR 0-form. Substituting (2.5), (3.1) and (3.2) we find
$$𝒲=e_0+h_0\tau +i\underset{i=1}{\overset{3}{}}(e_i+a_i\tau )\tau _iq_1\tau _2\tau _3q_2\tau _1\tau _3q_3\tau _1\tau _2+im\tau _1\tau _2\tau _3.$$
(3.4)
Upon mirror symmetry the $`\tau _i`$ go into Kähler moduli and $`\tau `$ becomes the IIA dilaton, $`\tau iS`$. Hence, we obtain the IIA superpotential
$$W_{ST}=e_0+ih_0S+\underset{i=1}{\overset{3}{}}(ie_ia_iS)T_iq_1T_2T_3q_2T_1T_3q_3T_1T_2+imT_1T_2T_3.$$
(3.5)
Notice that for $`a_i=0`$ this coincides with (2.14), with $`h_i=0`$, that was derived following the analysis of . For $`a_i0`$ it agrees with results of .
In type IIB the fluxes contribute to the $`C_4`$ tadpole with coefficient
$$N_{flux}=\overline{}_3\overline{}_3=h_0m+a_1q_1+a_2q_2+a_3q_3,$$
(3.6)
where we already substituted the fluxes at hand. Under mirror symmetry $`N_{flux}`$ transforms into a $`C_7`$ tadpole in the direction of the O6-planes. This tadpole also receives contributions from D6-branes. In general we introduce piles of $`N_a`$ intersecting D6-branes wrapping the factorizable 3-cycle
$$\mathrm{\Pi }_a=(n_a^1,m_a^1)(n_a^2,m_a^2)(n_a^3,m_a^3),$$
(3.7)
and the corresponding orientifold images wrapping $`_i(n_a^i,m_a^i)`$. Here $`n_a^i`$ $`(m_a^i)`$ are the wrapping numbers along the $`x^i`$ $`(y^i)`$ torus directions. Including the O6-planes, that wrap $`_i(1,0)`$, leads to the tadpole cancellation condition
$$\underset{a}{}N_an_a^1n_a^2n_a^3+\frac{1}{2}(h_0m+a_1q_1+a_2q_2+a_3q_3)=16.$$
(3.8)
This agrees with the result of . Tadpoles due to fluxes of the NS and the RR 0-form have been considered in . To our knowledge, tadpoles due to metric fluxes were first discussed in .
The tadpole condition can also be derived from the equation of motion for $`C_7`$ in type IIA. Let $`G_2=dC_1+mB_2+\overline{F}_2`$ and $`{}_{}{}^{}F_{2}^{}=F_8=dC_7`$, then the relevant piece in the action is
$$_{M_4\times Y}[C_7dF_2+C_7(m\overline{H}_3+d\overline{F}_2)]+\underset{a}{}N_a_{M_4\times \mathrm{\Pi }_a}C_7.$$
(3.9)
The first term arises from the kinetic energy $`G_2{}_{}{}^{}G_{2}^{}`$. The second term takes into account the coupling to D6-branes and O6-planes. The important point to notice is that $`d\overline{F}_20`$ due to the metric fluxes. For instance, using (2.21), we obtain $`(d\overline{F}_2)_{456}=(a_1q_1+a_2q_2+a_3q_3)`$, and thus recover (3.8). Moreover, from other components of $`C_7`$ there are further cancellation conditions
$`{\displaystyle \underset{a}{}}N_an_a^1m_a^2m_a^3+{\displaystyle \frac{1}{2}}(mh_1q_1b_{11}q_2b_{21}q_3b_{31})`$ $`=`$ $`0,`$
$`{\displaystyle \underset{a}{}}N_am_a^1n_a^2m_a^3+{\displaystyle \frac{1}{2}}(mh_2q_1b_{12}q_2b_{22}q_3b_{32})`$ $`=`$ $`0,`$ (3.10)
$`{\displaystyle \underset{a}{}}N_am_a^1m_a^2n_a^3+{\displaystyle \frac{1}{2}}(mh_3q_1b_{13}q_2b_{23}q_3b_{33})`$ $`=`$ $`0.`$
These also agree with the conditions found in .
We have just seen that the metric fluxes $`b_{ij}`$ create RR tadpoles. Recently it has been observed that they also generate superpotential terms involving the $`U_k`$, as expected from the fact that the fluxes $`a_i`$ produce terms involving $`S`$. Performing a generalized dimensional reduction on the twisted torus it has been shown in that the superpotential for the complex structure moduli is an extension of $`W_Q`$ that can be expressed as
$$W_Q=_Y\mathrm{\Omega }_c(\overline{H}_3+dJ_c).$$
(3.11)
Such an expression was already proposed in . The metric fluxes appear in $`dJ_c`$ that is computed using (2.21). A similar modification of the superpotential occurs in heterotic compactifications on non-Kähler manifolds . In our setup, computing $`W_Q`$ and combining with $`W_K`$ in (2.14) yields the full superpotential
$`W`$ $`=`$ $`e_0+ih_0S+{\displaystyle \underset{i=1}{\overset{3}{}}}[(ie_ia_iSb_{ii}U_i{\displaystyle \underset{ji}{}}b_{ij}U_j)T_iih_iU_i]`$ (3.12)
$``$ $`q_1T_2T_3q_2T_1T_3q_3T_1T_2+imT_1T_2T_3.`$
This is the result obtained in .
An obvious and interesting question is how the above superpotential changes when we include non-diagonal moduli related to parameters $`\tau _{ia}`$ and $`t_{ia}`$ that can appear in $`\mathrm{\Omega }`$ and $`J`$. In (3.12) only the diagonal parameters, i.e. $`\tau _i\tau _{i,i+3}`$ and $`t_it_{i,i+3}`$, are taken into account. Furthermore, we have only switched on diagonal fluxes, along (2.6) and (2.13), which we still continue to do. Now, from the general expressions of $`W_K`$ and $`W_Q`$, eqs. (2.10) and (3.11), it is easy to see that the diagonal fluxes $`h_i`$, $`e_i`$ and $`a_i`$, or $`h_0`$, do not excite the off-diagonal moduli, generically denoted $`T^{}`$ and $`U^{}`$. However, the fluxes $`m`$, $`q_i`$ and $`b_{ij}`$ can potentially generate a superpotential for $`T^{}`$ and $`U^{}`$ that is at least quadratic in these fields. The Kähler potential has also quadratic corrections. The upshot is that when we look for supersymmetric minima there is always a solution $`T^{}=U^{}=0`$ and the diagonal moduli fixed as when $`T^{}`$ and $`U^{}`$ are not included. We know that at the point $`T^{}=U^{}=0`$ there is a global $`_2\times _2`$ symmetry and furthermore, supersymmetry guarantees that this minimum is stable. In the following we will then disregard the off-diagonal moduli.
For future purposes we define
$`\stackrel{~}{T}_I`$ $`=`$ $`(i,T_1,T_2,T_3);A_{IJ}=\left(\begin{array}{cccc}h_0& h_1& h_2& h_3\\ a_1& b_{11}& b_{12}& b_{13}\\ a_2& b_{21}& b_{22}& b_{23}\\ a_3& b_{31}& b_{32}& b_{33}\end{array}\right)`$ (3.17)
$`\stackrel{~}{U}_I`$ $`=`$ $`(S,U_1,U_2,U_3).`$
The $`\stackrel{~}{U}_I`$ dependent superpotential, due to NS and metric fluxes, takes the simple form
$$W_Q=\underset{I,J=0}{\overset{3}{}}A_{IJ}\stackrel{~}{T}_I\stackrel{~}{U}_J.$$
(3.18)
The flux contribution to $`C_7`$ tadpoles can also be written in terms of the matrix $`A`$.
Recall that the metric fluxes are constrained by the Jacobi identities (2.23). For instance, there is a solution $`b_{ji}=b_i`$, $`b_{ii}=b_i`$, $`a_i=a`$. Further choosing RR fluxes $`q_i=c_2`$ and $`e_i=c_1`$, allows a configuration with $`T_1=T_2=T_3=T`$. Then the superpotential reduces to
$$W=e_0+3ic_1T+3c_2T^2+imT^3+ih_0S3aST\underset{k=0}{\overset{3}{}}(ih_k+b_kT)U_k,$$
(3.19)
If the fluxes $`h_k`$ and $`b_k`$ are independent of $`k`$, we can also set $`U_1=U_2=U_3=U`$.
Given the fluxes leading to (3.19), the tadpole condition (3.8) becomes
$$\underset{a}{}N_an_a^1n_a^2n_a^3+\frac{1}{2}(h_0m3ac_2)=16.$$
(3.20)
On the other hand,
$`{\displaystyle \underset{a}{}}N_an_a^1m_a^2m_a^3+{\displaystyle \frac{1}{2}}(h_1m+b_1c_2)`$ $`=`$ $`0,`$
$`{\displaystyle \underset{a}{}}N_am_a^1n_a^2m_a^3+{\displaystyle \frac{1}{2}}(h_2m+b_2c_2)`$ $`=`$ $`0,`$ (3.21)
$`{\displaystyle \underset{a}{}}N_am_a^1m_a^2n_a^3+{\displaystyle \frac{1}{2}}(h_3m+b_3c_2)`$ $`=`$ $`0.`$
In the next section we will see that to obtain a minimum of the moduli scalar potential the fluxes must satisfy some relations that will in particular fix the sign of the flux contributions to the tadpoles.
## 4 Vacuum structure in IIA orientifolds with fluxes
In this section we analyze the moduli scalar potential induced by the fluxes. For the Kähler potential we assume the usual tree-level result displayed in (2.15). The scalar potential is then simply
$$V=e^K\left\{\underset{\mathrm{\Phi }=S,T_i,U_i}{}(\mathrm{\Phi }+\mathrm{\Phi }^{})^2|D_\mathrm{\Phi }W|^23|W|^2\right\},$$
(4.1)
where $`D_\mathrm{\Phi }W=_\mathrm{\Phi }W+W_\mathrm{\Phi }K`$ and $`W`$ is given in eq.(3.12). We want to look for solutions of $`V/\mathrm{\Phi }=0`$. A well known and easy to show general result is that the simpler unbroken susy conditions $`D_\mathrm{\Phi }W=0`$ imply a minimum.
As it happens in the IIB case , it is quite complicated to perform a complete analysis of all possible minima induced by the flux superpotential. We have however explored several possibilities including $`𝒩`$=1 supersymmetric Minkowski vacua, no-scale models in Minkowski, $`𝒩`$=1 supersymmetric AdS vacua and non-supersymmetric AdS vacua. Before providing specific details let us make some general comments about these vacua and their comparison to IIB mirror orientifolds discussed in . In type IIB the flux-induced superpotential only depends on the axion-dilaton and complex structure fields, and all Kähler moduli remain undetermined. In the IIA case at hand all moduli appear in the superpotential and in principle may be fixed. There are however classes of type IIB minima whose mirror IIA dual is not included among our vacua. In particular, those arising when the IIB superpotential is linear in the axion-dilaton $`\tau `$ and at the same time quadratic/cubic in the complex structure fields. In fact, most of the examples in refs. belong to this class. A naive IIB/IIA mirror symmetry would suggest terms in the superpotential linear in $`S`$ and quadratic/cubic in the Kähler moduli which are not present in our superpotential (3.12). Otherwise, the IIA options for fluxes are richer and lead to many new possibilities, not only in the number of determined moduli but also in the contribution of fluxes to the different RR tadpoles, as we discuss below.
Constraining to the dependence on the 7 moduli, $`S`$ and the diagonal $`U_i,T_i`$, in general one finds that to get (non-runaway) minima the superpotential has to depend at least on four moduli. Note also that the fields $`S`$ and $`U_i`$ appear in a similar (linear) way in the superpotential so that, e.g. given a superpotential depending on $`S`$, one can obtain another model replacing $`SU_i`$ and properly relabelling the fluxes. Compared to the original model, the new one so derived has in general different contributions to the RR tadpoles. The same is true for the $`T_i`$ moduli in the $`m=0`$ case. One can obtain new models replacing $`T_iU_i`$ and exchanging appropriately the fluxes. This will be illustrated below.
In this section we then use the flux-induced superpotential to study $`𝒩`$=1 and $`𝒩`$=0 type IIA vacua. In appendix A we will also discuss how $`𝒩`$=1 vacua in the presence of metric fluxes can be analyzed in terms of compactifications on the twisted torus seen as a manifold with $`SU(3)`$ structure.
Let us now discuss the different types of vacua in turn. To describe the results we will denote $`\mathrm{Re}T_it_i`$, $`\mathrm{Im}T_iv_i`$, $`\mathrm{Re}Ss`$ and $`\mathrm{Re}U_iu_i`$.
### 4.1 Supersymmetric Minkowski vacua
Supersymmetric Minkowski vacua have $`D_\mathrm{\Phi }W=0`$ with $`V=0`$ so that the cosmological constant vanishes. Clearly, (4.1) then implies $`W=0`$ and the supersymmetry condition reduces to $`W/\mathrm{\Phi }=0`$. To simplify notation, in the following we stop writing vevs explicitly, it should be obvious when quantities are meant to be evaluated at the extremum.
When $`W`$ depends on only three fields it suffices to analyze $`W=W_K(T_1,T_2,T_3)`$, purely due to RR fluxes. Other cases can be easily translated. The general solution of $`D_{T_i}W_K=0`$ was found in in relation to BPS black holes, and applied in in the context of type IIA vacua with fluxes. If $`W_K=0`$ there are only pathological solutions in which e.g. $`t_1=0`$.
Going to superpotentials depending on four moduli or more it is easy to find models with $`𝒩`$=1 supersymmetry in Minkowski space. One can obtain examples of this kind of minima with the superpotential depending on a) one $`S/U_i`$ modulus and three $`T_i`$; b) two $`S/U_i`$ fields and two $`T_i`$ and c) two $`S/U_i`$ and three $`T_i`$. We describe these cases below. We have not found any models with superpotentials depending on more than 5 fields. In all the examples we find we need the presence of metric fluxes and in addition $`m=0`$. In all cases there is only a partial fixing of moduli, our examples fixing at most 3 complex linear combinations of moduli. The fluxes in these examples contribute to the RR tadpoles with the same sign as D6-branes do. However they may contribute in any of the four RR tadpole directions, depending on each case. This is explained below.
#### 4.1.1 Superpotentials depending on four moduli
With four fields it is enough to study in detail $`W=W(S,T_1,T_2,T_3)`$, c.f. (3.5), independent of the $`U_i`$. Other cases, e.g. $`W(S,U_1,T_2,T_3)`$, can be mapped into this by renaming fields and choosing parameters appropriately. We look for solutions of $`W=0`$, $`_SW=0`$ and $`_{T_i}W=0`$. If there are no metric fluxes we find that $`W`$ must vanish identically to avoid $`t_i=0`$. If instead $`a_i0`$ there are solutions provided $`m=0`$. Moreover, taking the real part of $`_\mathrm{\Phi }W=0`$ gives four homogeneous equations $`a_1t_1+a_2t_2+a_3t_3=0`$ and $`a_is+q_jt_k+q_kt_j=0`$, $`ijk`$. To have $`s,t_i0`$ the determinant of this system must vanish, this is
$$(a_1q_1+a_2q_2a_3q_3)^2=4a_1q_1a_2q_2.$$
(4.2)
Let us discuss now some particular cases.
i) Example SM-1
We can for example take $`a_1=0`$ which requires $`q_1=0`$ to avoid $`t_2,t_3=0`$. The superpotential has the general form
$$W=e_0+ih_0S+i\underset{i}{}e_iT_iS(a_2T_2+a_3T_3)T_1(q_2T_3+q_3T_2).$$
(4.3)
It must be $`a_2q_2=a_3q_3`$ and for consistency one also needs $`e_2a_3=e_3a_2`$, $`h_0q_2=a_3e_1`$, and $`h_0e_2=e_0a_2`$. Note that these conditions are easily satisfied, e.g. for $`h_0=e_I=0`$. Neither the imaginary nor the real parts of the moduli are fully determined, only
$$h_0=a_2v_2+a_3v_3;e_2=a_2\mathrm{Im}S+q_3v_1;t_3=\frac{q_3t_2}{q_2};s=\frac{q_3t_1}{a_2}.$$
(4.4)
For $`s,t_i>0`$ we must have $`q_2q_3<0`$ and $`a_2q_2>0`$. Now, the flux term in the (3.8) RR tadpole is equal to $`2a_2q_2`$ and is therefore positive.
ii) Example SM-2
The above example with $`a_1=0`$ can be used to analyze vacua in which we replace one Kähler, say $`T_1`$, by a complex structure modulus, say $`U_1`$. The $`W(S,U_1,T_2,T_3)`$ superpotential is now
$$W=T_2(a_2S+b_{21}U_1)T_3(a_3S+b_{31}U_1)+e_0+ih_0Sih_1U_1+ie_2T_2+ie_3T_3.$$
(4.5)
This is clearly equivalent to (4.3) after renaming $`T_1U_1`$, $`e_1h_1`$, $`q_2b_{31}`$ and $`q_3b_{21}`$. The physics is however different. In particular, since all $`q_i=0`$ and $`m=0`$, these fluxes do not contribute at all to the RR tadpoles. Thus, this is an example in which one can fix moduli without affecting tadpoles, the fluxes are essentially unconstrained, except from the consistency conditions mentioned above. We will see later that there are $`𝒩`$=1 supersymmetric AdS vacua in which one can fix all moduli without any restriction from RR tadpoles.
iii) Example SM-3
Another simple solution of (4.2) is to take $`a_2=a_3=a_1/2`$, and $`q_2=q_3=q_1/2`$. Further selecting $`h_0=0`$ and $`e_I=0`$ one has a superpotential
$$W=a_1ST_1q_1T_2T_3+\frac{1}{2}(T_2+T_3)(a_1S+q_1T_1).$$
(4.6)
One obtains a supersymmetric Minkowski minimum with $`T_1=T_2=T_3=a_1S/q_1`$. The fluxes contribute $`(3a_1q_1/2)>0`$ to the (3.8) RR tadpoles.
#### 4.1.2 Superpotentials depending on five or more moduli
With five fields there are solutions of $`W=0`$ and $`W/\mathrm{\Phi }=0`$. This is our next example.
iv) Example SM-4
Consider the superpotential
$$W=(q_3T_2+q_2T_3)T_1(b_{22}T_2+b_{32}T_3)U_2(b_{23}T_2+b_{33}T_3)U_3.$$
(4.7)
Observe that the non-zero $`b_{ij}`$ trivially satisfy the Jacobi identities (2.23). If the fluxes satisfy $`q_2b_{22}=q_3b_{32}`$, $`q_2b_{23}=q_3b_{33}`$ and $`q_2q_3<0`$, there is a solution with $`|q_2|t_3=|q_3|t_2`$. There is also a relation $`q_3t_1=b_{22}u_2+b_{23}u_3`$. To have $`t_1>0`$ for any $`u_2`$, $`u_3>0`$, we need $`q_2b_{22}>0`$ and $`q_2b_{23}>0`$. Hence, the flux piece in (3.10) is negative (same as D6-branes). In this particular case fluxes contribute $`2b_{22}q_2`$ to the last two RR tadpoles in (3.10).
As we have said, for the given class of $`W`$’s, in which the metric fluxes must satisfy the Jacobi identities (2.23), there cannot be supersymmetric Minkowski solutions when $`W`$ depends on more than five fields. To see this, first observe that without loss of generality we can always take the three $`T_i`$ to be among the fields in $`W`$. Next, from (3.18) we see that $`W/\stackrel{~}{U}_K=0`$ implies $`A_{IJ}\stackrel{~}{T}_I=0`$, and taking real part, $`A_{iJ}t_i𝒜_{Ji}t_i=0`$, where $`J`$ takes three or four values. Note that the $`t_i`$ correspond to the kernel of the metric flux matrix $`𝒜`$. Thus to have solutions with $`t_i0`$, $`\mathrm{rank}𝒜2`$. After using the Jacobi identities one is left with $`\mathrm{rank}𝒜=1`$. One can then check that the number of fields in $`W`$ can be at most five. However, we will see that for $`W`$ depending on all seven moduli, there are supersymmetric AdS minima in which $`D_\mathrm{\Phi }W=0`$ but $`W0`$.
As all examples so far show, fluxes allowing $`𝒩`$=1 supersymmetric Minkowski vacua only fix the moduli partially. Recall that in the type IIB case there are supersymmetric Minkowski minima with all complex structure fields fixed (but not the Kähler moduli). This is because, as already mentioned, in IIB there are extra superpotential couplings $`\tau \tau _i\tau _j`$ or $`\tau \tau _1\tau _2\tau _3`$, whose naive IIA mirror $`ST_iT_j`$, $`ST_1T_2T_3`$ we do not have. This situation changes when $`W`$ depends on the seven moduli. As we will see, supersymmetric AdS minima with all real moduli determined can then occur. This will be discussed in sections 4.3 and 4.4.
### 4.2 No-scale models in Minkowski space
As ‘no-scale’ we distinguish models in which the superpotential is independent of three of the moduli, so that the form (2.15) of the Kähler potential guarantees a cancellation of the cosmological constant . In fact, the scalar potential is positive definite and it is minimized with respect to all fields when $`D_\mathrm{\Phi }^{}W(\mathrm{\Phi }^{})=0`$. Since in general $`W0`$, supersymmetry is broken by the F-terms of the fields not appearing in $`W`$.
One easily finds no-scale minima with superpotentials depending on the moduli sets $`\{U_I,T_1,T_2,T_3\}`$ or $`\{U_I,U_j,T_k,T_{\mathrm{}}\}`$, with $`U_I=S,U_i`$. There are no minima coming from superpotentials of the form $`W(U_I,U_j,U_k,T_{\mathrm{}})`$ or $`W(S,U_1,U_2,U_3)`$. Unlike the previous case with $`𝒩`$=1 supersymmetry, one can find no-scale models with and without metric fluxes. In the latter case one necessarily has $`m0`$. As in the $`𝒩`$=1 supersymmetric examples, the moduli are typically only partially fixed and the fluxes contribute to RR tadpoles like D6-branes. In particular, if $`W`$ only depends on $`S,T_i`$, the fluxes only contribute to the (3.8) RR tadpoles. If it depends on one or two $`U_i`$’s, fluxes may contribute to other RR tadpoles, but always positively. We can consider a superpotential $`W=W(S,T_1,T_2,T_3)`$ as generic, since one can always replace $`S`$ or one or two $`T_i`$’s by $`U_i`$’s if appropriate fluxes are also exchanged.
Our task is then to look for solutions of $`D_SW=0`$ and $`D_{T_i}W=0`$, without imposing $`W=0`$, starting with the $`W(S,T_1,T_2,T_3)`$ given in eq.(3.5). We remind the reader that this superpotential may be obtained performing T-dualities on the superpotential $`𝒲`$ of type IIB generated by certain $`\overline{}_3`$ and $`\overline{}_3`$ fluxes. Now, in the type IIB case, solving $`\mathrm{D}_\tau 𝒲=0`$, $`\mathrm{D}_{\tau _i}𝒲=0`$, is equivalent to demanding that the flux $`\overline{𝒢_3}=(\overline{}_3\tau \overline{}_3)`$ be imaginary self dual (ISD) . Indeed, e.g. the conditions (4.13), (4.14) and (4.15) below amount to the ISD requirement.
#### 4.2.1 Examples with no metric fluxes
i) Example NS-1
In the case with no metric fluxes one can check that there are minima only if $`m0`$, $`h_00`$, and $`\gamma _1=\gamma _2=\gamma _3=0`$, where the $`\gamma _i`$ are the combination of RR fluxes
$$\gamma _i=me_i+q_jq_k;ijk.$$
(4.8)
The superpotential has the form
$$W=e_0+ih_0S+i\underset{i}{}e_iT_iq_1T_2T_3q_2T_1T_3q_3T_1T_2+imT_1T_2T_3.$$
(4.9)
For the moduli we determine $`v_i=q_i/m`$, $`\mathrm{Im}S=(e_0m^2q_1q_2q_3)/h_0m^2`$, whereas for the real parts $`h_0s=mt_1t_2t_3`$. Hence, $`h_0m>0`$ and again the flux contribution to tadpoles is positive. We also find that at the minimum the superpotential is $`W_0=2ih_0s`$ and the gravitino mass is $`m_{3/2}^2=h_0m/32u_1u_2u_3`$. Note that this class of background is mirror to an analogous class of no-scale models in type IIB discussed in .
#### 4.2.2 Examples with metric fluxes
With a superpotential of type $`W(S,T_1,T_2,T_3)`$ this amounts to having $`a_i0`$ for some $`i`$. Let us consider some simple models giving rise to a no-scale structure.
ii) Example NS-2
Let us study first the case in which $`m=0`$ and one of the $`a_i`$ vanishes, say $`a_1=0`$. We will then be able to translate the results for, say $`W(S,U_1,T_2,T_3)`$ if we wish. The superpotential has the form
$$W=e_0+ih_0S+i\underset{i}{}e_iT_iS(a_2T_2+a_3T_3)q_1T_2T_3q_2T_1T_3q_3T_1T_2.$$
(4.10)
If $`m=0`$, then $`t_2,t_30`$ require $`q_1=0`$. One finds that the axions are then fully determined to be
$$\mathrm{Im}S=\frac{e_2q_2e_3q_3}{q_2a_2q_3a_3};v_1=\frac{e_3a_2e_2a_3}{q_2a_2q_3a_3};v_2=\frac{h_0q_2e_1a_3}{q_2a_2q_3a_3};v_3=\frac{e_1a_3h_0q_3}{q_2a_2q_3a_3}.$$
(4.11)
There is a further relation $`e_0=h_0\mathrm{Im}S+e_1v_1`$. The real parts verify $`q_2t_1t_3=a_2st_2`$ and $`q_3t_1t_2=a_3st_3`$. Thus we must have $`a_2q_2>0`$ and $`a_3q_3>0`$, indicating a positive contribution to the (3.8) RR tadpoles. On the other hand only pairwise ratios of real moduli are fixed, namely
$$s^2=\frac{q_2q_3}{a_2a_3}t_1^2;t_3^2=\frac{a_2q_3}{a_3q_2}t_2^2.$$
(4.12)
At the minimum, $`W_0=2s(a_2t_2+a_3t_3)`$.
In a variant of this model one can further set $`a_2=0`$ and for consistency $`q_2=0`$. The imaginary parts are obtained substituting these values in (4.11). For the real parts it only follows that $`a_3st_3=q_3t_1t_2`$. In yet another variant one can set $`a_2q_2=a_3q_3`$. The imaginary parts are then given as in (4.4), while $`a_3t_3=\pm a_2t_2`$ and $`a_2s=\pm q_3t_1`$. This allows either $`a_2a_3>0`$ or $`a_2a_3<0`$ (so we could further impose $`W_0=0`$ as in the model SM-1).
iii) Example NS-3
Let us now consider the case with $`m0`$, still with $`a_1=0`$. One finds that in order to have a solution the fluxes must verify
$$\gamma _2=\frac{a_2\gamma _3}{a_3};h_0\gamma _3=a_3(e_1q_1+me_0).$$
(4.13)
For the imaginary parts we obtain
$$\mathrm{Im}S=\frac{me_0+q_1e_1}{mh_0};v_1=\frac{q_1}{m};v_2=\frac{q_2}{m}+\frac{a_2st_2}{mt_1t_3};v_3=\frac{q_3}{m}+\frac{a_3st_3}{mt_1t_2}.$$
(4.14)
The real parts instead satisfy
$$a_2a_3s^2=\gamma _1t_1^2;m^2t_1^2t_2^2t_3^2+a_2^2s^2t_2^2+a_3^2s^2t_3^2=(h_0m+a_2q_2+a_3q_3)st_1t_2t_3.$$
(4.15)
This shows that $`(h_0m+a_2q_2+a_3q_3)>0`$ and hence the flux contribution to tadpoles is again positive. Notice that the above solution simplifies upon taking $`a_2=0`$ which is consistent if $`\gamma _1=\gamma _2=0`$. In this case
$$t_3=\frac{(h_0m+q_3a_3)st_1t_2}{(a_3s)^2+(mt_1t_2)^2}.$$
(4.16)
We also find that at the mimimum
$$W_0=\frac{2(h_0m+q_3a_3)st_1t_2}{a_3s+imt_1t_2}.$$
(4.17)
The gravitino mass turns out to be
$$m_{3/2}^2=\frac{(h_0m+q_3a_3)}{32u_1u_2u_3}.$$
(4.18)
As expected for a no-scale model it only depends on the $`u_i`$.
iv) Example NS-4
Consider now an isotropic setup with $`a_1=a_2=a_3=a`$ and $`q_1=q_2=q_3=q`$. It is possible to obtain vacua with the four moduli fixed . At the minimum, $`\mathrm{Re}T_k=t`$, $`\mathrm{Im}T_k=v`$. After taking $`e_I=0`$ we find the equations
$`v[2amv^2(h_0m3aq)v2h_0q]`$ $`=`$ $`0,`$
$`(h_03av)(q+mv)`$ $`=`$ $`amt^2.`$ (4.19)
The dilaton is given by $`aS=(q+mv)T^{}iqv`$. There is a solution with $`v=0`$ which, since then $`T=\sqrt{h_0q/ma}`$, $`aS=qT`$, can occur only if $`h_0m>0`$ and $`aq>0`$. A solution with $`v0`$ can happen if $`(aq/h_0m)<1`$ and $`h_0m<0`$, or if $`1/9<(aq/h_0m)<0`$ and $`h_0m>0`$. In all cases the flux contribution to the RR tadpole is positive.
### 4.3 AdS vacua without metric fluxes
We now consider the superpotential, depending on all seven moduli, given by the sum of the two terms in (2.14). As we will see, as long as $`m0`$, $`T_i`$, $`s`$ and $`u_i`$ can be fixed in AdS, but only a linear combination of the axions $`\mathrm{Im}S`$ and $`\mathrm{Im}U_i`$ is determined. This can occur both for broken and unbroken supersymmetry but in all cases we prove that the contribution of fluxes to RR tadpoles is always positive.
We find that in the absence of metric fluxes the vacuum structure is determined by the combination of fluxes $`\gamma _i`$, c.f. (4.8). In particular, $`\gamma _i<0`$ is required to solve $`D_{T_i}W=0`$, $`D_SW=0`$, and $`D_{U_i}W=0`$ in the supersymmetric case. Then, the Kähler axions are fixed as $`v_i=q_i/m`$, whereas for the other axions
$$h_0\mathrm{Im}S\underset{i}{}h_i\mathrm{Im}U_i=\frac{1}{m^2}[e_0m^2q_1q_2q_3+\underset{i}{}q_i\gamma _i].$$
(4.20)
The real parts are instead determined to be
$$t_1=\sqrt{\frac{5|\gamma _2\gamma _3|}{3m^2|\gamma _1|}};t_2=\frac{\gamma _1t_1}{\gamma _2};t_3=\frac{\gamma _1t_1}{\gamma _3};s=\frac{2\gamma _1t_1}{3mh_0};u_i=\frac{2\gamma _1t_1}{3mh_i}.$$
(4.21)
Observe that in order to have $`s>0`$, $`u_k>0`$, it must be that $`mh_0>0`$ whereas $`mh_k<0`$. Hence, the flux contribution to the tadpoles is positive in (3.8) and negative in (3.10). Note that in the present case only $`m`$, $`h_0`$ and $`h_k`$ are restricted by RR tadpole conditions while the fluxes $`e_0`$, $`c_1`$ and $`c_2`$ are essentially unconstrained. This will allow us to find minima at arbitrarily large volume and small dilaton, see below. Type IIA supersymmetric AdS vacua without metric fluxes have been recently addressed in . We obtain similar results.
To go beyond supersymmetric minima and find all solutions of $`V/\mathrm{\Phi }=0`$ we will analyze the case $`T_1=T_2=T_3=T`$, so that $`W`$ is given in (3.19) with $`b_k=a=0`$. The vacuum structure now depends on $`\gamma =mc_1+c_2^2`$. In particular, there exists a supersymmetric AdS minimum only if $`\gamma <0`$. We also find that necessarily $`mh_0>0`$ and $`mh_k<0`$. Therefore, the flux contribution to the tadpoles is positive in (3.20) and negative in (3.21).
To summarize the results we use the shorthand $`\mathrm{Re}T=t`$ and $`\mathrm{Im}T=v`$. The extremum only fixes one linear combination of the imaginary parts of dilaton and complex structure fields which is given by
$$h_0\mathrm{Im}S\underset{k=1}{\overset{3}{}}h_k\mathrm{Im}U_k=e_03c_1v3c_2v^2+mv^3.$$
(4.22)
There are two branches for $`v`$, namely,
$$v_s=\frac{c_2}{m};v_{ns}=\frac{c_2\pm \sqrt{\gamma m^2t^2/2}}{m}.$$
(4.23)
For each value of $`v`$ there are various sub-branches according to the relation among the real parts of $`S`$ and the $`U_k`$. From now on we just look at solutions with $`h_1u_1=h_2u_2=h_3u_3`$. In this case there are two sub-branches characterized by
$`(\mathrm{I})`$ $`:`$ $`h_ku_k=h_0s;k=1,2,3`$
$`(\mathrm{II})`$ $`:`$ $`h_ku_k=h_0smt^3.`$ (4.24)
In the $`v_s`$ sub-branch I,
$$m^2t^2=\pm \frac{5}{3}\gamma ;h_0s=\frac{2}{5}mt^3;\mathrm{\Lambda }_s=\frac{\gamma ^2}{24m^2su_1u_2u_3t}.$$
(4.25)
For $`\gamma <0`$ this is the AdS supersymmetric minimum. For $`\gamma >0`$ it is a non-supersymmetric AdS extremum with same data for the moduli and the cosmological constant. In the $`v_s`$ sub-branch II,
$$m^2t^2=\pm \frac{5}{\sqrt{6}}\gamma ;h_0s=\frac{4}{5}mt^3;\frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_s}=\frac{32}{27}\left(\frac{6}{9}\right)^{1/4}1.071.$$
(4.26)
Both $`\gamma <0`$ and $`\gamma >0`$ are allowed. In either case it is a non-supersymmetric AdS extremum.
The $`v_{ns}`$ branch can occur only if $`\gamma >0`$. There are two non-supersymmetric AdS sub-branches according to (4.24). Their data are
$`(\mathrm{I})`$ $`:`$ $`m^2t^2={\displaystyle \frac{4}{3}}\gamma ;h_0s={\displaystyle \frac{2\gamma t}{3m}};{\displaystyle \frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_s}}={\displaystyle \frac{25\sqrt{5}}{48}}1.165,`$
$`(\mathrm{II})`$ $`:`$ $`m^2t^2={\displaystyle \frac{196}{99}}\gamma ;h_0s={\displaystyle \frac{14\gamma t}{9m}};{\displaystyle \frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_s}}={\displaystyle \frac{11^45^23^2\sqrt{55}}{2^47^7\sqrt{3}}}1.070.`$
This ends our list of non-supersymmetric AdS vacua.
Note that in all these examples without metric fluxes the fixed moduli scale with respect to the RR 4-form and 2-form fluxes $`c_1`$ and $`c_2`$ as
$$ts^{1/3}u_k^{1/3}\gamma ^{1/2}c_1^{1/2},c_2,$$
(4.28)
for large fluxes. Thus the compactification volume may be arbitrarily large for large $`c_1`$ and/or $`c_2`$. For large fluxes, the four- and ten-dimensional dilatons behave as
$$e^{\varphi _4}c_1^{3/2},c_2^3;e^\varphi c_1^{3/4},c_2^{3/2}$$
(4.29)
so that the vacua lie in a perturbative regime for sufficiently large RR 4-form and/or 2-form fluxes. Concerning the cosmological constant, one can check that for large fluxes $`c_1`$ and $`c_2`$ it scales as
$$V_0\gamma ^{9/2}c_1^{9/2},c_2^9.$$
(4.30)
Thus, for large $`c_1`$/$`c_2`$ the c.c. goes with the string dilaton like $`e^{6\varphi }`$. The density of RR fluxes is also suppressed. As pointed out in , to compute this density a factor of $`g_s=e^\varphi `$ must be included. Then, the flux density of $`\overline{F}_4`$ ($`\overline{F}_2`$) behaves like $`c_1^{3/2}`$ ($`c_2^3`$) for large $`c_1`$ ($`c_2`$) fluxes.
### 4.4 Supersymmetric AdS vacua with metric fluxes
We finally treat the full superpotential given in (3.12). We will see that all real moduli may be fixed at the minimum. Concerning the imaginary parts, some linear combinations of $`\mathrm{Im}S`$, $`\mathrm{Im}U_i`$ remain massless but, as we will discuss in the next section, in the presence of D6-branes those massless fields are in fact necessary for certain (potentially anomalous) brane $`U(1)`$ fields to get a Stückelberg mass. We will also see that to get these minima certain discrete relationships among the fluxes must be fulfilled.
There are two classes of models depending on whether $`m=0`$ or not. In the former case one finds that fluxes in general contribute to all RR tadpole directions with a sign which is opposite to that of D6-branes. This is important since it offers an alternative to orientifold planes to cancel RR tadpoles. In the second case with $`m0`$ one finds the interesting result that, depending on different flux choices, always including metric fluxes, the sign of the contribution to RR tadpoles may be arbitrary and the net contribution may vanish. In the latter case one has a cancellation of a positive RR-NS contribution $`h_0m`$ with a RR-metric flux contribution of type $`a_iq_i`$. This is interesting because in this class of backgrounds all real moduli are determined but the fluxes are unconstrained by RR tadpole cancellation conditions.
We will examine the case $`T_k=T`$ and look for supersymmetric minima for any $`m`$. From $`D_{U_k}W=0`$ and $`D_SW=0`$, with $`W`$ given in (3.19), we find
$$3as=b_ku_k.$$
(4.31)
Hence, $`a`$ and $`b_k`$ must be both non-zero and of the same sign. Moreover, there are consistency conditions
$$3h_ka+h_0b_k=0;k=1,2,3.$$
(4.32)
Therefore, either both $`h_0`$ and $`h_k`$ vanish or both are non-zero and of opposite sign. These conditions do not involve the moduli so at most we will have five equations for six unknowns, i.e. we will have at least one flat direction for the supersymmetric minima. In fact, only a combination of complex structure axions is fixed as
$$3a\mathrm{Im}S+\underset{k=1}{\overset{3}{}}b_k\mathrm{Im}U_k=3c_1+\frac{3c_2}{a}(3h_07av)\frac{3m}{a}v(3h_08av).$$
(4.33)
If $`h_0,h_k0`$ using (4.32) we can write the fixed axion combination as $`h_0\mathrm{Im}S_kh_k\mathrm{Im}U_k`$. We also find
$$as=2t(c_2mv).$$
(4.34)
Except for some axion directions, we have thus determined all the moduli in terms of $`T`$ which is found from the remaining equations. The solution depends on whether $`m`$ is different from zero or not. We now specify to these two possibilities.
i) $`m=0`$. Examples AdS-1
When $`m=0`$ we find the simple results
$$v=\frac{h_0}{3a};9c_2t^2=e_0\frac{h_0c_1}{a}\frac{h_0^2c_2}{3a^2}.$$
(4.35)
At the minimum, $`W_0=12c_2t^2`$ and the cosmological constant turns out to be
$$V_0=\mathrm{\Lambda }=\frac{ab_1b_2b_3}{128c_2^2t^3},$$
(4.36)
where $`t`$ is given in (4.35). It is important to notice that (4.34) fixes $`c_2a>0`$ so that in the supersymmetric minima with $`m=0`$ the metric fluxes give a negative contribution to the tadpole in (3.20). Similarly, $`c_2b_k>0`$ and the flux contribution to the tadpoles in (3.21) is positive. This is the result we advertised.
Let us now check whether we have enough freedom to locate all moduli at large volume and small dilaton so that one can trust the effective 4-dimensional field theory approximation being used. The fluxes of type $`a,b_k`$ and $`c_2`$ are constrained by the RR tadpole cancellation conditions and by the extra conditions (4.32). The values of $`h_0`$ and $`h_k`$ are constrained by the latter but in principle both may be large as long as $`h_0/h_k=3a/b_k`$. On the other hand, the fluxes of the RR 6-form ($`e_0`$) and 4-form ($`c_1`$) are unconstrained and may be arbitrarily large. Note then that for large $`e_0`$ and $`c_1`$ the moduli fields behave all like
$$tsu_ke_0^{1/2},c_1^{1/2}.$$
(4.37)
In order for our vacuum to remain in a perturbative regime we would like to have small values for the 4-dimensional coupling $`e^{\varphi _4}`$ and the 10-dimensional string coupling $`e^\varphi `$. They are found to be
$`e^{\varphi _4}`$ $`=`$ $`(su_1u_2u_3)^{1/4}=t^1{\displaystyle \frac{(ab_1b_2b_3)^{1/4}}{23^{3/4}c_2}},`$
$`e^\varphi `$ $`=`$ $`e^{\varphi _4}t^{3/2}=t^{1/2}{\displaystyle \frac{(ab_1b_2b_3)^{1/4}}{23^{3/4}c_2}}.`$ (4.38)
We thus see that for large $`t`$ (which may be obtained e.g. with a large 6-form flux $`e_0`$) the 4-dimensional dilaton is small. However the string dilaton grows with $`t`$. Only by appropriately choosing the fluxes, i.e. with large $`c_2`$ one can perhaps maintain it under control. On the other hand such fluxes are in general very much constrained by the RR tadpole conditions so it seems difficult having small string dilaton and large volume at the same time. We will see however that in the case with $`m0`$ one can easily stabilize the moduli in the perturbative regime.
ii) $`m0`$. Examples AdS-2
To deal with $`m0`$ it is convenient to introduce a new variable for $`\mathrm{Im}T`$. If $`h_00`$ we use
$$v=(\lambda +\lambda _0)\frac{h_0}{3a};\lambda _0=\frac{3c_2a}{mh_0}.$$
(4.39)
The value of $`\lambda `$ follows from the cubic equation
$$160\lambda ^3+186(\lambda _01)\lambda ^2+27(\lambda _01)^2\lambda +\lambda _0^2(\lambda _03)+\frac{27a^2}{mh_0^3}(e_0ac_1h_0)=0.$$
(4.40)
Clearly, we need a real solution for $`\lambda `$ and moreover, such that $`\lambda (\lambda +\lambda _01)>0`$ because now $`t`$ is determined from
$$3a^2t^2=5h_0^2\lambda (\lambda +\lambda _01).$$
(4.41)
Notice also that (4.34) takes the form $`3a^2s=2h_0m\lambda t`$. For the cosmological constant we find
$$V_0=\frac{ab_1b_2b_3\lambda _0^2(16\lambda +\lambda _01)}{1920c_2^2t^3\lambda ^3},$$
(4.42)
where $`\lambda `$ is the appropriate solution of (4.40) and $`t`$ is given in (4.41).
There is a variety of cases depending on the values of the different fluxes. One of the interesting features when $`m0`$ is that the contribution to the RR tadpoles may have either sign and even vanish. In fact, the flux-induced tadpoles in (3.20) and (3.21) are respectively $`\frac{1}{2}h_0m(1\lambda _0)`$ and $`\frac{1}{2}h_km(1\lambda _0)`$. Thus, the flux tadpoles vanish at the special value $`\lambda _0=1`$. This is important, as we mentioned above, since it allows to fix the moduli without any constraint from RR tadpole cancellations.
To analyze the equations that determine $`\lambda `$ and $`t`$ we can proceed in various ways. We could choose for example $`e_0`$ and $`c_1`$ and study the allowed values of $`\lambda _0`$. For instance, with $`e_0=c_1=0`$, we find that to have solutions for $`\lambda `$ with acceptable $`t`$ necessarily $`\frac{1}{3}<\lambda _0<3`$. We also need $`mh_0<0`$ and $`mh_k>0`$ so that $`s>0`$ and $`u_k>0`$. The special value $`\lambda _0=1`$ at which tadpoles vanish is allowed and leads in turn to
$$t=\sqrt{\frac{5}{3}}\lambda \left|\frac{h_0}{a}\right|;s=\frac{2mh_0\lambda }{3a^2}t;u_k=\frac{2mh_0\lambda }{ab_k}t.$$
(4.43)
From the cubic equation we find the value $`\lambda =(10)^{2/3}/20`$.
As we advanced, with $`m0`$ one can locate the minima in perturbative regions. Consider for instance the case $`e_0=c_1=0`$ and $`\lambda _0=1`$ so that the real moduli are given in (4.43). Note that one can have $`h_0`$, $`h_k`$ and $`c_2`$ arbitrarily large as long as $`\lambda _0=1`$ and eq.(4.32) is respected. Then one can check that
$$e^{\varphi _4}h_0^2;e^\varphi h_0^{1/2},$$
(4.44)
so that for sufficiently large $`h_0`$ the minima will be perturbative. Note also that the NS flux density is diluted for large fluxes since it goes like $`h_0/t^{3/2}h_0^{1/2}`$. Concerning the RR flux $`\overline{F}_2`$, its density also goes like $`h_0^{1/2}`$ for large fluxes, taking the factor of $`g_s`$ into account . The cosmological constant eq.(4.42) scales like
$$V_0h_0^5(e^\varphi )^{10}$$
(4.45)
and hence is substantially suppressed for large $`h_0`$. Similar results may be obtained for values of $`\lambda _0`$ sufficiently close to 1, which would allow for contributions to RR tadpoles with either sign and of arbitrary size. In section 6 we will consider this possibility to construct a semi-realistic intersecting D6-brane model with all diagonal closed string moduli stabilized.
Let us mention for completeness other solutions within this class of AdS minima with $`m0`$. We may start by choosing a preferred value for some of the moduli. For example, we can set $`v=0`$, and $`h_0\mathrm{Im}S_kh_k\mathrm{Im}U_k=0`$. Then, necessarily $`c_1=3h_0c_2/a`$ and, from the cubic equation with $`\lambda =\lambda _0`$, $`e_0=45h_0c_2^2/ma`$. This is the solution found in .
Another way to proceed with the analysis is to fix $`\lambda _0`$. For example, we can take $`c_2=0`$ so that $`\lambda _0=0`$. Obviously, $`t^2>0`$ then requires either $`\lambda <0`$ or $`\lambda >1`$. If $`\lambda <0`$, then $`s>0`$, $`u_k>0`$ demand $`h_0m>0`$ and $`h_km<0`$, thus the flux contribution to tadpoles is like that of D6-branes. It is more interesting to consider $`\lambda >1`$ so that $`h_0m<0`$ and $`h_km>0`$. Furthermore, to satisfy (4.40), it must be $`(e_0ac_1h_0)>0`$. Were it not for the fact that the fluxes are integers, we could always find solutions for some chosen $`\lambda `$. But still there is room to adjust the fluxes. For example, for $`\lambda =3/2`$ we just need $`(e_0ac_1h_0)=6h_0^3m/a^2`$ and this could be verified say for $`e_0=0`$, $`c_1=24m`$ and $`h_0=2a`$.
We can also set $`h_0=h_k=0`$, but to this end a different parametrization of $`v`$, amounting to $`\lambda \lambda _0\widehat{\lambda }`$, must be employed. Now the interesting case is $`\widehat{\lambda }<1`$ because $`s>0`$, $`u_k>0`$ require $`c_2a>0`$ and $`c_2b_k<0`$ so that the flux contribution to tadpoles could cancel that of D6-branes. Again we can choose some $`\widehat{\lambda }`$ and find values of $`e_0`$ to satisfy the cubic equation for $`v`$. For instance, for $`\widehat{\lambda }=3/2`$ we need $`e_0m^2=161c_2^3`$. One can check however that in this and the previous solution it is again hard to achieve at the same time a large value for the volume and a small value for the 10-dimensional dilaton, the reason being that now the value of fluxes $`h_0`$, $`c_2`$ and $`h_k`$ will be constrained by RR tadpole cancellation conditions.
One can also easily find non-susy AdS vacua. We will just show a particularly simple example. In general, there are solutions in which (4.31) and (4.32) are still satisfied. To go further let us set $`m=0`$. Then there are solutions with $`as=2c_2t`$ and $`v=h_0/3a`$, but with the novelty that
$$\pm 9c_2t^2=e_0\frac{h_0c_1}{a}\frac{h_0^2c_2}{3a^2}.$$
(4.46)
With plus sign this is the supersymmetric minimum, but we can also choose the minus sign depending on the fluxes. For instance, if $`e_0=c_1=0`$, only the non-supersymmetric choice is available. In this case the minimum is AdS and it is typically stable because the eigenvalues of the Hessian are positive or negative but above the Breitenlohner-Freedman bound .
## 5 D6-branes, fluxes and the Freed-Witten anomaly
We are going to consider now adding D6-branes to a IIA background with fluxes turned on. Besides the general RR tadpole cancellation constraints already mentioned, a number of points do also change. One apparent puzzle is the following. In the world-volume of a generic stack of D6-branes there is a $`U(1)`$ gauge field whose scalar partner parametrizes the D6-brane position in compact space. These $`U(1)`$’s often get Stückelberg masses by swallowing RR scalars and disappear from the low-energy spectrum. At the same time these scalars participate in the cancellation of $`U(1)`$ gauge anomalies through a variation of the Green-Schwarz mechanism . Let us specify now to the toroidal case considered in detail in the present paper. Consider a $`\mathrm{D6}_a`$ wrapping a factorizable torus with wrapping numbers as in eq. (3.7). Then the $`U(1)_a`$ field couples to RR fields in $`D`$=4 as follows
$$F^aN_a\underset{I=0}{\overset{3}{}}c_I^aC_I^{(2)}$$
(5.1)
with
$$c_0^a=m_a^1m_a^2m_a^3;c_1^a=m_a^1n_a^2n_a^3;c_2^a=n_a^1m_a^2n_a^3;c_3^a=n_a^1n_a^2m_a^3.$$
(5.2)
Here the $`C_I^{(2)}`$ are 2-forms which are Poincaré duals in $`D`$=4 to the $`\mathrm{Im}U_I`$ fields considered above. In terms of them the couplings have a Higgs-like form
$$A_\mu ^a^\mu (c_0^a\mathrm{Im}Sc_1^a\mathrm{Im}U_1c_2^a\mathrm{Im}U_2c_3^a\mathrm{Im}U_3).$$
(5.3)
We thus observe that certain linear combinations of $`\mathrm{Im}U_I`$ ($`U_0=S`$) fields get a mass by combining with open string vector bosons living on the branes. Moreover, these linear combinations of $`\mathrm{Im}U_I`$ fields will transform with a shift under $`U(1)_a`$ gauge transformations, like Goldstone bosons do. On the other hand, we have seen above that NS and metric backgrounds give rise to terms in the superpotential linear in $`\mathrm{Im}U_I`$, i.e.
$$W_Q=_Y\mathrm{\Omega }_c(\overline{H}_3+dJ_c)=\underset{I,J=0}{\overset{3}{}}A_{IJ}\stackrel{~}{T}_I\stackrel{~}{U}_J.$$
(5.4)
Such terms generically may give rise to potential terms for the $`\mathrm{Im}U_I`$ fields which would not be invariant under the shifts induced by $`U(1)_a`$ gauge transformations. The condition to restore consistency and gauge invariance would be to impose the constraint
$$_{\mathrm{\Pi }_a}(\overline{H}_3+\omega J_c)=0,$$
(5.5)
evaluated at the appropriate vacuum. Here $`\mathrm{\Pi }_a`$ denotes the 3-cycle wrapped by the D6-brane, $`\omega `$ are the metric fluxes and $`J_c`$ is the complexified Kähler 2-form of the torus. We have used $`dJ_c=\omega J_c`$, according to (2.21).
In absence of metric fluxes eq. (5.5) may be understood in terms of the Freed-Witten anomaly <sup>1</sup><sup>1</sup>1We thank A. Uranga for pointing out this connection to us.. A simple way to see the origin of this constraint starting from type IIB is as follows . Consider a D3-brane wrapping a 3-cycle $`\mathrm{\Pi }`$ through which there is some quantized NS flux $`\overline{H}_3`$. On the world-volume of the D3-brane there is a CS coupling of the form $`C_2F_2`$, $`C_2`$ being the RR 2-form and $`F_2`$ the open string gauge field strength. After performing a IIB S-duality transformation one gets a coupling
$$_{\mathrm{\Pi }\times 𝐑}H_3\stackrel{~}{A}_1,$$
(5.6)
where $`\stackrel{~}{A}_1`$ is the gauge field dual to the $`A_1`$ form living on the D3-brane. This shows that a background for $`H_3`$ gives rise to a tadpole for $`\stackrel{~}{A}_1`$ and hence to an inconsistency. Performing three T-dualities one expects for D6-branes the analogous term
$$_{D6}H_3\stackrel{~}{A}_4.$$
(5.7)
The resulting tadpole is avoided if $`_{\mathrm{\Pi }_a}\overline{H}_3=0`$. Equation (5.5) should thus be the extension of this constraint to the case including metric fluxes $`\omega 0`$.
Note that, ignoring for the moment the effect of metric fluxes, eq. (5.5) implies that all D6-branes wrapping the orientifold should obey
$$\underset{I=0}{\overset{3}{}}c_I^ah_I=0.$$
(5.8)
This is in general a strong constraint on the possible D6-branes one may add in specific models with flux, as we discuss below in specific examples. Remarkably, this constraint guarantees that combinations of axions getting masses by mixing with vector bosons are orthogonal to those becoming massive from fluxes, the latter being typically of the form $`h_0\mathrm{Im}S_kh_k\mathrm{Im}U_k`$.
Another interesting point is the following. We have seen in section 4.4 that adding fluxes one can fix the torus moduli in a supersymmetry-preserving AdS minimum. Now, for non-zero NS fluxes one finds at the minima that $`h_i/h_0=s/u_i`$. Substituting this in eq. (5.8) and multiplying by the torus volume one arrives at
$$m_a^1m_a^2m_a^3(R_y^1R_y^2R_y^3)m_a^1n_a^2n_a^3(R_y^1R_x^2R_x^3)n_a^1m_a^2n_a^3(R_x^1R_y^2R_x^3)n_a^1n_a^2m_a^3(R_x^1R_x^2R_y^3)=0.$$
(5.9)
This condition means that the D6-brane wraps a special Lagrangian cycle (slag). From the effective Lagrangian point of view this is proportional to a Fayet-Iliopoulos (FI) term and hence it imposes dynamically that the D6-brane configuration should be supersymmetric (i.e. all FI terms should vanish). Thus one concludes that, in this class of AdS supersymmetric minima the constraint (5.8) implies that the brane configuration should be also supersymmetric. Notice that including metric fluxes in this class of minima does not add extra constraints to be satisfied due to the relations (4.32).
The condition (5.9) in these AdS vacua is in fact
$$\mathrm{Im}\mathrm{\Omega }|_{\mathrm{\Pi }_a}=0.$$
(5.10)
It arises from $`_{\mathrm{\Pi }_a}\overline{H}_3=0`$ because at the AdS minimum $`\overline{H}_3\mathrm{Im}\mathrm{\Omega }`$. In turn, in our setup this is a simple consequence of $`h_i/h_0=s/u_i`$. This is also found in a more general analysis of type IIA susy AdS vacua . Likewise, if there are metric fluxes, $`dJ_c\mathrm{Im}\mathrm{\Omega }`$. For instance, in the models of section 4.4 this can be deduced using (4.31). Therefore, even if the NS fluxes vanish, in these models there is still a FW constraint of the form
$$3ac_0^ab_1c_1^ab_2c_2^ab_3c_3^a=0.$$
(5.11)
This guarantees that combinations of axions acquiring a mass from fluxes or from $`U(1)`$ mixing are orthogonal to each other
Recently models of type IIB orientifolds with fluxes and intersecting (or rather magnetized) D-branes with semi-realistic spectrum have been constructed . Some of them do not verify the (IIB version of) constraint (5.8) and hence would be in principle inconsistent. This possible problem with the FW anomaly was already pointed out in where it was suggested that it could be cured if additional D-branes were included. In the case of IIA orientifolds under consideration we would need to add D4-branes hanging between different sets of D6-branes and their orientifold mirrors. It may be argued that the chiral spectra from intersecting D6-brane models does not get affected by the presence of these extra D4-branes. However, no specific construction with this possible cancellation mechanism has been presented in the literature. In addition, it is not clear whether in the case of supersymmetric D6-brane configurations the addition of the extra D-branes does not spoil supersymmetry. Given this fact, it seems sensible to impose the constraint (5.8) on specific models with Minkowski vacua and that will be our approach below. In AdS vacua, in which the real parts of all moduli are determined, the FW anomaly cancels automatically if the brane configuration is supersymmetric.
## 6 Intersecting D6-brane models in the presence of fluxes
We have shown in previous sections that the addition of fluxes in type IIA theory leads to new properties not present in analogous IIB models. Some of the aspects we have found with potential model-building applications are: 1) fluxes may contribute to all four RR tadpoles, 2) one can have examples of fluxes fixing part or all of closed string moduli but not contributing to RR tadpoles, and 3) there are models with metric fluxes (as well as other NS and RR fluxes) in which one can obtain AdS supersymmetric vacua with all moduli stabilized and contribution to RR tadpoles opposite to that of D6-branes. In addition to these properties, since plenty of flux variables do not contribute to RR tadpoles, there is substantial freedom in the choice of the parameters of the vacua and in particular one can obtain minima at large volume and small dilaton values, in which the approximations inherent to a 4-dimensional effective Lagrangian approach hold.
To illustrate the possible applications to model-building of these results in previous sections we are going to consider here specific intersecting D6-branes models with semi-realistic spectrum. The first two examples correspond to Minkowski vacua both with unbroken $`𝒩`$=1 supersymmetry and with broken supersymmetry but no-scale structure. Although in these cases only some of the closed string moduli are fixed at the minima, it is interesting to consider them since other effects could perhaps stabilize the rest of the moduli. In these two cases the models will be left-right symmetric extensions of the MSSM, with gauge group $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}\times U(1)`$, rather than the MSSM. The third example has an AdS supersymmetric background and is particularly interesting since, to our knowledge, is the first semi-realistic three-generation model with all closed string moduli stabilized. In this case also the gauge group is closer to that of the MSSM, since it is that of the SM with some additional $`U(1)`$’s. We consider models with non-supersymmetric intersecting branes and all closed string moduli fixed in appendix B.
### 6.1 Minkowski MSSM-like
In the class of type IIB orientifold models with fluxes studied up to now, it has been shown that flux backgrounds with Minkowski geometry, either $`𝒩`$=1 supersymmetric or not, lead to positive contributions to RR tadpoles. This stems from the fact that ISD fluxes always contribute to RR tadpoles as D-branes do. In building semi-realistic models this leads to problems with RR tadpole cancellation conditions, since typically fluxes contribute too much to tadpoles. It was pointed out in that this problem may be cured if appropriate additional D9-anti-D9-brane pairs contributing negatively to some of the RR tadpoles are added. In any case, full cancellation of RR tadpoles in realistic toroidal models require considering orbifold generalizations like $`_2\times _2`$ . Semi-realistic $`𝒩`$=1 supersymmetric type IIB $`_2\times _2`$ orientifolds with flux backgrounds have been studied in . The class of models of has a brane content as given in table 1.
In the case of the IIB $`_2\times _2`$ orientifold the $`(n,m)`$ integers would be magnetic numbers whereas in the T-dual IIA orientifold they correspond to wrapping numbers along horizontal and vertical directions of each $`\mathrm{T}^2`$ in the factorized $`\mathrm{T}^6`$ respectively. Note that this set contains as a subset the MSSM-like model introduced in . We assume as in that the $`b`$ and $`c`$ D6-branes sit on top of the orientifold plane so that the corresponding gauge symmetries are enhanced to $`SU(2)_L`$ and $`SU(2)_R`$ respectively. The full initial gauge group is then $`U(4)\times SU(2)_L\times SU(2)_R\times [U(1)_1\times U(1)_2]`$. Separating one of the $`a`$-branes from the other three produces the breaking $`U(4)U(3)\times U(1)`$. Furthermore, two out of the three $`U(1)`$’s get a Stückelberg mass by combining with RR axion fields. We are thus left with a gauge group $`SU(3)_c\times SU(2)_L\times SU(2)_R\times U(1)_{BL}\times [U(1)]`$, which contains the left-right symmetric extension of the SM with an extra $`U(1)`$. The branes $`a,b,c`$, give rise to a 3-generation MSSM-like spectrum whereas the additional branes $`h_{1,2}`$ in table 1 are used to help in cancelling the RR tadpoles.
Note that in the case of the $`_2\times _2`$ IIA orientifold the RR tadpole cancellation conditions in the presence of fluxes will have the form
$`{\displaystyle \underset{a}{}}N_an_a^1n_a^2n_a^3+{\displaystyle \frac{1}{2}}(h_0m+a_1q_1+a_2q_2+a_3q_3)`$ $`=`$ $`16,`$
$`{\displaystyle \underset{a}{}}N_an_a^1m_a^2m_a^3+{\displaystyle \frac{1}{2}}(mh_1q_1b_{11}q_2b_{21}q_3b_{31})`$ $`=`$ $`16,`$
$`{\displaystyle \underset{a}{}}N_am_a^1n_a^2m_a^3+{\displaystyle \frac{1}{2}}(mh_2q_1b_{12}q_2b_{22}q_3b_{32})`$ $`=`$ $`16,`$ (6.1)
$`{\displaystyle \underset{a}{}}N_am_a^1m_a^2n_a^3+{\displaystyle \frac{1}{2}}(mh_3q_1b_{13}q_2b_{23}q_3b_{33})`$ $`=`$ $`16.`$
where the $`(16)`$ in the last three conditions is the RR tadpole contribution of the other 3 orientifold planes existing in the $`_2\times _2`$ case. Note that the branes $`h_1`$ and $`h_2`$ contribute negatively to all four RR tadpoles so that in principle one can use them to compensate for a too large contribution to the first tadpole condition from fluxes. Precisely this was the approach in ref. (see also ). Here we will use this class of models as our starting point for the IIA orientifold case. Here are some possibilities:
i) A 3 generation $`𝒩`$=1 MSSM-like model with some fixed moduli
Consider the above model in which we turn on non-vanishing fluxes as in one of the susy Minkowski examples of section (4.1) with non-vanishing $`b_{31},b_{21},a_2,a_3`$ (example SM-2). The addition of NS fluxes $`h_0,h_1`$ and RR $`e_0,e_2,e_3`$, is optional, but we set all the remaining backgrounds to zero. The superpotential has then the form
$$W=T_2(a_2S+b_{21}U_1)T_3(a_3S+b_{31}U_1)+e_0+ih_0Sih_1U_1+ie_2T_2+ie_3T_3.$$
(6.2)
As explained in section 4.1 this has a Minkowski supersymmetric minimum with
$$h_0=a_2v_2+a_3v_3;e_2=a_2\mathrm{Im}S+b_{21}\mathrm{Im}U_1;t_3=\frac{b_{21}t_2}{b_{31}};s=\frac{b_{21}u_1}{a_2}.$$
(6.3)
as long as $`e_2a_3=e_3a_2`$, $`h_0b_{31}=a_3h_1`$, $`h_0b_{21}=a_2h_1`$ and $`h_0e_2=e_0a_2`$. Thus in this supersymmetric Minkowski background two complex linear combinations of moduli are fixed at the minimum.
Note that, since $`m=q_i=0`$, in this background the fluxes do not contribute to the RR tadpole. Thus one can consider the addition of D6-branes as in the the case with $`N_f=5`$ in table 1. As pointed out in with this choice all RR tadpoles cancel without the addition of fluxes in type IIB theory. In the present IIA case we can rather add the background considered here and the RR tadpoles are not modified and hence cancel. However the moduli are partially fixed by eq.(6.3).
It is easy to check that the $`a,b`$ and $`c`$ branes where the SM lives trivially satisfy the FW constraint. However the branes of type $`h_{1,2}`$ may be problematic unless:
$$a_2(m_a^1m_a^2m_a^3)b_{21}(m_a^1n_a^2n_a^3)=a_212b_{21}=0$$
(6.4)
which on the other hand may be easily satisfied by appropriately choosing $`a_2`$, $`b_{21}`$. Note that this condition guarantees that the linear combination of $`\mathrm{Im}S`$, $`\mathrm{Im}U_1`$ getting masses through fluxes (eq.(6.3)) is orthogonal to the linear combination getting masses by mixing with the $`U(1)`$’s of branes $`h_{1,2}`$.
Note that in the IIB version of this orientifold with fluxes considered in , the latter contributes to RR tadpoles and one can only get a one-generation $`𝒩`$=1 supersymmetric model.
ii) A 3 generation no-scale model
One can also consider one of the no-scale backgrounds discussed in section (4.2), the variant of the NS-2 model, and include a set of D6-branes as in table 1. A simple example is as follows. Take non-vanishing $`a_3,q_3`$ with the remaining $`q_i=a_i=0`$. In addition one may include non-vanishing $`h_0,e_0,e_i`$ but set the remaining backgrounds to zero. The superpotential has then the form
$$W(S,T_i)=a_3ST_3q_3T_1T_2+e_0+ih_0S+i\underset{i}{}e_iT_i.$$
(6.5)
The imaginary part of $`S`$ and the $`T_i`$ are fixed as in eq.(4.11) whereas one has for the real parts the relationship $`a_3st_3=q_3t_1t_2`$. In addition one has the constraint $`e_0=h_0\mathrm{Im}S+e_1v_1`$. There is only a contribution equal to $`\frac{1}{2}a_3q_3`$ to the first RR tadpole. We consider fluxes quantized in units of 8 to avoid problems with flux quantization . One can then cancel tadpoles in a $`_2\times _2`$ orientifold with branes as in table 1 with $`N_f=1`$ and $`a_3=q_3=8`$.
One can also consider a no-scale model with a non-vanishing IIA mass parameter $`m`$ and with no metric fluxes, as described at the beginning of subsection 5.2. One takes non-vanishing $`m`$ and $`h_0`$. In addition one can have non-vanishing $`e_i,q_j`$, verifying $`\gamma _i=me_i+q_jq_k=0`$ ($`ijk`$). Setting $`h_0=m=8`$ and $`N_f=1`$ one cancels all tadpoles. Note that this model, which has no metric fluxes, is the IIA mirror of a similar no-scale model considered in .
One can check however that both these no-scale models as they stand have FW anomalies. The danger comes from the $`h_{1,2}`$ branes which have a non-vanishing product $`m_a^1m_a^2m_a^30`$. One possibility which might cure this problem is if, as suggested in , the brane $`h_1`$ recombines with the mirror of $`h_2`$ into a single (non-factorizable) D6-brane $`h_1+h_2^{}`$. One can in fact claim that this is the generic situation for branes like these which do intersect. In this case, since $`h_1`$ and $`h_2^{}`$ have equal and opposite $`m_a^1m_a^2m_a^3`$, the FW would cancel on the recombined brane. On the other hand it is not clear whether after the addition of fluxes a flat direction in the effective potential exists corresponding to the recombination of those branes. In the $`𝒩`$=1 supersymmetric AdS model which we describe next no such problem appears.
### 6.2 A $`𝒩`$=1 MSSM-like model with all closed string moduli stabilized in AdS
The previous intersecting brane models were able to combine a semi-realistic spectrum with a partial determination of some closed string moduli. We now show that all such moduli may be stabilized in the case of AdS vacua, thus providing, to our knowledge, the first semi-realistic string model with all closed string moduli stabilized at weak coupling. Note first that in the past it has been argued that it is impossible to construct semi-realistic $`𝒩`$=1 supersymmetric intersecting D6-brane models wrapping the IIA orientifold $`T^6/(\mathrm{\Omega }(1)^{F_L}I_3)`$. The reason for this was essentially the impossibility to cancel the 4 RR tadpole conditions simultaneously while maintaining supersymmetry. To obtain $`𝒩`$=1 supersymmetric models extra orbifold twisting (e.g. $`_2\times _2`$, as in previous examples) had to be added, giving rise to extra orientifold planes which help in the cancellation of RR tadpoles . We will show here that one can build $`𝒩`$=1 supersymmetric configurations in the purely toroidal orientifold in which those RR tadpoles may be cancelled by the addition of NS/RR and metric fluxes. The role played by additional orientifold planes in orbifold (e.g. $`_2\times _2`$) models is here played by the additional fluxes which contribute like orientifold planes. At the same time those fluxes stabilize all closed string moduli in AdS space. Moreover the complex structure moduli are fixed at values which render the D6-brane configuration supersymmetric. Notice that in the $`𝒩`$=1 supersymmetric models previously considered in the literature those moduli where not determined by the dynamics.
Let us consider the set of D6-branes wrapping factorizable cycles in the orientifold as in table 2.
Note that this set only differs from the previous examples in the form of the additional branes $`h_1,h_2`$. Another difference is that in our IIA case we have a purely toroidal (no $`_2\times _2`$) orientifold without further twisting. The corresponding chiral spectrum at the intersections is given in table 3.
In the table a prime indicates the mirror brane. The gauge group after separating branes and after two of $`U(1)`$’s get Stückelberg masses is $`SU(3)\times SU(2)_L\times U(1)_R\times U(1)_{BL}\times [U(1)\times SU(3)^2]`$. Note that, unlike the case of the $`_2\times _2`$ models above, one can make the breaking $`SU(2)_RU(1)_R`$ by brane splitting, and hence the gauge group is that of the MSSM supplemented by some extra $`U(1)`$’s. We have three generations of quarks and leptons, one Higgs multiplet $`H`$ and extra matter fields involving the auxiliary branes $`h_1`$, $`h_2`$ and $`o`$.<sup>2</sup><sup>2</sup>2One can check that if the branes $`h_1`$ and $`h_2^{}`$ recombine, most of the extra matter beyond the SM disappears from the massless spectrum, with only additional $`SU(2)_{L,R}`$ doublets remaining.
With this brane content (plus the mirrors) the RR tadpole cancellation conditions are
$`64+{\displaystyle \frac{1}{2}}(h_0m+a_1q_1+a_2q_2+a_3q_3)`$ $`=`$ $`16,`$
$`4+{\displaystyle \frac{1}{2}}(h_1mq_1b_{11}q_2b_{21}q_3b_{31})`$ $`=`$ $`0,`$
$`4+{\displaystyle \frac{1}{2}}(h_2mq_1b_{12}q_2b_{22}q_3b_{32})`$ $`=`$ $`0,`$ (6.6)
$`4+{\displaystyle \frac{1}{2}}(h_3mq_1b_{13}q_2b_{23}q_3b_{33})`$ $`=`$ $`0.`$
We see that to cancel tadpoles the sign of the flux contribution must be opposite to that of D6-branes. We will now consider a AdS background with metric fluxes and $`m0`$ discussed in section (4.4). The reader can check that choosing the fluxes as
$$q_i=q=h_i2;a_i=16;m=b_{ij}=b_{ii}=4,$$
(6.7)
all RR tadpoles are cancelled. Note that eq.(4.32) fixes $`h_0=12h_i`$, otherwise the values of $`q`$, $`h_0`$ and $`h_i`$ may be arbitrarily large still cancelling all RR tadpoles.
The above type of flux backgrounds does give rise to supersymmetric AdS vacua with all real moduli fixed. In fact, the fluxes in (6.7) are isotropic so that the superpotential is of the form (3.19) and the results of section 4.4 with $`m0`$ can be applied. One can easily check that with the above fluxes $`\lambda _0=1+(24/h_0)`$, which is arbitrarily close to 1 for large $`h_0`$. Substituting these fluxes yields for the real moduli
$$s=\frac{h_0\lambda }{96}t;u_k=\frac{h_0\lambda }{8}t;t=\sqrt{\frac{5}{3}}\frac{|h_0|}{16}\lambda ^{1/2}(\lambda +\frac{24}{h_0})^{1/2},$$
(6.8)
where $`\lambda `$ is the appropriate solution of eq.(4.40) for the $`\lambda _0`$ indicated above. For large $`h_0`$, $`\lambda _0`$ is close to 1 so that $`\lambda (10)^{2/3}/20`$ when $`e_0=c_1=0`$. In this case one needs $`h_0<0`$. The imaginary part of the Kähler moduli are fixed as in eq.(4.39) whereas only the linear combination of dilaton and complex structure axions $`12\mathrm{I}\mathrm{m}S+_{k=1}^3\mathrm{Im}U_k`$ is fixed, as in eq.(4.33). As discussed in section 4.4, for large $`h_0`$ (which also implies large $`h_k`$, $`q`$) all moduli are stabilized in a regime in which perturbation theory in $`D`$=4 is a good approximation.
Note that the FW conditions (5.8) for the D6-branes $`a,h_1`$ and $`h_2`$ read respectively
$$h_2=h_3;h_1=h_3;h_1=h_2,$$
(6.9)
which are automatically satisfied because $`h_1=h_2=h_3=h_0/12`$. As we mentioned, this will guarantee that the the supersymmetry preserving conditions at the brane intersections
$`\mathrm{tg}^1\left({\displaystyle \frac{\tau _2}{3}}\right)\mathrm{tg}^1\left({\displaystyle \frac{\tau _3}{3}}\right)`$ $`=`$ $`0,`$
$`\mathrm{tg}^1\left({\displaystyle \frac{\tau _1}{2}}\right)\mathrm{tg}^1\left({\displaystyle \frac{\tau _2}{2}}\right)`$ $`=`$ $`0,`$ (6.10)
$`\mathrm{tg}^1\left({\displaystyle \frac{\tau _1}{2}}\right)\mathrm{tg}^1\left({\displaystyle \frac{\tau _3}{2}}\right)`$ $`=`$ $`0,`$
where $`\tau _i=R_y^i/R_x^i`$, are satisfied, since $`u_1=u_2=u_3`$. This is no surprise, since as we mentioned in section 5, in this class of AdS vacua all branes should be calibrated which in turn implies that the FW anomaly automatically cancels.
In this particular model it is also interesting to look at the structure of $`U(1)`$’s and the $`\mathrm{Im}U_I`$ RR fields. It is easy to check that the couplings (5.1) give masses to two linear combinations of $`U(1)`$’s by combining with certain linear combinations of $`\mathrm{Im}U_I`$ fields. Only the generator $`Q_a2(Q_1Q_2)`$ remains massless at this level. On the other hand, the fields $`\mathrm{Im}S`$ and $`_k\mathrm{Im}U_k`$ do not mix with the $`U(1)`$’s at all, as expected, since FW anomalies cancel. Note that the combination $`12\mathrm{I}\mathrm{m}S+_k\mathrm{Im}U_k`$ is the one which gets a mass from fluxes (see eq.(4.33)). The orthogonal linear combination is massless and may be identified with an axion which may be of relevance for the strong CP problem.
Although we have studied here only the dilaton and the diagonal closed string moduli of the orientifold, we already mentioned that setting all off-diagonal moduli to zero solves the extremum conditions. Furthermore, since we are in a $`𝒩`$=1 supersymmetric AdS background, this guarantees that these off-diagonal moduli are also stable. Thus, the closed string background discussed is completely stable. We have then succeeded in building the first semi-realistic $`𝒩`$=1 supersymmetric model with all closed string moduli stabilized in a consistent perturbative regime. The vacuum is AdS with a c.c. which may be made small (although not arbitrarily small, see below) for large fluxes. Unlike previous flux constructions in the present case we have a simple toroidal orientifold, without any further orbifold twist. Furthermore, the $`𝒩`$=1 supersymmetry conditions on the brane angles are forced upon us by the Freed-Witten constraint plus the minimization. Recall in this respect that in the $`𝒩`$=1 supersymmetric brane configurations constructed up to now the angles were fine-tuned to verify the supersymmetry conditions, there was no dynamical explanation for that choice, since not all closed string moduli were fixed.
Let us make some complementary comments about this kind of MSSM-like AdS constructions:
i) Other MSSM-like models in AdS may be constructed along similar lines making use of the backgrounds with metric fluxes and $`m0`$ discussed in subsection 4.4. An easy way to proceed is to start with a tadpole free $`𝒩`$=1 MSSM-like D6-brane configuration and embed it in a AdS background of the type $`\lambda _0=1`$ in which fluxes do not contribute to RR tadpoles at all. For example, one can start again from the $`_2\times _2`$ orientifold example in table 1 with $`N_f=5`$ in which all RR tadpoles cancel. In the prescribed AdS background all moduli are stabilized in a perturbative regime for large enough fluxes.
ii) One could think of building analogous AdS MSSM-like models with a background free of metric fluxes, as in subsection 4.3. This turns out to be difficult because fluxes contribute negatively to the RR tadpoles in eq.(3.21) and hence additional orientifold planes have to be added wrapping those directions. One possibility is to use again the $`_2\times _2`$ orientifold which has such O6-planes, but the fluxes tend to overwhelm the contribution of orientifolds and this procedure does not look promising.
iii) Besides the chiral spectrum described above, this class of toroidal models has massless adjoint chiral fields corresponding to the open string moduli parametrizing the location and Wilson lines on the branes. In a supersymmetric AdS background as the one we are considering here, those open string moduli are in any case stable. It is an interesting question to study what happens to them when some supersymmetry-breaking effect is included. It has been shown that fluxes in type IIB stabilize some (but not all) of the open string moduli in the toroidal case . Additional ways to give masses to these degrees of freedom in toroidal models have been recently described in . It should be worth to study this question in the context of our type IIA AdS backgrounds.
iv) Once all moduli are fixed in a given model like this, one can compute a number of interesting physical quantities like gauge coupling constants and Yukawa couplings, since they will be known functions of the fluxes. For example, in the above model the gauge kinetic functions of the groups $`SU(4)`$, $`SU(2)_L`$, $`SU(2)_R`$ have
$$\mathrm{Re}f_{SU(4)}=9s+u_1;\mathrm{Re}f_{SU(2)_L}=\frac{u_2}{2};\mathrm{Re}f_{SU(2)_R}=\frac{u_3}{2}.$$
(6.11)
Since these are the values at the string scale, to make contact with experiment we should then consider the running to low energies. As we said, simplest toroidal models like this have, in addition to the chiral spectrum, adjoint chiral fields which will generically spoil the running of coupling constants. Let us nevertheless proceed and compute them in this example. Since $`u_k`$ and $`s`$ are related by (4.31) one has e.g.
$$\alpha _{SU(2)_L}=\frac{8}{\lambda |h_0|t}=\frac{128\sqrt{3}}{\sqrt{5}h_0^2\lambda ^{3/2}(\lambda +24/h_0)^{1/2}},$$
(6.12)
and $`\alpha _{SU(4)}=\frac{2}{7}\alpha _{SU(2)_L}`$. Thus, we see that the SM gauge coupling constants depend strongly on the fluxes. In particular we cannot make $`h_0`$ arbitrarily large (as one would naively do to decrease the value of the c.c.) since we would get then too small SM gauge couplings, inconsistent with experiment. For example, in the present model one can see that in order to get values $`\alpha _i1/51/30`$ which might be consistent with low energy physics one needs to have $`h_0100`$ but not much bigger. This seems to be a generic property and not a particular feature of this class of models. Thus, indeed we have an infinite ‘landscape’ of models depending on unconstrained fluxes, but only a narrow region of fluxes would lead to consistent low-energy physics. Something similar happens with the Yukawa couplings, which have been computed and can be neatly written in terms of products of Jacobi $`\vartheta `$-functions in this model . They scale like the gauge couplings and hence are equally suppressed for large fluxes.
v) These models are constructed in AdS and an obvious question is how one could promote this kind of vacua to dS. One possibility which comes to mind is to add anti-D6-branes. Indeed, if we add a pair of $`\mathrm{D6}\overline{\mathrm{D6}}`$ branes to the model, there is an extra contribution to the scalar potential which has the form
$$V_{\overline{D6}}\frac{1}{u_1u_2u_3}$$
(6.13)
and should be included in the complete minimization. This would be essentially the mirror of the approach in which was used in the type IIB case. More generally, one may consider sets of D6-branes with uncancelled NS-tadpoles. A potential is generated due to the missancellation of the tensions of the D6-branes against the orientifold tension, i.e. (in the string frame)
$$V_{\mathrm{D6}/\mathrm{O6}}=\frac{T_6}{g_s}\left(\underset{a}{}N_al_al_{ori}\right)>0$$
(6.14)
where $`l_a`$ $`(l_{ori})`$ are the volume of the 3-cycles wrapped by each D6-brane (orientifold). It remains to be seen whether such a procedure can be made to work. In any event, the kind of fluxes considered here will stabilize non-supersymmetric D6-brane configurations with non-vanishing NS-tadpoles that have been considered in recent years. Examples of such non-supersymmetric D6-brane configurations with the chiral content of the SM are presented in Appendix B.
## 7 Conclusions
In this paper we have studied the minima of the flux-induced effective moduli potential in a simple $`\mathrm{T}^6/(\mathrm{\Omega }(1)^{F_L}I_3)`$ IIA orientifold. We have focused on the dilaton and the diagonal Kähler and complex structure fiels, but we have nevertheless argued that the results found ignoring off-diagonal moduli still provide stable extrema of the full potential in relevant cases. We have considered RR, NS as well as metric background fluxes. Unlike the IIB case, the richness of the flux options leads to a full stabilization of all closed string moduli in AdS without the need of non-perturbative effects. Furthermore, the RR tadpole conditions, which are very restrictive in the IIB case, only constraint some flux combinations in the IIA case. Thus, there is enough freedom to adjust fluxes so that the minima are located in regions with large volume and small dilaton where the effective 4-dimensional supergravity approximations hold. The combination of metric fluxes with NS/RR fluxes leads to new possibilities such as fluxes fixing all moduli in $`𝒩`$=1 supersymmetric AdS but not contributing to RR tadpoles. This provides us with a rigid ‘corset-like’ background which can stabilize any RR tadpole-free D6-brane configuration in this toroidal setting. In general, if metric fluxes are turned on, the overall fluxes can contribute to RR tadpoles like O6-planes do, thereby providing the interesting possibility of disposing of orientifold planes in some cases.
In models with all real moduli fixed, only one linear combination of the axions of the dilaton and complex structure fields is determined at the minima. This, which at first sight appears to be a limitation of the approach, is in fact a blessing. Indeed, eventually we may like to add systems of D6-branes leading to chiral physics in the background. The RR axions which are never fixed by fluxes are in fact needed by the D6-branes to get rid of (potentially anomalous) open string $`U(1)`$’s. We have seen that cancellation of the Freed-Witten anomaly guarantees that sufficient axions remain to give Stückelberg masses to the $`U(1)`$’s. In the case of AdS $`𝒩`$=1 supersymmetric AdS vacua the cancellation of FW anomaly does in turn force the different sets of D6-branes to be calibrated.
One can construct explicit models with a chiral spectrum quite close to that of the MSSM with three generations and with all closed string moduli fixed in AdS. In its construction we make use of fluxes (including metric ones) contributing like orientifold planes to RR tadpoles. Other analogous models may also be built. The minima may be located at large volume and small dilaton so that we can trust our approximations. In such a model, with all moduli fixed, one can compute explicitly all gauge and Yukawa couplings as known functions of the fluxes undetermined by RR tadpoles. In the particular example of section 6.2, essentially only one flux (which may be identified with the NS 3-form flux $`h_0`$) fixes all couplings and scales. Thus, although one may talk about a landscape of models depending on a single flux parameter $`h_0`$, only a narrow region of integer values for $`h_0`$ would give rise to gauge couplings compatible with experimental constraints. In particular, in our concrete model there is not enough freedom to make the c.c. arbitrarily small by making $`h_0`$ large. We believe that this is quite a generic feature. The dilaton and complex structure fields which determine the (inverse of) gauge coupling constants grow like some power of the fluxes. If we make the fluxes too large in order to get e.g. a small c.c. the SM couplings would get far too small.
Our approach has been to consider the metric fluxes as a deformation added to the original torus. The resulting twisted torus is a non-Calabi-Yau manifold in which we still know the moduli and are able to introduce D6-branes. It would be interesting to go beyond the toroidal geometry in this spirit. In appendix A we have shown that the analysis of $`𝒩`$=1 vacua deduced from the effective flux-induced superpotential agrees with recent results on supersymmetric IIA compactifications on manifolds with $`SU(3)`$ structure . Although we have worked out many specific minima of the general fluxed potential in the simplest IIA toroidal orientifold, we cannot claim that we have done a complete analysis. Furthermore, we have not explored the possibilities offered by some of the solutions (e.g. AdS non-supersymmetric minima) that we have analyzed nor made a systematic search for MSSM-like D6-brane configurations. Presumably there are many other options beyond the ones that were discussed. All the models with all moduli stabilized are however AdS. An important problem is how to modify the premises in order to obtain models with dS vacua. A possible option is to add anti-D6-branes or, more generally, consider non-supersymmetric brane configurations. Positive definite contributions to the potential will then in general appear which might help in going to dS. We hope to come back to all these issues in the near future.
Acknowledgments
We thank J.F.García-Cascales, F. Marchesano, S. Theisen, and especially A. Uranga for useful discussions. A.F. thanks the Max-Planck-Institut für Gravitationsphysik for hospitality while preparing this paper. The work of P.G.C. is supported by the Ministerio de Educación y Ciencia (Spain) through a FPU grant. This work has been partially supported by the European Commission under the RTN European Program MRTN-CT-2004-503369 and the CICYT (Spain).
## Appendix A: $`SU(3)`$ structure of twisted torus
In this appendix we study the relation between the metric fluxes and the $`SU(3)`$ structure of the twisted torus. The idea is to generalize the analysis of by turning on all metric fluxes in (2.22) and not only the $`a_i`$ obtained from T-duality of NS fluxes. We will see that in this more general situation the twisted torus is still a half-flat manifold as it occurs when only the $`a_i`$ are present . Using the results in this appendix we will also be able to describe our Minkowski and AdS supersymmetric vacua in terms of torsion classes. Supersymmetric IIA compactifications on manifolds with $`SU(3)`$ structure have been recently considered in .
On the twisted torus one can build the fundamental 2-form $`J`$ and the holomorphic 3-form $`\mathrm{\Omega }`$ in the usual way. Including the sizes and complex structure parameters we have
$`J`$ $`=`$ $`t_1\eta ^1\eta ^4t_2\eta ^2\eta ^5t_3\eta ^3\eta ^6,`$
$`\mathrm{\Omega }`$ $`=`$ $`(\eta ^1+i\tau _1\eta ^4)(\eta ^2+i\tau _2\eta ^5)(\eta ^3+i\tau _3\eta ^6).`$ (A.1)
These forms define an $`SU(3)`$ structure. In particular, they satisfy
$$J\mathrm{\Omega }=0;JJJ=\frac{3i}{4}\frac{t_1t_2t_3}{\tau _1\tau _2\tau _3}\mathrm{\Omega }\mathrm{\Omega }^{}.$$
(A.2)
The torsion classes can be read from (see e.g. )
$`dJ`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{t_1t_2t_3}{\tau _1\tau _2\tau _3}}\mathrm{Im}(𝒲_1\mathrm{\Omega }^{})+𝒲_4J+𝒲_3,`$
$`d\mathrm{\Omega }`$ $`=`$ $`𝒲_1JJ+𝒲_2J+𝒲_5^{}\mathrm{\Omega },`$ (A.3)
where $`𝒲_1`$ is a complex 0-form, $`𝒲_2`$ is a primitive ($`𝒲_2JJ=0`$) complex 2-form, $`𝒲_3`$ is a primitive ($`𝒲_3J=0`$) real $`(2,1)(1,2)`$-form, $`𝒲_4`$ is a real 1-form, and $`𝒲_5`$ is a complex $`(1,0)`$-form. The unusual factor in the first term of $`dJ`$ is needed so that $`d(J\mathrm{\Omega })=0`$.
When only the metric fluxes (2.22) are turned on, using (2.16) we find
$`d\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{1}{s}}(a_1s+b_{11}u_1+b_{12}u_2+b_{13}u_3)\eta ^{2536}+{\displaystyle \frac{1}{s}}(a_2s+b_{21}u_1+b_{22}u_2+b_{23}u_3)\eta ^{1436}`$
$`+{\displaystyle \frac{1}{s}}(a_3s+b_{31}u_1+b_{32}u_2+b_{33}u_3)\eta ^{1425},`$
$`dJ`$ $`=`$ $`(a_1t_1+a_2t_2+a_3t_3)\eta ^{456}(b_{13}t_1+b_{23}t_2+b_{33}t_3)\eta ^{126}`$
$`(b_{12}t_1+b_{22}t_2+b_{32}t_3)\eta ^{153}(b_{11}t_1+b_{21}t_2+b_{31}t_3)\eta ^{423},`$
where $`\eta ^{153}=\eta ^1\eta ^5\eta ^3`$, etc.. In $`d\mathrm{\Omega }`$ we have used $`u_i=s\tau _j\tau _k`$, $`ijk`$.
Clearly, $`d(JJ)=0`$ and $`d(\mathrm{Im}\mathrm{\Omega })=0`$, thus the twisted torus with the given fluxes is a half-flat manifold. For the torsion classes we easily read $`𝒲_4=𝒲_5=0`$. The torsions $`𝒲_2`$ and $`𝒲_3`$ are different from zero, they can be readily obtained using (A.3) and
$`𝒲_1`$ $`=`$ $`{\displaystyle \frac{1}{6st_1t_2t_3}}(a_1st_1+a_2st_2+a_3st_3+b_{11}t_1u_1+b_{12}t_1u_2+b_{13}t_1u_3`$ (A.5)
$`+b_{21}t_2u_1+b_{22}t_2u_2+b_{23}t_2u_3+b_{31}t_3u_1+b_{32}t_3u_2+b_{33}t_3u_3).`$
Notice that this $`𝒲_1`$ is similar to $`W_Q`$ without the NS fluxes, as expected because it is basically computed as $`\mathrm{\Omega }dJ`$.
In it has been shown that supersymmetric Minkowski vacua of type IIA require $`𝒲_1=0`$. In our examples of this kind of vacua in section 4.1 we indeed find $`𝒲_1=0`$. This follows simply because taking real part of $`W/\stackrel{~}{U_J}=0`$ gives $`_iA_{iJ}t_i=0`$ which is enough to show $`𝒲_1=0`$. In fact, $`dJ=0`$ so that $`𝒲_3=0`$ as well. We also find that $`m=0`$ and then from the real part of $`W/T_i=0`$ we deduce that $`d\mathrm{\Omega }=\frac{1}{s}\overline{F}_2J`$. We have examples, such as $`W(S,T_1,T_2,T_3)`$, $`W(U_1,T_1,T_2,T_3)`$ or the $`W(T_1,T_2,T_3,U_2,U_3)`$ in (4.7), in which $`d\mathrm{\Omega }0`$ and $`𝒲_2=\frac{1}{s}\overline{F}_2`$. Another characteristic feature of these models with $`\overline{F}_20`$ is the existence of flux tadpoles for some component of $`C_7`$. There are other models, such as $`W(S,U_1,T_2,T_3)`$ or $`W(U_1,U_2,T_1,T_2)`$, in which $`d\mathrm{\Omega }=0`$ and $`C_7`$ tadpoles vanish because $`\overline{F}_2=0`$ is required to have non-zero real parts of the moduli. In all cases, the typical configuration has neither NS fluxes nor RR fluxes for $`F_6`$ ($`e_0=0`$) and $`F_4`$ ($`e_i=0`$). These results are in agreement with the analysis of .
Type IIA supersymmetric AdS compactifications have also been studied in terms of $`SU(3)`$ structures . It is interesting to see how the same type of results follows in our setup. To begin we notice that taking real part of $`D_{\stackrel{~}{U}_J}W=0`$ gives
$$\underset{i=1}{\overset{3}{}}A_{iJ}t_i\stackrel{~}{u}_J=\frac{1}{2}\mathrm{Re}W;J=0,\mathrm{},3.$$
(A.6)
We can use these relations to compute $`dJ`$ and also
$$𝒲_1=\frac{\mathrm{Re}W}{3st_1t_2t_3}.$$
(A.7)
From the explicit $`dJ`$ we further read $`𝒲_3=0`$. To calculate $`d\mathrm{\Omega }`$ we look instead at the real part of $`D_{T_i}W=0`$ and deduce
$$a_is+b_{i1}u_1+b_{i2}u_2+b_{i3}u_3=\frac{\mathrm{Re}W}{2t_i}\underset{jki}{}t_j(q_k+mv_k).$$
(A.8)
Moreover, combining with (A.6) yields the relation
$$t_1t_2(q_3+mv_3)+t_1t_3(q_2+mv_2)+t_2t_3(q_1+mv_1)=\frac{1}{4}\mathrm{Re}W.$$
(A.9)
It is then straightforward to determine $`d\mathrm{\Omega }`$ and from it obtain
$$𝒲_2J=\frac{1}{4}𝒲_1JJ+\frac{m}{s}B_2J\frac{1}{s}\overline{F}_2J,$$
(A.10)
which satisfies $`𝒲_2JJ=0`$ by virtue of (A.9). We conclude that supersymmetric AdS compactifications have $`𝒲_1`$ and $`𝒲_2`$ different from zero but $`𝒲_3=𝒲_4=𝒲_5=0`$, as found in in a more general setup. There is also a particular case in which $`𝒲_2=0`$ . Our example in section 4.4 is of this type. Indeed, using (4.31) gives $`𝒲_1=2a/t^2`$ and $`d\mathrm{\Omega }=𝒲_1JJ`$. We also find $`\overline{H}_3=\sqrt{|h_1h_2h_3/h_0|}\mathrm{Im}\mathrm{\Omega }`$.
## Appendix B: Stabilizing non-susy intersecting D-brane models
In the $`𝒩`$=1 supersymmetric AdS constructions in the main text, we have discussed examples in which all D6-branes preserve the same $`𝒩`$=1 supersymmetry. In this appendix we would like to study the non-supersymmetric class of semi-realistic intersecting D6-brane models of ref. . One of the known problems of these non-susy models is that they are unstable due to the existence of NS tadpoles. These appear from a miscancellation of the tensions of the D6-branes with the orientifold tension. We would like to point out here that those non-susy models may in general become stable in the presence of fluxes. We will see that the FW conditions in AdS will force the branes to preserve supersymmetry locally, i.e. any pair of intersecting D6-branes will preserve one unbroken supersymmetry although there is no overall $`𝒩`$=1 supersymmetry preserved simultaneously by all D6-branes. We will see that all closed string moduli will be also determined, although, as we argue at the end, a complete treatment would require taking into account D-term’s in the scalar potential.
In a general class of solutions was given for the wrapping numbers $`(n_a^i,m_a^i)`$ giving rise to a SM spectrum. These are shown in table 4. In this table we have several discrete parameters. First we consider $`\beta ^i=1,1/2`$. From the point of view of branes at angles $`\beta ^i=1`$ stands for a rectangular lattice for the $`i^{th}`$ torus, whereas $`\beta ^i=1/2`$ describes a tilted lattice allowed by the $`\mathrm{\Omega }I_3`$ symmetry. We also have two phases $`ϵ,\stackrel{~}{ϵ}=\pm 1`$ and the parameter $`\rho `$ which can only take the values $`\rho =1,1/3`$. Furthermore, each of these families of D6-brane configurations depend on four integers ($`n_a^2,n_b^1,n_c^1`$ and $`n_d^2`$). Any of these choices leads exactly to the same massless fermion spectrum of the SM with 3 generations.
Now, imposing the condition (5.8) for branes $`a`$, $`b`$, $`c`$, $`d`$, respectively, leads to the relations
$`h_2{\displaystyle \frac{ϵ\beta ^2}{\rho \beta ^1}}h_3{\displaystyle \frac{\stackrel{~}{ϵ}n_a^2}{2\beta ^1}}`$ $`=`$ $`0,`$
$`h_1{\displaystyle \frac{ϵ\stackrel{~}{ϵ}\beta ^1}{\beta ^2}}h_3{\displaystyle \frac{\stackrel{~}{ϵ}3\rho n_b^1}{2\beta ^2}}`$ $`=`$ $`0,`$ (B.1)
$`h_3{\displaystyle \frac{n_c^1}{\beta ^2}}`$ $`=`$ $`0,`$
$`h_2{\displaystyle \frac{ϵ\beta ^2}{\rho \beta ^1}}+h_3{\displaystyle \frac{\stackrel{~}{ϵ}3\rho n_d^2}{2\beta ^1}}`$ $`=`$ $`0.`$
These constraints have two solutions, depending on the value of $`n_c^1`$:
i) $`n_c^10`$
In this case necessarily $`h_3=0`$ and hence $`h_1=h_2=0`$. Only the flux $`h_0`$ may be added and only the field $`S`$ may be fixed by fluxes.
ii) $`n_c^1=0`$
In this case one can check that the constraints are solved as long as
$`h_3(n_a^2`$ $`+`$ $`3\rho n_d^2)=0,`$
$`h_1`$ $`=`$ $`{\displaystyle \frac{3\rho n_b^1}{2ϵ\beta ^1}}h_3,`$ (B.2)
$`h_2`$ $`=`$ $`{\displaystyle \frac{\rho \stackrel{~}{ϵ}n_a^2}{2ϵ\beta ^2}}h_3,`$
so that one can only have non-vanishing $`h_i`$ if $`n_a^2=3\rho n_d^2`$. One can check that, when this condition is verified, there are two massless $`U(1)`$’s in the spectrum, $`U(1)_R`$ and $`U(1)_{BL}`$, rather than just hypercharge, and only two linear combinations of the RR fields
$$\mathrm{Im}U_2\frac{\stackrel{~}{ϵ}\rho n_a^2}{2ϵ\beta ^2}\mathrm{Im}U_3;\mathrm{Im}U_1\frac{3\rho n_b^1}{2ϵ\beta ^1}\mathrm{Im}U_3$$
(B.3)
become massive by combining with the $`U(1)_{3B+L}`$ and $`U(1)_b`$ gauge bosons respectively. The orthogonal linear combination
$$\frac{3\rho n_b^1}{2ϵ\beta ^1}\mathrm{Im}U_1+\frac{\stackrel{~}{ϵ}\rho n_a^2}{2ϵ\beta ^2}\mathrm{Im}U_2+\mathrm{Im}U_3=\frac{1}{h_3}(\underset{I=1,2,3}{}h_I\mathrm{Im}U_I)$$
(B.4)
is precisely a piece of the combination appearing in the superpotential, which is expected to acquire a mass from fluxes upon minimization. The other RR field $`\mathrm{Im}S`$ may become massive depending on the presence or not of a non-vanishing $`h_0`$ background, which leads to no constraint in the model (since $`m_a^1m_a^2m_a^3=0`$ for all the D6-branes present).
If upon minimization one finds $`\mathrm{Re}U_I1/h_I`$ one obtains that the above conditions imply
$$\frac{\mathrm{Re}U_3}{\mathrm{Re}U_1}=\frac{3\rho n_b^1}{2ϵ\beta ^1};\frac{\mathrm{Re}U_3}{\mathrm{Re}U_2}=\frac{\rho \stackrel{~}{ϵ}n_a^2}{2ϵ\beta ^2}.$$
(B.5)
One can check that these conditions guarantee that at each brane intersection there is one unbroken supersymmetry, although in this model no overall supersymmetry generator is preserved by all intersections. This kind of local (but not global) supersymmetry was termed ‘Q-SUSY’ in . Thus we see that adding fluxes $`h_i0`$, $`i=1,2,3`$, is only possible if the brane configuration is locally supersymmetric.
Let us now be more explicit and chose the wrapping number parameters as follows:
$$n_a^2=n_d^2=\beta _1=ϵ=\stackrel{~}{ϵ}=1;\beta _2=1/2;\rho =1/3;n_b^1=2.$$
(B.6)
One can easily check that RR tadpoles cancel without the addition of any further D6-brane nor fluxes. The conditions (B.2) now read
$$h_1=h_3=3h_2.$$
(B.7)
Consider now again the same AdS vacua with $`m0`$ that we discussed in section 4.4. We saw there that for $`m0`$ one can find AdS vacua in which the NS/RR contribution $`mh_0`$ to RR tadpoles may be cancelled by the metric fluxes contribution $`3ac_2`$ ($`\lambda _0=1`$ case). Then we have the interesting possibility of fixing all closed string moduli without fluxes contributing to RR tadpoles at all. Consider backgrounds as follows
$$b_k=(6,2,6);h_k=rb_k,$$
(B.8)
where $`r=h_0/3a`$ must be a positive integer. Then, if we further chose $`e_0=c_1=0`$ one finds a minimum of the flux-induced potential as long as $`h_0>0`$, $`a<0`$, $`m<0`$. The flux contribution to tadpoles vanishes if we further take $`c_2=rm`$. Then the real parts of closed string moduli are fixed as
$$u_k=\frac{3as}{b_k};s=\frac{m}{a10^{1/3}}rt;t=\frac{\sqrt{15}10^{2/3}}{20}r.$$
(B.9)
Certain linear combinations of the imaginary parts of the moduli fields are fixed as discussed in subsection 4.4. Note that choosing $`r=h_0/3a`$ large one can fix the moduli at arbitrarily large values with small 4- and 10-dimensional dilatons.
As we mentioned, one has to be careful in applying the results obtained in the main text to a non-supersymmetric brane configuration like this. Indeed, in this case in addition to the F-term scalar potential one has to add the piece (6.14). Still one expects a full determination of all closed string moduli also in this non-supersymmetric example.
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# The Positivity Set of a Recurrence Sequence
## 1 Introduction and Main Results
A sequence $`(f_n)_{n0}`$ of real numbers is called a *recurrence sequence* if it satisfies a linear recurrence
$$f_{n+h}=c_1f_{n+h1}+\mathrm{}+c_{h1}f_{n+1}+c_hf_n,n0,$$
with constant coefficients $`c_k`$. (Since we are concerned with questions of positivity, we restrict attention to real sequences.) One of the most charming and celebrated results in the theory of recurrence sequences is the Skolem-Mahler-Lech theorem. It asserts that the zero set
$$\{n:f_n=0\}$$
of a recurrence sequence is the union of a finite set and finitely many arithmetic progressions. The recent comprehensive monograph by Everest et al. contains references to the substantial literature devoted to this result and to related questions. However, not much seems to be known about the positivity set
$$\{n:f_n>0\}.$$
(1)
In the following section we establish that the density of the set (1) always exists, where the (natural) density of a set $`A`$ is defined as
$$𝒅(A):=\underset{x\mathrm{}}{lim}x^1\mathrm{}\{nx:nA\},$$
provided that the limit exists.
###### Theorem 1.
Let $`(f_n)`$ be a recurrence sequence. Then the density of the set $`\{n:f_n>0\}`$ exists.
Recall that a recurrence sequence can be written as a generalized power sum
$$f_n=\underset{k=0}{\overset{m}{}}P_k(n)\gamma _k^n,n0,$$
(2)
with non-zero polynomials $`P_k[n]`$ and *roots* $`\gamma _k`$ that are roots of the *characteristic polynomial*
$$z^hc_1z^{h1}\mathrm{}c_{h1}zc_h$$
of the recurrence. The roots of largest modulus are called *dominating roots* of $`(f_n)`$.
It does not come as a surprise that recurrence sequences with no positive dominating root have oscillating behaviour. Indeed, in section 3 we prove that for such a sequence $`(f_n)`$ the densities of (1) and the negativity set
$$\{n:f_n<0\}$$
(3)
are always positive. This generalizes the following known result : If a recurrence sequence has at most four dominating roots, and none of them is real positive, then the sets (1) and (3) both have infinitely many elements.
###### Theorem 2.
Let $`(f_n)`$ be a nonzero recurrence sequence with no positive dominating characteristic root. Then the sets $`\{n:f_n>0\}`$ and $`\{n:f_n<0\}`$ have positive density.
In section 4 we investigate which numbers actually occur as density of the positivity set of some recurrence sequence. It turns out that all possible values occur, both for sequences with no dominating positive root and in general.
Finally (section 5), we return to the Skolem-Mahler-Lech theorem. Our approach yields the following weak version: The density of the zero set exists and is a rational number.
The conclusion hints at algorithmic aspects of the positivity of recurrence sequences.
## 2 The Density of the Positivity Set
Notation: We write $`f_n0`$ if $`f_n=0`$ for all $`n0`$. The Lebesgue measure of a set $`B^m`$ is denoted by $`𝝀(B)`$.
The goal of this section is to prove Theorem 1. Dividing $`f_n`$ by $`n^D|\gamma _1|^n`$, where $`\gamma _1`$ is a dominating root of $`f_n`$ and $`D`$ is the maximal degree of the $`P_k`$ with $`|\gamma _k|=|\gamma _1|`$, we obtain from (2)
$$n^D|\gamma _1|^nf_n=\underset{i=1}{\overset{d}{}}a_i\mathrm{cos}(2\pi \theta _in+\beta _i)+vr_n,$$
where $`r_n=\mathrm{O}(1/n)`$ is a recurrence sequence, $`\theta _1,\mathrm{},\theta _d`$ are in $`]0,1[`$ and $`a_i,\beta _i,v`$. From now on we will assume w.l.o.g. $`D=0`$ and $`|\gamma _1|=1`$. As a first step we get rid of any integer relations that the $`\theta _i`$’s might satisfy.
###### Lemma 3.
Let $`\theta _1,\mathrm{},\theta _d`$ be real numbers. Then there is a basis $`\{\tau _1,\mathrm{},\tau _{m+1}\}`$ of the $``$-module
$$M=+\theta _1+\mathrm{}\theta _d$$
such that $`1/\tau _{m+1}`$ is a positive integer and $`1,\tau _1,\mathrm{},\tau _m`$ are linearly independent over $``$.
###### Proof.
$`M`$ is finitely generated and torsion free, hence it is free \[7, Theorem III.7.3\]. Let $`\{\alpha _1,\mathrm{},\alpha _{m+1}\}`$ be a basis. Since $`1M`$, there are integers $`e_1,\mathrm{},e_{m+1}`$ such that
$$e_1\alpha _1+\mathrm{}+e_{m+1}\alpha _{m+1}=1.$$
We complete $`(e_1/g,\mathrm{},e_{m+1}/g)`$, where $`g:=\mathrm{gcd}(e_1,\mathrm{},e_{m+1})`$, to a unimodular integer matrix $`𝐂`$ with last row $`(e_1/g,\mathrm{},e_{m+1}/g)`$ \[7, §XXI.3\]. Then
$$(\tau _1,\mathrm{},\tau _{m+1})^T:=𝐂(\alpha _1,\mathrm{},\alpha _{m+1})^T$$
yields a basis of $`M`$ with $`\tau _{m+1}=1/g`$. Now suppose
$$u_1\tau _1+\mathrm{}+u_m\tau _m=u$$
for integers $`u_1,\mathrm{},u_m,u`$. Since $`u`$ has also the representation
$$ug\tau _{m+1}=u,$$
it follows $`u_1=\mathrm{}=u_m=u=0`$. ∎
Take $`\tau _1,\mathrm{},\tau _{m+1}`$ as in Lemma 3, with $`\tau _{m+1}=1/g`$. Roughly speaking, we have put all integer relations among the $`\theta _i`$ into the rational basis element $`\tau _{m+1}`$. There are integers $`b_{ij}`$ with
$$\theta _i=\underset{j=1}{\overset{m+1}{}}b_{ij}\tau _j.$$
Now we split the sequence $`(f_n)`$ into the subsequences $`(f_{gn+k})_{n0}`$ for $`0k<g`$. We have
$$f_{gn+k}=G_ns_n,$$
where $`s_n:=r_{gn+k}`$ and $`G_n`$ is the dominant part. Defining the integer matrix
$$𝐁:=(gb_{ij})_{\begin{array}{c}1id\\ 1jm\end{array}}^{d\times m}$$
and the real vector $`𝐜=(c_1,\mathrm{},c_d)`$ with
$$c_i:=2\pi k\underset{j=1}{\overset{m+1}{}}b_{ij}\tau _j+\beta _i,1id,$$
it can be written as (cos is applied component wise)
$$G_n=𝐚^T\mathrm{cos}(2\pi n𝐁𝝉+𝐜)+v.$$
We show that the density of $`\{n:f_{gn+k}>0\}`$ exists for each $`k`$. Since $`s_n`$ is a recurrence sequence with fewer characteristic roots than $`f_n`$, we may assume inductively that $`𝒅(\{n:s_n<0\})`$ exists. Thus, if $`G_n`$ is the zero sequence, we are done. Now let $`k`$ be such that $`G_n`$ is not the zero sequence. It is plain that $`G_n=H(n𝝉)`$, where
$$H(𝐭):=𝐚^T\mathrm{cos}(2\pi \mathrm{𝐁𝐭}+𝐜)+v,𝐭[0,1]^m.$$
The following theorem shows that the function $`H`$ can be used to evaluate the density of the positivity set of $`G_n`$, which equals, as we will see below, that of the set $`\{n:f_{gn+k}>0\}`$.
###### Theorem 4 (Kronecker-Weyl).
Let $`\tau _1,\mathrm{},\tau _m`$ be real numbers such that $`1,\tau _1,\mathrm{},\tau _m`$ are linearly independent over $``$. Then for every Jordan measurable set $`A[0,1]^m`$ we have
$$𝒅(\{n:n𝝉mod1A\})=𝝀(A).$$
###### Proof.
We refer to Cassels \[2, Theorems IV.I and IV.II\]. ∎
We define
$$L_\epsilon :=\{n:G_n\epsilon \}\text{and}S_\epsilon :=\{n:|G_n|<\epsilon \}.$$
(4)
The corresponding sets for the function $`H`$ are defined as
$$\stackrel{~}{L}_\epsilon :=\{𝐭[0,1]^m:H(𝐭)\epsilon \}\text{and}\stackrel{~}{S}_\epsilon :=\{𝐭[0,1]^m:|H(𝐭)|<\epsilon \}.$$
Since for all $`\epsilon 0`$
$$L_\epsilon =\{n:n𝝉mod1\stackrel{~}{L}_\epsilon \},$$
we have $`𝒅(L_\epsilon )=𝝀(\stackrel{~}{L}_\epsilon )`$ for all $`\epsilon 0`$ by Theorem 4. Similarly,
$$𝒅(S_\epsilon )=𝝀(\stackrel{~}{S}_\epsilon ),\epsilon >0.$$
(5)
Note that the boundary of the bounded set $`\stackrel{~}{S}_\epsilon `$ (respectively $`\stackrel{~}{L}_\epsilon `$) is a Lebesgue null set (as seen by applying the following lemma with $`F(𝐭)=H(𝐭)\epsilon `$), hence $`\stackrel{~}{S}_\epsilon `$ and $`\stackrel{~}{L}_\epsilon `$ are Jordan measurable, and Theorem 4 is indeed applicable. Lemma 5 seems to be known , but we could not find a complete proof in the literature.
###### Lemma 5.
Let $`F:^m`$ be a real analytic function. Then the zero set of $`F`$ has Lebesgue measure zero, unless $`F`$ vanishes identically.
The proof of Lemma 5 is postponed to the end of this section. Since $`G_n`$ is not the zero sequence, the function $`H`$ does not vanish identically on $`[0,1]^m`$. By the Lebesgue dominated convergence theorem and Lemma 5 we thus find
$$\underset{\epsilon 0}{lim}𝝀(\stackrel{~}{S}_\epsilon )=0\text{and}\underset{\epsilon 0}{lim}𝝀(\stackrel{~}{L}_\epsilon )=𝝀(\stackrel{~}{L}_0).$$
This yields $`𝒅(\{n:G_n>s_n\})=𝝀(\stackrel{~}{L}_0)`$ by the following lemma, which completes the proof of Theorem 1.
###### Lemma 6.
Let $`G_n`$ and $`s_n`$ be real sequences with $`s_n=\mathrm{o}(1)`$ and let $`L_\epsilon `$, $`S_\epsilon `$ be as in (4). Suppose that $`𝐝(L_\epsilon )`$ and $`𝐝(S_\epsilon )`$ exist for all $`\epsilon 0`$, and that
$$\underset{\epsilon 0}{lim}𝒅(L_\epsilon )=𝒅(L_0)and\underset{\epsilon 0}{lim}𝒅(S_\epsilon )=0.$$
Then
$$𝒅(\{n:G_n>s_n\})=𝒅(L_0).$$
###### Proof.
For any set $`A`$ we write $`A(x):=\{nx:nA\}`$. Define
$$P:=\{n:G_n>s_n\}.$$
Let $`\epsilon >0`$ be arbitrary. Take $`n_0`$ such that $`|s_n|<\epsilon `$ for $`n>n_0`$. It follows
$`\mathrm{}P(x)`$ $`=\mathrm{}\{nn_0:G_n>s_n\}+\mathrm{}\{n_0<nx:G_n\epsilon \}`$
$`+\mathrm{}\{n_0<nx:s_n<G_n<\epsilon \},`$
hence
$$|\mathrm{}P(x)\mathrm{}L_\epsilon (x)|\mathrm{}S_\epsilon (x)+\mathrm{o}(x)$$
as $`x\mathrm{}`$. Thus we have
$`|x^1\mathrm{}P(x)𝒅(L_0)|`$ $`|x^1\mathrm{}P(x)x^1\mathrm{}L_\epsilon (x)|+|x^1\mathrm{}L_\epsilon (x)𝒅(L_0)|`$
$`x^1\mathrm{}S_\epsilon (x)+|x^1\mathrm{}L_\epsilon (x)𝒅(L_0)|+\mathrm{o}(1).`$
The right hand side goes to
$$𝒅(S_\epsilon )+|𝒅(L_\epsilon )𝒅(L_0)|$$
as $`x\mathrm{}`$. By assumption, this can be made arbitrarily small, which implies $`𝒅(P)=𝒅(L_0)`$. ∎
###### Proof of Lemma 5.
For $`m=1`$ this is clear, since then the zero set is countable. Now assume that we have established the result for $`1,\mathrm{},m1`$. Put
$$V:=\{(t_2,\mathrm{},t_m)^{m1}:F(,t_2,\mathrm{},t_m)\text{ vanishes identically }\}.$$
Take a real number $`s`$ such that $`F(s,,\mathrm{},)`$ is not identically zero. Clearly, $`F(s,t_2,\mathrm{},t_m)=0`$ for all $`(t_2,\mathrm{},t_m)V`$. By the induction hypothesis, this implies $`𝝀(V)=0`$. Note that $`V`$ is closed, hence measurable. Since $`F`$ is real analytic in the first argument, we have
$$_{}\chi _Z(t_1,\mathrm{},t_m)d𝝀(t_1)=0$$
for all $`(t_2,\mathrm{},t_m)V`$, where $`\chi _Z`$ is the characteristic function of the zero set
$$Z:=\{(t_1,\mathrm{},t_m)^m:F(t_1,\mathrm{},t_m)=0\}.$$
Since $`V`$ has measure zero, this implies
$$_{}\mathrm{}_{}\chi _Z(t_1,\mathrm{},t_m)d𝝀(t_1)\mathrm{}d𝝀(t_m)=0.$$
This argument works for any order of integration, hence we obtain $`_^m\chi _Z=0`$ by Tonelli’s theorem. ∎
## 3 Sequences with no Positive Dominating Root
In this section we prove Theorem 2. We begin by settling the special cases where the $`\theta _i`$ are all irrational or all rational, and then put them together.
###### Lemma 7.
Let $`\theta _1,\mathrm{},\theta _d`$ be irrational numbers, and let $`a_i,\beta _i`$ be real numbers such that the sequence
$$u_n=\underset{i=1}{\overset{d}{}}a_i\mathrm{cos}(2\pi \theta _in+\beta _i)$$
is not identically zero. Let further $`r_n`$ be a recurrence sequence with $`r_n=\mathrm{o}(1)`$. Then the set $`\{n:u_n>r_n\}`$ has positive density.
###### Proof.
Proceeding as in the proof of Theorem 1, we can write
$$G_n:=u_{gn+k}=𝐚^T\mathrm{cos}(2\pi n𝐁𝝉+𝐜),$$
where $`𝐁`$ is an integer matrix no row of which is zero, $`𝐜`$ is a real vector and $`1,\tau _1,\mathrm{},\tau _m`$ are linearly independent over $``$. If $`k`$ is such that $`G_n=u_{gn+k}0`$, then the density of $`\{n:G_n>s_n\}`$, where $`s_n=r_{qn+k}`$, exists by Theorem 1, but may be zero. Now choose a $`k_0`$ such that the corresponding sequence $`G_n=u_{gn+k_0}`$ is not the zero sequence. We have $`G_n=H(n𝝉)`$, where
$$H(𝐭):=𝐚^T\mathrm{cos}(2\pi \mathrm{𝐁𝐭}+𝐜).$$
Moreover, with the notation of the proof of Theorem 1, we have
$$𝒅(\{n:G_n>s_n\})=𝝀(\stackrel{~}{L}_0).$$
The function $`H`$ is not identically zero on $`[0,1]^m`$. But
$$_0^1\mathrm{}_0^1H(t_1,\mathrm{},t_m)dt_1\mathrm{}dt_m=0,$$
(6)
because no row of $`𝐁`$ is the zero vector. Hence $`H`$ has a positive value on $`[0,1]^m`$, and since it is continuous, we have $`𝝀(\stackrel{~}{L}_0)>0`$. ∎
Observe that the integral in (6) need not vanish if $`𝐁`$ has a zero row, which can only happen if the $`\theta _i`$ corresponding to this row is a rational number. This is the reason why we consider rational $`\theta _i`$’s separately.
###### Lemma 8.
Let $`\theta _1,\mathrm{},\theta _d`$ be rational numbers in $`]0,1[`$, and let $`a_i,\beta _i`$ be real numbers such that the purely periodic sequence
$$u_n=\underset{i=1}{\overset{d}{}}a_i\mathrm{cos}(2\pi \theta _in+\beta _i)$$
is not identically zero. Then $`u_n`$ has a positive and a negative value.
###### Proof.
By the identity
$$\underset{k=0}{\overset{q1}{}}\mathrm{cos}\frac{2\pi kp}{q}+\mathrm{i}\underset{k=0}{\overset{q1}{}}\mathrm{sin}\frac{2\pi kp}{q}=\underset{k=0}{\overset{q1}{}}\mathrm{e}^{2\pi \mathrm{i}kp/q}=0,$$
valid for integers $`0<p<q`$, and the addition formula of cos we obtain
$$u_0+\mathrm{}+u_{q1}=0,$$
where $`q`$ is a common denominator of $`\theta _1,\mathrm{},\theta _d`$. ∎
###### Proof of Theorem 2.
It suffices to consider the positivity set. We may write
$$f_n=u_n+v_nr_n,$$
where $`r_n=\mathrm{o}(1)`$ is a recurrence sequence,
$`u_n`$ $`={\displaystyle \underset{i=1}{\overset{d}{}}}a_i\mathrm{cos}(2\pi \theta _in+\beta _i),`$
$`v_n`$ $`={\displaystyle \underset{i=d+1}{\overset{e}{}}}a_i\mathrm{cos}(2\pi \theta _in+\beta _i),`$
$`\theta _1,\mathrm{},\theta _d`$ are irrational, $`\theta _{d+1},\mathrm{},\theta _e`$ are rational numbers in $`]0,1[`$ with common denominator $`q>0`$ and $`u_n+v_n0`$. If $`v_n0`$, then the result follows from Lemma 7. Now suppose $`v_n0`$. Then for each $`k`$ the density of the set $`\{n:f_{qn+k}>0\}`$ exists by Theorem 1. By Lemma 8 there is $`k_0`$ such that $`v_{qn+k_0}=v>0`$. It suffices to show that the set $`\{n:f_{qn+k_0}>0\}`$ has positive density. This is clear if $`u_{qn+k_0}0`$. Otherwise, notice that
$$\{n:f_{qn+k_0}>0\}\{n:u_{qn+k_0}>r_{qn+k_0}\},$$
and the latter set has positive density by Lemma 7. ∎
## 4 The Possible Values of the Density
In this section we investigate which values from $`[0,1]`$ occur as density of some recurrence sequence. In its basic form, the question is readily answered:
###### Example 9.
Let $`w`$ be a real number and define
$$f_n:=\mathrm{sin}(2\pi n\sqrt{2})w.$$
Then, by Theorem 4,
$`𝒅(\{n:f_n>0\})`$ $`=𝝀(\{t[0,1]:\mathrm{sin}(2\pi t)>w\})`$
$`=\{\begin{array}{cc}1\hfill & w1\hfill \\ \frac{1}{2}\frac{1}{\pi }\mathrm{arcsin}w\hfill & 1w1\hfill \\ 0\hfill & w1\hfill \end{array}.`$
Since the range of $`\mathrm{arcsin}`$ is $`[\frac{\pi }{2},\frac{\pi }{2}]`$, for every $`\kappa [0,1]`$ this yields a recurrence sequence $`f_n`$ such that
$$𝒅(\{n:f_n>0\})=\kappa .$$
The following proposition generalizes this example. Note that the density of the zero set of a recurrence sequence is always a rational number by the Skolem-Mahler-Lech theorem.
###### Proposition 10.
Let $`\kappa `$ be a real number and $`r`$ be a rational number with $`0\kappa ,r1`$ and $`\kappa +r1`$. Then there is a recurrence sequence $`(f_n)`$ such that
$$𝒅(\{n:f_n>0\})=\kappa \text{and}𝒅(\{n:f_n=0\})=r.$$
###### Proof.
Suppose that $`r=p/q`$ for positive integers $`p`$ and $`q`$. As seen in Example 9, there is a recurrence sequence $`(g_n)`$ such that the density of the zero set of $`(g_n)`$ is zero and the density of its positivity set is $`\kappa /(1r)`$ (The case $`r=1`$ is trivial). The interlacing sequence
$$f_{bn+k}:=\{\begin{array}{cc}0\hfill & 0k<p\hfill \\ g_n\hfill & pk<q\hfill \end{array}$$
is a recurrence sequence \[3, section 4.1\]. Clearly, the density of its zero set is $`r`$, and the density of its positivity set is
$$𝒅(\{n:f_n>0\})=\frac{qp}{q}\times \frac{\kappa }{1r}=\kappa ,$$
as required. ∎
If we restrict attention to sequences without dominating real positive roots, then Theorem 2 tells us that the density of the positivity set can be neither zero nor one. Still, all values in between occur.
###### Theorem 11.
Let $`\kappa ]0,1[`$. Then there is a recurrence sequence $`(f_n)`$ with no positive dominating characteristic root and $`𝐝(\{n:f_n>0\})=\kappa `$.
###### Proof.
Let $`\epsilon >0`$ be arbitrary. We define a function $`H`$ on $`[0,\frac{1}{2}]`$ by
$$H(t):=\{\begin{array}{cc}\frac{(\epsilon 1)^2}{\epsilon }\left(1\frac{2t}{\epsilon }\right)\hfill & 0t\frac{\epsilon }{2}\hfill \\ \epsilon 2t\hfill & \frac{\epsilon }{2}t\frac{1}{2}\hfill \end{array}$$
and extend it to an even, $`1`$-periodic function $`H`$ on $``$ (see Figure 1). It is continuous and satisfies
$$_0^1H(t)dt=0\text{and}𝝀(\{t[0,1]:H(t)>0\})=\epsilon .$$
Expanding $`H`$ into a Fourier series, we find that there are real $`a_j`$ such that $`H`$ is the pointwise limit of
$$H_m(t):=\underset{j=1}{\overset{m}{}}a_j\mathrm{cos}(2\pi jt)$$
as $`m\mathrm{}`$. Since the zero set of $`H`$ is a null set, the Lebesgue dominated convergence theorem yields
$$\underset{m\mathrm{}}{lim}𝝀(\{t[0,1]:H_m(t)>0\})=\epsilon .$$
We fix an $`m`$ such that
$$𝝀(\{t[0,1]:H_m(t)>0\})2\epsilon .$$
The function
$$\varphi (A_1,\mathrm{},A_m):=𝝀\left(\{t[0,1]:\underset{j=1}{\overset{m}{}}A_j\mathrm{cos}(2\pi jt)>0\}\right)$$
is continuous on $`^m\backslash \{\mathrm{𝟎}\}`$. To see this, observe that $`\varphi `$ is continuous at all points $`(A_1,\mathrm{},A_m)`$ for which $`_{j=1}^mA_j\mathrm{cos}(2\pi jt)`$ is not identically zero and appeal to the uniqueness of the Fourier expansion. Since $`\varphi (1,0,\mathrm{},0)=\frac{1}{2}`$ and $`\varphi (a_1,\mathrm{},a_m)2\epsilon `$, the function $`\varphi `$ assumes every value from $`[2\epsilon ,\frac{1}{2}]`$ by the intermediate value theorem.
Hence the positivity sets of the sequences
$$f_n:=\underset{j=1}{\overset{m}{}}A_j\mathrm{cos}(2\pi jn\sqrt{2})$$
assume all densities from $`[2\epsilon ,\frac{1}{2}]`$ for appropriate choices of $`(A_1,\mathrm{},A_m)`$ by Theorem 4. Repeating the whole argument with $`H`$ instead of $`H`$ yields the desired result for $`\kappa [\frac{1}{2},12\epsilon ]`$. Since $`\epsilon `$ was arbitrary, the theorem is proved. ∎
## 5 A Weak Version of Skolem-Mahler-Lech
Without using the Skolem-Mahler-Lech theorem, it follows from Theorem 1 that the density of the zero set of a recurrence sequence $`(f_n)`$ exists. We can show a bit more with our approach. Recall, however, that we only deal with real sequences, whereas the Skolem-Mahler-Lech theorem holds for any field of characteristic zero.
###### Proposition 12.
The density of the zero set of a (real) recurrence sequence $`(f_n)`$ is a rational number.
###### Proof.
Let $`k`$ be a natural number, and let $`g`$, $`G_n`$ and $`s_n`$ be as in the proof of Theorem 1. If $`k`$ is such that $`G_n0`$, then the density of the zero set of $`f_{gn+k}`$ is rational, since we may assume inductively that the density of $`\{n:s_n=0\}`$ is rational.
Now suppose $`G_n0`$. The zero set of $`f_{gn+k}`$ can be partitioned as
$$\{n:G_n=s_n\}=\{n:G_n=s_n,|G_n|<\epsilon \}\{n:G_n=s_n,|G_n|\epsilon \},$$
where $`\epsilon 0`$ is arbitrary. The latter set is finite, and the first one is contained in $`S_\epsilon `$, defined in (4). Hence
$$𝒅(\{n:G_n=s_n\})𝒅(S_\epsilon )$$
for all $`\epsilon 0`$. But we know that $`lim_{\epsilon 0}𝒅(S_\epsilon )=0`$ from the proof of Theorem 1, which yields
$$𝒅(\{n:G_n=s_n\})=0.$$
Thus, the zero sets of all subsequences $`(f_{gn+k})_{n0}`$, $`0k<g`$, have rational density, which proves the desired result. ∎
## 6 Conclusion
There is no algorithm known for deciding whether $`f_n>0`$ for all $`n`$, nor has the problem been shown to be undecidable. When we are talking about algorithmics, it is natural to assume that the recurrence coefficients and the initial values are rational numbers. In this case Gourdon and Salvy have proposed an efficient method for ordering the characteristic roots w.r.t. to their modulus. Thus, the dominating characteristic roots can be identified algorithmically. If none of them is real positive, then we know that the sequence oscillates by Theorem 2. On the other hand, sequences where a positive dominating root is accompanied by complex dominating roots seem to pose difficult Diophantine problems. For instance, we do not know if the sequence
$$f_n:=\mathrm{cos}(2\pi \theta n)+1+\left(\frac{1}{2}\right)^n$$
(7)
is positive for $`\theta =\sqrt{2}`$, say. It can be shown, however, that the set of $`\theta `$’s for which the corresponding sequence $`(f_n)`$ (defined by (7)) is positive has measure zero \[4, Theorem 7.2\].
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# Phonon anomalies and charge dynamics in Fe1-xCuxCr2S4 single crystals
## I Introduction
The discovery of colossal magnetoresistance (CMR) in perovskite-type manganites has attracted considerable attention.Kusters378 ; Jirak374 ; Von360 ; Chahara371 ; Jin373 Double-exchange (DE) mechanism, Zener365 ; Millis362 strong electron-phonon coupling, Millis362 phase separation scenarios Mayr387 or a Griffiths singularity Salamon were suggested to clarify the origin of the CMR effect, but a conclusive microscopic model has not yet been established. Ever since, the occurrence of CMR effects has been reported for various other classes of materials, such as pyrochlores,Shimakawa96 rare-earths based compounds like GdI<sub>2</sub>,Felser99 and ternary chalcogenide spinels $`A`$Cr<sub>2</sub>S<sub>4</sub>.Ramirez255 These CMR materials have been classified in terms of spin-disorder scattering and a universal dependence of the magnetoresistence vs. carrier density has been suggested on theoretical grounds.Majumdar98 ; Gomez-Santos04
Ramirez *et al.* drew attention to the spinel system Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> in 1997.Ramirez255 In polycrystalline FeCr<sub>2</sub>S<sub>4</sub> with $`T_C=170`$ K, the CMR effect reaches values comparable to those observed in perovskite oxides. The substitution of Fe by Cu increases $`T_C`$ to temperatures above room temperature, and the CMR effect remains relatively strong ($``$7 %).Ramirez255 In addition, solid solutions of the ferrimagnetic semiconductor FeCr<sub>2</sub>S<sub>4</sub> and the metallic ferromagnet CuCr<sub>2</sub>S<sub>4</sub> show a number of puzzling properties: From the very beginning, a controversial discussion has been arising whether the Cu ions are mono- or divalent for $`x0.5`$.Lotgering380 ; Goodenough110 ; Kurmaev379 For $`x<0.5`$ it was established that only monovalent and hence diamagnetic ($`d^{10}`$) Cu exists in the mixed crystals.Tsurkan32 Moreover, Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> shows two metal-to-insulator transitions as a function of $`x`$, as the room-temperature resistivity reveals two minima at $`x=0.2`$ and $`x=1`$ and concomitantly the Seebeck coefficient changes sign two times.Lotgering380 ; Haacke Additionally, band-structure calculations predicted that the Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> system should exhibit a half-metallic nature.Pickett01 ; Park213 ; Kurmaev379
Recent experimental investigations of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> single crystals indicated a strong dependence of their magnetic and magneto-transport properties on hydrostatic pressure suggesting a strong magneto-elastic coupling.Tsurkan347 ; Fritsch395 Measurements on the ac susceptibility in pure FeCr<sub>2</sub>S<sub>4</sub> exhibited a cusp in the low-field magnetization and the onset of magnetic irreversibilities at 60 K was explained by domain-reorientation processes.Tsurkan346 Later on, ultrasonic studies indicated an anomaly in the temperature dependence of the shear modulus close to 60 K, and it was suggested that the onset of orbital order induces a structural distortion at this temperature.Maurer396 This result, however, is hardly compatible with the observation that orbital order is established in polycrystals close to 10 K, while an orbital glass state is found in single crystals.Fichtl04
Optical spectroscopy simultaneously probes the lattice and electronic degrees of freedom and is, therefore, ideally suited to investigate structural phase transitions and to clarify the importance of electron-phonon coupling for the CMR effect.Hartinger04 Earlier infrared (IR) studies in polycrystalline FeCr<sub>2</sub>S<sub>4</sub> reported, in accordance with the crystal-lattice symmetry $`Fd3m`$, the existence of four IR-active phonons, which strongly depend on temperature near and below $`T_C`$.Lutz375 ; Wakamura351 We performed measurements of the optical properties of single crystals of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> ($`x=0,0.2,0.4`$ and 0.5) to shed light on the interplay of structural and electronic properties in these compounds. Since the optical properties of the samples with $`x=0.4`$ and $`x=0.5`$ were found to be very similar, we only show and discuss the corresponding data for $`x=0.5`$ in the following.
## II Experimental details
Single crystals of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> were grown using a chemical transport-reaction method with chlorine as transport agent and the ternary polycrystals as starting material. Details of the sample preparation are described elsewhere.Kurmaev379 No indication for the existence of secondary phases was found by x-ray diffraction analysis of powdered single crystals. X-ray single-crystal analysis confirmed the high structural homogeneity of the samples. The composition and homogeneity of the samples were examined by electron-probe microanalysis. The samples were optically polished platelets with dimensions of about $`3\times 5\times 1`$ mm<sup>3</sup>. Structural, magnetic and electrical transport data are given in Ref. Fritsch395, .
Two Fourier-transform-infrared spectrometers with a full bandwidth of 10 to 8000 cm<sup>-1</sup> (Bruker IFS 113v) and 500 to 42000 cm<sup>-1</sup> (Bruker IFS 66v/S) together with a <sup>4</sup>He cryostat (Oxford Optistat) were used for measurements of the optical reflectivity in the energy range from 70 to 30000 cm<sup>-1</sup> due to small sample dimensions and for temperatures of $`5\text{K}<T<300\text{K}`$. In order to investigate small fractions of the sample surface in the range of 0.1 mm<sup>2</sup> we utilized an IR microscope (Bruker IRscope II), which works in the far- (FIR) and mid-infrared (MIR) range.
## III Experimental results and discussion
### III.1 Phonon excitations
Figure 1 shows the temperature dependence of the FIR reflectivity $`R`$ vs. wave number of pure FeCr<sub>2</sub>S<sub>4</sub>. In the upper panel $`R`$ is plotted for 5 and 300 K. The four visible phonon peaks are attributed to the four IR-active $`F_{1u}`$ modes (symmetry group $`Fd3m`$, #227).Lutz375 To analyze the spectra, we used a 4-parameter fit assuming frequency-dependent damping constants to account for the asymmetry of the phonon peaks. This fitting procedure infers a splitting of the longitudinal and transverse eigenfrequencies, $`\omega _L`$ and $`\omega _T`$, and the corresponding damping constants, $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_T`$. Berreman392 The resulting curves describe the measured reflectivity down to 100 cm<sup>-1</sup> very well, without assuming an additional contribution of free charge carriers. A representative result of these fits is shown by the solid line in the upper panel of Fig. 1 for $`T=5`$ K. The detailed temperature dependence of the reflectivity is visualized in the two-dimensional (2D) contour plot in the lower panel of Fig. 1. To enable a comparison of the phonon shift, the peak positions (maxima in $`R`$) for $`T=5`$ K are indicated as vertical lines. Around $`T_C=167`$ K a shift of the phonon frequencies can be observed, especially for the mode $`d`$ close to 100 cm<sup>-1</sup>. The intensity of this mode strongly depends on temperature, too (see upper frame of Fig. 1).
The resonance frequencies $`\omega _L`$ and $`\omega _T`$ (left frames) and the corresponding damping rates $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_T`$ (right frames) are shown in Fig. 2 as a function of temperature. Above the Curie temperature $`T_C=167`$ K, the resonance frequencies $`\omega _L`$ and $`\omega _T`$ of all modes reveal a similar quasi-linear increase with decreasing temperature, which can be fully ascribed to anharmonic contributions to the lattice potential.Wakamura383 In contrast to the rather usual behavior in the paramagnetic regime, modes $`a`$ and $`b`$ soften for temperatures below $`T_C`$, while $`\omega _L`$ and $`\omega _T`$ increase towards lower temperatures in the case of modes $`c`$ and $`d`$. These anomalous changes of the eigenfrequencies in the vicinity of $`T_C`$ suggest a correlation with the onset of magnetic order. However, it has to be stated that the size of the effect is different for the observed modes: $`\mathrm{\Delta }\omega =[\omega (T=T_C)\omega (T=0`$K$`)]/\omega (T=0`$K$`)`$ is of the order of $`+3`$ % for the internal mode $`d`$, $`+1`$ % for the bending mode $`b`$, approximately $`1.5`$ % for the bending mode $`c`$, and $`1`$ % for the stretching mode $`a`$. Longitudinal and transverse eigenfrequencies behave rather similar.
The influence of magnetic order on phonons in magnetic semiconductors has been proposed by Baltensperger and HelmanBaltensperger26 and BaltenspergerBaltensperger27 more than 30 years ago, and has recently been used by Sushkov et al. to describe the phonon spectra in ZnCr<sub>2</sub>O<sub>4</sub>.Sushkov04 Based on a model calculation, where superexchange interaction between the magnetic ions infers a spin-phonon coupling, relative frequency shifts up to 10<sup>-2</sup> have been predicted. The order of magnitude of this effect corresponds nicely to the experimentally observed values in FeCr<sub>2</sub>S<sub>4</sub> and, therefore, WakamuraWakamura351 considered this mechanism to dominate the phonons’ behavior for $`TT_C`$. Subsequently, Wakamura and coworkersWakamura383 ; Wakamura discussed the sign of the relative frequency shift in terms of nearest-neighbor FM exchange and next-nearest-neighbor AFM exchange for CdCr<sub>2</sub>S<sub>4</sub>, which exhibits phonon modes with a similar temperature dependence as FeCr<sub>2</sub>S<sub>4</sub>. Moreover, they could show that these anomalous changes in the phonon frequencies are absent in non-magnetic CdIn<sub>2</sub>S<sub>4</sub>, further corroborating their approach.Wakamura Thus, the positive shift of modes $`a`$ and $`b`$ would indicate that FM exchange (Cr-S-Cr) dominates in accordance with a strong influence of the (Cr-S) force constants on these modes, and, correspondingly, the negative shift of modes $`c`$ and $`d`$ favors AFM exchange (Cr-S-Cd-S-Cr) with a strong influence of the (Cd-S) force constants. Note that a more rigorous theoretical treatment of anharmonic spin-phonon and phonon-phonon interactions in cubic spinels by Wesselinova and ApostolovWesselinova96 confirms the above interpretation. In FeCr<sub>2</sub>S<sub>4</sub> the interpretation of the effect of magnetic ordering on the IR active phonon modes becomes even more complicated, because there exist, besides FM nearest-neighbor Cr-S-Cr bonds, additional exchange paths via AFM Fe-S-Fe and Fe-S-Cr-S-Fe bonds. Nevertheless, the overall temperature behavior of the phonon frequencies in FeCr<sub>2</sub>S<sub>4</sub> is similar to CdCr<sub>2</sub>S<sub>4</sub> and may be well interpreted, accordingly. Note, however, that a critical discussion of the above approach is given by Bruesch and d’Ambrogio.Bruesch
A straightforward interpretation of the temperature dependence of the damping constants (right panel of Fig. 2) is not obvious at all. Again, considering only the anharmonicity of ionic non-magnetic crystals, the damping is expected to show some residual low-temperature value and a quasi-linear increase in the high-temperature limit, just as observed for the longitudinal damping constants of modes $`a`$ and $`b`$ for $`T>T_C`$.Wakamura351 However, the temperature dependence of $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_T`$ in general deviates from such a behavior: In the case of mode $`d`$ both damping constants show a broad maximum just above $`T_C`$ and a steep decrease towards lower temperature for $`T<T_C`$. Mode $`c`$ follows a similar temperature dependence for $`T<T_C`$, but the reduction of the damping constants is slightly smaller, and in the paramagnetic regime $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_T`$ remain almost constant in contrast to the results of Wakamura.Wakamura351 The behavior of modes $`a`$ and $`b`$ for $`TT_C`$ appears even more complex, but one can identify the onset of enhancement damping close to $`T_C=170`$ K followed by broad cusp-like maxima close to 100 K, except for $`\omega _T`$ of mode $`a`$ that increases linearly with decreasing temperatures.
WakamuraWakamura351 argues that the maxima of mode $`d`$ (and $`c`$) are due to spin fluctuations of the Fe spins, in agreement with the strong influence of the corresponding force constant on this mode according to Bruesch and d’Ambrogio. Bruesch Furthermore, long range spin order assumingly leads to the anomalous changes of the damping constants for all modes below $`T_C`$. In comparison to the temperature dependences of the damping constants in CdCr<sub>2</sub>S<sub>4</sub>, one finds that modes $`c`$ and $`d`$ behave similar to the case of FeCr<sub>2</sub>S<sub>4</sub>.Wakamura383 On the other hand, modes $`a`$ and $`b`$ in FeCr<sub>2</sub>S<sub>4</sub> clearly reveal a more complex behavior than in CdCr<sub>2</sub>S<sub>4</sub>, indicating a significant influence of the iron sublattice and the additional effective exchange coupling between Fe-Fe and Fe-Cr ions on these modes.
Additionally, we want to mention the large increase in intensity (about 20 %) for mode $`d`$ (close to 120 cm<sup>-1</sup>) when cooling from room temperature to 5 K (see Fig. 1). The intensity remains almost constant above 200 K, while a linear increase with decreasing temperature is observed below 200 K. At this temperature, maxima appear in the temperature dependence of the damping constants, suggesting a correlation of the two phenomena with regard to the spin-fluctuation scenario discussed above.
When adopting the overall interpretation of the data in terms of spin-phonon coupling, one has to consider, however, that e.g. the appearance of the cusps in the damping constants may be connected to domain reorientation processes visible in the ac susceptibilityTsurkan346 and anomalies detected by ultrasonic investigations.Maurer396 Although the absence of significant changes of the phonon frequencies contradicts the scenario of a structural phase transition at 60 K driven by orbital ordering as suggested in Ref. Maurer396, , it becomes clear that the complex mechanisms dominating the damping effects demand further theoretical studies to single out the important contributions in detail.
Having discussed the phonon properties of pure FeCr<sub>2</sub>S<sub>4</sub> we now turn to the temperature dependence of the phonon modes for Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub>. Figure 3 shows the FIR reflectivity for $`x=0.2`$ (upper panel) and $`x=0.5`$ (lower panel) for temperatures 5 K and 300 K each. The results for $`x=0.4`$ are very similar to those obtained for $`x=0.5`$ and, hence, only the data for $`x=0.5`$ is shown and discussed. The reflectivity of both samples, $`x=0.2`$ and 0.5, shows a Drude-like contribution due to the presence of free charge carriers, while FeCr<sub>2</sub>S<sub>4</sub> can be described as an insulator. The highest Drude-like conductivity is found for $`x=0.2`$ and the phonon modes are on the verge of being fully screened. For both compounds the internal mode $`d`$ at $`120`$cm<sup>-1</sup> (see Fig. 1 for the pure compound) can hardly be detected. Focusing on the group of external modes, a new mode $`e`$ appears close to $`350`$cm<sup>-1</sup>, while on increasing Cu concentration $`x`$ mode $`a`$ at $`380`$cm<sup>-1</sup> becomes considerably reduced in intensity. Without an accompanying lattice dynamical calculation one cannot decide, if this new mode represents an impurity mode due to the doping with Cu or a symmetry change. There are reports in literatureLotgering380 ; Palmer384 claiming the reduction of symmetry to $`F\overline{4}3m`$ because of the ordering of Fe and Cu ions on the A sublattice.
In this case seven IR-active phonon modes are predicted. Assuming that the peak close to 275 cm<sup>-1</sup> (mode $`c`$) is generated by two single phonon modes (see inset in Fig. 3), five modes are visible, only, while the internal mode $`d`$ close to $`120`$cm<sup>-1</sup> remains screened. However, our results rather point toward a continuous evolution of the phonon modes on increasing $`x`$ and favor a statistical A-site distribution of Cu and Fe ions throughout the lattice instead of a symmetry reduction resulting from a superstructure due to an ordered A sublattice.
We tried to fit the complete spectra taking into account the reflectivity up to 10000 cm<sup>-1</sup>, using a 4-parameter fit for the phonon modes and a Lorentz oscillator for the mid-infrared excitation at about $`2500`$cm<sup>-1</sup> (see next section). Representative results of these fits at low wave numbers are shown in Fig. 3 as dashed lines. The temperature dependences of the longitudinal modes $`\omega _L`$ as derived from these fits are shown in Fig. 4 together with corresponding data for $`x=0`$. The transverse eigenfrequencies behave rather similar (not shown). Compared to the sample with $`x=0`$, the temperature dependence of the damping constants in the doped compounds is weak and therefore not shown here either. Regarding the sample with $`x=0.2`$ similar anomalies as in pure FeCr<sub>2</sub>S<sub>4</sub> can be seen in the vicinity of $`T_C`$ for the observable modes $`a,b`$, and $`e`$. Obviously, the temperature dependence of all phonon frequencies for $`x=0.5`$ is very weak and no clear anomalies around $`T_C`$ are visible. Within the experimental uncertainties one can detect a slight decrease of $`\omega _L`$ towards lower temperatures except for mode $`c`$, which behaves similarly to the case of FeCr<sub>2</sub>S<sub>4</sub> (compare Fig. 2).
Keeping in mind the influence of spin fluctuations and spin-phonon coupling on the phonon properties in FeCr<sub>2</sub>S<sub>4</sub>, Cu-doping seems to reduce these features significantly. This observation is in agreement with reduced spin-orbit coupling due to the substitution of Jahn-Teller active Fe<sup>2+</sup> by non Jahn-Teller active Fe<sup>3+</sup>. Therefore, for $`x=0.5`$ only Fe<sup>3+</sup> with a half-filled $`d`$-shell is present in the systemLang00 ; Kurmaev379 and the system becomes almost magnetically isotropic as it was confirmed by ferromagnetic resonance experiments.Fritsch395
### III.2 Dynamic conductivity and electronic excitations
When the reflectivities of the doped compounds with Cu concentrations $`x=0.2`$ and 0.5 (Fig. 3) are compared with that of pure FeCr<sub>2</sub>S<sub>4</sub> it becomes clear that contributions from free charge carriers have to be taken into consideration. The metallic-like behavior is most significant for $`x=0.2`$, but it becomes reduced again on further doping. For a consistent description of the Drude-type behavior of the doped compounds, it is important to measure the reflectivity spectra to higher energies. The room-temperature reflectivities of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> for $`x=0.2`$ and 0.5 are plotted in the upper panel of Fig. 5 up to 3$`\times `$10<sup>4</sup> cm<sup>-1</sup>, corresponding to almost 4 eV, and are compared to the reflectivity of insulating FeCr<sub>2</sub>S<sub>4</sub>. For the Kramers-Kronig analysis of the smoothed reflectivity data we used a low-frequency Hagen-Rubens extrapolation and a high-frequency extrapolation with a $`\nu ^{0.5}`$ power law up to 10<sup>6</sup> cm<sup>-1</sup> and a subsequent $`\nu ^4`$ high-frequency tail. The resulting dynamic conductivities $`\sigma (\nu )`$ are shown in the lower panel of Fig. 5. We carefully checked the high-frequency extrapolation, also trying smoother extrapolations, but found that the results are not influenced in the relevant energy range below 20000 cm<sup>-1</sup>. The use of a Hagen-Rubens extrapolation is justified by the fact that we have the complete information on the absolute values of the dc conductivities and the corresponding temperature dependences for all compounds, although we are aware of the additional uncertainties originating from the Hagens-Rubens extrapolation, specifically for the sample with $`x=0.5`$. However, the best fits of the reflectivity at room temperature, even in the limited spectral range, yielded dc conductivities of 150 ($`\mathrm{\Omega }`$cm)<sup>-1</sup> for $`x=0.2`$ and 35 ($`\mathrm{\Omega }`$cm)<sup>-1</sup> for $`x=0.5`$, close to the dc values derived from the 4-probe measurements on single crystals by Fritsch et al. Fritsch395
For $`x=0`$ a weak but well defined electronic transition is observed close to $`2000`$cm<sup>-1</sup> and a further transition appears close to 20000 cm$`{}_{}{}^{1}(2.5`$ eV). On substituting iron by copper, metallic behavior shows up and for Fe<sub>0.8</sub>Cu<sub>0.2</sub>Cr<sub>2</sub>S<sub>4</sub> the dc conductivity is of the order 150 ($`\mathrm{\Omega }`$cm)<sup>-1</sup>. The transition at $`2000`$cm<sup>-1</sup>, becomes almost fully suppressed for $`x=0.2`$. Obviously, the $`d`$-electrons become strongly delocalized. It is generally accepted that in an ionic picture monovalent Cu is substituted inducing trivalent Fe. Our results suggest that the system behaves as if holes are doped into an insulator driving the compound into a metallic regime. Unexpectedly, a broad peak appears again close to 2500 cm<sup>-1</sup> for $`x=0.5`$.
The observed doping dependence of the conductivity spectra as documented in Fig. 5 can be compared with band-structure calculations of these compounds. Lang00 ; Park213 ; Kurmaev379 Local spin-density approximation (LSDA) band-structure calculations predict a half-metallic ground state of FeCr<sub>2</sub>S<sub>4</sub>, with a partly filled $`e`$ band at the Fermi level. Correlation effects via LSDA+U yield a splitting of the Fe $`e`$ band into a lower and upper Hubbard band characterizing FeCr<sub>2</sub>S<sub>4</sub> as a Mott-Hubbard insulator.Park213 The splitting of the $`e`$ band is of the order of about 0.5 eV, and, hence, the peak close to 2000 cm<sup>-1</sup> may be interpreted as a transition between the lower and upper Hubbard band. Accordingly, the high-energy excitation can be attributed to a Cr(3d) to Fe(3d) transition.
Using an ionic picture with localized Fe $`d`$ states, alternatively, the transition at 2000 cm<sup>-1</sup> may correspond to a transition between the lower $`e`$ doublet and the $`t_2`$ triplet of the Fe $`d`$-states split in a tetrahedral crystal-field. The expected crystal-field splitting for Fe<sup>2+</sup> located in the tetrahedral site of the spinel structure is rather weakAbragam70 and a splitting of the order $`20003000`$ cm<sup>-1</sup> seems reasonable. Further support for this interpretation comes from the observation of crystal-field transitions as measured for diluted Fe<sup>2+</sup> in CdIn<sub>2</sub>S<sub>4</sub>. Here a crystal-field splitting of approximately 2500 cm<sup>-1</sup> has been reported by Wittekoek *et al.*Wittekoek397
The appearance of the broad excitation for $`x=0.5`$ in the mid-infrared region at about 2500 cm<sup>-1</sup>, however, cannot be explained easily. In an ionic picture only trivalent iron and monovalent copper are expected for $`x=0.5`$,Lotgering380 ; Goodenough110 and recent x-ray photoelectron spectroscopy Tsurkan32 strongly favors the existence of only monovalent Cu for $`x=0.5`$. Therefore, one can exclude the possibility that the broad excitation may be attributed to Fe<sup>2+</sup> similarly to the well-defined electronic excitation for $`x=0`$. Nevertheless, it has been concluded from Mössbauer experiments in Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub>, that an ionic picture is not applicable at all. Haacke2 For $`x=0.3`$ and $`T<T_C`$ the complicated Mössbauer spectra indicate two different Fe sites corresponding to Fe<sup>2+</sup> and Fe<sup>3+</sup>, while for $`T>T_C`$ a single line pointed towards a fast electron exchange between these two sites. For $`x=0.5`$ the line pattern for $`T>T_C`$ evidenced the existence of Fe<sup>3+</sup> and a strong delocalization of the Cu $`d`$-derived electrons. Hence, further studies beyond the scope of this paper are needed to clarify the nature of this mid-infrared excitation.
In the following we will discuss the optical conductivity results in the low frequency range in comparison with the dc conductivity data reported in Ref. Fritsch395, . The room-temperature spectra for the concentrations $`x=0.2`$ and 0.5, shown in Fig. 5 have been used to estimate the Drude-like conductivity. For all temperatures, the spectra could satisfactorily be described using a plasma frequency $`\omega _p=12000`$ cm<sup>-1</sup> and a dielectric constant $`ϵ_{\mathrm{}}=10.6`$ for $`x=0.2`$, which is close to the value $`ϵ_{\mathrm{}}=11.5`$ for $`x=0`$. For $`x=0.5`$ we used $`\omega _p=5000`$cm<sup>-1</sup> and an enhanced dielectric constant $`ϵ_{\mathrm{}}=15.5`$. The enhanced $`ϵ_{\mathrm{}}`$ indicates strong changes in the electronic excitation spectrum at higher frequencies, but due to the complexity of the spectrum in this energy region there is also a larger uncertainty in $`ϵ_{\mathrm{}}`$ for $`x=0.5`$. The decrease of the plasma frequency by a factor of 2.4 can be explained by a decrease of the charge carrier density, as Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> approaches a metal-to-insulator transition close to $`x=0.5`$. With these values, the conductivity below 500 cm<sup>-1</sup> could reasonably be fitted for all temperatures as indicated by the dashed lines in Fig. 3 for the spectra at 5 K and 300 K.
The resulting temperature dependences of the dc conductivity (upper panel) and relaxation rates $`\gamma \tau ^1`$ (lower panel) are shown in Fig. 6. The dc conductivities as derived from 4-probe measurementsFritsch395 are indicated by solid lines. The dc conductivities were scaled at room temperature, utilizing a factor of 1.6 for $`x=0.2`$ and a factor of 1.05 for $`x=0.5`$. Above 100 K the 4-probe dc results and the dc values as derived from the optical measurements follow a similar temperature dependence. However, at low temperatures the dc measurements are dominated by localization effects, which appear much weaker in the high-frequency ($`>100`$cm<sup>-1</sup>) derived optical data. That localization effects are most significant in the low-frequency (”dc”) transport measurements becomes clear from the fact that in doped semiconductors the conductivity below the FIR regime increases almost linearly with frequency. Lunkenheimer In the sample with $`x=0.5`$, which exhibits the lower conductivity, localization effects dominate already at higher temperatures. This may be attributed to a significant decrease of the charge-carrier density and concomitant increase of disorder due to the statistical distribution of the Cu ions in the lattice,Lang00 further discarding the possibility of $`A`$-site order of Fe and Cu for Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub>.
Finally, we want to draw attention to the temperature dependence of the relaxation rate $`\gamma `$ (lower panel of Fig. 6). In the magnetically ordered state below $`T_C`$, the relaxation rates become significantly reduced, e.g. the reduction amounts to almost 50% for $`x=0.2`$. We recall that the plasma frequency has been kept constant for each compound as a function of temperature. This indicates that the increase of the conductivity just below the magnetic ordering temperature results from the freezing-out of disorder scattering and not from a change of the carrier density via band-structure changes at the onset of ferrimagnetic order. Taking into account the classification of chalcogenide spinels $`A`$Cr<sub>2</sub>S<sub>4</sub> as systems where CMR originates from spin-disorder scattering,Majumdar98 the observed reduction of the relaxation rate below $`T_C`$ has to be regarded as direct evidence of such a scenario: In external fields the onset of ferrimagnetic order shifts to higher temperatures. Concomitantly, a reduction of the scattering rate and the anomalous increase of the conductivity arise. As a consequence, maximal CMR effects will show up just below $`T_C`$ as a function of an external magnetic field. A similar scenario has been reported for GdI<sub>2</sub>, where the magnetic and magneto-transport properties have been described successfully in terms of spin-fluctuations and their suppression by external magnetic fields in the vicinity of $`T_C`$.Eremin01 ; Deisenhofer04 We would like to point out, that at low temperatures the relaxation rates for $`x=0.2`$ and $`x=0.5`$ are of the same order of magnitude $`10^4`$cm<sup>-1</sup>, indicating a similar level of disorder for the Cu doped compounds.
## IV Summary
In summary, we investigated the optical properties of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> single crystals for Cu concentrations $`x=0`$, 0.2, 0.4 and 0.5. Phonon excitations and dynamic conductivity for $`x=0.4`$ are very similar to the results for Fe<sub>0.5</sub>Cu<sub>0.5</sub>Cr<sub>2</sub>S<sub>4</sub> and were not discussed separately. The phonon excitations were measured as a function of temperature between 5 K and room temperature. Pure FeCr<sub>2</sub>S<sub>4</sub> shows clear anomalies in the eigenfrequencies at the transition from the paramagnetic to the ferromagnetic state, which can be explained by spin-phonon coupling. Concerning the complex behavior of the damping constants, spin fluctuations in the vicinity of $`T_C`$ may describe many of the anomalous changes, but further theoretical studies are necessary to corroborate this interpretation. The influence of magnetic order on the eigenmodes is reduced with increasing $`x`$, and the appearance of a new phonon mode close to 350 cm<sup>-1</sup> is attributed to an impurity mode rather than to a symmetry reduction due to A-site order.
Morover, the charge dynamics of Fe<sub>1-x</sub>Cu<sub>x</sub>Cr<sub>2</sub>S<sub>4</sub> were investigated. FeCr<sub>2</sub>S<sub>4</sub> is an insulator, but becomes metallic when slightly doped with Cu. The conductivity of the free charge carriers can be described by a normal Drude-type behavior. The dc conductivity for $`x=0.2`$ is enhanced by a factor of four in comparison to $`x=0.5`$. The temperature dependence of the optically derived dc conductivity for both doped compounds is is in good agreement with resistivity measurements, but localization effects at lowest temperatures appear weaker in the optical measurements. The corresponding behavior of the scattering rate, which shows a strong decrease below the ferrimagnetic phase transition, evidences the freezing-out of disorder scattering below $`T_C`$. In accordance with the proposed classification of the ternary chalcogenide spinels as spin-disorder magnetoresistive materials, the reduction of the relaxation rate corroborates such a scenario and makes clear that spin-disorder has to be considered a necessary ingredient towards a theoretical description of this fascinating class of materials.
###### Acknowledgements.
It is a pleasure to thank H.-A. Krug von Nidda, J. Hemberger, and Ch. Hartinger for fruitful discussions. This work was partly supported by the DFG via the Sonderforschungsbereich 484 (Augsburg), by the BMBF/VDI via the Contract No. EKM/13N6917/0, by the U.S. Civilian Research & Development Foundation (CRDF) and by the Moldavian Research & Development Association (MRDA) via Grant No. MP2-3047.
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# Casimir force between planes as a boundary finite size effect
## I Introduction
The Casimir force between two neutral macroscopic bodies in vacuum is considered as a manifestation of zero point fluctuations intrinsic to any quantum field theory (QFT). It has received a lot of interest recently both from the theoretical and from the experimental side (cf. the recent reviews reviews and the book by Milton book ).
This force can be determined from the volume dependence of the ground state energy (Casimir energy). It is usually derived either by a summation of zero point fluctuations of the quantum field restricted to the domain in question, or by analyzing the statistical fluctuations of the interactions between the boundaries. Our new derivation is similar to the first approach in the sense that it uses the quantum fields living between the boundaries. However, we describe the quantum field theory in question in terms of its particle excitations rather than the field fluctuations. As a consequence we describe the Casimir energy as ’virtual particle’ exchange between the two boundaries and our formula is expressed in terms of the analytically continued reflection amplitudes of these particles.
The interaction between the quantum field and the macroscopic bodies is most often described by imposing appropriate boundary conditions on the field. The aim of this paper is to explore and make explicit the dependence of the Casimir force on the boundary conditions obeyed by the field(s) in a geometry when the boundary conditions are given on two parallel planes. We achieve this by regarding the force as a physical manifestation of finite size effects in a boundary QFT (BQFT).
In QFT, finite size effects in a periodic box were considered by Lüscher luscher , who showed how the correction to particle masses (infinite volume energies) can be expressed in terms of infinite volume scattering data ($`S`$ matrix elements). Following this general idea, we have recently determined the finite size correction to the ground state energy in $`2`$ dimensional BQFT in terms of scattering data on the half line (reflection amplitudes). This explicitly describes how the ground state energy depends on the volume and it is obvious from the results in bluscher that a generalization to BQFT in any spacetime dimension is straightforward and leads to a description of the Casimir effect.
Boundary conditions in quantum field theory can be equally well described by a boundary state in a crossed channel. This formalism was worked out in detail for integrable 1+1 dimensional BQFTs GZ , but it can be extended to nonintegrable field theories in any number of spacetime dimensions. We show that this formalism allows an alternative method to derive the Casimir energy, and it has the advantage of using only renormalized and phenomenologically meaningful field theoretic quantities. A further advantage of using the boundary state formalism is that it provides a systematic (large volume) expansion for the Casimir energy in the interacting case.
Our main result is a large volume expression for the Casimir energy/force, valid in any BQFT, which depends only on the reflection amplitudes on the boundaries (although the boundary conditions considered allow transmission as well). These reflection amplitudes can be determined in a straightforward way in the half-infinite (one boundary) geometry, and using them in the expression for the Casimir energy makes it unnecessary to carry out the finite volume quantization separately in the various applications.
The paper is organized as follows: in Section II we introduce BQFTs and derive our main results for the volume dependence of their ground state energy using the boundary state formulation. Section III contains a description of several applications of the result in various BQFTs: First we derive a formula for the Casimir energy of a free massive boson subject to Robin type boundary condition. This expression is tested against the well-known Dirichlet, Neumann and massless limits. We further check our formula by recomputing the Casimir effect for parallel dielectric slabs separated by a vacuum slot and for a massless fermion subject to “bag boundary condition”. We make our concluding remarks in Section IV. In the Appendix we confirm our result in a simple case by presenting the naive derivation based on mode summation tailor made for the framework of boundary QFTs.
## II Derivation of the main formula
Here we summarize the main aspects of BQFTs and present the novel derivation of the Casimir energy. We follow the description of BQFTs along the line of bLSZ extending the result of GZ for higher dimensions. In the process we show that this formalism makes it possible to establish a systematic (large volume) expansion of the Casimir energy even for interacting fields.
### II.1 Boundary state formalism
To simplify the presentation consider a quantum field theory of a scalar field, $`\mathrm{\Phi }(t,x,\stackrel{}{y})`$, in $`D+1`$ dimensions with Lagrangian
$$=\frac{1}{2}\left(_t\mathrm{\Phi }\right)^2\frac{1}{2}\left(_x\mathrm{\Phi }\right)^2\frac{1}{2}\left(\stackrel{}{}\mathrm{\Phi }\right)^2V(\mathrm{\Phi })$$
where $`\stackrel{}{y}`$ is a $`D1`$ dimensional position vector, $`\stackrel{}{}_i=\frac{}{y_i}`$. We also suppose that the spectrum consists of one particle type with mass $`m`$. The Hilbert space of the model can be spanned by asymptotic multi-particle states with momentum parameterized by $`(k_i,\stackrel{}{k}_i)`$ (where again $`\stackrel{}{k}_i`$ denotes a $`D1`$ dimensional momentum vector):
$$=\left\{|k_1,\stackrel{}{k}_1;k_2,\stackrel{}{k}_2;\mathrm{};k_n,\stackrel{}{k}_n\right\}$$
The energy of a one-particle state is
$$\omega (k,\stackrel{}{k})=\sqrt{m^2+k^2+\stackrel{}{k}^2}$$
In the asymptotic scattering configurations (large negative or positive time) the particles are distant from each other, and the spectrum can be described by free *in* or *out* particles so the energy of the multi-particle state is the sum of the individual energies. These two Hilbert spaces are connected by the scattering matrix
$`S(\{k_i^{^{}},\stackrel{}{k}_i^{^{}}\},\{k_i,\stackrel{}{k}_i\})`$ $`=`$
$`{}_{}{}^{out}k_1^{^{}},\stackrel{}{k}_1^{^{}};k_2^{^{}},\stackrel{}{k}_2^{^{}};\mathrm{};k_n^{^{}},\stackrel{}{k}_n^{^{}}|k_1,\stackrel{}{k}_1;k_2,\stackrel{}{k}_2;\mathrm{};k_m,\stackrel{}{k}_m_{}^{in}`$
which can be expressed in terms of the correlators
$`G(t_1,x_1,\stackrel{}{y}_1;\mathrm{};t_N,x_N,\stackrel{}{y}_N)`$ $`=`$
$`0\left|\mathrm{\Phi }(t_1,x_1,\stackrel{}{y}_1)\mathrm{}\mathrm{\Phi }(t_N,x_N,\stackrel{}{y}_N)\right|0`$
via the LSZ reduction formula. The time evolution of the state is generated by the Hamiltonian
$$H=_{\mathrm{}}^{\mathrm{}}𝑑x𝑑\stackrel{}{y}\left(\frac{1}{2}\mathrm{\Pi }^2+\frac{1}{2}\left(_x\mathrm{\Phi }\right)^2+\frac{1}{2}\left(\stackrel{}{}\mathrm{\Phi }\right)^2+V(\mathrm{\Phi })\right)$$
where $`\mathrm{\Pi }=_t\mathrm{\Phi }`$ is the conjugate momentum.
Let us suppose now that the theory is restricted to the half space $`x<0`$ and the boundary condition is given by specifying some boundary potential at $`x=0`$
$$V_B(\mathrm{\Phi }(0,\stackrel{}{y},t))$$
The essential observation is that we can have two different Hamiltonian descriptions of this system. We can take $`t`$ as the time variable and then the boundary is situated in space (i.e. it has a spacelike normal vector). In this case the Hilbert space consist of multi-particle states
$$_B=\left\{|k_{},\stackrel{}{k}_{};k_{}^{},\stackrel{}{k^{}}_{};\mathrm{}_B\right\}$$
where we indicate explicitly that the full momentum vector is composed of two components, one perpendicular and the other parallel to the boundary. The states are normalized in the following way:
$`{}_{B}{}^{in}k_{}^{},\stackrel{}{k}_{}^{}|k_{},\stackrel{}{k}_{}_{B}^{in}=`$
$`\left(2\pi \right)^D\omega (k_{},\stackrel{}{k}_{})\delta \left(k_{}k_{}^{}\right)\delta ^{(D1)}\left(\stackrel{}{k}_{}\stackrel{}{k}_{}^{}\right)`$
The subscript $`B`$ indicates these states satisfy the boundary condition. In the asymptotic past all the particles move towards the boundary: $`k_{}>0`$. For $`t\mathrm{}`$ they are separated far from each other and from the boundary forming the *in* state, which is a free state. For $`t\mathrm{}`$ all the scatterings and reflections are terminated they move backwards from the boundary are far from each other and from the boundary forming the free *out* state. These two Hilbert spaces are connected by the reflection matrix
$$R_{\alpha \beta }={}_{B}{}^{out}\alpha |\beta _{B}^{in}$$
The simplest case describes the one particle elastic reflection on the boundary and can be written as
$`R(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))\left(2\pi \right)^D\delta \left(\theta \theta ^{}\right)\delta ^{(D1)}\left(\stackrel{}{k}_{}\stackrel{}{k}_{}^{}\right)=`$
$`{}_{B}{}^{out}k_{}^{^{}},\stackrel{}{k}_{||}^{^{}}|k_{},\stackrel{}{k}_{||}_{B}^{in}`$
where we exploited the unbroken Poincaré symmetry (in the coordinates $`t,\stackrel{}{y}`$) and parameterized the perpendicular momentum as
$$k_{}=m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{sinh}\theta ;m_{\mathrm{eff}}(\stackrel{}{k}_{||})=\sqrt{m^2+\stackrel{}{k}_{||}^2}$$
and then
$$\delta \left(\theta \theta ^{}\right)=m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta \delta \left(k_{}k_{}^{}\right)$$
The reflection factors can be expressed in terms of the correlators
$`G(t_1,x_1,\stackrel{}{y}_1;\mathrm{};t_N,x_N,\stackrel{}{y}_N)`$ $`=`$
$`{}_{B}{}^{}0\left|\mathrm{\Phi }(t_1,x_1,\stackrel{}{y}_1)\mathrm{}\mathrm{\Phi }(t_N,x_N,\stackrel{}{y}_N)\right|0_{B}^{}`$
via the boundary reduction formula bLSZ . The time evolution of the state is generated by the Hamiltonian
$`H_B`$ $`=`$ $`{\displaystyle _{\mathrm{}}^0}dx{\displaystyle }d\stackrel{}{y}[{\displaystyle \frac{1}{2}}\mathrm{\Pi }^2+{\displaystyle \frac{1}{2}}\left(_x\mathrm{\Phi }\right)^2+{\displaystyle \frac{1}{2}}\left(\stackrel{}{}\mathrm{\Phi }\right)^2`$
$`+V(\mathrm{\Phi })+\delta (x)V_B(\mathrm{\Phi })]`$
In the second description of the boundary theory we can take $`\tau =ix`$ as time and $`\xi =it`$ as space variable in the Hamiltonian formalism <sup>1</sup><sup>1</sup>1In the Euclidean formalism this corresponds to swapping Euclidean time with the space variable $`x`$, which is manifestly a symmetry of the theory in the bulk.. The Hilbert space now is that of the bulk theory (without any boundary), and time evolution is given by the bulk Hamiltonian (this is the extension of the crossing symmetry of the bulk theory, and is analogous to the open-closed duality in string theory*)*. In this case the boundary condition appears in time and serves as an initial state in calculating correlators:
$`G(\tau _1,\xi _1,\stackrel{}{y}_1;\mathrm{};\tau _N,\xi _N,\stackrel{}{y}_N)`$ $`=`$
$`0\left|\mathrm{\Phi }(\tau _1,\xi _1,\stackrel{}{y}_1)\mathrm{}\mathrm{\Phi }(\tau _N,\xi _N,\stackrel{}{y}_N)\right|B`$
where now $`|B`$ is a state in the bulk Hilbert space $``$. The boundary state $`|B`$ is in fact defined by the equality of correlation functions calculated in the two pictures. Using asymptotic completeness it can be expanded in the basis of the asymptotic *in* states of the bulk theory
$`|B`$ $`=`$ $`\{1+K_1A_{in}^{}(0,0)+{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\theta }{2\pi }}{\displaystyle }{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{\left(2\pi \right)^{D1}}}`$
$`K_2(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_{in}^{}(\theta ,\stackrel{}{k}_{||})A_{in}^{}(\theta ,\stackrel{}{k}_{||})+\mathrm{}\}|0`$
where $`A^{}(\theta ,\stackrel{}{k}_{||})`$ is the *in* asymptotic creation operator and the ellipses denote the terms with higher particle number. Due to translational invariance in the spatial direction all the contributing states must have zero total momentum. Ghoshal and Zamolodchikov have shown GZ that the coefficient $`K_2`$ is related to the one-particle elastic reflection factor $`R`$ on the boundary by
$$K_2(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))=R(\frac{i\pi }{2}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))$$
(here we simply treat $`\stackrel{}{k}_{||}`$ as a label for infinitely many species of two-dimensional particles of mass $`m_{\mathrm{eff}}(\stackrel{}{k}_{||})`$, using the fact that $`\stackrel{}{k}_{||}`$ is conserved in the reflection off the boundary). The one-particle coefficient $`K_1`$ is only nonzero when the vacuum expectation value of the interpolating field $`\mathrm{\Phi }`$ of the particle does not vanish. In this work we suppose that $`{}_{B}{}^{}0\left|\mathrm{\Phi }(t,x,\stackrel{}{y})\right|0_{B}^{}=0`$ (which is the case for the usual applications). We plan to present the general case in a forthcoming publication, together with more detailed derivation of the form of the coefficient $`K_1`$ and $`K_2`$ using the boundary reduction formula obtained in bLSZ , and a generalization to more complex particle spectra with several different masses. Here we restrict ourselves to a vanishing $`K_1`$ to keep the discussion simple. We remark that when the theory in the bulk is free and the reflection is elastic, the boundary state can be written in a closed form <sup>2</sup><sup>2</sup>2In 1+1 dimensions this can be extended to any integrable QFT with integrable boundary condition GZ .
$`|B`$ $`=`$ $`\mathrm{exp}({\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\theta }{2\pi }}{\displaystyle }{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{\left(2\pi \right)^{D1}}}`$
$`K_2(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_{in}^{}(\theta ,\stackrel{}{k}_{||})A_{in}^{}(\theta ,\stackrel{}{k}_{||}))|0`$
We note also that if we have more then one particle species, with the same mass $`m`$ created by $`A_{in}^{}(\theta ,\stackrel{}{k}_{||})_j`$ (i.e. a multiplet), the formula for the boundary state changes as
$`|B`$ $`=`$ $`\{1+{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\theta }{2\pi }}{\displaystyle }{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{\left(2\pi \right)^{D1}}}`$
$`K_2^{ij}(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_{in}^{}(\theta ,\stackrel{}{k}_{||})_iA_{in}^{}(\theta ,\stackrel{}{k}_{||})_j+\mathrm{}\}|0`$
Boundary conditions considered in the context of the Casimir effect generally allow transmission as well, and such boundaries are called ’defects’. A suitable generalization of the above formalism can be obtained by a folding trick, which maps the defect into a boundary system defect . Suppose now that a defect is located at $`x_0`$. In the crossed channel picture it can be represented by a defect operator which acts from the bulk Hilbert space of the $`x<x_0`$ system into that of the $`x>x_0`$ system. Lets denote the operator creating the particle for the $`x<x_0`$ domain as $`A_1^{}`$ while for the $`x>x_0`$ domain as $`A_2^{}`$. There are now four one-particle reflection amplitudes: $`R^\pm `$ are the ones preserving the species number $`1,2`$, while $`T^\pm `$ are the ones changing $`1`$ into $`2`$ and $`2`$ into $`1`$, respectively. Upon the folding correspondence, $`R^\pm `$ correspond to reflection processes on the two sides of the defect, while $`T^\pm `$ describe the transmission amplitudes from one side to the other. Using the folding map to the boundary system we obtain the defect operator defect as
$`D`$ $`=`$ $`1+{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta }{4\pi }}{\displaystyle }{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{\left(2\pi \right)^{D1}}}(`$
$`R^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1^{}(\theta ,\stackrel{}{k}_{||})A_1^{}(\theta ,\stackrel{}{k}_{||})+`$
$`T^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1^{}(\theta ,\stackrel{}{k}_{||})A_2(\theta ,\stackrel{}{k}_{||})+`$
$`T^{}({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1(\theta ,\stackrel{}{k}_{||})A_2^{}(\theta ,\stackrel{}{k}_{||})+`$
$`R^{}({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_2(\theta ,\stackrel{}{k}_{||})A_2(\theta ,\stackrel{}{k}_{||}))+\mathrm{}`$
which (with the same conditions as for the boundary state, but now elasticity is required for the combined one-particle reflection/transmission amplitude mussardo <sup>3</sup><sup>3</sup>3We remark that integrable defects (which are exactly the ones for which the defect operator can be exponentiated) with nontrivial reflection and transmission at the same time are only possible when the bulk scattering is trivial mussardo . This only means a restriction in $`1+1`$ dimensions, where integrable theories with nontrivial bulk scattering matrices are possible. In higher space-time dimension the conditions for the exponential form of the defect operator include the triviality of the bulk $`S`$ matrix anyway.) can be exponentiated to the form
$`D`$ $`=`$ $`\mathrm{exp}\{{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta }{4\pi }}{\displaystyle }{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{\left(2\pi \right)^{D1}}}`$
$`(R^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1^{}(\theta ,\stackrel{}{k}_{||})A_1^{}(\theta ,\stackrel{}{k}_{||})+`$
$`T^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1^{}(\theta ,\stackrel{}{k}_{||})A_2(\theta ,\stackrel{}{k}_{||})+`$
$`T^{}({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_1(\theta ,\stackrel{}{k}_{||})A_2^{}(\theta ,\stackrel{}{k}_{||})`$
$`+R^{}({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))A_2(\theta ,\stackrel{}{k}_{||})A_2(\theta ,\stackrel{}{k}_{||}))\}`$
### II.2 Derivation of Casimir energy
Let us now turn to the derivation of Casimir energy of a $`D+1`$ dimensional scalar field $`\mathrm{\Phi }(t,x,\stackrel{}{y})`$ in a domain of width $`L`$ in $`x`$. To facilitate later applications, it is useful to consider the case when the field $`\mathrm{\Phi }(t,x,\stackrel{}{y})`$ is allowed to penetrate through the ends of the domain in such a way, that in the two domains of width $`R`$ adjoining $`L`$, the dispersion relation of the asymptotic particles may be different from the vacuum one in $`L`$. The totally reflecting boundaries at $`x=\pm (R+L/2)`$ are treated as auxiliary boundary conditions which are necessary in order to have a discrete spectrum; at the end of calculations we shall take the limit $`R\mathrm{}`$ and check that the results are independent of the choice of the auxiliary boundary conditions. So let us take the following situation in the coordinate $`x`$
where total reflection occurs at $`x_l=RL/2`$ and $`x_r=R+L/2`$ with reflection factors $`R_l(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))`$ and $`R_r(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))`$ and there are two “defects” (i.e. boundaries allowing transmission) located at $`x_1=L/2`$ and $`x_2=L/2`$ with one-particle defect matrices
$$D_i(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))=\left(\begin{array}{cc}R_i^+(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))& T_i^{}(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))\\ T_i^+(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))& R_i^{}(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))\end{array}\right)$$
where $`i=1,2`$. The Hamiltonian of the system is denoted by $`H_B`$ and its Hilbert space by $`_B.`$ It is the ground state energy of the system which is of interest for calculating the Casimir effect, i.e. the lowest eigenvalue of $`H_B`$, which can be evaluated via the partition function. Compactifying all infinite (temporal and spatial) dimensions on circles with perimeter $`T`$ we can calculate the partition function of our system in two different ways:
$`Z_R(L,T)`$ $`=`$ $`\mathrm{Tr}__B\mathrm{e}^{TH_B}`$
$`=`$ $`B_l\left|\mathrm{e}^{RH_x^{(1)}}D_1\mathrm{e}^{LH_x^{(2)}}D_2\mathrm{e}^{RH_x^{(3)}}\right|B_r`$
where $`H_x^{(i)}`$ is the periodic Hamiltonian in the $`x`$ channel in the three domains indexed by $`i=1,2,3`$. Inserting complete sets of bulk states, taking the limit $`R\mathrm{}`$ and normalizing the ground state energy in infinite space to $`E_0=0`$ we obtain
$`Z_{\mathrm{}}(L,T)`$ $`=`$ $`{\displaystyle \underset{n}{}}B_l|00\left|D_1\right|nn\left|D_2\right|00|B_r\mathrm{e}^{LE_n}`$
$`=`$ $`{\displaystyle \underset{n}{}}0\left|D_1\right|nn\left|D_2\right|0\mathrm{e}^{LE_n}`$
where the dependence on the auxiliary boundary conditions $`B_{l,r}`$ drops out (since in the $`R\mathrm{}`$ limit only the vacuum state contributes from the expansions of $`|B_{l,r}`$). The first few terms can be written explicitly as
$`1+{\displaystyle \underset{\theta ,\stackrel{}{k}_{||}}{}}{\displaystyle \underset{\theta ^{},\stackrel{}{q}_{||}}{}}0|D_1|\theta ,\stackrel{}{k}_{||};\theta ^{},\stackrel{}{q}_{||}\theta ,\stackrel{}{k}_{||};\theta ^{},\stackrel{}{q}_{||}|D_2|0\times `$
$`e^{L\left(m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta +m_{\mathrm{eff}}(\stackrel{}{q}_{||})\mathrm{cosh}\theta ^{}\right)}+O(e^{3mL})`$
The term $`1`$ is the contribution from the vacuum ($`|n=|0`$), the next term comes from two-particle terms in (II.1) and the higher-order corrections come from the higher multi-particle terms. This is a sort of cluster expansion similar to the one used in bluscher , valid for large values of the volume $`L`$. It is obvious from these expressions that the leading (two-particle) contribution depends only on $`R_1^{}`$ and $`R_2^+`$. The ground state (Casimir) energy (per unit transverse area) can be extracted from the partition function as
$$E(L)=\underset{T\mathrm{}}{lim}\frac{1}{T^D}\mathrm{log}Z_{\mathrm{}}(L,T)$$
The result is
$`E(L)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta }{4\pi }}\mathrm{cosh}\theta {\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{(2\pi )^{D1}}m_{\mathrm{eff}}(\stackrel{}{k}_{||})}`$ (5)
$`R_1^{}({\displaystyle \frac{i\pi }{2}}+\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))R_2^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))\times `$
$`\mathrm{e}^{2m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta L}+O(e^{3mL})`$
The correction terms correspond to higher particle terms in the expansion (II.1) of the defect operator $`D`$ and include the amplitudes of reflection/transmission processes involving more than one particle in at least one of the asymptotic states. These can be computed e.g. using a BQFT formulation as the one in bLSZ , but it is obvious that they are suppressed by a factor $`\mathrm{e}^{mL}`$ with respect to the leading order term due to the presence of at least one additional particle in the corresponding term of the expansion of the defect operator $`D`$.
The formula (5), describing Casimir effect to leading order at long distances, is the main result of this paper. Note that it is applicable in the presence of nontrivial bulk and boundary interactions: their effects at leading order are contained in the reflection factors $`R^\pm `$, so as long as there is some theoretical or experimental input from which these can be determined the leading order contribution can be evaluated. In integrable (2D) boundary theories $`R^\pm `$ are obtained as solutions of the boundary Yang Baxter equation, and the bulk interaction manifests itself thorough the bulk S matrix appearing in this equation. In case of nonintegrable interacting bulk theories e.g. a perturbative expansion can be given to determine $`R^\pm `$, (for more details see bLSZ ).
Another important point is that this approach formulates the Casimir effect from an infrared viewpoint. Standard derivations of the Casimir effect (such as the one presented in the appendix) solve the microscopic field theory. This necessitates tackling diverse issues such as renormalization, and also the possibility that the infrared (long distance behaviour) may be quite different from the microscopic description of the theory (as is the case for example in QCD). Formula (5) is expressed in terms of the asymptotic particles, the long distance degrees of freedom, and provides a long distance expansion for Casimir energy. It describes a finite size effect in a boundary quantum field theory close in spirit to our previous investigation of the boundary Lüscher formula bluscher , and it is not difficult to see that the latter is just a special case of (5). As a consequence our formula has already been tested in interacting (integrable) 2D quantum field theories in bluscher . (This is the only case to our knowledge where exact reflection factors of interacting BQFTs have been computed). Furthermore, the result (5) includes the contribution of states localized to the defects (called ’surface plasmons’ in the literature), since they are taken into account as poles at imaginary rapidity of the reflection factors. This fact is also demonstrated explicitely using the zero mode summation method in Appendix A.
A very appealing property of the boundary state approach is its universality and we shall see in the next section that it indeed reproduces all the results previously known for the planar situation. But in those cases we can go even further, because in most calculations of the Casimir effect the bulk is free and the boundary scattering is elastic, therefore one can use the exponentiated boundary state (II.1) together with the analogous defect operator (II.1) which makes it possible to sum up all the multi-particle terms of the cluster expansion. The calculation to be done is essentially identical to the derivation of the boundary Thermodynamic Bethe Ansatz equation with trivial bulk $`S`$ matrix and reflection factors $`R_1^{}`$ and $`R_2^+`$, which was already performed in lmss . The only difference is that lmss supposes fermionic statistics and that in our case the reflection factors $`R_1^{}`$ and $`R_2^+`$ are not unitary due to the existence of transmission, but these do not change the overall reasoning of the derivation. The result is
$`E(L)`$ $`=`$ $`\pm {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta }{4\pi }}\mathrm{cosh}\theta {\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{(2\pi )^{D1}}m_{\mathrm{eff}}(\stackrel{}{k}_{||})}`$ (6)
$`\mathrm{log}(1R_1^{}({\displaystyle \frac{i\pi }{2}}+\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))R_2^+({\displaystyle \frac{i\pi }{2}}\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))`$
$`\mathrm{e}^{2m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta L})`$
(the upper/lower signs correspond to the bosonic/fermionic case, respectively). Although formula (5) is more general and adequate enough for the large distance regime, the theoretical situations considered in the literature can be directly compared to (6). It is obvious that for $`Lm^1`$ (6) reduces to (5). The Casimir force is dominated by modes with momentum of the order $`1/L`$; for large enough separation the bulk interaction can be dropped due to its short-ranged nature, and at the same time the characteristic energy is lower than the threshold for inelastic boundary scattering.
## III Applications
To show the universal nature of these results we consider some applications. We demonstrate that taking the reflection factors from the literature and substituting into (6) immediately gives the Casimir energy without going through a detailed analysis of the microscopic quantum fluctuations.
As a first application we determine the Casimir energy for a massive free scalar field, $`V(\mathrm{\Phi })=\frac{m^2}{2}\mathrm{\Phi }^2`$, subject to Robin boundary conditions, $`V_B(\mathrm{\Phi })=\frac{c}{2}\mathrm{\Phi }^2`$, on two parallel hyperplanes. If the two planes are located at $`x=0`$ and $`x=L`$ the boundary conditions are given as
$$_x\mathrm{\Phi }c_1\mathrm{\Phi }|_{x=0}=0;_x\mathrm{\Phi }+c_2\mathrm{\Phi }|_{x=L}=0;c_1,c_20,$$
and the reflection amplitudes on these planes can be written as
$$R_j(\theta ,m_{\mathrm{eff}}(\stackrel{}{k}_{||}))=\frac{m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{sinh}\theta ic_j}{m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{sinh}\theta +ic_j};j=1,2.$$
Note that they have the same form as in the two dimensional case, but now they depend also on $`\stackrel{}{k}_{||}`$ via the rapidity parameterization $`k_{}=m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{sinh}\theta `$. When $`\theta `$ is continued to $`\theta \theta +i\frac{\pi }{2}`$, we have
$$R_j(\theta +i\frac{\pi }{2},m_{\mathrm{eff}}(\stackrel{}{k}_{||}))=\frac{m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta c_j}{m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta +c_j};j=1,2.$$
Introducing the variable $`q`$ by $`m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta =\sqrt{m^2+q^2}`$ and performing the angular integrations one obtains
$`E(L)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^{D/2}\mathrm{\Gamma }(D/2)}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑qq^{D1}`$
$`\mathrm{log}\left(1{\displaystyle \frac{\sqrt{m^2+q^2}c_1}{\sqrt{m^2+q^2}+c_1}}{\displaystyle \frac{\sqrt{m^2+q^2}c_2}{\sqrt{m^2+q^2}+c_2}}e^{2L\sqrt{m^2+q^2}}\right).`$
This formula is a new result of the present paper. To show its correctness we consider two already known limiting cases. First consider the limit when both $`c_j0`$ or $`c_j\mathrm{}`$ corresponding to Neumann or Dirichlet boundary conditions for $`\mathrm{\Phi }`$. In both cases the coefficient of the exponent in the logarithm becomes one, and we obtain the Ambjorn-Wolfram result for Dirichlet boundary conditions ambjorn as reported in Milton’s book book . In the second limit we let the mass of the scalar field vanish $`m0`$; then
$`E(L)`$ $`=`$
$`{\displaystyle \frac{1}{(4\pi )^{D/2}\mathrm{\Gamma }(D/2)}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑qq^{D1}\mathrm{log}\left(1{\displaystyle \frac{qc_1}{q+c_1}}{\displaystyle \frac{qc_2}{q+c_2}}e^{2Lq}\right).`$
It is straightforward to show that this result coincides with that obtained in robin for the Casimir energy per unit area of a massless scalar field with Robin boundary conditions.
The second application concerns the Casimir force between parallel dielectric slabs separated by a vacuum slot of width $`L`$, a problem first investigated by Lifshitz and collaborators almost fifty years ago lifshitz . In this geometry the dielectric constants in the three regions are:
$`ϵ(x)=ϵ_1;x<0;`$ $`ϵ(x)=ϵ_2;L<x;`$
$`ϵ(x)=1;0<x<L.`$
In this case $`R_{1,2}^\pm (i\frac{\pi }{2}+\theta )`$ are nothing else but the (appropriate analytic continuations of the) ordinary reflection amplitudes of electromagnetic waves incident from vacuum at the plane interface between the vacuum and the dielectric materials, given in many textbooks (e.g. Jackson ). Indeed denoting again by $`(k_{},\stackrel{}{k_{||}})`$, $`\omega `$ the wave vector and frequency of the electromagnetic radiation ($`\stackrel{}{k}_{||}`$ being the component parallel, while $`k_{}`$ the one perpendicular to the plane interface) the reflection amplitudes are
$$R_{\mathrm{perp}}^{(i)}(\omega ,\stackrel{}{k}_{||})=\frac{\sqrt{\omega ^2\stackrel{}{k}_{||}^2}\sqrt{ϵ_i\omega ^2\stackrel{}{k}_{||}^2}}{\sqrt{\omega ^2\stackrel{}{k}_{||}^2}+\sqrt{ϵ_i\omega ^2\stackrel{}{k}_{||}^2}},ϵ_i=ϵ_i(\omega ),$$
when the electric field $`\stackrel{}{E}`$ is perpendicular to the plane of incidence, and
$$R_{\mathrm{par}}^{(i)}(\omega ,\stackrel{}{k}_{||})=\frac{ϵ_i\sqrt{\omega ^2\stackrel{}{k}_{||}^2}\sqrt{ϵ_i\omega ^2\stackrel{}{k}_{||}^2}}{ϵ_i\sqrt{\omega ^2\stackrel{}{k}_{||}^2}+\sqrt{ϵ_i\omega ^2\stackrel{}{k}_{||}^2}},ϵ_i=ϵ_i(\omega ),$$
when $`\stackrel{}{E}`$ is parallel to the plane of incidence and $`i=1,2`$. (To obtain these expressions we assumed that the permeabilities of the slabs are unity, $`\mu _i=1`$). Continuing $`\theta \theta +i\frac{\pi }{2}`$ corresponds to continuing $`\omega `$ to purely imaginary values $`\omega m_{\mathrm{eff}}(\stackrel{}{k}_{||})i\mathrm{sinh}\theta =i\zeta ,`$ and the reflection amplitudes become ($`q^2=\stackrel{}{k}_{||}^2`$):
$`R_{\mathrm{perp}}^{(i)}(i\zeta ,q)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\zeta ^2+q^2}\sqrt{ϵ_i\zeta ^2+q^2}}{\sqrt{\zeta ^2+q^2}+\sqrt{ϵ_i\zeta ^2+q^2}}},`$
$`R_{\mathrm{par}}^{(i)}(i\zeta ,q)`$ $`=`$ $`{\displaystyle \frac{ϵ_i\sqrt{\zeta ^2+q^2}\sqrt{ϵ_i\zeta ^2+q^2}}{ϵ_i\sqrt{\zeta ^2+q^2}+\sqrt{ϵ_i\zeta ^2+q^2}}}.`$
Since the electromagnetic field can have both polarizations, one obtains the following form for the Casimir energy per unit area
$`E(L)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dq^2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\zeta [`$
$`\mathrm{log}\left(1R_{\mathrm{perp}}^{(1)}(i\zeta ,q)R_{\mathrm{perp}}^{(2)}(i\zeta ,q)\mathrm{e}^{2L\sqrt{q^2+\zeta ^2}}\right)+`$
$`\mathrm{log}(1R_{\mathrm{par}}^{(1)}(i\zeta ,q)R_{\mathrm{par}}^{(2)}(i\zeta ,q)\mathrm{e}^{2L\sqrt{q^2+\zeta ^2}})].`$
The Casimir force $`(L)=E(L)/L`$ computed from this expression agrees with that of Lifshitz et al. as reported in the reviews reviews ; book .
Next we consider the Casimir energy of a massless fermion field in $`1+3`$ dimensions subject to the “bag boundary condition” in our planar geometry:
$$(1i\gamma ^3)\psi |_{x=0}=0,(1+i\gamma ^3)\psi |_{x=L}=0.$$
Using the chiral representation of Dirac matrices one readily shows that these boundary conditions break chirality but commute with $`\mathrm{\Sigma }_3=\left(\begin{array}{cc}\sigma _3& 0\\ 0& \sigma _3\end{array}\right)`$. Furthermore they imply that for both the positive ($`u`$) and the negative ($`v`$) eigenvectors of $`\mathrm{\Sigma }_3`$ the reflection amplitudes are constants, satisfying $`R_{1(s)}^{}(i\pi /2+\theta )R_{2(s)}^+(i\pi /2+\theta )=1`$ for $`s=u,v`$. Thus (taking into account crossing-unitarity (23))
$$E(L)=2\frac{1}{4\pi ^2}\underset{0}{\overset{\mathrm{}}{}}𝑑qq^2\mathrm{log}(1+e^{2Lq}),$$
which indeed coincides with the known result book ; fermion .
## IV Concluding remarks
The main result of the paper is formula (5) which expresses the Casimir energy in large volume in a general form valid for any QFT - even for theories interacting in the bulk - in terms of the reflection amplitudes (infrared data). It passed the test of interacting but integrable two dimensional quantum field theories in bluscher . The boundary state approach provides a systematic infrared (large volume) expansion, which in the case of noninteracting bulk and elastic boundary scattering can be summed up leading to (6). In order to check the result for noninteracting bulk we present in the appendix an alternative derivation using the naive mode summation method.
In the usual calculations of the Casimir effect corrections for temperature and roughness effects must be made before comparing to the experiment. Both corrections can be introduced into (5,6) in a straightforward manner, following the usual procedure in the literature (cf. book ) see also rough . In the case of a massive theory with mass gap $`m`$ the terms coming from higher particle corrections to (II.1) and (II.1) are suppressed by $`\mathrm{exp}(mL)`$. For theories with a zero mass gap the situation is different and indeed the asymptotic state formalism is not well-founded if interaction with massless particles (such as photons) is taken into account, due to infrared divergencies. In the case of electromagnetic field the corrections to (6) are known to be suppressed by $`\alpha \lambda /L`$ where $`\lambda `$ is the Compton wave length of the electron and $`\alpha `$ is the fine structure constant (Chapter 13 of book ), due to the fact that the self-interaction of the electromagnetic field arises only through radiative corrections and the leading one is an electron loop.
The striking fact about the general formulae (5,6) is that they are manifestly finite, and only depend on physical quantities, such as the dispersion relation of asymptotic particles and their reflection amplitudes on the boundaries. This is even more emphasized in the boundary state formalism, where the bulk and boundary divergences appearing in usual derivations of the Casimir effect are manifestly absent, and the whole derivation is performed in terms of finite (renormalized) physical objects. The appearance of divergences in the naive method of summing zero-point energies is a general feature when one calculates the interaction energy of point-like sources by integrating the energy stored in the field generated by them. The divergences are then eliminated by noticing that they contribute to the self-energy of the sources themselves and are present even at infinite separation. These divergences are absorbed by renormalization, but they are entirely unphysical and eliminated when expressing the measurable energy differences in terms of physically meaningful, finite quantities. Further bonus of the general formulae (5,6) is that the reflection factors appearing in them can be calculated easily in the semi-infinite geometry, e.g. using perturbation theory along the lines of bLSZ .
## Acknowledgments
This research was partially supported by the EC network “EUCLID”, contract number HPRN-CT-2002-00325, and Hungarian research funds OTKA D42209, T037674, T043582 and TS044839. GT was also supported by a Széchenyi István Fellowship.
## Appendix A Casimir energy by mode summation
Here we want to support our derivation coming from the boundary state formalism by obtaining the analogous result using a new version of the mode summation method for fluctuating fields which are free in the bulk. Ideas somewhat similar to these have been used in a similar framework earlier in JR and their findings also support our general expression. The novel viewpoints of our approach are twofold: first we concentrate only on the finite size piece of the ground state energy that should be independent of regularization and renormalization thus we do not need to discuss the regularization/renormalization of the (infinite space) bulk/boundary ground state energies, since eventually these are subtracted. The second feature of our approach is that we do not solve the quantization condition for the zero point frequencies explicitly but use only the functional form of these equations to obtain a general expression for the ground state energy containing only the reflection amplitudes.
For simplicity let us start with the $`1+1`$ dimensional free scalar field $`\mathrm{\Phi }(t,x)`$ of mass $`m`$ in a domain described earlier. We assume that in all three domains the frequency of the modes of $`\mathrm{\Phi }`$ is the same, but the wave vectors $`k^{}`$, $`k^{\prime \prime }`$ may be different from $`k`$ (e.g. the mass of the scalar field is different in the three domains), but they can be regarded as functions of $`k`$. All the wave number parameters $`k,k^{},k^{\prime \prime }`$ are chosen positive. Then for the plane wave modes one has the situation displayed in the following figure
where total reflection occurs at $`x_l=RL/2`$ and $`x_r=R+L/2`$ with reflection factors $`R_l`$ and $`R_r`$
$`Ae^{ik^{}(R+\frac{L}{2})}`$ $`=`$ $`Be^{ik^{}(R+\frac{L}{2})}R_l(k^{})`$ (7)
$`Fe^{ik^{\prime \prime }(R+\frac{L}{2})}`$ $`=`$ $`Ee^{ik^{\prime \prime }(R+\frac{L}{2})}R_r(k^{\prime \prime }),`$
and the two “defects” (i.e. boundaries allowing transmission) located at $`x_1=L/2`$ and $`x_2=L/2`$ describe one-particle scattering with defect matrices
$$D_i(k)=\left(\begin{array}{cc}R_i^+(k)& T_i^{}(k)\\ T_i^+(k)& R_i^{}(k)\end{array}\right),i=1,2$$
In the case of elastic defect scattering their is no particle production and the defect matrices are unitary. They connect the incoming and outgoing amplitudes as follows
$`\left(\begin{array}{c}B\mathrm{e}^{+ik^{}L/2}\\ C\mathrm{e}^{ikL/2}\end{array}\right)`$ $`=`$ $`D_1(k)\left(\begin{array}{c}A\mathrm{e}^{ik^{}L/2}\\ D\mathrm{e}^{+ikL/2}\end{array}\right)`$ (12)
$`\left(\begin{array}{c}D\mathrm{e}^{ikL/2}\\ E\mathrm{e}^{+ik^{\prime \prime }L/2}\end{array}\right)`$ $`=`$ $`D_2(k)\left(\begin{array}{c}C\mathrm{e}^{+ikL/2}\\ F\mathrm{e}^{ik^{\prime \prime }L/2}\end{array}\right)`$ (17)
Unitarity and time-reversal gives the relations
$`D_i(k)^{}D_i(k)=1`$ , $`D_i(k)^1=D_i(k)`$ (18)
$`R_{r,l}(k)^{}R_{r,l}(k)=1`$ , $`R_{r,l}(k)^1=R_{r,l}(k).`$
Consistency of the homogenous linear system (7,12) is expressed by the quantization condition
$$Q(k_n)=0$$
(19)
where
$`Q\left(k\right)=\left(1R_l\left(k\right)R_1^+\left(k\right)e^{i2k^{}R}\right)\left(1R_r\left(k\right)R_2^{}\left(k\right)e^{i2k^{\prime \prime }R}\right)`$
$`e^{i2kL}R_1^{}\left(k\right)R_2^+\left(k\right)\left(1{\displaystyle \frac{R_l\left(k\right)e^{i2k^{}R}}{R_1^+\left(k\right)}}\right)\left(1{\displaystyle \frac{R_r\left(k\right)e^{i2k^{\prime \prime }R}}{R_2^{}\left(k\right)}}\right)`$
and we exploited (18) to write
$$detD_i=R_i^{}(k)/R_i^+(k)=R_i^+(k)/R_i^{}(k)$$
Due to unitarity the quantization condition can only have real and purely imaginary solutions. For each real solution $`k_n`$, $`k_n`$ is also a solution, and $`k_0=0`$ always solves (19). The purely imaginary solutions located at $`k_j=i\kappa _j`$ ($`0<\kappa _j<m`$) are related to the poles of the reflection and transmissions factors and correspond to defect and boundary bound states. Then the ground state energy of the system is given by the sum of zero modes, which can be written in an integral form
$`E(L,R)={\displaystyle \underset{\{k_n>0\}}{}}{\displaystyle \frac{1}{2}}\omega (k_n)+{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{2}}\omega (i\kappa _j)=`$
$`{\displaystyle \frac{1}{2}}\left\{{\displaystyle \frac{1}{2}}\left({\displaystyle _{C_+}}+{\displaystyle _C_{}}{\displaystyle _{C_0}}\right)+{\displaystyle \underset{j}{}}{\displaystyle _{C_j}}\right\}{\displaystyle \frac{dk}{2\pi i}}{\displaystyle \frac{Q^{}(k)}{Q(k)}}\omega (k)`$
where $`Q^{}(k)=\frac{dQ}{dk}`$, $`\omega (k)=\sqrt{k^2+m^2}`$, $`\mathrm{}`$ is set to one and the contours are shown on the figure
The poles on the figure denote the poles of the $`Q`$ function originating from the poles of the reflection and transmission amplitudes. Contrary to $`\kappa _j`$, which do depend on the volume, the location of these poles is independent of $`L,R`$.
We concentrate on the $`L`$ dependence so the $`L`$ independent $`C_0`$ integral can be dropped. Next we analyze the contribution of $`C_{}`$. Using (18) we can relate
$`{\displaystyle \frac{Q^{}(k)}{Q(k)}}{\displaystyle \frac{Q^{}(k)}{Q(k)}}`$ $`=`$ $`2iL2iR\left({\displaystyle \frac{dk^{}}{dk}}+{\displaystyle \frac{dk^{\prime \prime }}{dk}}\right)`$
$`{\displaystyle \frac{d}{dk}}\mathrm{log}R_r{\displaystyle \frac{d}{dk}}\mathrm{log}(detD_1D_2){\displaystyle \frac{d}{dk}}\mathrm{log}R_l.`$
The first two terms on the right hand side (proportional to the volumes $`L`$, $`R`$) correspond to divergent bulk energy contributions. The other terms here correspond to infinite boundary and defect energies which are independent of the volumes ($`L`$, $`R`$). The precise definition of these terms needs some regularization and renormalization, but they give no contribution to the Casimir force since they are present in the infinite system as well. Therefore they can be dropped and then the $`C_+`$ and $`C_{}`$ integral terms give identical results.
Thus the relevant contribution to the Casimir energy can be written in terms of the variable $`m\mathrm{sinh}\theta =k`$ as
$$E(L,R)=m_{\mathrm{}+i\eta }^{\mathrm{}+i\eta }\frac{d\theta }{4\pi i}\mathrm{cosh}\theta \frac{Q^{}(\theta )}{Q(\theta )}+\underset{j}{}\frac{1}{2}\omega (i\kappa _j).$$
(where $`\eta `$ is a small positive parameter). Next we shift the contour to $`\eta =\frac{\pi }{2}`$ which eliminates the discrete sum via the contributions from the purely imaginary zeros of $`Q(\theta )`$. There can be additional contributions from poles of $`Q\left(\theta \right)`$ on the imaginary axis, but their positions are independent of the volume and therefore they only contribute volume independent terms which can be dropped. Integrating by parts and introducing the new variable $`\theta \frac{i\pi }{2}+\theta `$ we obtain
$$E(L,R)=m_{\mathrm{}}^{\mathrm{}}\frac{d\theta }{4\pi }\mathrm{cosh}\theta \mathrm{log}Q(\frac{i\pi }{2}+\theta ).$$
In the limit $`R\mathrm{}`$
$`E(L)`$ $`=`$ $`m{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta }{4\pi }}\mathrm{cosh}\theta `$
$`\mathrm{log}\left(1R_1^{}\left({\displaystyle \frac{i\pi }{2}}+\theta \right)R_2^+\left({\displaystyle \frac{i\pi }{2}}+\theta \right)\mathrm{e}^{2m\mathrm{cosh}\theta L}\right).`$
It is straightforward to generalize these expressions to the $`D+1`$ dimensional case when the boundary conditions are specified on two parallel $`D1`$ dimensional hyperplanes. Introducing the component parallel ($`\stackrel{}{k}_{||}`$) respectively perpendicular ($`k_{}`$) to the hyperplanes one can write
$`\stackrel{}{k}`$ $`=`$ $`(k_{},\stackrel{}{k}_{||});\omega ^2k_{}^2=m^2+\stackrel{}{k}_{||}^2m_{\mathrm{eff}}^2(\stackrel{}{k}_{||});`$ (20)
$`k_{}=m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{sinh}\theta .`$
We assume that the system is translationally invariant in the direction of the hyperplanes, and so the reflection and transmission processes preserve $`\stackrel{}{k}_{||}`$. Therefore we can apply the previous considerations to the decoupled one-dimensional systems labeled by the parameter $`\stackrel{}{k}_{||}`$ and sum up their contributions. Note that both $`\omega `$ and the reflection/transmission amplitudes depend on $`\stackrel{}{k}_{||}`$, but we shall not write this out explicitely. The Casimir energy per unit transverse area can be written as
$`E(L,R)`$ $`=`$
$`{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{(2\pi )^{D1}}m_{\mathrm{eff}}(\stackrel{}{k}_{||})_{\mathrm{}}^{\mathrm{}}\frac{d\theta }{4\pi }\mathrm{cosh}\theta \mathrm{log}Q\left(\frac{i\pi }{2}+\theta \right)}.`$
From this, the Casimir energy between the two defects can be obtained by taking the limit $`R\mathrm{}`$
$`E(L)`$ $`=`$ $`{\displaystyle \frac{d^{D1}\stackrel{}{k}_{||}}{(2\pi )^{D1}}m_{\mathrm{eff}}(\stackrel{}{k}_{||})_{\mathrm{}}^{\mathrm{}}\frac{d\theta }{4\pi }\mathrm{cosh}\theta }`$
$`\mathrm{log}\left(1R_1^{}\left({\displaystyle \frac{i\pi }{2}}+\theta \right)R_2^+\left({\displaystyle \frac{i\pi }{2}}+\theta \right)\mathrm{e}^{2m_{\mathrm{eff}}(\stackrel{}{k}_{||})\mathrm{cosh}\theta L}\right).`$
Note that the latter result is independent of the auxiliary boundary conditions $`R_r`$ and $`R_l`$ which can be expected on physical grounds. Taking into account the crossing-unitarity property satisfied by the reflection amplitudes GZ ; defect which for free bosons/fermions takes the form
$$R\left(\frac{i\pi }{2}+\theta \right)=\pm R\left(\frac{i\pi }{2}\theta \right)$$
(23)
one can see that this result is in a complete agreement with the formula obtained form the boundary state formalism previously (6). (The results in (A, A) should be multiplied by $`(1)`$ in case of computing the Casimir energy of fermionic fields, since their vacuum energy is $`\frac{1}{2}\omega (k)`$).
If there are more than one fields in the problem, or the fields have more than one component one has to sum up the contribution of all the fields (components) to obtain the Casimir energy.
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# Quantum 𝐿𝐶 circuits with charge discreteness: normal and anomalous spectrum
(Departamento de Física, Universidad de Tarapacá, Casilla 7-D, Arica, Chile)
Abstract: A solvable quantum $`LC`$ circuit with charge discreteness is studied. Two discrete spectral branches are obtained: (i) the normal branch corresponding to a charged capacitor with integer effective charge $`k=q_en`$ ($`q_e`$ elementary charge, $`n`$ integer) and (ii) the anomalous branch where the energy is related to non-integer effective charge $`k=q_e(nx)`$ in the capacitor. For usual mesoscopics data like quantum point contact we found $`x0.7`$.
Published in Europhys.Lett. 69,116 (2005).
PACS: 73.23.-b, 73.63.-b, 73.21.-b
Keywords: Electronic transport. Nanoscale. Mesoscopic.
In mesoscopic systems the discrete nature of the electrical charge plays an important role ; phenomena like Coulomb blockage, current magnification, spectral properties of wires etc. are related to it. In addition, the concept of an elementary quantum flux $`h/e`$ arises naturally from charge discreteness. On the other hand, some of the above mentioned systems, and others, could be studied through the analogy with simple electrical circuits with capacitances and inductances \[2-7\]. From this point of view, the quantization of circuits with charge discreteness plays an important role in mesoscopic physics and this paper is related to this topic. We emphasize that charge discreteness is assumed in this paper as an experimental fact and therefore, it is not our goal to prove it from first principles.
Consider a quantum $`LC`$ circuit with charge operator $`\widehat{Q}`$, magnetic flux operator $`\widehat{\varphi }`$ and the commutation rule $`[\widehat{Q},\widehat{\varphi }]=i\mathrm{}`$. In close analogy with the known quantum problem of a particle in a box, where the momentum is quantized, boundary conditions on the state $`\psi (\varphi )`$ must be imposed. So, on the wave function we must impose the charge discreteness condition related to
$$\psi (\varphi +\frac{h}{q_e})=\psi (\varphi ),$$
(1)
where $`q_e`$ is the elementary charge. Note that (1) means that the system must be translation invariant in the $`\varphi `$-space. Moreover, the eigenvectors of the charge operator $`\widehat{Q}=i\mathrm{}/\varphi `$ are $`e^{iq_e\varphi j/\mathrm{}}`$, with eigenvalues $`jq_e`$ ($`j`$ integer ) and compatible with charge discreteness. In this way, the state of the systems in charge representation $`\psi _j`$ and flux representation $`\psi (\varphi )`$ are connected by the transformation
$$\psi (\varphi )=\underset{j}{}e^{i\varphi q_ej/\mathrm{}}\psi _j.$$
(2)
We note that, in the flux space, eqs. (1) and (2) are consequences of charge discreteness and the fact that the charge operator and flux operator are canonical conjugate.
The Hamiltonian $`\widehat{H}`$ for the $`LC`$ circuit with charge discreteness requires compatibility with the boundary condition (1). Note that when $`q_e=0`$ the $`LC`$-Hamiltonian is given by
$$\widehat{H}_{q_e=0}=\frac{1}{2L}\widehat{\varphi }^2+\frac{1}{2C}\widehat{Q}^2,$$
(3)
where electron-electron interaction is automatically considered in the capacitor properties.
The periodic boundary condition (1), tells us that the Hamiltonian has the generic form
$$\widehat{H}_{q_e}=\frac{2\mathrm{}^2}{q_e^2L}P(\frac{q_e}{h}\widehat{\varphi })+\frac{1}{2C}\widehat{Q}^2,$$
(4)
where $`P(x)`$ is a periodic function of period $`1`$ (translation invariance). The factor $`2\mathrm{}^2/q_e^2L`$ was incorporated to make contact with references \[3-7\], were charge discretization was considered. In fact, references \[3-7\] consider charge discreteness with the choice $`P(x)=\mathrm{sin}^2(\pi x/2)`$. This is directly related to the discretization of the operator $`i\mathrm{}/Q`$ as finite differences. Nevertheless, it is important to realize that there are other possible choices for discretization procedure which could be used to represent charge discreteness. Namely, any system described by the Hamiltonian (4), with the additional condition (1) imposed upon it, represents a discrete charge quantum $`LC`$ circuit. Note that the inductance $`L`$ and capacitance $`C`$ are fixed parameters in our theory. Nevertheless, an equivalent approach could be assumed if, for instance, the inductance in eq.(3) is flux-depending (effective inductance) in correspondence with (4).
In this paper we shall study a solvable model for the $`LC`$ circuit, described by a Hamiltonian with charge discreteness, and closely related to quantum point contact systems. We shall consider a Hamiltonian with the structure of (4) and related to the Schrödinger equation in flux representation ( $`\widehat{Q}=i\mathrm{}/\varphi `$):
$$E\psi \left(\varphi \right)=\frac{2\mathrm{}^2}{q_e^2L}\left\{\underset{l}{}\delta \left(\frac{q_e}{h}\varphi l\right)\right\}\psi \left(\varphi \right)\frac{\mathrm{}^2}{2C}\frac{^2}{\varphi ^2}\psi \left(\varphi \right).$$
(5)
The above equation is the basis of our calculations and is closely related to the Kronig-Penney model used in Solid State Physics, nevertheless, the condition (1) is more restrictive than the usual one (Bloch theorem). In fact, translation invariance of (5) tell us that $`\psi (\varphi +h/q_e)=\psi (\varphi )exp(iK2\pi /q_e)`$ (Bloch) and then from (1) we must consider only $`K=0`$. So, some points, or levels, of the band structure of (5) are related to the condition of charge discreteness.
As usual for an equation like (5), between two $`\delta `$barriers (at position $`l/hq_e`$ and $`\left(l+1\right)/hq_e`$, $`l`$ integer) we have the equation $`E\psi \left(\varphi \right)=\frac{\mathrm{}^2}{2C}\frac{^2}{\varphi ^2}\psi \left(\varphi \right)`$ and its solution is given by
$$\psi _l(\varphi )=A_le^{ik\varphi /\mathrm{}}+B_le^{ik\varphi /\mathrm{}},$$
(6)
where the effective charge $`k`$ is a parameter with dimension of electrical charge. In term of this parameter, the energy of the system becomes
$$E=\frac{1}{2C}k^2,$$
(7)
corresponding to the energy of a capacitor. As a consequence of the boundary condition (1), the parameter $`k`$ is not continuous. In fact, the matching conditions in a barrier $`\psi (\frac{h}{q_e}l)_+=\psi (\frac{h}{q_e}l)_{}`$ (by left and right side) and the step in the first derivative $`\psi ^{}(\frac{h}{q_e}l)_+\psi ^{}(\frac{h}{q_e}l)_{}=\left(4Ch/Lq_e^3\right)\psi `$ tell us that the coefficients $`A`$ and $`B`$ of (6) are not arbitrary. For instance, for the coefficient $`A`$ one obtains the recursive equation
$$A_{l+1}e^{2\pi i\frac{k}{q_e}}+A_{l1}e^{2\pi i\frac{k}{q_e}}=A_l\left(\alpha e^{2\pi i\frac{k}{q_e}}+\alpha ^{}e^{2\pi i\frac{k}{q_e}}\right),$$
(8)
where the coefficient $`\alpha =1i\frac{2C\mathrm{}h}{Lq_e^3k}`$. Defining the new variable $`x_l=e^{i2\pi lk/q_e}A_l`$ the above equation becomes
$$x_{l+1}+x_{l1}=\left\{2\mathrm{cos}\left(\frac{2\pi k}{q_e}\right)+\frac{4C\mathrm{}h}{Lq_e^3k}\mathrm{sin}\left(\frac{2\pi k}{q_e}\right)\right\}x_l.$$
(9)
Like to the Solid State case, equation (9) gives a band structure when the condition
$$\mathrm{cos}\left(\frac{2\pi }{q_e}K\right)=\mathrm{cos}\left(\frac{2\pi k}{q_e}\right)+\frac{2C\mathrm{}h}{Lq_e^3k}\mathrm{sin}\left(\frac{2\pi k}{q_e}\right),$$
(10)
holds for some values of $`K`$. Since the energy $`E`$ and the parameter $`k`$ are related by (7), the above equation gives formally the spectrum $`E(K)`$ of (5). Nevertheless, as said before, from the condition (1) for quantum circuits we must consider $`K=0`$ on (10). Therefore, the spectrum of the $`LC`$ circuit with charge discreteness is obtained from the solutions of
$$1=\mathrm{cos}\left(\frac{2\pi k}{q_e}\right)+\frac{2C\mathrm{}h}{Lq_e^3k}\mathrm{sin}\left(\frac{2\pi k}{q_e}\right).$$
(11)
Observe that $`k=0`$ is not a solution (zero point fluctuations). To solve (11), consider the definition
$$\mathrm{tan}\theta =\frac{2C\mathrm{}h}{Lq_e^3k},$$
(12)
then, the equation (11) becomes
$$\mathrm{cos}\left(\theta \frac{2\pi k}{q_e}\right)=\mathrm{cos}\theta .$$
(13)
There is two solutions:
(a) Normal spectrum ($`E^{(N)}`$): where $`\left(\theta 2\pi k/q_e\right)=+\left(\theta 2\pi n\right)`$, then we have the effective charge $`k=nq_e`$ $`(n0,`$ and integer$`)`$ and the spectrum is
$$E_n^{(N)}=\frac{1}{2C}\left(nq_e\right)^2,$$
(14)
with some similarities to the spectrum of a quantum dot .
(b) Anomalous spectrum ($`E^{(A)}`$): where $`\left(\theta 2\pi k/q_e\right)=\left(\theta 2\pi n\right)`$, then $`\theta =\pi \left(k/q_en\right)`$. Using the definition (12), in this case the equation for the effective charge $`k`$ becomes
$$\frac{k}{q_e}\mathrm{tan}\left(\frac{k}{q_e}\pi \right)=\frac{2C\mathrm{}h}{Lq_e^4},$$
(15)
which defines implicitly the anomalous spectrum $`E_n^{(A)}`$. This is an interesting result, because our technique produces a second branch different from the expected result (14) breaking degeneracy. For instance, in the approximation $`2C\mathrm{}h/Lq_e^41`$ we obtain for the effective charge $`kq_e\left(n1/2\right)`$, with its corresponding energy spectrum given by (7). In the anomalous case we have $`0<k/q_e<0.5`$ (mod.1). In general, from the above results we conclude that between two integer (normal) solutions of (11) (given by $`k=nq_e`$) there is another, (anomalous) related to the branch (15). In this way, the magnetic term in equation (5) removes degeneracy and one obtains the two branches at the spectrum. Moreover, for large values of $`k`$ this two branches approach together as shows figure 1.
Equation (15) can be simplified by assuming that the effective charge $`k=q_e(nx)`$ where the approximation $`nx`$ is considered ( $`0<x<1`$). In this approximation, equation (15) becomes
$$n\mathrm{tan}\pi x=\frac{2C\mathrm{}h}{Lq_e^4}.$$
(16)
As a possible application of the above results consider a two-dimensional electron gas in a mesoscopic device like a quantum point contact \[8-13\]. Since crudely $`Cϵ_od`$ and $`L\mu _od`$ then the right hand side of (15) is $`2C\mathrm{}h/Lq_e^41,399\times 10^3`$. On the other hand, this two-dimensional gas have a density of $`10^{15}m^2`$ , then the number $`n`$ of particles can be estimated as $`10^{15}\times \left(10^6\right)^2=10^3`$(with $`d10^6m,`$ a micrometric distance ). The explicit evaluation of (16) is
$$10^3\mathrm{tan}\pi x=1.399\times 10^3,$$
(17)
with the solution $`x=0.6980.7`$. Therefore, just below the normal state with effective charge $`k=q_en`$, there lies an anomalous state with effective charge $`k=q_e(nx)`$. Note that strictly the right hand of eqs. (15) or (16) are depending on the geometry (size, etc.) of the systems and the assumptions $`Cϵ_od`$ and $`L\mu _od`$ is a crude approximation. Namely, strictly we can expect a size dependence on $`x`$.
Decoherence and dissipation is a hard topic in quantum mechanics and then, also in quantum circuit. In fact, circuits are close related with effects like Ohm law or Joule dissipative effects. The standard way to consider decoherence and dissipation in quantum mechanics is related to the connection of the systems to a thermic bath. Closed equations for the system are found by tracing onto the bath degree. This standard technique was used in meso-particles ( and \[16-20\]). In our case of circuits, a similar procedure could be envisaged in the future.
In resume, we have presented a quantum solvable $`LC`$ circuit with charge discreteness. The spectrum was studied and it contains two branches, the normal (14) and the anomalous (15). It corresponds to the remotion of degeneracy due to the magnetic term in (5). For quantum point contact data, the spectrum was explicitly characterized (17). From a general point of view, our theory represents an easily and efficient approach to consider mesoscopic and nano-devices in the future.
This work was supported by FONDECYT (Grant 1040311). Useful discussions with C. A. Utreras (UACH), A. Perez (PUC) and S. Montecinos (UFRO) concerning charge discreteness and boundary conditions are acknowledged.
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# Spherically symmetric monopoles in noncommutative space
## I Introduction
Field theories in noncommutative space have received renewed interest since their emergence in certain low-energy limits of string and M-theory. In particular, non-perturbative soliton configurations have been the object of numerous investigations in recent years NS -Correa:2004xm . Concerning monopoles, they have been constructed mainly using an extension of Nahm equation in noncommutative space Gross1 , Hamanaka:2001dr . Other approach exploits the connection between soliton solutions in four dimensions and monopole configurations defined in a curved space Correa:2004xm . However, in all of these approaches the obtained configurations are not the natural extension of the well known ’t Hooft-Polyakov monopole solution. The reason is simple, the ’t Hooft ansatz has explicit spherical symmetry while the standard noncommutative three dimensional algebra
$$[x_i,x_j]=i\theta _{ij},\theta _{ij}:\text{constant matrix}$$
(1)
is not invariant under rotations, thus breaking explicitly the rotationally symmetry of any noncommutative field theory (NCFT) defined on it.
In this article we will construct an explicit rotationally invariant noncommutative space by deforming adequately the algebra (1). In particular we will show how to construct a gauge theory in this space by a extending the commutative-space theory written in terms of explicit rotationally invariant operators. The evident advantage is that in this formulation the equations of motion accept a spherically symmetric ansatz, resemblant to the ’t Hooft form. Moreover, we will show that in the small $`\theta `$ limit the solutions tends to the well-known Prasad-Sommerfield solutions.
The article is organized as follows. In section 2 we construct a rotationally invariant noncommutative space. We find that this deformation reduces to a foliation of the three-dimensional space with concentric 2-fuzzy spheres. In section 3 we show how to construct gauge fields in a manner consistent with the rotationally symmetry of the space. In section 4 we construct a Yang-Mills-Higgs theory and derive the equations of motion. Section 5 is devoted to the solution of the equations of motion. Finally in section 6 we summarize the paper and present some discussion.
## II Rotationally invariant noncommutative space
One of the main problems in finding noncommutative monopole solutions is that the simplest ansatz (’t Hooft) has explicit spherical symmetry whereas the standard noncommutative space in three dimensions breaks rotational invariance. Of course, spherical symmetry is not essential for the construction of monopole solutions and in fact several non-spherically symmetric solution has been found explicitly Gross1 , Hamanaka:2001dr , Correa:2004xm . However spherical symmetry greatly simplify the equations by reducing the number of degrees of freedom. So in order to the take advantage of this simplification let us modify the noncommutative structure of the space in order to preserve rotational symmetry.
Consider a three-dimensional noncommutative space with coordinates satisfying the commutator algebra
$$[x_i,x_j]=i\theta \epsilon _{ijk}f(r)x_k$$
(2)
with $`f(r)`$ a function to be determined and $`r^2=x_ix_i`$.
It can be shown that the Jacobi identity imposes the condition $`f(r)r`$ ($`r^2`$ is a Casimir of the algebra), so we have
$$[x_i,x_j]=i\theta r\epsilon _{ijk}x_k$$
(3)
with $`\theta `$ a dimensionless parameter (here, unlike fuzzy-sphere coordinates, the coordinates $`x_1,x_2,x_3`$ are all independent, there is no constraint between them). Then the operators $`x_i/(r\theta )`$ satisfy the $`SU(2)`$ algebra. Being the algebra (3) invariant under space rotations, it is natural to extend it with angular momentum operators $`L_i`$
$`[L_i,L_j]`$ $`=i\epsilon _{ijk}L_k`$
$`[L_i,x_j]`$ $`=i\epsilon _{ijk}x_k`$ (4)
We can find a representation of (3) and (4) by identifying the coordinate operators with $`\theta rL_i`$:
$$x_i=\theta rL_i$$
(5)
with $`L_i`$, $`SU(2)`$ operators. We have that $`1/\theta ^2=L_iL_i`$ and if restrict ourselves to finite-dimensional representations we have $`1/\theta ^2=l(l+1),l\frac{1}{2}`$ (in principle we allow spinor representations). In this representation $`r`$ is a continuos commutative variable. Notice that the algebra (3)-(4), for fixed $`r`$ describes a fuzzy sphere Madore:1991rr -Abe:2002in , so essentially what we are doing is foliating the three-dimensional noncommutative space with concentric fuzzy spheres.
Since algebra (3) is not invariant under space translations, is imposible to define momentum operators satisfying
$`[P_i,x_j]`$ $`=`$ $`i\delta _{ij}`$ (6)
These relations violate the Jacobi identity for three operators $`\{P_i,x_j,x_k\}`$. This is analogous to the fact that for constant non-commutative space
$$[x_i,x_j]=i\theta _{ij}$$
(7)
it is not possible to define angular momentum operators satisfying (4) since the algebra (7) is not rotationally invariant (the Jacobi identity fails for the triplet $`\{L_i,x_j,x_k\}`$).
In order to define a field theory in this noncommutative space we first define transversal and radial field components and write the appropriate lagrangian. In commutative space, given a vector field in Cartesian coordinates $`V_i,i=1,2,3`$, we can define transversal components $`V_i^T`$ and a radial component $`V_r`$ as
$`V_i^T`$ $`=\epsilon _{ijk}x_jV_k`$
$`V_r`$ $`=x_iV_i`$ (8)
The transversal part satisfies the constraint
$$x_iV_i^T=0$$
(9)
Cartesian coordinates can be recovered from the transversal and radial ones through the identity
$$r^2V_i=\epsilon _{ijk}x_jV_k^Tx_iV_r$$
(10)
Since we are working in a non-commutative space with explicit rotational invariance (3) it is natural consider the transversal and radial fields (8) as our primary fields and not the Cartesian components $`V_i`$. This is crucial because in noncommutative space there is no mapping as (10) to define Cartesian coordinates.
So, a transversal field in the noncommutative space is a field that satisfies the constraint $`x_iV_i^T+V_i^Tx_i=0`$, or in virtue of representation (5)
$$L_iV_i^T+V_i^TL_i=0$$
(11)
Is straightforward to check that any vector field of the form
$$[L_i,\mathrm{\Phi }]$$
(12)
is transversal. We will see that a slight modification has to be done in the case of a gauge theory.
## III Gauge fields
As we did for arbitrary vector fields we define transversal and radial gauge fields $`A_i^T`$ and $`A_r`$ in analogy with their commutative counterparts (from now on we will drop the superscript $`T`$ in $`A^T`$). That is $`A_i`$ and $`A_r`$ are fields that in commutative space take the form
$`A_i`$ $`=\epsilon _{ijk}x_j𝒜_k`$
$`A_r`$ $`=x_i𝒜_i`$ (13)
where $`𝒜_i`$ are the Cartesian components of the standard vector potential.
These fields transform, under gauge transformations, as follows:
$`A_i`$ $`g^1A_igg^1[L_i,g]`$
$`A_r`$ $`g^1A_rgg^1[P,g]`$ (14)
with
$$P=ir_r$$
(15)
Again, we are going to promote to noncommutative space the transversal and radial fields $`A_i`$ and $`A_r`$, and not the standard Cartesian gauge field $`𝒜_i`$.
We want to stress again that the map (13) between spherical and Cartesian coordinates is only possible in commutative space. In noncommutative space we are forced to work with spherical coordinates and we cannot recover the Cartesian coordinates. That is, in this space the fundamental fields are the variables $`A_i`$, $`A_r`$ and not $`𝒜_i`$.
But now we have a problem trying to impose the constraint (11). Clearly, the constraint is not invariant under gauge transformations and thus not well defined for gauge fields. In order to define a gauge invariant transversal constraint we introduce the gauge covariant distance $`X_i`$,
$$X_i=x_i\theta rA_i=\theta r\left(L_iA_i\right)$$
(16)
As its names suggests, this quantity transforms under gauge transformations as
$$X_ig^1X_ig$$
(17)
So the correct gauge invariant constraint is given by<sup>1</sup><sup>1</sup>1This equation was first proposed in Karabali:2001te . See also Nair:2001kr , Morariu:2002tx , Abe:2002in .
$$X_iX_i=x_ix_i=r^2$$
(18)
This can be written as
$$\{x_i,A_i\}=\theta rA_iA_i$$
(19)
which in the limit $`\theta 0`$ coincides with (11). It is useful at this point to introduce the transverse covariant derivative operator
$$D_i=L_iA_i$$
(20)
so $`X_i=\theta rD_i`$ and the constraint can be written as
$$D_iD_i=L_iL_i=\kappa $$
(21)
(we have defined $`\kappa =l(l+1)=1/\theta ^2`$), or
$$\{L_i,A_i\}A_iA_i=0$$
(22)
The field strength $`F_{ij}`$ and $`F_{ir}`$ are defined in analogy with the commutative case
$`F_{ij}`$ $`=i\left([L_i,A_j][L_j,A_i][A_i,A_j]i\epsilon _{ijk}A_k\right)`$
$`=i\left([D_i,D_j]i\epsilon _{ijk}D_k\right)`$ (23)
$`F_{ir}`$ $`=i\left([L_i,A_r][P,A_i][A_i,A_r]\right)`$
$`=i[D_i,D_r]`$ (24)
where $`D_r`$ is the radial covariant derivative
$$D_r=PA_r$$
(25)
For convenience we will work in the gauge $`A_r=0`$ so $`D_r=P`$.
As usual, the field $`F`$ is gauge covariant
$`F_{ij}`$ $`g^1F_{ij}g`$
$`F_{ir}`$ $`g^1F_{ir}g`$ (26)
and satisfy the transversality conditions
$$\{D_i,F_{ij}\}=\{D_i,F_{ir}\}=0$$
(27)
## IV Yang-Mills Higgs theory
### The action
To write an action in this geometry, we simply write the action in the commutative-space case in terms of transversal and radial fields using the definition (13) and then promote the fields to noncommutative space, respecting gauge invariance when needed. For a Yang-Mills and Higgs actions we have<sup>2</sup><sup>2</sup>2The integration is defined as $`𝑑x^3=\frac{4\pi }{2l+1}\text{tr}r^2𝑑r`$, where the trace is taken over the angular momentum representation indices.
$`S_{YM}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^3\frac{1}{r^4}\mathrm{tr}\left(F_{ij}F_{ij}+2F_{ir}F_{ir}\right)}`$
$`S_{Higgs}`$ $`={\displaystyle 𝑑x^3\left(\frac{1}{r^2}\mathrm{tr}\left([D_i,\varphi ][D_i,\varphi ]+[D_r,\varphi ][D_r,\varphi ]\right)+V[\varphi ]\right)}`$ (28)
That is, in commutative space $`S_{YM}`$ and $`S_{Higgs}`$ are the usual Yang-Mills and Higgs actions written in term of the transversal and radial fields $`A_i`$, $`A_r`$. Using eqs.(13) we can recover the standard form of the actions in terms of the standard gauge potential $`𝒜_i`$. However eqs.(13) are not valid in noncommutative space and expressions (28) have to be taken as the defining actions in this geometry.
### Equations of motion
From actions (28) we get the Euler-Lagrange equations of motion
$$[D_r,F_{ir}]iF_{ir}[D_j,F_{ji}]\frac{i}{2}\epsilon _{ijk}F_{jk}ir^2[[D_i,\varphi ],\varphi ]=\{\mu ,D_i\}$$
(29)
with $`\mu `$ a Lagrange multiplier enforcing the constraint (21). The r.h.s cancels the longitudinal part of the l.h.s so the resulting equation is transversal.
The remaining equations of motion are
$`[D_i,F_{ir}]+ir^2[[D_r,\varphi ],\varphi ]=0`$ (30)
$`[D_i,[D_i,\varphi ]]+[D_r,[D_r,\varphi ]]=r^2{\displaystyle \frac{\delta V}{\delta \varphi }}`$ (31)
We will concentrate in the case $`V0`$.
To eliminate the Lagrange multiplier we note that given an arbitrary vector $`V_i`$ we can write its transversal part as
$$V_i^T=V_i\frac{1}{2}\{\mu ,D_i\}$$
(32)
for some function $`\mu `$. Now imposing on $`V^T`$ the transversality condition for, $`\{V_i^T,D_i\}=0`$ we find the following equation for $`\mu `$:
$$\kappa \mu +D_i\mu D_i=\{V_i,D_i\}$$
(33)
The transversal part is obtained inserting the solution of eq. (33) in eq. (32).
Before ending this section we have to mention possible BPS equations of motion. In commutative space, in terms of the radial and transversal fields, the BPS equations read
$`D_r\varphi `$ $`={\displaystyle \frac{1}{2r^2}}\epsilon _{ijk}x^iF_{jk}`$ (34)
$`D_i\varphi `$ $`=\pm {\displaystyle \frac{1}{r^2}}\epsilon _{ijk}x^jF_{kr}`$ (35)
However we have been unable to construct a noncommutative version of them. The obvious modifications, replacing the coordinate $`x^i`$ by the covariant coordinate operator $`X^i`$ and the product of $`x^i`$ and $`F_{ab}`$ by the Moyal anticommutator $`\{X^i,F_{ab}\}`$, does not work. For example after this replacement equation (35) reads
$$[D_i,\varphi ]=\pm \frac{1}{2r^2}\epsilon _{ijk}\{X^j,F_{kr}\}$$
(36)
But while the l.h.s is transversal with respect to $`X^i`$, the r.h.s is not. Even projecting the r.h.s over the transverse components does not reproduce the equations of motion.
## V Monopole solutions
### Spherically symmetric ansatz
The most general spherically symmetric ansatz can be written using the operators
$$V_i^{(0)}=L_i,V_i^{(1)}=\sigma ^i,V_i^{(2)}=\{\alpha ,L_i\},V_i^{(3)}=[\alpha ,L_i]$$
(37)
where $`\sigma ^i`$ are the Pauli matrices and
$$\alpha =\underset{i}{\overset{3}{}}\sigma _iL_i$$
(38)
Although (37) is the most general set of rotationally covariant operators, it can be shown that the set remains consistent if we drop $`V_i^{(3)}`$. That is, when we expand the fields in the basis $`\{V_i^{(0)},V_i^{(1)},V_i^{(2)}\}`$ the equations of motions does not have components in the direction $`V_i^{(3)}`$. So, from now on we will work with the basis $`\{V_i^{(0)},V_i^{(1)},V_i^{(2)}\}`$.
Then we expand
$$D_i=\underset{a=0}{\overset{2}{}}v_aV_i^{(a)}$$
(39)
with $`v_a`$ arbitrary functions of the radial coordinate $`r`$, $`v_av_a(r)`$.
The constraint (21) implies the following two equations for the coefficients $`v_0`$, $`v_1`$ and $`v_2`$
$`\kappa v_{0}^{}{}_{}{}^{2}+3v_{1}^{}{}_{}{}^{2}+4\kappa v_1v_2+2\kappa \left(2\kappa 1\right)v_{2}^{}{}_{}{}^{2}=\kappa `$
$`\left(4\kappa 3\right)v_{2}^{}{}_{}{}^{2}2v_0\left(v_1\left(2\kappa 1\right)v_2\right)=0`$ (40)
That is, the field $`D_i`$ depends only on one function.
We have for the field strength
$`F_{ij}=\epsilon _{ija}`$ $`\{(v_0v_{0}^{}{}_{}{}^{2}+(34\kappa )v_{2}^{}{}_{}{}^{2})V_a^{(0)}+`$
$`\left(2v_{1}^{}{}_{}{}^{2}+2\kappa v_0v_24\kappa v_{2}^{}{}_{}{}^{2}+v_1\left(14\kappa v_2\right)\right)V_a^{(1)}+`$
$`(v_23v_0v_2+2v_1v_2+5v_{2}^{}{}_{}{}^{2})V_a^{(2)}\}`$ (41)
and
$`F_{kr}=r\left(v_0^{}V_k^{(0)}+v_1^{}V_k^{(1)}+v_2^{}V_k^{(2)}\right)`$ (42)
To write the equation of motion (29) we have first to solve equation (33) to find the Lagrange multiplier $`\mu `$ that projects the solution onto the tangential space. Spherical symmetry imposes that $`\mu `$ has the form
$$\mu =\mu _0+\alpha \mu _1$$
(43)
so equation (33) leads to algebraic equations for the coefficients $`\mu _0`$ and $`\mu _1`$. However we will see later that we can solve explicitly the constraint and thus work with the physical, unconstrained degrees of freedom, making unnecessary the Lagrange multiplier.
The Higgs field can be expanded as
$$\varphi =\varphi _0(r)+\varphi _1(r)\alpha $$
(44)
and the covariant derivatives takes the form
$`[D_i,\varphi ]`$ $`=\varphi _1(v_02v_1v_2)V_i^{(3)}`$
$`[D_r,\varphi ]`$ $`=ir\left(\varphi _0^{}(r)+\varphi _1^{}(r)\alpha \right)`$ (45)
It can be checked that the equation of motion (30) is trivially satisfied.
Finally, the last equations, (31) take the form
$`r{\displaystyle \frac{d^2}{dr^2}}\left(r\varphi _0\right)`$ $`=0\varphi _0={\displaystyle \frac{c_1}{r}}+c_0`$
$`r{\displaystyle \frac{d^2}{dr^2}}\left(r\varphi _1\right)`$ $`=2\varphi _1\left(v_02v_1v_2\right)^2`$ (46)
Notice that $`\varphi _0`$ is decoupled from the other fields and in fact it is an irrelevant constant.
In these variables the Hamiltonian takes the form
$`H=`$ $`8\pi {\displaystyle }dr\{r^2(\varphi _{0}^{}{}_{}{}^{2}+\varphi _{1}^{}{}_{}{}^{2}\kappa )+2\varphi _{1}^{}{}_{}{}^{2}\kappa (v_0+2v_1+v_2)^2)+`$
$`{\displaystyle \frac{1}{r^2}}(3(12v_1)^2v_{1}^{}{}_{}{}^{2}+\kappa (2v_{0}^{}{}_{}{}^{3}+v_{0}^{}{}_{}{}^{2}(1+8(4\kappa 3)v_{2}^{}{}_{}{}^{2})`$
$`2(4\kappa 3)v_0v_{2}^{}{}_{}{}^{2}(3+4v_1+10v_2)+v_2(32v_{1}^{}{}_{}{}^{3}+8v_{1}^{}{}_{}{}^{2}(4\kappa v_23)+`$
$`v_2\left(2+4\kappa +4\left(6\kappa 5\right)v_2+\left(41+4\kappa \left(11+4\kappa \right)\right)v_{2}^{}{}_{}{}^{2}\right)+`$
$`4v_1(1+v_2(3+2(8\kappa 5)v_2))))+v_{0}^{}{}_{}{}^{4})+(3v_{1}^{}{}_{}{}^{2}+\kappa (v_{0}^{}{}_{}{}^{2}+`$
$`2v_2^{}(2v_1^{}+(1+2\kappa )v_2^{}))\}`$ (47)
### Small $`\theta `$ expansion
Let us study first the small $`\theta `$ expansion of the monopole equations. First we note that the operator $`V_i^{(2)}`$ is already of order $`1/\theta ^2`$:
$$V_i^{(2)}=\{\alpha ,L_i\}=\frac{1}{\theta ^2}\{\widehat{X}\stackrel{}{\sigma },\widehat{X}^i\}$$
(48)
so the coefficient $`v_2`$ is of order $`\theta ^2`$. (Since the covariant derivative operator starts with $`L_i`$ at zero-order the coefficient $`v_0`$ is order zero).
To compare with the usual ’t Hofft-̃Polyakov-̃Julia-̃Zee Prasad Sommerfield solutions we write
$`v_01=\theta ^2v_0^{(2)}+\theta ^4v_0^{(4)}+\mathrm{}`$
$`v_1={\displaystyle \frac{k1}{2}}\theta ^2{\displaystyle \frac{k_1}{2}}+\mathrm{}`$
$`v_2=v_2^{(0)}\theta ^2+\theta ^4v_2^{(2)}+\mathrm{}`$
$`\varphi _1={\displaystyle \frac{\theta }{2r}}(h+\theta ^2h_1+\mathrm{}`$ (49)
The constraints (40) can be solved perturbatively in $`\theta `$ and we can write the coefficients of $`v_0`$ ($`v_0^{(2)},v_0^{(4)},\mathrm{}`$) and $`v_2`$ ($`v_2^{(0)},v_2^{(2)},\mathrm{}`$) as functions of the coefficients of $`v_1`$ ($`k,k_1,\mathrm{}`$). We have
$`v_01={\displaystyle \frac{\theta ^2}{4}}\left(k1\right)^2{\displaystyle \frac{\theta ^4}{32}}\left(k1\right)\left(5+k7k^2+k^3+16k_1\right)+\mathrm{}`$
$`v_1={\displaystyle \frac{1k}{2}}+{\displaystyle \frac{\theta ^2}{8}}\left(k^214k_1\right)\mathrm{}`$
$`v_2=\theta ^2{\displaystyle \frac{k1}{4}}+\theta ^4{\displaystyle \frac{1}{4}}k_1+\mathrm{}`$
At leading order we recover the standard monopole equations
$`r^2k^{\prime \prime }(r)=k(r)\left(k(r)^21+h^2(r)\right)`$
$`r^2h^{\prime \prime }(r)=2k(r)h(r)`$ (51)
with the well-known solutions Prasad:1975kr
$`k(r)`$ $`={\displaystyle \frac{r}{\mathrm{sinh}(r)}}`$
$`h(r)`$ $`=r\mathrm{coth}(r)1`$ (52)
The next order equations read
$`r^2`$ $`k_{1}^{}{}_{}{}^{\prime \prime }(r)+(1h(r)^23k(r)^2)k_1(r)={\displaystyle \frac{1}{4}}(1+8h(r)h_1(r)k(r)+`$
$`3k(r)^2+7k(r)^34k(r)^42k(r)^5+h(r)^2(1+k(r)`$
$`2k(r)^3)+4r^2k^{}(r)^2+k(r)(32r^2k^{}(r)^2+2r^2k^{\prime \prime }(r)))`$ (53)
$`r^2h_{1}^{}{}_{}{}^{\prime \prime }(r)2h_1(r)k(r)^2=h(r)k(r)\left(1+k(r)2k(r)^2+4k_1(r)\right)`$ (54)
and can be solved numerically.
### Solving the constraint
Instead of working with the “linear” variables $`v_0,v_1,v_2`$ and the constraints (40) we can try to reparametrize the fields and solve the constraint explicitly. Then the resulting fields are the physical degrees of freedom and the constraint is automatically incorporated in the equations of motion.
In fact, we can see that the replacement
$`v_0`$ $`(z_0z_1)/\sqrt{2}`$
$`v_1`$ $`(z_0+z_1)/\sqrt{2}{\displaystyle \frac{2\kappa 1}{\sqrt{4\kappa 3}}}z_2`$
$`v_2`$ $`z_2/\sqrt{4\kappa 3}`$ (55)
diagonalizes the second equation (40)
$$z_0^2z_1^2z_2^2=0$$
(56)
which is straightforwardly solved in term of two functions $`\rho `$ and $`u`$ (both are functions of $`r`$)
$`z_0`$ $`=\rho `$
$`z_1`$ $`=\rho c(u)`$
$`z_2`$ $`=\rho u`$ (57)
with
$$c(u)=\sqrt{1u^2}\text{and}1u1$$
(58)
(we chose the branch solution that matches the standard $`\theta 0`$ limit). Replacing this solution into the first of equations (40) we get a quadratic equation for $`\rho `$ that can be easily solved. Finally we have a parametrization that solves the constraint
$`v_0`$ $`={\displaystyle \frac{\sqrt{\kappa }}{\sqrt{2d(u)}}}\left(1+c(u)\right)`$
$`v_1`$ $`={\displaystyle \frac{\sqrt{\kappa }}{2\sqrt{d(u)}\sqrt{4\kappa 3}}}\left(2u(2\kappa 1)+\sqrt{8\kappa 6}\left(1c(u)\right)\right)`$
$`v_2`$ $`={\displaystyle \frac{\sqrt{\kappa }}{\sqrt{d(u)}\sqrt{4\kappa 3}}}u`$ (59)
where
$`d(u)=3+\kappa +{\displaystyle \frac{1}{2}}u^2\left(3\kappa 5\right)+\left(\kappa 3\right)c(u)+\sqrt{8\kappa 6}\left(1+c(u)\right)u`$ (60)
That is, we have parametrized the gauge fields in terms of only one function $`u`$, which together with the Higgs field $`\varphi _1`$ are the only nontrivial degrees of freedom. The next step is to write the equations of motion in terms of them. Actually, though we have reduced significantly the number of degrees of freedom, the equations of motion are very complicated in terms of these fields. We show the complete expression of the equations of motions in the appendix.
In this variables the small $`\theta `$ limit can be recovered through the identification
$$u(r)=\frac{\theta }{\sqrt{2}}\left(1k(r)\right)+O(\theta ^3)$$
(61)
### Boundary conditions
In order to get non-singular, finite energy solutions we have to impose appropriate boundary conditions. At the origin we have the usual conditions
$`u(0)=0`$
$`\varphi _1(0)=0`$ (62)
At $`r\mathrm{}`$ the situation is different from the commutative case. Notice that in the presence of a potential,
$$V=\lambda \left(\varphi ^2\eta ^2\right)^2$$
(63)
the Higgs field tends asymptotically to a minimum of the potential. That is, asymptotically, $`\varphi _0`$ and $`\varphi _1`$ are minima of
$$V=\lambda \left(\frac{1}{16}\frac{1}{2}\varphi _0^2+\varphi _0^4\frac{\kappa }{2}\varphi _1^2+6\kappa \varphi _0^2\varphi _1^24\kappa \varphi _0\varphi _1^3+\kappa \left(\kappa +1\right)\varphi _1^4\right)$$
(64)
$$V=12\varphi _{0}^{}{}_{}{}^{2}+\varphi _{0}^{}{}_{}{}^{4}2\varphi _{1}^{}{}_{}{}^{2}\kappa +6\varphi _{0}^{}{}_{}{}^{2}\varphi _{1}^{}{}_{}{}^{2}\kappa 4\varphi _0\varphi _{1}^{}{}_{}{}^{3}\kappa +\varphi _{1}^{}{}_{}{}^{4}\kappa +\varphi _{1}^{}{}_{}{}^{4}\kappa ^2$$
(65)
(we have rescaled the fields so $`\eta =1/2`$, consistent with the small $`\theta `$ expansion solution). Besides the trivial solution $`\varphi _0=1,\varphi _1=0`$ we have the solutions
$`\varphi _0={\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{1}{\sqrt{1+4\kappa }}}\right),`$ $`\varphi _1={\displaystyle \frac{1}{2\sqrt{1+4\kappa }}}`$ (66)
$`\varphi _0={\displaystyle \frac{1}{2\sqrt{1+4\kappa }}},`$ $`\varphi _1={\displaystyle \frac{1}{\sqrt{1+4\kappa }}}`$ (67)
(these correspond to absolute minima of the potential; there are other local minima but those will give infinite energy when integrated over the whole space).
The first of these equations gives a nontrivial $`U(1)`$ contribution in the $`\theta 0`$ limit so we discard it. The second one gives the correct small $`\theta `$ behavior so we take it as the asymptotic boundary condition:
$`\underset{r\mathrm{}}{lim}\varphi _1(r)={\displaystyle \frac{1}{\sqrt{1+4\kappa }}}`$ (68)
Of course, this is valid in the presence of a potential. For vanishing coupling constant, as it happens in commutative space, we can rescale the Higgs fields arbitrarily by rescaling appropriately the radial variable $`r`$.
For the gauge field we impose, as usual, that at infinity the Higgs kinetic term vanishes. This gives the behavior
$$\underset{r\mathrm{}}{lim}u(r)=\frac{4\sqrt{2}}{\sqrt{4\kappa 3}+3\sqrt{4\kappa +1}}$$
(69)
### Numerical solutions
We solved numerically the equations of motion for different values of $`\kappa =1/\theta ^2=l(l+1)`$. We found solutions for essentially any value of $`\kappa `$ allowed ($`\kappa >3/4`$). As expected, for large $`\kappa `$ (small $`\theta `$ ) the solution tends to the Prasad-Sommerfield configurations. Indeed, even for $`l=1`$, the profile of the solutions are very similar to the P-S solutions. In order to see the departure of the P-S solutions we considered continuos values of $`\kappa `$ (which correspond to infinite dimensional representation of the non-commutative algebra). It is remarkable that the Higgs field solution is not very sensitive to $`\kappa `$, even for extreme values ($`\kappa 3/4`$). On the other hand, the gauge field in very sensitive to $`\kappa `$. We show in figures 1 and 2 the solutions for the fields $`u`$ and $`\varphi _1`$ respectively, for various values of $`\theta `$.
We also studied the energy of the monopole solutions as a function of $`\theta `$. For small vales of $`\theta `$, the energy, in units of $`e^2/4\pi `$, tends to $`1`$ as expected (BPS bound). As $`\theta `$ increases, the energy also increases and diverges as $`\theta ^2`$ approaches to $`4/3`$. A plot of the energy as a function of $`\theta `$ is shown in figure 3. This behavior is another hint that for $`\theta `$ different from zero, the solutions obtained are not self-dual, since in that case we expect the energy of the configuration to be equal to some topological number (independent of $`\theta `$). This situation can be contrasted with the case of self-dual vortex solutions in NC space. In the later, while the profile of the solutions are dependent of the noncommutative parameter, the energy is $`\theta `$-independent and in particular equal to $`1`$ (in appropriate units), the Bogomolny energy bound Lozano:2000qf .
## VI Summary and conclusions
Previous analysis of monopole configuration in noncommutative space were done using the standard non-commutative relations
$$[x_i,x_j]=i\theta _{ij}\theta _{ij}:\text{constant}$$
(70)
Although this algebra is invariant under space translations, as is immediate from the definition, commutation relations (70) are not invariant under space rotations. In particular we cannot benefit from the simplifications, in structure and in number of degrees of freedom, that a spherically symmetric ansatz produce.
In contrast to relations (70), we can construct a different noncommutative algebra which is manifestly rotationally invariant
$$[x_i,x_j]=i\theta r\epsilon _{ijk}x_k$$
(71)
but at the expense of loosing translational invariance. In fact, the algebra (71) is incompatible with a momentum operator $`P_i`$ generating infinitesimal translations. However, this is not an impediment to construct a field theory in this geometry. A representation of this algebra can be constructed by identifying $`x_i=\theta rL_i`$. In this representation the value $`\theta `$ labels the representation through the relation $`1/\theta ^2=\stackrel{}{L}^2`$. So, although $`\theta `$ can take any positive value, for the special case $`1/\theta ^2=l(l+1),l`$, we have finite-dimensional representations (notice however that the radial variable $`r`$ takes continuous values).
In commutative space, a Poincaré invariant Lagrangian can be written in terms of momentum operators $`P_i`$,where translational invariance is manifest, or in terms of angular momentum operators $`L_i`$ (together with a radial scaling operator $`P`$), where rotational invariance is obvious. To construct a NCFT with the algebra (70) one choose the former and promote the variables (with some prescribed order) to noncommutative operators. Analogously, to construct a NCFT with the algebra (71), we can choose the later and again, promote the variables to noncommutative operators.
In particular we constructed a Yang-Mills Higgs Hamiltonian in this space, and also derived the equations of motions. A puzzling aspect is that we were unable to derive first order (BPS) equations of motion. Though we do not have a rigorous proof of this statement, there are several hints that suggest this property. Since the theory is manifestly invariant under rotations we tried a spherically symmetric ansatz, which is nothing but a noncommutative extension of the ’t Hooft monopole ansatz. Then, as it happens in the commutative case, the number of degrees of freedom is reduced to just two, one for the gauge field and other for the Higgs field. The final equations of motion are very complicated in form but not difficult to solve numerically. Moreover, we showed that in the limit $`\theta 0`$ the equations of motion (and in fact the whole Hamiltonian), reduces to the standard commutative Yang-Mills-Higgs theory, allowing then a perturbative solution in the noncommutative parameter. Another characteristic of this theory is that it blows up at $`\theta ^2=4/3`$, which incidentally is the maximum value of $`\theta `$ for which there is a finite-dimensional representation of the algebra.
We solved numerically the Euler equations of motion for different values of $`\theta `$. As expected for small values of $`\theta `$ the solutions is indistinguishable from the exact Prasad-Sommerfield solution. As we increase the value of $`\theta `$ the profile of the solution depart from the P-S solutions, and also the energy increases. In particular the energy diverges as $`\theta `$ approaches to its maximum value $`\theta ^2=4/3`$.
###### Acknowledgements.
We are grateful with Diego Correa and Gustavo Lozano for helpful comments. We would like to thank Fidel Schaposnik for valuable discussions and critical readings of the manuscript.
*
## Appendix
In this appendix we present the Euler-Lagrange equations of motion in terms of the unconstrained variable $`u`$ and the Higgs field $`\varphi _1`$. The equations read:
$`u^{\prime \prime }(r)+((2r^2\kappa c(u)^4d(u)^2\varphi _1(r)^2(\sqrt{8\kappa 6}3\sqrt{8\kappa 6}c(u)+`$
$`2(4\kappa 3)u\left)\right(6(4\kappa 3)c(u)^3+c(u)^2(2432\kappa +(2117\kappa )\times `$
$`\sqrt{8\kappa 6}u)+u((3\kappa )\sqrt{8\kappa 6}6\sqrt{8\kappa 6}d(u)`$
$`2(1219\kappa +4\kappa ^2)u+2\sqrt{2}(4\kappa 3)^{\frac{3}{2}}u^2)+2c(u)(4\kappa 3+`$
$`(68\kappa )d(u)+(7\kappa 10)\sqrt{8\kappa 6}u+(2441\kappa +12\kappa ^2)u^2)))\times `$
$`(4\kappa 3)^1+c(u)^4d(u)(\sqrt{2}\kappa (\left(\sqrt{2}\kappa (1+c(u))^3\right)+`$
$`3\sqrt{\kappa }\left(1+c(u)\right)^2\sqrt{d(u)}\sqrt{2}\left(1+c(u)\right)\left(d(u)+8\kappa u^2\right)+`$
$`(2\sqrt{\kappa }u^2(2\sqrt{2}\sqrt{\kappa }\sqrt{4\kappa 3}c(u)+3\sqrt{4\kappa 3}\sqrt{d(u)}+`$
$`2\sqrt{\kappa }(\sqrt{8\kappa 6}+(7+4\kappa )u)))(4\kappa 3)^{1/2})\times (2d(u)u+`$
$`(1+c(u))(\sqrt{8\kappa 6}c(u)^2c(u)(\sqrt{8\kappa 6}+(3\kappa 5)u)+`$
$`u(\kappa 3\sqrt{8\kappa 6}u)))+(2\kappa u(2\sqrt{4\kappa 3}c(u)^3+`$
$`c(u)^2(4\sqrt{4\kappa 3}+\sqrt{2}(3\kappa 5)u)+u(\sqrt{2}(\kappa 3)+2\sqrt{2}d(u)`$
$`2\sqrt{4\kappa 3}u)+2c(u)(\sqrt{4\kappa 3}2\sqrt{2}(\kappa 2)u+\sqrt{4\kappa 3}u^2))\times `$
$`(24\kappa (4\kappa 3)c(u)^2+(8\kappa 6)d(u)3\sqrt{\kappa }\sqrt{d(u)}(\sqrt{2}(4\kappa 3)+`$
$`2(74\kappa )\sqrt{4\kappa 3}u)+\kappa u((2720\kappa )\sqrt{8\kappa 6}+2(7164\kappa +`$
$`16\kappa ^2)u)+c(u)(15\sqrt{2}(34\kappa )\sqrt{\kappa d(u)}+3\kappa (8(4\kappa 3)+`$
$`(2912\kappa )\sqrt{8\kappa 6}u))))\times ((4\kappa 3)^{\frac{3}{2}})^1+((6(4\kappa 3)c(u)^3+`$
$`\sqrt{4\kappa 3}c(u)^2\left(12\sqrt{4\kappa 3}+\sqrt{2}\left(21+17\kappa \right)u\right)+`$
$`u(3(\kappa 3)\sqrt{8\kappa 6}+6\sqrt{8\kappa 6}d(u)+2(1827\kappa +4\kappa ^2)u+`$
$`2(34\kappa )\sqrt{8\kappa 6}u^2)+2c(u)(912\kappa +(8\kappa 6)d(u)`$
$`5(2\kappa 3)\sqrt{8\kappa 6}u+(24+41\kappa 12\kappa ^2)u^2))\times `$
$`(4\sqrt{2}\kappa (4\kappa 3)^{\frac{3}{2}}c(u)^3(4\kappa 3)d(u)(\sqrt{8\kappa 6}+2(2\kappa 1)u)+`$
$`6\sqrt{\kappa }\sqrt{4\kappa 3}\sqrt{d(u)}\times \left(4\kappa 3+\left(2+3\kappa \right)\sqrt{8\kappa 6}u+2u^2\right)`$
$`2\kappa (2\sqrt{2}(4\kappa 3)^{\frac{3}{2}}+4(34\kappa )^2u+3\sqrt{8\kappa 6}(47\kappa +4\kappa ^2)u^2`$
$`2(443\kappa +48\kappa ^216\kappa ^3)u^3)+6c(u)^2(\left(\sqrt{\kappa }(4\kappa 3)^{\frac{3}{2}}\sqrt{d(u)}\right)`$
$`2\kappa (4\kappa 3)\times (\sqrt{8\kappa 6}+2(1+2\kappa )u))+\sqrt{4\kappa 3}c(u)\times `$
$`(\sqrt{2}(4\kappa 3)d(u)6\sqrt{\kappa }\sqrt{d(u)}(8\kappa 6+(\kappa 2)\sqrt{8\kappa 6}u)+`$
$`6\kappa (\sqrt{8}(4\kappa 3)+(\kappa 1)\sqrt{4\kappa 3}u+\sqrt{2}(421\kappa +12\kappa ^2)u^2))))\times `$
$`\left((34\kappa )\right)^2)r^2\kappa (16\sqrt{8\kappa 6}(12+\kappa ^2+(12+\kappa ^2)c(u))+`$
$`16\left(11754\kappa 8\kappa ^2+2\kappa ^3+\left(117+58\kappa 8\kappa ^2+2\kappa ^3\right)c(u)\right)u+`$
$`4\sqrt{8\kappa 6}\left(33958\kappa +17\kappa ^2+\left(315+38\kappa +15\kappa ^2\right)c(u)\right)u^2`$
$`4\left(1101884\kappa 55\kappa ^2+30\kappa ^3\left(859836\kappa +103\kappa ^226\kappa ^3\right)c(u)\right)u^3`$
$`\sqrt{8\kappa 6}\left(1519596\kappa +93\kappa ^2+\left(977+432\kappa +65\kappa ^2\right)c(u)\right)u^4+`$
$`\left(22652612\kappa +213\kappa ^2+34\kappa ^3\left(385600\kappa +141\kappa ^2+14\kappa ^3\right)c(u)\right)u^5+`$
$`\sqrt{8\kappa 6}\left(41+84\kappa 43\kappa ^2+\left(109+48\kappa 69\kappa ^2\right)c(u)\right)u^6+`$
$`\left(26760\kappa 305\kappa ^2+54\kappa ^32\left(485872\kappa +361\kappa ^218\kappa ^3\right)c(u)\right)u^7+`$
$`4\sqrt{8\kappa 6}(99112\kappa +21\kappa ^2)u^8)u^{}(r)^2\left)\right(2r^2\kappa c(u)^3d(u)(8(6\kappa +`$
$`2\kappa ^2+(6\kappa +2\kappa ^2)c(u))+8\sqrt{8\kappa 6}(2+\kappa +(\kappa 2)c(u))u+`$
$`4\left(4+5\kappa +2\left(\kappa 2\right)\kappa c(u)\right)u^2+2\sqrt{8\kappa 6}\left(3\kappa +5\left(\kappa 2\right)c(u)\right)u^3+`$
$`(49+57\kappa 6\kappa ^2)u^4))^1=0`$ (72)
$`r{\displaystyle \frac{d^2}{dr^2}}\left(r\varphi _1\right)\left(\kappa \left(\left(3c(u)1\right)\sqrt{8\kappa 6}u\right)^2\varphi _1\right)d(u)^1=0`$ (73)
where the functions $`c(u)`$ and $`d(u)`$ are defined in equations (58) and (60) respectively.
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# 1 Introduction
## 1 Introduction
### 1.1 A brief review of nanomagnetism
In the last decade, there has been a considerable interest in atomic engineering, which makes it possible to create small-scale materials with the use of various methods (for reviews, see Refs. ). Small-scale magnetic systems ranging from grains (micros), nanosystems, molecular magnets and atomic clusters, display a variety of interesting properties. Magnetic nanosystems consist of small clusters of magnetic ions embedded within nonmagnetic ligands or on nonmagnetic substrates. Nanomagnetism shows interesting properties different from bulk magnetism. Nanoclusters consisting of transition metals such as $`\mathrm{Fe}_N`$ ($`N`$=15-650) , $`\mathrm{Co}_N`$ ($`N`$=20-200) , and $`\mathrm{Ni}_N`$ ($`N`$=5-740) have been synthesized by laser vaporization and their magnetic properties have been measured, where $`N`$ denotes the number of atoms per cluster. Magnitudes of magnetic moments per atom are increased with reducing $`N`$ . It is shown that magnetic moments in Co monatomic chains constructed on Pt substrates are larger than those in monolayer Co and bulk Co . Recently Au nanoparticles with average diameter of 1.9 nm (including 212 atoms), which are protected by polyallyl amine hydrochloride (PAAHC), are reported to show ferromagnetism while bulk Au is diamagnetic . This is similar to the case of gas-evaporated Pd fine particles with the average diameter of 11.5 nm which show the ferromagnetism whereas bulk Pd is paramagnetic . The magnetic property of four-Ni molecular magnets with the tetrahedral structure (abbreviated as Ni4) in metallo-organic substance $`[\mathrm{Mo}_{12}\mathrm{O}_{30}(\mu _2\mathrm{OH})_{10}\mathrm{H}_2\{\mathrm{Ni}(\mathrm{H}_2\mathrm{O}_3)\}_4]14\mathrm{H}_2\mathrm{O}`$ has been studied . Their temperature-dependent susceptibility and magnetization process have been analyzed by using the Heisenberg model with the antiferromagnetic exchange couplings between Ni atoms . Similar analysis has been made for magnetic molecules of Fe$`N`$ ($`N=6,\mathrm{\hspace{0.25em}8},\mathrm{\hspace{0.25em}10}`$ and 12) , and V6 . Extensive studies have been made for single molecule magnets of Mn12 in $`[\mathrm{Mn}_{12}\mathrm{O}_{12}(\mathrm{CH}_3\mathrm{COO})_{16}(\mathrm{H}_2\mathrm{O})_4]`$ and Fe8 in $`[\mathrm{Fe}_8(tanc)_6\mathrm{O}_2(\mathrm{OH})_{12}]\mathrm{Br}_99\mathrm{H}_2\mathrm{O}`$ . Both Mn12 and Fe8 behave as large single spins with $`S=10`$, and show quantum tunneling of magnetization and the square-root relaxation, which are current topics in namomagnetism. Much attention has been recently paid to single molecule magnets which are either dimers or behave effectively as dimers, due to their potential use as magnetic storage and quantum computing. The iron $`S=5/2`$ dimer (Fe2) in $`[\mathrm{Fe}(\mathrm{OMe})(dbm)_2]_2`$ has a nonmagnetic, singlet ground state and its thermodynamical property has been analyzed with the use of the Heisenberg model -. Similar analysis has been made for transition-metal dimers of V2 , Cr2 , Co2 , Ni2 and Cu2 .
### 1.2 Non-extensive statistics
As the size of systems becomes smaller, effects of fluctuations and contributions from surface play more important roles. There are currently three approaches to discussing nanothermodynamics for small-size systems: (1) a modification of the Boltzman-Gibbs statistics (BGS) adding subdivision energy , (2) non-equilibrium thermodynamics including work fluctuations , and (3) the non-extensive statistics (NES) generalizing the BGS as to take account of the non-extensive feature of such systems -. A comparison between these approaches have been made in Refs. .
Before discussing the NES, let’s recall the basic feature of the BGS for a system with internal energy $`E`$ and entropy $`S`$, which is immersed in a large reservoir with energy $`E_0`$ and entropy $`S_0`$. The temperature of the system $`T`$ is the same as that of the reservoir $`T_0`$ where $`T=\delta E/\delta S`$ and $`T_0=\delta E_0/\delta S_0`$. If we consider the number of possible microscopic states of $`\mathrm{\Omega }(E_0)`$ in the reservoir, its entropy is given by $`S_0=k_B\mathrm{ln}\mathrm{\Omega }(E_0)`$ where $`k_B`$ denotes the Boltzman constant. The probability of finding the system with the energy $`E`$ is given by $`p(E)=\mathrm{\Omega }(E_0E)/\mathrm{\Omega }(E)\mathrm{exp}(E/k_BT)`$ with $`EE_0`$. When the physical quantity $`Q`$ of a system containing $`N`$ particles is expressed by $`QN^\gamma `$, they are classified into two groups in the BGS: intensive ($`\gamma =0`$) and extensive ones ($`\gamma =1`$). The temperature and energy are typical intensive and extensive quantities, respectively. This is not the case in the NES, as will be shown below.
When a small-scale nanosystem is immersed in a reservoir, the temperature of the nanosystem is expected to fluctuate around the temperature of the reservoir $`T_0`$ because of the smallness of the nanosystem and its quasi-thermodynamical equilibrium states. Then the BGS distribution mentioned above has to be averaged over the fluctuating temperature. This idea has been expressed by
$`p(E)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\beta e^{\beta E}f^B(\beta )`$ (1)
$`=`$ $`[1(1q)\beta _0E]^{\frac{1}{1q}}\mathrm{exp}_q(\beta _0E),`$
with
$`q`$ $`=`$ $`1+{\displaystyle \frac{2}{N}},`$ (2)
$`f^B(\beta )`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }\left(\frac{N}{2}\right)}}\left({\displaystyle \frac{N}{2\beta _0}}\right)^{\frac{N}{2}}\beta ^{\frac{N}{2}1}\mathrm{exp}\left({\displaystyle \frac{N\beta }{2\beta _0}}\right),`$ (3)
$`\beta _0`$ $`=`$ $`{\displaystyle \frac{1}{k_BT_0}}={\displaystyle _0^{\mathrm{}}}𝑑\beta f(\beta )\beta E(\beta ),`$ (4)
$`{\displaystyle \frac{2}{N}}`$ $`=`$ $`{\displaystyle \frac{E(\beta ^2)E(\beta )^2}{E(\beta )^2}},`$ (5)
where $`\mathrm{exp}_q(x)`$ denotes the $`q`$-exponential function defined by
$`\mathrm{exp}_q(x)`$ $`=`$ $`[1+(1q)x]^{\frac{1}{1q}},\text{for }1+(1q)x>0`$ (6)
$`=`$ $`0.\text{otherwise}`$
In Eqs. (1)-(6), $`q`$ expresses the entropic index, $`f^B(\beta )`$ the $`\mathrm{\Gamma }`$ (or $`\chi ^2`$) distribution function of the order $`N`$, $`E(Q)`$ the expectation value of $`Q`$ averaged over $`f(\beta )`$, $`\beta _0`$ the average of the fluctuating $`\beta `$ and $`2/N`$ its variances. The $`\mathrm{\Gamma }`$ distribution of the order $`N`$ is emerging from the sum of squares of $`N`$ Gaussian random variables. In deriving Eqs. (1)-(5), we have assumed that $`N`$ particles are confined within a small volume of $`L^3`$ ($`L<\xi `$) where the variable $`\beta `$ uniformly fluctuates, $`\xi `$ standing for the coherence length .
The important consequence of the NES is that energy and entropy are not proportional to $`N`$ in nanosystems. The non-extensivity of the entropy was first demonstrated by Tsallis, who proposed the generalized entropy given by
$$S_q=k_B\left(\frac{_ip_i^q1}{1q}\right)=k_B\underset{i}{}p_i^q\mathrm{ln}_q(p_i),$$
(7)
where $`p_i`$ \[$`=p(ϵ_i)`$\] denotes the probability distribution for the energy $`ϵ_i`$ in the system and $`\mathrm{ln}_q(x)`$ \[$`=(x^{1q}1)/(1q)`$\] the $`q`$-logarithmic function, the inverse of the $`q`$-exponential function defined by Eq. (6). It is noted that in the limit of $`q=1`$, Eq. (7) reduces to the entropy of BGS, $`S_{BG}`$, given by
$$S_1=S_{BG}=k_B\underset{i}{}p_i\mathrm{ln}p_i.$$
(8)
The non-extensivity in the Tsallis entropy is satisfied as follows. Suppose that the total system containing $`2N`$ particles is divided into two independent subsystems, each of which contains $`N`$ particles, with the probability distributions, $`p_i^{(1)}`$ and $`p_i^{(2)}`$. The total system is described by the factorized probability distribution $`p_{ij}=p_i^{(1)}p_j^{(2)}`$. The entropy for the total system $`S(2N)`$ is given by
$$S(2N)=S(N)+S(N)+O\left(\frac{1}{N}\right),$$
(9)
where $`S(N)`$ stands for the entropy of the $`N`$-particle subsystem, the index $`q`$ given by Eq. (2) being employed. Similarly the energy of the total system is expressed by
$$E(2N)=E(N)+E(N)+O\left(\frac{1}{N}\right),$$
(10)
The difference of $`E(2N)2E(N)`$ is attributed to the surface contribution. This implies that the index $`\gamma `$ in $`QN^\gamma `$ is neither 0 nor 1 for $`Q`$ = $`S`$ and $`E`$ in nanosystems within the NES.
The functional form of the probability distribution $`p(E)`$ expressed by Eq. (1) was originally derived by the maximum-entropy method . The probability of $`p_i[=p(ϵ_i)]`$ for the eigenvalue $`ϵ_i`$ in the NES is determined by imposing the variational condition to the entropy given by Eq. (7) with the two constraints :
$`{\displaystyle \underset{i}{}}p_i`$ $`=`$ $`1,`$ (11)
$`{\displaystyle \frac{_ip_i^qϵ_i}{_ip_i^q}}`$ $`=`$ $`E_q.`$ (12)
The maximum-entropy method leads to the probability distribution $`p_i`$ given by
$`p_i`$ $``$ $`\mathrm{exp}_q\left[\beta _0\left(ϵ_iE_q\right)\right],`$ (13)
with
$`\beta _0`$ $`=`$ $`{\displaystyle \frac{\beta }{c_q}},`$ (14)
$`c_q`$ $`=`$ $`{\displaystyle \underset{i}{}}p_i^q,`$ (15)
where $`\beta `$ denotes the Lagrange multiplier relevant to the constraint given by Eq. (12). It has been shown that the physical temperature $`T`$ of the nanosystem is given by
$`T`$ $`=`$ $`{\displaystyle \frac{c_q}{k_B\beta }},\text{(AMP)}`$ (16)
In the limit of $`q=1`$, we get $`\mathrm{exp}_q[x]=e^x`$, $`c_q=1`$ and $`p_i`$ given by Eqs. (13)-(16) reduces to the results obtained in the BGS, related discussions being given in Sec. 4.
In previous papers , I have applied the NES to the Hubbard model, which is one of the most important models in solid-state physics (for a recent review, see Ref. ). The Hubbard model consists of the tight-binding term expressing electron hoppings and the short-range interaction between two electrons with opposite spins. The Hubbard model provides us with good qualitative description for many interesting phenomena such as magnetism, electron correlation, and superconductivity. In particular, the Hubbard model has been widely employed for a study on transition-metal magnetism. In the limit of strong interaction ($`U/t1`$), the Hubbard model with the half-filled electron occupancy reduces to the Heisenberg or Ising model. The two-site Hubbard model has been adopted for a study on some charge-transfer salts like tetracyanoquinodimethan (TCNQ) with dimerized structures -. Their susceptibility and specific heat were analyzed by taking into account the interdimer hopping within the BGS. The NES calculations have been made for thermodynamical properties of canonical and grand-canonical ensembles of Hubbard dimers, each of which is described by the two-site Hubbard model. It has been shown that the temperature dependences of the specific heat and susceptibility is significantly different from those calculated by the BGS when the entropic index $`q`$ departs from unity for small $`N`$ \[Eq. (2)\], the NES in the limit of $`q=1`$ reducing to the BGS.
The purpose of the present paper is to show (1) how thermodynamical property of a nanocluster containing a small number of Hubbard dimers is different from that of macroscopic systems, and (2) how thermodynamical property of a given nanocluster is changed when $`M`$, the number of Hubbard dimers contained in it, is varied. The paper is organized as follows. In Sec. 2, I apply the NES to nanoclustes, providing expressions for the energy, entropy, magnetization, specific heat and susceptibility. Numerical calculations of the temperature and magnetic-field dependences of thermodynamical quantities are reported for various $`M`$ values. The final Sec. 4 is devoted to discussions and conclusions. In the Appendix, the NES has been applied to a cluster containing spin dimers described by the Heisenberg model.
## 2 Nonextensive thermodynamics of Hubbard dimers
### 2.1 Energy and entropy
I have adopted a system consisting of sparsely distributed $`N_c`$ nanoclusters, each of which contains independent $`M`$ dimers. It has been assumed that the distance between nanoclusters is larger than $`\xi `$, the coherence length of the fluctuating $`\beta `$ field, and that the linear size of the clusters is smaller than $`\xi `$. Physical quantities such as the entropy and energy are extensive for $`N_c`$, but not for $`M`$ in general .
The Humiltonian of the cluster is given by
$`H`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=1}{\overset{M}{}}}H_{\mathrm{}}^{(d)},`$ (17)
$`H_{\mathrm{}}^{(d)}`$ $`=`$ $`t{\displaystyle \underset{\sigma }{}}(a_{1\sigma }^{}a_{2\sigma }+a_{2\sigma }^{}a_{1\sigma })+U{\displaystyle \underset{j=1}{\overset{2}{}}}n_jn_j\mu _BB{\displaystyle \underset{j=1}{\overset{2}{}}}(n_jn_j),`$ (18)
($`1,2\mathrm{}`$)
where $`H_{\mathrm{}}^{(d)}`$ denotes the two-site Hamiltonian for the $`\mathrm{}`$th dimer, $`n_{j\sigma }=a_{j\sigma }^{}a_{j\sigma }`$, $`a_{j\sigma }`$ the annihilation operator of an electron with spin $`\sigma `$ on a site $`j`$ ($`\mathrm{}`$), $`t`$ the hopping integral, $`U`$ the intraatomic interaction, $`\mu _B`$ the Bohr magneton, and $`B`$ an applied magnetic field. In the case of the half-filled occupancy, in which the number of electrons is $`N_e=2`$, six eigenvalues of $`H_{\mathrm{}}^{(d)}`$ are given by
$$ϵ_i\mathrm{}=0,\mathrm{\hspace{0.33em}2}\mu _BB,2\mu _BB,U,\frac{U}{2}+\mathrm{\Delta },\frac{U}{2}\mathrm{\Delta },\text{for }i=16,\mathrm{}=1M$$
(19)
where $`\mathrm{\Delta }=\sqrt{U^2/4+4t^2}`$ . The number of eigenvalues of the total Hamiltonian $`H`$ is $`6^M`$.
First we employ the BGS, in which the canonical partition function for $`H`$ is given by
$`Z_{BG}`$ $`=`$ $`\mathrm{Tr}\mathrm{exp}(\beta H),`$ (20)
$`=`$ $`{\displaystyle \underset{i_1=1}{\overset{6}{}}}{\displaystyle \underset{i_M=1}{\overset{6}{}}}\mathrm{exp}[\beta (ϵ_{i_1}++ϵ_{i_M})],`$ (21)
$`=`$ $`[Z_{BG}^{(d)}]^M,`$ (22)
$`Z_{BG}^{(d)}`$ $`=`$ $`1+2\mathrm{cosh}(2\beta \mu _BB)+e^{\beta U}+2e^{\beta U/2}\mathrm{cosh}(\beta \mathrm{\Delta }),`$ (23)
where $`\beta =1/k_BT`$, Tr denotes the trace and $`Z_{BG}^{(d)}`$ the partition function for a single dimer. By using the standard method in the BGS, we can obtain various thermodynamical quantities of the system . Because of a power expression given by Eq. (22), the energy and entropy are proportional to $`M`$: $`E_{BG}=ME_{BG}^{(d)}`$ and $`S_{BG}=MS_{BG}^{(d)}`$ where $`E_{BG}^{(d)}`$ and $`S_{BG}^{(d)}`$ are for a single dimer. This is not the case in the NES as will be discussed below.
Next we adopt the NES, where the entropy $`S_q`$ for the quantum system is defined by
$$S_q=k_B\left(\frac{\mathrm{Tr}(\rho _q^q)1}{1q}\right).$$
(24)
Here $`\rho _q`$ stands for the generalized canonical density matrix, whose explicit form will be determined shortly \[Eq. (27)\]. We will impose the two constraints given by
$`\mathrm{Tr}(\rho _q)`$ $`=`$ $`1,`$ (25)
$`{\displaystyle \frac{\mathrm{Tr}(\rho _q^qH)}{\mathrm{Tr}(\rho _q^q)}}`$ $``$ $`<H>_q=E_q,`$ (26)
where the normalized formalism is adopted . The variational condition for the entropy with the two constraints given by Eqs. (25) and (26) yields
$$\rho _q=\frac{1}{X_q}\mathrm{exp}_q\left[\left(\frac{\beta }{c_q}\right)(HE_q)\right],$$
(27)
with
$$X_q=\mathrm{Tr}\left(\mathrm{exp}_q\left[\left(\frac{\beta }{c_q}\right)(HE_q)\right]\right),$$
(28)
$$c_q=\mathrm{Tr}(\rho _q^q)=X_q^{1q},$$
(29)
where $`\mathrm{exp}_q(x)`$ is the $`q`$-exponential function given by Eq. (6) and $`\beta `$ is a Lagrange multiplier given by
$$\beta =\frac{S_q}{E_q}.$$
(30)
The trace in Eq. (28) and (29) is performed over the $`6^M`$ eigenvalues, for example, as
$`X_q`$ $`=`$ $`{\displaystyle \underset{i_1=1}{\overset{6}{}}}{\displaystyle \underset{i_M=1}{\overset{6}{}}}\left(\mathrm{exp}_q[\left({\displaystyle \frac{\beta }{c_q}}\right)(ϵ_{i_1}++ϵ_{i_M}E_q)]\right),`$ (31)
$``$ $`{\displaystyle \underset{i}{}}\left(\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)(ϵ_iE_q)\right]\right),`$ (32)
where the following conventions are adopted:
$`ϵ_i`$ $`=`$ $`ϵ_{i_1}++ϵ_{i_M},`$ (33)
$`{\displaystyle \underset{i}{}}`$ $`=`$ $`{\displaystyle \underset{i_1=1}{\overset{6}{}}}{\displaystyle \underset{i_M=1}{\overset{6}{}}}.`$ (34)
It is noted that in the limit of $`q=1`$, Eq. (31) reduces to
$$X_1=Z_{BG}\mathrm{exp}[\beta E_1]=[Z_{BG}^{(d)}\mathrm{exp}(\beta E_{BG}^{(d)})]^M.$$
(35)
For $`q1`$, however, $`X_q`$ cannot be expressed as a power form because of the property of the $`q`$-exponential function:
$`\mathrm{exp}_q(x+y)`$ $``$ $`\mathrm{exp}_q(x)\mathrm{exp}_q(y).\text{(for }q1\text{)}`$ (36)
It is necessary to point out that $`E_q`$ in Eq. (26) includes $`X_q`$ which is expressed by $`E_q`$ in Eq. (28). Then $`E_q`$ and $`X_q`$ have to be determined self-consistently by Eqs. (26)-(29) with the $`T\beta `$ relation given by Eq. (16) for a given temperature $`T`$. The calculation of thermodynamical quantities in the NES generally becomes more difficult than that in BGS.
### 2.2 Specific heat
The specific heat in the NES is given by
$$C_q=\left(\frac{d\beta }{dT}\right)\left(\frac{dE_q}{d\beta }\right).$$
(37)
Because $`E_q`$ and $`X_q`$ are determined by Eqs. (26)-(29), we get simultaneous equations for $`dE_q/d\beta `$ and $`dX_q/d\beta `$, given by
$`{\displaystyle \frac{dE_q}{d\beta }}`$ $`=`$ $`a_{11}\left({\displaystyle \frac{dE_q}{d\beta }}\right)+a_{12}\left({\displaystyle \frac{dX_q}{d\beta }}\right)+b_1,`$ (38)
$`{\displaystyle \frac{dX_q}{d\beta }}`$ $`=`$ $`a_{21}\left({\displaystyle \frac{dE_q}{d\beta }}\right)+a_{22}\left({\displaystyle \frac{dX_q}{d\beta }}\right),`$ (39)
with
$`a_{11}`$ $`=`$ $`q\beta X_q^{q2}{\displaystyle \underset{i}{}}w_i^{2q1}ϵ_i,`$ (40)
$`a_{12}`$ $`=`$ $`X_q^1E_q\beta q(q1)X_q^{q3}{\displaystyle \underset{i}{}}w_i^{2q1}ϵ_i(ϵ_iE_q),`$ (41)
$`a_{21}`$ $`=`$ $`\beta X_q^q,`$ (42)
$`a_{22}`$ $`=`$ $`0,`$ (43)
$`b_1`$ $`=`$ $`qX_q^{q2}{\displaystyle \underset{i}{}}w_i^{2q1}ϵ_i(ϵ_iE_q),`$ (44)
$`w_i`$ $`=`$ $`\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)(ϵ_iE_q)\right],`$ (45)
$`X_q`$ $`=`$ $`{\displaystyle \underset{i}{}}w_i.`$ (46)
The specific heat is then given by
$$C_q=\left(\frac{d\beta }{dT}\right)\left(\frac{b_1}{1a_{11}a_{12}a_{21}}\right).$$
(47)
with
$`{\displaystyle \frac{d\beta }{dT}}`$ $`=`$ $`\left({\displaystyle \frac{\beta ^2}{X_q^{1q}\beta (1q)X_q^q(dX_q/d\beta )}}\right),`$ (48)
In the limit of $`q1`$, Eqs. (38)-(46) yield the specific heat in the BGS, given by
$$C_{BG}=\frac{dE_{BG}}{dT}=k_B\beta ^2(<ϵ_i^2>_1<ϵ_i>_1^2),$$
(49)
where $`<>_1`$ is defined by Eq. (26) with $`q=1`$:
$$<Q_i>_1=X_1^1\underset{i}{}\mathrm{exp}[\beta (ϵ_iE_1)]Q_i=Z_{BG}^1\underset{i}{}\mathrm{exp}(\beta ϵ_i)Q_i.$$
(50)
### 2.3 Magnetization
The field-dependent magnetization $`m_q`$ in the NES is given by
$`m_q`$ $`=`$ $`{\displaystyle \frac{E_q}{B}}+(k_B\beta )^1{\displaystyle \frac{S_q}{B}},`$ (51)
$`=`$ $`{\displaystyle \frac{E_q}{B}}+\beta ^1X_q^q{\displaystyle \frac{X_q}{B}}.`$ (52)
By using Eqs. (26)-(29), we get the simultaneous equations for $`E_q/B`$ and $`X_q/B`$ given by
$`{\displaystyle \frac{E_q}{B}}`$ $`=`$ $`a_{11}{\displaystyle \frac{E_q}{B}}+a_{12}{\displaystyle \frac{X_q}{B}}+d_1,`$ (53)
$`{\displaystyle \frac{X_q}{B}}`$ $`=`$ $`a_{21}{\displaystyle \frac{E_q}{B}}+a_{22}{\displaystyle \frac{X_q}{B}}+d_2,`$ (54)
with
$`d_1`$ $`=`$ $`X_q^1{\displaystyle \underset{i}{}}w_i^q\mu _i+\beta qX_q^{q2}{\displaystyle \underset{i}{}}w_i^{2q1}ϵ_i\mu _i,`$ (55)
$`d_2`$ $`=`$ $`\beta X_q^{q1}{\displaystyle \underset{i}{}}w_i^q\mu _i,`$ (56)
where $`\mu _i=ϵ_i/B`$, and $`a_{ij}`$ ($`i,j=1,2`$) are given by Eqs. (40)-(43). From Eqs. (51)-(56), we obtain $`m_q`$ given by
$`m_q`$ $`=`$ $`\left({\displaystyle \frac{c_{12}+\beta ^1X_q^q(1c_{11})}{1c_{11}c_{12}c_{21}}}\right)d_2,`$ (57)
$`=`$ $`X_q^1{\displaystyle \underset{i}{}}w_i^q\mu _i=<\mu _i>_q.`$ (58)
In the limit of $`q1`$, Eqs. (55) and (56) reduce to
$`d_1`$ $`=`$ $`\mu _i_1+\beta ϵ_i\mu _i_1,`$ (59)
$`d_2`$ $`=`$ $`\beta X_1\mu _i_1,`$ (60)
where $`<>_1`$ is given by Eq. (50). By using Eq. (58), we get
$`m_{BG}`$ $`=`$ $`\mu _i_1,`$ (61)
$`=`$ $`{\displaystyle \frac{4\mu _B\mathrm{sinh}(2\beta B)}{Z_{BG}}},`$ (62)
where $`Z_{BG}`$ and $`<>_1`$ are given by Eqs. (20) and (50), respectively.
### 2.4 Susceptibility
The high-field susceptibility in the NES is given by
$$\chi _q(B)=\frac{m_q}{B}.$$
(63)
The zero-field susceptibility $`\chi _q(B=0)`$ is given by
$`\chi _q`$ $`=`$ $`\chi _q(B=0)=E_q^{(2)}+\beta ^1X_q^qX_q^{(2)},`$ (64)
where $`E_q^{(2)}=^2E_q/B^2_{B=0}`$ and $`X_q^{(2)}=^2X_q/B^2_{B=0}`$. With the use of Eqs. (26)-(29), we get simultaneous equations for $`E_q^{(2)}`$ and $`X_q^{(2)}`$ given by
$`E_q^{(2)}`$ $`=`$ $`a_{11}E_q^{(2)}+a_{12}X_q^{(2)}+f_1,`$ (65)
$`X_q^{(2)}`$ $`=`$ $`a_{21}E_q^{(2)}+a_{22}X_q^{(2)}+f_2,`$ (66)
with
$`f_1`$ $`=`$ $`2\beta qX_q^{q2}{\displaystyle \underset{i}{}}w_i^{2q1}\mu _i^2,`$ (67)
$`f_2`$ $`=`$ $`\beta ^2qX_q^{2(q1)}{\displaystyle \underset{i}{}}w_i^{2q1}\mu _i^2,`$ (68)
where $`a_{ij}`$ ($`i,j=1,2`$) are given by Eqs. (40)-(43). From Eqs. (64)-(68), we get
$`\chi _q`$ $`=`$ $`{\displaystyle \frac{f_2}{a_{21}}}=\beta qX_q^{q2}{\displaystyle \underset{i}{}}w_i^{2q1}\mu _i^2_{B=0}.`$ (69)
In the limit of $`q=1`$, Eq. (69) yields the susceptibility in BGS:
$`\chi _{BG}`$ $`=`$ $`\beta <\mu _i^2_{B=0}>_1,`$ (70)
$`=`$ $`\left({\displaystyle \frac{\mu _B^2}{k_BT}}\right){\displaystyle \frac{8}{3+e^{\beta U}+2e^{\beta U/2}\mathrm{cosh}(\beta \mathrm{\Delta })}}.`$ (71)
## 3 Calculated results
### 3.1 Temperature dependence
In order to study how thermodynamical quantities of a cluster containing $`M`$ Hubbard dimers depend on $`M`$, I have made some NES calculations, assuming the $`Mq`$ relation given by
$$q=1+\frac{1}{M},$$
(72)
which is derived from Eq. (2) with $`M=2N`$ for dimers. Simultaneous equations for $`E_q`$ and $`X_q`$ given by Eqs. (26)-(29) have been solved by using the Newton-Raphson method with initial values of $`E_1`$ and $`X_1`$ obtained from BGS ($`q=1`$) corresponding to $`M=\mathrm{}`$ in Eq. (72). Calculated quantities are given per dimer.
Figures 1(a), 1(b) and 1(c) show the temperature dependence of the specific heat $`C_q`$ for $`U/t=0`$, 5 and 10, respectively, with various $`M`$ values. The specific heat for $`M=\mathrm{}`$ shown by bold solid curves, expresses the result in BGS, and it has a peak at lower temperatures for the larger interaction, as previous BGS calculations showed . Note that the horizontal scales of Fig. 1(c) are enlarged compared to those of Figs. 1(a) and 1(b). The peak becomes broader for smaller $`M`$.
The temperature dependence of the susceptibility $`\chi _q`$ for $`U/t=0`$, 5 and 10 is plotted in Figs. 2(a), 2(b) and 2(c), respectively. The susceptibility for $`M=\mathrm{}`$ (BGS) shown by the bold solid curve, has a larger peak at lower temperatures for larger $`U`$ . Note that the horizontal and vertical scales of Fig. 2(c) are different from those of Figs. 2(a) and 2(b). We note that for smaller $`M`$, the peak in $`\chi _q`$ becomes broader, which is similar to the behavior of the specific heat shown in Figs. 1(a)-1(c).
When the $`M`$ value is varied, maximum values of the specific heat ($`C_q^{}`$) and the susceptibility ($`\chi _q^{}`$) are changed, and the temperatures ($`T_C^{}`$ and $`T_\chi ^{}`$) where these maxima are realized, are also changed. Figure 3(a) depicts $`T_C^{}`$ and $`T_\chi ^{}`$ for $`U/t=5`$ as a function of $`1/M`$. It is shown that with increasing $`1/M`$, $`T_\chi ^{}`$ is much increased than $`T_C^{}`$. Similarly, the $`1/M`$ dependences of $`C_q^{}`$ and $`\chi _q^{}`$ for $`U/t=5`$ are plotted in Fig. 3(b), which shows that maximum values of $`C_q`$ and $`\chi _q`$ are decreased with decreasing $`M`$. This trend against $`1/M`$ is due to the fact that a decrease in $`M(=2N)`$ yields an increase in fluctuations of $`\beta `$ fields, and then peaked structures of the specific heat and susceptibility realized in the BGS, are smeared out by $`\beta `$ in Eq. (1).
### 3.2 Magnetic-field dependence
Next I discuss the magnetic-field dependence of physical quantities. Figure 4 shows the $`B`$ dependence of the magnetization $`m_q`$ for $`U/t=0`$, 5 and 10 with $`M=2`$ at $`k_BT/t=1`$. For $`U/t=0`$, $`m_q`$ in the NES is smaller than that in the BGS at $`\mu _BB/t<1`$, but at $`\mu _BB/t>1`$ the former becomes larger than the latter. In contrast, in cases of $`U/t=5`$ and 10, $`m_q`$ in the NES is larger than that in the BGS for $`\mu _BB/t>0`$. In order to study the $`B`$ dependence in more details, I show in Fig. 5 the $`B`$ dependence of the six eigenvalues of $`ϵ_i`$ for $`U/t=5`$ \[Eq. (19)\]. We note the crossing of the lowest eigenvalues of $`ϵ_3`$ and $`ϵ_6`$ at the critical filed:
$$\mu _BB_c=\sqrt{\frac{U^2}{16}+t^2}\frac{U}{4},$$
(73)
leading to $`\mu _BB_c/t=0.351`$ for $`U/t=5.0`$. At $`B<B_c`$ ($`B>B_c`$), $`ϵ_6`$ ($`ϵ_3`$) is the ground state. At $`B=B_c`$ the magnetization $`m_q`$ is rapidly increased as shown in Figs. 6(a) and 6(b) for $`k_BT/t=1.0`$ and 0.1, respectively: the transition at lower temperatures is more evident than at higher temperatures. This level crossing also yields a peak in $`\chi _q`$ \[Figs. 6(c) and 6(d)\] and a dip in $`C_q`$ \[Figs. 6(e) and 6(f). It is interesting that the peak of $`\chi _q`$ for $`M=2`$ is more significant than that for $`M=\mathrm{}`$ whereas the dip of $`C_q`$ for $`M=2`$ is broader than that for $`M=\mathrm{}`$. When the temperature becomes higher, these peak structures become less evident as expected. Similar phenomenon in the field-dependent specific heat and susceptibility have been pointed out in the Heisenberg model within the BGS .
In the case of the quarter-filled occupancy ($`N_e=1`$), the eigenvalues are $`ϵ_i=t\mu _BB`$, $`t+\mu _BB`$, $`t\mu _BB`$, and $`t+\mu _BB`$ for $`i=14`$. Although the level crossing occurs between $`ϵ_2`$ and $`ϵ_3`$ at $`\mu _BB=t`$, it does not show any interesting behavior because the crossing occurs between the excited states. The case for the three-quarter-filled occupancy ($`N_e=3`$) is the same as that of the quarter-filled occupancy because of the electron-hole symmetry of the model.
Figure 6(b) reminds us the quantum tunneling of magnetization observed in magnetic molecular clusters such as Mn12 and Fe8 , which originates from the level crossings of magnetic molecules when a magnetic field is applied .
## 4 Discussions and conclusions
I have applied the NES to Hubbard dimers for a study of their thermodynamical properties. The current NES is, however, still in its infancy, having following unsettled issues.
(i) For relating the physical temperature $`T`$ to the Lagrange multiplier $`\beta `$, I have employed the $`T\beta `$ relation given by Eq. (16). There is an alternative proposal with the $`T\beta `$ relation given by
$`T`$ $`=`$ $`{\displaystyle \frac{1}{k_B\beta }},\text{(TMP)}`$ (74)
which is the same as in the BGS. At the moment, it has not been established which of the AMP and TMP methods given by Eqs. (16) and (74), respectively, is appropriate as the $`T\beta `$ relation in the current NES. It has been demonstrated that the negative specific heat of a classical gas model realized in the TMP method , is remedied in the AMP method . Recent theoretical analyses also suggest that the AMP method is better than the TMP method . The TMP method yields an anomalously large Curie constant of the susceptibility in the free spin model and in the Hubbard model . In my previous papers -, NES calculations have been made by using the TMP and AMP methods. It has been shown that both methods yield qualitatively similar results although there are some quantitative difference between the two: the non-extensivity in the TMP method generally appears more significant than that in the AMP method.
(ii) The $`Nq`$ relation given by Eq. (2) was obtained in Eqs. (1)-(5) with the $`\mathrm{\Gamma }`$ distribution $`f^B(\beta )`$ given by Eq. (3). Alternatively, by using the large-deviation approximation, Touchette has obtained the distribution function $`f^T(\beta )`$, in place of $`f^B(\beta )`$, given by
$`f^T(\beta )`$ $`=`$ $`{\displaystyle \frac{\beta _0}{\mathrm{\Gamma }\left(\frac{N}{2}\right)}}\left({\displaystyle \frac{N\beta _0}{2}}\right)^{\frac{N}{2}}\beta ^{\frac{N}{2}2}\mathrm{exp}\left({\displaystyle \frac{N\beta _0}{2\beta }}\right).`$ (75)
For $`N\mathrm{}`$, both $`f^B(\beta )`$ and $`f^T(\beta )`$ distribution functions reduce to the delta-function densities, and for a large $`N(>100)`$, both distribution functions lead to similar results. For a small $`N(<10)`$, however, there is a clear difference between the two distribution functions (see Fig. 4 of Ref.). It should be noted that $`f^T`$ cannot lead to the $`q`$-exponential function which plays a crucial role in the NES. For a large $`ϵ`$, the $`\mathrm{\Gamma }`$ distribution $`f^B`$ in Eq. (1) yields the power form of $`w(ϵ)ϵ^{\frac{1}{q1}}`$ while $`f^T`$ substituted to Eq. (1) leads to the stretched exponential form of $`w(ϵ)e^{c\sqrt{ϵ}}`$. This issue of $`f`$ versus $`f^T`$ is related to the superstatistics, which is currently studied with much interest .
To summarize, I have discussed thermodynamical properties of a nanocluster containing $`M`$ dimers, applying the NES to the Hubbard model. It has been demonstrated that the thermodynamical properties of a nanocluster with a small $`M`$ calculated by the NES may be considerably different from those obtained by the BGS. It is interesting to compare our theoretical prediction with experimental results for samples containing a small number of transition-metal dimers. Unfortunately samples with such a small number of dimers have not been reported: samples having been so far synthesized include macroscopic numbers of dimers, to which the present analysis cannot be applied. I expect that it is possible to form a dimer assembly by STM manipulation of individual atoms . Scanning probes may be used also as dipping pens to write small dimerized structures . Theoretical and experimental studies on nanoclusters with changing $`M`$ could clarify a link between the behavior of the low-dimensional infinite systems and finite-size nanoscale systems. I hope that the unsettle issues (i) and (ii) in the current NES mentioned above are expected to be resolved by future experiments on nanosystems with changing their sizes. It would be interesting to adopt quantum-master-equation and quantum-Langevin-equation approaches, and/or to perform large-scale molecular-dynamical simulations, for nanoclusters described by the Hubbard model.
## Acknowledgements
It is my great pleasure that on the occasion of the 60th birthday of Professor David G. Pettifor, I could dedicate the present paper to him, with whom I had an opportunity of collaborating in Imperial College London for one year from 1980 to 1981.
Appendix: NES for Heisenberg dimers
I have considered a cluster containing $`M`$ spin dimers (called Heisenberg dimers) described by the Heisenberg model ($`s=1/2`$) given by
$`H`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=1}{\overset{M}{}}}H_{\mathrm{}}^{(d)},`$ (76)
$`H_{\mathrm{}}^{(d)}`$ $`=`$ $`J𝐬_1𝐬_2g\mu _BB(s_{1z}+s_{2z}),\text{(}1,2\mathrm{}\text{)}`$ (77)
where $`J`$ stands for the exchange interaction, $`g`$ (=2) the g-factor, $`\mu _B`$ the Bohr magneton, and $`B`$ an applied magnetic field. Four eigenvalues of $`H_{\mathrm{}}^{(d)}`$ are given by
$`ϵ_i\mathrm{}`$ $`=`$ $`{\displaystyle \frac{J}{4}}g\mu _BBm_i,\text{with }m_1=1,\mathrm{\hspace{0.25em}0},1\text{ for }i=1,2,3,`$ (78)
$`=`$ $`{\displaystyle \frac{3J}{4}}g\mu _BBm_i.\text{with }m_4=0\text{ for }i=4,`$
In the BGS the canonical partition function is given by -
$`Z_{BG}`$ $`=`$ $`[Z_{BG}^{(d)}]^M,`$ (79)
$`Z_{BG}^{(d)}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\beta J}{4}}\right)[1+2\mathrm{c}\mathrm{o}\mathrm{s}\mathrm{h}(g\mu _B\beta B)]+\mathrm{exp}\left({\displaystyle \frac{3\beta J}{4}}\right),`$ (80)
with which thermodynamical quantities are easily calculated. The susceptibility is, for example, given by
$`\chi _{BG}`$ $`=`$ $`M\chi _{BG}^{(d)},`$ (81)
$`\chi _{BG}^{(d)}`$ $`=`$ $`{\displaystyle \frac{\mu _B^2}{k_BT}}\left({\displaystyle \frac{8}{3+\mathrm{exp}(J/k_BT)}}\right).`$ (82)
The calculation of thermodynamical quantities in the NES for the Heisenberg model goes parallel to that discussed in Sec. 2 if we employ eigenvalues given by Eq. (78). For example, by using Eq. (69), we get the susceptibility for the Heisenberg model, given by
$`\chi _q`$ $`=`$ $`g^2\mu _B^2\left({\displaystyle \frac{q\beta }{c_q}}\right){\displaystyle \frac{1}{X_q}}{\displaystyle \underset{i}{}}w_i^{2q1}m_i^2.`$ (83)
In the case of $`M=1`$ (a single dimer), we get
$`\chi _q^{(d)}`$ $`=`$ $`g^2\mu _B^2\left({\displaystyle \frac{q\beta }{c_q}}\right)\left({\displaystyle \frac{2}{X_q}}\right)\left(\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)\left({\displaystyle \frac{J}{4}}+E_q\right)\right]\right)^{2q1},`$ (84)
with
$`X_q`$ $`=`$ $`3\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)\left({\displaystyle \frac{J}{4}}+E_q\right)\right]+\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)\left({\displaystyle \frac{3J}{4}}E_q\right)\right],`$ (85)
$`E_q`$ $`=`$ $`{\displaystyle \frac{1}{X_q}}\{\left({\displaystyle \frac{3J}{4}}\right)\left(\mathrm{exp}_q\left[\left({\displaystyle \frac{\beta }{c_q}}\right)({\displaystyle \frac{J}{4}}+E_q)\right]\right)^q`$ (86)
$`+`$ $`\left({\displaystyle \frac{3J}{4}}\right)\left(\mathrm{exp}_q\left[({\displaystyle \frac{\beta }{c_q}})({\displaystyle \frac{3J}{4}}E_q)\right]\right)^q\}.`$
In the limit of $`q=1`$, Eq. (84) reduces to $`\chi _{BG}^{(d)}`$ given by Eq. (82).
The Curie constant $`\mathrm{\Gamma }_q`$ defined by $`\chi _q=(\mu _B^2/k_B)(\mathrm{\Gamma }_q/T)`$ for $`TJ`$ is given by
$`\mathrm{\Gamma }_q`$ $`=`$ $`2Mq,\text{(AMP)}`$ (87)
$`=`$ $`2Mq\mathrm{\hspace{0.25em}4}^{M(q1)}.\text{(TMP)}`$ (88)
Equations (87) and (88) are derived with the use of the $`T\beta `$ relation given by Eqs. (16) and (74), respectively. These are consistent with results obtained for Hubbard dimes .
Figures 7(a) and 7(b) show the temperature dependence of the specific heat $`C_q`$ and susceptibility $`\chi _q`$ of Heisenberg dimers calculated with the use of Eq. (83) for $`M=`$ 1, 2, 3 and $`\mathrm{}`$ ($`M=\mathrm{}`$ corresponding to the BGS with $`q=1.0`$). We note that the results of Heisenberg dimers are quite similar to those of the Hubbard dimer for $`U/t=5`$ and 10 shown in Figs. 2(b), 2(c), 2(e) and 2(f). This is not surprising because the Hubbard model with the half-filled electron occupancy in the strong-coupling limit reduces to the Heisenberg model.
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# Influence of Chaos on the fusion enhancement by electron screening
## 1 Introduction
The relation between the tunneling process and dynamical chaos has been discussed with great interests in recent years.chaos ; kb Though the tunneling is completely quantum mechanical phenomenon, it is influenced by classical chaos. In the sense that the the chaos causes the fluctuation of the classical action which essentially determines the tunneling probability. We study the phenomenon by examining the screening effect by bound electrons in the low energy fusion reaction. In the low energy region the experimental cross sections with gas targets show an increasing enhancement with decreasing bombarding energy with respect to the values obtained by extrapolating from the data at high energies krauss . Many studies attempted to attribute the enhancement of the reaction rate to the screening effects by bound target electrons. In this context one often estimates the screening potential as a constant decrease of the barrier height in the tunneling region through a fit to the data. A puzzle has been that the screening potential obtained by this procedure exceeds the value of the so called adiabatic limit, which is given by the difference of the binding energies of the united atoms and of the target atom and it is theoretically thought to provide the maximum screening potential rolfs95 . Over these several years, the redetermination of the bare cross sections has been proposed theoretically barker and experimentally junker , using the Trojan Horse Method thm ; thm2 ; thm3 . The comparison between newly obtained bare cross sections, i.e., astrophysical S-factors, and the cross sections by the direct measurements gives a variety of values for the screening potential. These values are often smaller than the sudden limit or larger than the adiabatic limit. Theoretical studies performed using the time-dependent Hartree-Fock(TDHF) scheme skls ; ktab suggest that the screening potential is between the sudden and the adiabatic limits.
One of the aims of this paper is to try to assess the effect of the screening quantitatively. Up to now, the dynamical effects of bound electrons have been studied only in some limited cases with a few bound electrons(the D+$`d`$ with atomic target skls ; ktab and molecular D<sub>2</sub> target smkls , the <sup>3</sup>He+$`d`$ skls ) with the TDHF method. We investigate here the dynamical effects, including the tunneling region, for other systems with many bound electrons; D+D, <sup>3</sup>He+D, looking the effect of the electron capture of projectile. We see also some reactions including Li isotopes; <sup>6</sup>Li+$`d`$, <sup>6</sup>Li+D, <sup>7</sup>Li+$`p`$ and <sup>7</sup>Li+H.
To simulate the effects of many electrons, we use the constrained molecular dynamics (CoMD) model kb ; pmb ; kb-2 . At very low energies fluctuations are anticipated to play a substantial role. Such fluctuations are beyond the TDHF scheme. Not only TDHF calculations are, by construction, cylindrically symmetric around the beam axis. Such a limitation is not necessarily true in nature and the mean field dynamics could be not correct especially in presence of large fluctuations. Molecular dynamics contains all possible correlations and fluctuations due to the initial conditions(events). For the purpose of treating quantum-mechanical systems like target atoms and molecules, we use classical equations of motion with constraints to satisfy the Heisenberg uncertainty principle and the Pauli exclusion principle for each event pmb . In extending the study to the lower incident energies, we would like to stress the connection between the motion of bound electrons and chaos. In fact, depending on the dynamics, the behavior of the electron(s) is unstable and influences the relative motion of the projectile and the target. The feature is caused by the nonintegrablility of the $`N`$-body system($`N3`$) and it is well known that the tunneling probability can be modified by the existence of chaotic environment. We discuss the enhancement factor of the laboratory cross section in connection with the integrability of the system by looking the inter-nuclear and electronic oscillational motion. More specifically we analyze the frequency shift of the target electron due to the projectile and the small oscillational motion induced by the electron to the relative motion between the target and the projectile. We show that the increase of chaoticity in the electron motion decreases the fusion probability.
The paper is organized as follows. In sect. 2 we introduce the enhancement factor $`f_e`$ and describe the essence of the Constrained molecular dynamics approach briefly. In sect. 3 we apply it to asses the effect of the bound electrons during the nuclear reactions. We discuss also the relation between the amplitudes of the inter-nuclear oscillational motion and the enhancement factor. We summarize the paper in sect. 4.
## 2 Formalism
### 2.1 Enhancement Factor
We denote the reaction cross section at incident energy in the center of mass $`E`$ by $`\sigma (E)`$ and the cross section obtained in absence of electrons by $`\sigma _0(E)`$. The enhancement factor $`f_e`$ is defined as
$$f_e\frac{\sigma (E)}{\sigma _0(E)}.$$
(1)
If the effect of the electrons is well represented by the constant shift $`U_e`$ of the potential barrier, following alr ; skls , ($`U_eE`$):
$$f_e\mathrm{exp}\left[\pi \eta (E)\frac{U_e}{E}\right],$$
(2)
where $`\eta (E)`$ is the Sommerfeld parameter clayton .
### 2.2 Constrained Molecular Dynamics
We estimate the enhancement factor $`f_e`$ numerically using molecular dynamics approach;
$$\frac{d𝐫_i}{dt}=\frac{𝐩_ic^2}{_i},\frac{d𝐩_i}{dt}=_𝐫U(𝐫_i),$$
(3)
where (r<sub>i</sub>,p<sub>i</sub>) are the position, momentum of the particle $`i`$ at time $`t`$. $`_i=\sqrt{𝐩_i^2c^2+m_i^2c^4}`$, $`U(𝐫_i)`$ and $`m_i`$ are its energy, Coulomb potential and mass, respectively. We set the starting point of the reaction at 10Å inter-nuclear separation. In Eqs. (3) we do not take into account the quantum effect of Pauli exclusion principle and Heisenberg principle. As it is well known that these classical equations (3) can be derived by using the variational calculus of Lagrangian $``$ of the classical system as well. So as to take the feature of the Pauli blocking into account in this framework, we use the Lagrange multiplier method for constraints.
Our constraints which correspond to the Pauli blocking is $`\overline{f}_i1`$ in terms of phase space density, note that the phase space density can be directly related to the distance of two particles, i.e., $`𝐫_{ij}𝐩_{ij}`$, in the phase space. Here $`𝐫_{ij}=|𝐫_i𝐫_j|`$ and $`𝐩_{ij}=|𝐩_i𝐩_j|`$. The relation $`\overline{f}_i1`$ is fulfilled, if $`𝐫_{ij}𝐩_{ij}\xi _P\mathrm{}\delta _{S_i,S_j}`$, where $`\xi _P=2\pi (3/4\pi )^{2/3}`$. $`i,j`$ refer only to electrons and $`S_i,S_j(=\pm 1/2)`$ are their spin projection. For the Heisenberg principle $`𝐫_{ij}𝐩_{ij}\xi _H\mathrm{}`$, where $`\xi _H=1`$, $`i`$ and $`j`$ refer to not only electrons but the nucleus. It is determined to reproduce the correct energy of hydrogenic atoms. Obviously the conditions $`𝐫_{ij}𝐩_{ij}=\xi _{H(P)}\mathrm{}`$ must be fulfilled in the ground state configuration rather than $`𝐫_{ij}𝐩_{ij}>\xi _{H(P)}\mathrm{}`$.
Using these constraints, the Lagrangian of the system can be written down as
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{𝐩_i^2c^2}{_i}}{\displaystyle \underset{i,j(i)}{}}U(𝐫_{ij})+{\displaystyle \underset{i,j(i)}{}}\lambda _i^H\left({\displaystyle \frac{𝐫_{ij}𝐩_{ij}}{\mathrm{}}}1\right)`$
$`+`$ $`{\displaystyle \underset{i,j(i)}{}}\lambda _i^P\left({\displaystyle \frac{𝐫_{ij}𝐩_{ij}}{\xi _P\mathrm{}}}\delta _{S_i,S_j}1\right),`$ (4)
where $`\lambda _i^P`$ and $`\lambda _i^H`$ are Lagrange multipliers. The variational calculus leads
$`{\displaystyle \frac{d𝐫_i}{dt}}`$ $`=`$ $`{\displaystyle \frac{𝐩_ic^2}{_i}}+{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \underset{j(i)}{}}\left({\displaystyle \frac{\lambda _i^H}{\xi _H}}+{\displaystyle \frac{\lambda _i^P}{\xi _P}}\delta _{S_i,S_j}\right)𝐫_{ij}{\displaystyle \frac{𝐩_{ij}}{𝐩_i}},`$ (5)
$`{\displaystyle \frac{d𝐩_i}{dt}}`$ $`=`$ $`_𝐫U(𝐫_i){\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \underset{j(i)}{}}\left({\displaystyle \frac{\lambda _i^H}{\xi _H}}+{\displaystyle \frac{\lambda _i^P}{\xi _P}}\delta _{S_i,S_j}\right)𝐩_{ij}{\displaystyle \frac{𝐫_{ij}}{𝐫_i}}.`$ (6)
In order to obtain the atomic ground-state configuration, We perform the time integration of the eqs. (5) and (6). The value of $`\lambda _i^H`$ and $`\lambda _i^P`$ are determined depending on the magnitude of $`𝐫_{ij}𝐩_{ij}`$. If $`𝐫_{ij}𝐩_{ij}`$ is (smaller)larger than $`\xi _{H(P)}\mathrm{}`$, $`\lambda `$ has positive(negative) sign. Thus we change the phase space occupancy of the system. The constraints restrict us to variations $`\mathrm{\Delta }=0`$ that keep the constraints always true kb-2 . In this way we obtain many initial conditions which occupy different points in the phase space microscopically.
In order to treat the tunneling process, we define the collective coordinates $`𝐑^{coll}`$ and the collective momentum $`𝐏^{coll}`$ as
$$𝐑^{coll}𝐫_P𝐫_T;𝐏^{coll}𝐩_P𝐩_T,$$
(7)
where $`𝐫_T,𝐫_P`$ ($`𝐩_T,𝐩_P`$) are the coordinates(momenta) of the target and the projectile nuclei, respectively. When the collective momentum becomes zero, we switch on the collective force, which is determined by $`𝐅_P^{coll}\dot{𝐏}^{coll}`$ and $`𝐅_T^{coll}\dot{𝐏}^{coll}`$, to enter into imaginary time bk . We follow the time evolution in the tunneling region using the equations,
$$\frac{d𝐫_{T(P)}^{\mathrm{}}}{d\tau }=\frac{𝐩_{T(P)}^{\mathrm{}}}{_{T(P)}};\frac{d𝐩_{T(P)}^{\mathrm{}}}{d\tau }=_𝐫U(𝐫_{T(P)}^{\mathrm{}})2𝐅_{T(P)}^{coll},$$
(8)
where $`\tau `$ is used for imaginary time to be distinguished from real time $`t`$. $`𝐫_{T(P)}^{\mathrm{}}`$ and $`𝐩_{T(P)}^{\mathrm{}}`$ are position and momentum of the target (the projectile) during the tunneling process respectively. Adding the collective force corresponds to inverting the potential barrier which becomes attractive in the imaginary times. The penetrability of the barrier is given by bk
$$\mathrm{\Pi }(E)=\left(1+\mathrm{exp}\left(2𝒜(E)/\mathrm{}\right)\right)^1,$$
(9)
where the action integral $`𝒜(E)`$ is
$$𝒜(E)=_{r_b}^{r_a}𝐏^{coll}𝑑𝐑^{coll},$$
(10)
$`r_a`$ and $`r_b`$ are the classical turning points. The internal classical turning point $`r_b`$ is determined using the sum of the radii of the target and projectile nuclei. Similarly from the simulation without electron, we obtain the penetrability of the bare Coulomb barrier $`\mathrm{\Pi }_0(E)`$.
Since nuclear reaction occurs with small impact parameters on the atomic scale, we consider only head on collisions. The enhancement factor is thus given by eq. (1),
$$f_e=\mathrm{\Pi }(E)/\mathrm{\Pi }_0(E)$$
(11)
for each event in our simulation. Thus we have an ensemble of $`f_e`$ values at each incident energy.
## 3 Application to the Electron Screening Problem
#### 3.0.1 D+$`d`$ and D+D reactions
Fig. 1 shows the incident energy dependence of the enhancement factor for the reactions D+$`d`$ and D+D, where the systems involve 1 and 2 electrons respectively.
The open and closed squares show the average enhancement factors $`\overline{f_e}`$ over events for the reactions D+$`d`$ and D+D, respectively. The variances $`\mathrm{\Sigma }=\sqrt{\overline{f_e^2}(\overline{f_e})^2}`$ are shown with error bars. The dotted and dash-dotted curves show the enhancement factors in the adiabatic limit $`f_e^{(AD)}`$ for an atomic deuterium target and it is obtained by assuming equally weighted linear combination of the lowest-energy gerade and ungerade wave function for the electron, reflecting the symmetry in the D+$`d`$, i.e.,
$$f_e^{(AD)}=\frac{1}{2}\left(e^{\pi \eta (E)\frac{U_e^{(g)}}{E}}+e^{\pi \eta (E)\frac{U_e^{(u)}}{E}}\right),$$
(12)
where $`U_e^{(g)}=`$ 40.7 eV and $`U_e^{(u)}=`$ 0.0 eV ktab ; skls for D+$`d`$ case. If we take into account the electron capture of the projectile, i.e., in the case of D+D, the enhancement factor in the adiabatic limit is
$$f_e^{(AD)}=\frac{1}{4}e^{\pi \eta (E)\frac{U_e^{(g.s.)}}{E}}+\frac{3}{4}e^{\pi \eta (E)\frac{U_e^{(1es)}}{E}},$$
(13)
where $`U_e^{(g.s.)}=`$ 51.7 eV and $`U_e^{(1es)}=`$ 31.9 eV kt . The solid curve and dashed curve show the enhancement factors in the dissipative limit $`f_e^{(DL)}`$ for the reactions D+$`d`$ and D+D respectively. Notice how the calculated enhancement factor with their variances nicely ends up between the adiabatic and the dissipative limits. We performed also a fit of our data using eq. (2) including the very low energy region and obtained $`U_e=`$ 15.9 $`\pm `$ 2.0 eV for D+$`d`$ case and $`U_e=`$21.6 $`\pm `$ 0.3 eV for D+D.
Now we look at the oscillational motions of the particle’s coordinates as the projection on the $`z`$-axis (the reaction axis). We denote the $`z`$-component of $`𝐫_T,𝐫_P`$ and $`𝐫_e`$ as $`z_T,z_P`$ and $`z_e`$, respectively. Practically, we examine the oscillational motion of the electron around the target $`z_{Te}=z_ez_T`$ and the oscillational motion of the inter-nuclear motion, i.e., the motion between the target and the projectile, $`z_s=z_T+z_P`$, which essentially would be zero due to the symmetry of the system in the absence of the perturbation.
In Fig. 2 these two values are shown for 2 events, which have the enhancement factor $`f_e=`$ 170.8 (ev. A), and $`f_e=`$ 6.5 (ev. B), at the incident energy $`E_{cm}=0.15`$ keV. The panels show the $`z_s,z_{Te}`$ as a function of time. The stars indicate the time at which the system reaches the classical turning point. It is clear that in the case of event B the orbit of the electron is much distorted from the unperturbed one than in event A. Characteristics of $`z_s`$ are that (1) its value often becomes zero, as it is expected in the un-perturbed system, and (2) the component of the deviation from zero shows periodical behavior. It is remarkable that the amplitude of the deviation becomes quite large at some points in the case of event B which shows the small enhancement factor. Note that in event B one observes clear beats, i.e., resonances. Thus for two events, with the same macroscopic initial conditions, we have a completely different outcome, which is a definite proof of chaos in our 3-body system. We can understand these results in first approximation by considering the motion of the ions to be much slower than the rapidly oscillating motion of the electrons. kb From the Fig.2 we can deduce the following important fact. If the motion of the electron is initially in the plane perpendicular to the reaction axis, the enhancement factor is large, event A(notice $`|z_{Te}|R_B`$, i.e., the Bohr radius, at $`t0`$). On the other hand if there is a substantial projection of the electron motion, as in event B(the amplitude of $`|z_{Te}|R_B`$ at $`t0`$), on the reaction axis the enhancement factor is relatively small because of the increase of chaoticity. The fact suggests that if one performs experiments at very low bombarding energies with polarized targets, the enhancement factor can be controlled by changing the polarization. The largest enhancement would be gained with targets polarized perpendicularly to the beam axis.
In order to test this estimation, we prepared ensembles of target atoms which are polarized perpendicular(P) and parallel(P) to the beam axis, numerically. In Fig. 3 we show the incident energy dependence of the average enhancement factor for the P and P targets with pluses and crosses, particularly in the low energy region. The enhancement factors from the P targets are always larger than that from the P targets. In contrast to the average enhancement from the P targets, which increases monotonically as the incident energy becomes smaller, the average enhancement from the P targets fluctuates. It has also large variances at low energies. Remarkable thing is that with the parallel targets the enhancement factor often becomes less than 1. It means that in this case the bound electron gives the effect of hindrance to the tunneling probability.
#### 3.0.2 <sup>3</sup>He+$`d`$ and <sup>3</sup>He+D reactions
An excess of the screening potential was reported for the reactions <sup>3</sup>He+$`d`$ with atomic gas <sup>3</sup>He target, and D<sub>2</sub> \+ <sup>3</sup>He with deuterium molecular gas target, for the first time in the reference krauss . Since then various experiments have been performed for these reactions. The incident energy covers from 5 keV to 50 keV for <sup>3</sup>He+$`d`$. Though once the problem of the discrepancy between experimental data and theoretical prediction seemed to be solved by considering the correct energy loss data lsbr , recent measurements using measured energy loss data aliotta report larger screening potentials than in the adiabatic limit for both reactions.
The electron capture by the projectile plays a minor role in the case of <sup>3</sup>He+d, since electrons are more bound in helium targets. However in the recent measurement Aliotta et al. was performed using molecular D$`{}_{}{}^{+}{}_{2}{}^{}`$ and D$`{}_{}{}^{+}{}_{3}{}^{}`$ targets aliotta . Thus we assess the contribution from the reaction <sup>3</sup>He+D, as well.
The enhancement factor in the adiabatic limit give $`U_e`$=119 eV for <sup>3</sup>He+$`d`$ and $`U_e`$=110 eV for <sup>3</sup>He+D, respectively. These are shown in the figure 4 with the solid curve for <sup>3</sup>He+$`d`$ and with the dashed curve for <sup>3</sup>He+D. The comparison of these two adiabatic limits implies that the electron capture of projectile would give a hindrance compared with the bare deuteron projectile. Meanwhile the latest analysis of the experimental data using $`R`$-matrix two level fit barker suggests the screening potential $`U_e=`$ 60 eV(corresponding enhancement factor is shown with dotted curve). The comparison between direct measurement and an indirect method, the Trojan Horse method, suggests the screening potential $`U_e=`$ 180$`\pm `$40 eV ( the corresponding enhancement factor is shown with dot-dashed curve) thm3 . The average enhancement factors $`\overline{f_e}`$ over events in our simulations using the CoMD are shown with
the open and closed squares for the reactions <sup>3</sup>He+$`d`$ and <sup>3</sup>He+D, respectively. The enhancement factors of the both reactions <sup>3</sup>He+$`d`$ and <sup>3</sup>He+D are in agreement with the extracted values using the $`R`$-matrix approach within the variances over all the events. Notice that our calculated enhancement factors for the two systems display an opposite trend as compared to the adiabatic limits. The average enhancement factor of the reaction <sup>3</sup>He+D agrees with the estimation of the adiabatic limit and the reaction <sup>3</sup>He+d is below the corresponding adiabatic limit. The paradoxical feature comes from the fact that an electron between the two ions is often kicked out during the reaction process, i.e., the electron configuration seldom settles down the <sup>5</sup>Li<sup>+</sup> ground state in the reaction <sup>3</sup>He+d. It is known as autoionization in the context of the Classical Trajectory Monte Carlo method gr . Instead in the case of the <sup>3</sup>He+D, the deuterium projectile brings its bound electron in a tight bound state around the unified nuclei of <sup>3</sup>He and $`d`$, practically it ends up with a ground state configuration of the <sup>5</sup>Li atom. The fits of the obtained enhancement factors suggests the screening potentials $`U_e=`$ 82.4 $`\pm `$ 1.9 eV for the <sup>3</sup>He+$`d`$ and $`U_e=`$ 102.8 $`\pm `$ 3.0 eV for the <sup>3</sup>He+D.
#### 3.0.3 <sup>6</sup>Li+$`d`$, <sup>6</sup>Li+D, <sup>7</sup>Li+$`p`$ and <sup>7</sup>Li+H
The S-factors for the reactions <sup>6</sup>Li+$`d`$, <sup>6</sup>Li+$`p`$ and <sup>7</sup>Li+$`p`$ were measured over the energy range 10 keV $`<E_{cm}<`$ 500 keV by Engstler,et al. eknrsl . They used LiF solid targets and deuteron projectiles as well as deuterium molecular gas targets and Li projectiles.
In the case of LiF target which is a large band gap insulator, one often approximates the electronic structure of the target <sup>6</sup>Li(<sup>7</sup>Li) state by the <sup>6</sup>Li<sup>+</sup>(<sup>7</sup>Li<sup>+</sup>) with only two innermost electrons. Thus for all three reactions one expects the screening potential in the adiabatic limit $`U_e^{(AD)}=371.8198.2174`$ eV. Instead if one uses the ground state of the <sup>6</sup>Li(<sup>7</sup>Li) atom and of the bare deuteron target as the initial state, $`U_e^{(AD)}=`$186 eV bfmmq , which is given by the solid curve in Fig. 5 .
However one should be aware that the deuteron or hydrogen projectile plausibly moves with a bound electron in LiF solid insulator target eder . Under such an assumption we could estimate the screening potential $`U_e^{(AD)}=389.9198.2192`$ eV. In the case of molecular D<sub>2</sub> or H<sub>2</sub> gas targets, as well, we should consider the electron capture by the lithium projectile.
The bare S-factors for the same reaction have been extracted using an indirect method, the Trojan-Horse Method through the reaction <sup>6</sup>Li(<sup>6</sup>Li,$`\alpha \alpha `$)<sup>4</sup>He thm2 . The comparison between direct and the indirect methods gives the screening potential $`U_e=`$ 320$`\pm `$50 eV. The corresponding enhancement factors are shown with the dash-dotted curve. The contrast between the direct measurement data and the theoretical estimation for the bare S-factor using the $`R`$-matrix theory gives $`U_e`$=240 eV. It is shown with dotted line. The extracted $`U_e`$ with the two different methods are larger than the adiabatic limit.
We simulate the reactions <sup>6</sup>Li+$`d`$, <sup>6</sup>Li+D, <sup>7</sup>Li+$`p`$ and <sup>7</sup>Li+H. In the figure 5(and 6) the open and closed squares show the enhancement factor for the reactions <sup>6</sup>Li+$`d`$ and <sup>6</sup>Li+D,(and <sup>7</sup>Li+$`p`$ and <sup>7</sup>Li+H) respectively.
Again the average enhancement factors of the reaction <sup>6</sup>Li+D(<sup>7</sup>Li+H) are larger than those of the <sup>6</sup>Li+$`d`$(<sup>7</sup>Li+$`p`$). The enhancement factors of the reaction <sup>6</sup>Li+D are in agreement with the extracted values using the $`R`$-matrix approach within the variances over all the events. The fit of the obtained average enhancement factors suggests the screening potentials $`U_e=`$ 152.0 $`\pm `$ 9.9 eV for <sup>6</sup>Li+$`d`$ and $`U_e=`$ 214.4$`\pm `$18.5 for <sup>6</sup>Li+D. The screening potential for the reaction <sup>6</sup>Li+$`d`$ in our simulation does not exceed the adiabatic limit nor extracted values using the $`R`$-matrix theory and THM, but one for <sup>6</sup>Li+D verges on the extracted values using the $`R`$-matrix approach.
## 4 Summary
We discussed the effect of the screening by the electrons in nuclear reactions at the astrophysical energies. We performed molecular dynamics simulations with constraints and imaginary time for the reactions D+$`d`$, D+D, <sup>3</sup>He+$`d`$, <sup>3</sup>He+D, <sup>6</sup>Li+$`d`$, <sup>6</sup>Li+D, <sup>7</sup>Li+$`p`$, <sup>7</sup>Li+H. For all the reactions it is shown that both the average enhancement factors and their variances increase as the incident energy becomes lower. Using bare projectiles we obtained the average screening potential smaller than the value in the adiabatic limit for all reactions. It is because of the excitation or emission of several bound electrons during the reactions. The comparison between bare and atomic projectile cases for each reactions revealed that the electron capture of the projectile guides to larger enhancements. The derived enhancement factors in our simulation are in agreement with those extracted within the $`R`$-matrix approach including the variances over all the events.
We report also the results of the numerical experiments using polarized targets for the reaction D+d. Using P targets we obtained relatively large enhancements with small variances, instead P target gives large variances of the enhancement factors and relatively small averaged enhancement factors. It is because with the P targets the force exerted from the electron to the relative motion of the nuclei is oscillational, in the direction of the beam axis, and the motion of the electron becomes often excited or unstable. It is the case where the chaoticity of the electron motion affects the tunneling probability and at the same time the enhancement factor of the cross section. This suggests that if one performs experiments at very low bombarding energies with polarized targets, the enhancement factor can be controlled by changing the polarization. The largest enhancement with targets polarized perpendicularly to the beam direction.
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# Flux Modulation from Non-Axisymmetric Stuctures in Accretion Discs
## 1 Introduction
X-ray emission from accreting black holes in binary stellar systems varies on time scales ranging from milliseconds to years. Variability on time scales longer than days appears to be driven by changes in the accretion rate onto the black hole, and is often manifested as transient outbursts in which the luminosity of a source changes by a million-fold Lasota (2001). At the shortest time scales, quasi-periodic oscillations (QPOs) are observed in the X-ray emission. The highest frequency QPOs ($`>100`$Hz) are consistent with those expected from general relativistic orbits near the innermost stable orbit around the black hole, and they are likely caused by inhomogeneities in the inner accretion flow (Remillard et al., 2002, Stella & Vietri, 1998). However, the cause of low frequency QPO (LFQPO)s 0.1–20 Hz is still a mystery. This motivates the work herein. For neutron star X-ray binaries, low frequency QPOs can be sub-divided into more precise categories (e.g. Psaltis et al.,1999). Here we focus on a basic paradigm for the modulation and leave explanations of possible harmonic relations and phase lag behavior Varnière (2005) for further work. Black hole microquasars avoid the role of any solid stellar surface, and offer a purer probe of the accretion process without interaction with a stellar surface. The basic properties of LFQPOs that any model needs to explain are these:
(1) LFQPOs often appear to be present in widely spaced observations, over a period of weeks, with fractionally narrow frequency widths ($`\mathrm{\Delta }\nu /\nu 1/30`$) (Morgan et al., 1997, Remillard et al., 1999). Although the observations are not continuous, the fact that the LFQPO are observed to be the same (within a very small variation) over several observations within the same week/month, suggest that the QPO are likely always present during that time, at least when averaged over such long time scales within the hard state; there are periods where the LFQPO appear and disappear on timescale of a few seconds. If LFQPOs result from inhomogeneities orbiting at a Keplerian speed in the accretion disc they are constrained to a narrow range of radii.
(2) LFQPO amplitudes are typically 5-10% RMS (root mean squared), but can reach 20% RMS. Given that the X-ray luminosities of black hole binaries approach $`10^{39}`$ erg s<sup>-1</sup>, LFQPOs involve an enormously energetic fraction of the accretion flow.
(3) On the other hand, LFQPOs are transient features, so the spectral properties of the X-ray emission when LFQPOs are present and absent can be used to constrain the LFPQO origin (Muno et al., 1999, Sobczak et al., 2000).
(4) The 0.1–20 Hz LFQPOs only appear when the non-thermal component of the X-ray spectrum is strong, and their fractional amplitudes increase with energy between 2–20 keV. However, a thermal component that contributes $`10\%`$ of the X-ray emission at lower energies must also be present, and the frequencies of the LFQPO appear to be correlated best with this. The energy output of accreting black holes can generally be decomposed into a $`1`$ keV thermal component (thought to originate from the optically thick accretion disc) and a non-thermal component with a power-law spectrum extending beyond 100 keV (thought to result from inverse-Compton scattering of cool photons by a corona of hot electrons.) Thus, it appears that the 0.1–20 Hz QPOs either originate in the boundary between the disc and the corona, or are part of the mechanism which accelerates the hot Comptonizing electrons.
As mentioned, LFQPOs can be further subdivided (Psaltis et al., 1999), but the above represents a basic set of characteristics that a zeroth order LFQPO flux modulation model should explain.
Two different conceptual paradigms for LFQPOs have been proposed in this regard: (1) In the centrifugal pressure supported boundary layer model (CENBOL, Chakrabarti & Manickam, 2000), a QPO is produced by a shock in the accretion flow where it makes a transition from a Keplerian disc to a hot Comptonizing region. If the cooling time of the post-shock region is resonant with the free-fall time at the shock, the shock can oscillate radially with a frequency on order a Hz. The QPO is, in principle, produced because the shock modulates the flux of seed photons that reach the Comptonizing post-shock region. In this model, the LFQPO would represent a global radial oscillation that modifies the emitted flux. (2) A second paradigm for LFQPOs is non-axisymmetric structures in the disk. Their motion around the black hole would lead to flux modualtion. These structures could move at the Keplerian speed, a precession speed (such as Lense-Thirring) or a phase velocity associated with a structural instability. One example occurs in the accretion-ejection instability model (AEI, Tagger & Pellat, 1999, Varnière & Tagger, 2002). Here a spiral shock forms in an accretion disc threaded by a vertical magnetic field. The dispersion relation is similar to that of a galactic disc but with the gravitational potential replaced by a magnetic one. At the co-rotation radius, at which the Keplerian velocity matches that of the pattern, a Rossby vortex forms and emits vertical Alfvén waves into a corona. The LFQPO would result from the orbit of this non-axisymmetric structure in the disc. 2D numerical simulations of this instability, for example, demonstrated (Varnière et al., 2003) that the X-ray flux was modulated. However, in 2D, the constraints on the cause of the modulation and how it can be stronger and weaker were limited.
In this paper, we focus on the paradigm exemplified by the latter type of model, namely the production of LFQPOs by the orbit of a non-axisymmetric structure. Here we do not focus on the specific origin of the structure but rather investigate the consequences of the presence of such a structure. We investigate the specific question of whether a stable pattern in an accretion disc can reproduce the observed characteristics of LFQPO simply by its motion around the central engine. In Sec. 2 we give analytical formulae that can be used to model arbitrarily shaped blobs and spirals as non-axisymmetric disc features. In Sec. 3 we compute the emission from an accretion disc with such non-axisymmetric structures and discuss the influence of the shape of the particular non-axisymmetric structures on the RMS LFQPO amplitude. We also discuss the parameter choices which provide the most observationally consistent non-axisymmetric structures for the microquasar example discussed above. We also note that the formalism is generically applicable to other non-axisymmetric accretion disc systems. We conclude in Sec. 5.
## 2 Modeling non-axisymmetric discs
Blobs and spirals are two useful categories of non-axisymmetric structures. We use blob to indicate a generic localized feature and spirals to indicate the more specifically identifiable global structure produced by gravitational or MHD instabilities–the latter via the accretion-ejection instability (AEI) i.e. Tagger & Pellat 1999. Here we do not detail the formation of these non-axisymmetric structures but focus on providing an analytic framework that characterizes their shape for practical use and allows observational implications for flux modulation to be quantified.
Because discs of observed systems such as microquasars and active galactic nuclei (AGN) are not resolved, the disc structure cannot be directly imaged. But from studying the timing evolution of the flux, especially the presence of the LFQPO with its frequency and RMS amplitude, we can infer what structures might be present. We concentrate on the RMS amplitude of the modulation, assuming that the frequency is already matched by the presumed location of the given structure.
Taking into account the disc thickness is a particularly important aspect of our endeavor. For discs viewed at highly inclined angles, the shadowing from a local thickening of the disc can be important. The height and temperature profiles are coupled in the hydrostatic equilibrium approximation. Although we make this approximation here, detailed $`3`$D simulations will ultimately be needed to get an exact profile of $`h`$ and $`T`$ for a more realistic disc.
### 2.1 Analytic expression for the disc thickness
For the disc thickness, we write
$`h(r,\varphi )=h_o(r)+h_1(r(\varphi ),\varphi ))=h_o(r)+s(rr_s)d(r)`$ (1)
where $`h_1`$ is a perturbation in thickness around the unperturbed thickness $`h_o`$, $`r`$ is the radial location, and $`\varphi `$ is the azimuthal angle. We have used $`h_1=s(rr_s)d(r)`$, where $`s`$ is a “shape” function. The latter is finite only near the disc structure causing the non-axisymmetry. The blob or spiral wave feature is localized by $`r_sr_ce^{\alpha (r)\varphi }`$, where $`r_c`$ is the point where the structure begins and $`\alpha `$ is the opening angle of the structure. The quantity $`d(r)`$ is a thickness function which is defined as the height of the disc at each point.
We now take $`s(rr_s)`$ to be a Gaussian and $`d`$ to be a power-law in $`r`$. This provides a simple but useful framework to model non-axisymmetric structures. From (1) we then have
$`h(r,\varphi )=h_o(r)+\stackrel{~}{\gamma }\left({\displaystyle \frac{r_c}{r}}\right)^\beta e^{0.5\left(\frac{rr_s}{\delta }\right)^2},`$ (2)
where the constant $`\beta `$ measures how fast the thickness decreases from the maximum, $`\delta `$ parameterizes the radial extent of the structure, and $`\stackrel{~}{\gamma }`$ defines the disc thickness at $`r_c`$, We allow the maximum height $`\stackrel{~}{\gamma }`$ to be a function of the unperturbed thickness at that point, that is $`\stackrel{~}{\gamma }=\gamma h_o(r_c)`$. This allows consideration of cases with similar $`h_0/r`$ but different $`r_c`$. We also consider the number of times the non-axisymmetric structure winds around the disc.
The influence of the parameters in (2) is illustrated in Fig. 1. The role of $`\alpha `$ is seen in Fig. 1a which shows a top view of two spirals’ “spines” (i.e. the line defined by $`r_s(\alpha ,\varphi ,r_c)`$ tracing the maximal height above the unperturbed disc at each $`\varphi `$.). The larger the $`\alpha `$ the more open the spiral. The role of parameters $`\stackrel{~}{\gamma }`$, $`\delta `$, and $`\beta `$ are seen in Fig. 1b which shows cross sections of the disc height profile for different $`\varphi `$ projected into the $`r,z`$ plane for $`r_c=1.5`$. Increasing $`\delta `$ would increase the full width-half-maximum of the peaks, which in practice has a larger effect on the rise to the peak than the fall because the Gaussian perturbation is superimposed upon a positively sloped disk. Increasing $`\stackrel{~}{\gamma }`$ would increase the maximum height above the $`h_o`$, and the downward slope of the line that connects the peaks at different $`\varphi `$ values would be steeper for larger $`\beta `$. Fig. 1c shows a 3-D close-up of the inner region of the disc with a one-armed spiral characterized by $`h_o/r=0.01`$, $`r_c=3`$, $`\alpha =0.07`$, $`\stackrel{~}{\gamma }=0.3`$, $`\beta =2`$, and $`\delta =0.2`$ at an viewing angle of $`70^{}`$ from the normal.
### 2.2 Emission
Having obtained an expression for the disc height, we compute the temperature using the approximation that $`c_s=h\mathrm{\Omega }`$, which gives $`T=\mu /R_{gas}c_s^2`$. This means, for example, that a change in temperature by $`20\%`$ is related to a change in thickness of about $`45\%`$.
In order to determine the observational effect of a non-axisymmetric disc thickness, we compute the flux as a function of azimuth that an observer would received as time-dependent flux modulation when the disc rotates. Assuming that the spectrum from each point is a blackbody, we use the height-temperature relation above to compute the photon flux
$$f(E)=\frac{2}{c^2h^3}\frac{E^2}{\mathrm{exp}[E/kT]1}.$$
(3)
We then sum the flux from each cell, multiply by $`cos([\pi /2\theta ]\zeta )`$, where $`\theta `$ is the disc inclination angle from the normal and
$$\zeta =atan\left(\frac{dh}{dr}\mathrm{cos}\varphi \frac{dh}{rd\varphi }\mathrm{sin}\varphi \right).$$
(4)
We trace rays from the individual cells to the observer in order to determine whether the flux is intercepted by the outer portions of the disc. We compute the observed spectrum as the structure rotates by viewing a single snapshot at all angles $`\varphi `$. The spectrum is taken as multi-temperature blackbody, and we sum the number of photons received. We compute the Fourier transform of the profile in order to estimate the amplitudes of any modulation at multiples of the pattern frequency.
Our simple spectral modeling does not include scattering in the disc atmosphere, the presence of an electron corona, or special and general relativistic (GR) effects in part because the LFQPOs observed in microquasars could be coming from shadowing by structures at $`>100`$ gravitational radii. Also, the non-relativistic formalism applies at all radii for generic disks systems around stars. We note however, that non-Keplerian origins of the non-axisymmetries leading to LFQPOs are also possible, such as binary induced precession, spiral wave phase velocities that differ from Kelperian, or Lense-Thirring precession, The latter is an intrinsically general relativistic phenomenon. While relativistic effects should be considered in future more detailed applications to compact systems, here we focus on qualifying the effect of the key parameters in (2) on the modulation, not on the particular source of the the non-axisymmetry.
## 3 Simulated flux modulation
Table 1 shows a subset of our numerically solved cases for the blob and spiral. Cases 1-15 can be considered spirals, whilst cases 16-19 can be considered blobs because their structures extend less than $`2\pi `$ radians. For all cases we have taken an inclination angle of $`70^{}`$, motivated by the inferred observation angle of GRS 1915+105 (Mirabel & Rodr guez, 1994).
There is degeneracy in the relative influences of the various parameters of (2) on the RMS amplitude of modulation of Table 1, but it is instructive to discuss the influence of each of the parameters separately.
### 3.1 Spirals and physical interpretation of the geometric parameters
First notice the lack of influence of the range of $`\varphi >2\pi `$ in Table 1: For each $`\varphi `$, most of the modulation comes from the most inward “bump” in $`h(r)`$. Allowing for more than 1 bump does not modify the RMS modulation significantly, as seen in cases $`\mathrm{\#}1,\mathrm{\#}2,\mathrm{\#}3`$. For the majority of the other cases we will therefore consider only the range of $`\varphi 2\pi `$.
The parameter $`\stackrel{~}{\gamma }`$, measures the maximum thickness of the perturbation. Its effect is evident in comparing cases $`\mathrm{\#}3,\mathrm{\#}6,\mathrm{\#}7`$. The greater $`\stackrel{~}{\gamma }`$, all else being equal, the stronger the modulation. This is because a thicker structure more strongly shadows the inner disc. A change in thickness of $`70\%`$ means a temperature change of $`50\%`$. This somewhat extreme case is taken to illustrate the role of $`\stackrel{~}{\gamma }`$ on the RMS amplitude. Constraining the amplitude modulation from first principles requires a $`3`$D MHD disc simulation that exhibits a spiral wave instability. Here the non-axisymmetry would emerge rather than be imposed. Previous $`2.5`$-D (i.e. no dynamical vertical structure) disc simulations unstable to the AEI exhibit spiral waves, and lead to a temperature variation typically of order of $`2030\%`$ in the hydrostatic equilibrium approximation (Varnière et al, 2003). But in 3-D, hydrostatic equilibrium would underestimate the perturbation thickness of an AEI generated spiral because of additional heating from spiral shock dissipation. Thus a higher $`\stackrel{~}{\gamma }`$ than that obtained from the $`2.5`$D simulation can be expected, provided the AEI survives in 3-D.
Because $`\stackrel{~}{\gamma }\gamma h_o(r_c)`$ is the maximum height of the spiral arm, a change in $`h_o`$ (the zeroth order disc thickness) also increases the modulation. By comparing cases $`\mathrm{\#}14`$ and $`\mathrm{\#}15`$ we see that changing $`h_o`$ by a factor of two is not exactly the same as changing $`\stackrel{~}{\gamma }`$ by a factor of two because increasing $`\stackrel{~}{\gamma }`$ also increases the difference between the maximum height of the perturbation and $`h_o`$.
As discussed at the end of Sec. 2, the parameter $`\beta `$ measures how fast the maximum height along the spiral decreases with radius. Its effect is revealed in cases $`\mathrm{\#}3,\mathrm{\#}4,\mathrm{\#}5`$. The RMS amplitude of the modulation increases with $`\beta `$ for an $`\alpha >0`$ spiral. A rapidly decreasing height perturbation with radius means that the height also rapidly decreases along the spiral. This causes a stronger modulation by producing a stronger azimuthal variation. This is particularly important for tight spirals (small $`\alpha 0.05`$): were it not for a large $`\beta `$, little non-axisymmetry would otherwise arise.
The simulations relevant for studying the influence of $`r_c`$, the initial radius where the perturbation begins $`(r_c)`$, are $`\mathrm{\#}6`$, $`\mathrm{\#}9`$, $`\mathrm{\#}10`$ and $`\mathrm{\#}11`$. We see a maximum in the RMS amplitude for $`r_c`$ around $`2r_{in}`$, where $`r_{in}`$ is the inner disc radius: When the spiral is too near $`r_{in},`$ it obscures less of the inner disc which is the most luminous part. On the other hand, if the spiral is too far out, it can only obscure a less luminous outer region giving a smaller modulation. The combination of these effects is an intermediate $`r_c`$ for maximal modulation. Since the position of the spiral is chosen based on the LFQPO frequency, a correlation between the RMS amplitude and the frequency of the QPO is expected. In practice this may be hard to disentangle from variations in the other parameters.
The influence of the parameter $`\alpha `$, which determines the opening of the spiral wave can be seen from simulations $`\mathrm{\#}6`$, $`\mathrm{\#}12`$ and $`\mathrm{\#}13`$. There we Ssee that the more open the spiral, the higher the RMS amplitude. This is because a more open spiral means a more non-axisymmetric thickness profile, leading to a higher RMS amplitude of modulation.
Finally, consider the parameter $`\delta `$, which measures the width of the spiral or blob at its base. From cases $`\mathrm{\#}6`$ and $`\mathrm{\#}8`$ we see that a larger $`\delta `$, implies a larger RMS amplitude. For our choice of a highly inclined system, this trend results because a larger $`\delta `$ makes more of the region inner to the peak of the perturbation more perpendicular to the line of sight. This produces a larger observed flux when the observer is looking at the disk from an azimuth for which the line of sight intersects the inner part of the perturbation. But as the disk rotates, the outer edge of the perturbation comes into view, and the shadowing of the inner region occurs similarly for large or small $`\delta `$. The contrast in flux (and thus the RMS amplitude of the modulation) is thus larger for larger $`\delta `$, explaining the trend.
### 3.2 Blobs or hot spots
Because the difference between spirals and blobs in our model is just whether the structure extends for a range of $`\varphi 2\pi `$ (spiral) or $`\varphi <2\pi `$ (blob), the parameter influences are similar for the two cases. The relevant blob simulations are $`\mathrm{\#}16`$, $`\mathrm{\#}17`$, $`\mathrm{\#}18`$, and $`\mathrm{\#}19`$. There we see that a blob can create a non-negligible modulation of the X-ray flux but to reach a large RMS amplitude, the blob needs a large azimuthal extent–making it more banana or spiral shaped. The reason is that a localized structure does not provide much shadowing over a disc rotation period and thus offers only a weak RMS modulation. In contrast to spirals, such blobs probably cannot account for observed LFQPOs in microquasars.
## 4 Expected dependence of RMS amplitude on inclination and energy
Because the modulation comes from shadowing, the disc inclination angle $`\theta `$ is important in determining the maximum shadowing and thus the maximum RMS amplitude that a given choice of parameters can produce. Motivated by a comparison with the LFQPOs of GRS1915+105 (see e.g. the review Mclintock & Remillard 2004), we again focus on the parameter choices of simulation $`\mathrm{\#}8`$ and vary the inclination angle to obtain different values of the RMS amplitude. Fig. 2 shows the result. As expected, the more edge-on the view, the higher the RMS amplitude. Present observational data are insufficient to definitively confirm or contradict the predicted behavior. Such a trend could explain the extremely weak LFQPO in Cyg X-1, as the RMS amplitude expected from the inferred inclination angle is very small<sup>1</sup><sup>1</sup>1Cyg X-1 seems to show a weaker and broader structure than the “usual” LFQPO that could be a weak LFQPO in aggrement with this simple model.. More objects with a wider range of inclination angles are needed. Additional techniques of determining the disc inclination angle in microquasars besides using the jet propagation direction would also be desirable. In the case of neutron stars, data from eclipsing binaries and dippers provides good information on the binary and thus disk inclinations (e.g Frank et al., 1987).
It is important to note that the dependence of the rms amplitude on the inclination angle shown in Fig. 2 is displayed only for the case in which all other parameters are identical. Therefore the trend shown should be seen inside the same “family” of objects–those having similar disk properties. It was shown in table 1 that some parameters can have a greater influence on the rms amplitude than the inclination. These paramaters could give rise to low inclination source with a strong spiral having a higher rms amplitude than a high inclination source with a very small spiral
Using the relation between $`T`$ and $`h`$, we can also study how the RMS amplitude behaves as function of energy. Several studies show that LFQPO amplitudes rise with energy up to 15-20keV, although the behaviour at higher energies seems to depend on the source, for example in GRS $`1915`$+$`105`$ (Tomsick & Kaaret, 2001). We have not included a corona (high energy component) in our simple spectral modeling herein, but we can make a qualitative prediction of the LFQPO amplitude vs. energy trend by studying the RMS amplitude as function of our predicted disc flux, supplemented by knowledge of the disc/corona flux ratio taken from observation. We find that the RMS amplitude of the modulation gets bigger at higher energy. his is simply because a spiral with $`r_c/r_{in}2`$ is hotter than the surrounding disc and therefore its blackbody contribution to the spectrum peaks at a higher temperature. When a corona is present, this trend would apply at low energies, but we expect a critical energy above which the corona strongly dominates and the trend reverses.
We emphasize that in order to predict the exact position of the modification in the trend at high energies, we need to take into account not only the disk but also the corona self-consistently. This is beyond the scope of the present paper which is intended to show the basics of how flux can be modulated by non-axisymmetric structures.
## 5 Summary and Discussion
We have presented a simple analytical framework for modeling and interpreting how non-axisymmetric disk features such as spirals or blobs can create modulation of the observed flux from an accretion disk. We connect the disk geometry to the flux via the hydrostatic equilibrium assumption. Our focus on spirals is motivated by their production in disc simulations of the AEI and the efficacy with which such spirals seem to account for the observed LFQPO properties of microquasars, see e.g. the review McClintock & Remillard, 2004. Our formalism does not depend on the origin of the non-axisymmetric structures. Rather, it is used to quantify the influence of a range of generic spirals and hotspot geometries.
Using our formalism (with parameters guided by a fit to the LFQPO properties of GRS1915+105) we expect that the RMS amplitude increases with disc inclination because the modulation from a spiral or blob results from shadowing of the inner disc. Because the spiral is the hottest part of the disc for typical $`r_c2r_{in}`$, we also find that the RMS amplitude should increase with energy, except above energies where the emission becomes strongly corona dominated. More quantitative modeling and more observations are needed.
If the time scale of a LFQPO in a black holes system is indeed determined by an orbit or pattern speed of a non-axisymmetric feature, then at a given number of gravitational radii, the LFQPO frequency should scale inversely with the black hole mass. If accretion disks around AGN and microquasars incur similar instabilities, then a $`110`$Hz LFQPO in microquasars would correspond to a $`<10^6`$Hz LFQPO in AGN. Testing this prediction is presently difficult because it pushes the present limits of continuous observation times for individual AGN.
Finally note that although we have focused on black hole systems because of the LFQPO observations in microquasars, the analytic formalism for modeling the influence of spirals and hot spots herein does not depend on the compactness of the central object. It can also be applied to discs around neutron stars, white dwarfs, or young stellar objects.
More work is needed to incorporate our basic framework into detailed models of individual LFQPO sourceds.
Acknowledgments: We acknowledge support from NSF grant AST-0406799 and NASA grant ATP04-0000-0016. PV thanks M. Tagger for discussions and we thank M. Muno for discussions and allowing use of some code.
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# 1. Introduction
## 1. Introduction
This is an expanded version of a course given by the author at the University of Chicago in Winter 2003. A preliminary draft of lecture notes was prepared by Daniel Hoyt at the time of the course. The present version is about twice as large as the original notes and it contains a lot of additional material. However, the reader should be warned at the outset that the version at hand is still a very rough draft which is by no means complete. I decided to make the text public ‘as is’ since there is a real danger that a more perfect and more complete version will not appear in a foreseeable future.
In these lectures we will not attempt to present a systematic treatment of noncommutative geometry since we don’t think such a theory presently exists. Instead, we will try to convey an (almost random) list of beautiful concrete examples and general guiding principles which seem certain to be part of any future theory, even though we don’t know what that theory is going to be. Along the way, we will try to formulate many open questions and problems.
To avoid misunderstanding, we should caution the reader that the name ‘noncommutative geometry’ is quite ambiguous; different people attach to it different meanings. Typically, by noncommutative (affine) algebraic geometry one understands studying noncommutative algebras from the point of view of their similarity to coordinate rings of affine algebraic varieties. More generally, (a not necessarily affine) noncommutative geometry studies (some interesting) abelian, resp. triangulated, categories which share some properties of the abelian category of coherent sheaves on a (not necessarily affine) scheme, resp. the corresponding derived category.
It is important to make a distinction between what may be called noncommutative geometry ‘in the small’, and noncommutative geometry ‘in the large’. The former is a generalization of the conventional ‘commutative’ algebraic geometry to the noncommutative world. The objects that one studies here should be thought of as noncommutative deformations, sometimes referred to as quantizations, of their commutative counterparts. A typical example of this approach is the way of thinking about the universal enveloping algebra of a finite dimensional Lie algebra $`𝔤`$ as a deformation of the symmetric algebra $`S(𝔤)`$, which is isomorphic to the polynomial algebra.
As opposed to the noncommutative geometry ‘in the small’, noncommutative geometry ‘in the large’ is not a generalization of commutative theory. The world of noncommutative geometry ‘in the large’ does not contain commutative world as a special case, but is only similar, parallel, to it. The concepts and results that one develops here, do not specialize to their commutative analogues. Consider for instance the notion of smoothness that exists both in commutative algebraic geometry and in noncommutative algebraic geometry ‘in the large’. A commutative algebra $`A`$ may be smooth in the sense of commutative algebraic geometry, and at the same time be non-smooth from the point of view of noncommutative geometry ‘in the large’. An explantation of this phenomenon comes from operad theory, see e.g. \[MSS\], \[GiK\], \[Ka1\]. Each of the mathematical worlds that we study is governed by an appropriate operad. Commutative geometry is governed by the operad of commutative (associative) algebras, while noncommutative geometry ‘in the large’ is governed by the operad of associative not necessarily commutative algebras. In this sense, it would be more appropriate to speak of ‘associative geometry’ instead of what we call noncommutative geometry ‘in the large’. There are other geometries arising from operads of Lie algebras, Poisson algebras, etc.
Many interesting and important topics of noncommutative geometry are completely left out in these notes. For example, I have not discussed noncommutative projective geometry at all. The interested reader is referred to \[St\] and \[SvB\] for an excellent reviews. Some additional references are given in the bibliography at the end of the lectures.
### Basic notation.
Throughout this text we fix $`\mathrm{𝕜}`$, an algebraically closed field of characteristic zero, which may be assumed without loss of generality to be the field of complex numbers. By an algebra we always mean an associative, not necessarily commutative, unital $`\mathrm{𝕜}`$-algebra. If $`A`$ is an algebra, we denote by $`A\text{-}\mathrm{𝗆𝗈𝖽}`$, $`\mathrm{𝗆𝗈𝖽}\text{-}A`$ and $`A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}`$ the categories of left $`A`$-modules, right $`A`$-modules and $`A`$-bimodules, respectively. We write $`=_\mathrm{𝕜}`$.
### Acknowledgments.
I would like to thank D. Boyarchenko for providing proofs of several results claimed in the course without proof and to D. Hoyt for making preliminary notes of the course.
I am very much indebted to Maxim Kontsevich for generously sharing with me his unpublished ideas. I have benefited from many useful discussions with Bill Crawley-Boevey, Pavel Etingof, Mikhail Kapranov, Toby Stafford, Boris Tsygan, and Michel Van den Bergh.
This work was partially supported by the NSF.
## 2. Morita Equivalence
### 2.1. Categories and functors.
We remind the reader some basic concepts involving categories and functors.
Recall that a category where $`\mathrm{Hom}`$-spaces are abelian groups and such that the notion of direct sum (also called coproduct) of a family of objects is defined is called an additive category.
A functor $`F:𝒞_1𝒞_2`$ is *fully faithful* if $`F`$ is an isomorphism on every set of morphisms, and that $`F`$ is *essentially surjective* if for every object $`X𝒞_2`$, there is some $`Y𝒞_1`$ such that $`X`$ and $`F(Y)`$ are isomorphic.
The most commonly used way to establish an equivalence of categories is provided by the following
###### Lemma 2.1.1.
Let $`𝒞_1`$ and $`𝒞_2`$ be two abelian categories, and let $`F:𝒞_1𝒞_2`$ be an exact, fully faithful, essentially surjective functor. Then $`F`$ is an equivalence of categories. ∎
Let $`\mathrm{𝖲𝖾𝗍𝗌}`$ be the category of sets. For any category $`𝒞`$, functors from $`𝒞`$ to $`\mathrm{𝖲𝖾𝗍𝗌}`$ form a category $`\mathrm{Fun}(𝒞,\mathrm{𝖲𝖾𝗍𝗌})`$. Now, any object $`X𝒞`$, gives rise to the functor $`\mathrm{Hom}_𝒞(X,):𝒞\mathrm{𝖲𝖾𝗍𝗌}.`$ It is straightforward to check that the assignment $`X\mathrm{Hom}_𝒞(X,)`$ extends to a natural contravariant functor $`𝒞\mathrm{Fun}(𝒞,\mathrm{𝖲𝖾𝗍𝗌})`$.
###### Lemma 2.1.2 (Yoneda lemma).
The functor $`𝒞\mathrm{Fun}(𝒞,\mathrm{𝖲𝖾𝗍𝗌})`$ induces isomorphisms on $`\mathrm{Hom}`$’s, in other words, it is a fully faithful functor.
Write $`\mathrm{𝖠𝖻}`$ for the category of abelian groups.
###### Definition 2.1.3.
An object $`P`$ of an abelian category $`𝒞`$ is said to be *projective* if the functor $`\mathrm{Hom}_𝒞(P,):𝒞\mathrm{𝖠𝖻}`$ is exact.
In other words, $`P`$ is projective if given a short exact sequence
$$0M^{}MM^{\prime \prime }0$$
in $`𝒞`$, we have that
$$0\mathrm{Hom}_𝒞(M^{},P)\mathrm{Hom}_𝒞(M,P)\mathrm{Hom}_𝒞(M^{\prime \prime },P)0$$
is exact in $`\mathrm{𝖠𝖻}`$. An object $`G`$ of $`𝒞`$ is called a *generator* if $`\mathrm{Hom}_𝒞(G,A)`$ is nonzero for every nonzero object $`A`$ of $`𝒞`$.
###### Definition 2.1.4.
Let $`𝒞`$ be an additive category with arbitrary direct sums (also referred to as coproduct). An object $`X`$ of $`𝒞`$ is called *compact* if, for an arbitrary set of objects of $`𝒞`$ and a morphism $`f:X_{\alpha I}M_\alpha ,`$ there exists some finite set $`FI`$ such that $`\mathrm{Im}f`$ is a subobject of $`_{\alpha F}M_\alpha `$.
An easy consequence of the definition of compactness is the following
###### Lemma 2.1.5.
An object $`X`$ in an abelian category $`𝒞`$ (with arbitrary direct sums) is compact if and only if the functor $`\mathrm{Hom}_𝒞(X,)`$ commutes with arbitrary direct sums, that is
$$\mathrm{Hom}_𝒞(X,_{\alpha \mathrm{\Lambda }}Y_\alpha )=_{\alpha \mathrm{\Lambda }}\mathrm{Hom}_𝒞(X,Y_\alpha ).\mathrm{}$$
###### Lemma 2.1.6.
Let $`A`$ be a ring and $`M`$ an $`A`$-module.
$`(𝗂)`$If $`M`$ is a finitely generated $`A`$-module, then $`M`$ is a compact object of $`A\text{-}\mathrm{𝗆𝗈𝖽}`$.
$`(\mathrm{𝗂𝗂})`$If $`M`$ is projective and is a compact object of $`A\text{-}\mathrm{𝗆𝗈𝖽}`$, then $`M`$ is finitely generated.
###### Proof.
$`(𝗂)`$is obvious. To prove (ii), assume $`M`$ is projective, and choose any surjection $`p:A^IM`$, where $`I`$ is a possibly infinite set. There exists a section $`s:MA^I`$. If $`M`$ is compact, the image of $`s`$ must lie in a submodule $`A^JA^I`$ for some finite subset $`JI`$. Then $`p|_{A^J}`$ is still surjective, which shows that $`M`$ is finitely generated. ∎
The following result provides a very useful criterion for an abelian category to be equivalent to the category of left modules over a ring.
###### Proposition 2.1.7.
Let $`𝒞`$ be an abelian category with arbitrary direct sums. Let $`P𝒞`$ be a compact projective generator and set $`B=(\mathrm{End}_𝒞P)^{\mathrm{op}}`$. Then the functor $`\mathrm{Hom}_𝒞(P,)`$ yields an equivalence of categories between $`𝒞`$ and $`B\text{-}\mathrm{𝗆𝗈𝖽}`$.
###### Proof of Proposition 2.1.7.
We will show that $`F(X)=\mathrm{Hom}_𝒞(P,X)`$ is fully faithful and then apply Lemma 2.1.1.
We wish to show that there is an identification between $`\mathrm{Hom}_𝒞(X,M)`$ and $`\mathrm{Hom}_{B\text{-}\mathrm{𝗆𝗈𝖽}}(F(X),F(M))`$ for all $`M𝒞`$. Since $`P`$ is a generator, we deduce that $`\mathrm{Hom}_𝒞(P,M)0`$. Define
$$\phi :\underset{f\mathrm{Hom}_𝒞(P,M)}{}PM$$
by $`\phi (p_f):=_{f\mathrm{Hom}_𝒞(P,M)}f(p_f)`$ (this sum is finite). Let $`L=\mathrm{Im}\phi `$. Then $`L`$ is a submodule of $`M`$, and if it is not all of $`M`$ then $`M/L`$ is nonzero, hence there is some nonzero map $`PM/L`$. But this map lifts to a map to $`M`$, and the image of this lift must include points not in $`L`$, which is a contradiction. So, every $`MA\text{-}\mathrm{𝗆𝗈𝖽}`$ can be written as a quotient of $`P^T`$ for some cardinal $`T`$. Let $`K`$ denote the kernel, and take $`P^S`$ surjecting onto $`K`$. Then composing this with the inclusion of $`K`$ into $`P^T`$ yields the exact sequence
$$P^SP^TM0.$$
Since $`P`$ is projective, $`\mathrm{Hom}_𝒞(P,)`$ is exact. Hence,
$$\mathrm{Hom}_𝒞(P,P^S)\mathrm{Hom}_𝒞(P,P^T)\mathrm{Hom}_𝒞(P,M)0$$
is exact. Since $`P`$ is finitely generated (i.e., compact), it commutes with arbitrary direct sums (see Lemma 2.1.5. So, using exactness of the above sequence it suffices to check that $`\mathrm{Hom}_𝒞(X,P)=\mathrm{Hom}(F(X),F(P))`$. Since $`F(P)=\mathrm{End}_AP=B^{\mathrm{op}}`$, this is automatic. ∎
We need two more definitions.
###### Definition 2.1.8.
Let $`𝒞_1`$ and $`𝒞_2`$ be two categories, and let $`F,G:𝒞_1𝒞_2`$ be two functors. A *morphism* $`\varphi :FG`$ is a natural transformation, i.e., a collection of morphisms
$$\varphi _X:F(X)G(X)$$
for each $`X𝒞_1`$ such that for any morphism $`f:XY`$ ($`Y𝒞_1`$), the following diagram commutes:
$$\text{}.$$
In particular, if $`𝒞`$ is an abelian category, we have the identity morphism $`\mathrm{id}_𝒞:𝒞𝒞`$. We define the *center* $`𝖹(𝒞)=\mathrm{End}(\mathrm{id}_𝒞)`$.
###### Example 2.1.9.
It is a worthwhile exercise to check that $`𝖹(\mathrm{Coh}X)𝒪(X)`$ for any algebraic variety (where $`𝒪(X)`$ is the ring of global regular functions on $`X`$). $`\mathrm{}`$
###### Lemma 2.1.10.
Let $`A`$ be a associative algebra. Then $`𝖹(A\text{-}\mathrm{𝗆𝗈𝖽})=𝖹_A`$.
###### Proof.
Choose an element $`z𝖹_A`$. Define an endomorphism of $`\mathrm{id}_𝒞`$ by setting $`\varphi _X`$ to be the action of $`z`$ on the module $`X`$. Since $`z`$ is central, this is a module homomorphism and it commutes with all module maps, that is, the diagram
$$\text{}.$$
Conversely, suppose an endomorphism $`\varphi :\mathrm{id}_𝒞\mathrm{id}_𝒞`$ is given. Set $`z=\varphi _A(1_A)`$. We need to check that this is indeed central. Choose any $`aA`$, and define the left $`A`$-module map $`f:AA`$ by $`f(x)=xa`$. Then since $`\varphi `$ is a morphism, we know that $`f\varphi _A=\varphi _Af`$. So,
$`za`$ $`=\varphi _A(1_A)a=(f\varphi _A)(1_A)`$
$`=(\varphi _Af)(1_A)=\varphi _A(a)=a\varphi _A(1_A)=az.`$
So, $`z𝖹_A`$. It is clear that this association is an algebra homomorphism. ∎
### 2.2. Algebras and spaces.
One of the cornerstones of geometry is the equivalence of categories of spaces and categories of algebras. For example, the Gelfand theorem asserts an (anti)-equivalence between the category of locally compact Hausdorff spaces with proper maps and the category of commutative $`C^{}`$-algebras. Similarly, in algebra, one has an (anti)-equivalence between the category of affine algebraic varieties and the category of finitely generated commutative algebras without nilpotent elements. A first step to such an equivalence in algebraic geometry is by associating to each scheme $`X`$ its structure sheaf $`𝒪_X`$. However, this is unsatisfactory since $`𝒪_X`$ explicitly refers to the space $`X`$, since it is a sheaf on $`X.`$
One way to resolve this difficulty is to forget about scheme $`X`$ altogether, and to consider instead the abelian category $`\mathrm{Coh}(X)`$ of coherent sheaves over $`X`$. The space $`X`$ can be reconstructed in a natural way from this category.
A first step in considering such abelian categories is to look at the category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ of left modules over an associative algebra $`A`$. One would hope that the category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ should uniquely determine $`A`$ up to isomorphism. It turns out that this is the case for commutative associative algebras, but not for arbitrary associative algebras.
###### Definition 2.2.1.
Let $`A`$ and $`B`$ be associative, not necessarily commutative, algebras. Then we say that $`A`$ and $`B`$ are *Morita equivalent*, if there is an equivalence of categories between $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ and $`B\text{-}\mathrm{𝗆𝗈𝖽}`$.
Morita equivalence of two commutative algebras is particularly simple.
###### Proposition 2.2.2.
Commutative algebras $`A`$ and $`B`$ are Morita equivalent if and only if they are isomorphic.
###### Proof of Proposition 2.2.2.
Clearly isomorphic algebras are Morita equivalent. Now suppose that $`A`$ and $`B`$ are Morita equivalent associative algebras. Then $`A\text{-}\mathrm{𝗆𝗈𝖽}B\text{-}\mathrm{𝗆𝗈𝖽}`$, so certainly we have that $`𝖹(A\text{-}\mathrm{𝗆𝗈𝖽})𝖹(B\text{-}\mathrm{𝗆𝗈𝖽})`$. If $`A`$ and $`B`$ are both commutative, then by Lemma 2.1.10 we have $`A=𝖹_A`$ and $`B=𝖹_B`$. ∎
###### Example 2.2.3.
Given an algebra $`A`$ and an integer $`n1`$, let $`\mathrm{Mat}_n(A)`$, be the algebra of $`n\times n`$-matrices with entries in $`A`$. A typical example of Morita equivalence involving noncommutative algebras is provided by the following easy result.
###### Lemma 2.2.4.
For any algebra $`A`$ and any integer $`n1`$, the algebras $`A`$ and $`\mathrm{Mat}_n(A)`$ are Morita equivalent.
Note that even if $`A`$ is a commutative algebra, the algebra $`\mathrm{Mat}_n(A)`$ is not commutative for any $`n>1`$.
###### Proof.
Let $`e_{11}\mathrm{Mat}_n(A)`$ be the matrix with entry $`1`$ at the $`1\times 1`$ spot, and zeros elsewhere. It is clear that the algebra $`e_{11}\mathrm{Mat}_n(A)e_{11}`$ is isomorphic to $`A`$. On the other hand, it is clear that multiplying the matrices with arbitrary entry $`aA`$ at the $`1\times 1`$ spot, and zeros elsewhere by other matrices from $`\mathrm{Mat}_n(\mathrm{𝕜})\mathrm{Mat}_n(A)`$, one can obtain, taking linear combinations, every $`A`$-valued $`n\times n`$-matrix. Thus, we have shown that $`\mathrm{Mat}_n(A)=\mathrm{Mat}_n(A)e_{11}\mathrm{Mat}_n(A)`$ and that $`Ae_{11}\mathrm{Mat}_n(A)e_{11}`$. At this point, the result follows from the general criterion of Corollary 2.3.4 below. ∎
Thus, in general, the algebra $`A`$ can not be recovered from the corresponding abelian category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$. Therefore, in order for a concept in noncommutative geometry to have an intrinsic meaning, that concept must be Morita invariant. In particular, the question of which properties of an algebra are Morita invariant becomes very important.
### 2.3. Morita theorem.
The main result about Morita equivalent algebras is provided by the following<sup>1</sup><sup>1</sup>1The exposition below follows the notes prepared by M. Boyarchenko.
###### Theorem 2.3.1.
Let $`A`$ and $`B`$ be two rings, and $`F:A\text{-}\mathrm{𝗆𝗈𝖽}B\text{-}\mathrm{𝗆𝗈𝖽}`$ an additive right exact functor. Then there exists a $`(B,A)`$-bimodule $`Q`$, unique up to isomorphism, such that $`F`$ is isomorphic to the functor
$$A\text{-}\mathrm{𝗆𝗈𝖽}B\text{-}\mathrm{𝗆𝗈𝖽},MQ_AM.$$
###### Proof.
The uniqueness of $`Q`$ (if it exists) is clear, since we have $`Q=Q_AA=F(A)`$. To prove existence, let $`Q=F(A)`$; by assumption, this is a left $`B`$-module. Moreover, for every $`aA`$, the operator $`\rho _a`$ of right multiplication by $`a`$ is an endomorphism of $`A`$ as a left $`A`$-module, whence we obtain a ring homomorphism $`A^{op}\mathrm{End}_B(Q)`$, $`aF(\rho _a)`$. This homomorphism makes $`Q`$ into a $`(B,A)`$-bimodule.
Now, for every $`MA\text{-}\mathrm{𝗆𝗈𝖽}`$, we define, functorially, a $`B`$-module homomorphism $`Q_AMF(M)`$. Let us first define a $``$-bilinear map $`\varphi _M:Q\times MF(M)`$. An element $`mM`$ gives rise to an $`A`$-module homomorphism $`\rho _m:AM`$, $`aam`$. We define $`\varphi _M(q,m)=F(\rho _m)(q)`$. Now since $`F(\rho _m)`$ is a $`B`$-module homomorphism, the map $`\varphi _M`$ commutes with left multiplication by elements of $`B`$. Also, if $`aA`$, then
$$\varphi _M(qa,m)=F(\rho _m)\left(F(\rho _a)(q)\right)=F(\rho _m\rho _a)(q)=F(\rho _{am})(q)=\varphi _M(a,qm),$$
whence $`\varphi _M`$ descends to a left $`B`$-module homomorphism $`\psi _M:Q_AMF(M)`$.
It is obvious that $`\psi _M`$ is functorial with respect to $`M`$. Moreover, by construction, $`\psi _M`$ is an isomorphism whenever $`M`$ is free. In general, we use the fact that both $`F`$ and the functor $`Q_A`$ are exact. Given any left $`A`$-module $`M`$, choose an exact sequence $`F_1F_0M0`$, where $`F_0`$ and $`F_1`$ are free $`A`$-modules, and apply both functors to this sequence. Using the morphism of functors, we get a commutative diagram, and the Five Lemma finishes the proof. ∎
###### Corollary 2.3.2.
Two rings, $`A`$ and $`B`$, are Morita equivalent if and only if there exist an $`(A,B)`$-bimodule $`P`$ and a $`(B,A)`$-bimodule $`Q`$ such that $`P_BQA`$ as $`A`$-bimodules and $`Q_APB`$ as $`B`$-bimodules. Under this assumption, we have
$$\mathrm{End}_{A\text{-}\mathrm{𝗆𝗈𝖽}}(P)=B^{op},\mathrm{op}and\mathrm{End}_{B\text{-}\mathrm{𝗆𝗈𝖽}}(Q)=A^{op}.$$
Moreover, $`P`$ is projective as an $`A`$-module and $`Q`$ is projective as a $`B`$-module.
###### Proof.
Equivalences of categories are exact functors and preserve projective objects. ∎
###### Corollary 2.3.3.
If $`A`$ and $`B`$ are Morita equivalent rings, then the categories $`\mathrm{𝗆𝗈𝖽}\text{-}A`$ and $`\mathrm{𝗆𝗈𝖽}\text{-}B`$ are also equivalent. Moreover, there is a natural equivalence of categories $`A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}B\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}`$ which takes $`A`$ to $`B`$ (with their natural bimodule structures).
Proof. Let $`P`$ and $`Q`$ be as in the previous corollary. For the first statement, use the functors $`_AP`$ and $`_BQ`$. For the second statement, use the functor
$$Q_A_AP:A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}B\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}.\mathrm{}$$
Let $`A`$ be a ring and $`e=e^2A`$ an idempotent. Clearly, $`eAe`$ is a subring in $`A`$ and $`Ae`$ is naturally an $`(A,eAe)`$-bimodule and $`eA`$ is an $`(eAe,A)`$-bimodule. Note that $`e`$ is the unit of the ring $`A`$. We see that the inclusion map $`eAeA`$ is usually not a ring homomorphism, and $`A`$ is not an $`eAe`$-module in a natural way.
For any left $`A`$-module $`M`$, the space $`eM`$ has a natural left $`eAe`$-module structure.
###### Corollary 2.3.4.
Let $`A`$ be a ring and $`eA`$ an idempotent. The functor $`A\text{-}\mathrm{𝗆𝗈𝖽}B\text{-}\mathrm{𝗆𝗈𝖽},MeM`$ is a Morita equivalence if and only if $`AeA=A`$. In this case, an inverse equivalence is given by the functor $`NAe_{eAe}N.`$
###### Remark 2.3.5.
Observe also that there exist rings $`R`$ such that $`R\times RR`$ (for example, take $`R`$ to be the product of an infinite number of copies of some fixed ring). Now for such a ring, let $`A=R\times R`$ and $`e=(1,0)A`$. Then $`eAe=RA`$, and in particular, $`eAe`$ is Morita equivalent to $`A`$, but we have $`AeA=R\times \{0\}A`$.
###### Proof.
Assume that $`AeA=A`$. We claim that $`eA_AAeeAe`$ as $`eAe`$-bimodules (in fact, this holds regardless of the assumption on $`e`$). Note that $`eA`$ is a direct summand of $`A`$ as a right $`A`$-module: $`A=eA(1e)A`$. Similarly, $`A=AeA(1e)`$. Now the multiplication map $`A_AAA`$ is obviously an $`A`$-bimodule isomorphism, and it clearly restricts to an isomorphism $`eA_AAe\stackrel{}{}eAe`$ of $`eAe`$-bimodules.
Now we also claim that $`Ae_{eAe}eAA`$ as $`A`$-bimodules. We have the natural $`A`$-bimodule map $`Ae_{eAe}eA\stackrel{m}{}A`$ given by multiplication. We can write down an explicit inverse for this map. Namely, since $`AeA=A`$, there exist elements $`a_j,b_jA`$ such that $`1=_ja_jeb_j`$. We define a map of abelian groups $`f:AAe_{eAe}eA`$ by $`f(a)=_jaa_jeeb_j`$. It is obvious that $`mf=\mathrm{id}_A`$. We must check that $`fm`$ is also the identity. We have
$`(fm)`$ $`\left({\displaystyle \underset{k}{}}c_keed_k\right)=f\left({\displaystyle \underset{k}{}}c_ked_k\right)={\displaystyle \underset{j,k}{}}(c_ked_ka_je)eb_j`$
$`={\displaystyle \underset{j,k}{}}c_ke(ed_ka_jeeb_j)={\displaystyle \underset{k}{}}c_ke\left[ed_k{\displaystyle \underset{j}{}}a_jeb_j\right]={\displaystyle \underset{k}{}}c_keed_k.`$
We leave the rest of the proof to the reader. ∎
###### Remark 2.3.6.
Let us write $`A\text{-}\mathrm{𝗆𝗈𝖽}_f`$ for the category of finitely generated $`A`$-modules. In the situation of Theorem 2.3.1, it is clear that the functor $`F`$ takes $`A\text{-}\mathrm{𝗆𝗈𝖽}_f`$ into $`B\text{-}\mathrm{𝗆𝗈𝖽}_f`$ if and only if $`Q`$ is finitely generated as a $`B`$-module. We claim that this is always the case whenever $`F`$ is an equivalence; in particular, if $`A`$ and $`B`$ are Morita equivalent rings, then the categories $`A\text{-}\mathrm{𝗆𝗈𝖽}_f`$ and $`B\text{-}\mathrm{𝗆𝗈𝖽}_f`$ are equivalent. The proof is based on the notion of a compact object. In particular, if $`F:A\text{-}\mathrm{𝗆𝗈𝖽}B\text{-}\mathrm{𝗆𝗈𝖽}`$ is an equivalence of abelian categories, then $`F(A)`$ must be a compact object of $`B\text{-}\mathrm{𝗆𝗈𝖽}`$, but it is also projective, whence finitely generated.
## 3. Derivations and Atiyah algebras
### 3.1.
We recall the definitions of derivations and super-derivations.
Let $`A`$ be an algebra, and let $`M`$ be an $`A`$-bimodule. A $`\mathrm{𝕜}`$-linear map $`\delta :AM`$ is called a *derivation* if it satisfies the Leibniz rule, that is, if
$$\delta (a_1a_2)=a_1\delta (a_2)+\delta (a_1)a_2\text{for all}a_1,a_2A.$$
We let $`\mathrm{Der}(A,M)`$ denote the $`\mathrm{𝕜}`$-vector space of all derivations from $`A`$ to $`M`$. For any $`mm`$, the map $`\mathrm{ad}m:MM,amaam`$ gives a derivation of $`M`$. The derivations of the form $`\mathrm{ad}m,mM,`$ are called inner derivations. We write $`\mathrm{Inn}(A,M)\mathrm{Der}(A,M)`$ for the space of inner derivations, and $`𝖹(M)=\{mMam=ma,aA\},`$ for the ‘center’ of $`A`$-bimodule $`M`$. Thus, we have the following exact sequence
$$0𝖹(M)M\stackrel{\mathrm{ad}}{}\mathrm{Der}(A,M)\mathrm{Der}(A,M)/\mathrm{Inn}(A,M)0.$$
(3.1.1)
###### Remark 3.1.2.
If the algebra $`A`$ is commutative, then any left $`A`$-module $`M`$ may be viewed as an $`A`$-bimodule, with right action being given by $`ma:=am,aA,mM`$ (this formula only gives a right $`A`$-module structure if $`A`$ is commutative). Bimodules of that type are called symmetric. Thus, given a left $`A`$-module viewed as a symmetric $`A`$-bimodule, one may consider derivations $`AM`$.
In the special case that $`M=A`$, we abbreviate $`\mathrm{Der}(A,A)`$ to $`\mathrm{Der}(A)`$, resp. $`\mathrm{Inn}(A,A)`$ to $`\mathrm{Inn}(A)`$. It is an easy calculation that $`\mathrm{Der}(A)`$ is a Lie algebra under the commutator, and that inner derivations form a Lie ideal $`\mathrm{Inn}(A)`$ in $`\mathrm{Der}(A)`$. It is also straightforward to check that the assignment $`a\mathrm{ad}a`$ is a Lie algebra map (with respect to the commutator bracket on $`A`$), whose kernel is the center, $`𝖹_A`$, of the algebra $`A`$.
The inter-relationships between $`A,𝖹_A,`$ and $`\mathrm{Der}(A)`$ are summirized in the following result which says that $`\mathrm{Der}(A)`$ is a Lie algebroid on $`\mathrm{Spec}𝖹_A`$, see Sect. 6.5.
###### Proposition 3.1.3.
$`(𝗂)`$The Lie algebra $`\mathrm{Der}(A)`$ acts on $`A`$, and this action preserves the center $`𝖹_AA`$.
$`(\mathrm{𝗂𝗂})`$For any $`z𝖹_A`$ and $`\theta \mathrm{Der}(A)`$, the map $`z\theta :az\theta (a)`$ is again a derivation of $`A`$. The assignment $`\theta z\theta `$ makes $`\mathrm{Der}(A)`$ a $`𝖹_A`$-module.
$`(\mathrm{𝗂𝗂𝗂})`$For any $`\theta ,\delta \mathrm{Der}(A)`$ and $`z𝖹_A`$, one has
$$[z\theta ,\delta ]=z[\theta ,\delta ]\delta (z)\theta .\mathrm{}$$
###### Example 3.1.4.
Let $`A=\mathrm{𝕜}[X]`$, be the coordinate ring of a smooth affine variety $`X`$. For each $`xX`$, let $`𝔪_x=\{fAf(x)=0\}`$ be the corresponding maximal ideal in $`A`$. Then $`\mathrm{𝕜}_x=A/𝔪_x`$ is a 1-dimensional $`A`$-module (which we will also view as an $`A`$-bimodule).
Given $`xX`$, let $`T_xX`$ be the tangent space at $`x`$. For any tangent vector $`\xi T_xX`$, differentiating the function $`f`$ with respect to $`\xi `$ gives a map $`Af(\xi f)(x)=df(\xi )|_x\mathrm{𝕜}`$. This map is a derivation $`_\xi :A\mathrm{𝕜}_x.`$ It is an elementary result of commutative algebra that this way one gets an isomorphism
$$T_xX\stackrel{_{}}{}\mathrm{Der}(A,\mathrm{𝕜}),\xi _\xi .$$
On $`X`$, we have the tangent bundle $`TXX`$. We write $`𝒯_X`$ for the tangent sheaf, the sheaf of algebraic sections of the tangent bundle. This sheaf is locally free since $`X`$ is smooth, the sections of $`𝒯_X`$ are nothing but algebraic vector fields on $`X`$. Write $`𝒯(X):=\mathrm{\Gamma }(X,𝒯_X)`$ for the vector space of (globally defined) algebraic vector fields on $`X`$. Commutator of vector fields makes $`𝒯(X)`$ a Lie algebra.
For any algebraic vector field $`\xi `$ on $`X`$, the map $`\mathrm{𝕜}[X]f\xi f`$ gives a derivation of $`A=\mathrm{𝕜}[X]`$. This way one obtains a canonical Lie algebra isomorphism
$$𝒯(X)\stackrel{_{}}{}\mathrm{Der}(A),\text{where}A=\mathrm{𝕜}[X].$$
Note that the product of a function and a vector field is again a well-defined vector field, in accordance with part (ii) of Proposition 3.1.3. $`\mathrm{}`$
###### Remark 3.1.5.
Let $`A`$ be an associative algebra. A derivation $`\delta :AA`$ may be thought of, heuristically, as a generator of an ‘infinitesimal’ one-parameter group $`\epsilon \mathrm{exp}(\epsilon \delta )=\mathrm{id}_A+\epsilon \delta +\frac{1}{2}\epsilon ^2\delta {}_{^{^{}}}{}^{}\delta +\mathrm{},`$ of automorphisms of $`A`$. To formalize this, introduce the ring $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2,`$ called the ring of ‘dual numbers’, and form the tensor product algebra $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2A.`$ To any linear map $`f:AA`$ we associate the map
$$F:A\mathrm{𝕜}[\epsilon ]/\epsilon ^2A,aF(a):=a+\epsilon f(a).$$
Here, we think of $`\epsilon `$ as an ‘infinitesimally small’ parameter, so, up to higher powers of $`\epsilon `$, one has $`F=\mathrm{Id}+\epsilon f\mathrm{exp}(\epsilon f).`$ Thus, the map $`F`$ may be thought of as a family of maps $`AA`$ ‘infinitesimally close’ to the identity. Now the equation $`F(aa^{})=F(a)F(a^{}),a,a^{}A`$, saying that our family is a family of algebra homomorphisms, when expressed in terms of $`f`$, reads (recall that $`\epsilon ^2=0`$):
$$aa^{}+\epsilon f(aa^{})=aa^{}+\epsilon (af(a^{})+f(a)a^{})),a,a^{}A.$$
(3.1.6)
Equating the coefficients in front of $`\epsilon `$, we see that (3.1.6) reduces to the condition for $`f`$ to be a derivation, as promised.
Let $`\mathrm{Aut}(A)`$ denote the group of all (unit preserving) automorphisms of the algebra $`A`$. Thus, if we think of $`\mathrm{Aut}(A)`$ as some sort of Lie group, then the ‘Lie algebra’ of that group is given by
$$\mathrm{Lie}\mathrm{Aut}(A)=\mathrm{Der}(A).$$
Moreover, one can argue that the Lie bracket on the left-hand side of this formula corresponds to the commutator bracket on the space of derivations on the right. $`\mathrm{}`$
### 3.2. Square-zero construction.
We would like to extend the intuitive point of view explained in Remark 3.1.5 to a more general case where $`\theta \mathrm{Der}(A,M)`$ for an arbitrary $`A`$-bimodule $`M`$. This can be achieved by the following general construction.
Suppose we are considering some class of algebraic structure, be it associative algebras, Lie algebras, Poisson algebras, etc. We also wish to discuss modules over these algebras. There is a natural way of defining what the “correct” module structure is for a given type of algebra called the *square zero* construction. Suppose that $`A`$ is some sort of algebra over $`\mathrm{𝕜}`$ and that $`M`$ is a $`\mathrm{𝕜}`$-vector space. We wish to give $`M`$ the type of module structure appropriate to the structure of $`A`$. Consider the vector space $`AM`$. Then giving a correct “bimodule” structure on $`M`$ is equivalent to giving an algebra structure on $`AM`$ such that
1. the projection $`AMA`$ is an algebra map, and
2. $`M^2=0`$ and $`M`$ is an ideal.
It is easy to see this principle at work. If $`A`$ is an associative algebra and $`M`$ is an $`A`$-bimodule, then $`AM`$ is an associative algebra under the multiplication $`(am)(a^{}m^{}):=aa^{}(am^{}+ma^{})`$. Indeed, $`M^2=0`$, $`M`$ is an ideal, and the projection $`AMA`$ is an algebra map. This is a rather trivial case, however.
###### Lemma 3.2.1.
A linear map $`\theta :AM`$ is a derivation if and only if the map $`A\mathrm{}MA\mathrm{}M`$ given by $`(a,m)(a,m+\theta (a))`$ is an algebra automorphism.
Thus, for an arbitrary bimodule $`M`$ we may think of derivations $`\theta \mathrm{Der}(A,M)`$ as ‘infinitesimal automorphisms’ of the algebra $`A\mathrm{}M`$.
### 3.3. Super-derivations.
Now suppose that $`A`$ is $``$-graded, that is, there is a direct sum decomposition (as $`\mathrm{𝕜}`$-vector spaces) $`A=_iA_i`$ such that $`A_iA_jA_{i+j}`$ for all $`i,j`$. We put
$$A_{\mathrm{ev}}=_iA_{2i}\text{and}A_{\mathrm{odd}}=_iA_{2i+1}.$$
A linear map $`f:AA`$ is said to be even, resp., odd, if $`f(A_\pm )A_\pm ,`$ resp., $`f(A_\pm )A_{}`$, where the plus sign stands for ‘even’ and the minus sign stands for ‘odd’.
###### Definition 3.3.1.
An odd $`\mathrm{𝕜}`$-linear map $`f:AA`$ is called either a *super-derivation* or an *odd derivation*, if it satisfies the graded Leibniz rule, that is,
$$f(a_1a_2)=f(a_1)a_2+(1)^{\mathrm{deg}a_1}a_1f(a_2),$$
where $`a_1`$ is a homogeneous element of degree $`k=\mathrm{deg}a_1`$, that is, $`a_1A_k`$.
It is an easy calculation to check that the $`/(2)`$-graded vector space of all even/odd derivations forms a Lie super-algebra under the super-commutator:
$$[f,g]:=f{}_{^{^{}}}{}^{}gg{}_{^{^{}}}{}^{}f,$$
where the plus sign is taken if both $`f`$ and $`g`$ are odd, and the minus sign in all other cases. In particular, for an odd derivation $`f:AA`$, we see that
$$ff=\frac{1}{2}[f,f]$$
is an even derivation.
Below, we will frequently use the following result, which follows from the Leibniz formula by an obvious induction.
###### Lemma 3.3.2.
Let $`f,g:AA`$ be to derivations (both even, or odd) of an associative algebra $`A`$. Let $`SA`$ be a set of algebra generators for $`A`$. Then, we have
$$f(s)=g(s),sSf=g.\mathrm{}$$
The most commonly used application of the Lemma is the following
###### Corollary 3.3.3.
Let $`d:AA`$ be an odd derivation of a graded algebra $`A`$. If $`SA`$ is a set of algebra generators for $`A`$, and $`d^2(s)=0`$, for any $`sS`$, then $`d^2=0`$.∎
### 3.4. The tensor algebra of a bimodule.
Let $`A`$ be an associative algebra. Given two $`A`$-bimodules $`M`$ and $`N`$ one has a well-defined $`A`$-bimodule $`M_AN`$. In particular, for an $`A`$-bimodule $`M`$ we put $`T_A^nM:=M_AM_A\mathrm{}_AM`$ ($`n`$ times), in particular, $`T_A^1M=M`$, and we put formally $`T_A^0M:=A`$. The direct sum
$$T_A^{\text{}}M:=_{i0}T_A^iM$$
acquires an obvious graded associative (unital) algebra structure, called the tensor algebra of an $`A`$-bimodule $`M`$.
Given a $`\mathrm{𝕜}`$-vector space $`V`$, we will often write $`V^n`$ instead of $`T_\mathrm{𝕜}^nV`$.
The tensor algebra construction may be usefully interpreted as an adjoint functor. Specifically, fix an associative unital algebra $`A`$, and consider the category $`A\text{-algebras},`$ whose objects are pairs $`(B,f)`$, where $`B`$ is an associative unital algebra and $`f:AB`$ is an algebra morphism such that $`f(1)=1`$. Morphisms in $`A\text{-algebras}`$ are defined as algebra maps $`\phi :BB^{}`$ making the following natural diagram commute
Note that an algebra morphism $`AB`$ makes $`B`$ an $`A`$-bimodule. Thus, we get a functor $`A\text{-algebras}A\text{-bimodules}`$. It is straighforward to check that the assignment $`MT_A^{\text{}}M`$ gives a right adjoint to that functor, i.e., there is a canonical adjunction isomorphism
$$\mathrm{Hom}_{_{A\text{-bimodules}}}(M,B)\stackrel{_{}}{}\mathrm{Hom}_{_{A\text{-algebras}}}(T_A^{\text{}}M,B),$$
for any $`A`$-bimodule $`M`$ and an $`A`$-algebra $`B`$.
Now, fix an algebra morphism $`AB`$, and view $`B`$ as an $`A`$-bimodule. Let $`M`$ be another $`A`$-bimodule. Any morphism $`\delta :MB,`$ of $`A`$-bimodules induces, by the adjunction isomorphism above, an algebra homomorphism $`T_A^{\text{}}MB`$. This way, the algebra $`B`$ may be regarded as a $`T_A^{\text{}}M`$-bimodule.
We leave to the reader to prove the following
###### Lemma 3.4.1.
The $`A`$-bimodule map $`\delta :MB`$ can be uniquely extended to a derivation, resp., super-derivation,
$$T(\delta ):T_A^{\text{}}MB\text{such that}T(\delta )|_{T_A^0M}=f,\text{and}T(\delta )|_{T_A^1M}=\delta .\mathrm{}$$
### 3.5. Picard group of a category.
Given a category $`𝒞`$ we let $`\mathrm{Pic}(𝒞)`$ be the group of all autoequivalences of $`𝒞`$, the group structure being given by composition.
For example, let $`A`$ be an associative algebra. We need the following
###### Definition 3.5.1.
A finitely-generated $`A`$-bimodule $`Q`$ is said to be invertible if there exists a finitely-generated $`A`$-bimodule $`P`$ such that $`Q_APAP_AQ`$. Then, $`P`$ is called an inverse of $`Q`$.
It is clear that (the isomorphism classes of) invertible $`A`$-bimodules form a group under the tensor product operation $`P,P^{}P_AP^{}`$. The unit element of this group is the isomorphism class of $`A`$, viewed as an $`A`$-bimodule.
Now, let $`𝒞:=A\text{-}\mathrm{𝗆𝗈𝖽}`$ be the category of left $`A`$-modules. We know by Morita theory, see Theorem 2.3.1, that any equivalence $`A\text{-}\mathrm{𝗆𝗈𝖽}A\text{-}\mathrm{𝗆𝗈𝖽}`$ has the form $`MQ_AM`$, for an invertible $`A`$-bimodule $`Q`$. Thus, we obtain
###### Corollary 3.5.2.
The group $`\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$ is canonically isomorphic to the group of (isomorphism classes of) invertible $`A`$-bimodules.∎
###### Remark 3.5.3.
In the traditional commutative algebraic geometry, given an algebraic variety $`X`$, one writes $`\mathrm{Pic}(X)`$ for the abelian group formed by the (isomorphism classes of) line bundles on $`X`$. Now assume $`X`$ is irreducible and affine, and put $`A=\mathrm{𝕜}[X]`$. Then any line bundle $``$ on $`X`$ gives a rank 1 projective left $`A`$-module $`L:=\mathrm{\Gamma }(X,)`$. This way, the group $`\mathrm{Pic}(X)`$ gets identified with the group formed by rank 1 projective left $`A`$-modules equipped with the tensor product structure. That gives, in view of Corollary 3.5.2, a certain justification for the name ‘Picard group’ of a category. $`\mathrm{}`$
Recall the group $`\mathrm{Aut}(A)`$ of automorphisms of the algebra $`A`$. Given an algebra automorphism $`g\mathrm{Aut}(A)`$ and an $`A`$-bimodule $`P`$, we define a new $`A`$-bimodule $`P^g`$ by ”twisting” the natural left action on $`P`$ via $`g`$, i.e., by letting $`aa^{}A^\mathrm{e}`$ act on $`pP`$ by the formula $`pg(a)pa^{}`$. It is clear that, given two automorphisms $`f,g:AA`$, there is a canonical isomorphism $`P^{fg}(P^g)^f.`$
Let $`\mathrm{Aut}_{\mathrm{𝗂𝗇}}(A)\mathrm{Aut}(A)`$ denote the (normal) subgroup formed by inner automorphisms of $`A`$, i.e., automorphisms of the form $`auau^1`$, where $`uA^\times `$ is an invertible element. We have the following result.
###### Proposition 3.5.4.
The assignment $`gA^g`$ yields a group homomorphism $`\mathrm{Aut}(A)`$ $`\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$. This homomorphism descends to a well-defined and injective homomorphism $`\mathrm{Aut}(A)/\mathrm{Aut}_{\mathrm{𝗂𝗇}}(A){}_{}{}^{}_{}^{}\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$.
###### Proof.
Let $`uA^\times `$ be an invertible element and $`g=g_u:auau^1`$, the corresponding inner automorphism of $`A`$. It is straightforward to verify that the map $`\phi :xux`$ yields an isomorphism of $`A`$-bimodules $`\phi :A\stackrel{_{}}{}A^{g_u}`$. Conversely, given $`g\mathrm{Aut}(A)`$ and an $`A`$-bimodule isomorphism $`\phi :A\stackrel{_{}}{}A^g`$, put $`u:=\phi (1)`$. Then we have $`\phi (a)=\phi (1a)=\phi (1)a=ua`$. It follows that $`g=g_u`$, moreover, $`u`$ is invertible since $`u^1=g^1(1)`$. This completes the proof. ∎
It is instructive to think of the group $`\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$ as some sort of Lie group. The ‘Lie algebra’ of that group should therefore be formed by $`A`$-bimodules that are ‘infinitesimally close’ to $`A`$, the unit of $`\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$. It may be argued that any right $`A`$-module that is ‘infinitesimally close’ to a rank one free right $`A`$-module is itself isomorphic to a rank one free right $`A`$-module. Thus, any object in $`\mathrm{Lie}\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$ may be viewed as being the right $`A`$-module $`A`$ on which the standard left multiplication-action is ‘infinitesimally deformed’. Denote this ‘deformed’ left action of an element $`aA`$ by $`bab`$. The deformed left action must commute with the standard right action, hence, for any $`a,b,cA`$, we must have $`(ab)c=a(bc)`$. This forces the deformed action to be of the form $`ab=F(a)b`$, where $`F:AA`$ is a certain linear map that should be ‘infinitesimally close’ to the identity $`\mathrm{Id}:AA`$. We express the latter condition by introducing the ring $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$ of dual numbers and writing the map $`F`$ in the form $`F(a)=a+\epsilon \psi (a)`$, as we have already done earlier. Now the condition that the assignment $`A\times A\mathrm{𝕜}[\epsilon ]/\epsilon ^2A,a,bab=F(a)b`$ gives a (left) action becomes the equation $`F(aa^{})=F(a)F(a^{}),a,a^{}A`$. The last equation is equivalent to the condition that $`\psi `$ is a derivation, see (3.1.6).
Furthermore, arguing similarly, one finds that any inner derivation $`\mathrm{ad}a\mathrm{Inn}(A)`$ gives rise to a bimodule of the form $`A^g`$ where $`g:b(1+\epsilon a)b(1+\epsilon a)^1`$ is an ‘infinitesimal inner automorphism’. Thus, we conclude that the ‘Lie algebra’ of the group $`\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})`$ is given by the formula
$$\mathrm{Lie}\mathrm{Pic}(A\text{-}\mathrm{𝗆𝗈𝖽})=\mathrm{Der}(A)/\mathrm{Inn}(A).$$
### 3.6. Atiyah algebra of a vector bundle.
Let $`X`$ be an affine variety, and $``$ an algebraic vector bundle on $`X`$, locally trivial in the Zariski topology. Giving $``$ is equivalent to giving $`E:=\mathrm{\Gamma }(X,)`$, the vector space of global sections of $``$, which is a finitely generated projective $`\mathrm{𝕜}[X]`$-module.
Write $`\mathrm{End}()`$ for the (noncommutative) associative algebra of $`\mathrm{𝕜}[X]`$-linear endomorphisms of the vector bundle $``$. The algebra $`\mathrm{End}()`$ acts naturally on $`E:=\mathrm{\Gamma }(X,)`$. The center of the algebra $`\mathrm{End}()`$ is the subalgebra $`\mathrm{𝕜}[X]=\mathrm{𝕜}[X]\mathrm{Id}_{}{}_{}{}^{}_{}^{}\mathrm{End}()`$, formed by scalar endomorphisms.
The following result provides a basic example of Morita equivalence
###### Theorem 3.6.1.
The algebras $`\mathrm{𝕜}[X]`$ and $`\mathrm{End}()`$ are Morita equivalent. Specifically, the following functors
$`S:ME_{\mathrm{𝕜}[X]}M,\text{and}T:F\mathrm{Hom}_{_{\mathrm{End}()}}(E,F)`$
provide mutually inverse equivalences.
###### Definition 3.6.2.
The Lie algebra $`𝒜():=\mathrm{Der}(\mathrm{End}()),`$ the derivations of the associative algebra $`\mathrm{End}(),`$ is called the Atiyah algebra of $``$.
To obtain a more explicit description of the Atiyah algebra, assume that the affine variety $`X`$ is smooth. Let $`𝒟()`$ be the (associative) algebra of algebraic differential operators acting on sections of $``$. This algebra has a natural increasing filtration $`0=𝒟_1()𝒟_0()𝒟_1()\mathrm{},`$ by the order of differential operator. In particular, we have $`𝒟_0()=\mathrm{End}()`$.
Write $`\mathrm{gr}_{\text{}}𝒟()=_{i0}𝒟_i()/𝒟_{i1}()`$ for the associated graded algebra. Assigning the principal symbol to a differential operator gives rise to a canonical graded algebra isomorphism
$$\sigma :\mathrm{gr}𝒟()\stackrel{_{}}{}\mathrm{End}()_{_{\mathrm{𝕜}[X]}}\mathrm{Sym}^{\text{}}𝒯(X).$$
###### Remark 3.6.3.
Note that the algebra $`\mathrm{gr}𝒟()`$ is not commutative unless $``$ has rank one, i.e., unless $``$ is a line bundle. $`\mathrm{}`$
The top row of the diagram below is a natural short exact sequence involving the principal symbol map on the space of first order differential operators.
The space $`𝒟_1^{\mathrm{}}()`$ in the bottom row is formed by first order differential operators with scalar principal symbol.
The Theorem below shows that the space $`𝒟_1^{\mathrm{}}()`$ is closely related to the Atiyah algebra $`𝒜()`$.
###### Theorem 3.6.4.
$`(𝗂)`$The space $`𝒟_1^{\mathrm{}}()`$ is a Lie subalgebra in the associative algebra $`𝒟()`$ (with respect to the commutator bracket).
$`(\mathrm{𝗂𝗂})`$The bottom row in the diagram above is an extension of Lie algebras; in particular, $`𝒟_0()=\mathrm{End}()𝒟_1^{\mathrm{}}()`$ is a Lie ideal.
$`(\mathrm{𝗂𝗂𝗂})`$The adjoint action of an element $`u𝒟_1^{\mathrm{}}()`$ in the ideal $`\mathrm{End}()`$ gives a derivation, $`\mathrm{ad}u,`$ of $`\mathrm{End}()`$ viewed as an associative algebra.
$`(\mathrm{𝗂𝗏})`$The assignment $`u\mathrm{ad}u`$ is a Lie algebra homomorphism whose kernel is the space $`\mathrm{𝕜}[X]\mathrm{End}()=𝒟_0(),`$ of scalar endomorphisms. In particular, this space $`\mathrm{𝕜}[X]`$ is a Lie ideal in $`𝒟_1^{\mathrm{}}()`$.
$`(𝗏)`$The adjoint action described in (iii) gives rise to the following Lie algebra exact sequence
$$0\mathrm{𝕜}[X]𝒟_1^{\mathrm{}}()\stackrel{u\mathrm{ad}u}{}\mathrm{Der}\mathrm{End}()0.$$
Thus, the map $`u\mathrm{ad}u`$ induces an isomorphism $`𝒟_1^{\mathrm{}}()/\mathrm{𝕜}[X]\stackrel{_{}}{}𝒜().`$
###### Proof.
Observe first that the Theorem may be seen as a unification of two special cases:
$``$ For $`=𝒪_X,`$ the trivial rank one bundle, we have $`\mathrm{End}()=\mathrm{𝕜}[X]`$, and the derivations of the latter algebra are given by vector fields on $`X`$, by definition.
$``$ If $`X=\{pt\}`$ is a single point, then we have $`\mathrm{End}()=\mathrm{End}(E)\mathrm{Mat}_n\mathrm{𝕜},`$ is a matrix algebra, and any derivation of the matrix algebra is well-known to be inner.
Thus, our Theorem says, essentially, that in the general case of an arbitrary vector bundle on a variety $`X`$, the Lie algebra of derivations of $`\mathrm{End}()`$ is an extention of the Lie algebra of vector fields by inner derivations.
To prove the Theorem, we use Morita invariance of Hochschild cohomology of an algebra, cf. §5 below. In particular, applying this for the 1-st Hochschild cohomology of two Morita equivalent algebras, $`A`$ and $`B`$, we get $`\mathrm{Der}(A)/\mathrm{Inn}(A)\mathrm{Der}(B)/\mathrm{Inn}(B).`$
Now, take $`A=\mathrm{𝕜}[X]`$ and $`B=\mathrm{End}`$. Since $`A`$ is commutative, we have $`\mathrm{Inn}(A)=0`$, hence $`\mathrm{Der}(A)/\mathrm{Inn}(A)=\mathrm{Der}(A)=𝒯(X)`$. On the other hand, inner derivations of the algebra $`B`$ form the Lie algebra $`B/𝖹_B\mathrm{End}()/\mathrm{𝕜}[X]`$. Thus, the isomorphism of the previous paragraph yields a short exact sequence
$$\frac{\mathrm{End}()}{\mathrm{𝕜}[X]}=\mathrm{Inn}(B){}_{}{}^{}_{}^{}\mathrm{Der}(B){}_{}{}^{}_{}^{}\frac{\mathrm{Der}(B)}{\mathrm{Inn}(B)}=\frac{\mathrm{Der}(A)}{\mathrm{Inn}(A)}=𝒯(X).$$
We leave to the reader to verify that this short exact sequence coincides with the bottom row of the diagram preceeding Theorem 3.6.4. ∎
## 4. The Bar Complex
### 4.1. Free product of algebras
Given two associative algebras, $`A`$ and $`B`$, let $`AB`$ be the $`\mathrm{𝕜}`$-vector space whose basis is formed by words in elements of $`A`$ and $`B`$, with additional relations that adjacent elements from the same algebra are multiplied together. If, in addition, the algebras have units $`1_AA`$ and $`1_BB`$, then we impose the relation that multiplication by either unit acts as identity.
More generally, given two algebras $`A,B`$, and algebra imbeddings $`ı:C{}_{}{}^{}_{}^{}A,`$ and $`ȷ:C{}_{}{}^{}_{}^{}B,`$ such that $`ı(1_C)=1_A,`$ and $`ȷ(1_C)=1_B`$, one defines $`A__CB,`$ the free product of $`A`$ and $`B`$ over $`C`$, as the following unital associative algebra
$$A__CB:=\frac{T_\mathrm{𝕜}(AB)}{\begin{array}{c}aa^{}=aa^{},bb^{}=bb^{}\\ ı(c)=ȷ(c),\mathrm{\hspace{0.17em}1}_A=1=1_B\end{array}}_{a,a^{}A,b,b^{}B,cC}$$
where $`\mathrm{}`$ denotes the two-sided ideal generated by the indicated relations.
The operation of free product of associative algebras plays a role somewhat analogous to the role of tensor product for commutative associative algebras.
### 4.2.
Throughout, we let $`A`$ be an associative $`\mathrm{𝕜}`$-algebra (with $`1`$ as usual). We wish to associate to $`A`$ a sequence of homology groups which will play the noncommutative role of (co-)homology of a space. In order to do this, we wish to construct a particular resolution of $`A`$ by free $`A`$-bimodules. Before beginning, we remark that an $`A`$-bimodule is the same thing as a left $`A_\mathrm{𝕜}A^{\mathrm{op}}`$-module; if $`M`$ is an $`A`$-bimodule, then we define the action of $`AA^{\mathrm{op}}`$ on $`M`$ by $`(a_1a_2^{\mathrm{op}})m:=a_1ma_2.`$
We consider the following complex of $`A`$-bimodules
$$\text{},$$
where $`m:AAA`$ is the multiplication on $`A`$ and $`b:A^{(n+1)}A^n`$ is given by
$`b(a_0a_1\mathrm{}a_n)`$ $`:=a_0a_1a_2\mathrm{}a_n+`$
$`+{\displaystyle \underset{j=1}{\overset{n1}{}}}(1)^ja_0\mathrm{}(a_ja_{j+1})\mathrm{}a_n.`$
It is a tedious (but simple) calculation that $`b^2=0`$.
###### Definition 4.2.1.
We set $`A^\mathrm{e}:=AA^{\mathrm{op}}`$, and for any $`i=0,1,2,\mathrm{},`$ put $`𝖡_iA:=AA^iA`$, a free $`A^\mathrm{e}`$-module generated by the $`\mathrm{𝕜}`$-vector space $`A^i`$. This way, the complex above can be writen as the following *bar complex*
$$𝖡_{\text{}}A:\left[\text{}\right]{}_{}{}^{}_{}^{}A.$$
We claim that this sequence is exact, i.e., that the bar complex provides a free $`A`$-bimodule resolution of $`A`$, viewed as an $`A`$-bimodule. To show this, we will construct a chain homotopy $`h:A^iA^{(i+1)}`$ such that $`bh+hb=\mathrm{id}`$. Usual homological algebra then implies exactness. We will first construct $`h`$ “by hand,” then give alternate descriptions of the bar complex which will make this definition more natural (and the proof of the homotopy easier). In particular, we define
$$h(a_1\mathrm{}a_i)=1_Aa_1\mathrm{}a_i.$$
We check that $`h`$ indeed satisfies $`bh+hb=\mathrm{id}`$ on the degree $`2`$ piece. That is, we will show that $`h(b(a_1a_2))+b(h(a_1a_2))=a_1a_2`$. Now, by definition $`b(a_1a_2)=a_1a_2`$, and $`h(a_1a_2)=1_Aa_1a_2`$. For the second term, $`h(a_1a_2)=1_Aa_1a_2`$, and by the definition of $`b`$ we have
$$b(1_Aa_1a_2)=1_Aa_1a_21_Aa_1a_2=a_1a_21_Aa_1a_2.$$
So,
$$(hb+bh)(a_1a_2)=1_Aa_1a_2+a_1a_21_Aa_1a_2=a_1a_2,$$
as desired.
### 4.3. Second construction of the bar complex (after Drinfeld).
Let $`A\mathrm{𝕜}[\epsilon ]`$ be the free product of the algebra $`A`$ and the polynomial algebra $`\mathrm{𝕜}[\epsilon ]`$ in one variable $`\epsilon `$. An element of this free product can be written in the form $`a_1\epsilon ^{n_1}a_2\epsilon ^{n_2}\mathrm{}a_k`$ for elements $`a_1,\mathrm{},a_kA`$ and non-negative integers $`n_1,\mathrm{},n_{k1}`$ (if the last factor is from $`\mathrm{𝕜}[\epsilon ]`$, simply “pad” with $`1_A`$, and similarly for the first factor). However, we can rewrite $`\epsilon ^{n_j}=1_A\epsilon 1_A\mathrm{}1_A\epsilon 1_A`$, where we have $`n_j`$ factors of $`\epsilon `$. So, we can always write any element of $`A\mathrm{𝕜}[\epsilon ]`$ in the form $`a_1\epsilon a_2\epsilon \mathrm{}\epsilon a_k`$ for some $`a_1,\mathrm{},a_kA`$. So, the $`\epsilon `$ plays no role other than a separator, and we shall replace it by a bar, that is, we define $`a_1a_2\mathrm{}a_k:=a_1\epsilon a_2\mathrm{}\epsilon a_k`$. This notation is the genesis of the name “bar complex.”
Now, we put a grading on $`A\mathrm{𝕜}[\epsilon ]`$ by declaring $`\mathrm{deg}a=0`$ for all $`aA`$ and $`\mathrm{deg}\epsilon =1`$. We will now make $`A\mathrm{𝕜}[\epsilon ]`$ a differential graded algebra (DGA from now on) by defining a super-differential $`d:A\mathrm{𝕜}[\epsilon ]A\mathrm{𝕜}[\epsilon ]`$. Recall that a super-differential is simply a super-derivation $`d`$ satisfying $`d^2=0`$. We define $`d`$ on generators, namely, we set $`da=0`$ for all $`aA`$ and $`d\epsilon =1_A`$, and we extend it to $`A\mathrm{𝕜}[\epsilon ]`$ uniquely by requiring that it obey the graded Leibniz rule. If we now identify $`A\epsilon A\epsilon \mathrm{}\epsilon A`$ ($`n`$ factors of $`A`$) with $`A^n`$ in the obvious fashion, we obtain an identification of $`𝖡_{\text{}}A`$ with $`A\mathrm{𝕜}[\epsilon ]`$, and the super-differential $`d`$ on $`A\mathrm{𝕜}[\epsilon ]`$ becomes the bar differential.
Observe next that $`d^2(\epsilon )=d(1)=0`$, and also $`d^2(a)=0`$, for any $`aA`$, by definition. Hence, Corollary 3.3.3 yields $`d^2=0`$. We conclude that $`d`$ is a differential on $`A\mathrm{𝕜}[\epsilon ]`$.
In this context, the proof of exactness becomes trivial. Since $`A\mathrm{𝕜}[\epsilon ]`$ is a DGA, we can calculate its cohomology in the usual fashion. The claim that $`A\mathrm{𝕜}[\epsilon ]`$ is exact (i.e., that $`𝖡_{\text{}}A`$ is exact) is equivalent to claiming that $`H^{\text{}}(A\mathrm{𝕜}[\epsilon ])=0`$. Now, by definition we have that $`d\epsilon =1_A`$, hence $`1_A`$ is a coboundary and therefore zero in cohomology. But since $`A\mathrm{𝕜}[\epsilon ]`$ is a DGA, $`H^{\text{}}(A\mathrm{𝕜}[\epsilon ])`$ is an algebra (it’s even graded, but that is not important for our purposes). In particular, the cohomology class of $`1_A`$ acts as a multiplicative unit for $`H^{\text{}}(A\mathrm{𝕜}[\epsilon ])`$, and since $`[1_A]=0`$, we see that $`H^{\text{}}(A\mathrm{𝕜}[\epsilon ])=0`$. We can of course mimic the previous proof and construct a (co-)chain homotopy. In the notation of $`A\mathrm{𝕜}[\epsilon ]`$, the homotopy $`h`$ is defined by $`h(u)=\epsilon u`$ for all $`uA\mathrm{𝕜}[\epsilon ]`$. Then $`d(\epsilon u)=d\epsilon u\epsilon du=1_Au\epsilon du`$. So, $`(dh+hd)(u)=u\epsilon du+\epsilon du=u`$, as required.
### 4.4. Third construction of the bar complex.
The following construction only applies (as presented) in the case where $`A`$ is finite dimensional. Let $`A^{}=\mathrm{Hom}_\mathrm{𝕜}(A,\mathrm{𝕜})`$ be the dual vector space of $`A`$. Then we form its non-unital tensor algebra
$$T^+(A^{}):=A^{}(A^{}A^{})\mathrm{}(A^{})^k\mathrm{}.$$
This is the free associative algebra on $`A^{}`$ with no unit.
###### Proposition 4.4.1.
Giving an associative algebra structure on $`A`$ is equivalent to giving a map $`d:T^+(A^{})T^+(A^{})`$ such that
$`(𝗂)`$$`d`$ is a super-derivation of degree $`1`$, i.e., $`d((A^{})^k)(A^{})^{(k+1)}`$; and
$`(\mathrm{𝗂𝗂})`$$`d^2=0`$.
###### Proof.
First, suppose we have any linear map $`d:A^{}A^{}A^{}`$. Then we can always extend this map to a super-derivation on $`T^+(A^{})`$ by applying the super-Leibniz rule. Now, if we are given a multiplication map $`m:AAA`$, then by taking transposes we obtain a map $`m^{_{}}:A^{}(AA)^{}`$. Since $`A`$ is finite dimensional, $`(AA)^{}`$ and $`A^{}A^{}`$ are canonically isomorphic. So, we can regard the transpose of multiplication as a map $`m^{_{}}:A^{}A^{}A^{}`$. We let $`d=m^{_{}}`$ and extend this to a super-derivation as described above. Of course, if we are already given a super-derivation of $`T^+(A^{})`$, then we can transpose its restriction to $`A^{}A^{}`$ to obtain a multiplication on $`A`$ (since we can identify $`A^{}`$ with $`A`$ and $`(A^{}A^{})^{}`$ with $`AA`$).
We will now show that the associative law for $`m`$ is equivalent to $`d^2=0`$ where $`m`$ and $`d`$ are related as in the previous paragraph. Since $`T^+(A^{})`$ is generated by elements of $`A^{}`$, we need only show that $`d^2:A^{}A^{}A^{}A^{}`$ is the zero map–the super-Leibniz rule will take care of the rest. So, suppose we are given some linear functional $`\lambda A^{}`$. Then $`d\lambda A^{}A^{}`$, which we have canonically identified with $`(AA)^{}`$. In particular, $`d\lambda (ab)=m^{_{}}\lambda (ab)=\lambda (m(a,b))=\lambda (ab)`$ (where we will write $`m(a,b)=ab`$ for simplicity’s sake). Further, since $`d\lambda A^{}A^{}`$, we can find $`\mu _1,\mathrm{},\mu _n,\nu _1,\mathrm{},\nu _nA^{}`$ so that $`d\lambda =_{i=1}^n(\mu _i\nu _i)`$. Then we have that
$$d\lambda (ab)=\underset{i=1}{\overset{n}{}}(\mu _i\nu _i)(ab)=\underset{i=1}{\overset{n}{}}\mu _i(a)\nu _i(b).$$
Now, we wish to consider $`d(d\lambda )`$. Since this lies in $`A^{}A^{}A^{}`$, we can view this is a linear functional on $`A^3`$. Then, using the fact that $`d`$ is a super-derivation, we obtain
$`d(d\lambda )(abc)`$ $`=d\left[{\displaystyle \underset{i=1}{\overset{n}{}}}(\mu _i\nu _i)\right](abc)`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}(d\mu _i\nu _i\mu _id\nu _i)(abc)`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}d\mu _i(ab)\nu _i(c){\displaystyle \underset{i=1}{\overset{n}{}}}\mu _i(a)d\nu _i(bc)`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}\mu _i(ab)\nu _i(c){\displaystyle \underset{i=1}{\overset{n}{}}}\mu _i(a)\nu _i(bc)`$
$`=\lambda ((ab)c)\lambda (a(bc))=\lambda ((ab)ca(bc)).`$
So, if $`d^2=0`$, we see that $`\lambda ((ab)ca(bc))=0`$ for every functional $`\lambda `$, hence $`(ab)c=a(bc)`$ (i.e., $`m`$ is associative). Conversely, if $`m`$ is associative, we see that $`d^2=0`$. ∎
This is not quite the bar complex, since for one instance the differential goes in the wrong direction. But since $`d:T^i(A^{})T^{i+1}(A^{})`$, we can consider the transpose of $`d`$, $`d^{}:(T^{i+1}(A^{}))^{}(T^i(A^{}))^{}`$. Since $`A`$ was assumed finite dimensional, we can identify $`T^i(A^{})^{}`$ and $`T^i(A)`$, so finally $`d^{}:T^{i+1}(A)T^i(A)`$ is the desired complex. One advantage of this approach is that similar constructions can be made for different algebraic structures.
The reader is invited to consider how such a construction can be performed for, say Lie algebras.
### 4.5. Reduced Bar complex.
It turns out that the Bar complex $`𝖡_{\text{}}(A)`$ contains a large acyclic subcomplex. Specifically, for each $`n>1`$, in $`𝖡_nA=AA^nA`$ consider the following $`A`$-subbimodule
$$\mathrm{Triv}_nA:=\underset{i=1}{\overset{n}{}}A\left(A^{(i1)}\mathrm{𝕜}A^{(ni)}\right)AAA^nA.$$
It is easy to see from the formula for the bar-differential that $`b(\mathrm{Triv}_nA)\mathrm{Triv}_{n1}A`$. Thus, $`\mathrm{Triv}_{\text{}}A`$ is a subcomplex in $`𝖡_{\text{}}A`$.
Define the reduced bar complex to be the quotient $`\overline{𝖡}_{\text{}}A:=𝖡_{\text{}}A/\mathrm{Triv}_{\text{}}A`$. By definition, the $`n^{\text{th}}`$ term of the reduced bar complex is
$$\overline{𝖡}_nA=𝖡_nA/\mathrm{Triv}_nA=A\overline{A}^nA,$$
where $`\overline{A}=A/\mathrm{𝕜}`$ as a vector space. The differential is the one induced by $`d`$ on the bar complex.
The reduced bar complex has the following interpretation in terms of the free product construction $`𝖡_{\text{}}A=A\mathrm{𝕜}[\epsilon ]`$, see §4.3. Observe that since $`\mathrm{deg}\epsilon =1`$, the Leibniz formula for an odd derivation yields
$$d(\epsilon ^2)=(d\epsilon )\epsilon \epsilon (d\epsilon )=1\epsilon \epsilon 1=0.$$
It follows readily that the two-sided ideal $`\epsilon ^2A\mathrm{𝕜}[\epsilon ]`$, generated by the element $`\epsilon ^2`$ is $`d`$-stable, i.e., we have $`d\left(\epsilon ^2\right)\epsilon ^2`$. Hence, the differential $`d`$ on $`A\mathrm{𝕜}[\epsilon ]`$ descends to a well-defined differential on the graded algebra $`(A\mathrm{𝕜}[\epsilon ])/\epsilon ^2`$. We claim that, under the identification of §4.3, we have $`\mathrm{Triv}^{\text{}}A=\epsilon ^2`$, and therefore
$$\overline{𝖡}_{\text{}}A=(A\mathrm{𝕜}[\epsilon ])/\epsilon ^2=A(\mathrm{𝕜}[\epsilon ]/\epsilon ^2).$$
(4.5.1)
To see this, notice that an element of $`AA^nA=𝖡_nA`$ belongs to $`\mathrm{Triv}_nA`$ if and only if it is a $`\mathrm{𝕜}`$-linear combination of terms, each involving a subexpression like $`(\mathrm{}|1_A|\mathrm{})`$. But such an expresion, when translated into the free product construction, reads: $`(\mathrm{}\epsilon 1_A\epsilon \mathrm{})=(\mathrm{}\epsilon ^2\mathrm{})`$, and our claim follows.
Observe next that the argument proving acyclicity of the complex $`(A\mathrm{𝕜}[\epsilon ],d)`$ applies verbatim to yield acyclicity of $`(A\mathrm{𝕜}[\epsilon ]/\epsilon ^2,d)`$. We conclude
###### Corollary 4.5.2.
The reduced bar complex provides a free $`A`$-bimodule resolution of $`A`$.∎
## 5. Hochschild homology and cohomology
### 5.1.
Given an associative $`\mathrm{𝕜}`$-algebra $`A`$, let $`A^{\mathrm{op}}`$ denote the opposite algebra, and write $`A^\mathrm{e}=AA^{\mathrm{op}}`$. There is a canonical isomorphism $`(A^\mathrm{e})^{\mathrm{op}}A^\mathrm{e}`$. Thus, an $`A`$-bimodule is the same thing as a left $`A^\mathrm{e}`$-module, and also the same thing as a right $`A^\mathrm{e}`$-module. Recall the notation $`A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}`$ for the abelian category of all $`A`$-bimodules.
Given an $`A`$-bimodule $`M`$, we write $`[M,A]`$ for the commutator space, the $`\mathrm{𝕜}`$-vector subspace of $`M`$ spanned by all commutators $`maam,aA,mM`$. Let $`M/[M,A]`$ denote the corresponding commutator quotient of $`M`$.
Equivalently, any $`A`$-bimodule $`M`$ may be viewed as either right or left $`A^\mathrm{e}`$-module. In particular, we view $`M`$ as a right $`A^\mathrm{e}`$-module and view $`A`$ as a left $`A^\mathrm{e}`$-module, and form the tensor product $`M_{A^\mathrm{e}}A`$. One has a canonical vector space identification $`M_{A^\mathrm{e}}A=M/[M,A]`$. The assignment $`MM_{A^\mathrm{e}}A`$ clearly gives a right exact functor from $`A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}`$ to the category of $`\mathrm{𝕜}`$-vector spaces. So, $`M_{A^\mathrm{e}}`$ has left derived functors, the Tor functors $`\mathrm{Tor}_i^{A^\mathrm{e}}(M,A)`$. For ease of notation, we denote $`\mathrm{Tor}_i^{A^\mathrm{e}}(M,A)`$ by $`HH_i(M)`$, the *$`i^{\text{th}}`$ Hochschild homology group* (which is really a $`\mathrm{𝕜}`$-vector space) of $`M`$ over $`A`$.
By definition, for any $`A`$-bimodule $`M`$, we have
$$HH_0(M)=\mathrm{Tor}_0^{A^\mathrm{e}}(M,A)=M/[M,A].$$
In the particular case that $`M=A`$, we obtain $`HH_0(A)=A/[A,A]`$. Notice that $`[A,A]`$ is not an ideal in $`A`$, simply a $`\mathrm{𝕜}`$-linear subspace of $`A`$.
Computing higher degree Hochschild homology groups requires a choice of some projective resolution of $`A`$ as $`A`$-bimodules (i.e., left $`A^\mathrm{e}`$-modules). The bar complex provides a canonical choice of such resolution. So, to compute the groups $`HH_i(M)`$ we need only apply the functor $`M_{A^\mathrm{e}}`$ to the bar complex
$$𝖡_{\text{}}A:\text{}.$$
We then tensor this on the left with $`M`$ over $`A^\mathrm{e}`$ to yield
$$M_{A^\mathrm{e}}𝖡_{\text{}}:\text{}.$$
To simplify this, pick some $`m(a_0\mathrm{}a_n)M_{A^\mathrm{e}}A^n`$ and write $`a_0a_1\mathrm{}a_{n1}a_n=(a_0a_n^{\mathrm{op}})(1_Aa_1\mathrm{}a_{n1}1_A)`$. Then
$`m(a_0a_1\mathrm{}a_{n1}a_n)`$ $`=m(a_0a_n^{\mathrm{op}})(1_Aa_1\mathrm{}a_{n1}1_A)`$
$`=a_nma_0(1_Aa_1\mathrm{}a_{n1}1_A).`$
By dropping the two $`1_A`$’s and observing that only scalars on the intermediate $`a_j`$’s “commute” past $``$, we can identify $`M_{A^\mathrm{e}}A^n`$ with $`MA^{(n2)}`$. So, the complex $`M_{A^\mathrm{e}}𝖡_{\text{}}A`$ becomes
$$\text{}.$$
Examining the identification of $`M_{A^\mathrm{e}}A^n`$ with $`MA^{(n2)}`$, we find that the differential for this complex is given by (again, using the bar notation)
$`d(ma_1`$ $`\mathrm{}a_n)=ma_1a_2\mathrm{}a_n+`$
$`{\displaystyle \underset{i=1}{\overset{n1}{}}}(1)^ima_1\mathrm{}(a_ia_{i+1})\mathrm{}a_n+(1)^na_nma_1\mathrm{}a_{n1}.`$
We remark that the bar complex, resp., the reduced bar complex, for $`A`$ can be recovered from the Hochschild chain complex, resp., reduced Hochschild chain complex, of the free rank one $`A^\mathrm{e}`$-module $`A^\mathrm{e}=AA^{\mathrm{op}}`$ (viewed as an $`A`$-bimodule) as follows
$$𝖡_{\text{}}A=C_{\text{}}(A,AA^{\mathrm{op}}),\text{resp.,}\overline{𝖡}_{\text{}}A=\overline{C}_{\text{}}(A,AA^{\mathrm{op}}).$$
(5.1.1)
### 5.2. Hochschild cohomology
As before, we take $`A`$ to be an associative $`\mathrm{𝕜}`$-algebra and $`M`$ is an $`A`$-bimodule. For Hochschild homology, we considered $`\mathrm{Tor}_{A^\mathrm{e}}^{\text{}}(A,M)`$. Now we wish to consider $`\mathrm{Ext}_{A^\mathrm{e}}^{\text{}}(A,M)`$.
Recall that $`A^\mathrm{e}=AA^{\mathrm{op}}`$. We define the *Hochschild cohomology* of $`A`$ with coefficients in $`M`$ by $`HH^{\text{}}(M):=\mathrm{Ext}_{A^\mathrm{e}}^{\text{}}(A,M)`$.
###### Proposition 5.2.1.
The functors $`HH_{\text{}}`$ and $`HH^{\text{}}`$ are both Morita invariant. In particular,
$$HH_{\text{}}(A)=HH_{\text{}}(\mathrm{Mat}_r(A)),\text{and}HH^{\text{}}(A)=HH^{\text{}}(\mathrm{Mat}_r(A))$$
where $`\mathrm{Mat}_r(A)`$ denotes $`r\times r`$-matrices over $`A`$.
###### Proof.
This is immediate from the definition:
$$HH^i(A,A)=\mathrm{Ext}_{A^e\text{-}\mathrm{𝗆𝗈𝖽}}^i(A,A)\text{and}HH_i(A,A)=\mathrm{Tor}_{A^e\text{-}\mathrm{𝗆𝗈𝖽}}^i(A,A),$$
since, for Morita equivalent algebras, the corresponding categories of bimodules are equivalent as abelian categories, hence give rise to isomorphic Ext and Tor groups.
It is instructive, however, to give a direct computational proof for the zeroth order homology, that is, to show that $`HH_0(A)=HH_0(\mathrm{Mat}_r(A))`$. We wish to construct a canonical isomorphism
$$A/[A,A]\mathrm{Mat}_r(A)/[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)].$$
We will define an isomorphism $`\mathrm{Tr}:\mathrm{Mat}_r(A)/[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]A/[A,A]`$ by using, as the notation suggests, the standard trace on $`\mathrm{Mat}_r(A)`$. In particular, $`\mathrm{Tr}`$ sends a matrix $`(a_{ij})`$ to the element
$$\mathrm{tr}(a_{ij})mod[A,A]=\underset{i=1}{\overset{n}{}}a_{ii}mod[A,A].$$
First, observe that $`\mathrm{Tr}(XY)=\mathrm{Tr}(YX)`$ for $`X,Y\mathrm{Mat}_r(A)`$, so $`\mathrm{Tr}`$ factors through a map (also called $`\mathrm{Tr}`$) $`\mathrm{Mat}_r(A)/[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]A/[A,A]`$. We wish to show that the kernel of $`\mathrm{Tr}`$ is precisely $`[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$. So, choose any matrix $`X=(a_{ij})\mathrm{Ker}\mathrm{Tr}`$. We will let $`E_{ij}(a)`$ denote the elementary matrix with $`a`$ in the $`ij`$-entry and zeroes elsewhere. Then for $`ij`$, we easily find that $`E_{ij}(a)=[E_{ij}(a),E_{jj}(1_A)]`$. So, all matrices with only off diagonal entries lie in $`[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$. Hence we can write
$$X=(a_{ij})=\mathrm{diag}(a_1,a_2,\mathrm{},a_r)mod[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)].$$
Also for $`ij`$, we can directly compute that $`[E_{ij}(a),E_{ji}(1_A)]=E_{ii}(a)E_{jj}(a)`$. So, for each $`j=2,\mathrm{},r`$, the matrix $`E_{11}(a_j)E_{jj}(a_j)[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$, hence
$$\mathrm{diag}(_{j=1}^ra_j,0,\mathrm{},0)\mathrm{diag}(a_1,\mathrm{},a_r)[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)].$$
So, we can write any matrix in $`\mathrm{Mat}_r(A)`$ as $`E_{11}(a)mod[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$ for some $`aA`$. Since,
$$0=\mathrm{Tr}(X)=\mathrm{Tr}(E_{11}(a))=amod[A,A],$$
we find that $`a[A,A]`$ and $`X[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$ to begin with. So, $`\mathrm{Ker}\mathrm{Tr}=[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]`$, and since $`\mathrm{Tr}`$ is clearly surjective we obtain that
$$\mathrm{Tr}:\mathrm{Mat}_r(A)/[\mathrm{Mat}_r(A),\mathrm{Mat}_r(A)]A/[A,A]$$
is an isomorphism. ∎
To calculate this cohomology, we will again use the bar complex $`𝖡_{\text{}}A`$:
$$\text{}.$$
Applying the functor $`\mathrm{Hom}_{A^\mathrm{e}}(,M)`$ to this complex (and accounting for contravariance), we obtain the sequence
$$\text{},$$
Recall that $`A^2`$ is free of rank one as an $`A`$-bimodule (which is just a left $`AA^\mathrm{e}`$-module). So, $`\mathrm{Hom}_{A^\mathrm{e}}(A^2,M)M`$. Similarly, if $`\phi :A^3M`$ is an $`A`$-bimodule map, then
$$\phi (a_1a_2a_3)=a_1\phi (1_Aa_21_A)a_3.$$
The association of $`\phi `$ to the $`\mathrm{𝕜}`$-linear map $`AM`$, $`a\phi (1_Aa1_A)`$ gives an isomorphism between $`\mathrm{Hom}_{A^\mathrm{e}}(A^3,M)`$ and $`\mathrm{Hom}_\mathrm{𝕜}(A,M)`$. Indeed, we find that $`\mathrm{Hom}_{A^\mathrm{e}}(A^n,M)\mathrm{Hom}_\mathrm{𝕜}(A^{(n2)},M)`$ for all $`n3`$. So, the complex whose cohomology we wish to compute reduces to
$$\text{}.$$
An explicit formula for $`b`$ in this interpretation is as follows. Recall that for the calculation of Hochschild homology, we used the differential
$`b(a_0a_1\mathrm{}a_n)=a_0a_1\mathrm{}a_n`$
$`+{\displaystyle \underset{j=1}{\overset{n1}{}}}a_0a_1\mathrm{}(a_ja_{j+1})\mathrm{}a_n+(1)^na_na_0a_1\mathrm{}a_{n1},`$
where we have $`a_0M`$ and $`a_jA`$ for $`j1`$. For the first differential, $`b:M\mathrm{Hom}_\mathrm{𝕜}(A,M)`$, we find that we must have $`bm(a)=amma`$. For $`n=1`$, we obtain for $`f\mathrm{Hom}_\mathrm{𝕜}(A,M)`$ that $`bf(a_1,a_2)=a_1f(a_2)f(a_1a_2)+f(a_1)a_2`$. Similarly, for all $`f\mathrm{Hom}_\mathrm{𝕜}(A^n,M)`$, we find that
$`bf(a_1,a_2,\mathrm{},a_{n+1})=a_1f(a_2,\mathrm{},a_{n+1})`$ $`+{\displaystyle \underset{i=1}{\overset{n}{}}}(1)^if(a_1,\mathrm{},a_ia_{i+1},\mathrm{},a_{n+1})`$
$`+(1)^{n+1}f(a_1,\mathrm{},a_n)a_{n+1}.`$
We will often write $`C^n(A,M):=\mathrm{Hom}_\mathrm{𝕜}(A^n,M)`$.
With these formulae for $`b`$ in hand, we can now explicitly calculate the first few cohomology groups.
### 5.3. Interpretation of $`HH^0`$.
For degree zero, we see that $`HH^0(M)=\mathrm{Ker}b`$. But $`m\mathrm{Ker}b`$ if and only if $`bm(a)=amma=0`$ for all $`aA`$. It is natural to call $`\{mMam=ma\}`$ the *center* of the module $`A`$. In particular, if $`M=A`$, then we in fact have that $`HH^0(A)=𝖹_A`$. Notice that there is a fair bit of “duality” in the degree zero setting between Hochschild homology and cohomology. For $`HH^0(M)`$, we obtain the center of $`M`$, and for $`HH_0(M)`$ we obtain the “cocenter” of $`M`$, namely $`M/[A,M]`$.
### 5.4. Interpretation of $`HH^1`$.
Moving on to $`HH^1(A,M)`$, we claim that the kernel of $`d`$ on $`C^1(A,M)`$ consists precisely of all derivations $`AM`$. Indeed, if $`fC^1(A,M)`$, then $`df(a_1,a_2)=0`$ is equivalent to requiring that
$$a_1f(a_2)f(a_1a_2)+f(a_1)a_2=0,$$
that is, that $`f(a_1a_2)=a_1f(a_2)+f(a_1)a_2`$. So, $`f`$ is a derivation. The image of $`d:M\mathrm{Hom}_\mathrm{𝕜}(A,M)`$ is precisely the set of all *inner* derivations, $`\mathrm{Inn}(A,M)`$, that is, derivations of the form $`aamma`$ for some $`mM`$. So,
$$HH^1(M)\mathrm{Der}(A,M)/\mathrm{Inn}(A,M),$$
the *outer derivations* from $`A`$ to $`M`$. In particular, we can rewrite (3.1.1) in the following way
$$0HH^0(A,M)M\stackrel{\mathrm{ad}}{}\mathrm{Der}(A,M)HH^1(A,M)0.$$
(5.4.1)
Observe that if $`M=A`$, then $`\mathrm{Der}(A)`$ is a Lie algebra and $`\mathrm{Inn}(A)`$ is a Lie ideal. So, $`HH^1(M)`$ inherits the structure of a Lie algebra.
### 5.5. Interpretation of $`HH^2`$.
This is given by considering algebra extensions. Suppose now that $`A`$ is an associative algebra and consider an extension of algebras
$$\text{},$$
where $`M`$ is an ideal of $`\stackrel{~}{A}`$ satisfying $`M^2=0`$. Then $`M`$ is an $`A`$-bimodule. Define $`am`$ by choosing a lift $`\stackrel{~}{a}`$ of $`a`$ in $`\stackrel{~}{A}`$ and setting $`am=\stackrel{~}{a}m`$. If $`\stackrel{~}{a}^{}`$ is another such lift, then since $`\stackrel{~}{a}\stackrel{~}{a}^{}`$ maps to zero in $`A`$, it must be an element of $`M`$. Hence $`(\stackrel{~}{a}\stackrel{~}{a}^{})mM^2=0`$, so $`\stackrel{~}{a}m=\stackrel{~}{a}^{}m`$.
Choose some $`\mathrm{𝕜}`$-linear splitting (as a vector space) $`c:A{}_{}{}^{}_{}^{}\stackrel{~}{A}`$. Then we get a vector space direct sum decomposition $`\stackrel{~}{A}c(A)M`$. Further, for any $`a_1,a_2A,`$ we have $`c(a_1)c(a_2)c(a_1a_2)M`$. Therefore we can write the product on $`\stackrel{~}{A}`$ as
$$(a_1m_1)(a_2m_2)=a_1a_2(a_1m_2+m_1a_2+\beta (a_1,a_2)),$$
where $`\beta `$ is an arbitrary bilinear map $`A^2M`$, i.e., an element of $`C^2(A,M)`$. This formula gives an associative product on $`\stackrel{~}{A}`$ if and only if $`d\beta =0`$. Indeed,
$$d\beta (a_1,a_2,a_3)=a_1\beta (a_2,a_3)\beta (a_1a_2,a_3)+\beta (a_1,a_2a_3)\beta (a_1,a_2)a_3.$$
Then by writing out $`[(a_1m_1)(a_2m_2)](a_3m_3)`$ and $`(a_1m_1)[(a_2m_2)(a_3m_3)]`$ and equating them, we see that $`d\beta `$ must equal zero. We can also check that if $`\beta _1`$ and $`\beta _2`$ define associative structures on $`\stackrel{~}{A}`$, then they define isomorphic extensions if and only if $`\beta _1\beta _2`$ is a Hochschild coboundary. So, $`HH^2(M)`$ classifies extensions of $`A`$ by $`M`$.
###### Remark 5.5.1.
If we consider the case where $`A`$ is a unital algebra and we only consider the square-zero extensions where $`\stackrel{~}{A}`$ is also a unital algebra, and the homomorphism $`\stackrel{~}{A}A`$ takes the unit to the unit, then it is easy to check that such extensions are also classified by $`HH^2(M),`$ but where $`HH^2(M)`$ is computed using the reduced bar complex. $`\mathrm{}`$
### 5.6. Interpretation of $`HH^3`$.
The group $`HH^3(M)`$ classifies so-called crossed-bimodules. A crossed-bimodule is a map $`\phi :CB`$ of associative algebras where $`B`$ is unital, $`C`$ is nonunital and maps to a $`2`$-sided ideal of $`B`$, the cokernel of $`\phi `$ is $`A`$, the kernel is $`M`$, and we require
$$\phi (bcb^{})=b\phi (c)b^{}\text{and}\phi (c)c^{}=cc^{}=c\phi (c^{}),$$
for all $`b,b^{}B`$ and $`c,c^{}C`$.
### 5.7. Reduced cochain complex
One can use the reduced bar complex to compute Hochschild homology and cohomology. Specifically, for any $`A`$-bimodule $`M`$, define the reduced Hochschild cochain complex $`\overline{C}^{\text{}}(A,M)`$ by
$`\overline{C}^n(A,M):=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\overline{𝖡}_nA,M)`$ $`=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A\overline{A}^nA,M)`$
$`\mathrm{Hom}_\mathrm{𝕜}(\overline{A}^n,M),n=0,1,\mathrm{}.`$
## 6. Poisson brackets and Gerstenhaber algebras
### 6.1. Polyvector fields
Let $`X`$ be an affine variety with coordinate ring $`A=\mathrm{𝕜}[X],`$ and let $`,,`$ be locally free coherent sheaves on $`X`$. We write $`E=\mathrm{\Gamma }(X,),`$ resp., $`F=\mathrm{\Gamma }(X,),`$ for the corresponding (projective) $`A`$-modules of global sections. Then, one has canonical isomorphisms
$$E_AF\mathrm{\Gamma }(X,_{_{𝒪_X}}),\mathrm{\Lambda }_A^pE\mathrm{\Gamma }(X,\mathrm{\Lambda }^p),\mathrm{Sym}_A^pE\mathrm{\Gamma }(X,\mathrm{Sym}^p).$$
In particular, for the tangent sheaf $`𝒯_X`$ we have $`\mathrm{\Gamma }(X,\mathrm{Sym}𝒯_X)=\mathrm{𝕜}[T^{}X],`$ is the algebra of regular functions on the total space of the cotangent bundle on $`X`$.
We introduce the notation $`\mathrm{\Theta }_p(X):=\mathrm{\Gamma }(X,\mathrm{\Lambda }^p𝒯_X)`$ for the vector space of $`p`$-polyvector fields on $`X`$. The graded-commutative algebra $`\mathrm{\Theta }_{\text{}}(X):=_p\mathrm{\Theta }_p(X)=\mathrm{\Gamma }(X,\mathrm{\Lambda }^{\text{}}𝒯_X)`$ may be thought of as a odd analogue of the commutative algebra $`\mathrm{\Gamma }(X,\mathrm{Sym}𝒯_X)=\mathrm{𝕜}[T^{}X]`$.
For any polyvector $`\pi \mathrm{\Theta }_p`$, one defines a natural contraction operator $`i_\pi :\mathrm{\Omega }^k(X)\mathrm{\Omega }^{kp}(X),\alpha i_\pi \alpha ,`$ where, for $`\pi =\xi _1\mathrm{}\xi _p`$, the $`(kp)`$-form $`i_\pi \alpha `$, is given by $`\eta _1,\mathrm{},\eta _{kp}\alpha (\xi _1,\mathrm{},\xi _p,\eta _1,\mathrm{},\eta _{kp}).`$ For $`p=1`$ this reduces to the standard contraction of a differential form with respect to a vector field.
Further, there is a natural *Schouten bracket* on $`\mathrm{\Theta }_{\text{}}(X)`$:
$$\{,\}:\mathrm{\Theta }_p(X)\times \mathrm{\Theta }_q(X)\mathrm{\Theta }_{p+q1}(X).$$
given by the following formula
$`\{\xi _1\mathrm{}\xi _p,\eta _1\mathrm{}\eta _q\}=`$ (6.1.1)
$`={\displaystyle \underset{i,j=1}{\overset{i=p,j=q}{}}}(1)^{i+j}[\xi _i,\eta _j]\xi _1\mathrm{}\widehat{\xi _i}\mathrm{}\xi _p\eta _1\mathrm{}\widehat{\eta _j}\mathrm{}\eta _q.`$
If $`p=q=1`$, this formula reduces to the usual bracket of vector fields.
To make the bracket $`\{,\}`$ compatible with the gradings, we note that $`(p+1)`$ $`+(q+1)1=(p+q)1`$. Thus, shifting the natural grading on $`\mathrm{\Theta }_{\text{}}(X)`$ by $`1`$, we obtain a new graded strucre, to be denoted $`\mathrm{\Theta }_1(X)`$, such that the bracket (6.1.1) is compatible with this new graded structure.
Below, we summarize the various natural structures on $`\mathrm{\Theta }_{\text{}}(X)`$ and $`\mathrm{\Omega }^{\text{}}(X)`$.
$$\begin{array}{c}\text{Wedge-product of polyvector fields makes }\mathrm{\Theta }_{\text{}}(X)\text{ a super-commu-}\text{}\text{tative algebra.}\hfill \\ \text{The bracket (}\text{6.1.1}\text{) makes }\mathrm{\Theta }_1(X)\text{ a graded Lie super-algebra.}\hfill \\ \text{The contraction operators }i_\pi ,\pi \mathrm{\Theta }_{\text{}}(X),\text{ make }\mathrm{\Omega }^{\text{}}(X)\text{ a graded}\text{}\text{module over }\mathrm{\Theta }_{\text{}}(X)\text{, viewed as a super-commutative algebra.}\hfill \\ \text{The De Rham differential }d:\mathrm{\Omega }^{\text{}}(X)\mathrm{\Omega }^{+1}(X)\text{ is an odd derivation}\text{}\text{of the super-commutative algebra }\mathrm{\Omega }^{\text{}}(X)\text{.}\hfill \end{array}$$
(6.1.2)
Further, one extends the definition of Lie derivative from vector fields to polyvector fields. Specifically, given a polyvector field $`\pi \mathrm{\Theta }_p(X),`$ we define the Lie derivative operator
$$_\pi :=[d,i_\pi ]:\mathrm{\Omega }^{\text{}}(X)\mathrm{\Omega }^{p+1}(X).$$
This formula reduces, in the special case $`p=1`$, to the classical Cartan homotopy formula for the Lie derivative.
Using Lemma 3.3.2, one verifies the following standard identities
$`i_{\pi \phi }=i_\pi i_\phi ,_{\pi \phi }=_\pi i_\phi +(1)^{\mathrm{deg}\pi }i_\pi _\phi ,`$ (6.1.3)
$`[_\pi ,_\phi ]=_{\{\pi ,\phi \}},[i_\pi ,_\phi ]=i_{\{\pi ,\phi \}},[i_\pi ,i_\phi ]=0.`$
### 6.2. Poisson brackets.
Let $`A`$ be an associative (not necessarily commutative) algebra. A (skew-symmetric) Lie bracket $`\{,\}:AAA`$ is said to be a Poisson bracket if, for any $`aA`$, the map $`\{a,\}:AA`$ is a derivation, i.e., the following Leibniz identity holds
$$\{a,bc\}=b\{a,c\}+\{a,b\}c,a,b,cA.$$
(6.2.1)
In view of skew-symmetry, the Leibniz identity says that the bracket $`\{,\}:AAA`$ is a bi-derivation on $`A`$.
In case the associative product on $`A`$ is commutative we will sometimes call $`A`$ a commutative Poisson algebra. In the non-commutative case, we will say that $`A`$ has a NC-Poisson structure.
Let $`A=\mathrm{𝕜}[X]`$ be the coordinate ring of a smooth affine variety. Then, it is easy to show that any bi-derivation on $`\mathrm{𝕜}[X]`$ is given by a bi-vector, that is, there is a regular section $`\pi \mathrm{\Gamma }(X,^2𝒯_X)`$ such that the bi-derivation has the form $`(f,g)dfdg,\pi .`$ Thus, giving a Poisson structure on $`\mathrm{𝕜}[X]`$ amounts to giving a bracket
$$\{f,g\}=dfdg,\pi ,\text{where}\pi \mathrm{\Gamma }(X,^2𝒯_X)\text{is such that}[\pi ,\pi ]=0$$
(6.2.2)
(the condition $`[\pi ,\pi ]=0`$ on the Schouten bracket is equivalent to the Jacobi identity for the Poisson bracket (6.2.2)). In this case we refer to $`\pi `$ as a Poisson bivector.
Fix a smooth variety $`X`$ with a A Poisson bivector $`\pi `$ on a smooth variety $`X`$ gives rise to a canonical Lie algebra structure on $`\mathrm{\Omega }^1(X)`$, the space of 1-forms on $`X`$. The corresponding Lie bracket on $`\mathrm{\Omega }^1(X)`$ is given by the formula
$$[\alpha ,\beta ]:=L_{ı_\pi \alpha }\beta L_{ı_\pi \beta }\alpha d\alpha \beta ,\pi .$$
(6.2.3)
This Lie bracket has the following properties (which uniquely determine the bracket):
$``$ The De Rham differential $`d:A=\mathrm{𝕜}[X]\mathrm{\Omega }^1(X)`$ is a Lie algebra map, i.e.:
$$[da,db]=d(\{a,b\}),a,bA;$$
$``$ The following Leibniz identity holds:
$$[\alpha ,f\beta ]=f[\alpha ,\beta ]+(ı_\pi \beta )f\alpha ,fA,\alpha ,\beta \mathrm{\Omega }^1(X).$$
### 6.3. Gerstenhaber algebras.
The notation of Gerstenhaber algebra is an odd analogue of the notion of Poisson algebra. More explicitly,
###### Definition 6.3.1.
A Gerstenhaber algebra is a graded super-commutative algebra $`G^{\text{}}=_iG^i`$ with a bracket
$$\{,\}:G^p\times G^qG^{p+q1}$$
which makes $`G^{\text{}}`$ a Lie super-algebra so that for every $`aG^{\text{}}`$, the map $`\{a,\}`$ is a super-derivation with respect to the product, i.e., we have
$$\{a,bc\}=\{a,b\}c+(1)^{(\mathrm{deg}a1)\mathrm{deg}b}b\{a,c\}.$$
(6.3.2)
This definition is motivated by the following basic
###### Example 6.3.3.
The Schouten bracket makes $`\mathrm{\Theta }_{\text{}}(X)`$, the graded space of polyvector fields on a smooth manifold $`X`$, a Gerstenhaber algebra.
As another example, let $`X`$ be a Poisson manifold with Poisson bivector $`\pi `$, see (6.2.2). Then, one proves the following
###### Lemma 6.3.4.
The bracket on 1-forms given by formula (6.2.3) extends uniquely to a bracket $`[,]:\mathrm{\Omega }^i(X)\times \mathrm{\Omega }^j(X)\mathrm{\Omega }^{i+j}(X),`$ such that the wedge product and the bracket give $`_{i1}\mathrm{\Omega }^i(X)`$ the structure of Gerstenhaber algebra.∎
###### Remark 6.3.5.
Suppose $`(X,\omega )`$ is a symplectic manifold. The isomorphism $`𝒯_X\stackrel{_{}}{}`$ $`𝒯_X^{},\xi ı_\xi \omega `$ induces a graded algebra isomorphism $`^{\text{}}𝒯_X\stackrel{_{}}{}^{\text{}}𝒯_X^{}`$. Therefore, one can transport the Gerstenhaber algebra structure on $`^{\text{}}𝒯_X`$ given by the Schouten bracket to a Gerstenhaber algebra structure on $`\mathrm{\Omega }^{\text{}}(X)`$. The latter one turns out to be the same as the Gerstenhaber algebra structure of Lemma 6.3.4. $`\mathrm{}`$
###### Remark 6.3.6.
Given a manifold $`X`$ and a bivector $`\pi \mathrm{\Gamma }(X,^2𝒯_X),`$ such that $`[\pi ,\pi ]=0`$ one can also introduce a different Lie bracket $`B:\mathrm{\Omega }^p(X)\times \mathrm{\Omega }^q(X)\mathrm{\Omega }^{p+q2}(X),`$ which has degree $`2`$. For $`\alpha =\alpha _1\mathrm{}\alpha _p\mathrm{\Omega }^p(X)`$ and $`\beta =\beta _1\mathrm{}\beta _q\mathrm{\Omega }^q(X)`$, this new bracket is defined by
$$B(\alpha ,\beta ):=(1)^{i+j}\alpha _i\beta _j,\pi \alpha _1\mathrm{}\widehat{\alpha _i}\mathrm{}\alpha _p\eta _1\mathrm{}\widehat{\eta _j}\mathrm{}\eta _q.$$
$`\mathrm{}`$
Let $`G^{\text{}}`$ be a Gerstenhaber algebra, and $`M^{\text{}}`$ a graded vector space. We say that $`M^{\text{}}`$ is a Gerstenhaber module over $`G^{\text{}}`$ if the square-zero construction, $`G^{\text{}}\mathrm{}M^{\text{}}`$, is equipped with a Gerstenhaber algebra structure such that $`G^{\text{}}`$ is a Gerstenhaber subalgebra in $`G^{\text{}}\mathrm{}M^{\text{}}`$ and, moreover, we have $`M^{\text{}}M^{\text{}}=0=\{M^{\text{}},M^{\text{}}\}`$.
### 6.4. $`\epsilon `$-extension of a Gerstenhaber algebra.
Let $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$ be the ring of dual numbers. Given a Gerstenhaber algebra $`G^{\text{}}=_iG^i,`$ with operations $`()()`$ and $`\{,\}`$, one defines, c.f., \[TT\], a new Gerstenhaber algebra $`G_\epsilon ^{\text{}}=_iG_\epsilon ^i`$, over the ground ring $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$. Specifically, put
$$G_\epsilon ^i:=G^i\epsilon G^{i1},$$
and introduce $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$-bilinear operations $`()_\epsilon ()`$ and $`\{,\}_\epsilon `$, defined for any homogeneous elements $`a,bG`$ by
$$a_\epsilon b:=ab+(1)^{\mathrm{deg}a}\epsilon \{a,b\},\text{and}\{a,b\}_\epsilon :=\{a,b\}.$$
(6.4.1)
###### Remark 6.4.2.
Note that in order to have the new product $`()_\epsilon ()`$ be graded commutative, one has the additional term $`(1)^{\mathrm{deg}a}\epsilon \{a,b\}`$ to be “symmetric” in the graded sense. Such a construction would have been impossible for an ordinary Poisson algebra, where the dot-product is always symmetric and the Poisson bracket is always skew-symmetric, so that the linear combination of the two does not have a defininite symmetry. $`\mathrm{}`$
Applying the ‘$`\epsilon `$-construction’ above to $`\mathrm{\Theta }_{\text{}}(X)`$, the Gerstenhaber algebra of polyvector fields on a manifold $`X`$, one obtains a Gerstenhaber algebra $`\mathrm{\Theta }_{\text{}}(X)_\epsilon .`$
The identities in (6.1.3) are conveniently encoded in the following result.
###### Proposition 6.4.3 (\[TT\]).
The following formulas
$$(\pi +\epsilon \phi )_\epsilon \alpha :=(1)^{\mathrm{deg}\pi }i_\pi \alpha ,\text{and}\{\pi +\epsilon \phi ,\alpha \}_\epsilon :=_\pi \alpha +\epsilon i_\phi \alpha $$
make $`\mathrm{\Omega }^{\text{}}(X)_\epsilon `$ a Gerstenhaber module over $`\mathrm{\Theta }_{\text{}}(X)_\epsilon .`$
### 6.5. Lie algebroids.
Let $`X`$ be an algebraic variety with structure sheaf $`𝒪_X`$ and tangent sheaf $`𝒯_X`$. Let $`𝒜`$ be a coherent sheaf of $`𝒪_X`$-modules equipped with a (not necessarily $`𝒪_X`$-bilinear) Lie bracket $`[,]:𝒜\times 𝒜𝒜`$.
###### Definition 6.5.1.
The data of a sheaf $`𝒜`$ as above and an $`𝒪_X`$-linear map $`\tau :𝒜𝒯_X,v\tau _v,`$ (called anchor map) is said to give a Lie algebroid on $`X`$ if the following holds:
$``$ The map $`\tau `$ is a Lie algebra map;
$``$ We have $`[fv,u]=f[v,u]+\tau _u(f)v`$ for any $`v,u𝒜,f𝒪_X`$.
###### Examples 6.5.2.
(1) The tangent sheaf $`𝒯_X`$ equipped with the standard Lie bracket on vector fields and with the identity anchor map $`\tau =\mathrm{id}:𝒯_X𝒯_X`$ is a Lie algebroid.
(2) Any coherent sheaf of $`𝒪_X`$-modules equipped with an $`𝒪_X`$-bilinear Lie bracket is a Lie algebroid with zero anchor map. In particular, a vector bundle on $`X`$ whose fibers form an algebraic family of Lie algebras is a (locally-free) Lie algebroid on $`X`$.
(3) For any Lie algebroid $`𝒜`$, the kernel of the anchor map $`𝒦:=\mathrm{Ker}[\tau :𝒜𝒯_X]`$ is a sub Lie algebroid in $`𝒜`$. This Lie algebroid $`𝒦`$ is of the type decribed in example (2) above.
(4) Proposition 3.1.3 may be conveniently expressed by saying that, for any algebra $`A`$ with center $`𝖹_A`$, the space $`\mathrm{Der}(A)`$ is (the space of global sections of) a Lie algebroid on $`\mathrm{Spec}𝖹_A`$.
(5) As a special case of (4) we get: the Atiyah algebra of a vector bundle $``$ on $`X`$ has a natural structure of Lie algebroid on $`X`$, since $`𝖹_{_{\mathrm{End}()}}\mathrm{𝕜}[X]`$. In particular, the sheaf $`𝒟_1`$ of (ordinary) first order differential operators on $`X`$ is a Lie algebroid on $`X`$.
(6) A Poisson bivector $`\pi `$ on a smooth variety $`X`$ gives the cotangent sheaf $`𝒯_X^{}`$ a Lie algebroid structure. The corresponding Lie bracket of sections of $`𝒯_X^{}`$ is defined by formula (6.2.3), and the anchor map $`𝒯_X^{}𝒯_X`$ is given by contraction $`\alpha ı_\pi \alpha `$.
Fix a locally free sheaf (a vector bundle) $`𝒜`$ on $`X`$ and an $`𝒪_X`$-linear map $`\tau :𝒜𝒯_X`$. Performing the Symmetric, resp. Exterior, algebra construction to the sheaf $`𝒜`$ and to the map $`\tau `$ we obtain the following graded algebras and graded algebra morphisms:
$$\mathrm{Sym}^{\text{}}𝒜\mathrm{Sym}^{\text{}}𝒯_X,\mathrm{\Lambda }^{\text{}}𝒜\mathrm{\Lambda }^{\text{}}𝒯_X,\mathrm{\Lambda }^{\text{}}𝒜^{}\mathrm{\Lambda }^{\text{}}𝒯_X^{}.$$
A Lie algebroid structure on $`𝒜`$ (with the anchor map $`\tau `$) gives rise to the following additional structures on the graded algebras above:
$``$ A Poisson algebra structure on $`\mathrm{Sym}^{\text{}}𝒜`$ with Poisson bracket given by
$$\{a_1\mathrm{}a_p,b_1\mathrm{}b_q\}:=\underset{i,j=1}{\overset{i=p,j=q}{}}[a_i,b_j]a_1\mathrm{}\widehat{a_i}\mathrm{}a_pb_1\mathrm{}\widehat{b_j}\mathrm{}b_q.$$
(6.5.3)
$``$ A Gerstenhaber algebra structure on $`\mathrm{\Lambda }^{\text{}}𝒜`$ given by formula (6.1.1).
$``$ A differential graded algebra structure on $`\mathrm{\Lambda }^{\text{}}𝒜^{}`$ with differential $`d:\mathrm{\Lambda }^p𝒜^{}\mathrm{\Lambda }^{p+1}𝒜^{}`$ defined as follows.
For $`p=0`$: the differential $`𝒪_X=\mathrm{\Lambda }^0𝒜^{}\mathrm{\Lambda }^1𝒜^{}=𝒜^{}`$ is the composite $`𝒪_X\mathrm{\Omega }_X^1𝒜^{}`$, where the first map is the de Rham differential and the second map is obtained by dualizing the morphism $`\tau :𝒜𝒯_X`$;
For $`p=1`$: the differential $`d:𝒜^{}\mathrm{\Lambda }^2𝒜^{}`$ is given by
$$d\alpha ,ab:=\tau _a\alpha ,b\tau _b\alpha ,a\alpha ,[a,b],\alpha 𝒜^{},a,b𝒜.$$
For $`p>1`$: the differential is given by
$$d(\alpha _1\mathrm{}\alpha _p):=\underset{i=1}{\overset{p}{}}(1)^{i1}\alpha _1\mathrm{}d\alpha _i\mathrm{}\alpha _p$$
Given a Lie algebroid $`𝒜`$, we define a right $`𝒪_X`$-action on $`𝒜`$ by the formula $`uf:=fu+\tau _u(f)`$. Lie algebroid axioms insure that the right action so defined commutes with the natural left $`𝒪_X`$-action on $`𝒜`$ and provides $`𝒜`$ with the structure of a (not necessarily symmetric) $`𝒪_X`$-bimodule. This $`𝒪_X`$-bimodule may be thought of as a coherent sheaf on $`X\times X`$ set-theoretically concentrated on the diagonal $`X_\mathrm{\Delta }X\times X`$.
### 6.6. Gerstenhaber structure on Hochchild cochains.
Let $`G^{\text{}}`$ be a graded associative, not necessarily commutative, algebra equipped with a bracket
$$\{,\}:G^p\times G^qG^{p+q1}.$$
We say that this bracket makes $`G^{\text{}}`$ a noncommutative Gerstenhaber algebra provided it gives $`G^{\text{}}`$ the structure of Lie super-algebra (in particular, the bracket is skew-symmetric/symmetric depending on the parity of its arguments) and the super-Leibniz identity (6.3.2) holds (with the order of factors in the various dot-products being as indicated in (6.3.2)).
Now, let $`A`$ be an associative not necessarily commutative algebra and write $`C^{\text{}}(A,A)=`$$`_i\mathrm{Hom}_\mathrm{𝕜}(A^i,A)`$ for the Hochschild cochain complex of $`A`$.
There is a natural associative (non-commutative) graded algebra structure on $`C^{\text{}}(A,A)`$ given by the so-called cup-product. It is defined, for any $`fC^p(A,A),gC^q(A,A),`$ by the formula
$$fg:a_1,\mathrm{},a_{p+q}f(a_1,\mathrm{},a_p)g(a_{p+1}\mathrm{},a_{p+q}).$$
The Hochschild differential is a super-derivation with respect to the cup-product, that is, we have
$$d(fg)=(df)g+(1)^{\mathrm{deg}f}f(dg).$$
This formula shows that the cup-product of Hochschild cocycles is again a cocycle, and the cup-product of a cocycle and a coboundary is a coboundary. Thus the cup-product descends to a well-defined associative product on $`HH^{\text{}}(A,A)`$, the Hochschild cohomology. It is not difficult to verify that the resulting graded algebra structure on $`HH^{\text{}}(A,A)`$ gets identified, under the isomorphism $`HH^{\text{}}(A,A)\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^{\text{}}(A,A)`$, with the standard Yoneda product on the Ext-groups.
In addition to the cup-product, there is a much deeper structure on Hochschild cochains, revealed by the following result
###### Theorem 6.6.1 (Gerstenhaber).
$`(𝗂)`$There exists a canonical Lie superalgebra structure $`\{,\}:C^p(A,A)\times C^q(A,A)C^{p+q1}(A,A)`$, called the Gerstenhaber bracket.
$`(\mathrm{𝗂𝗂})`$The cup-product and the Gerstenhaber bracket make $`C^{\text{}}(A,A)`$ a noncommutative Gerstenhaber algebra.
###### Proof/Construction.
For $`fC^p(A,A)`$ and $`gC^q(A,A)`$, define
$`(fg)(a_1,`$ $`\mathrm{},a_{p+q1})`$ (6.6.2)
$`={\displaystyle \underset{i=1}{\overset{p}{}}}(1)^{(i1)(q1)}f(a_1,\mathrm{},a_{i1},g(a_i,\mathrm{},a_{i+q}),\mathrm{},a_{p+q1}).`$
Notice that if we regrade the cochains by $`(C^{})^p:=C^{p1}`$, then the degrees are additive in this product. This product is not the cup-product that we have discussed earlier, it is not associative. However, one has the following key identity in $`C^{\text{}}(A,A)`$ due to Gerstenhaber
$$fg(1)^{(\mathrm{deg}f)(\mathrm{deg}g)}gf=d(fg)dfg(1)^{\mathrm{deg}f}fdg.$$
(6.6.3)
Now, following Gerstenhaber, introduce a bracket on Hochschild cohains by the formula
$$\{f,g\}:=fg(1)^{(p1)(q1)}gf,$$
It is straightforward to verify that this gives us a Lie super-algebra structure claimed in part (i) of the Theorem. ∎
The remarkable fact discovered by Gerstenhaber is that noncommutative Gerstenhaber algebra structure on Hochschild cochains gives rise to a commutative Gerstenhaber algebra structure on Hochschild cohomology. That is, one has
###### Proposition 6.6.4.
The cup-product and the Gerstenhaber bracket make$`HH^{\text{}}(A,A)`$ a (super-commutative) Gersenhaber algebra.∎
###### Proof.
We need to show that both the cup-product and Gerstenhaber bracket descend to Hochschild cohomology. The case of cup-product follows from an easy identity
$$d(fg)=dfg+(1)^{\mathrm{deg}f}fdg.$$
Observe that formula (6.6.3) insures that the resulting cup-product on Hochschild cohomology is graded commutative.
Observe further that the super-commutator on the LHS of (6.6.3) is clearly skew super-symmetric with respect to $`fg`$. Hence, super-symmetrization of the LHS, hence of the RHS of (6.6.3), must vanish. This yields
$$0=d\{f,g\}\{df,g\}(1)^{\mathrm{deg}f}\{f,dg\}$$
(6.6.5)
The identity insures that the Gerstenhaber bracket descends to a well-defined bracket on Hochschild cohomology. ∎
More conceptual approach to the Gerstenhaber bracket. We will take $`A`$ to be finite-dimensional so we can take duals, otherwise we would need to use coalgebras. Consider $`T(A^{})`$. The comultiplication $`\delta :A^{}A^{}A^{}`$ extends uniquely to a superderivation on $`T(A^{})`$. Now, he space $`\mathrm{Der}(T^{}A)`$ of all super-derivations on $`T(A^{})`$ is a Lie super-algebra. Now, the Leibniz rule implies that every superderivation $`\theta \mathrm{Der}(T(A^{}))`$ is determined by where it sends each generator. Hence it determines a $`\mathrm{𝕜}`$-linear map $`A^{}T(A^{})`$ (simply see where each basis element of $`A^{}`$ is sent). Clearly the correspondence is reversible, so
$`\mathrm{Der}(T(A^{}))=\mathrm{Hom}_\mathrm{𝕜}(A^{},T(A^{}))`$ $`=(A^{})^{}T(A^{})`$
$`=AT(A^{})=\mathrm{Hom}_\mathrm{𝕜}(TA,A),`$
where we consider the *graded* dual of $`T(A^{})`$. Notice that the $`n^{\text{th}}`$ degree component of $`\mathrm{Hom}_\mathrm{𝕜}(TA,A)`$ is given by $`\mathrm{Hom}_\mathrm{𝕜}(A^n,A)=C^n(A,A)`$. Since $`\mathrm{Der}(T(A^{}))`$ is a Lie super-algebra, we obtain a super-bracket on $`C^n(A,A)`$, which is the Gerstenhaber bracket.
###### Remark 6.6.6.
Recall that $`HH^0(A,A)=𝖹_A,`$ is the center of the algebra $`A`$. For this reason, one may think of the algebra $`HH^{\text{}}(A,A)`$ as a kind of “derived center” of $`A`$. The corollary above confirms that this “derived center” is indeed a commutative algebra.
Recall that the center of $`A`$ may be identified further with $`\mathrm{Hom}(\mathrm{Id}_{A\text{-}\mathrm{𝗆𝗈𝖽}},`$ $`\mathrm{Id}_{A\text{-}\mathrm{𝗆𝗈𝖽}})`$, the endomorphism algebra of the identity functor on the category of left $`A`$-modules. In this spirit, one may think of $`HH^{\text{}}(A,A)`$ to be $`\mathrm{Ext}^{\text{}}(\mathrm{Id}_{A\text{-}\mathrm{𝗆𝗈𝖽}},\mathrm{Id}_{A\text{-}\mathrm{𝗆𝗈𝖽}})`$, the appropriately defined Ext-algebra of the identity functor. $`\mathrm{}`$
### 6.7. Noncommutative Poisson algebras.
In this subsection, given an associative algebra $`A`$, we always write $`[a,b]:=abba`$, for the commutator with respect to the associative product. To avoid confusion, we will use the notation $`\{,\}`$ for the Lie bracket on a Lie algebra.
###### Example 6.7.1.
Let $`A`$ be any associative algebra, and $`t\mathrm{𝕜}`$ a fixed number. For any $`aA`$, the map $`[a,]:AA`$ is a derivation, the inner derivation corresponding to $`a`$. Hence, the map $`t[a,]`$ is also a derivation. Therefore, setting $`\{a,b\}_t:=t[a,b]`$, one obtains a Poisson algebra structure on $`A`$. Note that, for $`t=0`$, the corresponding Poisson bracket vanishes identically. $`\mathrm{}`$
Write $`\mathrm{𝖠𝗌𝗌𝗈𝖼}`$, $`\mathrm{𝖢𝗈𝗆𝗆}`$ and $`\mathrm{𝖫𝗂𝖾}`$ for the operads of associative, commutative, and Lie algebras, respectively, cf. e.g. \[GiK\] for more information about operads. Also, let $`\mathrm{𝖯𝗈𝗂𝗌𝗌}`$, and $`\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}`$ denote respectively the operads of commutative and not necessarily commutative Poisson algebras.
Taking the zero-bracket on a commutative associative algebra, one obtains a functor commutative $`\text{algebras}\text{commutative Poisson algebras}.`$ Further, forgetting the associative product on a commutative Poisson algebra, gives a functor $`\text{Forget}:\text{Poisson algebras}\text{Lie algebras}`$. These functors give rise to the following canonical sequence of morphisms of operads
$$\mathrm{𝖫𝗂𝖾}\mathrm{𝖯𝗈𝗂𝗌𝗌}\mathrm{𝖢𝗈𝗆𝗆}.$$
The forgetful functor $`\text{Forget}:\text{Poisson algebras}\text{Lie algebras}`$ has a left adjoint functor $`{}_{}{}^{}\text{Forget}:\text{Lie algebras}\text{Poisson algebras}.`$ It is given by associating to a Lie algebra $`𝔤`$ the symmetric algebra $`\mathrm{Sym}𝔤`$, equipped with Kirillov-Kostant bracket, cf. section 15. Thus, the commutative Poisson algebga $`\mathrm{Sym}𝔤`$ may be thought of as the (commutative) Poisson envelope of the Lie algebra $`𝔤`$.
Similarly to the above, taking the zero-bracket on an associative algebra, gives a functor $`\text{associative algebras}\text{NC}\text{-}\text{Poisson algebras}.`$ Also, forgetting the associative product on an NC-Poisson algebra, gives a functor $`\text{Forget}:\text{NC}\text{-}\text{Poisson algebras}`$ $`\text{Lie algebras}`$. These functors give rise to the following canonical sequence of morphisms of operads
$$\mathrm{𝖫𝗂𝖾}\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}\mathrm{𝖠𝗌𝗌𝗈𝖼}.$$
###### Question 6.7.2.
Is it true that the operad maps above induce an isomorphism $`\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}/(\mathrm{𝖫𝗂𝖾})\stackrel{_{}}{}\mathrm{𝖠𝗌𝗌𝗈𝖼},`$ where $`(\mathrm{𝖫𝗂𝖾})`$ denotes the operad ideal in $`\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}`$ generated by (the image of) $`\mathrm{𝖫𝗂𝖾}`$.
It is easy to see that the operad $`\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}`$ is quadratic. Therefore, since $`\mathrm{𝖠𝗌𝗌𝗈𝖼}^!=\mathrm{𝖠𝗌𝗌𝗈𝖼}`$ and $`\mathrm{𝖫𝗂𝖾}^!=\mathrm{𝖢𝗈𝗆𝗆},`$ the sequence above induces the dual sequence
$$\mathrm{𝖠𝗌𝗌𝗈𝖼}\mathrm{𝖭𝖢}\text{-}\mathrm{𝖯𝗈𝗂𝗌𝗌}^!\mathrm{𝖢𝗈𝗆𝗆}.$$
We recall that the operad $`\mathrm{𝖯𝗈𝗂𝗌𝗌}`$ is known to be Koszul, see \[GiK\] and also \[MSS\].
###### Theorem 6.7.3.
The operad NC-Poiss is Koszul.
###### Proof.
<sup>2</sup><sup>2</sup>2This proof was kindly communicated to me by Martin Markl.
The theorem would immediately follow from Theorem 4.5 of \[Ml\] once we check that the following distributivity rule
$$[ab,c]=a[b,c]+[a,b]c.$$
(6.7.4)
that ties up the associative product with the Lie bracket is a distributive law. This can be done in either of the following two ways.
(1) It is easy to verify directly that (6.7.4) satisfies the condition of Definition 2.2 of the above mentioned paper, that is, it is indeed a distributive law.
(2) A less direct way is the following. It follows from general theory (see, for example, Theorem 3.2 of \[FM\]) that (6.7.4) is a distributive law if, for any vector space $`V,`$ the free noncommutative Poisson algebra $`NCP(V)`$ on $`V`$ is isomorphic to $`T(L(V)),`$ the free associative algebra generated by the free Lie algebra on $`V.`$
It is obvious that $`NCP(V)=P_{nc}(L(V))`$, the NC-Poisson envelope of the free Lie algebra $`L(V)`$ which you defined in your paper. It is immediate to see that the ideal $`I`$ defined on page 26 of your paper is trivial if $`g=L(V),`$ therefore $`P_{nc}(L(V))=T(L(V))`$ and the result follows.
Therefore the Koszulity of NC-Poiss follows from the same arguments as the Koszulity of the usual commutative $`\mathrm{𝖯𝗈𝗂𝗌𝗌}`$. ∎
The forgetful functor $`\text{Forget}:\text{NC}\text{-}\text{Poisson algebras}\text{Lie algebras}`$ has a left adjoint functor $`{}_{}{}^{}\text{Forget}:\text{Lie algebras}\text{NC}\text{-}\text{Poisson algebras}.`$ In other words, given a Lie algebra $`𝔤`$, there is a uniquely defined NC-Poisson algebra $`𝒫_{\mathrm{nc}}(𝔤)`$, called the NC-Poisson envelope of $`𝔤`$, that comes equipped with a Lie algebra map $`\iota :𝔤𝒫_{\mathrm{nc}}(𝔤)`$ (with respect to the Poisson bracket on $`𝒫_{\mathrm{nc}}(𝔤)`$) and such that the following universal property holds:
$``$ For any NC-Poisson algebra $`P`$ and a Lie algebra map $`\varphi :𝔤P`$ there exists a unique morphism $`𝒫_{\mathrm{nc}}(\varphi ):𝒫_{\mathrm{nc}}(𝔤)P,`$ of NC-Poisson algebras, that makes the following diagram commute
(6.7.5)
The construction of universal universal NC-Poisson envelope $`𝒫_{\mathrm{nc}}(𝔤)`$ is due to Th. Voronov in \[Vo\]. It is based on the following
###### Lemma 6.7.6 (\[Vo\]).
Let $`P`$ be a not necessarily commutative Poisson algebra with Poisson bracket $`\{,\}`$. Then, for any $`a,b,c,d,uP`$, one has
$$\{a,b\}u[c,d]=[a,b]u\{c,d\},$$
where $`[x,y]=xyyx`$ stands for the commutator for the associative product.
###### Proof.
Consider the expression $`\{ac,bud\}`$. Compute this in two different ways, first by applying the Leibniz rule for $`\{ac,\}`$, and second, by applying the Leibniz rule for $`\{,bud\}`$. ∎
Now, given a Lie algebra $`𝔤`$ with Lie bracket $`\{,\}`$, Voronov considers a two-sided ideal $`IT(𝔤)`$, in the tensor algebra of the vector space $`𝔤`$, generated by the elements indicated below
$$I:=\{a,b\}u[c,d][a,b]u\{c,d\}_{a,b,c,d𝔤,uT(𝔤)},$$
where $`[,]`$ stands for the commutator in the associative algebra $`T(𝔤)`$ and $`\{,\}`$ stands for the Lie bracket in $`𝔤`$.
###### Theorem 6.7.7 (\[Vo\]).
The assignment
$`(a_1\mathrm{}a_k)`$ $`\times (b_1\mathrm{}b_l)\{a_1\mathrm{}a_k,b_1\mathrm{}b_l\}:=`$
$`{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{s=1}{\overset{l}{}}}b_1\mathrm{}b_{s1}a_1\mathrm{}a_{s1}\{a_r,b_s\}a_{r+1}`$
$`\mathrm{}a_kb_{s+1}\mathrm{}b_l`$
gives rise to a well-defined noncommutative Poisson structure on the associative algebra $`T(𝔤)/I`$. The universal property (6.7.5) holds for the NC-Poisson algebra $`𝒫_{\mathrm{nc}}(𝔤):=T(𝔤)/I`$ thus defined. ∎
###### Example 6.7.8.
Let $`𝔤`$ be a Lie algebra and $`𝒰𝔤`$ its enveloping algebra. We may consider the associative algebra $`𝒰𝔤`$ as a NC-Poisson algebra with Poisson bracket $`\{,\}=[,]`$. Therefore, the canonical Lie algebra map $`𝔤{}_{}{}^{}_{}^{}𝒰𝔤`$ gives rise, via the universality property of $`𝒫_{\mathrm{nc}}(𝔤)`$, to a natural morphism $`𝒫_{\mathrm{nc}}(𝔤)𝒰𝔤`$ of NC-Poisson algebras. The latter morphism is easily seen to be surjective. Thus, $`𝒰𝔤`$ is a quotient of $`𝒫_{\mathrm{nc}}(𝔤)`$.
A similar argument shows that the commutative Poisson algebra, $`\mathrm{Sym}𝔤`$, is also a quotient of $`𝒫_{\mathrm{nc}}(𝔤)`$. $`\mathrm{}`$
## 7. Deformation quantization
### 7.1. Star products.
Let $`A`$ be an associative $`\mathrm{𝕜}`$-algebra. We wish to define a twisted “product” on $`A`$. For $`a,bA`$, define
$$a_tb=ab+t\beta _1(a,b)+t^2\beta _2(a,b)+\mathrm{}A[[t]]$$
for maps $`\beta _j:A\times AA`$. We wish for this map to be “associative,” and ask what conditions this places on the $`\beta _j`$’s. Define
$$a_t(b_tc)=(a_t(bc))+[a_t\beta _1(b,c)]t+[a_t\beta _2(b,c)]t^2+\mathrm{},$$
and similarly for $`(a_tb)_tc`$. Clearly there is no issue for the constant term. Consider the coefficients of the $`t`$ term. From $`(a_tb)_tc`$ we obtain $`\beta _1(ab,c)+\beta _1(a,b)c`$, and from $`a_t(b_tc)`$ we obtain $`a\beta _1(b,c)+\beta _1(a,bc)`$. Equating these two we see that $`\beta _1:AAA`$ must be a Hochschild cocycle in $`C^2(A,A)`$. In particular, we begin with a commutative algebra $`A`$ (which we can think of as the affine coordinate ring of a variety $`X`$), and we will require the deformed products $`_t`$ to be associative only. It is interesting to examine how noncommutative $`_t`$ is, so we define
$$[a_tb]=a_tbb_ta=t[\beta _1(a,b)\beta _1(b,a)]+O(t^2).$$
Define $`\{a,b\}=\beta _1(a,b)\beta _1(b,a)`$, which is clearly a skew-symmetric product on $`A`$. So, we can write
$$[a_tb]=t\{a,b\}+O(t^2).$$
Now, recall that commutator in any associative algebra satisfies the Jacobi identity. So, for $`(A[[t]],_t)`$, we have the identities
$`[a_tb_tc]`$ $`=[a_tb]_tc+b_t[a_tc]`$
$`[[a_tb]_tc]`$ $`=[[b_tc]_ta]+[[c_ta]_tb].`$
Now, we insert the formula $`[a_tb]=t\{a,b\}+O(t^2)`$ into Leibniz’s rule and take only the first order terms. Since every term of the bracket on $`A[[t]]`$ introduces a factor of $`t`$, we need only take the first order terms of each term. For example, we find that
$$[a_tb_tc]=[a_tbc]+O(t^2)=t\{a,bc\}+O(t^2).$$
The other brackets simplify similarly, so we obtain
$$\{a,bc\}=\{a,b\}c+b\{a,c\}.$$
Similarly, we examine Jacobi’s rule and take the coefficient of the lowest power of $`t`$, namely $`t^2`$. Arguing similarly, we find that
$$\{\{a,b\},c\}=\{\{b,c\},a\}+\{\{c,a\},b\}.$$
So, $`\{,\}`$ is a Lie bracket on $`A`$, and combined with the associative product on $`A`$ it satisfies the Leibniz rule. This defines the structure of a *Poisson algebra* on $`A`$. So, we can reformulate the problem of classifying deformation quantizations as a problem regarding Poisson algebra structures that can be placed on $`A`$.
### 7.2.
The calculations above have one large deficiency: if $`\beta _1`$ is identically zero then the bracket $`\{,\}`$ will be as well. So, we repeat the above reasoning for the general case of formal deformations.
Thus, let $`A`$ be a commutative associative algebra, and $`\stackrel{~}{A}`$ be a formal flat deformation of $`A`$, i.e., a topologically free (not necessarily commutative) associative $`\mathrm{𝕜}[[t]]`$-algebra equiped with an algebra isomorphism $`\stackrel{~}{A}/t\stackrel{~}{A}A.`$
Since $`A`$ is commutative, this implies that for each $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}`$, $`\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{b}\stackrel{~}{a}t\stackrel{~}{A}`$. For each pair $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}`$, define a natural number $`m(\stackrel{~}{a},\stackrel{~}{b})`$ to be the maximum integer such that
$$[\stackrel{~}{a},\stackrel{~}{b}]t^{m(\stackrel{~}{a},\stackrel{~}{b})}\stackrel{~}{A}.$$
If $`[\stackrel{~}{a},\stackrel{~}{b}]`$ is contained in every $`t^i\stackrel{~}{A}`$, we set $`m(\stackrel{~}{a},\stackrel{~}{b})=\mathrm{}`$. Let $`N=\mathrm{min}\{m(\stackrel{~}{a},\stackrel{~}{b})|`$ $`\stackrel{~}{a},\stackrel{~}{b}A\},`$ which is necessarily greater than or equal to $`1`$. Now, by Krull’s theorem, we know that $`_{i=1}^{\mathrm{}}t^i\stackrel{~}{A}=0,`$ so if $`N=\mathrm{}`$ this means that every $`[\stackrel{~}{a},\stackrel{~}{b}]t^iA`$ for each $`i1`$, hence $`[\stackrel{~}{a},\stackrel{~}{b}]=0`$. This would force $`\stackrel{~}{A}`$ to be commutative.
Since we are interested in noncommutative deformations, we will only consider the case $`N<\mathrm{}`$. For any $`a,bA`$, choose $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}`$ such that
$$\stackrel{~}{a}modt\stackrel{~}{A}=a,\stackrel{~}{b}modt\stackrel{~}{A}=b.$$
(7.2.1)
Then by the definition of $`N`$, we know that $`[\stackrel{~}{a},\stackrel{~}{b}]t^N\stackrel{~}{A}`$, so $`t^N(\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{b}\stackrel{~}{a})`$ is a well-defined element of $`\stackrel{~}{A}`$. We define
$$\{a,b\}=t^N[\stackrel{~}{a},\stackrel{~}{b}]modt\stackrel{~}{A}.$$
(7.2.2)
It is then easy to check that $`\{a,b\}`$ is independent of the choice of $`\stackrel{~}{a}`$ and $`\stackrel{~}{b}`$, and that $`\{,\}`$ gives $`A`$ the structure of a Poisson algebra.
### 7.3. A Lie algebra associated to a deformation.
Let now $`A`$ be an associative, not necessarily commutative, algebra and write $`𝖹_A`$ for the center of $`A`$.
Hayashi observed, see \[Hay\], that given $`\stackrel{~}{A}`$, a formal flat deformation of $`A`$, the construction of the previous subsection can be adapted to produce a Poisson structure on $`𝖹_A`$, a commutative algebra. In more detail, let $`a,b𝖹_A`$, and choose $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}`$ as in (7.2.1). Then, put
$$\{a,b\}:=\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}]modt\stackrel{~}{A}$$
It is straightforward to verify that
* $`\{a,b\}𝖹_A`$;
* $`\{a,b\}`$ is independent of the choices of $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}`$;
* The assignment $`a,b\{a,b\},`$ makes $`𝖹_A`$ a Poisson algebra.
###### Remark 7.3.1.
The reader should be alerted that, in contrast with the case of §7.2, the Poisson bracket on $`𝖹_A`$ thus defined may turn out to be identically zero. The reason is that unlike formula (7.2.2), we now divide by the first power of $`t`$ rather than the $`N^{th}`$ power (one can check that dividing by the $`N^{th}`$ power does not give rise to a well-defined bracket on $`𝖹_A`$). Now, given $`a,b𝖹_A`$, put $`N(a,b):=\mathrm{min}_{\stackrel{~}{a},\stackrel{~}{b}A}m(\stackrel{~}{a},\stackrel{~}{b}),`$ where the minimum is taken over all possible lifts of $`a`$ and $`b`$. Then, it is clear from the construction that if $`N(a,b)>1`$ then the elements $`a,b`$ Poisson commute. Therefore, the Poisson bracket on $`𝖹_A`$ vanishes whenever one has $`N(a,b)>1,a,b𝖹_A`$. $`\mathrm{}`$
The Hayashi construction has been further refined in \[BD\] as follows. We introduce the following vector subspace
$$\stackrel{~}{A}^{}:=\{\stackrel{~}{a}|\stackrel{~}{a}\stackrel{~}{A}\text{such that}\stackrel{~}{a}modt\stackrel{~}{A}𝖹_A\}.$$
We put $`𝒜:=\stackrel{~}{A}^{}/t\stackrel{~}{A}^{}`$. Observe next that, for any $`\stackrel{~}{a},\stackrel{~}{b}\stackrel{~}{A}^{},`$ we have $`[\stackrel{~}{a},\stackrel{~}{b}]t\stackrel{~}{A},`$ We claim that the element $`\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}]`$ belongs to $`\stackrel{~}{A}^{}`$. Indeed, for any $`\stackrel{~}{c}\stackrel{~}{A}`$, using Jacobi identity in $`\stackrel{~}{A}`$, we find
$`[\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}],\stackrel{~}{c}]=\frac{1}{t}[[\stackrel{~}{a},\stackrel{~}{c}],\stackrel{~}{b}]+\frac{1}{t}[\stackrel{~}{a},[\stackrel{~}{b},\stackrel{~}{c}]]`$ $`\frac{1}{t}[t\stackrel{~}{A},\stackrel{~}{b}]+\frac{1}{t}[\stackrel{~}{a},t\stackrel{~}{A}]`$
$`=[\stackrel{~}{A},\stackrel{~}{b}]+[\stackrel{~}{a},\stackrel{~}{A}]t\stackrel{~}{A}+t\stackrel{~}{A}=t\stackrel{~}{A}.`$
Hence, the expression in the LHS of the top line vanishes modulo $`t`$. Thus, we have proved that $`\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}]\stackrel{~}{A}^{}`$. Further, it is easy to see that the class $`\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}]modt\stackrel{~}{A}^{}`$ depends only on the classes $`\stackrel{~}{a}modt\stackrel{~}{A}^{}`$ and $`\stackrel{~}{b}modt\stackrel{~}{A}^{}`$.
This way, one proves the following
###### Proposition 7.3.2.
$`(𝗂)`$The assignment $`\stackrel{~}{a},\stackrel{~}{b}\frac{1}{t}[\stackrel{~}{a},\stackrel{~}{b}]modt\stackrel{~}{A}^{}`$ induces a Lie algebra structure on $`𝒜`$.
$`(\mathrm{𝗂𝗂})`$The projection $`𝒜𝖹_A,\stackrel{~}{a}\stackrel{~}{a}modt\stackrel{~}{A},`$ is a Lie algebra map (with respect to the Hayashi bracket on $`𝖹_A`$), that gives rise to a Lie algebra extension:
$$0A/𝖹_A𝒜𝖹_A\mathrm{\hspace{0.17em}0}.$$
$`(\mathrm{𝗂𝗂𝗂})`$The adjoint action of $`\stackrel{~}{A}^{}`$ on $`\stackrel{~}{A}`$ descends to a well-defined Lie algebra action of $`𝒜`$ on $`A`$, i.e., gives a Lie algebra map $`𝒜\mathrm{Der}(A),x_x`$; moreover, for any $`xA/𝖹_A𝒜`$ and $`aA`$, we have $`_x(a)=[\stackrel{~}{x},\stackrel{~}{a}]modt\stackrel{~}{A},`$ where $`\stackrel{~}{x},\stackrel{~}{a}\stackrel{~}{A}`$ are any lifts of $`x`$ and $`a`$, respectively.∎
### 7.4. Example: deformations of the algebra $`\mathrm{End}()`$.
Let $``$ be an algebraic vector bundle on a affine variety $`X`$, and $`A=\mathrm{End}()`$ the endomorphism algebra of this vector bundle. The center of this algebra is $`𝖹_A=\mathrm{𝕜}[X]\mathrm{End}()`$, the subalgebra of ‘scalar’ endomorphisms.
Now let $`\stackrel{~}{A}`$ be a formal deformation of the algebra $`A`$. By Hayashi construction, this deformation gives rise to a Poisson bracket on $`𝖹_A=\mathrm{𝕜}[X]`$, thus makes $`X`$ a Poisson variety. In particular, assigning to each $`z\mathrm{𝕜}[X]`$ the derivation $`\{z,\}`$ yields a Lie algebra map $`\xi :\mathrm{𝕜}[X]\mathrm{Der}(\mathrm{𝕜}[X])=𝒯(X).`$
Recall the Atiyah algebra $`𝒜()=\mathrm{Der}\mathrm{End}()`$ of first order differential operators on $``$ with scalar principal symbol, introduced in §3.6. The following result gives an explicit description of the Lie algebra $`𝒜`$ of Proposition 7.3.2 in the special case at hand.
###### Proposition 7.4.1.
There is a natural Lie algebra map $`\xi __𝒜:𝒜𝒜()`$ making the second row of the diagram below the pull-back (via $`\xi `$) of the standard Lie algebra extension in the first row of the diagram
###### Proof.
By definition, the Lie algebra $`𝒜()`$Theorem 3.6.4, is the Lie algebra of derivations of the endomorphism algebra, which is our algebra $`A`$. Therefore, producing a map $`\xi __𝒜:𝒜𝒜()`$ amounts to constructing a Lie algebra map $`𝒜\mathrm{Der}A`$. But the latter map has been already constructed in part (iii) of Proposition 7.3.2. It is straightforward to verify, using part (ii) of that Proposition that the map $`\xi __𝒜`$ arising in this way indices the map $`\xi `$, the vertical arrow on the right of the diagram above. ∎
### 7.5.
###### Definition 7.5.1.
Let $`B`$ be a $`\mathrm{𝕜}[[t]]`$-algebra with $`t`$-linear associative, not necessarily commutative product $`_t`$ and a $`t`$-linear Lie bracket $`[,]_t`$ such that
* $`[b,]_t`$ is a derivation with respect to $`_t`$; and
* $`a_tbb_ta=t[a,b]_t`$.
Then we say that $`B`$ is a *$`t`$-algebra*.
###### Example 7.5.2.
$`(𝗂)`$If $`B`$ is a flat deformation $`\widehat{A}=A[[t]]`$ of an algebra $`A`$ such that $`A`$ is commutative, then $`[a,b]_t=\frac{a_tbb_ta}{t}`$ makes $`\widehat{A}`$ into a $`t`$-algebra.
$`(\mathrm{𝗂𝗂})`$Let $`A`$ be a Poisson algebra with Poisson bracket $`\{,\}`$. Take $`_t`$ to be the given multiplication and $`[,]_t=\{,\}`$, and let $`t`$ act by zero on $`A`$. Then $`B`$ is a $`t`$-algebra. $`\mathrm{}`$
We let $`𝒫_A`$ denote the moduli space of flat deformations as $`t`$-algebra of a given Poisson algebra $`A`$ (recall that we assume that all Poisson algebras are commutative). The quantization problem can be summarized as an attempt to understand this space. Recall the moduli space $`_A`$ of deformations of $`A`$ as an associative algebra. Giving a Poisson structure on $`A`$ is equivalent to giving an element of the second Hochschild cohomology of $`A`$, which we identify with the tangent space of $`_A`$ at $`A`$. That is, an element of $`𝒫_A`$ corresponds to an element of $`T_A_A`$. So, we can view $`𝒫_A`$ is the germ of the blow-up of $`_A`$ at the point given by $`A`$ and tangent direction specified by the given Poisson structure on $`A`$.
### 7.6. Moyal-Weyl quantization.
Let $`(V,\omega )`$ be a symplectic manifold (recall that this means that $`\omega \mathrm{\Omega }^2V`$ is nondegenerate and closed, $`d\omega =0`$). Let $`A=\mathrm{𝕜}[V]`$.
The symplectic form $`\omega `$ gives rise to a Poisson bracket in the following way. Since $`\omega `$ is nondegenerate, it induces an isomorphism of bundles $`TVT^{}V`$. Given some section $`\alpha \mathrm{\Gamma }(V,T^{}V)`$, this isomorphism provides a vector field $`\xi _\alpha \mathrm{\Gamma }(V,TV)`$. So, if $`f\mathrm{𝕜}[V]`$ is a regular function, $`df`$ is a one-form, that is, $`df\mathrm{\Gamma }(V,T^{}V)`$. This gives rise to a vector field $`\xi _{df}\mathrm{\Gamma }(V,TV)`$. So, if we have two functions $`f,g\mathrm{𝕜}[V]`$, we define
$$\{f,g\}=\xi _{df}g.$$
This is clearly skew-symmetric.
We consider a special case where $`V`$ is a symplectic vector space of dimension $`dimV=2n`$, so $`\omega ^2V^{}`$ is a nondegenerate skew-symmetric bilinear form $`\omega :VV\mathrm{𝕜}`$. Recall that given such a symplectic form $`\omega `$ on $`V`$, we can find coordinates $`p_1,\mathrm{},p_n,q_1,\mathrm{},q_n`$ on $`V`$ such that in these coordinates $`\omega `$ is the standard symplectic form, i.e.,
$$\omega =\underset{i=1}{\overset{n}{}}dp_idq_i.$$
Then if $`f`$ is a regular function on $`V`$, we find that the vector field associated to $`df`$ is given by
$$\xi _{df}=\underset{i=1}{\overset{n}{}}\left[\frac{f}{p_i}\frac{}{q_i}\frac{f}{q_i}\frac{}{p_i}\right],$$
so that
$$\{f,g\}=\underset{i=1}{\overset{n}{}}\left[\frac{f}{p_i}\frac{g}{q_i}\frac{f}{q_i}\frac{g}{p_i}\right].$$
Let $`\pi ^2V`$ be the bivector obtained by transporting the 2-form $`\omega `$ via the isomorphism $`^2V^{}\stackrel{_{}}{}^2V`$ induced by the symplectic form $`\omega `$. Using $`\pi `$ we can rewrite the above Poisson bracket on $`\mathrm{𝕜}[V]`$ as $`f,g\{f,g\}:=dfdg,\pi `$ on $`\mathrm{𝕜}[V],`$ the polynomial algebra on $`V`$. The usual commutative product $`m:\mathrm{𝕜}[V]\mathrm{𝕜}[V]\mathrm{𝕜}[V]`$ and the Poisson bracket $`\{,\}`$ make $`\mathrm{𝕜}[V]`$ a Poisson algebra.
This Poisson algebra has a distinguished Moyal-Weyl quantization (\[Mo\], see also \[CP\]). This is an associative star-product depending on a formal quantization parameter $`t`$, defined by the formula
$$f_tg:=m{}_{^{^{}}}{}^{}e_{}^{\frac{1}{2}t\pi }(fg)\mathrm{𝕜}[V][t],f,g\mathrm{𝕜}[V][t].$$
(7.6.1)
To explain the meaning of this formula, view elements of $`\mathrm{Sym}V`$ as constant-coefficient differential operators on $`V`$. Hence, an element of $`\mathrm{Sym}V\mathrm{Sym}V`$ acts as a constant-coefficient differential operator on the algebra $`\mathrm{𝕜}[V]\mathrm{𝕜}[V]=\mathrm{𝕜}[V\times V].`$ Now, identify $`^2V`$ with the subspace of skew-symmetric tensors in $`VV`$. This way, the bivector $`\pi ^2VVV`$ becomes a second order constant-coefficient differential operator $`\pi :\mathrm{𝕜}[V]\mathrm{𝕜}[V]\mathrm{𝕜}[V]\mathrm{𝕜}[V].`$ Further, it is clear that for any element $`fg\mathrm{𝕜}[V]\mathrm{𝕜}[V]`$ of total degree $`N`$, all terms with $`d>N`$ in the infinite sum $`e^{t\pi }(fg)=_{d=0}^{\mathrm{}}\frac{t^d}{d!}\pi ^d(fg)`$ vanish, so the sum makes sense. Thus, the symbol $`m{}_{^{^{}}}{}^{}e_{}^{t\pi }`$ in the right-hand side of formula (7.6.1) stands for the composition
$$\mathrm{𝕜}[V]\mathrm{𝕜}[V]\stackrel{e^{t\pi }}{}\mathrm{𝕜}[V]\mathrm{𝕜}[V]\mathrm{𝕜}[t]\stackrel{m\mathrm{Id}_{\mathrm{𝕜}[t]}}{}\mathrm{𝕜}[V]\mathrm{𝕜}[t],$$
where $`e^{t\pi }`$ is an infinite-order formal differential operator.
In down-to-earth terms, choose coordinates $`x_1,\mathrm{},x_n,y_1,\mathrm{},y_n`$ on $`V`$ such that the bivector $`\pi `$, resp., the Poisson bracket $`\{,\}`$, takes the canonical form
$$\pi =\underset{i}{}\frac{}{x_i}\frac{}{y_i}\frac{}{y_i}\frac{}{x_i},\text{resp.,}\{f,g\}=\underset{i}{}\frac{f}{x_i}\frac{g}{y_i}\frac{f}{y_i}\frac{g}{x_i}.$$
(7.6.2)
Thus, in canonical coordinates $`x=(x_1,\mathrm{},x_n),y=(y_1,\mathrm{},y_n),`$ formula (7.6.1) for the Moyal product reads
$`(f_tg)(x,y)`$ $`={\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^d}{d!}}\left({\displaystyle \underset{i}{}}{\displaystyle \frac{}{x_i^{}}}{\displaystyle \frac{}{y_i^{\prime \prime }}}{\displaystyle \frac{}{y_i^{}}}{\displaystyle \frac{}{x_i^{\prime \prime }}}\right)^df(x^{},y^{})g(x^{\prime \prime },y^{\prime \prime })|_{\genfrac{}{}{0pt}{}{x^{}=x^{\prime \prime }=x}{y^{}=y^{\prime \prime }=y}}`$
$`={\displaystyle \underset{𝐣,𝐥_0^n}{}}(1)^{𝐥|}{\displaystyle \frac{t^{|𝐣|+|𝐥|}}{𝐣!𝐥!}}{\displaystyle \frac{^{𝐣+𝐥}f(x,y)}{x^𝐣y^𝐥}}{\displaystyle \frac{^{𝐣+𝐥}g(x,y)}{y^𝐣x^𝐥}},`$ (7.6.3)
where for $`𝐣=(j_1,\mathrm{},j_n)_0^n`$ we put $`|𝐣|=_ij_i`$ and given $`𝐣,𝐥_0^n,`$ write
$$\frac{1}{𝐣!𝐥!}\frac{^{𝐣+𝐥}}{x^𝐣y^𝐥}:=\frac{1}{j_1!\mathrm{}j_n!l_1!\mathrm{}l_n!}\frac{^{|𝐣|+|𝐥|}}{x_1^{j_1}\mathrm{}x_n^{j_n}y_1^{l_1}\mathrm{}y_n^{l_n}}.$$
### 7.7. Weyl Algebra
A more conceptual approach to formulas (7.6.1)–(7.6) is obtained by introducing the Weyl algebra $`A_t(V)`$. This is a $`\mathrm{𝕜}[t]`$-algebra defined by the quotient
$$A_t(V):=(TV^{})[t]/I(uu^{}u^{}ut\pi ,uu^{})_{u,u^{}V^{}},$$
where $`TV^{}`$ denotes the tensor algebra of the vector space $`V^{}`$, and $`I(\mathrm{})`$ denotes the two-sided ideal generated by the indicated set.
For instance, if $`dimV=2`$ and $`p,q`$ are canonical coordinates on $`V`$, then we have
$$A_t=A_t(V)=\mathrm{𝕜}p,q/(pqqp=t).$$
By scaling, there are essentially only two different cases for $`t`$, namely $`t=0`$ or $`t0`$. However, it will be convenient to have a continuous parameter. Also, the algebra $`A_1`$ will be called simply the Weyl algebra. Notice that even though a monomial does not have a well-defined degree (for example, $`pq=qpt`$, the left-hand side has degree two and the right is not even homogeneous if $`t0`$), we can see that any homogeneous monomial has a highest degree in which it may be expressed. We filter $`A_t(V)`$ by this highest degree for each monomial.
Now, a version of the Poincaré-Birkhoff-Witt theorem says that the natural symmetrization map yields a $`\mathrm{𝕜}[t]`$-linear bijection $`\sigma _W:\mathrm{𝕜}[V][t]\stackrel{_{}}{}A_t(V)`$. In the special case of a 2-dimensional space $`V`$, the linear bijection $`\sigma _t:\mathrm{𝕜}[p,q]A_t`$ is given by sending $`p^mq^n\mathrm{𝕜}[p,q]`$ to the average of all possible permutations of the $`m`$ $`p`$’s and $`n`$ $`q`$’s. This becomes the identity when we pass to the graded case. That is, if $`\phi `$ is a homogeneous polynomial, then the principal symbol of $`\sigma _t(\phi )`$ in $`\mathrm{gr}A_t`$ is precisely $`\phi `$.
Thus, transporting the multiplication map in the Weyl algebra $`A_t(V)`$ via the bijection $`\sigma _W,`$ one obtains an associative product
$$\mathrm{𝕜}[V][t]_{\mathrm{𝕜}[t]}\mathrm{𝕜}[V][t]\mathrm{𝕜}[V][t],fg\sigma _W^1(\sigma _W(f)\sigma _W(g)).$$
It is known that this associative product is equal to the one given by formulas (7.6.1)–(7.6).
The easiest way to see the last assertion is to argue heuristically as follows. We assume $`dimV=2,`$ for simplicity.
First of all, one verifies that $`fg=\sigma _W^1(\sigma _W(f)\sigma _W(g))`$ admits an expansion with the initial term $`fg`$ plus terms of higher order in $`t`$ whose coefficients are all composed of differential operators with polynomial coefficients applied to $`f`$ and $`g`$. Given this claim, we can then extend the product formula to all smooth functions, not just polynomials. In particular, we choose to take
$$\phi (p,q)=e^{\alpha p+\beta q}\text{and}g(p,q)=e^{\gamma p+\delta q},$$
where none of $`\alpha ,\beta ,\gamma ,\delta `$ equal one another and none of them equal zero or one. Then a differential operator on $`\phi `$ and $`g`$ is determined completely by its action on the above $`\phi `$ and $`g`$.
So, we have essentially moved from the problem of computing the product in $`A_t`$ to computing its logarithm. In the case of Lie algebras, we can invoke the Campbell-Hausdorff theorem to obtain some partial information. In this case, we know that both $`e^x`$ and $`e^y`$ are elements of an associated Lie group, hence so is $`e^x_te^y`$. So, we should be able to express $`e^x_te^y=e^{z_t(x,y)}`$ for some function $`z_t(x,y)`$. In particular, we always know the first two terms:
$$z_t(x,y)=(x+y)+\frac{t}{2}[x,y]+O(t^2).$$
Notice that the first term is $`x+y`$. Here is a little exercise: Check that the $`\beta _j`$’s (the coefficient $`t^j`$) can be expressed as a differential operator as we claim if and only if the above expansion starts with $`x+y`$.
We now specialize the general Campbell-Hausdorff theorem in the case of the 3-dimensional Heisenberg Lie algebra $`𝔥`$ with basis $`\{x,y,c\}`$ and commutation relations
$$[x,y]=c,[x,c]=[y,c]=0,$$
in particular, $`c`$ is central.
It is clear that the Weyl algebra $`A_t`$ is a quotient of the enveloping algebra of $`𝔥`$, specifically, we have $`A_t=𝒰(𝔥)/(c=t)`$. So, we can apply the discussion regarding the Campbell-Hausdorff theorem to $`𝔥`$. Notice that since every bracket is central, there can be no nontrivial interated brackets. So, we can check that the Campbell-Hausdorff theorem yields the simple relation
$$z_t(x,y)=(x+y)+\frac{t}{2}[x,y].$$
We find that
$$e^x_te^y=e^{x+y+\frac{t}{2}[x,y]}.$$
We can generalize the above to the case $`dimV>2`$. This will allow us to better understand some of the symmetries of the situation. The Heisenberg Lie algebra $`𝔥`$ for dimension $`V`$ is given by a central extension of $`\mathrm{𝕜}`$ by $`V`$ treated as an abelian Lie group (i.e., $`[v,w]=0`$ for all $`v,wV`$). That is,
$$\text{}.$$
The above exact sequence splits as a vector space, so we may write $`𝔥=\mathrm{𝕜}c+V`$, where $`c`$ is a non-zero element of $`\mathrm{𝕜}`$. To determine the Lie bracket on $`𝔥`$, we need only compute $`[x,y]`$ for $`x,yV`$ (since $`c`$ is central). We set $`[x,y]=\omega (x,y)c`$ for all $`x,yV`$. Notice that this is invariant under the action of the symplectic group $`Sp(V,\omega )`$ on $`V`$. In our situation, we obtain the formula
$$e^x_te^y=e^{x+y+\frac{t}{2}\omega (x,y)}=e^{\frac{t}{2}\omega (x,y)}e^{x+y},$$
which is a simple scalar correction term. If we now apply this calculation for $`\phi (p,q)=e^{\alpha p+\beta q}`$ and $`\psi (p,q)=e^{\gamma p+\delta q}`$, we find
$`\phi _t\psi (p,q)`$ $`=\mathrm{exp}\left[{\displaystyle \frac{t}{2}}\omega (\alpha p+\beta q,\gamma p+\delta q)\right]\phi \psi `$
$`=\mathrm{exp}\left\{{\displaystyle \frac{t}{2}}[\alpha \gamma \omega (p,p)+\alpha \delta \omega (p,q)+\beta \gamma \omega (q,p)+\beta \delta \omega (q,q)]\right\}\phi \psi `$
$`=\mathrm{exp}\left[{\displaystyle \frac{t}{2}}(\alpha \delta \beta \gamma )\right]\phi \psi `$
$`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^n}{2^nn!}}(\alpha \delta \beta \gamma )^n\phi (p,q)\psi (p,q).`$
Of course, $`\alpha `$ corresponds to differentiating $`\phi `$ with respect to its first argument, $`\beta `$ is differentiation of $`\phi `$ with respect to the second, etc. So, for any $`f,g\mathrm{𝕜}[p,q]`$, we obtain
$$f_tg=\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^n}{2^nn!}\left(\frac{}{p^{}}\frac{}{q^{\prime \prime }}\frac{}{p^{\prime \prime }}\frac{}{q^{}}\right)\phi (p^{},q^{})\psi (p^{\prime \prime },q^{\prime \prime })|_{\begin{array}{c}p^{}=p^{\prime \prime }=p\\ q^{}=q^{\prime \prime }=q\end{array}}.$$
The right hand side here is exactly the same expression as given by formula (7.6).
## 8. Kähler differentials
### 8.1.
In order to make some connections to geometry, we are going to discuss a construction for commutative algebras. So, until further notice, $`A`$ is a commutative algebra. Kähler differentials for $`A`$ are $`A`$-linear combinations of the symbols $`db`$, where $`bA`$ and $`d(ab)=adb+bda`$. More formally:
###### Definition 8.1.1.
Set $`\mathrm{\Omega }_{\mathrm{com}}^1(A):=AA/(abcabc+acb).`$ The elements of $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ are called the *(commutative) Kähler differentials* of $`A`$ (intuitively, $`ab`$ should be thought of as a differential form $`adb`$).
Let us examine the connection of $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ to Hochschild homology. Recall that Hochschild homology is the homology of the complex $`A^3A^2A0`$. An element $`ab`$ of $`AA`$ is automatically a $`1`$-cycle, since $`d(ab)=abba=0`$ as $`A`$ is commutative. Since
$$d(abc)=abcabc+cab=abcabc+acb,$$
we see that the relation for $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ is derived precisely from the consideration of $`1`$-boundaries. So,
$$\mathrm{\Omega }_{\mathrm{com}}^1(A)HH_1(A).$$
If we think of $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ as the free $`A`$-module generated by symbols $`db`$ where $`d`$ acts as a derivation from $`A`$ to $`A`$, we find that $`d1_A=0`$ through the usual argument that a derivation is zero on constants. Indeed, we find that $`d(\lambda 1_A)=0`$ for all $`\lambda \mathrm{𝕜}`$.
The space of Kähler differentials for $`A`$ plays an important universal role relative to derivations. Define $`:A\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ by $`a=a1_A1_Aa`$. Then $``$ is indeed a derivation, and we will often denote it symbolically by $`ada`$.
###### Theorem 8.1.2.
Let $`M`$ be an $`A`$-module and let $`\theta :AM`$ be a derivation. Then the assignment $`\mathrm{\Omega }_{\mathrm{com}}^1(\theta ):(adb)a\theta (b)`$ gives an $`A`$-module map $`\mathrm{\Omega }_{\mathrm{com}}^1(\theta ):\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ $`M`$, which is uniquely defined by the requirement that the following diagram commutes:
$$\text{}.$$
The proof is a routine calculation.
###### Example 8.1.3.
Consider the case $`A=\mathrm{Sym}V`$, the symmetric algebra for the $`\mathrm{𝕜}`$-vector space $`A`$. Since a derivation on an algebra is uniquely defined by specifying its values on generators and applying the Leibniz rule, we find that $`\mathrm{Der}(\mathrm{Sym}V,M)\mathrm{Hom}_\mathrm{𝕜}(V,M)`$. Then the theorem asserts that there is an isomorphism between $`\mathrm{Der}(\mathrm{Sym}V,M)`$ and $`\mathrm{Hom}_A(\mathrm{\Omega }_{\mathrm{com}}^1(A),M)`$, so $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ is the unique $`\mathrm{Sym}V`$-module such that
$$\mathrm{Hom}_\mathrm{𝕜}(V,M)\mathrm{Hom}_A(\mathrm{\Omega }_{\mathrm{com}}^1(A),M).$$
Clearly, $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ is the free $`\mathrm{Sym}V`$-module with $`V`$ being the space of generators that is, $`\mathrm{\Omega }_{\mathrm{com}}^1(A)(\mathrm{Sym}V)V`$. With this definition, we can explicitly calculate that $`:\mathrm{Sym}V=A\mathrm{\Omega }_{\mathrm{com}}^1(A)=(\mathrm{Sym}V)V`$ is given by
$$(v_1\mathrm{}v_n)=\underset{i=1}{\overset{n}{}}v_1\mathrm{}\widehat{v}_i\mathrm{}v_nv_i.$$
$`\mathrm{}`$
Suppose now $`A=\mathrm{𝕜}[X]`$, the algebra of regular functions on some affine variety $`X`$. We would like to identify $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ with $`𝒯^{}(X)`$, the space of global sections of the cotangent bundle on $`X`$. This can be done in the following way.
Consider the diagonal embedding $`XX\times X`$. Then we can view $`𝒯_X`$ as the normal bundle to $`X`$ in $`X\times X`$, that is, $`TX=T_X(X\times X)`$. Similarly, we view $`𝒯^{}(X)`$ as the conormal bundle $`T_X^{}(X\times X)`$. We would like to be able to identify $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ with $`\mathrm{\Gamma }(X,𝒯_X^{}(X\times X))`$, the space of regular sections of the corresponding conormal sheaf.
In general, if we have an embedding $`XY`$ of affine varieties, then we obtain an embedding $`I(X)\mathrm{𝕜}[Y]`$, where $`I(X)`$ denotes the ideal of regular functions vanishing on $`X`$. We can view $`I(X)/I(X)^2`$ as a linear form on $`T_XY`$ which is zero in the $`X`$-direction. Intuitively, if $`fI(X)`$, then the Taylor expansion of $`f`$ around a point $`xX`$ begins with the first derivative. By quotienting out $`I(X)^2`$, we are killing the higher derivative terms, and only considering the derivative of $`f`$ applied to a tangent vector to $`Y`$. So, the definition $`\mathrm{\Gamma }(X,𝒯_X^{}Y):=I(X)/I(X)^2`$ seems to be a legitimate one.
In our case, $`Y=X\times X`$ and $`XX\times X`$ is the diagonal map, so $`I(X)=\mathrm{Ker}(AAA)`$, where the map is multiplication. Notice that multiplication is an algebra map if and only if $`A`$ is commutative. We can now formulate the following proposition.
###### Proposition 8.1.4.
There is a canonical isomorphism of $`A`$-modules,
$$\mathrm{\Omega }_{\mathrm{com}}^1(A)I(X)/I(X)^2.$$
###### Proof.
First, we will set $`I=I(X)=\mathrm{Ker}(AAA)`$ to simplify the notation.
Given $`A`$ a commutative algebra, the multiplication map $`m:AAA`$ is an algebra homomorphism, and we let $`I`$ denote its kernel, which is therefore an ideal of $`AA`$. Let $`M`$ be a left $`A`$-module. We will view $`M`$ as an $`A`$-bimodule by equipping it with the symmetric $`A`$-bimodule structure, that is, $`ma=am`$ for all $`mM`$, $`aA`$. Then the short exact sequence of $`A`$-bimodules
induces a long exact sequence of $`\mathrm{Ext}`$-groups, the first few terms of which we produce below
$`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,M)`$ $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(AA,M)\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(I,M)`$
$`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(A,M)\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(AA,M)\mathrm{}.`$
Now, since an $`A`$-bimodule is a left $`AA^{\mathrm{op}}`$-module, and $`A`$ is commutative, we see that $`AA`$ is the free $`A`$-bimodule of rank one. Hence $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(AA,M)=0`$ and $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(AA,M)M`$. Also, we have by definition that $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(A,M)HH^1(M)`$, the first Hochschild cohomology of $`A`$ with coefficients in $`M`$, which is precisely the set of all outer derivations from $`A`$ to $`M`$. However, $`M`$ has a symmetric $`A`$-bimodule structure, so there are no inner derivations. Hence, $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(A,M)\mathrm{Der}(A,M)`$.
Finally, we observe that $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,M)=\{mMam=ma\}`$, since this is just the zeroth-degree Hochschild cohomology. But $`M`$ is symmetric, so this is all of $`M`$. Therefore, the map $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,M)\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(AA,M)`$ is an isomorphism, and $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(I,M)\mathrm{Der}(A,M)`$ is a surjection. A calculation then shows that in fact $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(I,M)\mathrm{Hom}_A(I/I^2,M)`$, so we see that $`\mathrm{\Omega }_{\mathrm{com}}^1(A)I/I^2`$, which is essentially the definition we gave before.
Then given an $`A`$-bimodule (i.e., an $`AA`$-module), the $`A`$-bimodule $`M/IM`$ has a symmetric bimodule structure, that is, $`am=ma`$ for all $`aA`$ and $`mM`$. Second, we recall that the bar complex is exact, and that it was (initially) given by
$$\text{}.$$
So, $`I=\mathrm{Im}b`$. So, the following sequence is exact
$$\text{}.$$
Now, both $`\mathrm{Tor}`$ and $`\mathrm{Ext}`$ are homological functors, so they associate long exact sequences to short exact sequences such as the one above. Since $`AA`$ is a free $`A`$-bimodule, $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(AA,)=0`$. By standard results in homological algebra, we obtain that for any $`A`$-bimodule $`M`$,
$$\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(A,M)\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^0(I,M).$$
But we know that $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(A,M)HH^1(M)\mathrm{Der}(A,M)`$, and also that $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^0(I,M)HH^0(I,M)=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(I,M)`$. Now recall that $`A`$ is commutative and $`M`$ is a symmetric bimodule. Then
$$\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(I,M)=\mathrm{Hom}_{A\text{-}\mathrm{𝗆𝗈𝖽}}(I/I^2,M).$$
Following the line of isomorphisms, we find that indeed $`I/I^2\mathrm{\Omega }_{\mathrm{com}}^1(A)`$. ∎
This proof relied heavily on the symmetry of the module actions. Indeed, this is the first point where we will see noncommutative and commutative geometry diverging.
###### Remark 8.1.5.
Identify $`\mathrm{\Omega }_{\mathrm{com}}(A)`$ with $`I/I^2`$ where $`I=\mathrm{Ker}[m:AAA]`$. It is easy to check that the map $`\mathrm{\Omega }_{\mathrm{com}}(\theta ):I/I^2M`$ corresponding to a derivation $`\theta :AM`$ is induced by a map $`AAIM`$ given by a similar formula
$$\mathrm{\Omega }_{\mathrm{com}}(\theta ):_{i=1}^na_ib_i_{i=1}^na_i\theta (b_i).$$
The Leibniz formula for $`\theta `$ insures that this map indeed vanishes on $`I^2`$.
## 9. The Hochschild-Kostant-Rosenberg Theorem
### 9.1. Smoothness.
Let $`X\mathrm{𝕜}^n`$ be an algebraic set defined by a system of polynomial equations
$$X=\{x=(x_1,\mathrm{},x_n)\mathrm{𝕜}^n|f_1(x)=0,\mathrm{},f_r(x)=0\}.$$
(9.1.1)
Thus, $`X`$ is the zero variety of the ideal $`I:=f_1,\mathrm{},f_r\mathrm{𝕜}[x_1,\mathrm{},x_n]`$. We call $`\mathrm{𝕜}[X]:=\mathrm{𝕜}[x_1,\mathrm{},x_n]/I`$ the (scheme-theoretic) coordinate ring of $`X`$. Hilbert’s Nullstellensatz says that the ideal $`I`$ is radical, i.e., the algebra $`\mathrm{𝕜}[X]`$ has no nilpotents, if and only if $`I`$ is equal to the ideal of all polynomials $`f\mathrm{𝕜}[x_1,\mathrm{},x_n]`$ that vanish on the set $`X`$ pointwise. In this case, the coordinate ring $`\mathrm{𝕜}[X]`$ is said to be reduced, and the algebraic set $`X`$ is called an affine algebraic variety.
We would like to discuss the notion of smoothness of algebraic varieties. To this end, fix an algebraic set $`X`$ as in (9.1.1). For any point $`aX`$, we introduce the following Jacobian $`n\times r`$-matrix:
$$J_a(f_1,\mathrm{},f_r):=\left(\begin{array}{ccc}\frac{f_1(x)}{x_1}& \mathrm{}& \frac{f_1(x)}{x_n}\\ \frac{f_2(x)}{x_1}& \mathrm{}& \frac{f_2(x)}{x_n}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{f_r(x)}{x_1}& \mathrm{}& \frac{f_r(x)}{x_n}\end{array}\right)|_{x=a}$$
Further, let $`𝒪_a`$ denote the local ring at $`a`$ (the localization of $`\mathrm{𝕜}[X]`$ with respect to the multiplicative set of all polynomials $`f`$ such that $`f(a)0`$), and write $`𝔪_a𝒪_a`$ for the maximal ideal of the local ring $`𝒪_a`$. Thus $`𝒪_a/𝔪_a=\mathrm{𝕜}`$, and $`𝔪_a/𝔪_a^{\mathrm{\hspace{0.17em}2}}`$ is a finite dimensional vector space over $`𝒪_a/𝔪_a=\mathrm{𝕜}`$. We consider the graded algebra $`\mathrm{Sym}^{\text{}}(𝔪_x/𝔪_x^{\mathrm{\hspace{0.17em}2}}):=_{i0}\mathrm{Sym}^i(𝔪_x/𝔪_x^{\mathrm{\hspace{0.17em}2}})`$, and also the graded algebra $`_{i0}𝔪^i/𝔪_x^{i+1}`$. The imbedding $`𝔪_x/𝔪_x^{\mathrm{\hspace{0.17em}2}}{}_{}{}^{}_{}^{}_{i0}𝔪^i/𝔪_x^{i+1}`$ extends by multiplicativity to a graded algebra homomorphism $`\mathrm{Sym}^{\text{}}(𝔪_x/𝔪_x^{\mathrm{\hspace{0.17em}2}})_{i0}𝔪^i/𝔪_x^{i+1}`$.
One has the following basic result.
###### Theorem 9.1.2.
For an irreducible algebraic set $`X`$ the following properties $`(𝗂)`$$`(\mathrm{𝗂𝗏})`$are equivalent:
$`(𝗂)`$For any $`xX,`$ the map $`\mathrm{Sym}^{\text{}}(𝔪_x/𝔪_x^{\mathrm{\hspace{0.17em}2}})_{i0}𝔪^i/𝔪_x^{i+1}`$ is an isomorphism;
$`(\mathrm{𝗂𝗂})`$The module $`\mathrm{\Omega }_{\mathrm{com}}^1(\mathrm{𝕜}[X])`$ of Kähler differentials is a projective $`\mathrm{𝕜}[X]`$-module;
$`(\mathrm{𝗂𝗂𝗂})`$The coordinate ring $`\mathrm{𝕜}[X]`$, viewed as a module over $`\mathrm{𝕜}[X]^\mathrm{e}=\mathrm{𝕜}[X]\mathrm{𝕜}[X]`$, has a finite projective resolution;
$`(\mathrm{𝗂𝗏})`$For any point $`aX`$ one has $`\mathrm{rk}J_a(f_1,\mathrm{},f_r)=ndimX`$.
The algebraic set $`X`$ satisfying the equivalent conditions (i)–(iv) of the Theorem is called smooth. It is easy to see that condition (i) above implies that $`\mathrm{𝕜}[X]`$ is reduced, i.e., a smooth algebraic set is necessarily an affine algebraic variety.
###### Theorem 9.1.3 (HKR).
Let $`A=\mathrm{𝕜}[X]`$, where $`X`$ is a smooth affine variety. Then
$`HH_k(A)`$ $`=\mathrm{\Gamma }(X,\mathrm{\Lambda }^k𝒯^{}(X))=\mathrm{\Lambda }_A^k\mathrm{\Omega }_{\mathrm{com}}^1(A)`$
$`HH^k(A)`$ $`=\mathrm{\Gamma }(X,\mathrm{\Lambda }^k𝒯(X))=\mathrm{\Lambda }_A^k\mathrm{Der}(A),`$
where $`𝒯(X)`$ is the tangent bundle of $`X`$, $`𝒯^{}(X)`$ is the cotangent bundle, and $`\mathrm{\Gamma }(X,)`$ denotes global sections.
###### Remark 9.1.4.
First, observe that the HKR theorem shows that both Hochschild homology and cohomology are commutative algebras (i.e., the total homology $`HH_{\text{}}(A)=_{n=0}^{\mathrm{}}HH_n(A)`$ is an algebra, as is the cohomology). In the case of cohomology, this is not too terribly surprising. Recall that we defined
$$HH^{\text{}}(A)=\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^{\text{}}(A,A),$$
which has a commutative algebra structure induced by the diagonal map $`AAA`$. We can see in a more elementary way that $`HH_{\text{}}(A)=\mathrm{Tor}_{\text{}}^{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,A)`$ has a commutative algebra structure by observing that the multiplication map $`AAA`$ is an algebra map if and only if $`A`$ is commutative. In particular, we would not expect $`HH_{\text{}}(A)`$ to be an algebra if $`A`$ is not commutative. $`\mathrm{}`$
### 9.2. From Hochschild complex to Chevalley-Eilenberg complex.
For any commutative algebra $`A`$ and a left $`A`$-module $`M`$, viewed as a symmetric $`A`$-bimodule, we are going to construct the following natural maps:
(9.2.3)
(9.2.6)
The proof involves Lie algebra homology, which we recall here. If $`𝔤`$ is a Lie algebra and $`M`$ is a $`𝔤`$-module, then the Lie algebra homology $`H_p^{\mathrm{Lie}}(𝔤,M)`$ is computed in terms of the following complex: the $`p^{\text{th}}`$ term $`C_p^{\mathrm{Lie}}(𝔤,M)`$ is given by $`M\mathrm{\Lambda }^p𝔤`$, and the differential $`d:C_p^{\mathrm{Lie}}(𝔤,M)C_{p1}^{\mathrm{Lie}}(,M)`$ is given by
$`d_{\mathrm{Lie}}(m`$ $`(x_1\mathrm{}x_p))={\displaystyle }_{i=1}^p(1)^ix_im(x_1\mathrm{}\widehat{x}_i\mathrm{}x_p)`$
$`+{\displaystyle \underset{j<k}{}}(1)^{j+k}m([x_j,x_k]x_1\mathrm{}\widehat{x}_j\mathrm{}\widehat{x}_k\mathrm{}x_p).`$
Now suppose that $`A`$ is an associative algebra and $`M`$ is an $`A`$-bimodule. Consider $`A`$ as a Lie algebra under the commutator bracket, and make $`M`$ a Lie $`A`$-module via the action $`(a,m)amma`$. Notice that if $`A`$ is commutative and $`M`$ is symmetric, then both the Lie algebra structure on $`A`$ and the Lie $`A`$-module structure of $`M`$ are trivial.
Consider the following diagram
$$\text{},$$
where $`\mathrm{alt}:\mathrm{\Lambda }_\mathrm{𝕜}^pAA^p`$ is the completely alternating map given by
$$\mathrm{alt}(a_1\mathrm{}a_p)=\underset{\sigma S_p}{}(1)^{\mathrm{𝚜𝚒𝚐𝚗}\sigma }a_{\sigma (1)}\mathrm{}a_{\sigma (p)}.$$
In the case that $`A`$ is commutative and $`M`$ symmetric, then the image of the alternating map on $`C_p^{\mathrm{Lie}}(A,M)`$ is zero. In any case, it is tedious but easy to check that $`\mathrm{alt}(C_p^{\mathrm{Lie}}(A,M))Z_p(A,M)`$. By passing to the quotient by $`B_p(A,M)`$, we obtain a map $`\mathrm{alt}:M\mathrm{\Lambda }_\mathrm{𝕜}^pAHH_p(A,M)`$.
We will now construct the inverse map and leave it to the reader to check the remainder of the proof. Choose any $`m(a_1\mathrm{}a_p)C_p(A,M)`$. Define $`\pi :C_p(A,M)M_A\mathrm{\Lambda }_A^p\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ by
$$\pi (m(a_1\mathrm{}a_p))=mda_1\mathrm{}da_p.$$
It is then easy to see that $`\pi \mathrm{alt}`$ is given by
$$m(a_1\mathrm{}a_p)p!m(da_1\mathrm{}da_p).$$
We claim that in the case that $`M=A`$, where $`A`$ is the coordinate ring of a smooth affine variety, $`\pi `$ is an isomorphism from $`HH_p(A)`$ to $`\mathrm{\Lambda }_A^p\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ as desired. For simplicity’s sake, we will let $`\mathrm{\Omega }_{\mathrm{com}}^p(A)=\mathrm{\Lambda }_A^p\mathrm{\Omega }_{\mathrm{com}}^1(A)`$.
Indeed, we will show that $`\pi `$ is an isomorphism of $`A`$-modules. It is clear that $`\mathrm{\Omega }_{\mathrm{com}}^p(A)`$ is an $`A`$-module by its definition. Observe that for any associative algebra, $`HH_p(M)`$ has the structure of a $`𝖹_A`$-module for all $`p`$ (recall that $`𝖹_A`$ is the center of $`A`$). Indeed, if $`z𝖹_A`$, the action is given by
$$z[m(a_1\mathrm{}a_p)]=zm(a_1\mathrm{}a_p).$$
It is easy to see that the $`z`$-action commutes with $`d`$ since $`z`$ commutes with all elements of $`A`$. Since $`A`$ is commutative, $`𝖹_A=A`$, hence $`HH_p(M)`$ is an $`A`$-module. The proof that $`\pi `$ is an isomorphism then follows a standard argument. If $`\phi :MN`$ is a map of $`R`$-modules for some commutative ring $`R`$, then $`\phi `$ is an isomorphism if and only if $`\phi _𝔪:M_𝔪N_𝔪`$ is an isomorphism for every maximal ideal $`𝔪R`$, where $`M_𝔪`$ is $`M`$ localized at $`𝔪`$.
### 9.3. Proof of Theorem 9.1.3.
We will only prove the result for the homology, since the result for cohomology is analogous.
Now, recall that $`HH_p(A)=\mathrm{Tor}_p^{AA}(A,A)`$. We now localize at some maximal ideal $`𝔪A`$. We then obtain the localized map
$$\text{}.$$
Ideally, we would hope that $`\mathrm{Tor}_p^{AA}(A,A)_𝔪=\mathrm{Tor}_p^{A_𝔪A_𝔪}(A_𝔪,A_𝔪)`$. Indeed, this is the case. If $`M`$ is any $`A`$-bimodule, then $`\mathrm{Tor}_0^{AA}(A,M)=A_{AA}M`$. Then we see that
$$\mathrm{Tor}_p^{AA}(A,A)_𝔪=(A_{AA}M)_𝔪=A_𝔪_{A_𝔪A_𝔪}M_𝔪=\mathrm{Tor}_p^{A_𝔪𝔪}(A_𝔪,M_𝔪).$$
Recall that localization is an exact functor. So, it commutes with derived functors, in particular, it commutes with the derived functors of $``$, which are precisely the $`\mathrm{Tor}`$-groups. So, we are reduced to the case of a map
$$\text{}.$$
So, it suffices to consider a local ring.
###### Lemma 9.3.1.
Let $`R`$ be a commutative $`\mathrm{𝕜}`$-algebra, $`JR`$ an ideal finitely generated by $`x_1,\mathrm{},x_nR`$, where $`x_1,\mathrm{},x_n`$ is a regular sequence (that is, $`x_{i+1}`$ is not a zero divisor in $`R/(x_1,\mathrm{},x_i)`$ for all $`i`$). Then
$$\mathrm{Tor}_1^R(R/J,R/J)J/J^2.$$
Indeed we have an isomorphism for all $`p`$,
$$\mathrm{Tor}_p^R(R/J,R/J)\mathrm{\Lambda }_{R/J}^p(J/J^2).$$
We will omit the proof of this lemma. With this result in hand, we set $`R=AA`$, which we view as $`\mathrm{𝕜}[X\times X]`$, and we let $`J`$ be the ideal of regular functions on $`X\times X`$ vanishing on the diagonal. If $`X`$ is smooth, then $`J`$ is indeed generated by elements forming a regular sequence. An application of the lemma finishes the proof.
### 9.4. Digression: Applications to formality.
In this section we are going to use the relation between Hochschild and Chevalley-Eilenberg complexes to obtain some (non-obvious) formality results in the algebraic geometry.
We fix $`X`$, a smooth projective variety with the structure sheaf $`𝒪_X`$. We write $`D_{\mathrm{coh}}^b(X)`$ for the bounded derived category of complexes of $`𝒪_X`$-modules with coherent cohomology sheaves, cf. e.g. \[???\].
Let $`ı:X{}_{}{}^{}_{}^{}X\times X`$ be the diagonal imbedding. Associated to this imbedding, one has a direct image functor $`ı_{}:D_{\mathrm{coh}}^b(X)D_{\mathrm{coh}}^b(X\times X)`$, and an inverse image functor $`ı^{}:D_{\mathrm{coh}}^b(X\times X)D_{\mathrm{coh}}^b(X)`$. We will be interested in the composite functor
$$ı^{}ı_{}:D_{\mathrm{coh}}^b(X)D_{\mathrm{coh}}^b(X\times X)D_{\mathrm{coh}}^b(X).$$
It is known, for instance, that for any coherent sheaf $``$ on $`X`$, the cohomology sheaves of the object $`ı^{}ı_{}`$ are given by
$$^i(ı^{}ı_{})=\mathrm{\Omega }_X^i_{_{𝒪_X}}.$$
The Proposition below implies that the object $`ı^{}ı_{}D_{\mathrm{coh}}^b(X)`$ is actually quasi-isomorphic to a direct sum of its cohomology sheaves.
###### Proposition 9.4.1.
There is a quasi-isomorphism
$$ı^{}ı_{}𝒪_X_{i=0}^{dimX}\mathrm{\Omega }_X^i[i]\text{in}D_{\mathrm{coh}}^b(X).$$
###### Proof.
Assume first that $`X=\mathrm{Spec}A`$, is affine. We have $`AA=\mathrm{𝕜}[X\times X]`$. The kernel $`I:=\mathrm{Ker}(AAA)`$ of the multiplication map may be identified with the defining ideal of the diagonal $`XX\times X`$.
We know that the bar complex $`𝖡_{\text{}}A`$ provides a resolution of the $`\mathrm{𝕜}[X\times X]`$-module $`A=\mathrm{𝕜}[X\times X]/I`$ by free $`\mathrm{𝕜}[X\times X]`$-modules. For each $`i=0,1,\mathrm{},`$ let $`\widehat{𝖡}_iA`$ be the $`I`$-adic completion of $`𝖡_iA`$. The standard homotopy on the bar complex, that shows that the complex $`𝖡_{\text{}}A`$ is acyclic, extends to the completions. Hence, the completed bar complex provides a resolution of $`\widehat{A^\mathrm{e}/I}=\mathrm{𝕜}[X\times X]/I`$ of the form
$$\mathrm{}\widehat{𝖡}_iA\widehat{𝖡}_{i1}A\mathrm{}\widehat{𝖡}_1A\widehat{𝖡}_0A{}_{}{}^{}_{}^{}\mathrm{𝕜}[X\times X]/I=\mathrm{𝕜}[X],$$
where each term $`\widehat{𝖡}_iA`$ is a flat $`\mathrm{𝕜}[X\times X]`$-module.
The resolution above globalizes naturally to an arbitrary, not necessarily affine, variety $`X`$. Specifically, let $`X`$ be any smooth variety. For each $`i0`$, we let $`\widehat{𝖡}_i(X):=\widehat{𝒪}_{X^{i+2}}`$ be the completion of the structure sheaf of the Cartesian power $`X^{i+2}`$ along the principal diagonal $`X{}_{}{}^{}_{}^{}X^{i+2}`$. The projection $`pr_{1,i+2}:X^{i+2}X\times X`$, on the first and last factors, makes $`\widehat{𝒪}_{X^{i+2}}`$ a sheaf of $`𝒪_{X\times X}`$-modules. Thus, we have constructed a complex $`\widehat{𝖡}_{\text{}}(X)`$ such that, for each $`i0`$, we have
$``$ $`\widehat{𝖡}_i(X)`$ is a (not quasi-coherent) sheaf of flat $`𝒪_{X\times X}`$-modules;
$``$ $`\widehat{𝖡}_i(X)`$ is set-theoretically supported on the diagonal $`X_\mathrm{\Delta }X\times X`$, so may be regarded as a sheaf on $`X_\mathrm{\Delta }`$;
$``$ For any Zariski-open affine subset $`U=\mathrm{Spec}AX_\mathrm{\Delta }`$, we have $`\mathrm{\Gamma }(U,\widehat{𝖡}_{\text{}}(X))=\widehat{𝖡}_{\text{}}A`$.
We deduce that the object $`ı^{}ı_{}𝒪_XD_{\mathrm{coh}}^b(X)`$ is represented by the following complex of $`𝒪_{X\times X}`$-modules
$$\widehat{𝖡}_{\text{}}(X)_{𝒪_{X\times X}}𝒪_{X_\mathrm{\Delta }}=\left[\mathrm{}\widehat{𝒪}_{X^3}\widehat{𝒪}_{X^2}𝒪_{X_\mathrm{\Delta }}0\right].$$
Now, the assignment sending, for each Zariski-open affine subset $`UX_\mathrm{\Delta }`$, $`a_0a_1\mathrm{}a_i\mathrm{\Gamma }(U,\widehat{𝒪}_{X^{i+1}})`$ to $`a_0da_1\mathrm{}da_i`$ yields a quasi-isomorphism of the above complex with $`_{i0}\mathrm{\Omega }_X^i[i].`$
###### Remark 9.4.2 (Kapranov).
Let $`X`$ be a smooth algebraic variety. Consider the complex $`RHom`$ (over $`X\times X`$) from $`𝒪_\mathrm{\Delta }`$ to itself, where $`\mathrm{\Delta }`$ is the diagonal. Its cohomology sheaves, i.e., the Ext’s are just the exterior powers of $`𝒯_X,`$ the tangent bundle. One may ask if this complex is quasiisomorphic to the direct sum of its cohomology sheaves.
Further, given a coherent sheaf $`F`$ on $`X,`$ let $`Quot_h(F)`$, be the quot-scheme that parametrizes quotient sheaves of $`F`$ with Hilbert polynomial $`h`$ or, equivalently, subsheaves $`K`$ in $`F`$ with Hilbert polynomial $`h_Fh.`$ The tangent space to $`Quot_h(F)`$ at the point represented by a sheaf $`K`$ is the space $`Hom_X(K,F/K).`$ In \[CK1\], the authors have constructed a derived quot-scheme, a smooth dg-manifold $`RQuot`$ whose tangent dg-space (i.e. complex) at a point $`K`$ as above is $`RHom(K,F/K).`$
In the special case where $`h=1`$ and $`F=𝒪_X`$, we have $`Quot_h=X`$, the quotient sheaves being skyscrapers at points of $`X.`$ Yet, the corresponding derived quot-scheme is different, its tangent space at $`xX`$ is a complex concentrated in nonnegative degrees whose degree $`i0`$ cohomology is equal to $`^{i+1}T_x(X)`$. According to Proposition 9.4.1, this dg-manifold is split, i.e., is quasiisomorphic to $`X`$ with the structure sheaf being the symmetric algebra of the graded algebra sheaf formed by the direct sum of $`\mathrm{\Omega }_X^i`$ in degree $`(i+1),i>0.`$ $`\mathrm{}`$
###### Remark 9.4.3.
For any dg-manifold $`Y`$ we have the scheme $`\pi _0(Y),`$ the spectrum of the $`0`$th cohomology of $`𝒪_Y,`$ and on $`\pi _0(Y)`$ we have the sheaf of $`\mathrm{Lie}_{\mathrm{}}`$ algebras obtained by restricting the tangent dg-bundle $`TY`$ onto $`\pi _0(Y)`$ in the sense of dg-schemes (note that $`\pi _0(Y)`$ is a dg-subscheme in $`Y`$). Denote this sheaf by $`t_Y`$ (small $`t`$ to avoid confusion). This sheaf of $`\mathrm{Lie}_{\mathrm{}}`$ algebras defines $`Y`$ up to quasiisomorphisms (duality between commutative and Lie algebras, sheafified along $`\pi _0(Y)`$). To be more precise, it determined the formal completion of $`Y`$ along $`\pi _0(Y)`$ (which is all that is needed in practice). Taking $`Y=RQuot`$ as before, we get $`\pi _0(Y)=X`$ and $`t_Y=`$ the quotient of $`R\mathrm{Hom}(𝒪_\mathrm{\Delta },𝒪_\mathrm{\Delta })`$ by the actual Hom in degree $`0,`$ then shifted by $`1.`$ $`\mathrm{}`$
## 10. Noncommutative differential forms
### 10.1.
For an associative commutative $`\mathrm{𝕜}`$-algebra $`A`$ and a left $`A`$-module $`M`$, we may consider the space $`\mathrm{Der}(A,M)`$ of all derivations $`\theta :AM`$. We have seen that the functor $`M\mathrm{Der}(A,M)`$ is representable by the $`A`$-module $`\mathrm{\Omega }_{\mathrm{com}}(A)`$ of Kähler differentials.
If $`A`$ is an associative not necessarily commutative $`\mathrm{𝕜}`$-algebra, the space $`\mathrm{Der}(A,M)`$ is defined provided $`M`$ is an $`A`$-bimodule. We are going to show that the functor $`M\mathrm{Der}(A,M)`$, defined on the category of $`A`$-bimodules, is also representable.
To this end, we let $`m:AAA`$ denote the multiplication map viewed as a map of $`A`$-bimodules. The kernel of $`m`$ is a sub-bimodule in $`AA`$.
###### Definition 10.1.1.
We denote $`\mathrm{\Omega }_{\mathrm{nc}}^1(A):=\mathrm{Ker}[m:AAA],`$ and call it the $`A`$-bimodule of *noncommutative $`1`$-forms* on $`A`$. Thus, one has the fundamental exact sequence of $`A`$-bimodules
$$0\mathrm{\Omega }_{\mathrm{nc}}^1(A)AA\stackrel{\mathrm{𝚖𝚞𝚕𝚝}}{}A\mathrm{\hspace{0.17em}0}.$$
###### Proposition 10.1.2.
For every $`MA\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}`$, there is a canonical isomorphism
$$\mathrm{Der}(A,M)\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^1(A),M).$$
Thus, the functor $`M\mathrm{Der}(A,M)`$ is representable by the $`A`$-bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$.
###### Proof.
Recall the bar resolution of $`A`$:
$$𝖡_{\text{}}A:\text{},$$
where $`𝖡_j=𝖡_jA=AA^jA`$, and the final map $`𝖡_0=AAA`$ is the multiplication map. Recall also that this complex is acyclic. Now let $`\delta :AM`$ be a derivation. We will define an $`A`$-bimodule map $`\stackrel{~}{\delta }:A^3M`$. We set
$$\stackrel{~}{\delta }(a^{}aa^{\prime \prime })=a^{}\delta (a)a^{\prime \prime }.$$
We claim that $`\delta `$ is a derivation if and only if $`\stackrel{~}{\delta }`$ is a Hochschild 2-cocycle. Indeed,
$`(b\stackrel{~}{\delta })(a_0`$ $`a_1a_3a_4)`$
$`=\stackrel{~}{\delta }(a_0a_1a_2a_3)\stackrel{~}{\delta }(a_0(a_1a_2)a_3)+\stackrel{~}{\delta }(a_0a_1(a_2a_3))`$
$`=a_0a_1\delta (a_2)a_3a_0\delta (a_1a_2)a_3+a_0\delta (a_1)a_2a_3.`$
Clearly, the latter expression vanishes for all $`a_1,a_2,a_3`$ if and only $`\delta `$ is a derivation.
Since $`\stackrel{~}{\delta }`$ is a $`2`$-cocycle, it is zero on all $`3`$-boundaries, that is, $`\stackrel{~}{\delta }|_{d𝖡_2}=0`$. So, $`\stackrel{~}{\delta }`$ descends to a well defined $`A`$-bimodule homomorphism $`\stackrel{~}{\delta }:𝖡_1/d𝖡_2M`$. Since the bar complex is exact, $`𝖡_1/d𝖡_2d𝖡_1`$, and exactness again implies that $`d𝖡_1`$ is precisely the set of all $`2`$-cocycles, that is, $`d𝖡_1`$ is the kernel of the differential $`𝖡_0A`$, which is precisely the multiplication map. So,
$$𝖡_1/d𝖡_2d𝖡_1=\mathrm{Ker}[𝖡_0A]=\mathrm{Ker}[m:AAA]=\mathrm{\Omega }_{\mathrm{nc}}^1(A).$$
Hence $`\stackrel{~}{\delta }\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^1(A),M).`$ Clearly we can reverse this argument and produce from every $`A`$-bimodule map $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)M`$ a derivation $`AM`$ in an analogous manner. ∎
Here is another picture of the $`1`$-forms on $`A`$. Set $`I=\mathrm{Ker}[m:AAA]`$. Let $`\overline{A}`$ denote the (vector space) quotient $`A/\mathrm{𝕜}=A/\mathrm{𝕜}1_A`$, and write $`x\overline{x}`$ for the projection $`A\overline{A}=A/\mathrm{𝕜}.`$
###### Proposition 10.1.3.
$`(𝗂)`$The map $`d:AI,xdy=y1_A1_Ay`$ is a derivation.
$`(\mathrm{𝗂𝗂})`$The map $`A\overline{A}I,x\overline{y}xyxy1_A=xdy`$ is well-defined, and it is an isomorphisms of left $`A`$-bimodules.
###### Proof.
Endow $`A\overline{A}`$ with an $`A`$-bimodule structure by equipping it with the obvious left action and by setting
$$(x\overline{y})z=x\overline{yz}xy\overline{z}.$$
We let $`\psi :x\overline{y}xyxy1_A`$ denote the map considered in the Proposition. Observe that $`\psi `$ is well-defined. Indeed, if $`y`$ and $`y+\lambda 1_A`$ are two representatives of $`\overline{y}\overline{A}`$, then
$`x(y+\lambda 1_A)x(\lambda 1_A+y)1`$ $`=xy+x\lambda 1_A\lambda 1_Ax1_Axy1_A`$
$`=xyxy1_A+x\lambda 1_Ax\lambda 1_A`$
$`=xyxy1_A.`$
To show that $`\psi `$ is surjective, observe that $`I`$ is generated by terms of the form $`xyxy1_A`$. Indeed, this follows since $`I=d𝖡_1`$ (where again $`𝖡_1=A^3`$ from the bar complex). But then $`\psi (x\overline{y})`$ is nothing more than $`d(xy1_A)`$. Since $`d`$ is an $`A`$-bimodule map, $`d(xyz)=d(xy1_A)z=\psi (x\overline{y})z`$, so $`I`$ is generated by the image of $`\psi `$ as an $`A`$-bimodule.
To see that $`\psi `$ is an injection, consider the multiplication map $`m:AAA`$. Since $`\mathrm{Im}\psi I=\mathrm{Ker}m`$, we can descend to an $`A`$-bimodule map $`\overline{m}:AA/\mathrm{Im}\psi A`$. Define $`\phi :AAA/\mathrm{Im}\psi `$ by $`\phi (a)=a1_A+\mathrm{Im}\psi `$. Then
$`\overline{m}(\phi (a))`$ $`=\overline{m}(a1_A+\mathrm{Im}\psi )=m(a1_A)=a`$
$`\phi (\overline{m}(xy+\mathrm{Im}\psi ))=\phi (xy)=xy1_A+\mathrm{Im}\psi .`$
Since $`xyxy1_A\mathrm{Im}\psi `$, we see that $`xy1_A+\mathrm{Im}\psi =xy+\mathrm{Im}\psi `$. So, $`\phi `$ is an inverse to $`\overline{m}`$, hence $`\overline{m}`$ is an isomorphism. But then $`\mathrm{Im}\psi =\mathrm{Ker}m=I`$, as claimed. ∎
Further, combining Propositions 10.1.2 and 10.1.3, we obtain the following result, which is completely analogous to a similar result for Kähler differentials proved in §8.
###### Corollary 10.1.4.
Let $`A`$ be an associative not necessarily commutative algebra with unit and $`M`$ an $`A`$-bimodule. For any derivation $`\theta :AM`$, the assignment $`xdyx\theta (y)`$ gives a well-defined $`A`$-bimodule map $`\mathrm{\Omega }_{\mathrm{nc}}^1(\theta )`$ that makes the following diagram commute
Next we observe that
$$H_0(A,A^\mathrm{e})=A^\mathrm{e}/[A,A^\mathrm{e}]=A^\mathrm{e}_{A^\mathrm{e}}A=A,$$
(10.1.5)
where the isomorphism $`A^\mathrm{e}/[A,A^\mathrm{e}]\stackrel{_{}}{}A`$ is induced, explicitly, by the assignment $`A^\mathrm{e}uvvu`$. One verifies unrevelling the definitions, that the canonical map $`H_0(A,A^\mathrm{e})H_0(A,A)`$ induced by the multiplication morphism $`A^\mathrm{e}A`$ is nothing but the natural projection $`A{}_{}{}^{}_{}^{}A/[A,A]`$.
We have the following maps
(10.1.6)
One verifies that the composite map is given by the formula $`x\overline{y}yxxy.`$
### 10.2. The differential envelope.
Our goal is to construct an analogue of de Rham differential on noncommutative differential forms. To this end, we recall
###### Definition 10.2.1.
A *differential graded algebra* (DGA) is a graded algebra $`D=_nD^n`$ equipped with a super-derivation $`d:D^nD^{n+1}`$ such that $`d^2=0`$.
Given an associative algebra $`A`$, its *differential envelope* $`D(A)`$ is the solution to the following lifting problem. Consider the category whose objects consist of all algebra maps from $`A`$ into the zeroth degree of some DGA $`D^0`$. If $`\phi :AD^0`$ and $`\psi :AE^0`$ are two such objects, a morphism from $`\phi `$ to $`\psi `$ is a map $`\theta :DE`$ such that $`\theta (D^n)E^n`$ and $`d_E\theta =d_D`$. The differential envelope $`D(A)`$ is then defined to be the initial object in this category, that is, there is a map of algebras $`i:AD(A)^0`$ such that for any algebra map $`\phi :AD^0`$ ($`D`$ a DGA), there is a unique map of DGA’s $`\psi :D(A)D`$ such that $`\psi i=\phi `$.
Explicitly, $`D(A)`$ is generated by the algebra $`A`$ and all symbols of the form $`\overline{a}=da`$, where $`aA`$ and $`d(ab)=(da)b+a(db)`$.
###### Remark 10.2.2.
It is clear that if $`A`$ is generated, as an algebra, by $`a_1,\mathrm{},a_n,`$ then the elements $`a_1,\mathrm{},a_n,da_1,\mathrm{},da_n,`$ generate $`D(A)`$ as an algebra. Thus, if $`A`$ is a finitely generated algebra then so is the algebra $`D(A)=_{i0}D^i(A)`$. Moreover, one verifies similarly that in that case each homogeneous component $`D^i(A)`$ is finitely generated as an $`A`$-bimodule. $`\mathrm{}`$
###### Definition 10.2.3.
Define the algebra $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ of *noncommutative differential forms* on $`A`$ to be the tensor algebra (over $`A`$) of the $`A`$-bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1(A),`$ that is
$`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ $`:=T_A\mathrm{\Omega }_{\mathrm{nc}}^1(A)=A\mathrm{\Omega }_{\mathrm{nc}}^1(A)T_A^2\mathrm{\Omega }_{\mathrm{nc}}^1(A)T_A^3\mathrm{\Omega }_{\mathrm{nc}}^1(A)\mathrm{}.`$
It turns out that there is a canonical graded algebra isomorphism
$$\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)D(A).$$
(10.2.4)
To construct this isomorphism, we first consider the canonical algebra map $`i:AD(A)`$ (this turns out to be an injection). We also have a derivation $`d:AD(A)`$ given by $`ada=\overline{a}`$. By our definition of $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$, there is then a unique $`A`$-bimodule map $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)D(A)`$. This $`A`$-bimodule morphism extends canonically to an algebra map (10.2.4).
###### Proposition 10.2.5.
With the notation as above, the map $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)D(A)`$ in (10.2.4) is a graded algebra isomorphism.
Thus, using the isomorphism of the Proposition, we transport the differential on $`D(A)`$ to obtain a degree one super-differential $`d:\mathrm{\Omega }_{\mathrm{nc}}^n(A)\mathrm{\Omega }_{\mathrm{nc}}^{n+1}(A)`$ making $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ a DGA.
###### Corollary 10.2.6.
The graded algebra $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ can be given the structure of DGA with differential transported from the one on $`D(A)`$.
Proposition 10.2.5 will be proved later as a special case of Theorem 10.7.1 below. One also has the following important result proved in \[CQ1\].
###### Theorem 10.2.7.
There is a canonical isomorphism of left $`A`$-modules
$$\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\underset{n=0}{\overset{\mathrm{}}{}}A\overline{A}^n$$
With this isomorphism, the differential $`d:\mathrm{\Omega }_{\mathrm{nc}}^n(A)\mathrm{\Omega }_{\mathrm{nc}}^{n+1}(A)`$ is given by the formula
$$d(a_0\overline{a_1}\mathrm{}\overline{a_n})=1\overline{a_0}\overline{a_1}\mathrm{}\overline{a_n}.$$
We denote the summand $`A\overline{A}^n`$ by $`\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$. We denote the element $`a_0\overline{a_1}\mathrm{}\overline{a_n}\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$ by $`a_0da_1\mathrm{}da_n`$. so that the differential reads
$$d(a_0da_1\mathrm{}da_n)=da_0da_1\mathrm{}da_n.$$
### 10.3. The universal square-zero extension.
The algebra imbedding $`A{}_{}{}^{}_{}^{}\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ makes, for each $`k`$, the graded component $`\mathrm{\Omega }_{\mathrm{nc}}^k(A)`$ an $`A`$-bimodule. We define an algebra structure on the vector space $`A\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$ by the formula
$$(a,\omega )(a^{},\omega ^{})=(aa^{},a\omega ^{}+\omega a^{}+c(a,a^{})),\text{where}c(a,a^{})=dada^{}.$$
It is straightforward to check that this way we get an associative algebra, to be denoted $`A\mathrm{}_c\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$. Clearly, $`\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$ is a two-sided ideal in $`A\mathrm{}_c\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$. Moreover, the natural imbedding $`i:\omega (0,\omega )`$, and the projection $`(a,m)a,`$ give rise to a square-zero extension
$$0\mathrm{\Omega }_{\mathrm{nc}}^2(A)\stackrel{i}{}A\mathrm{}_c\mathrm{\Omega }_{\mathrm{nc}}^2(A)\stackrel{p}{}A0.$$
(10.3.1)
###### Lemma 10.3.2.
The square-zero extension (10.3.1) is universal, i.e., for any square-zero extension $`I{}_{}{}^{}_{}^{}\stackrel{~}{A}{}_{}{}^{}_{}^{}A`$, there exists a unique $`A`$-bimodule map $`\phi `$ making the following diagram commute
###### Proof.
Choose a $`\mathrm{𝕜}`$-linear splitting $`c:A{}_{}{}^{}_{}^{}\stackrel{~}{A}`$, and note that, for any $`a_1,a_2A`$, we have $`c(a_1)c(a_2)c(a_1a_2)I`$. Define the map
$$\stackrel{~}{\phi }:A\overline{A}^2I,aa_1a_2a[c(a_1)c(a_2)c(a_1a_2)].$$
We use the isomorphism $`:\mathrm{\Omega }_{\mathrm{nc}}^2(A)A\overline{A}^2`$ to view the map above as a map $`\stackrel{~}{\phi }::\mathrm{\Omega }_{\mathrm{nc}}^2(A)I,ada_1da_2\stackrel{~}{\phi }(ada_1da_2).`$ We leave to the reader to verify that this is an $`A`$-bimodule map.
Further, write $`\stackrel{~}{A}AI`$ for the vector space decomposition corresponding to the splitting $`c`$. One shows that the following assignment
$$A\mathrm{}_c\mathrm{\Omega }_{\mathrm{nc}}^2(A)AI=\stackrel{~}{A},a^{}ada_1da_2a^{}\stackrel{~}{\phi }(ada_1da_2)$$
gives a map with all the required properties. ∎
### 10.4. Hochschild differential on non-commutative forms.
Computing Hochschild homology of the $`A`$-bimodule $`M=A`$ using the bar resolution, we get $`\overline{C}_{\text{}}(A,A)=\overline{𝖡}_{\text{}}(A)_{A^\mathrm{e}}A`$, where $`\overline{𝖡}_{\text{}}(A)`$ is the reduced bar complex and $`A^\mathrm{e}=AA^{\mathrm{op}}`$. Thus
$$\overline{C}_n(A,A)=\mathrm{\Omega }_{\mathrm{nc}}^n(A)A_{A^\mathrm{e}}A\mathrm{\Omega }_{\mathrm{nc}}^n(A).$$
The Hochschild differential $`\overline{C}_{n+1}(A,A)\overline{C}_n(A,A)`$ is then given by the differential $`b:\mathrm{\Omega }_{\mathrm{nc}}^{n+1}(A)\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$ mentioned before. So, $`HH_{\text{}}(A)=H_{\text{}}(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),b)`$.
###### Proposition 10.4.1.
The Hochschild differential $`b:\overline{C}_{n+1}(A,A)\overline{C}_n(A,A)`$, viewed as a map $`:\mathrm{\Omega }_{\mathrm{nc}}^{n+1}(A)\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$, is given by the following explicit formula
$$\delta (a_0da_1\mathrm{}da_na^{})=\delta (\alpha da_na^{})=(1)^{\mathrm{deg}\alpha }(\alpha a_na^{}\alpha a_na^{}).$$
###### Proof.
Direct calculation, see \[CQ1\]. ∎
As an application, we are going to construct the following exact sequence
$$HH_1(A)\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[A,\mathrm{\Omega }_{\mathrm{nc}}^1(A)]\stackrel{b}{}[A,A]0.$$
(10.4.2)
where the map $`b`$ is given by the assignment
$$b:\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[A,\mathrm{\Omega }_{\mathrm{nc}}^1(A)]A,udv[u,v]$$
(10.4.3)
We first verify that the map $`b`$ in (10.4.3) is well-defined, i.e., we have $`b([xdy,a])=0`$ for any $`[xdy,a][\mathrm{\Omega }_{\mathrm{nc}}^1(A),A]`$. Indeed, we compute
$`b([xdy,a])`$ $`=b(xdya)b(axdy)`$
$`=b(xd(ya))b(xyda)[ax,y]`$
$`=[x,ya][xy,a][ax,y]=([ya,x]+[xy,a]+[ax,y])`$
$`=x(ya)(ya)x(xy)a+a(xy)(ax)y+y(ax)=0.`$
Now, to construct (10.4.2) is to use a long exact sequence for Hochschild homology arising from the fundamental short exact sequence $`0\mathrm{\Omega }_{\mathrm{nc}}^1(A)A^\mathrm{e}A0`$:
$$\mathrm{}H_1(A,A^\mathrm{e})H_1(A,A)H_0(A,\mathrm{\Omega }_{\mathrm{nc}}^1(A))H_0(A,A^\mathrm{e})H_0(A,A)0.$$
We have $`H_0(A,A)=A/[A,A],`$ and $`H_0(A,\mathrm{\Omega }_{\mathrm{nc}}^1(A))=\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[A,\mathrm{\Omega }_{\mathrm{nc}}^1(A)].`$ Observe further that $`H_1(A,A^\mathrm{e})=0`$, since $`A^\mathrm{e}`$ is a free $`A`$-bimodule. Thus, since $`H_0(A,A^\mathrm{e})=A`$ by (10.1.5), the long exact sequence above reduces to a short exact sequence as in (10.4.2).
### 10.5. Relative differential forms
Let $`A`$ be an associative algebra, $`BA`$ a subalgebra, and $`M`$ an $`A`$-bimodule.
###### Definition 10.5.1.
A derivation $`f:AM`$ such that $`f(b)=0`$ for all $`bB`$ is said to be a derivation from $`A`$ to $`M`$ *relative to $`B`$*. We write $`\mathrm{Der}_B(A,M)\mathrm{Der}(A,M)`$ for the subspace of all such derivations, which form a Lie subalgebra.
Notice that a derivation is a relative one, namely it is relative to the subalgebra $`\mathrm{𝕜}1_AA`$. For ease of notation, we will denote this subalgebra simply by $`\mathrm{𝕜}`$.
### 10.6. The Commutative Case
Let $`A`$ be a commutative algebra, and led let $`M`$ be a left $`A`$-module. We make $`M`$ into an $`A`$-bimodule by equipping it with the symmetric bimodule structure. Then $`M\mathrm{Der}_B(A,M)`$ defines a functor $`\mathrm{Der}_B(A,):A\text{-}\mathrm{𝗆𝗈𝖽}\mathrm{𝖵𝖾𝖼𝗍}`$, where $`\mathrm{𝖵𝖾𝖼𝗍}`$ denotes the category of $`\mathrm{𝕜}`$-vector spaces. As in the case where $`B=\mathrm{𝕜}`$, $`\mathrm{Der}_B(A,)`$ can be represented by an $`A`$-module $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)`$, which we call the *relative Kähler differentials*.
If we regard $`A`$ and $`B`$ as the coordinate rings of affine varieties and set $`X=\mathrm{Spec}A`$ and $`Y=\mathrm{Spec}B`$, then the embedding $`BA`$ induces a surjection $`\pi :XY`$. Recall that we view $`\mathrm{Der}(A)`$ as the global sections of $`𝒯_X`$, the tangent bundle of $`X`$, that is, as the space $`\mathrm{\Gamma }(X,𝒯_X)`$ of algebraic vector fields on $`X`$. In this geometric picture, $`\mathrm{Der}_B(A)=\mathrm{Der}_B(A,A)`$ corresponds to the subspace in $`\mathrm{\Gamma }(X,𝒯_X)`$ formed by the vector fields on $`X`$ which are parallel to fibers of the projection $`\pi :XY`$. More precisely, the differential of the map $`\pi `$ may be viewed as a sheaf map $`\pi _{}:𝒯_X\pi ^{}𝒯_Y`$. Let $`𝒯_{X/Y}:=\mathrm{Ker}[𝒯_X\pi ^{}𝒯_Y],`$ be the sheaf whose sections are the vector fields tangent to the fibers. If $`\pi `$ has surjective differential, then $`\pi _{}`$ is a surjective map of locally free sheaves, hence, its kernel is again locally free, and we have a short exact sequence
$$0𝒯_{X/Y}𝒯_X\pi ^{}𝒯_Y\mathrm{\hspace{0.17em}0}.$$
Similarly, we view $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)`$ as the space of differential $`1`$-forms along the fibers, that is of sections of the relative cotangent bundle $`𝒯_{X/Y}^{}`$. Explicitly, dualizing the short exact sequence above, one gets the dual exact sequence
$$0\pi ^{}𝒯_Y^{}𝒯_X^{}𝒯_{X/Y}^{}\mathrm{\hspace{0.17em}0}.$$
(10.6.1)
This presents the space $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)=\mathrm{\Gamma }(X,𝒯_{X/Y}^{})`$ as a quotient of $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$.
We can also realize $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)`$ by relative Hochschild homology. In particular, it is easy to see (by the same arguments for absolute differential forms) that we can write
$$\mathrm{\Omega }_{\mathrm{com}}^1(A/B)=\mathrm{Tor}_1^{A_BA}(A,A).$$
This can be expressed in terms of the kernel $`I`$of the multiplication map $`A_BAA`$. Indeed, we find that $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)I/I^2`$, as before. This is defined to be the first relative Hochschild group, $`HH_1(A;B)`$. It is also easy to see (as in the absolute case) that $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)=HH_1(A;B)A\overline{A}`$, where now $`\overline{A}=A/B`$.
As in the absolute case, we define $`\mathrm{\Omega }_{\mathrm{com}}^{\text{}}(A/B)=T_B\mathrm{\Omega }_{\mathrm{com}}^1(A/B)`$. We then obtain the following result, analogous to the absolute case.
###### Theorem 10.6.2.
$`\mathrm{\Omega }_{\mathrm{com}}^{\text{}}(A/B)`$ is a graded commutative differential envelope of $`A`$, specifically, it is the universal object in the category whose objects are algebra maps $`f:AD^0D`$ such that $`d_Df|_B=0`$, where $`D`$ is a graded commutative DGA.
### 10.7. The Noncommutative Case
We now assume only that $`A`$ is associative. Consider the functor
$$\mathrm{Der}_B(A,):A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}\mathrm{𝖵𝖾𝖼𝗍}.$$
Again, this functor is representable, and we call its representing object $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$. Once again, $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$ is the kernel of the multiplication map, but this time viewed as a map $`A_BAA`$. That is, $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$ fits in the exact sequence of $`A`$-bimodules:
$$\text{}.$$
It is easy to see that $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)A_BA/B`$ as a $`B`$-bimodule. Indeed, the isomorphism is given by sending an element $`x\overline{y}`$ of $`A_BA/B`$ to $`xyxy1_A`$ (of course, one needs to check that this is even well-defined).
Recall that in the commutative case, we thought of $`\mathrm{\Omega }_{\mathrm{com}}^1(A/B)`$ as a relative cotangent bundle $`\mathrm{\Gamma }(X,𝖳_{X/Y}^{})`$, where $`X=\mathrm{Spec}A`$ and $`Y=\mathrm{Spec}B`$, and $`\pi :XY`$ is the surjection induced by the inclusion $`BA`$. Similarly to the exact sequence (10.6.1), we have the following exact sequence in the noncommutative case
$`0\mathrm{Tor}_1^B(A,A)A_B\mathrm{\Omega }_{\mathrm{nc}}^1(B)_BA\mathrm{\Omega }_{\mathrm{nc}}^1(A)\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)0`$
There is a connection between $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$ and differential envelopes of $`A`$. Indeed, we have the important following result.
###### Theorem 10.7.1.
Let $`D^{\text{}}(A/B)`$ denote the relative differential envelope of $`A`$. That is, $`D^{\text{}}(A/B)`$ is universal in the category whose objects are algebra maps $`f:AD^0D`$ where $`f|_B=0`$ and $`D`$ is a DGA. Then
$$D^{\text{}}(A/B)T_A^{\text{}}\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)A_BT_B^{\text{}}(A/B).$$
As usual, we set $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B)=T_A\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$.
###### Proof.
We will prove that $`D^{\text{}}(A/B)\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B)`$. Of course, this proof will also be valid in the case $`B=\mathrm{𝕜}`$. We will complete the proof in two steps. First, we will observe that $`D^1(A/B)`$ and $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$ are isomorphic. We will then show that $`D^{\text{}}(A/B)T_AD^1`$.
For the first step, notice that $`D^1`$ is an $`A`$-bimodule, since we have a map $`i:AD^0`$. If we compose $`i`$ with the differential $`d`$ on $`D`$ (and denote this map by $`d`$), we obtain an $`A`$-bimodule derivation $`d:AD^1`$. Since $`i`$ vanishes on $`B`$, so does $`d`$. We wish to show that $`d:AD^1`$ represents $`\mathrm{Der}_B(A,)`$, which will show that $`D^1`$ and $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$ are isomorphic by the latter’s universal property. Let $`\delta \mathrm{Der}_B(A,M)`$, where $`M`$ is any $`A`$-bimodule. We wish to complete the following diagram
$$\text{}.$$
We shall perform this by using the square zero construction. Recall that we set $`A\mathrm{}M`$ to be, as a $`\mathrm{𝕜}`$-vector space, the direct sum $`AM`$. The product on $`A\mathrm{}M`$ is defined in such a way as to make $`A`$ a subalgebra and $`M^2=0`$. We make $`A\mathrm{}M`$ into a DGA by declaring that $`\mathrm{deg}A=0`$, $`\mathrm{deg}M=1`$, and by setting $`d(a,m)=(0,\delta m)`$. Define $`f:AA\mathrm{}M`$ by $`f(a)=(a,0)`$. Observe that $`df(b)=0`$ for all $`bB`$, hence we obtain a DGA map $`f^{}:D(A/B)A\mathrm{}M`$ from the universal property of $`D(A/B)`$. Consider the restriction of $`f^{}`$ to $`D^1(A/B)`$. This must map $`D^1(A/B)`$ into $`M=\{0\}\mathrm{}M`$, since only $`M`$ has degree one. Then using the fact that $`f^{}`$ is a DGA map, we find for all $`aA`$, $`f^{}(d(a))=\delta (a)`$. So, $`D^1(A/B)`$ satisfies the universal property for $`\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$, hence these two $`A`$-bimodules are isomorphic.
Now we wish to show that $`D`$ is generated as the free algebra over $`D^1`$. The inclusion of $`D^1D`$ as the degree one elements induces an algebra map $`\phi :T_AD^1D`$. We need only check that $`\phi `$ is an isomorphism. Here are two proofs of this fact. First, we adopt the notation
$$T_A^nD^1=(D^1)^n.$$
Then
$$T_A^nD^1=T_A^{n1}D^1_AD^1T_A^{n1}D^1_A(A_BA/B),$$
since we know that $`D^1\mathrm{\Omega }_{\mathrm{nc}}^1(A/B)`$. But then
$$T_A^{n1}D^1_A(A_BA/B)=T_A^{n1}D^1_BA/BD^n.$$
The second proof is in some way more elementary, as we simply construct an inverse $`\psi :DT_AD^1`$ to $`\phi `$. We recall that $`D`$ contains $`A`$ (since $`D`$ comes with an embedding $`AD^0`$), and $`D`$ also contains $`\overline{A}=dA=A/B`$. So, we consider the algebra
$$D^{}=T_\mathrm{𝕜}(A+\overline{A})/\{aa^{}aa^{}1_A,\overline{aa^{}}=\overline{a}a^{}+a\overline{a^{}}\}.$$
So $`D^{}`$ is simply the free algebra over $`A`$ and $`\overline{A}`$, modulo the relations which give the desired product (as in $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B)`$) and make $`a\overline{a}`$ a derivation. Clearly, $`D^{}`$ is precisely $`D`$. It is also easy to see that there is a map $`D^{}T_AD^1`$, and all relations are mapped to zero. So, $`D`$ and $`T_AD^1`$ are isomorphic. ∎
## 11. Noncommutative Calculus
### 11.1.
Let $`A`$ be an associative unital $`\mathrm{𝕜}`$-algebra, thought of as the coordinate ring of a ‘noncommutative space’. Then, any automorphism $`F:AA`$ may be thought of as an automorphism of that ‘noncommutative space’. Similarly, a derivation $`\theta :AA`$ may be viewed as an ‘infinitesimal automorphism’ of $`A`$, in the sense that if the linear map $`\theta :AA`$ could have been exponentiated, i.e., if the infinite series $`\mathrm{exp}(\theta )=\mathrm{id}_A+\theta +\frac{1}{2}\theta {}_{^{^{}}}{}^{}\theta +\frac{1}{3}\theta {}_{^{^{}}}{}^{}\theta {}_{^{^{}}}{}^{}\theta +\mathrm{}`$ made sense as a map $`AA`$, then a formal computation shows that the map $`\mathrm{exp}(\theta )`$ would have been an automorphism of $`A`$. Geometrically, one thinks of $`\theta `$ as a ‘vector field’ on a noncommutative space; then $`t\mathrm{exp}(t\theta )`$ is the one-parameter flow of automorphisms of that noncommutative space generated by our vector field.
Accordingly, an algebra automorphism $`F:AA,`$ resp., a derivation $`\theta :AA`$, induces an automorphism, resp., a Lie derivative endomorphism $`_\theta `$, of any of the objects $`\mathrm{Der}(A)`$, $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ and $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$.
We provide more details. Fix a derivation $`\theta :AA`$. Since the functor $`\mathrm{Der}(A,)`$ is represented by the bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$, we deduce that there is a unique $`A`$-bimodule homomorphism $`i_\theta :\mathrm{\Omega }_{\mathrm{nc}}^1(A)A`$ corresponding to $`\theta `$. Observe further that since $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)=T_A\mathrm{\Omega }_{\mathrm{nc}}^1(A),`$ is a free algebra of the bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1(A),`$ there is a unique way to extend the map $`i_\theta :\mathrm{\Omega }_{\mathrm{nc}}^1(A)A`$ to a degree $`(1)`$ super-derivation $`i_\theta :\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ of the algebra $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$. Explicitly, we have
$$i_\theta (a_0da_1\mathrm{}da_n)=\underset{j=1}{\overset{n}{}}(1)^{j1}a_0da_1\mathrm{}\theta (a_j)\mathrm{}da_n.$$
Recall next that, for any two super-derivations $`_1,_2`$ of odd degree, their super-commutator $`[_1,_2]:=_1{}_{^{^{}}}{}^{}_{2}^{}+_2{}_{^{^{}}}{}^{}_{1}^{}`$ is an even degree derivation. In particular, for $`_1=d`$ and $`_2=i_\theta `$, we obtain a degree zero derivation $`d{}_{^{^{}}}{}^{}i_{\theta }^{}+i_\theta {}_{^{^{}}}{}^{}d`$. Similarly, for $`_1=_2=d`$, resp., for $`_1=_2=i_\theta `$, we obtain a degree $`(+2)`$-derivation $`d^2`$, resp., a degree $`(2)`$-derivation $`(i_\theta )^2`$.
Next we apply Lemma 3.3.2 to $`S=\mathrm{\Omega }_{\mathrm{nc}}^1(A)\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)=R`$. It is straightforward to get from definitions that, for any $`s\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ one has $`d^2(s)=0=(i_\theta )^2(s).`$ It follows that $`d^2=0`$ and $`(i_\theta )^2=0`$ identically on $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$.
Similarly, we define the Lie derivative with respect to $`\theta \mathrm{Der}(A)`$ as a map $`_\theta :\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ given by the formula
$$_\theta (a_0da_1\mathrm{}da_n)=\theta (a_0)da_1\mathrm{}da_n+\underset{j=1}{\overset{n}{}}a_0da_1\mathrm{}d\theta (a_j)\mathrm{}da_n.$$
A direct calculation shows that for any $`s\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ one has $`_\theta (s)=[d,i_\theta ](s)`$. It follows by Lemma 3.3.2 that the following Cartan formula holds on $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$
$$_\theta =di_\theta +i_\theta d.$$
Similarly, one verifies the following identities
$$[_\theta ,_\gamma ]=_{[\theta ,\gamma ]},[_\theta ,i_\gamma ]=i_{[\theta ,\gamma ]},i_\theta ^2=0.$$
(11.1.1)
To prove these identities, observe that in each case, the identity in question is obvious on the elements of $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$. Hence, using the same argument as above, we deduce that it holds on the whole of $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$.
### 11.2. Operations on Hochschild complexes.
Let $`C_{\text{}}(A,A),C_k(A,A)=A^k,`$ be the Hochschild chain complex for $`A`$. For any Hochschild $`p`$-cochain $`cC^p(A,A)`$ and $`kp`$, define a contraction operator $`i_c:C_k(A,A)C_{kp}(A,A)`$ by the formula
$$i_c:a_0\mathrm{}a_kc(a_1,\mathrm{},a_p)a_{p+1}\mathrm{}a_k.$$
Now suppose we have a derivation $`\delta :AA`$. Then we can extend $`\delta `$ to a derivation on each $`𝖡_n`$ in the bar complex, namely
$$\delta (a_1\mathrm{}a_n)=\underset{j=1}{\overset{n}{}}a_1\mathrm{}\delta (a_j)\mathrm{}a_n.$$
Then $`\delta :\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A^n,A)\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A^n,A)`$. It is not hard to see that $`\delta `$ commutes with the bar differential, hence it induces a derivation on the Hochschild cohomology of $`A`$.
We can generalize the above to all $`p1`$, and define the Lie derivative with respect to a cochain $`cC^p(A,A)`$ as an operator $`_c:C_k(A,A)C_{kp+1}(A,A)`$ by the formula
$`_c:a_0\mathrm{}a_k`$ $`{\displaystyle \underset{i=0}{\overset{kp}{}}}(1)^{(p1)(i+1)}a_0\mathrm{}a_ic(a_{i+1},\mathrm{},a_{i+p})\mathrm{}a_k`$
$`+{\displaystyle \underset{i=kp}{\overset{k}{}}}(1)^{k(j+1)}c(a_{j+1},\mathrm{},a_0,\mathrm{})a_{p+jk}\mathrm{}a_j.`$
###### Remark 11.2.1.
For the 2-cochain $`m:a,bab,`$ given by the product in the algebra $`A`$, the operator $`_m`$ is nothing but the Hochschild differential on $`C_{\text{}}(A,A)`$. $`\mathrm{}`$
Recall the canonical noncommutative Gerstenhaber algebra structure on the Hochschild cochain complex $`C^{\text{}}(A,A)`$. Recall further that to any (not necessarily commutative) Gerstenhaber algebra $`G^{\text{}}`$ one can associate another Gerstenhaber algebra $`G_\epsilon ^{\text{}},`$ called the $`\epsilon `$-construction. We apply the $`\epsilon `$-construction to the Gerstenhaber algebra $`C^{\text{}}(A,A)`$.
We have the following noncommutative analogue of Proposition 6.4.3.
###### Theorem 11.2.2.
For any $`b,cC^p(A,A)`$, the following formulas
$$(b+\epsilon c)_\epsilon \alpha :=(1)^{\mathrm{deg}b}i_b\alpha ,\text{and}\{b+\epsilon c,\alpha \}_\epsilon :=_b\alpha +\epsilon i_c\alpha $$
make $`C_{\text{}}(A,A)_\epsilon `$ a Gerstenhaber module over $`C^{\text{}}(A,A)_\epsilon .`$
### 11.3. The functor of ‘functions’
Let $`A`$ be a not necessarily commutative associative algebra thought of as the coordinate ring of a ‘noncommutative scheme’. We would like to introduce a vector space $`𝖱(A)`$ playing the role of ‘the space of regular functions’ on that scheme. Of course, if $`A=\mathrm{𝕜}[X]`$ is the coordinate ring of an ordinary commutative scheme $`X`$, then regular functions on $`X`$ are by definition the elements of $`A`$. So, one might guess that, in the noncommutative case, the equality $`𝖱(A)=A`$ still holds. This cannot be quite right, however. Indeed, one expects the space $`𝖱(A)`$ to be a Morita invariant of $`A`$ since only Morita invariant notions are ‘geometrically meaningful’. Thus, starting from an ordinary commutative scheme $`X`$, for any $`n=1,2,\mathrm{},`$ we may form the algebra $`A=\mathrm{Mat}_n\mathrm{𝕜}\mathrm{𝕜}[X]`$, which is Morita equivalent to $`\mathrm{𝕜}[X]`$. Thus, we would like our definition of the space $`𝖱(A)`$ be such that the following holds
$$𝖱(\mathrm{Mat}_n\mathrm{𝕜}\mathrm{𝕜}[X])=𝖱(\mathrm{𝕜}[X])=\mathrm{𝕜}[X].$$
By Morita invariance of Hochschild homology, see Proposition 5.2.1, this requirement is satisfied if we introduce the following
###### Definition 11.3.1.
We define the *space of functions* associated to an associative algebra $`A`$ to be the vector space
$$𝖱(A):=A/[A,A].$$
### 11.4. Karoubi-de Rham complex.
We are going to show that for any associative algebra $`A`$ with unit, the complex $`(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),d)`$ of noncommutative differential forms has trivial cohomology:
$$H^i(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),d)=\{\begin{array}{cc}\mathrm{𝕜}\hfill & \text{if}i=0\hfill \\ 0\hfill & \text{if}i>0.\hfill \end{array}$$
(11.4.1)
To see this, recall the isomorphism $`\mathrm{\Omega }_{\mathrm{nc}}^p(A)=A\overline{A}^p,`$ where $`\overline{A}=A/\mathrm{𝕜}`$. The differential $`d`$ corresponds to the natural projection
$$A\overline{A}^p\overline{A}\overline{A}^p\mathrm{𝕜}\overline{A}^{(p+1)}A\overline{A}^{(p+1)}.$$
The kernel of the latter projection clearly equals $`\mathrm{𝕜}\overline{A}^p`$, which is exactly the image of the differential $`d:A\overline{A}^{p1}\mathrm{𝕜}\overline{A}^p`$. This proves (11.4.1).
Therefore, the differential $`d`$ on $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ does not give rise interesting cohomology theory. Things become better with the following definition.
###### Definition 11.4.2.
The *noncommutative de Rham complex* of $`A`$ is a graded vector space defined by
$$\mathrm{DR}^{\text{}}(A):=𝖱\left(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\right)=\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)/[\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)],$$
where $`[,]`$ denotes the super-commutator.
The differential $`d:\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\mathrm{\Omega }_{\mathrm{nc}}^{+1}(A)`$ descends to a well-defined differential $`d:\mathrm{DR}^{\text{}}(A)\mathrm{DR}^{+1}(A)`$, making the de Rham complex a differential graded $`\mathrm{𝕜}`$-vector space.
For example, $`\mathrm{DR}^0(A)=A/[A,A]`$, since only the degree zero terms of $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ contribute. Similarly, we have
$$\mathrm{DR}^1(A)=\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[A,\mathrm{\Omega }_{\mathrm{nc}}^1(A)]=HH_0(A,\mathrm{\Omega }_{\mathrm{nc}}^1(A)).$$
(11.4.3)
Hence, from (10.4.2), we obtain a canonical short exact sequence
$$0HH_1(A)\mathrm{DR}^1(A)\stackrel{b}{}[A,A]0.$$
(11.4.4)
For $`k>1`$, the relation between de Rham complex and Hochschild homology is more complicated.
One can check that the operations $`d,_\theta ,`$ and $`i_\theta `$ on noncommutative differential forms all descend to the de Rham complex.
Let $`A`$ be a smooth associative and commutative algebra. Then, by Hochschild-Kostant-Rosenberg theorem we have a vector space (but not necessarily algebra) isomorphism $`\mathrm{\Omega }_{\mathrm{com}}^{\text{}}(A)=HH_{\text{}}(A)`$. Obviously, the Hochschild differential is not the de Rham differential, since the de Rham differential increases degree while the Hochschild homology differential decreases degree. Indeed, there is another differential (the Connes differential) which yields the de Rham differential (and cyclic homology).
Suppose now that $`BA`$ is a subalgebra. The *relative de Rham complex* of $`A`$ is
$$\mathrm{DR}(A/B)=𝖱(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B))=\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B)/[\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B),\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A/B)],$$
where now $`[,]`$ is the graded commutator.
Strictly speaking, to be consistent with this notation, one has to write $`\mathrm{DR}(A/\mathrm{𝕜})`$ rather than $`\mathrm{DR}(A)`$ in the ‘absolute’ case.
The de Rham complex is Morita invariant in the following sense.
###### Proposition 11.4.5.
Let $`A`$ be an associative algebra. Then for any $`n`$, there is a canonical isomorphism
$$\mathrm{DR}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})\mathrm{DR}(A).$$
###### Proof.
Let $`A`$ and $`B`$ be arbitrary associative algebras. Then it is clear that there is a canonical isomorphism $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}\left((BA)/(1B)\right)B\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$. Also, it is clear that $`\mathrm{Mat}_nA\mathrm{Mat}_n\mathrm{𝕜}A`$. So, we find that
$$\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})\mathrm{Mat}_n\mathrm{𝕜}\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\mathrm{Mat}_n(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)).$$
But our comments regarding $`𝖱(\mathrm{Mat}_n\mathrm{𝕜})`$ show that
$$𝖱(\mathrm{Mat}_n(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)))𝖱(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A))=\mathrm{DR}(A).$$
Thus, $`𝖱\left(\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})\right)=\mathrm{DR}(A),`$ and we are done. ∎
###### Remark 11.4.6.
One can see by looking through the proof above that the map $`\mathrm{DR}(A)`$ $`\mathrm{DR}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})`$ that yields the isomorphism of the Proposition is induced by the algebra imbedding $`A{}_{}{}^{}_{}^{}\mathrm{Mat}_nA,a\left(\begin{array}{cc}a& 0\\ 0& 1_{n1}\end{array}\right).`$ The inverse isomorphism $`\mathrm{DR}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})\stackrel{_{}}{}`$ $`\mathrm{DR}(A)`$ is induced by the ‘trace-map’ sending $`x=x_{ij}\mathrm{Mat}_nA`$ to $`\mathrm{Tr}(x)=x_{ii}`$.
Next, let $`A=_{i0}A_i`$ be a graded $`\mathrm{𝕜}`$-algebra. Write $`𝔾_𝐦`$ for the multiplicative group, viewed as an algebraic group over $`\mathrm{𝕜}`$.
Giving a $``$-grading on an associative $`\mathrm{𝕜}`$-algebra $`A`$ is the same thing as giving an algebraic $`𝔾_𝐦`$-action on $`A`$ by algebra automorphisms. Specifically, given a grading $`A=_iA_i`$ one defines a $`𝔾_𝐦`$-action by the formula $`𝔾_𝐦\times A_it,at^ia,`$ for any $`i`$. Conversely, it is easy to see that any $`𝔾_𝐦`$-action on $`A`$ by algebra automorphisms arises in this way from a certain $``$-grading on $`A`$.
Assume now that $`A=_{i0}A_i`$ is graded by nonnegative integers. Geometrically, this means that the corresponding $`𝔾_𝐦`$-action is a ‘contraction’ of $`A`$ to the subalgebra $`A_0`$.
The result below says that de Rham cohomology is ‘invariant under contraction’.
###### Theorem 11.4.7 (Poincaré lemma).
For a graded algebra $`A=_{i0}A_i,`$ the algebra imbedding $`A_0{}_{}{}^{}_{}^{}A`$ induces isomorphisms
$$H^j(\mathrm{DR}(A_0))\stackrel{_{}}{}H^j(\mathrm{DR}(A)),j0.$$
###### Proof.
The assignment $`A_iaia,i=0,1,\mathrm{},`$ gives a derivation of $`A`$, called the *Euler derivation*. This derivation may be thought of as an infinitesimal generator of the $`𝔾_𝐦`$-action on $`A`$ corresponding to the grading.
Associated with the Euler derivation, one has $`_{\mathrm{𝖾𝗎}}`$, the Lie derivative with respect to $`\mathrm{𝖾𝗎}`$, acting on $`\mathrm{DR}(A)`$. The action of $`_{\mathrm{𝖾𝗎}}`$ on $`\mathrm{DR}(A)`$ is diagonalizable with nonnegative integral eigenvalues, and we write $`\mathrm{DR}(A)=_{m0}\mathrm{DR}(A)m`$ for the corresponding eigenspace direct sum decomposition. It is clear that we have $`\mathrm{DR}(A)0=\mathrm{DR}(A_0).`$
The de Rham differential $`d`$ commutes with $`_{\mathrm{𝖾𝗎}}`$, hence preserves the direct sum decomposition above. Further, the homotopy formula $`_{\mathrm{𝖾𝗎}}=di_{\mathrm{𝖾𝗎}}+i_{\mathrm{𝖾𝗎}}d`$, shows that $`i_{\mathrm{𝖾𝗎}}`$ is a chain homotopy between the map $`_{\mathrm{𝖾𝗎}}`$ and the zero map. Hence, the complex $`(\mathrm{DR}(A)m,d)`$ is acyclic for all $`m`$ except $`m=0`$. This proves the Theorem. ∎
### 11.5. The Quillen sequence.
Recall that $`\mathrm{DR}^1(A)=\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[\mathrm{\Omega }_{\mathrm{nc}}^1(A),A]`$, see (11.4.3), and define a map $`b:\mathrm{DR}^1(A)A`$ by $`b:xdy[x,y]`$. We have shown that this map is well-defined, see (10.4.3).
It is easy to check that $`bd=db=0`$. Thus $`d`$ and $`b`$ yield a $`2`$-periodic complex
$$\text{}.$$
(11.5.1)
This complex gives an approximation to the cyclic homology of the algebra $`A`$, see \[CQ2\]
Put $`\overline{\mathrm{DR}}^0(A)=A/([A,A]+\mathrm{𝕜})`$, let $`\mathrm{pr}:\overline{A}\overline{A}/[A,A]\overline{\mathrm{DR}}^0(A)`$ be the natural projection.
Consider the following sequence of maps, called the Quillen sequence:
$$\text{},$$
(11.5.2)
where the maps $`d`$ and $`b`$ have been introduced above. Since $`bd=0`$ and the image of $`b`$ is contained in $`[A,A],`$ it is clear that the composite of any two consequtive maps in the sequence is equal to zero, i.e., the sequence is a complex.
Consider the case $`A=T(V^{})`$, where $`V`$ is a finite-dimensional $`\mathrm{𝕜}`$-vector space.
###### Lemma 11.5.3.
If $`A=T(V^{})`$, then Quillen’s complex (11.5.2) is an exact sequence.
###### Proof.
We already know that, for $`A=T(V^{})`$, the de Rham complex $`\mathrm{DR}^{\text{}}(A)`$ is acyclic in positive degrees, and $`\mathrm{DR}^0(A)=\mathrm{𝕜}`$. It follows that the map $`\overline{\mathrm{DR}}^0(A)\mathrm{DR}^1(A)`$ in (11.5.2) is injective. Further, we have $`\mathrm{Im}(b)=[A,A]`$, hence the complex is exact at $`\overline{A}`$ and at $`\overline{\mathrm{DR}}^0(A)`$. Thus, it remains only to check that the complex is exact at $`\mathrm{DR}^1(A)`$.
To this end, it suffices to show that the class
$$\left[\overline{\mathrm{DR}}^0(A)\right]\left[\mathrm{DR}^1(A)\right]+\left[\overline{A}\right]\left[\overline{\mathrm{DR}}^0(A)\right]$$
vanishes in the Grothendieck group of graded spaces. The two terms of $`\left[\overline{\mathrm{DR}}^0(A)\right]`$ cancel to leave
$$\left[\mathrm{DR}^1(A)\right]+\left[\overline{A}\right].$$
But $`\mathrm{DR}^1(A)AV^{}=T(V^{})V^{}`$, while $`\overline{A}T(V^{})V^{}`$, as well. Hence, the remaining two terms cancel as well and the sequence is exact. ∎
Let $`x_i`$ and $`x^i`$ be dual bases in $`V`$ and $`V^{}`$.
###### Proposition 11.5.4.
Let $`A=T(V^{})`$.
$`(𝗂)`$For any $`f\mathrm{DR}^0(A)`$, one has the identity $`_{i=1}^n[\frac{f}{x^i},x_i]=0,`$ in $`[A,A]`$.
$`(\mathrm{𝗂𝗂})`$The space of closed forms in $`\mathrm{DR}_{closed}^2\mathrm{DR}^2(A)`$ is canonically isomorphic to $`[A,A]`$.
Proof. The proof uses the Quillen sequence,
$$0\overline{\mathrm{DR}}^0(A)\mathrm{DR}^1(A)[A,A]0,$$
which is exact for $`A=T(V^{})`$. We also have the following sequence, which is exact by the Poincaré lemma,
$$0\overline{\mathrm{DR}}^0(A)\mathrm{DR}^1(A)\mathrm{DR}_{closed}^2(A)0.$$
Combining these two exact sequences, we have
$$\text{}.$$
But then we can use standard diagram chase arguments to construct a map $`[A,A]\mathrm{DR}^2(A)`$, and this must be an isomorphism. This shows claim (2). As for (1), observe that
$$_{i=1}^n[\frac{f}{x^i},x_i]=bd(f)=0.\mathrm{}$$
### 11.6. The Karoubi Operator.
Define the *Karoubi operator* $`\kappa :\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ by
$$\kappa (\alpha da)=(1)^{\mathrm{deg}\alpha }da\alpha ,\text{and}\kappa (\alpha )=\alpha \text{if}\alpha =a\mathrm{\Omega }_{\mathrm{nc}}^0(A)=A.$$
The Karoubi operator provides the following the relation between $`b`$ and $`d`$:
$$db+bd=\mathrm{id}\kappa .$$
If we compare this to differential geometry and think of $`b`$ as $`d^{}`$, the adjoint of the de Rham differential with respect to a euclidean structure on the space of differential forms, then $`\kappa `$ is playing the role of the Laplace operator.
From a different viewpoint, observe that there is an obvious $`n+1`$-cycle action on the bar complex, namely given by the obvious cyclic action on $`A^{(n+1)}`$. Since the reduced bar complex is $`A\overline{A}^n\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$, there is no clear $`(n+1)`$-cycle action, and the operation $`\kappa `$ “approximates” an $`n`$-cycle action.
###### Proposition 11.6.1.
$`(𝗂)`$$`\kappa ^{n+1}d=d`$.
$`(\mathrm{𝗂𝗂})`$$`\kappa ^n=\mathrm{id}+b\kappa ^nd`$.
$`(\mathrm{𝗂𝗂𝗂})`$$`\kappa ^{n+1}=1db`$.
Proof.$`(𝗂)`$This is trivial. To prove $`(\mathrm{𝗂𝗂})`$we calculate
$`\kappa ^n(a_0da_1\mathrm{}da_n)`$ $`=(da_1\mathrm{}da_n)a_0`$
$`=a_0da_1\mathrm{}da_n+[da_1\mathrm{}da_n,a_0]`$
$`=a_0da_1\mathrm{}da_n+(1)^nb(da_1\mathrm{}da_nda_0)`$
$`=(\mathrm{id}+b\kappa ^nd)(a_0da_1\mathrm{}da_n).`$
$`(\mathrm{𝗂𝗂𝗂})`$Multiply (ii) on the right by $`\kappa `$ and observe that $`\kappa `$ commutes with $`b`$ and $`d`$.
$`\kappa ^{n+1}`$ $`=\kappa +b\kappa ^nd\kappa =\kappa +b\kappa ^{n+1}d`$
$`=\kappa +bd\text{by (1)}`$
$`=\kappa +(\mathrm{id}\kappa db)=\mathrm{id}db.\mathrm{}`$
###### Proposition 11.6.2.
On $`\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$, we have $`(\kappa ^n1)(\kappa ^{n+1}1)=0.`$
###### Proof.
First, we have that $`\kappa ^n1=b\kappa ^nd`$ by (ii). By (iii), $`\kappa ^{n+1}1=db`$. So,
$$(\kappa ^n1)(\kappa ^{n+1}1)=b\kappa ^nd^2b=0,$$
since $`d^2=0`$. ∎
### 11.7. Harmonic decomposition.
Notice that $`(n,n+1)=1`$. So, the polynomial $`(t^n1)`$ $`(t^{n+1}1)`$ has only simple roots except for a double root at one. The identity $`(\kappa ^n1)(\kappa ^{n+1}1)=0`$ implies that the action of $`\kappa `$ is locally-finite, and all of its eigenvalues have multiplicity one except for the eigenvalue $`1`$, which has multiplicity 2. So,
$$\mathrm{\Omega }_{\mathrm{nc}}^n(A)=\left[\mathrm{Ker}(\kappa 1)^2\right]\left[_{\lambda \mathrm{Spec}(\kappa )\{1\}}\mathrm{Ker}(\kappa \lambda )\right].$$
The space $`\mathrm{Ker}(\kappa 1)^2`$ is called the space of *harmonic* forms, denoted by $`\mathrm{Harm}`$. The remaining summand is denoted by $`\mathrm{Harm}^{}\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$.
###### Proposition 11.7.1.
$$\mathrm{Harm}^{}=d(\mathrm{Harm}^{})b(\mathrm{Harm}^{}).$$
###### Proof.
Observe that $`db+bd=1\kappa `$ is invertible on $`\mathrm{Harm}^{}`$. Let $`G`$ be its inverse. Then $`G`$ commutes with $`b`$ and $`d`$. So, $`Gdb:\mathrm{Harm}^{}d(\mathrm{Harm}^{})`$ and $`Gbd:\mathrm{Harm}^{}b(\mathrm{Harm}^{})`$ are both projectors. This yields the direct sum decomposition. ∎
### 11.8. Noncommutative polyvector fields.
Recall the notation $`\overline{A}=A/\mathrm{𝕜}`$; we have the reduced bar complex
$$\mathrm{}A\overline{A}^{(n+1)}AA\overline{A}^nAA\overline{A}^{(n1)}A\mathrm{}.$$
We now consider the reduced cochain complex, that is, for each $`n`$ and each $`A`$-bimodule $`M`$ we set
$$\overline{C}^n(A,M)=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A\overline{A}^nA,M),$$
where $`\overline{C}`$ denotes the reduced cochains. We let $`\overline{Z}^n(A,M)`$ denote the space of reduced $`n`$-cocycles.
###### Proposition 11.8.1.
Suppose $`fC^n(A,M)`$. Then $`f`$ is a cocycle if and only if the map $`da_1\mathrm{}da_nf(1a_1\mathrm{}a_n1)`$ extends to an $`A`$-bimodule map $`\mathrm{\Omega }_{\mathrm{nc}}^n(A)M`$.
###### Proof.
Suppose we are given some $`f\overline{C}^n(A,M)`$. Write $`\overline{f}\mathrm{Hom}_\mathrm{𝕜}(\overline{A}^n,M)`$ for the map $`\overline{f}(\omega )=f(1\omega 1)`$ for all $`\omega \overline{A}^n`$. If $`f`$ is a cocycle (i.e., if $`df=0`$), then
$$f^{}(da_1\mathrm{}da_n)=\overline{f}(a_1\mathrm{}a_n)$$
extends $`\overline{f}`$ uniquely to an $`A`$-bimodule map $`f^{}:\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)M`$. The left $`A`$-linearity of $`f^{}`$ is trivial, so the cocyclicity condition is then becomes equivalent to right $`A`$-linearity. ∎
Write $`\overline{Z}^p(A,M)`$ for the $`\mathrm{𝕜}`$-vector space of cocycles in the reduced cochain complex $`\overline{C}^p(A,M)`$.
###### Corollary 11.8.2.
There is a canonical vector space isomorphism
$$\overline{Z}^p(A,M)=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^p(A),M),$$
i.e., the $`A`$-bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$ represents the functor $`M\overline{Z}^n(A,M)`$.
In particular, if $`\theta \mathrm{Der}A\mathrm{Der}(T(A^{}))`$, then it extends uniquely to yield a $`1`$-cocycle $`\theta Z^1(A,A)`$. So, every derivation yields a cocycle. Thus, the above Corollary is a generalization to $`p1`$ of the interpretation of the space $`\overline{Z}^1(A,M)`$ as the space $`\mathrm{Der}(A,M)=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^1(A),M)`$, of derivations from $`A`$ to $`M`$ (where the latter equality follows from the universal property of $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$).
###### Definition 11.8.3.
For any $`p1`$, we set $`\mathrm{\Theta }_{\mathrm{nc}}^p(A):=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^n(A),A)`$, and call elements of the graded space $`\mathrm{\Theta }_{\mathrm{nc}}^{\text{}}(A):=_{p1}\mathrm{\Theta }_{\mathrm{nc}}^p(A)`$ *noncommutative polyvector fields* on $`A`$.
By definition, there is a natural pairing $`\mathrm{\Theta }_{\mathrm{nc}}^p(A)\mathrm{\Omega }_{\mathrm{nc}}^p(A)A`$.
###### Proposition 11.8.4.
For any associative algebra $`A`$, there is a natural graded Lie super-algebra structure on $`\mathrm{\Theta }_{\mathrm{nc}}^1(A)`$, i.e., a super-bracket such that
$$[\mathrm{\Theta }_{\mathrm{nc}}^p(A),\mathrm{\Theta }_{\mathrm{nc}}^q(A)]\mathrm{\Theta }_{\mathrm{nc}}^{p+q1}(A).$$
###### Remark 11.8.5.
The degrees given above are precisely those for the usual Schouten bracket. $`\mathrm{}`$
###### Proof of Proposition 11.8.4..
Inside the Hochschild cochain complex $`C^p(A,A)=`$ $`\mathrm{Hom}_\mathrm{𝕜}(A^p,A)`$ we have a subcomplex of reduced cochains $`\overline{C}^p(A,A)=\mathrm{Hom}_\mathrm{𝕜}(\overline{A}^p,A)`$. Thus, a cochain is a reduced cochain provided it vanishes whenever at least one entry is a scalar. It follows easily the reduced cochain complex is preserved by the Gerstenhaber bracket. It is also true, although not quite so transparent, that the Gerstenhaber bracket is compatible with the Hochschild differential. In particular, it preserves cocycles. So, if we let $`\overline{Z}^p`$ denote the reduced $`p`$-cocycles, we obtain a bracket
$$\overline{Z}^p\times \overline{Z}^q\overline{Z}^{p+q1}.$$
But by Proposition 11.8.1 we can extend a reduced $`p`$-cocycle $`\omega `$ to a bimodule map from $`\mathrm{\Omega }_{\mathrm{nc}}^p(A)`$ to $`A`$. An element of $`\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(\mathrm{\Omega }_{\mathrm{nc}}^p(A),A)`$ is called a *$`p`$-vector field*. The Gerstenhaber bracket then yields a Lie super-algebra structure on $`p`$-vector fields. ∎
###### Question 11.8.6.
$`(𝗂)`$Does the natural $`\mathrm{\Theta }_{\text{}}(A)`$-action on $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ descend to $`\mathrm{DR}^{\text{}}(A)`$ ?
$`(\mathrm{𝗂𝗂})`$Given $`c\mathrm{\Theta }_p(A)`$, when does the map $`i_c:\mathrm{\Omega }_{\mathrm{nc}}^n(A)\mathrm{\Omega }_{\mathrm{nc}}^{np}(A)`$ descend to a well-defined map $`\mathrm{DR}^n(A)\mathrm{DR}^{np}(A)`$ ?
## 12. The Representation Functor
### 12.1.
It is believed that the noncommutative geometry of an associative algebra $`A`$ is ‘approximated’ (in certain cases) by the (commutative) geometry of the scheme $`\mathrm{Rep}_n^A`$ of $`n`$-dimensional representations of $`A`$. Moreover, it is expected that this approximation becomes ‘better’ once the integer $`n`$ gets larger.
###### Remark 12.1.1.
There exist noncommutative associative algebras $`A`$ that do not have any finite dimensional representations at all. Thus the idea of looking at finite dimensional representations has obvious limitations. $`\mathrm{}`$
Suppose $`A`$ is a finitely generated associative algebra and $`E`$ is a finite dimensional $`\mathrm{𝕜}`$-vector space. Let $`\mathrm{Rep}_E^A`$ denote the affine (not necessarily reduced) scheme of all $`\mathrm{𝕜}`$-linear algebra maps $`A\mathrm{End}_\mathrm{𝕜}E`$, see below for a rigorous definition of the scheme structure.
Given $`aA`$, for each representation $`\rho \mathrm{Rep}_E^A`$, the element $`\rho (a)`$ is a $`\mathrm{𝕜}`$-linear endomorphism of $`E`$. The assignment $`\rho \rho (a)`$ is an $`\mathrm{End}_\mathrm{𝕜}E`$-valued regular algebraic function on $`\mathrm{Rep}_E^A`$, to be denoted $`\widehat{a}`$. Equivalently, the function $`\widehat{a}`$ may be viewed as an element of $`\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]`$, a tensor product of the finite dimensional simple algebra $`\mathrm{End}_\mathrm{𝕜}E`$ with $`\mathrm{𝕜}[\mathrm{Rep}_E^A]`$, the coordinate ring of the scheme $`\mathrm{Rep}_E^A.`$
Let $`GL(E)`$ be the group of invertible linear transformations of $`E`$. Then we have an action of $`GL(E)`$ on $`\mathrm{Rep}_E^A`$ by conjugation. That is, if $`\phi :A\mathrm{End}_\mathrm{𝕜}E`$ is a representation and $`gGL(E)`$, we define $`(g\phi )(a):=g\phi (a)g^1`$ for all $`aA`$.
The trivial bundle $`E\times \mathrm{Rep}_E^A\mathrm{Rep}_E^A`$ has a natural structure of $`GL(E)`$-equivariant vector bundle on $`\mathrm{Rep}_E^A`$ (with respect to the diagonal $`GL(E)`$-action on $`E\times \mathrm{Rep}_E^A`$). We call this vector bundle the tautological vector bundle, to be denoted $`E_{\mathrm{Rep}}.`$ The algebra $`\mathrm{End}E_{\mathrm{Rep}}`$, of vector bundle endomorphisms of $`E_{\mathrm{Rep}}`$ is clearly identified with $`\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]`$.
Observe that, for any $`aA`$, the element $`\widehat{a}\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]`$ is $`GL(E)`$-invariant with respect to the simultaneous $`GL(E)`$-action on $`\mathrm{End}_\mathrm{𝕜}E`$ (by conjugation) and on $`\mathrm{𝕜}[\mathrm{Rep}_E^A]`$. This way, the assignment $`a\widehat{a}`$ gives a canonical algebra map
$$\mathrm{𝗋𝖾𝗉}:A\left(\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]\right)^{GL(E)}=\left(\mathrm{End}E_{\mathrm{Rep}}\right)^{GL(E)}.$$
(12.1.2)
To make the scheme structure on $`\mathrm{Rep}_E^A`$ explicit, we first consider the special case $`A=\mathrm{𝕜}x_1,\mathrm{},x_r`$, a free algebra on $`r`$ generators.
An algebra homomorphism $`\rho :\mathrm{𝕜}x_1,\mathrm{},x_r\mathrm{End}_\mathrm{𝕜}(E)`$ is specified by an arbitrary choice of an $`r`$-tuple of endomorphisms $`X_1:=\rho (x_1),\mathrm{},X_n:=\rho (x_r)\mathrm{End}_\mathrm{𝕜}(E)`$. Thus, we have an isomorphism of algebraic varieties
$$\mathrm{Rep}_E^{\mathrm{𝕜}x_1,\mathrm{},x_r}\stackrel{_{}}{}\underset{r\text{ factors}}{\underset{}{\mathrm{End}_\mathrm{𝕜}(E)\times \mathrm{}\mathrm{End}_\mathrm{𝕜}(E)}},\rho (\rho (x_1),\mathrm{},\rho (x_r)).$$
(12.1.3)
Now, given any finitely generated associative algebra $`A`$, choose a finite set $`\{x_1,\mathrm{},x_r\}`$ of algebra generators for $`A`$. Then, we have
$$A=\mathrm{𝕜}x_1,\mathrm{},x_n/I$$
for some two-sided ideal $`I\mathrm{𝕜}x_1,\mathrm{},x_n`$. It is clear that, giving an algebra map $`A\mathrm{End}_\mathrm{𝕜}E`$ is the same thing as giving an algebra map $`\mathrm{𝕜}x_1,\mathrm{},x_n\mathrm{End}_\mathrm{𝕜}E`$ that vanishes on the ideal $`I`$. Put another way, the projection $`\mathrm{𝕜}x_1,\mathrm{},x_n{}_{}{}^{}_{}^{}A`$ induces a closed imbedding $`\mathrm{Rep}_E^A{}_{}{}^{}_{}^{}\mathrm{Rep}_E^{\mathrm{𝕜}x_1,\mathrm{},x_r}`$, and we have
$$\mathrm{Rep}_E^A=\{f\mathrm{Rep}_E^{\mathrm{𝕜}x_1,\mathrm{},x_r}\mathrm{End}_\mathrm{𝕜}(E)\times \mathrm{}\mathrm{End}_\mathrm{𝕜}(E)|f(I)=0\}.$$
(12.1.4)
The RHS of this formula is clearly an algebraic subset of a finite dimensional vector space defined by algebraic equations, that is, an affine subscheme. This puts an affine scheme structure on the LHS of the equation, that will be shown below to be independent of the choice of the generators $`x_1,\mathrm{},x_r`$ of the algebra $`A`$.
In order to make the construction of the scheme structure on $`\mathrm{Rep}_E^A`$ manifestly independent of the choice of generators, it is coventient to use the functor of points.
Recall that for any category $`𝒞`$, an object $`S𝒞`$ gives rise to a functor
$$S():𝒞\mathrm{𝖲𝖾𝗍𝗌},XS(X):=\mathrm{Hom}_𝒞(S,X).$$
Further, the Yoneda lemma says that, for any $`S,S^{}𝒞`$, every isomorphism of functors $`S()\stackrel{_{}}{}S^{}()`$ is necessarily induced by an isomorphism of objects $`S^{}\stackrel{_{}}{}S`$. In particular, the functor $`S()`$ determines the object $`S`$ uniquely, up to a unique isomorphism.
Recall next that the category of affine schemes of finite type over $`\mathrm{𝕜}`$ is equivalent, via the ‘coordinate ring’ functor, to the category of finitely generated commutative $`\mathrm{𝕜}`$-algebras. Thus, any affine scheme $`S`$ is completely determined by the corresponding functor
$`S():`$ $`\text{fin. gen. Commutative Alg.}\mathrm{𝖲𝖾𝗍𝗌},`$
$`BS(B):=\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(\mathrm{𝕜}[S],B)\mathrm{Hom}_{\text{Schemes}}(\mathrm{Spec}B,S).`$
The set $`S(B)`$ is usually referred to as the set of $`B`$-points of $`S`$; for $`B=\mathrm{𝕜}[X]`$, it is just the set of algebraic maps $`XS`$.
Now, fix a finitely generated associative (not necessarily commutative) algebra $`A`$, and an integer $`n1`$. Given a finitely generated commutative algebra $`B`$, write $`\mathrm{Mat}_nB`$ for the associative algebra of $`n\times n`$-matrices with entries in $`B`$. We define a functor on the category of finitely generated commutative algebras as follows:
$`\mathrm{Rep}_n^A():`$ $`\text{fin. gen. Commutative Alg.}\mathrm{𝖲𝖾𝗍𝗌},`$
$`B\mathrm{Rep}_n^A(B):=\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB).`$ (12.1.5)
For $`B=\mathrm{𝕜}[X]`$, one may think of the set $`\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_n(\mathrm{𝕜}[X])`$ as the set of families $`\{\rho _x:A\mathrm{Mat}_n\mathrm{𝕜}\}_{xX}`$ of $`n`$-dimensional representations of the algebra $`A`$ parametrized by points of the scheme $`X`$. Thus, if $`E=\mathrm{𝕜}^n`$ and $`X`$ is a point, we get back to the original definition of $`\mathrm{Rep}_n^A:=\mathrm{Rep}_E^A.`$ Observe also that if $`E=\mathrm{𝕜}^n`$, then we have
$$\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_n^A]=\mathrm{Mat}_n\mathrm{𝕜}\mathrm{𝕜}[\mathrm{Rep}_n^A]=\mathrm{Mat}_n\left(\mathrm{𝕜}[\mathrm{Rep}_n^A]\right).$$
The above discussion can be summed up in the following result that shows at the same time that the description of the scheme $`\mathrm{Rep}_E^A`$ given in (12.1.4) is independent of the presentation $`A=\mathrm{𝕜}x_1,\mathrm{},x_r/I`$. We restrict ourselves to the case $`E=\mathrm{𝕜}^n`$, in wich case we have
###### Proposition 12.1.6.
$`(𝗂)`$For any finitely generated associative algebra $`A`$, the corresponding functor (12.1) is representable by an affine scheme, which is called $`\mathrm{Rep}_n^A`$.
$`(\mathrm{𝗂𝗂})`$The coordinate ring $`R:=\mathrm{𝕜}[\mathrm{Rep}_n^A]`$ is a finitely generated commutative algebra equipped with a canonical algebra map $`\mathrm{𝗋𝖾𝗉}:A\mathrm{Mat}_nR,a\widehat{a},`$ see (12.1.2), such that the following universal property holds:
$``$ Given a finitely generated commutative algebra $`B`$ and an algebra map $`\rho :A\mathrm{Mat}_nB`$, there exists a unique algebra homomorphism $`\widehat{\rho }:RB`$ making the following diagram commute
Here and below, given a linear map $`f:VU`$ of vector spaces, we write $`\mathrm{Mat}_n(f):\mathrm{Mat}_n(V)\mathrm{Mat}_n(U)`$ for the map between the spaces of $`V`$-valued and $`U`$-valued matrices whose entries are all equal to $`f`$ (observe also that $`\mathrm{Mat}_n(V)\mathrm{Mat}_n\mathrm{𝕜}V`$).
###### Proof of Proposition.
To prove the representability of the functor one has to choose a set of algebra generators for $`A`$, and write $`A=\mathrm{𝕜}x_1,\mathrm{},x_r/I`$. Then, the set $`\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB)`$ may be identified, as in (12.1.4), with $`B`$-points of a subscheme in $`\mathrm{Mat}_n\mathrm{𝕜}\times \mathrm{}\times \mathrm{Mat}_n\mathrm{𝕜}`$. ∎
###### Remark 12.1.7.
Note that we have defined the scheme $`\mathrm{Rep}_n^A`$ through its functor of points, without describing the coordinate ring $`\mathrm{𝕜}[\mathrm{Rep}_n^A]`$. The latter ring does not have a simple description: below, we will produce a finite set of generators of the ring $`\mathrm{𝕜}[\mathrm{Rep}_n^A]`$, but determining all the relations among these generators is a formidable task. $`\mathrm{}`$
###### Corollary 12.1.8 (Functoriality).
Any algebra map $`f:AA^{},f(1_A)=1_A^{}`$ induces an algebra map $`\widehat{f}:\mathrm{𝕜}[\mathrm{Rep}_n^A]\mathrm{𝕜}[\mathrm{Rep}_n^A^{}],`$ hence, a morphism of algebraic varieties $`\widehat{f}^{}:\mathrm{Rep}_n^A^{}\mathrm{Rep}_n^A.`$
###### Proof.
Set $`B:=\mathrm{𝕜}[\mathrm{Rep}_n^A^{}]`$, and apply Proposition 12.1.6 to to the map
$$\rho :A\stackrel{f}{}A^{}\stackrel{\mathrm{𝗋𝖾𝗉}}{}\mathrm{Mat}_n\left(\mathrm{𝕜}[\mathrm{Rep}_n^A^{}]\right).$$
The universal property from the Proposition implies the existence of an algebra map $`\widehat{\rho }:\mathrm{𝕜}[\mathrm{Rep}_n^A]\mathrm{𝕜}[\mathrm{Rep}_n^A^{}]`$ that makes the following diagram commute:
Thus, we may put $`\widehat{f}:=\widehat{\rho }`$. ∎
### 12.2. Traces.
Let $`A`$ be a $`\mathrm{𝕜}`$-algebra and $`M`$ an $`A`$-module which is finite dimensional over $`\mathrm{𝕜}`$. Then, the action in $`M`$ of an element $`aA`$ gives a $`\mathrm{𝕜}`$-linear map $`a:MM`$. We write $`\mathrm{tr}_M(a)`$ for the trace of this map.
It is clear that if $`N`$ is an $`A`$-submodule in $`M`$, then we have
$$\mathrm{tr}_M(a)=\mathrm{tr}_{M/N}(a)+\mathrm{tr}_N(a)\text{(additivity of the trace)}.$$
Any finite-dimensional $`A`$-module clearly has finite length, hence has a finite Jordan-Hölder series $`M=M^0M^1\mathrm{}M^kM^{k+1}=0`$. The isomorphism classes of simple composition factors $`M^i/M^{i+1}`$ in this series are defined uniquely, up to permutation. Hence, the $`A`$-module $`\mathrm{𝗌𝗌}M:=M^i/M^{i+1}`$ is independent of the choice of Jordan-Hölder series. By construction, $`\mathrm{𝗌𝗌}M`$ is a semisimple $`A`$-module (i.e., a direct sum of simple $`A`$-modules), called the semi-simplification of $`M`$.
The additivity property of the trace implies that $`\mathrm{tr}_M(a)=\mathrm{tr}_N(a),aA`$, whenever $`M`$ and $`N`$ have the same semi-simplification. Conversely, one has
###### Theorem 12.2.1.
Let $`A`$ be a finitely generated $`\mathrm{𝕜}`$-algebra and $`M`$ and $`N`$ be finite-dimensional $`A`$-modules such that $`\mathrm{tr}_M(a)=\mathrm{tr}_N(a)`$ for any $`aA`$. Then, $`\mathrm{𝗌𝗌}M\mathrm{𝗌𝗌}N`$.
###### Proof.
Both the assumptions and the conclusion of the Theorem are unaffected by replacing $`M`$ by $`\mathrm{𝗌𝗌}M`$, and $`N`$ by $`\mathrm{𝗌𝗌}N`$. Therefore, we may reformulate the Theorem as follows: if $`M`$ and $`N`$ are semisimple $`A`$-modules such that $`\mathrm{tr}_M(a)=\mathrm{tr}_N(a)`$ for any $`aA`$, then $`MN`$.
Thus, from now on we assume that $`M`$ and $`N`$ are finite-dimensional semisimple $`A`$-modules.
Let $`\mathrm{𝙰𝚗𝚗}(MN)A`$ be the annihilator of $`MN`$, that is the set of all elements $`aA`$ that act by zero on $`MN`$. This is a two-sided ideal in $`A`$, and the action of $`A`$ gives an algebra imbedding $`A^{}:=A/\mathrm{𝙰𝚗𝚗}(MN){}_{}{}^{}_{}^{}\mathrm{End}_\mathrm{𝕜}(MN).`$ We see that $`A^{}`$ is a finite dimensional algebra and that $`MN`$ is an $`A^{}`$-module, which is moreover semisimple, since it is semisimple as an $`A`$-module.
Let $`\mathrm{Rad}A^{}`$ be the radical of $`A^{}`$, the intersection of the annihilators of all simple $`A^{}`$-modules. Thus, $`\mathrm{Rad}A^{}`$ annihilates any semisimple $`A^{}`$-module, in particular, annihilates $`MN`$. Furthermore, the structure theory of finite dimensional algebras over $`\mathrm{𝕜}`$ says that
$$\overline{A^{}}:=A^{}/\mathrm{Rad}A^{}_i\mathrm{End}_\mathrm{𝕜}(E_i)\text{where}E_i=\mathrm{𝕜}^{r_i}.$$
Thus, each $`\mathrm{End}_\mathrm{𝕜}(E_i)\mathrm{Mat}_{r_i}(\mathrm{𝕜}),`$ is a simple matrix algebra, and any simple $`\overline{A^{}}`$-module is isomorphic to some $`E_i`$, viewed as an $`\overline{A^{}}`$-module via the projection $`\overline{A^{}}{}_{}{}^{}_{}^{}\mathrm{End}_\mathrm{𝕜}(E_i).`$
We conclude that our semisimple $`\overline{A^{}}`$-modules $`M`$ and $`N`$ have the form $`M=E_i`$ and $`N=E_j`$, respectively. With this understood, the assumption of the Theorem reads:
$$\mathrm{tr}_M(a)=\mathrm{tr}_N(a)\text{for any}a_i\mathrm{End}_\mathrm{𝕜}(E_i).$$
This clearly implies that the direct sums $`M=E_i`$ and $`N=E_j`$ involve the same summands with the same multiplicities. Hence $`MN`$, and we are done. ∎
We now fix a finite dimensional vector space $`E`$ and consider the scheme $`\mathrm{Rep}_E^A`$. For any $`\rho \mathrm{Rep}_E^A`$, let $`𝕆(\rho )`$ denote the $`GL(E)`$-orbit of the point $`\rho `$. The orbit $`𝕆(\rho )`$ corresponds to the isomorphism class of the $`A`$-module $`E_\rho `$, and we write $`𝕆(\mathrm{𝗌𝗌}(\rho ))`$ for the orbit corresponding to the semi-simplification of $`E_\rho `$.
The following result is a standard application of Geometric Invariant Theory, see \[GIT\].
###### Theorem 12.2.2.
$`(𝗂)`$The orbit $`𝕆(\rho )`$ is closed in $`\mathrm{Rep}_E^A`$ if and only if $`E_\rho `$ is a semisimple $`A`$-module.
$`(\mathrm{𝗂𝗂})`$For each $`\rho \mathrm{Rep}_E^A`$, the orbit $`𝕆(\mathrm{𝗌𝗌}(\rho ))`$ is the unique closed $`GL(E)`$-orbit contained in $`\overline{𝕆(\rho )}`$, the closure of $`𝕆(\rho )`$.
Taking the trace of the $`\mathrm{End}_\mathrm{𝕜}E`$-valued function $`\widehat{a}`$, one obtains an element $`\mathrm{tr}(\widehat{a})\mathrm{𝕜}[\mathrm{Rep}_E^A]^{GL(E)}`$, the $`GL(E)`$-invariant regular functions on $`\mathrm{Rep}_E^A`$. Equivalently, $`\mathrm{tr}(\widehat{a})`$ is the image of $`\widehat{a}`$ under the linear map $`\mathrm{tr}\mathrm{id}:\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]\mathrm{𝕜}\mathrm{𝕜}[\mathrm{Rep}_E^A]=\mathrm{𝕜}[\mathrm{Rep}_E^A].`$
###### Theorem 12.2.3.
$`(𝗂)`$The functions $`\{\mathrm{tr}(\widehat{a}),aA\}`$ generate $`\mathrm{𝕜}[\mathrm{Rep}_E^A]^{GL(E)}`$ as an algebra.
$`(\mathrm{𝗂𝗂})`$The natural map $`\mathrm{Rep}_E^A\mathrm{Spec}\mathrm{𝕜}[\mathrm{Rep}_E^A]^{GL(E)}`$ induces a bijection between the isomorphism classes of semisimple $`A`$-modules of the form $`E_\rho ,\rho \mathrm{Rep}_E^A`$, and maximal ideals of the algebra $`\mathrm{𝕜}[\mathrm{Rep}_E^A]^{GL(E)}`$.
###### Proof.
Part (i) can be deduced from a theorem of LeBruyn-Procesi \[LP\]. Part (ii) follows from the general result, cf. \[GIT\] saying that the elements of the algebra $`\mathrm{𝕜}[\mathrm{Rep}_E^A]^{GL(E)}`$ separate closed $`GL(E)`$-orbits in $`\mathrm{Rep}_E^A`$. ∎
### 12.3. Noncommutative Rep-scheme (following \[LBW\]).
We have seen that the functor $`B\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB)`$ is representable, as a functor on the category of finitely generated commutative algebras. L. Le Bruyn and G. Van de Weyer have observed, see \[LBW\], that this functor is in effect representable on the much larger category of all finitely generated non-commutative associative algebras. The non-commutative algebra, $`\sqrt[n]{A},`$ that represents the functor $`\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_n())`$ maps surjectively onto $`\mathrm{𝕜}[\mathrm{Rep}_n^A]`$, the coordinate ring of the scheme $`\mathrm{Rep}_n^A`$ constructed in Proposition 12.1.6. Thus, $`\sqrt[n]{A}`$ may be thought of as the coordinate ring of a non-commutative space, the canonical ‘non-commutative thickening’ of the commutativescheme $`\mathrm{Rep}_n^A`$.
The algebra $`\sqrt[n]{A}`$ is constructed as follows, see \[LBW\] for more details.
First, one forms the free-product algebra $`A\mathrm{Mat}_n\mathrm{𝕜}`$, which contains the matrix algebra $`\mathrm{Mat}_n\mathrm{𝕜}`$ as a subalgebra. Let
$`\sqrt[n]{A}`$ $`:=\left[A\mathrm{Mat}_n\mathrm{𝕜}\right]^{\mathrm{Mat}_n\mathrm{𝕜}}`$ (12.3.1)
$`=\{rA\mathrm{Mat}_n\mathrm{𝕜}|xr=rx,x\mathrm{Mat}_n\mathrm{𝕜}\}.`$
be the centraliser of the subalgebra $`\mathrm{Mat}_n\mathrm{𝕜}A\mathrm{Mat}_n\mathrm{𝕜}.`$
###### Example 12.3.2 (\[LBW\]).
One can show that $`\sqrt[n]{\mathrm{𝕜}x_1,\mathrm{},x_r}\mathrm{𝕜}x_{_{11,1}},\mathrm{},`$ $`x_{_{rr,n}},`$ is a free algebra on $`nr^2`$ generators.∎ $`\mathrm{}`$
Following \[LBW\], we are going to construct, for any (not necessarily commutative), associative algebra $`B`$, a canonical bijection
$$\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB)\stackrel{}{}\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(\sqrt[n]{A},B)$$
(12.3.3)
The bijection would clearly yield the representability claim mentioned above.
The construction of this bijection is based on the elementary Lemma below. To formulate the Lemma, fix an associative algebra $`R`$, and an algebra map $`ı:\mathrm{Mat}_n\mathrm{𝕜}R`$, such that $`ı(1)=1`$. Let $`R^{\mathrm{Mat}_n}`$ denote the centralizer of $`ı\left(\mathrm{Mat}_n\mathrm{𝕜}\right)`$ in $`R`$.
###### Lemma 12.3.4.
In the above setting, the map $`\mathrm{Mat}_n\mathrm{𝕜}R^{\mathrm{Mat}_n}R,mrı(m)r,`$ is an algebra isomorphism. Moreover, if $`R`$ is a finitely generated algebra, then so is the algebra $`R^{\mathrm{Mat}_n}`$. ∎
###### Proof (by D. Boyarchenko)..
The map is clearly an algebra homomorphism. To prove that this map is bijective, we form the algebra $`S:=\mathrm{Mat}_n\mathrm{𝕜}__\mathrm{𝕜}\mathrm{Mat}_n\mathrm{𝕜}^{\mathrm{op}}`$. This is clearly a simple algebra of dimension $`n^4`$, which is moreover isomorphic to $`\mathrm{Mat}_{n^2}(\mathrm{𝕜})`$. Thus, by the well-known results about modules over finite dimensional simple $`\mathrm{𝕜}`$-algebras, the algebra $`S`$ has a unique, up to isomorphism, simple $`S`$-module $`L`$ which has dimension $`n^2`$. Furthermore, any faithful $`S`$-module $`M`$ (i.e., such that $`1_Sm=0m=0`$, for any $`mM`$) is isomorphic to a (possibly infinite) direct sum of copies of $`L`$, that is, has the form $`LV`$, for some vector space $`V`$.
The algebra $`\mathrm{Mat}_n\mathrm{𝕜}`$ has an obvious $`\mathrm{Mat}_n\mathrm{𝕜}__\mathrm{𝕜}\mathrm{Mat}_n\mathrm{𝕜}^{\mathrm{op}}`$-module structure via left and right multiplication. Also, $`dim\mathrm{Mat}_n\mathrm{𝕜}=n^2`$. We conclude that $`L=\mathrm{Mat}_n\mathrm{𝕜}`$, as a module over $`S=\mathrm{Mat}_n\mathrm{𝕜}__\mathrm{𝕜}\mathrm{Mat}_n\mathrm{𝕜}^{\mathrm{op}}`$.
Now, the map $`\iota :\mathrm{Mat}_n\mathrm{𝕜}R`$ makes $`R`$ a faithful (since $`\iota (1_{\mathrm{Mat}_n})=1_R`$) $`\mathrm{Mat}_n\mathrm{𝕜}`$-bimodule, hence, a left $`S`$-module. Thus, there is an $`S`$-module isomorphism $`RSV`$, where $`V`$ is some $`\mathrm{𝕜}`$-vector space. Under this isomorphism, $`R^{\mathrm{Mat}_n}`$ clearly corresponds to $`(\mathrm{𝕜}1_S)V`$ because the center of $`S`$ is $`\mathrm{𝕜}1_S`$. So the statement of the lemma becomes obvious. ∎
Applying the Lemma to the tautological imbedding $`ı:\mathrm{Mat}_n\mathrm{𝕜}R=A\mathrm{Mat}_n\mathrm{𝕜}`$, and using the notation (12.3.1), we obtain a canonical algebra isomorphism
$$\phi :\mathrm{Mat}_n\mathrm{𝕜}\sqrt[n]{A}=\mathrm{Mat}_n\mathrm{𝕜}R^{\mathrm{Mat}_n}\stackrel{_{}}{}A\mathrm{Mat}_n\mathrm{𝕜}.$$
Now, for any algebra $`B`$, an algebra morphism $`f:\sqrt[n]{A}B`$, induces an algebra morphism $`\stackrel{~}{f}:A\mathrm{Mat}_nB`$ given by the composition:
$$A{}_{}{}^{}_{}^{}A\mathrm{Mat}_n\mathrm{𝕜}\stackrel{\phi ^1}{\stackrel{_{}}{}}\mathrm{Mat}_n\mathrm{𝕜}\sqrt[n]{A}\stackrel{\mathrm{id}f}{}\mathrm{Mat}_n\mathrm{𝕜}B=\mathrm{Mat}_nB.$$
Conversely, given an algebra map $`g:A\mathrm{Mat}_nB`$, the universal property of free products yields an algebra map $`gȷ:A\mathrm{Mat}_n\mathrm{𝕜}\mathrm{Mat}_nB`$, where $`ȷ:\mathrm{Mat}_n\mathrm{𝕜}{}_{}{}^{}_{}^{}\mathrm{Mat}_nB`$ is the natural imbedding. Observe that the centralizer of the image of $`ȷ`$ is formed by “scalar matrices”, i.e., we have $`\left(\mathrm{Mat}_nB\right)^{\mathrm{Mat}_n\mathrm{𝕜}}=B`$. Therefore, the subalgebra $`gȷ(\sqrt[n]{A})`$, the image of the restriction of $`gȷ:A\mathrm{Mat}_n\mathrm{𝕜}\mathrm{Mat}_nB`$ to the centralizer of $`\mathrm{Mat}_n\mathrm{𝕜}`$ in $`A\mathrm{Mat}_n\mathrm{𝕜}`$, is contained in the subalgebra $`\left(\mathrm{Mat}_nB\right)^{\mathrm{Mat}_n\mathrm{𝕜}}=B`$. This way, one obtains a map $`\sqrt[n]{g}:\sqrt[n]{A}B`$.
It is straightforward to check that the assignments $`f\stackrel{~}{f}`$ and $`g\sqrt[n]{g},`$ are mutually inverse bijections. This completes the construction of the bijection in (12.3.3).
### 12.4. The Rep-functor on vector fields.
Below, we are going to relate various ‘non-commutative’ constructions on $`A`$ with their commutative counterparts for $`\mathrm{Rep}_E^A`$.
First, we claim that any derivation of $`A`$ gives rise to a $`GL(E)`$-invariant vector field on the scheme $`\mathrm{Rep}_E^A`$.
###### Example 12.4.1.
For instance, let $`\theta =\mathrm{ad}a\mathrm{Inn}(A)\mathrm{Der}(A)`$ be an inner derivation. The corresponding vector field $`\widehat{\theta }`$ on $`\mathrm{Rep}_E^A`$ is then going to be tangent to the orbits of the $`GL(E)`$-action on $`\mathrm{Rep}_E^A`$, and it is constructed as follows. Write $`𝔤=\mathrm{Lie}^{}GL(E)=\mathrm{End}_\mathrm{𝕜}E`$ for the Lie algebra of the group $`GL(E)`$. To $`aA`$, we have associated the element $`\widehat{a}\mathrm{End}_\mathrm{𝕜}E\mathrm{𝕜}[\mathrm{Rep}_E^A]`$, which we now view as a $`𝔤`$-valued regular function on $`\mathrm{Rep}_E^A`$. For point $`\rho \mathrm{Rep}_E^A`$, let $`\widehat{a}(\rho )𝔤`$ be the value of the function $`\widehat{a}`$ at $`\rho `$. Further, the infinitesimal $`𝔤`$-action on $`\mathrm{Rep}_E^A`$ (that is, the differential of the action map $`GL(E)\times \mathrm{Rep}_E^A\mathrm{Rep}_E^A`$ at the unit element of $`GL(E)`$) associates to any $`x𝔤`$ and $`\rho \mathrm{Rep}_E^A`$ a vector $`x_\rho `$ in $`T_\rho (\mathrm{Rep}_E^A),`$ the tangent space to $`\mathrm{Rep}_E^A`$ at $`\rho `$. In particular, we have the vector $`\widehat{a}(\rho )_\rho T_\rho (\mathrm{Rep}_E^A).`$ The assignment $`\rho \widehat{a}(\rho )_\rho `$ gives the required vector field, $`\widehat{\theta }`$, on $`\mathrm{Rep}_E^A`$. $`\mathrm{}`$
To study the general case, we first need to understand the tangent space of $`\mathrm{Rep}_E^A`$ at some point $`\rho \mathrm{Rep}_E^A`$. So, $`\rho :A\mathrm{End}_\mathrm{𝕜}E`$ is an algebra map, hence it induces an $`A`$-module structure on $`E`$. We denote this $`A`$-module by $`E_\rho `$. Then the tangent space $`T_\rho \mathrm{Rep}_E^A`$ is given by all $`\mathrm{𝕜}`$-linear maps $`\phi :A\mathrm{End}_\mathrm{𝕜}(E)`$ such that $`(\rho +\epsilon \phi ):A\mathrm{End}_{\mathrm{𝕜}[\epsilon ]}\left(\mathrm{𝕜}[\epsilon ]/(\epsilon ^2)E\right)`$ is an algebra map. Expanding both sides of $`(\rho +\epsilon \phi )(aa^{})=(\rho +\epsilon \phi )(a)(\rho +\epsilon \phi )(a^{})`$, we have
$$\rho (aa^{})+\epsilon \phi (aa^{})=\rho (a)\rho (a^{})+\epsilon \phi (a)\phi (a^{})+\rho (a)\epsilon \phi (a^{}).$$
Since $`\rho `$ is an algebra map, $`\rho (aa^{})=\rho (a)\rho (a^{})`$. Cancelling them and dividing by the common $`\epsilon `$ factor, we find that
$$\phi (aa^{})=\phi (a)\rho (a^{})+\rho (a)\phi (a^{}).$$
If we regard $`\mathrm{End}_\mathrm{𝕜}(E_\rho )`$ as an $`A`$-bimodule in the obvious fashion, this equation implies that $`\phi \mathrm{Der}(A,\mathrm{End}_\mathrm{𝕜}(E_\rho ))`$, i.e., $`T_\rho \mathrm{Rep}_E^A\mathrm{Der}(A,\mathrm{End}_\mathrm{𝕜}(E_\rho ))`$.
Now, suppose we have a derivation $`\theta \mathrm{Der}(A)`$. We wish to generate a vector field $`\widehat{\theta }`$ on $`\mathrm{Rep}_E^A`$. Indeed, for each $`\rho \mathrm{Rep}_E^A`$, we set
$$\widehat{\theta }_\rho (a)=\rho (\theta (a))=\widehat{\theta (a)}(\rho ).$$
It is not hard to check that this is a derivation, hence $`\widehat{\theta }_\rho \mathrm{Der}(A,\mathrm{End}_\mathrm{𝕜}(E_\rho ))`$ $`=T_\rho \mathrm{Rep}_E^A`$.
###### Example 12.4.2.
By Lemma 3.4.1, a derivation $`\theta `$ of the free algebra $`A=\mathrm{𝕜}x_1,\mathrm{},x_r`$ is determined by sending each generator $`x_i,i=1,\mathrm{},r,`$ to an arbitrarily chosen element $`\theta (x_i)=f_i(x_1,\mathrm{},x_r)A,`$ which we regard as a non-commutative polynomial in the variables $`x_1,\mathrm{},x_r`$. Given such a derivation $`\theta `$ we describe the corresponding vector field $`\widehat{\theta }`$ on $`\mathrm{Rep}_E^A`$ as follows.
Fix $`\rho \mathrm{Rep}_E^A`$ and let $`X_1:=\rho (x_1),\mathrm{},X_n:=\rho (x_r)\mathrm{End}_\mathrm{𝕜}(E)`$ be as above. We have $`T_\rho \mathrm{Rep}_E^A=\mathrm{Der}(A,\mathrm{End}E_\rho )`$, where $`E_\rho `$ denotes the $`A`$-module $`E`$ endowed with the $`A`$-module structure given by $`\rho `$. Then composing $`\theta `$ with $`\rho `$, we obtain a derivation $`\rho \theta :A\mathrm{End}E_\rho `$. Thus, writing $`\widehat{\theta }|_\rho `$ for the value of the vector field $`\widehat{\theta }`$ at the point $`\rho `$, we have $`\widehat{\theta }|_\rho =\rho \theta `$. Hence, for each $`i=1,\mathrm{},r,`$ we find $`\widehat{\theta }|_\rho :x_if_i(X_1,\mathrm{},X_r)`$. Thus, using the identification (12.1.3) we get
$$\widehat{\theta }|_{_{(X_1,\mathrm{},X_r)}}=(f_1(X_1,\mathrm{},X_r),\mathrm{},f_r(X_1,\mathrm{},X_r)),$$
where the $`r`$-tuple on the RHS is viewed as a vector in $`\mathrm{End}_\mathrm{𝕜}(E)\times \mathrm{}\mathrm{End}_\mathrm{𝕜}(E)`$. Now, a vector field $`\xi `$ on any vector space $`R`$ is nothing but a map $`RR`$ that sends each vector $`rR`$ to the vector $`\xi _r`$, the value of $`\xi `$ at $`r`$. With this understood, we see that the vector field $`\widehat{\theta }`$ is nothing but the following self-map of $`\mathrm{End}_\mathrm{𝕜}(E)\times \mathrm{}\mathrm{End}_\mathrm{𝕜}(E)`$
$$\widehat{f}:(X_1,\mathrm{},X_r)(f_1(X_1,\mathrm{},X_r),\mathrm{},f_r(X_1,\mathrm{},X_r)).$$
(12.4.3)
$`\mathrm{}`$
###### Proposition 12.4.4.
For any finitely generated associative algebra $`A`$, the assignment $`\theta \widehat{\theta }`$ gives a Lie algebra homomorphism $`\mathrm{Der}(A)𝒯(\mathrm{Rep}_E^A).`$
Proof (by D. Boyarchenko). Let $`\theta \mathrm{Der}(A),`$ be a derivation. Set $`R:=\mathrm{𝕜}[\mathrm{Rep}_n^A]`$, and recall the notation of Proposition 12.1.6. We have the following characterization of the vector field $`\widehat{\theta }`$:
###### Claim 12.4.5.
$`(𝗂)`$There exists a unique derivation $`\stackrel{~}{\theta }\mathrm{Der}(R)`$ that makes the following square commute
$`(\mathrm{𝗂𝗂})`$For any point $`\rho \mathrm{Rep}_n^A`$, the value at $`\rho `$ of the vector field on $`\mathrm{Rep}_n^A`$ induced by the derivation $`\stackrel{~}{\theta }`$ from $`(𝗂)`$, is equal to the vector $`\widehat{\theta }|_\rho `$ constructed above. Thus, $`\widehat{\theta }=\stackrel{~}{\theta }`$ is indeed a regular algebraic vector field.
Assuming the Claim, we complete the proof of the Proposition as follows. Let $`\theta ,\delta \mathrm{Der}(A)`$, and $`\stackrel{~}{\theta },\stackrel{~}{\delta }\mathrm{Der}(R)`$, be the corresponding derivations of $`R`$ arising via Claim 12.4.5. Then, clearly $`[\stackrel{~}{\theta },\stackrel{~}{\delta }]`$ is again a derivation of $`R`$ and, moreover, the diagrams of part (i) of the Claim for $`\theta `$ and $`\delta `$, respectively, yield commutativity of the diagram
Now, the uniqueness statement in Claim 12.4.5(i) combined with (ii) implies that $`\widehat{[\theta ,\delta ]}=[\widehat{\theta },\widehat{\delta }]`$, and the Proposition is proved.
###### Proof of Claim..
Observe that for $`A`$ viewed as an $`A`$-bimodule, the square-zero construction produces an algebra $`A\mathrm{}A`$ which is nothing but $`A_\epsilon :=A\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$. Further, let $`A{}_{}{}^{}_{}^{}A_\epsilon `$ be the tautological algebra imbedding. A linear map $`\theta :AA`$ is a derivation of $`A`$ if and only if the assignment $`\theta _\epsilon :aa+\epsilon \theta (a),`$ gives an algebra homomorphism $`AA_\epsilon `$.
Recall the canonical algebra map $`\mathrm{𝗋𝖾𝗉}:A\mathrm{Mat}_nR`$. Tensoring with $`\mathrm{𝕜}[\epsilon ]/\epsilon ^2`$, one gets a homomorphism $`\mathrm{𝗋𝖾𝗉}_\epsilon :A_\epsilon \left(\mathrm{Mat}_nR\right)_\epsilon .`$ But we have
$$\left(\mathrm{Mat}_nR\right)_\epsilon =\mathrm{Mat}_nR\mathrm{𝕜}[\epsilon ]/\epsilon ^2=\mathrm{Mat}_n\left(R\mathrm{𝕜}[\epsilon ]/\epsilon ^2\right)=\mathrm{Mat}_n(R_\epsilon ).$$
Now, fix $`\theta \mathrm{Der}(A)`$ and let $`\theta _\epsilon :AA_\epsilon `$ be the corresponding homomorphism. Composing $`\theta _\epsilon `$ with the morphism $`\mathrm{𝗋𝖾𝗉}_\epsilon `$, we obtain the top row of the following diagram
The universal property of $`R`$ explained in Proposition 12.1.6 guarantees the existence and uniqueness of an algebra map $`\varphi :RR_\epsilon `$ that makes the above diagram commute. By the discussion at the beginning of the proof, the map $`\varphi `$ thus defined gives a derivation $`\stackrel{~}{\theta }:RR`$ such that $`\varphi =(\stackrel{~}{\theta })_\epsilon .`$ This proves the Claim. ∎
Let $`𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}})`$ denote the space of first order differential operators on $`E_{\mathrm{Rep}}`$ with scalar principal symbol, cf. §3.6. This is a Lie algebra, and the group $`GL(E)`$ acts naturally on $`𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}})`$ by Lie algebra automorphisms. We consider the Lie subalgebra $`𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}})^{GL(E)}𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}}),`$ of $`GL(E)`$-invariant elements.
###### Remark 12.4.6.
Since the vector bundle $`E_{\mathrm{Rep}}`$ is canonically and $`GL(E)`$-equivariantly trivialized, the symbol map $`\sigma :𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}})𝒯(\mathrm{Rep}_E^A)`$ has a natural $`GL(E)`$-equivariant splitting $`𝒯(\mathrm{Rep}_E^A)𝒟_1^{\mathrm{}}(E_{\mathrm{Rep}}).`$ This splitting allows to lift the the Lie algebra map $`\theta \widehat{\theta }`$ of Proposition 12.4.4 to a Lie algebra map indicated by the dotted arrow in the following diagram
### 12.5. Rep-functor and the de Rham complex.
Given a noncommutative differential $`n`$-form, $`\omega =`$ $`a_0da_1\mathrm{}da_n\mathrm{\Omega }_{\mathrm{nc}}^n(A)`$, we define
$$\rho (\omega )=\mathrm{tr}(\widehat{a}_0d\widehat{a}_1\mathrm{}d\widehat{a}_n)\mathrm{\Omega }^n(\mathrm{Rep}_E^A)^{GL(E)},$$
where the $`\mathrm{\Omega }^n`$ on the right-hand side denotes the usual differential $`n`$-forms. Since the trace is symmetric (that is, $`\mathrm{tr}(ab)=\mathrm{tr}(ba)`$), it vanishes on $`[\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)]`$. Hence the $`\rho `$ descends to a map $`\rho :\mathrm{DR}^{\text{}}(A)\mathrm{\Omega }^{\text{}}(\mathrm{Rep}_E^A)^{GL(E)}`$.
We would like to deal with equivariant cohomology rather than this invariant part of the (usual) de Rham complex.
### 12.6. Equivariant Cohomology
Let $`X`$ be an algebraic variety, and let $`𝔤`$ be a Lie algebra. Suppose we have a Lie algebra map $`𝔤𝒯(X)`$. The $`𝔤`$-equivariant algebraic de Rham complext of $`X`$ is the complex
$$((\mathrm{\Omega }^{\text{}}(X)\mathrm{𝕜}[𝔤])^𝔤,d_𝔤),$$
(12.6.1)
where $`𝔤`$ acts on $`\mathrm{\Omega }^{\text{}}(X)`$, by the Lie derivative $``$, the $`𝔤`$-action on $`\mathrm{𝕜}[𝔤]`$ is induced by the adjoint action of $`𝔤`$ on, and the differential (called the Koszul differential) $`d_𝔤`$ is defined in the following way. Choose a basis $`e_r`$ for $`𝔤`$ and let $`e_r^{}`$ denote the dual basis. Then we set
$$d_𝔤=d+\underset{r=1}{\overset{n}{}}e_r^{}i_{e_r},$$
where $`d`$ is the usual de Rham differential on $`\mathrm{\Omega }^{\text{}}(X)`$ and $`i_{e_r}`$ denotes contraction by $`e_r`$.
###### Remark 12.6.2.
The cohomology of the complex (12.6.1) are called the equivariant cohomology of $`X`$. Assume $`G`$ is a Lie group with Lie algebra $`𝔤`$ such that the map $`𝔤𝒯(X)`$ can be exponentiated to a free $`G`$-action on $`X`$ and, moreover, the orbit space $`X/G`$ is a well-defined algebraic variety. Then, the complex (12.6.1) computes the ordinary de Rham cohomology of $`X/G`$. In the general case, the leaves of the vector fields coming from the image of $`𝔤𝒯(X)`$ allow to consider the stack quotient $`X/𝔤`$, and the cohomology of the complex (12.6.1) should be thought of as the de Rham cohomology of that stack quotient. $`\mathrm{}`$
Choose any element $`\omega f(\mathrm{\Omega }^{\text{}}(X)\mathrm{𝕜}[𝔤])^𝔤`$. Then $`\omega f(0)\mathrm{\Omega }^{\text{}}(X)^𝔤`$, which yields a surjection
$$(\mathrm{\Omega }^{\text{}}(X)\mathrm{𝕜}[𝔤])^𝔤\mathrm{\Omega }^{\text{}}(X)^𝔤$$
that intertwines the differentials $`d_𝔤`$ and $`d`$. However, this map is usually neither surjective nor injective on the level of cohomology.
Now, we let $`A`$ be an associative algebra and $`E`$ a finite-dimensional vector space. Let $`X=\mathrm{Rep}_E^A`$ be the representation variety of $`A`$ with the natural $`GL(E)`$-action. Then we obtain a map
$$\mathrm{𝗋𝖾𝗉}_{\mathrm{DR}}:\mathrm{DR}(A)\mathrm{\Omega }^{\text{}}(\mathrm{Rep}_E^A)^𝔤,\alpha \mathrm{tr}(\widehat{\alpha }),$$
where $`𝔤=𝔤𝔩(E)`$, given by $`a_0da_1\mathrm{}da_n\mathrm{tr}(\widehat{a}_0d\widehat{a}_1\mathrm{}d\widehat{a}_n)`$, where for each $`aA`$, $`\widehat{a}:\mathrm{Rep}_E^A\mathrm{End}_\mathrm{𝕜}E`$ is given by $`\widehat{a}(\phi )=\phi (a)`$.
###### Problem 12.6.3.
Is there a natural map making the following diagram commute ?
Here is a simple example of a possible lift (dotted arrow above) in a very special case. Let $`\alpha =da_1da_2\mathrm{DR}^2(A)`$, so $`\mathrm{tr}(\widehat{\alpha })=\mathrm{tr}(d\widehat{a}_1d\widehat{a}_2)`$. Finding the image of $`\alpha `$ under the dotted map amounts to finding a degree two element $`\stackrel{~}{\alpha }\left(\mathrm{\Omega }^{\text{}}(\mathrm{Rep}_E^A)\mathrm{𝕜}[𝔤]\right)^𝔤`$ of the form $`\stackrel{~}{\alpha }=\alpha 1+\mathrm{\Phi }`$, where $`\mathrm{\Phi }\left(\mathrm{𝕜}[\mathrm{Rep}_E^A]𝔤^{}\right)^𝔤`$. Since $`\alpha `$ is closed, the compatibility of the dotted map with the differentials reads
$$i_x\mathrm{tr}(\widehat{\alpha })=d(\mathrm{\Phi },x),x𝔤.$$
Using the trace pairing $`x,y\mathrm{tr}(xy)`$ on $`𝔤=\mathrm{End}_CE`$, we may identify $`𝔤`$ with $`𝔤^{}`$, and view $`\mathrm{\Phi }`$ as an element of $`\mathrm{𝕜}[\mathrm{Rep}_E^A]𝔤`$. With these identifications, we have $`\mathrm{\Phi },x=\mathrm{tr}(\mathrm{\Phi }x).`$ It is easy to check that the equation above is satisfied if one puts $`\mathrm{\Phi }:=\widehat{a}_1\widehat{a}_2`$ (product in the associative algebra $`\mathrm{𝕜}[\mathrm{Rep}_E^A]\mathrm{End}_\mathrm{𝕜}E`$).
## 13. Double-derivations and the double-tangent bundle.
### 13.1.
Given an associative algebra $`A`$, we define the $`A`$-bimodule of double-derivations of $`A`$ by
$$𝔻\mathrm{𝐞𝐫}(A):=\mathrm{Der}(A,A^\mathrm{e})=\mathrm{Hom}_{A^\mathrm{e}}(\mathrm{\Omega }_{\mathrm{nc}}^1(A),A^\mathrm{e}),$$
(13.1.1)
where the $`A`$-bimodule structure comes from the $`A^\mathrm{e}`$-action on the entry $`A^\mathrm{e}`$ by right multiplication.
The fundamental exact sequence (10.1.1) gives rise to an exact sequence
$$0\mathrm{End}_{A^\mathrm{e}}\left(\mathrm{\Omega }_{\mathrm{nc}}^1(A)\right)𝔻\mathrm{𝐞𝐫}(A)\mathrm{Der}(A)\mathrm{Ext}_{A^\mathrm{e}}^1(\mathrm{\Omega }_{\mathrm{nc}}^1(A),\mathrm{\Omega }_{\mathrm{nc}}^1(A)).$$
(13.1.2)
###### Remark 13.1.3.
Note that the space $`\mathrm{End}_{A^\mathrm{e}}\left(\mathrm{\Omega }_{\mathrm{nc}}^1(A)\right)`$ on the left is an associative algebra, in particular, a Lie algebra with respect to the commutator bracket. The space $`\mathrm{Der}(A)`$ is also a Lie algebra. Furthermore, the sequence above resembles the exact sequence for the Atiyah algebra, cf. §3.6, of the ‘vector bundle’ $`\mathrm{\Omega }_{\mathrm{nc}}^1(A).`$ $`\mathrm{}`$
We observe that if $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is a projective $`A^\mathrm{e}`$-module (in which case the algebra $`A`$ is said to be formally smooth, see §22), the $`\mathrm{Ext}^1`$-term on the right of (13.1.2) vanishes, hence the map $`𝔻\mathrm{𝐞𝐫}(A)\mathrm{Der}(A)`$ becomes surjective.
It is immediate to check that the map
$$\mathrm{\Delta }:AA^\mathrm{e}=AA,aa11a$$
(13.1.4)
gives a derivation, i.e., we have $`\mathrm{\Delta }𝔻\mathrm{𝐞𝐫}(A)`$. The derivation $`\mathrm{\Delta }`$ maps to zero under the projection $`𝔻\mathrm{𝐞𝐫}(A)\mathrm{Der}(A)`$, cf. (13.1.2), hence, $`\mathrm{\Delta }`$ may be identified with an element of $`\mathrm{End}_{A^\mathrm{e}}\left(\mathrm{\Omega }_{\mathrm{nc}}^1(A)\right)`$. It is easy to verify that the latter element is nothing but the identity map $`\mathrm{Id}_{\mathrm{\Omega }_{\mathrm{nc}}^1(A)}`$.
###### Example 13.1.5.
Let $`A=\mathrm{𝕜}[X]`$ be the coordinate ring of an affine algebraic variety. Then, $`A^\mathrm{e}=\mathrm{𝕜}[X\times X]`$, and the bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is the ideal of the diagonal $`X_\mathrm{\Delta }X\times X.`$ We see that if $`dimX>1`$, then $`X_\mathrm{\Delta }`$ has codimension $`2`$ in $`X\times X`$, It is easy to deduce, by Hartogs type theorem, that in this case the map $`\mathrm{ad}:A^\mathrm{e}𝔻\mathrm{𝐞𝐫}(A)`$ is an isomorphism.
Assume now that $`dimX=1`$, i.e., $`X`$ is a curve. Then, we have $`𝔻\mathrm{𝐞𝐫}(A)=\mathrm{\Gamma }(X\times X,𝒪_{X\times X}(X_\mathrm{\Delta })),`$ is the space of rational functions on $`X\times X`$ with simple pole along $`X_\mathrm{\Delta }.`$ On the other hand, if $`X`$ is a curve ($`dimX=1`$), then the exact sequence in (13.1.2) reduces to
$$0\mathrm{𝕜}[X\times X]𝔻\mathrm{𝐞𝐫}(A)𝒯(X)0.\mathrm{}$$
(13.1.6)
The “meaning” of the double-derivation bimodule may be seen in terms of the $`\mathrm{Rep}`$-functor, via the following construction due to Kontsevich-Rosenberg \[KR\].
Fix a finite dimensional vector space $`E`$, and let $`\rho :A\mathrm{End}_\mathrm{𝕜}E`$ be a representation of $`A`$ in $`E`$ (to be denoted $`E_\rho `$). The action map $`AE_\rho E_\rho `$ is surjective since $`A`$ contains the unit. This gives a surjective map of $`A`$-modules $`\mathrm{𝚊𝚌𝚝}:AE{}_{}{}^{}_{}^{}E_\rho `$, where $`AE`$ is regarded as a free left $`A`$-module generated by the vector space $`E`$. We set $`K_\rho :=`$ $`\mathrm{Ker}(AE{}_{}{}^{}_{}^{}E_\rho ),`$ a left $`A`$-module.
For any left $`A`$-module $`M`$, we have $`\mathrm{Hom}_A(AE,M)=\mathrm{Hom}_\mathrm{𝕜}(E_\rho ,M)`$ and $`\mathrm{Ext}_A^1(AE,M)=0`$. Hence, the long exact sequence of Ext-groups arising from the short exact sequence $`K_\rho {}_{}{}^{}_{}^{}AE{}_{}{}^{}_{}^{}E_\rho `$ reads:
$`0\mathrm{Hom}_A(E_\rho ,M)\stackrel{\mathrm{𝖾𝗏}}{}\mathrm{Hom}_\mathrm{𝕜}(E_\rho ,M)`$ $`\mathrm{Hom}_A(K_\rho ,M)`$ (13.1.7)
$`\mathrm{Ext}_A^1(E_\rho ,M)0.`$
It is easy to see that if $`M`$ is finite dimensional over $`\mathrm{𝕜}`$, then the space $`\mathrm{Ext}_A^1(E_\rho ,M)`$ is also finite dimensional over $`\mathrm{𝕜}`$. It follows that
$$dim_\mathrm{𝕜}M<\mathrm{}dim_\mathrm{𝕜}\mathrm{Hom}_A(K_\rho ,M)<\mathrm{}.$$
(13.1.8)
### 13.2.
We consider the scheme $`\mathrm{Rep}_E^A`$. To simplify notation, write $`\mathrm{End}:=\mathrm{End}_\mathrm{𝕜}E`$, a simple associative algebra isomorphic to $`\mathrm{Mat}_n\mathrm{𝕜}`$ where $`n=dimE`$. We set $`R:=\mathrm{𝕜}[\mathrm{Rep}_E^A]`$, the coordinate ring of the scheme $`\mathrm{Rep}_E^A`$, and write $`\mathrm{𝕜}_\rho `$ for the 1-dimensional $`R`$-module such that $`f\mathrm{𝕜}[\mathrm{Rep}_E^A]`$ acts in $`\mathrm{𝕜}_\rho `$ via multiplication by the number $`f(\rho )\mathrm{𝕜}`$. We identify the algebra $`R^\mathrm{e}=RR`$ with $`\mathrm{𝕜}[\mathrm{Rep}_E^A\times \mathrm{Rep}_E^A].`$ We have the tautological algebra map $`\mathrm{𝗋𝖾𝗉}:AR\mathrm{End},a\widehat{a}`$.
Further, we consider the algebra $`R\mathrm{End}R`$. There are two algebra maps $`\mathrm{𝗋𝖾𝗉}_l,\mathrm{𝗋𝖾𝗉}_r:AR\mathrm{End}R,`$ given by the formulas $`a\widehat{a}1,`$ and $`a1\widehat{a},`$ respectively. These two maps make $`R\mathrm{End}R`$ an $`A`$-bimodule.
We define the double-tangent $`R`$-bimodule by
$$𝒯_E^\mathrm{e}(A):=\mathrm{Der}(A,R_\mathrm{𝕜}\mathrm{End}_\mathrm{𝕜}E_\mathrm{𝕜}R).$$
(13.2.1)
This is clearly a (not necessarily free) $`\mathrm{𝕜}[\mathrm{Rep}_E^A\times \mathrm{Rep}_E^A]`$-module whose geometric fiber at a point $`(\rho ,\phi )\mathrm{Rep}_E^A\times \mathrm{Rep}_E^A`$ is defined as
$$𝒯_E^\mathrm{e}(A)|_{(\rho ,\phi )}:=(\mathrm{𝕜}_\rho \mathrm{𝕜}_\phi )_{RR}𝒯_E^\mathrm{e}(A).$$
The corresponding coherent sheaf on $`\mathrm{Rep}_E^A\times \mathrm{Rep}_E^A`$ will be referred to as the double-tangent bundle on $`\mathrm{Rep}_E^A`$.
###### Lemma 13.2.2.
For any $`(\rho ,\phi )\mathrm{Rep}_E^A\times \mathrm{Rep}_E^A`$, one has a canonical isomorphism
$$𝒯_E^\mathrm{e}(A)|_{(\rho ,\phi )}\mathrm{Hom}_A(K_\rho ,E_\phi ).$$
###### Proof.
We observe that this $`A`$-bimodule structure on $`R\mathrm{End}R`$ commutes with the obvious $`R`$-bimodule structure. It follows that the latter induces an $`R`$-bimodule structure on each Hochschild cohomology group $`HH^p(A,R\mathrm{End}R),p0.`$ In particular, the sequence (5.4.1) becomes, in our present situation, the following exact sequence of $`R`$-bimodules (to be compared with (13.1.7)):
$`0HH^0(A,R`$ $`\mathrm{End}R)R\mathrm{End}R`$ (13.2.3)
$`\mathrm{Der}(A,R\mathrm{End}R)HH^1(A,R\mathrm{End}R)0.`$
We leave to the reader to verify that, for any two representations $`\rho ,\phi \mathrm{Rep}_E^A`$, the geometric fibers at $`(\rho ,\phi )`$ of the $`R`$-bimodules occurring in (13.2.3) are given by
$`HH^0(A,R\mathrm{End}R)|_{(\rho ,\phi )}=\mathrm{Hom}_A(E_\rho ,E_\phi ),`$
$`(R\mathrm{End}R)|_{(\rho ,\phi )}=\mathrm{Hom}_\mathrm{𝕜}(E_\rho ,E_\phi ),`$
$`\mathrm{Der}(A,R\mathrm{End}R)|_{(\rho ,\phi )}=\mathrm{Hom}_A(K_\rho ,E_\phi ),`$
$`HH^1(A,R\mathrm{End}R)|_{(\rho ,\phi )}=\mathrm{Ext}_A^1(E_\rho ,E_\phi ).`$
The statement of the Lemma is now clear. ∎
Next, we consider the following maps
where $`m:\mathrm{End}\mathrm{End}\mathrm{End}`$ denotes the multiplication in the algebra $`\mathrm{End}`$. We observe that the two maps above are morphisms of $`A`$-bimodules, i.e., of $`A^\mathrm{e}`$-modules. Let $`\tau `$ denote the composite map $`A^\mathrm{e}=AAR\mathrm{End}R,`$ which is an $`A^\mathrm{e}`$-module map again.
Thus, from formula (13.2.1) we deduce that the $`A^\mathrm{e}`$-module map $`\tau `$ gives rise to a canonical map
$$\mathrm{𝗋𝖾𝗉}^\mathrm{e}:𝔻\mathrm{𝐞𝐫}(A):=\mathrm{Der}(A,A^\mathrm{e})\stackrel{\tau }{}\mathrm{Der}(A,R\mathrm{End}R)=:𝒯_E^\mathrm{e}(A).$$
### 13.3. Double-derivations for a free algebra.
Throughout this subsection we let $`A=\mathrm{𝕜}x_1,\mathrm{},x_r`$ be a free associative algebra on $`r`$ generators. For each $`i=1,\mathrm{},r`$, we introduce an element $`D_i𝔻\mathrm{𝐞𝐫}(A)`$ uniquely defined by the following conditions
$$D_i(x_i)=11\text{and}D_i(x_j)=0ji.$$
Now, let $`F:AA`$ be an algebra homomorphism, such that $`x_1F_1,\mathrm{},x_rF_r,`$ where $`F_1,\mathrm{},F_r\mathrm{𝕜}x_1,\mathrm{},x_r`$. The $`r`$-tuple $`F_1,\mathrm{},F_r`$ determines $`F`$ uniquely, and we will think of the elements $`F_i\mathrm{𝕜}x_1,\mathrm{},x_r`$ as functions $`F_i=F_i(x_1,\mathrm{},x_r)`$ in $`r`$ non-commutative variables $`x_1,\mathrm{},x_r`$.
We define the Jacobi matrix for the map $`F`$ to be the following $`AA`$-valued $`r\times r`$-matrix
$$DF=D_i(F_j)_{i,j=1,\mathrm{},r}\mathrm{Mat}_r(AA),$$
(13.3.1)
where $`D_i(F_j)`$ denotes the image of $`F_jA`$ under the derivation $`D_i:AAA`$ introduced above.
The following result is proved by a straightforward computation, see \[HT, Lemma 6.2.1\].
###### Proposition 13.3.2 (Chain rule).
For any two algebra homomorphisms $`F,G:AA,`$ in $`\mathrm{Mat}_r(AA)`$ one has
$$D(G{}_{^{^{}}}{}^{}F)=DG{}_{^{^{}}}{}^{}DF,\text{where}|DG{}_{^{^{}}}{}^{}DF|_{kl}:=\underset{i}{}D_k(G(F(x_i)))D_i(x_l).\mathrm{}$$
Next, we fix a finite dimensional vector space $`E`$ and consider the scheme $`\mathrm{Rep}_E^A`$, which is canonically isomorphic to $`\mathrm{End}_\mathrm{𝕜}(E)\times \mathrm{}\mathrm{End}_\mathrm{𝕜}(E)`$ via the identification, cf. (12.1.3)
$$\mathrm{Rep}_E^A\rho \left(X_1:=\rho (x_1),\mathrm{},X_n:=\rho (x_r)\right).$$
Analogously to the case of derivations considered in Example 12.4.2, any automorphism $`F:AA`$ induces an automorphism $`\widehat{F}`$ of the vector space $`\mathrm{Rep}_E^A`$. The map $`\widehat{F}`$ is given, in terms of the $`r`$-tuple $`F_1,\mathrm{},F_r\mathrm{𝕜}x_1,\mathrm{},x_r,`$ by the formula
$$\widehat{F}:(X_1,\mathrm{},X_n)(F_1(X_1,\mathrm{},X_n),\mathrm{},F_r(X_1,\mathrm{},X_n)).$$
(13.3.3)
For each $`i=1,\mathrm{},r`$, and any $`aA`$, we have an element $`D_i(a)=D_i^{}(a)D_i^{\prime \prime }(a)AA`$. Thus, given a point $`\rho \mathrm{Rep}_E^A`$, we have a well-defined element $`\rho (D_i(a))=\rho (D_i^{}(a))\rho (D_i^{\prime \prime }(a))\mathrm{End}_\mathrm{𝕜}(E)\mathrm{End}_\mathrm{𝕜}(E)`$. In particular, for each $`i,j=1,\mathrm{},r`$, there is an element $`\rho (D_i(F_j))\mathrm{End}_\mathrm{𝕜}(E)\mathrm{End}_\mathrm{𝕜}(E)`$. These elements give rise to a linear map
$$(DF)_\rho :\mathrm{Rep}_E^A\mathrm{Rep}_E^A,$$
given by the formula
$`(Z_1,\mathrm{},`$ $`Z_r)`$ (13.3.4)
$`({\displaystyle \rho (D_1^{}(F_j))Z_j\rho (D_1^{\prime \prime }(F_j))},\mathrm{},{\displaystyle \rho (D_r^{}(F_j))Z_j\rho (D_r^{\prime \prime }(F_j))}).`$
The differential of the map (13.3.3) can be read off from the Jacobi matrix (13.3.1) by means of the following result, proved in \[HT, Lemma 6.2.2\].
###### Proposition 13.3.5.
The differential of the map $`\widehat{F}`$ at a point $`\rho =(X_1,\mathrm{},X_n)\mathrm{Rep}_E^A`$ is a linear map $`d_\rho \widehat{F}:T_\rho \mathrm{Rep}_E^AT_\rho \mathrm{Rep}_E^A`$ given by formula (13.3.4), that is, we have $`d_\rho \widehat{F}=(DF)_\rho ,\rho \mathrm{Rep}_E^A.`$
### 13.4. The Crawley-Boevey construction.
Recall next that we have an isomorphism $`(A^\mathrm{e})^{\mathrm{op}}A^\mathrm{e}`$. Therefore, right multiplication in the algebra $`A^\mathrm{e}`$ makes $`A^\mathrm{e}`$, hence $`𝔻\mathrm{𝐞𝐫}(A)`$ a left $`A^\mathrm{e}`$-module, that is, an $`A`$-bimodule.
We consider $`T_A𝔻\mathrm{𝐞𝐫}(A)`$, the tensor algebra of the $`A`$-bimodule $`𝔻\mathrm{𝐞𝐫}(A)`$, and view the derivation $`\mathrm{\Delta }`$, see (13.1.4), as an element of $`T_A^1𝔻\mathrm{𝐞𝐫}(A)=𝔻\mathrm{𝐞𝐫}(A).`$ Given an element $`cA`$, viewed as an element of $`T_A^0𝔻\mathrm{𝐞𝐫}(A)=A`$, we introduce, following \[CrB\], an associative algebra $`\mathrm{\Pi }^c(A):=T_A𝔻\mathrm{𝐞𝐫}(A)/\mathrm{\Delta }c,`$ where $`\mathrm{\Delta }c`$ stands for the two-sided ideal generated by $`\mathrm{\Delta }c`$.
We first consider the case: $`c=0`$. To this end, we use formula (5.4.1) in the special case $`M=A^\mathrm{e}`$ and get a canonical exact sequence of $`A`$-bimodules:
$$0HH^0(A,A^\mathrm{e})A^\mathrm{e}\stackrel{\mathrm{ad}}{}𝔻\mathrm{𝐞𝐫}(A)HH^1(A,A^\mathrm{e})0.$$
(13.4.1)
This way, one proves
###### Proposition 13.4.2.
$`(𝗂)`$The above map $`A^\mathrm{e}𝔻\mathrm{𝐞𝐫}(A)`$ sends the element $`1A^\mathrm{e}`$ to the derivation $`\mathrm{\Delta }`$, see (13.1.4). Thus, $`HH^1(A,A^\mathrm{e})`$ is isomorphic to a quotient of $`𝔻\mathrm{𝐞𝐫}(A)`$ by the sub $`A`$-bimodule generated by $`\mathrm{\Delta }`$.
$`(\mathrm{𝗂𝗂})`$The isomorphism in $`(𝗂)`$induces a graded algebra isomorphism $`\mathrm{\Pi }^0(A)=T_A𝔻\mathrm{𝐞𝐫}(A)/\mathrm{\Delta }T_AHH^1(A,A^\mathrm{e})`$.∎
In the general case of an arbitrary $`cA`$, the grading on the tensor algebra induces a natural increasing filtration on the algebra $`T_A𝔻\mathrm{𝐞𝐫}(A)/\mathrm{\Delta }c`$. Furthermore, part (ii) of Proposition 13.4.2 yields a canonical surjective graded algebra map
$$T_AHH^1(A,A^\mathrm{e}){}_{}{}^{}_{}^{}\mathrm{gr}T_A𝔻\mathrm{𝐞𝐫}(A)/\mathrm{\Delta }c.$$
It would be interesting to find some sufficient conditions for the map above to be an isomorphism (a Poincaré-Birkhoff-Witt type isomorphism).
The case $`c=1`$ is especially important. To explain the geometric meaning of the algebra $`\mathrm{\Pi }^1(A)`$, we return to the setup of Example 13.1.5 and let $`A=\mathrm{𝕜}[X]`$ be the coordinate ring of a smooth affine algebraic curve.
###### Theorem 13.4.3.
The algebra $`\mathrm{\Pi }^1(\mathrm{𝕜}[X])`$ is canonically isomorphic, as a filtered algebra, to $`𝒟(X,\mathrm{\Omega }_X^{1/2})`$, the filtered algebra of twisted differential operators acting on half-forms on $`X`$.
###### Remark 13.4.4.
It has been shown in \[CrB\], that there is a natural graded algebra isomorphism $`\mathrm{\Pi }^0(\mathrm{𝕜}[X])\mathrm{𝕜}[T^{}X],`$ where $`T^{}X`$ stands for the total space of the cotangent bundle on the curve $`X`$.
We refer to \[BB\] for the basics of the theory of twisted differential operators on an algebraic variety. According to this theory, sheaves of twisted differential operators on a smooth affine algebraic variety $`X`$ are classified by the group $`H_{\mathrm{DR}}^1(X,)`$. This group vanishes if $`X`$ is a curve. Thus, in the case at hand, one can replace the algebra $`𝒟(X,\mathrm{\Omega }_X^{1/2})`$ in Theorem 13.4.3 by the algebra $`𝒟(X,𝒪_X)`$ of usual (not twisted) differential operators on $`X`$. In this form, an algebra isomorphism $`\mathrm{\Pi }^1(A)𝒟(X,𝒪_X)`$ has been already established in \[CrB\].
The advantage of our present version of Theorem 13.4.3, involving twisted differential operators, is that the isomorphism of the Theorem becomes canonical. In particular, both the statement and proof of the Theorem generalize easily to the case of sheaves of algebras of twisted differential operators on a not necessarily affine smooth curve $`X`$. To explain this, observe that in the case of an affine curve we have $`A^\mathrm{e}=\mathrm{𝕜}[X\times X]`$, hence the bimodule $`\mathrm{\Omega }_{\mathrm{nc}}^1AAA`$ is the ideal of the diagonal divisor $`X_\mathrm{\Delta }X\times X.`$ Therefore, we have
$$𝔻\mathrm{𝐞𝐫}A=\mathrm{Hom}_{A^e}(\mathrm{\Omega }_{\mathrm{nc}}^1A,AA)=\mathrm{\Gamma }(X\times X,𝒪_{X\times X}(X_\mathrm{\Delta })),$$
(13.4.5)
is the space of regular functions on $`(X\times X)X_\mathrm{\Delta }`$ with at most simple pole along the diagonal divisor $`X_\mathrm{\Delta }.`$
Now, in the case of an arbitrary smooth, not necessarily affine, curve $`X`$ it is natural to define the sheaf $`\mathrm{\Omega }_{\mathrm{nc}}^1(𝒪_X)`$ to be the ideal sheaf of the diagonal divisor $`X_\mathrm{\Delta }X\times X,`$ and to put
$$𝔻\mathrm{𝐞𝐫}(𝒪_X)=\mathrm{Hom}_{𝒪_{X\times X}}(\mathrm{\Omega }_{\mathrm{nc}}^1(𝒪_X),𝒪_{X\times X})𝒪_{X\times X}(X_\mathrm{\Delta }),$$
see (13.4.5). Repeating the definitions above, on constructs a sheaf $`\mathrm{\Pi }_X`$ of filtered algebras on $`X`$ that corresponds, locally in the Zariski topology, to the algebra $`\mathrm{\Pi }(A)`$. The sheaf-theoretic version of Theorem 13.4.3 says that there is a canonical isomorphism between $`\mathrm{\Pi }_X`$, viewed as a sheaf of filtered algebras on $`X`$ in the Zariski topology, and the sheaf of twisted differential operators on $`X`$ acting on half-forms. We note that, for any complete curve of genus $`1`$, the sheaf of twisted differential operators acting on half-forms is not isomorphic to the sheaf of usual, not twisted, differential operators.
### 13.5. Sketch of proof of Theorem 13.4.3.
First of all, for any associative not necessarily commutative algebra $`A`$ we have the tautological $`A`$-bimodule imbedding $`\mathrm{\Omega }_{\mathrm{nc}}^1A{}_{}{}^{}_{}^{}AA`$ and an $`A`$-bimodule map $`\mathrm{ad}:AA𝔻\mathrm{𝐞𝐫}A`$. Composing these two maps, yields a canonical $`A`$-bimodule morphism
$$\varphi :\mathrm{\Omega }_{\mathrm{nc}}^1A𝔻\mathrm{𝐞𝐫}A=\mathrm{Hom}_{A^\mathrm{e}}(\mathrm{\Omega }_{\mathrm{nc}}^1A,A^\mathrm{e}).$$
(13.5.1)
We observe that the morphism above is self-dual, i.e., applying the functor $`\mathrm{Hom}_{A^\mathrm{e}}(,A^\mathrm{e})`$ to $`\varphi `$ one gets the same morphism $`\varphi `$ again.
###### Lemma 13.5.2.
Assume that the bimodule $`AA`$ has trivial center, that is for $`xAA`$ we have $`ax=xa,aAx=0.`$ Then, the map $`\varphi `$ in (13.5.2) is injective and one has a canonical $`A`$-bimodule isomorphism
$$\mathrm{\Pi }_1(A)\mathrm{Coker}(\varphi )=𝔻\mathrm{𝐞𝐫}A/\mathrm{\Omega }_{\mathrm{nc}}^1A.$$
###### Proof.
For any three objects $`UVW`$ of an abelian category, one has a short exact sequence
$$0V\stackrel{f}{}W(V/U)\stackrel{g}{}W/U0,$$
(13.5.3)
where the map $`f`$ is the direct sum of the imbedding $`V{}_{}{}^{}_{}^{}W`$ with the projection $`V{}_{}{}^{}_{}^{}V/U`$ and the map $`g`$ is given by $`g(w(v\mathrm{mod}U))=w\mathrm{mod}U.`$
Now, let $`A`$ be an associative algebra satisfying the assumptions of the lemma. We let the triple $`UVW`$ to be the following triple of $`A`$-bimodules $`U=\mathrm{\Omega }_{\mathrm{nc}}^1AV=AA𝔻\mathrm{𝐞𝐫}A,`$ where the rightmost imbedding induced by the map $`\mathrm{ad}`$, which is injective by our assumptions. From the fundamental short exact sequence we get $`A=AA/\mathrm{\Omega }_{\mathrm{nc}}^1A=V/U`$. Further, the bimodule $`\mathrm{\Pi }_1^1(A)`$ is by definition the quotient $`𝔻\mathrm{𝐞𝐫}AA=WV/U`$ by the image of $`AA=U`$. Hence, the short exact sequence (13.5.3) yields the desired isomorphism
$$\mathrm{\Pi }_1^1(A)(𝔻\mathrm{𝐞𝐫}AA)/AA(WV/U)/VW/U=𝔻\mathrm{𝐞𝐫}A/\mathrm{\Omega }_{\mathrm{nc}}^1A.$$
We can now proof Theorem 13.4.3. By standard arguments, see e.g., \[BB\], to prove the theorem it suffices to construct a canonical isomorphism of Atiyah algebras:
(13.5.4)
The extension in the bottom row of this diagram can be computed explicitly. Specifically, one shows that this extension is canonically isomorphic, to (the spaces of global sections of) the following extension of sheaves
$$0𝒪_{X\times X}/𝒪_{X\times X}(X_\mathrm{\Delta })𝒪_{X\times X}(X_\mathrm{\Delta })/𝒪_{X\times X}(X_\mathrm{\Delta })𝒪_{X\times X}(X_\mathrm{\Delta })/𝒪_{X\times X}0.$$
The quotient sheaves on both sides are nothing but the cotangent and tangent sheaf on $`X_\mathrm{\Delta }=X`$, respectively. So the above extension reads
$$0𝒪_X𝒪_{X\times X}(X_\mathrm{\Delta })/𝒪_{X\times X}(X_\mathrm{\Delta })𝒯_X0.$$
(13.5.5)
Further, we have the following diagram of short exact sequences, cf. (13.1.6).
(13.5.6)
In the bottom row of the diagram above we have used shorthand notation $`\mathrm{\Gamma }()`$ for $`\mathrm{\Gamma }(X\times X,)`$; this row is obtained by applying the global sections functor to the natural extension of sheaves on $`X\times X`$. The vertical isomorphism $`\mathrm{\Phi }`$, in the diagram, follows from the identification $`\mathrm{Der}A=𝒯(X)`$, with the space of regular vector fields on $`X`$. The vertical isomorphism $`\mathrm{\Psi }`$ comes from (13.4.5).
Observe that the function $`1AA`$ corresponds under the above identifications to the element $`\mathrm{Id}_\mathrm{\Omega }\mathrm{Hom}_{A^e}(\mathrm{\Omega }_{\mathrm{nc}}^1A,\mathrm{\Omega }_{\mathrm{nc}}^1A)`$. Therefore, in the diagram we have $`j(1)=\mathrm{\Delta }`$, and the map $`j`$ is nothing but the imbedding $`\mathrm{ad}:AA{}_{}{}^{}_{}^{}\mathrm{Der}(A,AA),`$ of inner derivations.
Thus, we use the isomorphism $`𝒯(X)=\mathrm{Der}A=𝔻\mathrm{𝐞𝐫}A/A^\mathrm{e},`$ see (13.5.6), to identify the short exact sequence in (13.5.5) with the canonical short exact sequence
$$0A𝔻\mathrm{𝐞𝐫}(A)/\mathrm{\Omega }_{\mathrm{nc}}^1A\mathrm{Der}A0.$$
This short exact sequence can be identified with the extension in the top row of diagram (13.5.4) using Lemma 13.5.2. ∎
## 14. Noncommutative Symplectic Geometry
### 14.1.
Let $`A`$ be an associative algebra and $`\omega \mathrm{DR}^2(A)`$ a noncommutative $`2`$-form. Contraction with $`\omega `$ gives a linear map $`i_\omega :\mathrm{Der}(A)\mathrm{DR}^1(A),\theta i_\theta \omega .`$
###### Definition 14.1.1.
The pair $`(A,\omega )`$ is called a *noncommutative symplectic manifold* if $`d\omega =0`$ in $`\mathrm{DR}^3(A)`$, i.e., the $`2`$-form $`\omega \mathrm{DR}^2(A)`$ is closed, and furthermore, $`\omega `$ is *nondegenerate*, i.e., the map $`i_\omega :\mathrm{Der}(A)\mathrm{DR}^1(A)`$ is a bijection.
Fix a noncommutative symplectic manifold $`(A,\omega )`$.
###### Definition 14.1.2.
A derivation $`\theta \mathrm{Der}(A)`$ is called *symplectic* if $`_\theta \omega =0`$. We denote by $`\mathrm{Der}_\omega (A)`$ the Lie algebra of all symplectic derivations, i.e., $`\mathrm{Der}_\omega (A)=\{\theta \mathrm{Der}(A)_\theta \omega =0\}`$.
The space $`\mathrm{Der}_\omega (A)`$ inherits the Lie algebra structure from $`\mathrm{Der}(A)`$ given by commutators.
###### Lemma 14.1.3.
A derivation $`\theta `$ is symplectic if and only if $`i_\theta \omega `$ is closed in $`\mathrm{DR}^1(A)`$.
###### Proof.
This is simply an application of Cartan’s formula, namely
$$_\theta \omega =i_\theta d\omega +di_\theta \omega =d(i_\theta \omega ),$$
since $`d\omega =0`$ by assumption. ∎
Recall that, viewing $`𝖱(A)`$ as $`\mathrm{DR}^0(A)`$, we have a map $`d:𝖱(A)\mathrm{DR}^1(A)`$. For every $`f𝖱(A)`$, $`df`$ is exact in $`\mathrm{DR}^1(A)`$, hence closed. The previous lemma shows that the closed forms in $`\mathrm{DR}^1(A)`$ are identified with the symplectic vector fields. We let $`\theta _f`$ denote the symplectic vector field associated to $`df`$ under this identification. As in the classical theory, $`\theta _f`$ is called the Hamiltonian vector field associated to $`f`$. We then define a Poisson bracket on $`𝖱(A)`$ by
$$\{f,g\}=i_{\theta _f}(dg).$$
Notice that since $`dg\mathrm{DR}^1(A)`$, $`\{f,g\}`$ is indeed contained in $`\mathrm{DR}^0(A)=𝖱(A)`$. It is clear that we have the following several equalities
$$\{f,g\}:=i_{\theta _f}(dg)=i_{\theta _f}i_{\theta _g}\omega =i_{\theta g}i_{\theta f}\omega =i_{\theta _g}(df)=_{\theta _f}g=_{\theta _g}f.$$
###### Theorem 14.1.4.
(i) $`\{,\}`$ makes $`𝖱(A)`$ a Lie algebra.
(ii) The map $`f\theta _f`$ is a Lie algebra homomorphism from $`𝖱(A)`$ to $`\mathrm{Der}_\omega (A)`$.
###### Proof.
The skew symmetry of $`\{,\}`$ is immediate. We will first establish (ii), which will then prove (i). Let $`\theta ,\gamma `$ be arbitrary derivations from $`A`$ to $`A`$. Then by our standard identities we have
$$i_{[\theta ,\gamma ]}=_\theta i_\gamma i_\gamma _\theta =di_\theta i_\gamma +i_\theta di_\gamma i_\gamma di_\theta i_\gamma i_\theta d.$$
Now we specialize to the case $`\theta =\theta _f`$ and $`\gamma =\theta _g`$ and consider $`i_{[\theta _f,\theta _g]}\omega `$. By the definition of the Hamiltonian vector field, we know that $`i_{\theta _f}\omega =df`$ and $`i_{\theta _g}\omega =dg`$. Also, $`d\omega =0`$ by definition, so we are left with
$$i_{[\theta _f,\theta _g]}\omega =di_{\theta _f}(dg)+i_{\theta _f}d(dg)i_{\theta _g}d(df).$$
The latter two terms are zero, since $`d^2=0`$. Therefore,
$$i_{[\theta _f,\theta _g]}\omega =di_{\theta _f}(dg)=d\{f,g\}.$$
But by definition, $`\theta _{\{f,g\}}`$ is the unique symplectic derivation such that $`i_{\theta _{\{f,g\}}}\omega `$ $`=d\{f,g\}`$. Hence $`\theta _{\{f,g\}}=[\theta _f,\theta _g]`$, which establishes (ii). Since $`f\theta _f`$ is an isomorphism (*a priori* only of vector spaces), this shows that $`\{f,g\}`$ must satisfy the Jacobi identity (finishing (i)). Indeed, if we choose $`h𝖱(A)`$, then the equality $`\theta _f\theta _gh\theta _g\theta _fh=\theta _{\{f,g\}}h`$ becomes $`\{f,\{g,h\}\}\{g,\{f,h\}\}=\{\{f,g\},h\}`$, which is precisely Jacobi’s identity after some rearranging. ∎
###### Example 14.1.5.
Let $`(V,\omega _V)`$ be a finite dimensional symplectic vector space. We claim that the symplectic structure on $`V`$ induces a noncommutative symplectic structure on $`A=T^{\text{}}(V^{})`$, the tensor algebra. Explicitly, let $`x_1,\mathrm{},x_n,y_1,\mathrm{},y_nV^{}`$ be canonical coordinates in $`V`$, so that $`\omega _V=_{i=1}^ndx_idy_i`$. We will see in the next section that one has
$$\mathrm{Der}(A)=AV,\mathrm{DR}^1(A)=AV^{},\mathrm{\Omega }_{\mathrm{nc}}^i(A)=AT^i(V^{})A,i0.$$
We put $`\omega _A:=_{i=1}^n\mathrm{\hspace{0.17em}1}(x_iy_i)1AT^2(V^{})A=\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$. Further, the nondegeneracy of $`\omega _V`$ implies that the assignment $`v\omega _V(v,)`$ yields a vector space isomorphism $`V\stackrel{_{}}{}V^{}`$. The latter induces an isomorphism
$$\mathrm{id}_A\omega _V:\mathrm{Der}(A)=AV\stackrel{_{}}{}AV^{}=\mathrm{DR}^1(A).$$
It is easy to verify that the last isomorphism is nothing but the map $`\theta i_\theta \omega _A`$ arising from the noncommutative 2-form $`\omega _A\mathrm{DR}^2(A)`$. Thus, $`(A=T(V^{}),\omega _A)`$ is a noncommutative symplectic structure. $`\mathrm{}`$
###### Question 14.1.6.
$`(𝗂)`$Given a noncommutative symplectic structure on an associative algebra $`A`$, can one define a Lie super-algebra structure on $`_{i1}\mathrm{\Omega }_{\mathrm{nc}}^i(A)`$ which is a noncommutative analogue of the Lie super-algebra structure of Lemma 6.3.4 ?
$`(\mathrm{𝗂𝗂})`$In case of a positive answer to part $`(𝗂)`$, does the Lie super-algebra structure on $`_{i1}\mathrm{\Omega }_{\mathrm{nc}}^i(A)`$ combined with the associative product on $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ give rise to the structure of noncommutative Gerstenhaber algebra ?
### 14.2. Noncommutative ‘flat’ space
In Noncommutative Geometry, the free associative algebra $`\mathrm{𝕜}x_1,\mathrm{},`$ $`x_n`$ plays the role of coordinate ring of an $`n`$-dimensional affine space. It will be convenient to introduce an $`n`$-dimensional $`\mathrm{𝕜}`$-vector space $`V`$ with coordinates $`x_1,\mathrm{},x_n`$ (thus, the elements $`x_1,\mathrm{},x_n`$ form a base in $`V^{}`$, the dual space). This allows to adopt a ‘coordinate free’ point of view and to identify the algebra $`\mathrm{𝕜}x_1,\mathrm{},x_n`$ with $`A=T^{\text{}}(V^{})`$, the tensor algebra of $`V^{}`$.
The derivations from $`A`$ to an $`A`$-bimodule $`M`$ are specified precisely by a linear map from $`V^{}`$ to $`M`$–it is then extended to a derivation uniquely by the Leibniz rule. So, $`\mathrm{Der}(A,M)VM`$. Therefore, the functor $`\mathrm{Der}(A,)`$ is represented by the free $`A`$-bimodule generated by the space $`V^{}`$. Thus, we find
$$\mathrm{\Omega }_{\mathrm{nc}}^1(A)AV^{}AA(_{i>0}T^i(V^{}))A\overline{A},$$
Hence, for any $`p1,`$ we obtain
$$\mathrm{\Omega }_{\mathrm{nc}}^p(A)=T_A^p\mathrm{\Omega }_{\mathrm{nc}}^1(A)=AV^{}A\mathrm{}V^{}A\text{(}p\text{ factors }V\text{)}.$$
Observe that the assignment $`a_1v_1\mathrm{}v_pa_{p+1}a_1v_1\mathrm{}v_pa_{p+1}`$ gives an imbedding $`T_A^p\mathrm{\Omega }_{\mathrm{nc}}^1(A){}_{}{}^{}_{}^{}AA`$ (= free product of two copies of $`A`$). Further, an element of the form $`\mathrm{}v_11_Av_21_A\mathrm{}1_Av_m\mathrm{}`$ goes under this imbedding to the element $`\mathrm{}(v_1v_2\mathrm{}v_m)\mathrm{}`$. Thus we deduce that the imbedding above yields a graded algebra isomorphism
$$\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)\stackrel{_{}}{}AA=T(V^{})T(V^{}),$$
where the grading on the left-hand side corresponds to the total grading with respect to the second factor $`A=T(V^{})`$ on the right-hand side.
To describe $`\mathrm{DR}^0(A)`$, it is instructive to identify $`A=T(V^{})`$ with $`\mathrm{𝕜}x_1,\mathrm{},`$ $`x_n`$. Then, $`A`$ may be viewed as $`\mathrm{𝕜}`$-vector space whose basis is formed by all possible words in the alphabet formed by the $`x_i`$’s. The algebra structure is given by concatenation of words. Further, let $`A_{\mathrm{𝖼𝗒𝖼𝗅𝗂𝖼}}A`$ be the $`\mathrm{𝕜}`$-linear span of all cyclic words. It is clear that the composite map $`A_{\mathrm{𝖼𝗒𝖼𝗅𝗂𝖼}}{}_{}{}^{}_{}^{}A{}_{}{}^{}_{}^{}A/[A,A]`$ is a bijection. Thus, we may identify
$$\mathrm{DR}^0(A)=𝖱(A)=A/[A,A]=A_{\mathrm{𝖼𝗒𝖼𝗅𝗂𝖼}}$$
via this bijection.
It is more difficult to analyze the $`k^{\text{th}}`$ degree of $`\mathrm{DR}(A)`$ where $`k>0`$. In general we can only identify it as some quotient of $`ATV`$. However, in the particular case $`k=1`$, we find that
$$\mathrm{DR}^1(A)=\mathrm{\Omega }_{\mathrm{nc}}^1(A)/[\mathrm{\Omega }_{\mathrm{nc}}^1(A),A]AV^{},$$
as a complex vector space, since no other combinations of $`\mathrm{\Omega }_{\mathrm{nc}}^k(A)`$’s in the “denominator” will yield degree one.
So, for all $`xV^{}`$, we have an element $`dx=1x\mathrm{DR}^1(A)`$. Similarly, for every $`vV`$, we have $`_v:\mathrm{DR}^0(A)A`$ given by
$$_v(f)=df(v)$$
for all $`f\mathrm{DR}^0(A)=𝖱(A)=A/[A,A]`$. If $`x_i`$ is a basis of $`V^{}`$ and $`x^i`$ is the corresponding dual basis of $`V`$, consider the map $`d`$ which sends every $`f\mathrm{DR}^0(A)`$ to
$$\underset{j=1}{\overset{n}{}}\frac{f}{x^j}dx_j,$$
where we write $`\frac{}{x^j}=_{x^j}`$.
###### Example 14.2.1.
Let $`dim_\mathrm{𝕜}V=2`$, and equip $`V`$ with symplectic basis $`x`$ and $`y`$ and the standard symplectic 2-form $`dxdy`$. Let us calculate the Poisson bracket $`\{f,g\}`$ of two elements of $`𝖱(A)`$, where $`A=TV`$. View an element $`f𝖱(A)=TV/[TV,TV]`$ as a cyclic word in $`x`$ and $`y`$. We have already seen that $`\mathrm{DR}^1(A)=AV`$, so we can write $`df=f_xx+f_yy`$, where $`f_x,f_yA`$. This defines two maps $`\frac{}{x},\frac{}{y}:𝖱(A)𝖱(A)`$, the *partial derivative* maps, given by
$$\frac{}{x}:ff_xmod[A,A]\text{and}\frac{}{y}:ff_ymod[A,A].$$
Now, the correspondence $`AV=\mathrm{DR}^1(A)\mathrm{Der}(A)=AV^{}`$ given by the symplectic structure $`\omega \mathrm{DR}^2(A)`$ is nothing more than the canonical map $`AVAV^{}`$ that is the identity on $`A`$ and given by the identification $`VV^{}`$ given by the symplectic structure $`\omega _V`$ on $`V`$. Hence,
$$\theta _f=\frac{f}{x}x^{}+\frac{f}{y}y^{},$$
where $`x^{},y^{}`$ is the basis dual to $`x,y`$ under the correspondence induced by $`\omega _V`$ (that is, $`x^{}(v)=\omega (x,v)`$ for all $`vV`$, and similarly for $`y^{}`$). Then
$`\{f,g\}=\theta _f(dg)`$
$`=\left[{\displaystyle \frac{f}{x}}x^{}+{\displaystyle \frac{f}{y}}y^{}\right]\left[{\displaystyle \frac{g}{x}}x+{\displaystyle \frac{g}{y}}y\right]={\displaystyle \frac{f}{x}}{\displaystyle \frac{g}{y}}\omega _V(x,y)+{\displaystyle \frac{g}{x}}{\displaystyle \frac{f}{y}}\omega _V(y,x).`$
Since $`\omega _V=dxdy`$, $`\omega _V(x,y)=1`$ and $`\omega _V(y,x)=1`$. So we get the familiar formula
$$\{f,g\}=\frac{f}{x}\frac{g}{y}\frac{f}{y}\frac{g}{x}.$$
$`\mathrm{}`$
The next Proposition gives a non-commutative analogue of the classical Lie algebra exact sequence:
$$0\left[\begin{array}{c}\text{constant}\\ \text{functions}\end{array}\right]\left[\begin{array}{c}\text{regular}\\ \text{functions}\end{array}\right]\left[\begin{array}{c}\text{symplectic}\\ \text{vector fields}\end{array}\right]\mathrm{\hspace{0.17em}0},$$
associated with a connected and simply-connected symplectic manifold.
Let $`\omega =_idx_idy_i.`$ This is a symplectic structure on $`A=T(E^{}).`$
###### Proposition 14.2.2.
There is a natural Lie algebra central extension:
$$0\mathrm{𝕜}A/[A,A]\mathrm{Der}_\omega (A)\mathrm{\hspace{0.17em}0}.$$
###### Proof.
It is immediate from Lemma 14.1.3 and Theorem 11.4.7 that for the map: $`f\theta _f`$ we have: $`\mathrm{Ker}\{A/[A,A]\mathrm{Der}_\omega (A)\}`$ $`=\mathrm{Ker}d`$. Further, by Theorem 11.4.7 we get: $`\mathrm{Ker}d=\mathrm{𝕜}`$, and every closed element in $`\mathrm{DR}^1A`$ is exact. This yields surjectivity of the map: $`A/[A,A]\mathrm{Der}_\omega (A).`$
## 15. Kirillov-Kostant Bracket
In this section, we fix a finite-dimensional Lie algebra $`𝔤`$ over $`\mathrm{𝕜}`$. We are going to define a Poisson bracket on the polynomial algebra $`\mathrm{𝕜}[𝔤^{}]\mathrm{Sym}(𝔤)`$.
### 15.1. Coordinate formula.
Fix a basis $`e_1,\mathrm{},e_n`$ of $`𝔤`$, and let $`c_{ij}^k`$ denote the structure constants for this basis. That is, for all $`i,j=1,\mathrm{},n`$, we have that $`[e_i,e_j]=_{k=1}^nc_{ij}^ke_k.`$
Since $`𝔤`$ is finite-dimensional, it is isomorphic to its second dual. So, each $`e_i`$ gives rise to a linear functional on $`𝔤^{}`$, which we denote by $`x_i`$ (that is, for all $`\phi 𝔤^{}`$, $`x_i(\phi )=\phi (e_i)`$). Now, we can identify $`\mathrm{Sym}(𝔤)`$ with the polynomial algebra $`\mathrm{𝕜}[𝔤^{}]`$. Then if $`\phi ,\psi \mathrm{𝕜}[𝔤^{}]`$, we have that
$$\{\phi ,\psi \}=\underset{i,j,k=1}{\overset{n}{}}c_{ij}^k\frac{\phi }{x_i}\frac{\psi }{x_j}x_k.$$
Notice that in this case the Poisson bracket reduces the degree by $`1`$ (that is, $`\mathrm{deg}\{\phi ,\psi \}=\mathrm{deg}\phi +\mathrm{deg}\psi 1`$) since two derivatives are taken and a factor of $`x_k`$ is multiplied in. In the Poisson bracket associated to the Weyl quantization, the degree is reduced by two because no factor of $`x_k`$ is introduced.
There is also an explicit coordinate free way of writing the Poisson bracket. To this end, it is convenient to use the $`\mathrm{Sym}(𝔤)`$-realization of our algebra. Specifically, given two monomials $`a=a_1\mathrm{}a_n\mathrm{Sym}^n𝔤`$ and $`b=b_1\mathrm{}b_m\mathrm{Sym}^m𝔤`$, we have
$$\{a,b\}=\underset{i,j}{}[a_i,b_j]a_1\mathrm{}\widehat{a}_i\mathrm{}a_nb_1\mathrm{}\widehat{b}_j\mathrm{}b_m.$$
(15.1.1)
### 15.2. Geometric approach.
Choose two functions $`\phi ,\psi \mathrm{𝕜}[𝔤^{}]`$ (or even two smooth function on $`𝔤^{}`$). Fix some $`\lambda 𝔤^{}`$. Then $`d_\lambda \phi `$ and $`d_\lambda \psi `$ are linear functionals from $`T_\lambda 𝔤^{}𝔤^{}`$ to $`\mathrm{𝕜}`$. Since $`𝔤`$ is finite-dimensional, we identify $`𝔤`$ and $`𝔤^{}`$. By abuse of notation, we let $`d_\lambda \phi `$ denote the element of $`𝔤`$ corresponding to the linear function $`d_\lambda \phi `$ on $`𝔤^{}`$ under this identification. Then we set
$$\{\phi ,\psi \}(\lambda )=\lambda ,[d_\lambda \phi ,d_\lambda \psi ].$$
(15.2.1)
That is, we take the elements $`d_\lambda \phi `$ and $`d_\lambda \psi `$ of $`𝔤`$ and compute their Lie bracket. We then evaluate the linear functional on $`𝔤^{}`$ associated to this element of $`𝔤`$ on $`\lambda `$.
### 15.3. Symplectic structure on coadjoint orbits.
Let $`G`$ denote any connected Lie group such that $`𝔤=\mathrm{Lie}(G)`$ (in the future, we will call this a Lie group associated to $`𝔤`$). Consider the adjoint action of $`G`$ on $`𝔤`$. By transposing, this gives rise to the coadjoint action on $`𝔤^{}`$. We can then decompose $`𝔤^{}`$ into the disjoint union $`𝔤^{}=𝕆`$ of $`G`$-orbits.
According to a theorem of Kirillov and Kostant, every coadjoint orbit $`𝕆`$ admits a canonical symplectic structure. Explicitly, for any $`\lambda 𝕆_k`$ we have a natural isomorphism $`T_\lambda 𝕆=𝔤/𝔤(\lambda )`$, where $`𝔤(\lambda )`$ denotes the Lie algebra of the isotropy group of the point $`\lambda `$. Define the pairing
$$𝔤/𝔤(\lambda )\times 𝔤/𝔤(\lambda )\mathrm{𝕜}\text{by}(x,y)\lambda ,[x,y],x,y𝔤.$$
(15.3.1)
###### Proposition 15.3.2 (Kirillov-Kostant).
The pairing above descends to a well-defined skew-symmetric 2-form $`\omega _𝕆`$ on the coadjoint orbit $`𝕆`$.∎
The symplectic form $`\omega _𝕆`$ gives rise to a Poisson bracket $`\{,\}_𝕆`$ on the space of functions on the orbit $`𝕆`$. Since $`𝔤^{}`$ is the disjoint union of the $`𝕆`$’s, the formula
$$\{\phi ,\psi \}|_𝕆=\{\phi |_𝕆,\psi |_𝕆\}\text{for any coadjoint orbit}𝕆𝔤^{}$$
defines a Poisson structure on the whole of $`𝔤^{}`$. It is clear from (15.3.1) that this Poisson structure is nothing but the one given by formula (15.3.1).
###### Example 15.3.3.
Let $`E`$ be a finite dimensional vector space and $`𝔤=\mathrm{End}E`$ the Lie algebra of endomorphisms of $`E`$ with the commutator bracket. The trace provides an invariant bilinear form on $`𝔤`$, and this allows us to identify $`𝔤`$ with its dual $`𝔤^{}`$. Each conjugacy class $`𝕆𝔤`$ becomes a coadjoint orbit in $`𝔤^{}`$. For $`p𝕆𝔤`$, the Lie algebra $`𝔤(p)`$ is nothing but the centralizer of $`p`$ in $`𝔤=\mathrm{End}E`$. Thus, Kirillov-Kostant 2-form on the tangent space at $`p`$ is given by
$$\omega _p:𝔤/𝔤(p)\times 𝔤/𝔤(p)\mathrm{𝕜},(x,y)\mathrm{Tr}(p[x,y]),x,y𝔤.$$
$`\mathrm{}`$
We now give an interesting example of a noncommutative Kirillov-Kostant structure. Let $`A=\mathrm{𝕜}[e]/(e^2e)=\mathrm{𝕜}1\mathrm{𝕜}e`$. It is then possible to calculate $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ and $`\mathrm{DR}(A)`$ concretely. In particular, we find that
$$H^j(\mathrm{DR}(A))=\{\begin{array}{cc}\mathrm{𝕜},\hfill & \text{if }j\text{ is even;}\hfill \\ 0,\hfill & \text{if }j\text{ is odd.}\hfill \end{array}$$
Let $`\omega =edede\mathrm{\Omega }_{\mathrm{nc}}^2(A).`$ be a noncommutative 2-form on $`A`$.
###### Lemma 15.3.4.
The pair $`(A,\omega )`$ is a noncommutative symplectic structure.∎
This symplectic structure may be thought of as a ‘universal’ noncommutative Kirillov-Kostant structure. Indeed, fix a vector space $`E`$ of dimension $`dimE=n`$, and let $`\mathrm{Rep}_E^A`$ denote the variety of all algebra homomorphisms $`A\mathrm{End}_\mathrm{𝕜}E.`$ of $`A`$. Each representation is uniquely determined by the image of $`e`$, which maps to some idempotent, whose image has some rank. So, we can write
$`\mathrm{Rep}_E^A={\displaystyle \underset{kn}{}}\{p=p^2\mathrm{End}Edim(\mathrm{Im}p)=k\}={\displaystyle \underset{kn}{}}𝕆_k\mathrm{End}E,`$
where each $`𝕆_k`$ denotes the conjugacy class (under $`GL(E)`$) of idempotents with rank $`k`$. It is not difficult to see that the canonical map $`\mathrm{DR}^2(A)\mathrm{\Omega }_{\mathrm{com}}^2(\mathrm{Rep}_n^A)`$ sends $`edede`$ to the ordinary Kirillov-Kostant form (Proposition 15.3.2) on $`𝕆_k,k=0,1,2,\mathrm{}`$.
### 15.4. The algebra $`𝒪_𝒜`$
Let $`A`$ be an associative algebra, and as usual let $`𝖱(A)=A/[A,A]`$, which is a vector space. It is sometimes useful to form the commutative algebra $`𝒪_𝒜=\mathrm{Sym}𝖱(𝒜)`$. Consider $`\mathrm{Rep}_E^A`$, the variety of all representations of $`A`$ on the vector space $`V`$. Then we have defined a map $`𝖱(A)\mathrm{𝕜}[\mathrm{Rep}_E^A]`$, which we denote by $`a\mathrm{Tr}\widehat{a}`$. Extend this to a map $`𝒪_𝒜\mathrm{𝕜}[\mathrm{Rep}_{}^𝒜]`$.
Suppose now that $`A`$ has a noncommutative symplectic structure $`\omega `$. This makes $`𝖱(A)`$ into a Lie algebra, as we have seen before. Then $`𝒪_𝒜=\mathrm{Sym}𝖱(𝒜)`$ becomes a Poisson algebra with respect to the Kirillov-Kostant bracket (15.1.1).
### 15.5. Drinfeld’s bracket.
This section is taken from Drinfeld’s paper \[Dr\].
Let $`𝔞`$ be a Lie algebra. Define
$$𝖱(𝔞)=𝔞𝔞/\mathrm{𝕜}\text{-span of}\{xyyx,[x,y]zx[y,z]x,y,z𝔞\}.$$
As the notation suggests, $`𝖱(𝔞)`$ is meant to be a Lie analogue of $`𝖱(A)`$ for an associate algebra $`A`$. The set we are quotienting out by is meant to reproduce the key properties of the span of the commutators $`[A,A]`$ in the definition of $`𝖱(A)`$ for $`A`$ associative. Indeed, there is a striking similarity embodied in the following observation. Suppose $`\phi `$ is a linear functional $`𝖱(𝔞)\mathrm{𝕜}`$. Then this induces a symmetric invariant bilinear form. Simply define $`\tau _\phi :𝔞\times 𝔞\mathrm{𝕜}`$ by $`\tau _\phi (x,y)=\phi (xy)`$. Then $`\tau _\phi (x,y)=\tau _\phi (y,x)`$ and $`\tau _\phi ([x,y],z)=\tau _\phi (x,[y,z])`$ both follow directly from the definition of $`𝖱(𝔞)`$.
This is analogous to the associative case. Suppose we are given a linear functional $`\mathrm{Tr}:𝖱(A)\mathrm{𝕜}`$ (here $`A`$ is associative). Then this induces a symmetric bilinear form $`(a,a^{})=\mathrm{Tr}(aa^{})`$. Indeed, since $`𝖱(A)`$ consists of all cyclic words in $`A`$, $`(a,b)=\mathrm{Tr}(ab)=\mathrm{Tr}(ba)=(b,a)`$. It is also clear that $`(ab,c)=\mathrm{Tr}((ab)c)=\mathrm{Tr}(a(bc))=(a,bc)`$. So, we see that the second condition in the definition of $`𝖱(𝔞)`$ replaces associativity.
Let us apply these considerations to the case where $`𝔞`$ is the free Lie algebra on $`m`$ generators $`x_1,\mathrm{},x_m`$. Then $`𝖱(𝔞)`$ can be realized as pairs of words $`(w_1,w_2)`$, where each $`w_j`$ is composed of “Lie expressions” in the $`x_i`$’s. That is, each $`w_j`$ is a nested collection of Lie brackets of various generators $`x_i`$.
Claim. $`𝖱(𝔞)`$ has a natural Lie algebra structure.
Indeed, the Lie algebra structure is a noncommutative version of the Kirillov-Kostant Poisson bracket on $`\mathrm{Sym}(𝔤)`$, where $`𝔤`$ is a Lie algebra with a nondegenerate, invariant inner product $`(,)`$. In this context, $`\mathrm{Sym}(𝔤)`$ plays the role of $`𝖱(𝔞)`$.
As mentioned, we let $`𝔤`$ be a Lie algebra with nondegenerate, invariant inner product $`(,)`$. This inner product induces an isomorphism $`𝔤\times 𝔤𝔤^{}\times 𝔤^{}`$, which is the Lie algebra associated to $`(𝔤\times 𝔤)^{}`$. A pair of Lie words in $`𝔤`$ then yield a formula for a Lie word in the Lie algebra $`\mathrm{Lie}((𝔤\times 𝔤)^{})`$ via this isomorphism. This rule determines a formula for a Poisson bracket $`\{,\}`$ on $`𝖱(𝔞)`$.
We now define maps $`\frac{}{x_i}:𝖱(𝔞)𝖱(𝔞)`$ in the following way. Let $`f=(w_1,w_2)`$ be a pair of Lie words in the basis $`x_n`$ of $`𝔞`$–this is the form of an element of $`𝖱(𝔞)`$. Consider the substitution $`x_jx_j+z`$. This yields some expression in $`z`$ and the $`x_j`$’s, and take the $`z`$-linear part of it. Using the properties of $`𝖱(𝔞)`$, we can rewrite this linear part as
$$(z,w_3),$$
where $`w_3`$ is a word in the $`x_j`$’s. This word $`w_3`$ is defined to be $`\frac{f}{x_j}`$.
###### Lemma 15.5.1 (Poincaré Lemma).
Let $`f𝖱(𝔞)`$, and $`P_1,\mathrm{},P_n𝖱(𝔞)`$.
$`(𝗂)`$One has $`{\displaystyle _{j=1}^n}\{x_j,{\displaystyle \frac{f}{x_j}}\}=0.`$
$`(\mathrm{𝗂𝗂})`$If $`_{j=1}^n\{x_j,P_j\}=0,`$ then there exists $`f𝖱(𝔞)`$ such that $`P_j=\frac{f}{x_j}`$.
## 16. Review of (commutative) Chern-Weil Theory
### 16.1.
Let $`𝔤`$ be a finite-dimensional Lie algebra. Then the exterior algebra $`\mathrm{\Lambda }𝔤^{}`$ is equipped with a differential $`d`$ of degree $`+1`$. This differential is called the Koszul-Chevalley-Eilenberg differential, but we will usually shorten the name to the Koszul differential. The Koszul differential is defined on $`\mathrm{\Lambda }^1(𝔤^{})`$ by $`d(\lambda )(x,y)=\lambda ([x,y])`$. It is then extended to all of $`\mathrm{\Lambda }(𝔤^{})`$ by the Leibniz rule. Consider the Weyl algebra of $`𝔤`$, namely
$$W(𝔤)=\mathrm{Sym}(𝔤)^{}\mathrm{\Lambda }𝔤^{}.$$
Thus $`W(𝔤)`$ is a super-commutative algebra, a tensor product of a commutative algebra ($`\mathrm{Sym}(𝔤)^{}`$) with a super-commutative algebra ($`\mathrm{\Lambda }𝔤^{}`$). We wish to equip $`W(𝔤)`$ with a differential. To accomplish this, consider the graded Lie super-algebra $`\stackrel{~}{𝔤}=𝔤_0𝔤_1`$, where $`𝔤_0=𝔤_1=𝔤`$ as vector spaces, and the subscript indicates the degree of each component. The Lie super-algebra structure is given in the following way. We set $`[x,y]_{\stackrel{~}{𝔤}}=[x,y]_{𝔤_0}`$ for all $`x,y𝔤_0`$. We declare $`𝔤_1`$ to be an abelian subalgebra, that is, $`[z,w]_{\stackrel{~}{𝔤}}=0`$ for all $`z,w𝔤_1`$. So, we need only define $`[x,w]_{\stackrel{~}{𝔤}}`$, where $`x𝔤_0`$ and $`w𝔤_1`$. Define a map $`:\stackrel{~}{𝔤}\stackrel{~}{𝔤}`$ to be zero on $`𝔤_0`$ and the identity isomorphism $`𝔤_1𝔤_0`$. Then define $`[x,w]_{\stackrel{~}{𝔤}}=[x,w]_{𝔤_0}`$.
With $`\stackrel{~}{𝔤}`$ defined as above, we find that
$$\mathrm{\Lambda }(\stackrel{~}{𝔤}^{})=\mathrm{Sym}(𝔤)^{}\mathrm{\Lambda }𝔤^{}=W(𝔤).$$
Indeed, this is an isomorphism of graded algebras, if we give $`W(𝔤)`$ the grading
$$W_n(𝔤)=\underset{n=2p+q}{}\mathrm{Sym}^p(𝔤^{})\mathrm{\Lambda }^q(𝔤^{}).$$
So, if $`\lambda :𝔤\mathrm{𝕜}`$ is a linear functional, its image in $`\mathrm{Sym}(𝔤)^{}`$ has degree two, while its image in $`\mathrm{\Lambda }𝔤^{}`$ has degree one. We can use the identification $`W(𝔤)=\mathrm{\Lambda }\stackrel{~}{𝔤}^{}`$ to define a differential on $`W(𝔤)`$. Let $`:\mathrm{\Lambda }\stackrel{~}{𝔤}^{}\mathrm{\Lambda }\stackrel{~}{𝔤}^{}`$ denote the extension of the transpose of the map $`:\stackrel{~}{𝔤}\stackrel{~}{𝔤}`$ defined above, and let $`d:\mathrm{\Lambda }\stackrel{~}{𝔤}^{}\mathrm{\Lambda }\stackrel{~}{𝔤}^{}`$ denote the Koszul differential on $`\mathrm{\Lambda }\stackrel{~}{𝔤}^{}`$. Then we define $`d_W=d+`$.
A useful alternative picture of $`\stackrel{~}{𝔤}`$ is provided by the following. Define $`𝔤_\epsilon =𝔤(\mathrm{𝕜}[\epsilon ]/(\epsilon ^2))=𝔤\epsilon 𝔤`$. Clearly, $`𝔤`$ corresponds to $`𝔤_0`$ and $`\epsilon 𝔤`$ corresponds to $`𝔤_1`$ in $`\stackrel{~}{𝔤}`$. Define $`_\epsilon :𝔤_\epsilon 𝔤_\epsilon `$ by $`_\epsilon (\epsilon )=1`$ and $`_\epsilon (x)=0`$ for all $`x𝔤\{0\}`$.
###### Remark 16.1.1.
Notice that formally, $`_\epsilon =\frac{}{\epsilon }`$. $`\mathrm{}`$
Clearly, $`(𝔤_\epsilon ,_\epsilon )`$ is isomorphic to $`(\stackrel{~}{𝔤},)`$, and we will from here on identify these two constructions.
###### Proposition 16.1.2.
For the cohomology of the complex $`(W(𝔤),d_W)`$, we have
$$H^j(W(𝔤))=\{\begin{array}{cc}\mathrm{𝕜},\hfill & \text{if }j=0\text{;}\hfill \\ 0,\hfill & \text{if }j>0\text{.}\hfill \end{array}$$
###### Proof.
Since $`d_W=d+`$ is a sum of two anti-commuting differentials, we may view $`(W(𝔤),d+)`$ as a bicomplex. Therefore, we can compute the cohomology of the total differential via the standard spectral sequence for a bicomplex. Now, the differential $``$ is induced, by definition, by the differential on $`\stackrel{~}{𝔤}`$, which is the identity map $`\mathrm{id}:𝔤𝔤,`$ by definition. The two term complex $`𝔤𝔤,`$ given by this latter map is clearly acyclic. It follows that the induced differential $`:W(𝔤)W(𝔤)`$ is acyclic as well. Hence, the spectral sequence implies that the total differential $`d_W=d+`$ has trivial cohomology in all positive degrees. ∎
We now introduce a useful piece of notation. Let $`\lambda 𝔤^{}`$ be a linear functional. Then we can view $`\lambda `$ as both an element of $`\mathrm{Sym}(𝔤^{})`$ and an element of $`\mathrm{\Lambda }𝔤^{}`$. We denote the image of $`\lambda `$ in $`\mathrm{Sym}(𝔤^{})`$ by $`\lambda _+`$, and its image in $`\mathrm{\Lambda }𝔤^{}`$ is denoted by $`\lambda _{}`$.
Clearly, the elements $`\lambda _+`$ and $`\lambda _{}`$ generate $`\mathrm{Sym}(𝔤^{})`$ and $`\mathrm{\Lambda }𝔤^{}`$, respectively, so if we can calculate the action of $`d_W`$ on them we have a complete description of $`d_W`$ on $`W(𝔤^{})`$. Indeed, we find that
$$d_W\lambda _{}=\lambda _++\lambda _{}([,])\mathrm{Sym}^1𝔤^{}\mathrm{\Lambda }^2𝔤^{}=W_2(𝔤),$$
and
$$d_W\lambda _+=\underset{j=1}{\overset{n}{}}\mathrm{ad}^{}x_j(\lambda )x_j^{},$$
where $`\{x_1,\mathrm{},x_n\}`$ is a basis of $`𝔤`$, and $`x_j^{}`$ is the dual basis of $`𝔤^{}`$.
Using this characterization of $`W(𝔤)`$ and $`d_W`$, we are able to deduce the following.
###### Proposition 16.1.3 (Universal property).
If $`D`$ is a super-commutative DGA, and $`\phi :𝔤^{}D`$ is a $`\mathrm{𝕜}`$-linear map, then there exists a unique map of super-commutative DGA’s $`\phi _W:W(𝔤)D`$ such that $`\phi =\phi _W|_{\mathrm{\Lambda }^1𝔤^{}}`$.
###### Proof.
Extend $`\phi `$ to $`\stackrel{~}{\phi }:\mathrm{\Lambda }𝔤^{}D`$ as an algebra map. Notice that $`\stackrel{~}{\phi }`$ will not, in general, commute with the differential. But it is then possible to extend $`\stackrel{~}{\phi }`$ to $`\mathrm{Sym}(𝔤^{})`$ so that it kills this difference. This extension $`\phi _W`$ then commutes with differentials, and is the desired extension. Uniqueness is clear. ∎
For all $`x𝔤`$, we define the *Lie derivative* with respect to $`x`$, $`_x:W(𝔤)W(𝔤)`$ by $`L_x=\mathrm{ad}^{}x`$. This is a super-derivation of degree zero. We also define *contraction* $`i_x:W(𝔤)W(𝔤)`$ by $`i_x(f)=0`$ for all $`f\mathrm{Sym}(𝔤^{})`$ and $`i_x\alpha `$ is simply the contraction of $`\alpha `$ by $`x`$ for all $`\alpha \mathrm{\Lambda }(𝔤^{})`$. Notice that this is a degree $`1`$ map (hence it should be zero on $`\mathrm{Sym}(𝔤^{})`$ since these have only even degrees in $`W(𝔤)`$). Indeed, $`i_x`$ is a super-derivation.
###### Proposition 16.1.4.
The Cartan formula holds on $`W(𝔤)`$, that is,
$$_x=d_Wi_x+i_xd_W,\text{for all }x𝔤\text{.}$$
###### Proof.
This is immediate to verify on the generators of $`W(𝔤)`$. The result then follows from Lemma 3.3.2. ∎
###### Definition 16.1.5.
An element $`uW(𝔤)`$ is *basic* if $`_xu=i_xu=0`$ for all $`x𝔤`$. We let $`W(𝔤)_{\mathrm{basic}}`$ denote the set of all basic elements of $`W(𝔤)`$
###### Lemma 16.1.6.
$`W(𝔤)_{\mathrm{basic}}=(\mathrm{Sym}(𝔤^{}))^𝔤`$, where $`𝔤`$ acts on $`\mathrm{Sym}(𝔤^{})`$ by $`_x`$.
###### Proof.
Both $`_x`$ and $`i_x`$ vanish on $`(\mathrm{Sym}(𝔤^{}))^𝔤`$ by definition for all $`x𝔤`$. Conversely, $`i_x\alpha =0`$ for all $`x𝔤`$ and some $`\alpha \mathrm{\Lambda }𝔤^{}`$ forces $`\alpha =0`$, hence $`_{x𝔤}\mathrm{Ker}i_x\mathrm{Sym}(𝔤^{})`$. ∎
### 16.2. Connections on $`G`$-bundles.
We will now place the Weyl algebra in a geometric context. Suppose we have a principal $`G`$-bundle
$$\text{},$$
where $`G`$, $`P`$, and $`M`$ are all connected, $`G`$ is a Lie group, and the Lie algebra of $`G`$ is $`𝔤`$. A *connection* on $`P`$ is a $`𝔤`$-valued $`1`$-form $``$ on $`P`$, i.e., it is an element of $`\mathrm{\Omega }^1(P)𝔤`$, satisfying
* $``$ is $`𝔤`$-equivariant with respect to the diagonal action on $`\mathrm{\Omega }^1(P)`$ and $`𝔤`$, i.e., $`_x=0`$;
* For each $`x𝔤`$, let $`\xi _x`$ be the vector field on $`P`$ associated to $`x`$ by the $`G`$-action. Then $`(\xi _x)=x`$.
We call an element of $`\mathrm{\Omega }^{\text{}}(P)`$ *basic* if $`_x\omega =i_x\omega =0`$. We then have the following lemma.
###### Lemma 16.2.1.
Assume the group $`G`$ is connected. Then, the pullback along the bundle map yields a canonical isomorphism $`\mathrm{\Omega }^{\text{}}(P)_{\mathrm{basic}}\mathrm{\Omega }^{\text{}}(M)`$.
###### Proof.
It is clear that a differential form on the total space $`P`$ descends to a well-defined differential form on $`M`$ if and only if it is $`G`$-invariant, and annihilates all vectors tangent to the fibers of the bundle projection. But if $`G`$ is connected, the $`G`$-invariance of $`\alpha `$ is equivalent to $`_x\alpha =0,`$ for any $`x𝔤`$. ∎
We observe further that a connection gives rise to a linear map $`\mathrm{\Phi }^{}:𝔤^{}\mathrm{\Omega }^1(P)`$. Namely, $`\mathrm{\Phi }^{}(\lambda )=\lambda `$.
By the universal property of $`W(𝔤)`$, $`\mathrm{\Phi }^{}`$ extends to a DGA map $`\mathrm{\Phi }_W^{}:W(𝔤)\mathrm{\Omega }^{\text{}}(P)`$. By the connection conditions placed on $``$, we see that $`\mathrm{\Phi }_W^{}`$ commutes with $`_x`$ and $`i_x`$. So,
$`\mathrm{\Phi }_W^{}(W(𝔤)_{\mathrm{basic}})\mathrm{\Omega }^{\text{}}(P)_{\mathrm{basic}}\mathrm{\Omega }^{\text{}}(M).`$
But $`\mathrm{\Phi }_W^{}(W(𝔤)_{\mathrm{basic}})=\mathrm{\Phi }_W^{}((\mathrm{Sym}(𝔤^{}))^𝔤)`$, and $`(\mathrm{Sym}(𝔤^{}))^𝔤\mathrm{𝕜}[𝔤]^𝔤`$. The map $`\mathrm{\Phi }_W^{}`$ is called the *Chern character map*. We will also denote this map $`\mathrm{\Phi }_W^{}`$ by $`\mathrm{ch}`$.
###### Definition 16.2.2.
Let $``$ be a connection. Then the *curvature* of $``$, $`K()`$ is defined by
$$K()=d+\frac{1}{2}[,]\mathrm{\Omega }^2(P)𝔤.$$
If $`=_{j=1}^n\nu _jx_j`$, then
$$[,]=\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{n}{}}\nu _j\nu _k[x_j,x_k]\text{and}\frac{1}{2}[,]=\underset{1j<kn}{}\nu _j\nu _k[x_j,x_k].$$
It follows from the definition of curvature and the Chern character that for all $`\lambda 𝔤^{}`$, $`\mathrm{ch}(\lambda _+)=\lambda K()`$.
Next, we consider the ideal $`(\mathrm{Sym}^1(𝔤^{})1)W(𝔤)`$ of $`W(𝔤)`$. For simplicity, we let $`𝔤_+^{}=\mathrm{Sym}^1𝔤^{}1W(𝔤)`$. Since $`d_W(𝔤_+^{})g_+^{}W(𝔤)`$, we see that $`𝔤_+^{}W(𝔤)`$ is a differential ideal of $`W(𝔤)`$.
###### Definition 16.2.3.
The *Hodge filtration* of $`W(𝔤)`$ is a decreasing filtration $`W(𝔤)g_+^{}W(𝔤)\mathrm{}F^{p1}W(𝔤)F^pW(𝔤)`$ given by
$$F^pW(𝔤)=(𝔤_+^{})^pW(𝔤).$$
Notice that $`W(𝔤)/F^1W(𝔤)=\mathrm{\Lambda }𝔤^{}W(𝔤)`$.
###### Lemma 16.2.4.
For each $`p`$,
$$H^{2p}(F^pW(𝔤))=(\mathrm{Sym}^p𝔤^{})^𝔤.$$
### 16.3. Transgression map.
The short exact sequence
gives rise to a long exact sequence in cohomology. Recall that the cohomology of $`W(𝔤)`$ is the cohomology of a point (i.e., $`\mathrm{𝕜}`$ in degree zero, and zero elsewhere). The inclusion of $`F^pW(𝔤)F^1W(𝔤)`$ yields a homomorphism $`H^{2p}(F^pW(𝔤))H^{2p}(F^1W(𝔤))`$. Recalling that $`H^{2p}(F^pW(𝔤))=(\mathrm{Sym}^p𝔤^{})^𝔤`$, we obtain the exact sequence
$$\text{},$$
where $`d`$ denotes the Koszul differential.
But then $`H^{2p1}(\mathrm{\Lambda }𝔤^{},d)`$ is isomorphic to the $`2p1`$-Lie cohomology of $`𝔤`$, $`H_{\mathrm{Lie}}^{2p1}(𝔤)`$, which is isomorphic to $`H^{2p1}(G)`$. It is well know that the cohomology of $`G`$ can be calculated using only invariant differential forms on $`G`$.
A geometric meaning behind this can be found by considering the universal $`G`$-bundle $`EGBG`$. Then $`EG`$ is contractible, and $`W(𝔤)\mathrm{\Omega }^{\text{}}(EG)`$ and $`W_{\mathrm{basic}}(𝔤)\mathrm{\Omega }^{\text{}}(BG)`$. Then the $`2p`$-cohomology of $`F^1W(𝔤)`$ is calculating $`H^{2p}(BG)`$, so we obtain a homomorphism $`H^{2p}(BG)H^{2p1}(G)`$.
### 16.4. Chern-Simons formalism.
Let $`D`$ be a DGA, with differential $`d:D^nD^{n+1}`$. We set $`𝖱(D)=D/[D,D]`$, where $`[,]`$ here denotes the super-commutator, i.e.,
$$[x,y]=xy(1)^{(\mathrm{deg}x)(\mathrm{deg}y)}yx.$$
Then $`d`$ descends to a super-differential $`d`$ on $`𝖱(D)`$. Fix any $`aD^1`$. We define its *curvature* $`F:=da+a^2`$.
###### Proposition 16.4.1 (Bianchi Identity).
With $`D`$, $`a`$, and $`F`$ as above, the following identity holds in $`D`$:
$$(d+\mathrm{ad}a)F=0.$$
Proof. Observe that $`dF=d(da+a^2)=d^2ada+daa=ada+daa`$. Also,
$`(\mathrm{ad}a)F=[a,F]`$ $`=aF(1)^{(\mathrm{deg}a)(\mathrm{deg}F)}Fa`$
$`=aFFa=a(da+a^2)(da+a^2)a`$
$`=adadaa.\mathrm{}`$
Using the Bianchi identity, we see that in $`𝖱(D)`$ we have
$$d(F^n)=[a,F^n]=0,$$
since all super-commutators are zero in $`𝖱(D)`$. Hence the elements $`\frac{F^n}{n!}`$ are cocycles in $`𝖱(D)`$.
Consider the algebra $`D[t]=D\mathrm{𝕜}[t]`$. Take an element $`a_tD^1[t]=D^1\mathrm{𝕜}[t]`$. Define $`F_t=da_t+a_t^2`$, the curvature of $`a_t`$. Then a simple computation shows that
$$\frac{}{t}\left(\frac{F_t^n}{n!}\right)=d\left[\frac{1}{(n1)!}\frac{a_t}{t}F_t^{n1}\right].$$
In particular, take $`a_t=ta`$. Then $`F_t=tda+t^2a^2`$. Then we can integrate the above equation in $`t`$. Indeed, this yields that
$$\frac{F^n}{n!}=d\mathrm{cs}_{2n1}(a),$$
where
$$\mathrm{cs}_{2n1}(a)=_0^1a\frac{F_t^{n1}}{(n1)!}𝑑t.$$
The element $`\mathrm{cs}_{2n1}(a)`$ is called the $`2n1`$-*Chern-Simons class* of $`a`$. The map $`\frac{F^n}{n!}\mathrm{cs}_{2n1}(a)`$ is a trangression map.
### 16.5. Special case: $`𝔤=𝔤𝔩_n`$.
Let $`G=\mathrm{GL}_n`$. Then, giving a principal $`G`$-bundle $`PM`$ is the same thing as giving an ordinary vector bundle on $`M`$. Let $``$ be a connection on $`P`$. Then we define $`\mathrm{ch}_k=\frac{1}{k!}\mathrm{Tr}(K()^k)`$, which is
$$\mathrm{ch}_k=\frac{1}{k!}\mathrm{Tr}((d+\frac{1}{2}[,])^k).$$
Now, suppose that $`_0,\mathrm{},_N`$ are $`N+1`$ connections on the same bundle $`P`$. We wish to show that $`\mathrm{ch}_k`$ is independent of the connection used. This will follow essentially from the fact that the space of connections is convex. Let $`\mathrm{\Delta }`$ be the standard $`N`$-simplex in $`^{N+1}`$, that is,
$$\mathrm{\Delta }=\{(t_0,\mathrm{},t_N)^{N+1}t_j0,\underset{j=0}{\overset{n}{}}t_j=1\}.$$
For each $`t\mathrm{\Delta }`$, define $`(t)=_{j=0}^Nt_j_j`$. Now, it is a well-known fact that a connection $`_j`$ may be written in a local trivialization of $`P`$ as a sum of the usual differential and a matrix of $`1`$-forms, $`_j=d+A_j`$. Define
$$\mathrm{ch}_m^N(_0,\mathrm{},_N)=_\mathrm{\Delta }\mathrm{Tr}(K((t))^k)𝑑t.$$
This is a complicated expression in $`A_0,\mathrm{},A_N`$ and $`dA_0,\mathrm{},dA_n`$.
### 16.6. Quantized Weil algebra.
Let $`(𝔤,B)`$ be a Lie algebra equipped with an invariant, nondegenerate bilinear form $`B:𝔤\times 𝔤\mathrm{𝕜}`$. Then this form induces a canonical isomorphism $`𝔤𝔤^{}`$. In the standard Chern-Weil theory, we set
$$W(𝔤)=\mathrm{Sym}𝔤\mathrm{\Lambda }𝔤.$$
(The above formula is correct, since $`𝔤𝔤^{}`$.)
Following Alekseev-Meinrenken \[AM\], we would like to ‘quantize’ the algebra $`W(𝔤)`$ by replacing $`\mathrm{Sym}𝔤`$ by the universal enveloping algebra and $`\mathrm{\Lambda }𝔤`$ by the Clifford algebra $`\mathrm{Cliff}𝔤`$. Recall that
$$\mathrm{Cliff}𝔤=T𝔤/(xy+yx2B(x,y)).$$
We set
$$𝒲(𝔤)=U𝔤\mathrm{Cliff}𝔤.$$
It turns out that $`𝒲(𝔤)`$ has a differential. It is this algebra that acts as the quantized version of $`W(𝔤)`$.
This quantized Weyl algebra is connected to the work of Alekseev-Meinrenken in the following fashion. Consider the family of Lie algebras $`𝔤_t`$, where $`𝔤_t=𝔤`$ as a vector space and $`[x,y]_t=t[x,y]`$ for all $`x,y𝔤_t`$. Define
$$𝒲_t(𝔤)=𝒲(𝔤_t)=U𝔤_t\mathrm{Cliff}(𝔤_t).$$
Then $`𝒲_t(𝔤)`$ is a flat family of DGA’s, and $`𝒲_0(𝔤)W(𝔤)`$. Based on our previous work with deformations of associative algebras, we conclude that the family $`𝒲_t(𝔤)`$ induces a Poisson bracket on $`W(𝔤)`$ making it a Poisson DGA.
Let $`\mathrm{\Delta }\mathrm{Sym}^2(𝔤)W(𝔤)`$ be the canonical element corresponding to the inverse of the nondegenerate bilinear form $`B`$ on $`𝔤`$.
Using the deformation argument above, one proves
###### Theorem 16.6.1.
An invariant, nondegenerate bilinear form on $`𝔤`$ gives rise to a Poisson (super) algebra structure $`\{,\}`$ on $`W(𝔤)`$ such that
$$d_W(u)=\frac{1}{2}\{\mathrm{\Delta },u\},uW(𝔤).$$
## 17. Noncommutative Chern-Weil theory
### 17.1.
In the previous section, we discussed quantized Chern-Weil theory, which could be considered a part of noncommutative geometry “in the small.” That is, it is simply a deformation of the usual Chern-Weil theory. We now want to begin with a noncommutative algebra which will replace the Lie algebra $`𝔤`$.
Let $`A`$ be a (possibly noncommutative) associative algebra. Let $`A^{}`$ denote the dual of $`A`$. This is a coalgebra, with comultiplication map $`\mathrm{\Delta }:A^{}A^{}A^{}`$. For the sake of simplicity, we will use the Sweedler notation, that is, we write
$$\mathrm{\Delta }(\lambda )=\lambda ^{}\lambda ^{\prime \prime }$$
for all $`\lambda A^{}`$.
Following the paper \[C\], we define $`W_{\mathrm{nc}}(A)=T(A_+^{}A_{}^{})`$, where $`A_+^{}=A_{}^{}=A^{}`$. We make $`W_{\mathrm{nc}}(A)`$ graded algebra (under the usual multiplication for the tensor algebra) by taking
$$W_{\mathrm{nc}}(A)_p:=_{2n+m=p}(A_+^{})^n(A_{}^{})^m.$$
Define a differential $`d_W`$ on $`W_{\mathrm{nc}}(A)`$ by
$`d_W\lambda _{}=\lambda _++{\displaystyle \lambda _{}^{}\lambda _{}^{\prime \prime }}\text{and}d_W\lambda _+={\displaystyle (\lambda _+^{}\lambda _{}^{\prime \prime }\lambda _{}^{}\lambda _+^{\prime \prime })}.`$
As usual, $`\lambda `$ is a linear functional in $`A^{}`$, and $`\lambda _+`$ (respectively, $`\lambda _{}`$) represents its image in $`A_+^{}`$ (respectively $`A_{}^{}`$). This differential makes $`W_{\mathrm{nc}}(A)`$ a DGA.
Similarly to the construction of Bar-complex as a free product, it is sometimes useful to have the following alternative definition
$$W_{\mathrm{nc}}(A)=T\left((A_\epsilon )^{}\right),\text{where}A_\epsilon :=A\mathrm{𝕜}[\epsilon ]/(\epsilon ^2).$$
(17.1.1)
Let $`:W_{\mathrm{nc}}(A)W_{\mathrm{nc}}(A)`$ be the $`\mathrm{𝕜}`$-linear map sending $`\epsilon 1`$ and $`1_A0`$. Define
$$d_W=+\delta ,$$
where $`\delta `$ is essentially the differential dual to the multiplication map, as before. The two differentials $``$ and $`\delta `$ anti-commute, hence $`d_W{}_{^{^{}}}{}^{}d_{W}^{}=0`$. Observe that, in the presentation (17.1.1), the differential $``$ can be suggestively written as $`=/\epsilon `$.
As in the commutative and quantized cases, there is a Poincaré lemma for $`W_{\mathrm{nc}}(A)`$.
###### Lemma 17.1.2.
The cohomology of the complex $`(W_{\mathrm{nc}}(A),d_W)`$ are given by
$$H^j(W_{\mathrm{nc}}(A))=\{\begin{array}{cc}\mathrm{𝕜},\hfill & \text{if }j=0\text{;}\hfill \\ 0,\hfill & \text{if }j>0\text{.}\hfill \end{array}$$
###### Proof.
Suppose that the multiplication on $`A`$ is trivial. Then it is easy to see that $`H^j(W_{\mathrm{nc}}(A))`$ is zero for $`j>1`$, and $`\mathrm{𝕜}`$ for $`j=0`$. This follows, since the differential $`d_W`$ involved the transpose $`\delta `$ of the multiplication map, which would also be zero.
We construct an isomorphism of DGA’s from $`(W_{\mathrm{nc}}(A),d_W)`$ to $`(W_{\mathrm{nc}}(A),\stackrel{~}{d})`$, where $`\stackrel{~}{d}`$ is the differential given when the multiplication is trivial. Indeed, observe that $`W_{\mathrm{nc}}(A)`$ is independent of the multiplication on $`A`$. The map is given by
$$(W_{\mathrm{nc}}(A),d_W)(W_{\mathrm{nc}}(A),\stackrel{~}{d}),\lambda _{}\lambda _{},\lambda _+\lambda _++\lambda _{}^{}\lambda _{}^{\prime \prime }.$$
It is easy to check this is a DGA isomorphism.
An alternative proof of the Lemma can be obtained as follows. It is clear that the two term complex $`/_\epsilon :\epsilon A^{}A^{}`$ has trivial homology. Thus, the spectral sequence associated to the bicomplex given by the differentials $``$ and $`\delta `$ implies the result.∎
###### Lemma 17.1.3 (Universal Property).
Given a DGA $`D`$ and a $`\mathrm{𝕜}`$-linear map $`\phi :A^{}D`$, there is a unique extension $`\phi _W:W_{\mathrm{nc}}(A)D`$ such that $`\phi _W|_A^{}=\phi `$.
As usual, we set $`RW_{\mathrm{nc}}(A)=W_{\mathrm{nc}}(A)/[W_{\mathrm{nc}}(A),W_{\mathrm{nc}}(A)]`$, where the commutators are graded. This still has trivial cohomology (i.e., the cohomology of a point). The same calculation as before suffices–we simply take $`A`$ to have the trivial product.
###### Theorem 17.1.4 (\[C\]).
There exists a canonical transgression map
$$\text{}.$$
###### Definition 17.1.5.
Define $`I_+=W_{\mathrm{nc}}(A)A_+^{}W_{\mathrm{nc}}(A)W_{\mathrm{nc}}(A)`$. Then this is a differential ideal of $`W_{\mathrm{nc}}(A)`$, i.e., $`dI_+I_+`$. The *Hodge filtration* of $`W_{\mathrm{nc}}(A)`$ is given by $`F^pW_{\mathrm{nc}}(A)=I_+^p`$.
### 17.2. Example: case $`A=\mathrm{𝕜}`$.
In this case, $`W_{\mathrm{nc}}(A)`$ is the free algebra with generators $`\alpha `$ in degree one and $`\alpha `$ in degree two. Extend $``$ by the Leibniz rule to be a super-derivation of $`T(\alpha ,\alpha )`$ such that $`^2=0`$. Let $`F=\alpha +\alpha ^2`$. Then one finds that $`F=[\alpha ,F]`$. On $`W_{\mathrm{nc}}(A)`$, we have the Hodge filtration $`F^pW_{\mathrm{nc}}(A)`$, the two-sided ideal generated by terms involving $`\alpha `$ at least $`p`$ times. As before, one sees that $`(F^pW_{\mathrm{nc}}(A))W_{\mathrm{nc}}(A)`$. Finally, we consider the algebra $`𝖱(W_{\mathrm{nc}}(A)/F^pW_{\mathrm{nc}}(A))`$.
###### Theorem 17.2.1.
Let $`H^{\text{}}(𝖱(W_{\mathrm{nc}}(A)/F^pW_{\mathrm{nc}}(A)))`$ denote the cohomology of the complex$`𝖱(W_{\mathrm{nc}}(A)/F^pW_{\mathrm{nc}}(A))`$ with differential $`d_W=+\delta `$ (where $`\delta `$ is the differential dual to the multiplication). Then $`H^{\text{}}(𝖱(W_{\mathrm{nc}}(A)/F^pW_{\mathrm{nc}}(A)))`$ has a $`\mathrm{𝕜}`$-basis formed by the Chern-Simons classes $`\mathrm{cs}_{2n1}`$ for $`np`$.
###### Proof.
A $`\mathrm{𝕜}`$-basis of $`𝖱(W_{\mathrm{nc}}(A)/F^pW_{\mathrm{nc}}(A))`$ is given by $`\alpha ^{2k1}`$, $`k1`$, and $`(\alpha )^{\mathrm{}}`$, $`1\mathrm{}<p`$. Calculating the $`\delta `$-cohomology, we see that $`a^{2k1}`$ transgresses to $`(\alpha )^k`$ for $`1kp1`$. The remaining $`\alpha ^{2k1}`$’s are sent to the Chern-Simons classes, $`\mathrm{cs}_{2k1}`$ ($`kp`$). ∎
### 17.3. Gelfand-Smirnov bracket.
We fix a finite-dimensional associative algebra $`A`$ equipped with an invariant trace $`\mathrm{Tr}:A\mathrm{𝕜}`$, that is, $`\mathrm{Tr}(a_1a_2)=\mathrm{Tr}(a_2a_1)`$ for all $`a_1,a_2A`$. Write $`\mathrm{\Delta }\mathrm{Sym}^2(A^{})`$ for the canonical element corresponding to the bilinear form, i.e., such that $`B(x,y)=\mathrm{\Delta },xy,`$ for any $`x,yA`$.
The following result is a noncommutative analogue of Theorem 16.6.1
###### Theorem 17.3.1.
If the trace pairing $`a_1\times a_2\mathrm{Tr}(a_1a_2)`$ is non-degenerate, then the graded vector space $`𝖱(W_{\mathrm{nc}}(A))`$ has a canonical Lie super-algebra structure such that
$$d_W(u)=\frac{1}{2}\{\mathrm{\Delta },u\},u𝖱(W_{\mathrm{nc}}(A)).$$
###### Hint of Proof:.
The trace pairing on $`A`$ induces an isomorphism of vector spaces $`A\stackrel{_{}}{}A^{},a\mathrm{Tr}(a()).`$ Hence, we get an isomorphism $`\kappa :AA\stackrel{_{}}{}A^{}A^{}`$. Further, the trace pairing also gives rise to a non-degenerate skew-symmetric $`\mathrm{𝕜}`$-bilinear 2-form $`\omega `$ on the vector space $`AA`$ defined by the formula
$$(aa^{})\times (bb^{})\mathrm{Tr}(ab^{})\mathrm{Tr}(a^{}b).$$
Transporting this 2-form from $`AA`$ to $`A^{}A^{}`$ via the isomorphism $`\kappa `$ makes $`A^{}A^{}`$ a symplectic vector space. Therefore, the tensor algebra $`T(A^{}A^{})`$ acquires a natural structure of noncommutative symplectic manifold, see Example 14.1.5. Thus, we get a Lie bracket on $`𝖱\left(T(A^{}A^{})\right)`$.
Now, the algebra $`W_{\mathrm{nc}}(A))`$ is just $`T(A^{}A^{})`$, as an associative algebra. So, the above construction can be adapted, by inserting suitable signs (due to the fact that the first copy $`A^{}W_{\mathrm{nc}}(A)`$ is placed in degree 1 and the second copy is placed in degree 2), to produce the required Lie super-algebra structure on $`𝖱\left(W_{\mathrm{nc}}(A)\right)`$. ∎
Gelfand and Smirnov considered in \[GeSm\] a very special case of this situation where $`A=\mathrm{𝕜}\mathrm{}\mathrm{𝕜}`$ ($`n`$ copies) is a semisimple algebra equiped with the natural trace $`\mathrm{Tr}:(x_1\mathrm{}x_n)_ix_i`$.
It is clear that for $`A=\mathrm{𝕜}\mathrm{}\mathrm{𝕜}`$, one has
$$W_{\mathrm{nc}}(A)=\mathrm{𝕜}a_1,\mathrm{},a_n,b_1,\mathrm{},b_n,$$
is the free graded algebra on $`2n`$ generators, $`a_1,\mathrm{},a_n`$ and $`b_1,\mathrm{},b_n`$ where $`\mathrm{deg}a_j=1`$ and $`\mathrm{deg}b_j=2`$.
The differential $`d_W`$ is a super-derivation such that on generators we have $`d_W(a_j)=b_j`$, and $`d_W(b)=0.`$
Set $`𝖱^{\text{}}(W_{\mathrm{nc}}(A))=W_{\mathrm{nc}}(A)/[W_{\mathrm{nc}}(A),W_{\mathrm{nc}}(A)]`$, where $`[W_{\mathrm{nc}}(A),W_{\mathrm{nc}}(A)]`$ is the linear span of all super-commutators of elements of $`W_{\mathrm{nc}}(A)`$, and the grading on the RHS is induced from that on $`W_{\mathrm{nc}}(A)`$. We can view $`𝖱(W_{\mathrm{nc}}(A))`$ as all cyclic words in $`a_j`$’s and $`b_j`$’s (where $`\mathrm{deg}a_j=1`$ and $`\mathrm{deg}b_j=2`$).
Define maps $`\frac{}{a_j}:𝖱(W_{\mathrm{nc}}(A))𝖱(W_{\mathrm{nc}}(A))`$ in the following way. Choose any word $`x_1\mathrm{}x_k`$ in $`W_{\mathrm{nc}}(A)`$. Set $`\delta _{a_j}(x_1\mathrm{}x_k)`$ to be $`x_2\mathrm{}x_k`$ if $`x_1=a_j`$ and zero otherwise. Then $`\frac{}{a_j}(x_1\mathrm{}x_kmod[W_{\mathrm{nc}}(A),W_{\mathrm{nc}}(A)])`$ is defined to be the sum of $`\delta _{a_j}`$ applied to all cyclic permutations of $`x_1\mathrm{}x_n`$. Maps $`\frac{}{b_j}:𝖱(W_{\mathrm{nc}}(A))𝖱(W_{\mathrm{nc}}(A))`$ are defined analogously. Then set
$$\{P,Q\}=\underset{j=1}{\overset{n}{}}\frac{P}{a_j}\frac{Q}{b_j}+(1)^{(\mathrm{deg}P)(\mathrm{deg}Q)}\frac{Q}{a_j}\frac{P}{b_j}.$$
###### Proposition 17.3.2.
$`(𝗂)`$The above map
$$\{,\}:𝖱^p(W_{\mathrm{nc}}(A))\times 𝖱^q(W_{\mathrm{nc}}(A))𝖱^{p+q3}(W_{\mathrm{nc}}(A))$$
is well-defined, and gives a Lie super-algebra structure on $`𝖱(W_{\mathrm{nc}}(A))`$.
$`(\mathrm{𝗂𝗂})`$The above bracket is equal to the one arising from Theorem 17.3.1 in the special case $`A=\mathrm{𝕜}\mathrm{}\mathrm{𝕜}`$. in particular,
$`(\mathrm{𝗂𝗂𝗂})`$The differential $`d`$ on $`W_{\mathrm{nc}}(A)`$ descends to a well-defined Lie (super)algebra super-derivation of $`𝖱(W_{\mathrm{nc}}(A))`$ such that, for all $`P𝖱(W_{\mathrm{nc}}(A))`$, one has
$$dP=\{b_1^2+\mathrm{}+b_n^2,P\}.$$
Assume, for simplicity, that $`n=1`$, so $`W_{\mathrm{nc}}(A)=\mathrm{𝕜}a,b`$, where $`b=da`$. If $`a`$ were giving a connection on a bundle, then we would consider the quantity $`da+a^2=b+a^2`$, which is the curvature of the connection. This motivates the following definition
###### Definition 17.3.3.
For $`W_{\mathrm{nc}}(A)=\mathrm{𝕜}a,b`$ as above, define for each $`k=0,1\mathrm{}`$
$$\mathrm{ch}_k=(a^2+b)^k𝖱(W_{\mathrm{nc}}(A)).$$
Then we have the following lemma.
###### Lemma 17.3.4.
The following identities hold in $`𝖱(W_{\mathrm{nc}}(A))`$.
1. $`d(\mathrm{ch}_k)=0`$.
2. $`\{\mathrm{ch}_k,\mathrm{ch}_l\}=0`$ for all $`k,l\{0\}`$.
By the Poincaré lemma for $`𝖱(W_{\mathrm{nc}}(A))`$, every closed element of $`𝖱(W_{\mathrm{nc}}(A))`$ of degree $`>0`$ is exact. Hence, there is some element $`\mathrm{ch}_k^1𝖱(W_{\mathrm{nc}}(A))`$ such that $`d(\mathrm{ch}_k^1)=\mathrm{ch}_k`$. Indeed, an explicit formula for $`\mathrm{ch}_k^1`$ can be found. To simplify its expression, we introduce the following notation. For any $`x,yW_{\mathrm{nc}}(A)`$ and $`k,l\{0\}`$, we set $`\sigma _{k,l}(x,y)`$ to be (the image under the quotient of) the sum of all noncommutative words in elements of $`W_{\mathrm{nc}}(A)`$ that contains exactly $`k`$ symbols $`x`$ and $`l`$ symbols $`y`$. Then we have the following result.
Following \[GeSm\], for each $`k`$, set
$`\mathrm{ch}_k^1:={\displaystyle \frac{a}{(k1)!}}\left[{\displaystyle \frac{b^{k1}}{k}}+{\displaystyle \frac{\sigma _{1,k2}(a^2,b)}{k1}}+{\displaystyle \frac{\sigma _{2,k3}(a^2,b)}{k+2}}+\mathrm{}+{\displaystyle \frac{a^{2(k1)}}{2k1}}\right].`$
###### Proposition 17.3.5.
For any $`k`$, in $`𝖱(W_{\mathrm{nc}}(A))`$ one has $`d(\mathrm{ch}_k^1)=\mathrm{ch}_k`$.
## 18. Chern Character on $`K`$-theory
### 18.1. Infinite matrices.
Fix an associative algebra $`A`$. For each integer $`n1`$ we have the algebra $`\mathrm{Mat}_nA`$ of $`n\times n`$-matrices with entries in $`A`$. The assignment
$$\mathrm{Mat}_nA\mathrm{Mat}_{n+1}A,x\left(\begin{array}{cc}x& 0\\ 0& 0\end{array}\right)$$
gives an algebra imbedding (note that the unit $`1_n\mathrm{Mat}_nA`$ does not go to $`1_{n+1}\mathrm{Mat}_{n+1}A`$). We let $`\mathrm{Mat}_{\mathrm{}}(A):=\underset{n\mathrm{}}{lim}\mathrm{Mat}_nA`$ denote the corresponding direct limit under the “upper left hand corner” inclusions. Thus $`\mathrm{Mat}_{\mathrm{}}(A)`$ is an associative algebra without unit that can be identified with the algebra of infinite matrices with finitely many nonzero entries.
Further, let $`\mathrm{GL}_n(A)\mathrm{Mat}_nA`$ be the group of invertible $`n\times n`$-matrices with entries in $`A`$. The map
$$g\left(\begin{array}{cc}g& 0\\ 0& 1\end{array}\right)$$
gives a group imbedding $`\mathrm{GL}_n(A){}_{}{}^{}_{}^{}\mathrm{GL}_{n+1}(A)`$ (note that this time the map does take the unit into unit). We let $`\mathrm{GL}_{\mathrm{}}(A):=\underset{n\mathrm{}}{lim}\mathrm{GL}_n(A)`$ denote the corresponding direct limit. Thus $`\mathrm{GL}_{\mathrm{}}(A)`$ is a group that can be identified with the group of infinite matrices $`g=g_{ij}`$ such that the matrix $`g\mathrm{Id}=g_{ij}\delta _{ij}`$ has only finitely many nonzero entries (here $`\delta _{ij}`$ denotes the Kronecker delta).
For $`i,j\{1,\mathrm{},n\}`$ and $`aA`$, we let $`E_{ij}(a)`$ denote the elementary $`n\times n`$ matrix with $`a`$ in the $`ij`$-position and zero elsewhere.
Also, for any group $`G`$, let $`[G,G]G`$ denote the (normal) subgroup generated by the elements $`ghg^1h^1,g,hG`$. We will denote $`ghg^1h^1`$ by $`[[g,h]]`$.
###### Lemma 18.1.1.
The group $`[GL_{\mathrm{}}(A),GL_{\mathrm{}}(A)]`$ is generated by matrices of the form $`E_{ij}(a)`$ where $`ij`$.
###### Proof.
First, observe that
$$[[E_{ij}(a),E_k\mathrm{}(b)]]=\{\begin{array}{cc}1,\hfill & \text{if }jk\text{}i\mathrm{}\text{;}\hfill \\ E_i\mathrm{}(ab),\hfill & \text{if }j=k\text{}i\mathrm{}\text{;}\hfill \\ E_{kj}(ba),\hfill & \text{if }jk\text{}i=\mathrm{}\text{.}\hfill \end{array}$$
It is an easy computation to see that
$$\left(\begin{array}{cc}X& 0\\ 0& X^1\end{array}\right)=\left(\begin{array}{cc}1& X\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ X& 1\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
Hence $`\left(\begin{array}{cc}X& 0\\ 0& X^1\end{array}\right)`$ is a product of elementary matrices
$$\left(\begin{array}{cc}[[Y,Z]]& 0\\ 0& 1\end{array}\right)=\left(\begin{array}{cc}Y& 0\\ 0& Y^1\end{array}\right)\left(\begin{array}{cc}Z& 0\\ 0& Z^1\end{array}\right)\left(\begin{array}{cc}(ZY)^1& 0\\ 0& ZY\end{array}\right).$$
Thus, $`[[Y,Z]]`$ is a product of elementary matrices, hence any element in $`[GL_{\mathrm{}}(A),`$ $`GL_{\mathrm{}}(A)]`$ may be written as a product of elementary matrices. ∎
### 18.2. The group $`K^0(A)`$.
Fix an associative algebra $`A`$. Recall that $`K^0(A)`$ is defined to be an abelian group which is a quotient of the free abelian group $``$-generated by the isomorphism classes $`[P]`$ of all finite rank projective (left) $`A`$-modules $`P`$ modulo the subgroup generated by the following relations:
$$[P]+[Q][PQ].$$
In other words, $`K^0(A)`$ is the Grothendieck group of (the exact category of) finite rank projective $`A`$-modules, equipped with a semigroup structure by direct sum.
Each finite rank projective $`A`$-module $`P`$ is a direct summand of a free $`A`$-module $`A^n`$, that is, one has an $`A`$-module direct sum decomposition $`A^n=PQ`$. The projection to $`P`$ along $`Q`$ gives a map $`A^n=PQP{}_{}{}^{}_{}^{}A^n,`$ which is given by an $`n\times n`$-matrix, i.e., by an element $`e\mathrm{Mat}_nA`$. It is clear that $`e`$ is an idempotent, i.e., $`e^2=e`$. Direct sum of modules corresponds to direct sum of idempotents, where
$$ee^{}:=\left(\begin{array}{cc}e& 0\\ 0& e^{}\end{array}\right).$$
Let $`𝒫(A)`$ denote the set of idempotents in $`\mathrm{Mat}_{\mathrm{}}(A)`$. We view any $`n\times n`$-matrix as being imbedded into $`\mathrm{Mat}_{\mathrm{}}(A)`$, thus any idempotent $`e\mathrm{Mat}_nA`$ becomes an idempotent in $`\mathrm{Mat}_{\mathrm{}}(A)`$.
###### Lemma 18.2.1.
Two idempotents $`e,e^{}\mathrm{Mat}_{\mathrm{}}(A)`$ give rise to isomorphic projective $`A`$-modules if and only if $`e^{}=geg^1`$ for some $`g\mathrm{GL}_{\mathrm{}}(A).`$
###### Proof.
Let $`e\mathrm{Mat}_nA`$ and $`e^{}\mathrm{Mat}_mA`$ be two idempotents such that one has an $`A`$-module isomorphism $`A^neA^me^{}`$. Put $`P_1=A^ne,Q_1=A^n(1e),`$ and $`P_2=A^me^{},Q_2=A^m(1e^{}).`$ Thus, $`A^n=P_1Q_1,A^m=P_2Q_2,`$ and we are given an isomorphism $`\phi :P_1\stackrel{_{}}{}P_2.`$ We consider the following chain of isomorphisms
$`A^{n+m}=`$ $`A^nA^m\stackrel{_{}}{}P_1Q_1P_2Q_2\text{}`$
$`P_2Q_1P_1Q_2\text{}P_1Q_1P_2Q_2\stackrel{_{}}{}A^nA^m=A^{n+m}.`$
The composite map $`A^{n+m}\stackrel{_{}}{}A^{n+m}`$ is an isomorphism, hence, it is given by an invertible matrix $`g\mathrm{Mat}_{n+m}A`$. It is clear that $`g^1(e0_m)g=0_ne^{},`$ and we are done. ∎
We introduce an equivalence relation $`ee^{}`$ on $`𝒫(A)`$ by $`e^{}=geg^1`$ for some $`g\mathrm{GL}_{\mathrm{}}(A)`$, and denote by $`[e]`$ the equivalence class of $`e`$. We define a semigroup structure on the equivalence classes of idempotents by $`[e]+[e^{}]=[ee^{}]`$ and $`[0]=0`$. This way, one can rephrase the definition of $`K^0(A)`$ in terms of idempotents as follows:
$$K^0(A)(𝒫(A)/_{},).$$
### 18.3. Chern class on $`K^0`$ and $`K^1`$.
Recall that for any associative algebra $`A`$ there is a trace map $`\mathrm{tr}:\mathrm{Mat}_nAA/[A,A]`$ given by
$$\mathrm{tr}(a_{ij})=\underset{i=1}{\overset{n}{}}a_{ii}mod[A,A].$$
Since $`\mathrm{tr}(xy)=\mathrm{tr}(yx)`$, we see that if $`e`$ and $`e^{}`$ are equivalent idempotents (suppose $`e^{}=geg^1`$), then
$$\mathrm{tr}(e^{})=\mathrm{tr}(geg^1)=\mathrm{tr}(e).$$
So, $`\mathrm{tr}`$ descends to a well defined map $`\mathrm{tr}:𝒫(A)/_{}A/[A,A]`$. Notice that $`A/[A,A]=\mathrm{DR}^0(A)`$. This map is additive.
###### Proposition 18.3.1.
The assignment $`[e]\mathrm{tr}(e)`$ extends to a group homomorphism:
$$𝖼_0:K^0(A)\mathrm{Ker}\left[\mathrm{DR}^0(A)\mathrm{DR}^1(A)\right],$$
called the *Chern character*.
###### Proof.
Additivity of the map is clear. So, we must show that for all $`[e]K^0(A)`$, $`𝖼_0[e]`$ is a closed form. Indeed, choose some representative $`e\mathrm{Mat}_nA`$. Then $`e^2=e`$. Applying $`d`$ to both sides yields $`ede+(de)e=de`$. This way, one proves
$$ede=de(1e)\text{and}(de)e=(1e)de.$$
Therefore, we calculate
$`\mathrm{tr}(ede)`$ $`=\mathrm{tr}(e^2de)=\mathrm{tr}(e(de)(1e))=\mathrm{tr}((1e)e(de))=0,`$
since $`(1e)e=ee^2=ee=0`$. Similarly, we see that $`\mathrm{tr}((1e)de)=0`$. So,
$$\mathrm{tr}(de)=\mathrm{tr}(ede)+\mathrm{tr}((1e)de)=0.$$
Clearly, $`\mathrm{tr}`$ and $`d`$ commute, so $`𝖼_0([e])`$ is closed. ∎
Next, we define
$$K^1(A):=GL_{\mathrm{}}(A)/[GL_{\mathrm{}}(A),GL_{\mathrm{}}(A)].$$
We are going to construct a Chern character for $`K^1(A)`$.
###### Proposition 18.3.2.
There is a natural group homomorphism
$$𝖼_1:K^1(A)\mathrm{Ker}\left[\mathrm{DR}^1(A)\stackrel{b}{}A\right].$$
###### Proof.
Choose any $`[g]K^1(A)`$. Choose some representative $`g`$ of $`[g]`$, and define
$$𝖼_1[g]=\mathrm{tr}(g^1dg).$$
First, let us check that this is a group homomorphism. Indeed
$`\mathrm{tr}((g_1g_2)^1d(g_1g_2))`$ $`=\mathrm{tr}[(g_1g_2)^1dg_1g_2]+\mathrm{tr}[g_2^1g_1^1g_1dg_2]`$
$`=\mathrm{tr}(g_2^1g_1^1dg_1g_2)+\mathrm{tr}[g_2^1dg_2]`$
$`=\mathrm{tr}[g_1^1dg_1g_2g_2^1]+\mathrm{tr}(g_2^1dg_2)mod[GL_{\mathrm{}}(A),GL_{\mathrm{}}(A)]`$
$`=\mathrm{tr}[g_1^1dg_1]+\mathrm{tr}[g_2^1dg_2],`$
as desired. Again, $`𝖼_1[g]`$ is a $`b`$-cycle, since
$$b\left[\mathrm{tr}(g^1dg)\right]=\mathrm{tr}(g^1g)\mathrm{tr}(gg^1)=0.$$
### 18.4. Chern classes via connections.
Given a finite rank projective (left) module $`M`$ over an associative algebra $`A`$, one can associate to $`M`$ its de Rham characteristic classes $`\mathrm{ch}_k(M)\mathrm{DR}^{2k}(A),k=1,2,\mathrm{},`$ as follows.
Choose a direct sum decomposition $`MN=A^r`$ and let $`e\mathrm{Mat}_rA`$ be the corresponding projector $`A^rM`$. Then, one has a well-defined non-commutative 1-form $`de\mathrm{\Omega }_{\mathrm{nc}}^1(\mathrm{Mat}_rA)`$. For each $`k=1,2,\mathrm{},`$ we consider the differential form
$$e(de)^{2k}:=e\underset{2k\text{factors}}{\underset{}{dede\mathrm{}de}}\mathrm{\Omega }_{\mathrm{nc}}^{2k}(\mathrm{Mat}_rA),$$
(18.4.1)
and the corresponding class in $`\mathrm{DR}^{2k}(\mathrm{Mat}_rA/\mathrm{Mat}_r\mathrm{𝕜})`$. Let $`\mathrm{Tr}\left(e(de)^{2k}\right)`$ be the image of that class under the canonical ‘trace’-isomorphism $`\mathrm{DR}^{2k}(\mathrm{Mat}_rA/\mathrm{Mat}_r\mathrm{𝕜})\stackrel{_{}}{}`$ $`\mathrm{DR}^{2k}(A)`$, cf. Remark 11.4.6.
###### Proposition 18.4.2.
$`(𝗂)`$The class $`\mathrm{Tr}\left(e(de)^{2k}\right)`$ is independent of the choice of presentation of $`M`$ as a direct summand in a free $`A`$-module, hence is intrisically attached to $`M`$. The assignment
$$[M]\mathrm{ch}_k([M):=\frac{1}{k!}\mathrm{Tr}\left(e(de)^{2k}\right)\mathrm{DR}^{2k}(A)$$
gives a group homomorphism $`K^0(A)\mathrm{DR}^{2k}(A)`$.
$`(\mathrm{𝗂𝗂})`$The class $`\mathrm{Tr}\left(e(de)^{2k}\right)\mathrm{DR}^{2k}(A)`$ is closed, i.e., $`d\left[\mathrm{Tr}\left(e(de)^{2k}\right)\right]=0.`$
###### Proof.
It is clear that if $`e^{}=geg^1,g\mathrm{GL}_{\mathrm{}}(A),`$ is another projector then $`\mathrm{Tr}\left(e(de)^{2k}\right)=\mathrm{Tr}\left(e^{}(de^{})^{2k}\right)`$, due to the invariance of the trace. Therefore, Lemma 18.2.1 implies independence of presentation of $`M`$ as a direct summand in a free $`A`$-module.
Further, given two idempotents $`e_1,e_2\mathrm{Mat}_{\mathrm{}}(A),`$ we clearly have $`d(e_1e_2)=(de_1)(de_2)`$, hence, $`(e_1e_2)(d(e_1e_2))^{2k}=e_1(de_1)^{2k}e_2(de_2)^{2k}`$. Therefore, additivity of the trace implies that, for any finite rank projective $`A`$-modules $`P`$ and $`Q`$, one has $`\mathrm{ch}_k([P][Q])=\mathrm{ch}_k([P])+\mathrm{ch}_k([Q])`$. This completes the proof of (i).
Part (ii) may be verified by a direct computation. Instead of doing so, below we will give an alternative, more conceptual, construction of the characteristic classes in terms of connections, and then prove an analogue of Proposition 18.4.2(ii) in that more general framework. ∎
Following A. Connes \[Co\], one introduces
###### Definition 18.4.3.
A connection on a left $`A`$-module $`M`$ is a linear map $`:M\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM`$ such that
$$(am)=a(m)+dam,aA,mM.$$
###### Lemma 18.4.4.
A left $`A`$-module $`M`$ admits a connection if and only if it is projective.
###### Proof.
Assume $`M`$ is projective and write it as a direct summand of a free $`A`$-module $`AE`$, for some $`\mathrm{𝕜}`$-vector space $`E`$. Thus there is an $`A`$-module imbedding $`i:M{}_{}{}^{}_{}^{}AE`$ and a projection $`AE{}_{}{}^{}_{}^{}M`$ such that $`p{}_{^{^{}}}{}^{}i=\mathrm{Id}_M`$. We define $``$ to be the following composite map
$$M\stackrel{i}{{}_{}{}^{}_{}^{}}AE\stackrel{d\mathrm{Id}_E}{}\mathrm{\Omega }_{\mathrm{nc}}^1(A)E\stackrel{_{}}{}\mathrm{\Omega }_{\mathrm{nc}}^1(A)_A(AE)\stackrel{\mathrm{Id}_\mathrm{\Omega }p}{}\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM.$$
(18.4.5)
It is easy to see that this map gives a connection, sometimes called the Grassmannian connection induced from $`AE`$.
Conversely, let $`M`$ be any left $`A`$-module. Observe that since $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)\overline{A}A`$ is a free right $`A`$-module, the functor $`()_AM`$ takes the fundamental exact sequence $`0\mathrm{\Omega }_{\mathrm{nc}}^1(A)A^\mathrm{e}A0`$ to an exact sequence that looks as follows:
$$0\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM\stackrel{j}{}AM\stackrel{\text{act}}{}M0.$$
(18.4.6)
Here the map $`j`$ takes $`dam`$ to $`am1(am)`$, and the map $`\text{act}:AMM`$ is the action map.
Now, given a connection $`:M\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM`$, we define a map
$$s:MAM,m1mj{}_{^{^{}}}{}^{}(m).$$
For any $`aA`$, using the definition of $`j`$ and Definition 18.4.3, we compute
$`s(am)as(m)`$ $`=1(am)j{}_{^{^{}}}{}^{}(am)a[1mj{}_{^{^{}}}{}^{}(m)]`$
$`=1(am)amj[a(m)+dam]+j[a(m)]`$
$`=j[dam]j[a(m)+dam]+j[a(m)]=0.`$
Hence, $`s`$ is an $`A`$-module map, moreover, one finds
$$\text{act}{}_{^{^{}}}{}^{}s(m)=\text{act}[1mj{}_{^{^{}}}{}^{}(m)]=m(\text{act}{}_{^{^{}}}{}^{}j){}_{^{^{}}}{}^{}(m)=m0=m.$$
We see that the map $`s`$ provides a splitting of the projection act in (18.4.6), therefore $`M`$ is a direct summand in $`AM`$, hence it is projective. ∎
It is easy to see that any connection $`:M\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM`$ has a unique extension to a map $`:\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)_AM\mathrm{\Omega }_{\mathrm{nc}}^{\text{}+1}(A)_AM`$ such that
$$(\alpha \mu )=d\alpha \mu +(1)^{\mathrm{deg}\alpha }\alpha (\mu ),\alpha \mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A),\mu \mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)_AM.$$
(18.4.7)
In particular, one defines the curvature of the connection $``$ as the composite map
$$R_{}:\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)_AM\stackrel{}{}\mathrm{\Omega }_{\mathrm{nc}}^{\text{}+1}(A)_AM\stackrel{}{}\mathrm{\Omega }_{\mathrm{nc}}^{\text{}+2}(A)_AM.$$
(18.4.8)
From formula (18.4.7), we compute
$`R_{}(\alpha \mu )`$ $`={}_{^{^{}}}{}^{}(\alpha \mu )=[d\alpha \mu +(1)^{\mathrm{deg}\alpha }\alpha (\mu )]`$
$`=d^2(\alpha )\mu +(1)^{\mathrm{deg}d\alpha }d\alpha (\mu )+(1)^{\mathrm{deg}\alpha }d\alpha (\mu )`$
$`+(1)^{\mathrm{deg}\alpha +\mathrm{deg}d\alpha }\alpha {}_{^{^{}}}{}^{}(\mu )=\alpha R_{}(\mu ).`$
Thus, the curvature is an $`\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$-linear map.
We can now proceed to the construction of characteristic classes. Fix a finite rank projective left $`A`$-module $`M`$ and choose an imbedding of $`M`$ into a free module $`A^r`$ as a direct summand. Let $`e\mathrm{Mat}_rA`$ be the idempotent that projects $`A^r`$ to $`M`$, and let $`_e:=e{}_{^{^{}}}{}^{}d{}_{^{^{}}}{}^{}e`$ be the corresponding Grassmannian connection on $`M`$. We may view the curvature $`R__e=e{}_{^{^{}}}{}^{}d{}_{^{^{}}}{}^{}e{}_{^{^{}}}{}^{}d{}_{^{^{}}}{}^{}e`$ as an element of $`\mathrm{Mat}_rA\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$. For each $`k=1,2,\mathrm{},`$ we consider the class of the element $`(R__e)^k\mathrm{Mat}_rA\mathrm{\Omega }_{\mathrm{nc}}^{2k}(A)`$ in $`\mathrm{DR}^{2k}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})`$. Thus, applying the canonical isomorphism $`\mathrm{Tr}:\mathrm{DR}^{\text{}}(\mathrm{Mat}_nA/\mathrm{Mat}_n\mathrm{𝕜})\stackrel{_{}}{}\mathrm{DR}^{\text{}}(A)`$ of Proposition 11.4.5 to $`\frac{1}{n!}(R__e)^k`$ we obtain an element
$$\mathrm{ch}_k(M,_e):=\frac{1}{n!}\mathrm{Tr}(R__e)^k\mathrm{DR}^{2k}(A),$$
(18.4.9)
called the $`k`$-th de Rham Chern character class.
## 19. Formally Smooth Algebras
### 19.1.
We are going to study the concept of ‘smoothness’ in noncommutative geometry. Throughout this section $`A`$ denotes a finitely generated associative algebra.
Recall that a two-sided ideal $`I`$ of an associative algebra $`B`$ is said to be nilpotent if there exists $`n>0`$ such that $`b_1\mathrm{}b_n=0,`$ for any $`b_1,\mathrm{},b_nI`$.
###### Definition 19.1.1.
A finitely generated associative algebra $`A`$ is called *formally smooth* if the following lifting property holds. For every algebra $`B`$ and a nilpotent two-sided ideal $`IB`$, given a map $`AB/I`$ there is a lift $`AB`$ such that the following diagram commutes:
(19.1.2)
where $`B{}_{}{}^{}_{}^{}B/I`$ is the quotient map.
To build some intuition for formally smooth algebras we consider commutative case first.
###### Theorem 19.1.3.
For a finitely-generated commutative $`\mathrm{𝕜}`$-algebra the following conditions are equivalent.
1. Let $`m:AAA`$ be the multiplication map. Then $`\mathrm{Ker}m`$ has the locally complete intersection property (this is basically used in the proof of the Hochschild-Kostant-Rosenberg theorem).
2. $`\mathrm{\Omega }_{\mathrm{com}}^1(A)`$ is a projective $`A`$-module.
3. $`A`$ satisfies the lifting property (19.1.2) for any commutative algebra $`B`$. ∎
Recall that $`\mathrm{Rep}_E^A`$ denotes the algebraic variety of all representations of $`A`$ on $`E`$.
###### Proposition 19.1.4.
If $`A`$ is formally smooth and finitely generated, then for every finite-dimensional $`\mathrm{𝕜}`$-vector space $`E`$, the scheme $`\mathrm{Rep}_E^A`$ is smooth.
###### Proof.
Take $`E=\mathrm{𝕜}^n`$, and let $`\mathrm{Rep}_n^A:=\mathrm{Rep}_{\mathrm{𝕜}^n}^A=\mathrm{Rep}_E^A`$. For any scheme $`X`$ and any finitely generated commutative algebra $`B`$, we set
$$X(B)=\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(\mathrm{𝕜}[X],B).$$
Giving an element of $`X(B)`$ is equivalent to giving an algebraic map $`\mathrm{Spec}BX`$. The elements of $`X(B)`$ are called the $`B`$-points of $`X`$. In the case of $`\mathrm{Rep}_n^A`$, we see that for any such $`B`$,
$$\mathrm{Rep}_n^A(B)=\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB).$$
Observe that if $`IB`$ is a nilpotent ideal, then so is $`\mathrm{Mat}_n(I)\mathrm{Mat}_nB`$. Let $`R=\mathrm{𝕜}[\mathrm{Rep}_n^A]`$. We will check that $`R`$ is formally smooth. In other words, we wish to see if the obvious map
$$\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(R,B)\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(R,B/I)$$
is a surjection. By definition, $`\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(R,B)\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB)`$. By the formal smoothness of $`A`$,
$$\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB)\mathrm{Hom}_{\mathrm{𝖺𝗅𝗀}}(A,\mathrm{Mat}_nB/\mathrm{Mat}_n(I))$$
is surjective since $`\mathrm{Mat}_n(I)\mathrm{Mat}_nB`$. The proof then follows. ∎
###### Proposition 19.1.5.
The algebra $`A`$ is formally smooth if and only if $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is a projective $`A`$-bimodule. Equivalently, $`A`$ is formally smooth if and only if the functor $`\mathrm{Der}(A,)`$, on the category of $`A`$-bimodules, is exact.
###### Lemma 19.1.6.
If $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is a projective $`A^\mathrm{e}`$-module, then the categories $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ and $`A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽}`$ both have homological dimension less than or equal to 1, i.e., we have
$$\mathrm{Ext}_{A\text{-}\mathrm{𝗆𝗈𝖽}}^j(M,N)=0\text{resp.}\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^j(M,N)=0,\text{for all}j>1.$$
###### Proof.
We prove the statement for left $`A`$-modules; the proof for $`A`$-bimodules is similar.
Recall the fundumental exact sequence
$$0\mathrm{\Omega }_{\mathrm{nc}}^1(A)A^\mathrm{e}A0.$$
Now, if $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is a projective $`A^\mathrm{e}`$-module, then $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projective as a left and right module. So tensoring the fundamental sequence by any left $`A`$-module $`M`$ preserves exactness, hence yields an exact sequence of left $`A`$-modules
$$0\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AMAMM0.$$
Here, the left $`A`$-module $`AM`$ is projective since it is free. Further, the projectivity of $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ implies that $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projective as a left module. Indeed, if $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is a direct summand of the free bimodule $`AEA`$, then $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)_AM`$ is a direct summand of the free left module $`AEM`$. Thus, we have constructed a length two resolution of $`M`$ by projective $`A`$-modules. Computing the Ext-groups via this resolution we conclude that $`\mathrm{Ext}_{A\text{-}\mathrm{𝗆𝗈𝖽}}^j(M,N)=0`$ for any $`j>1`$. ∎
We observe that if $`I^2=0`$, then $`B=A\mathrm{}I`$, the square zero construction, in the lifting problem if $`A`$ is formally smooth.
###### Lemma 19.1.7.
The following are equivalent.
1. $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projective.
2. Lifting property holds for any square zero extension.
3. $`HH^2(A,M)=0`$ for any $`A`$-bimodule $`M`$.
###### Proof of Lemma 19.1.7..
We will prove that each of the claims (1) and (2) is equivalent to (3).
Recall that the square zero extensions the algebra $`A=B/I`$ by $`I`$ are classified by $`HH^2(B/I,I)`$.
Suppose that $`HH^2(A,M)=0`$ for all $`M`$. Suppose we wish to lift a map $`\alpha :AB/I`$. The pull-back of extension $`IBB/I`$ via $`\alpha `$ gives a commutative diagram:
where the left-hand vertical map is the identity. The algebra $`E`$ is given by the fiber product of $`A`$ and $`B`$, that is, $`E=\{(b,a)BA\alpha (a)=\beta (b)\}`$. Since $`HH^2(A,M)`$, the top row is split by some $`\sigma :AE`$. Then, letting $`EB`$ be denoted by $`\tau `$, the composition $`\tau \sigma `$ is the desired lift.
Now, suppose that the lifting property for square zero extensions holds. We wish to show that $`HH^2(A,M)=0`$. Now, an element of $`HH^2(A,M)`$ gives a square zero extension $`0MEA0`$. Now, since lifting holds for the square zero case, we can lift the identity map $`AA`$ to a map $`AE`$, which splits the extension. So, (2) and $`HH^2(A,M)=0`$ are equivalent.
Finally, an $`A`$-bimodule $`P`$ is projecive if and only if $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(P,M)=0`$ for all $`A`$-bimodules $`M`$. Hence, $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projecive if and only if $`\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(\mathrm{\Omega }_{\mathrm{nc}}^1(A),M)=0`$ for all $`A`$-bimodules $`M`$. But the long exact sequence for $`\mathrm{Ext}`$ arising from the fundamental short exact sequence $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)A^\mathrm{e}A`$ yields
$$\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^1(\mathrm{\Omega }_{\mathrm{nc}}^1(A),M)=\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^2(A,M)HH^2(A,M).$$
Thus, $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projecive if and only if $`HH^2(A,M)=0`$. This completes the proof. ∎
###### Proof of Proposition 19.1.5.
Let $`IB`$ be any nilpotent ideal, $`I^n=0`$, of an algebra $`B`$. We will proceed by induction on $`n`$. Consider the exact sequence
$$0I^{n1}/I^nB/I^nB/I^{n1}0.$$
This is then a square zero extension of $`B/I^{n1}`$ if $`n2`$. Take a map $`AB/I^{n1}`$. Then let $`E`$ be the fiber product of $`A`$ and $`B/I^n`$ to obtain the commutative diagram (with exact rows)
$$\text{}.$$
Then by assumption (and the lemma), there is a splitting $`AEB`$. So, by inducting on $`n`$ until $`I^n=0`$, we obtain lifting.
The implication that there exists a lifting implies $`\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ is projective follows from Lemma 19.1.7. ∎
### 19.2. Examples of formally smooth algebras.
Here are a few examples:
1. The free associative algebra $`\mathrm{𝕜}x_1,\mathrm{},x_n`$.
2. $`\mathrm{Mat}_n\mathrm{𝕜}`$.
3. $`\mathrm{𝕜}[X]`$ where $`X`$ is a smooth affine curve.
4. The path algebra of a quiver.
5. The upper triangular matrices. This is a special case of (4), since the algebra of upper triangular matrices is nothing but the path algebra of the quiver $`\mathrm{}`$.
6. If $`A`$ and $`B`$ are formally smooth, then so are $`AB`$ and $`AB`$.
The reader should be warned that the (commutative) polynomial algebra $`\mathrm{𝕜}[x_1,`$ $`\mathrm{},x_n]`$ is not formally smooth, for any $`n>1`$.
### 19.3. Coherent modules and algebras.
It is perhaps clear from discussion in the previous sections that a formally smooth finitely generated associative (not necessarily commutative) algebra $`A`$ should be viewed as a ‘noncommutative analogue’ of the coordinate ring of a smooth affine algebraic variety $`X`$. Accordingly, the category of finitely generated $`A`$-bimodules should be viewed as a ‘noncommutative analogue’ of the abelian category $`\mathrm{Coh}(X)`$.
An immediate problem that one encounters with such an analogy is that a finitely generated formally smooth algebra $`A`$ is typically not Noetherian, hence, neither the category of finitely generated left $`A`$-modules nor the category of finitely generated $`A`$-bimodules, are abelian categories, in general. This is so, for instance, in the ‘flat’ case where $`A=\mathrm{𝕜}x_1,\mathrm{},x_n,`$ is a free associative algebra on $`n`$ generators. This difficulty can be dealt with by replacing the notion of a finitely generated module by a more restrictive notion of coherent module.
In general, let $`A`$ be an associative algebra. We introduce, see \[Po\].
###### Definition 19.3.1.
A (left) $`A`$-module $`M`$ is called coherent if $`M`$ is finitely generated and, moreover, the kernel of any $`A`$-module map $`A^rM`$ is also a finitely generated $`A`$-module.
It is straightforward to verify that if $`f:MN`$ is a morphism between two coherent modules, then both the kernel and cokernel of $`f`$ are again coherent modules. Thus, coherent $`A`$-modules form an abelian category.
In order for the concept of coherent module to be useful one has to know that, for a given algebra $`A`$, there are “sufficiently many” coherent modules. We will see below that this is indeed the case for formally smooth algebras.
First, we recall that an algebra $`A`$ is called *hereditary* if $`\mathrm{Ext}_{A\text{-}\mathrm{𝗆𝗈𝖽}}^j(M,N)=0`$ for all $`A`$-modules $`M`$ and $`N`$ and $`j2`$, in other words, $`A`$ is hereditary if the category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ has homological dimension less than or equal to 1, c.f. Lemma 19.1.6.
###### Lemma 19.3.2.
An algebra $`A`$ is hereditary if and only if every submodule of a projective $`A`$-module is again projective.
###### Proof.
Let $`P`$ be a projective $`A`$-module, and let $`P^{}P`$ be a submodule. Then we have the short exact sequence
$$0P^{}PP/P^{}0,$$
where the middle term is projective. If $`M`$ is any $`A`$-module, we obtain a long exact sequence in $`\mathrm{Ext}`$ groups:
$$\mathrm{}\mathrm{Ext}^1(P,M)\mathrm{Ext}^1(P^{},M)\mathrm{Ext}^2(P/P^{},M)\mathrm{Ext}^2(P,M)\mathrm{}.$$
Since $`P`$ is projective, $`\mathrm{Ext}^1(P,M)=\mathrm{Ext}^2(P,M)=0`$, hence we find that
$$\mathrm{Ext}^1(P^{},M)\mathrm{Ext}^2(P/P^{},M)$$
for all $`A`$-modules $`M`$. But $`A`$ is hereditary, hence $`\mathrm{Ext}^2(P/P^{},M)=0`$. Therefore, $`\mathrm{Ext}^1(P^{},M)=0`$ for all $`MA\text{-}\mathrm{𝗆𝗈𝖽}`$, hence $`P^{}`$ is projective.
Conversely, let $`M`$ be any $`A`$-module. Then we have the resolution
$$0PA^sM0,$$
where $`P=\mathrm{Ker}(A^sM)`$. Since $`A^s`$ is free, it is projective. Therefore $`P`$ is also projective by assumption. This shows that every $`A`$-module $`M`$ has a projective resolution of length at most two, hence $`\mathrm{Ext}^j(M,N)=0`$ for all $`j2`$. ∎
It is known that $`A`$ is hereditary if and only if all left ideals of $`A`$ are projective.
We observe next that if $`A`$ is formally smooth, then by Lemma 19.1.6 the category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ has homological dimension less than or equal to 1, so $`A`$ is a *hereditary* algebra. Furthermore, the proposition below insures that the category of coherent $`A`$-modules, reps., coherent $`A`$-bimodules, is sufficiently “large”.
###### Lemma 19.3.3.
Let $`A`$ be a finitely generated hereditary algebra. Then, any finite rank free $`A`$-module is coherent.
###### Proof (by D. Boyarchenko).
Let $`M=A^n`$ be a free left $`A`$-module of finite rank $`n`$. Clearly, $`M`$ is finitely generated, so we only have to prove that if $`f:A^rM`$ is a homomorphism, then $`\mathrm{Ker}(f)`$ is also finitely generated. Let $`K=\mathrm{Ker}(f),`$ and $`Q=\mathrm{Im}(f).`$ We have a short exact sequence $`0KA^rQ0.`$ By construction, $`Q`$ is a submodule of the free module $`M.`$ Hence $`Q`$ is projective, since $`A`$ is hereditary. Hence the above exact sequence splits. In particular, this yields a surjection $`A^r{}_{}{}^{}_{}^{}K,`$ which implies that $`K`$ is finitely generated. It follows that any finite rank free left $`A`$-module is coherent. ∎
###### Corollary 19.3.4.
If $`A`$ is a formally smooth algebra, then the cokernel of any $`A`$-module map $`A^mA^n`$ is a coherent $`A`$-module. Also, any finite rank free $`A^\mathrm{e}`$-module, is coherent.
###### Proof.
The first claim follows from the previous Lemma, since a formally smooth algebra is hereditary.
Further, one proves easily that if $`M`$ is a coherent left $`A`$-module and $`N`$ is a coherent left $`B`$-module, then $`MN`$ is a coherent left $`AB`$-module. The claim on $`A^\mathrm{e}`$-modules follows. ∎
### 19.4. Smoothness via torsion-free connection
For any associative algebra $`A`$, we define
$$D(A)=T(A+\overline{A})/(\overline{ab}=a\overline{b}+\overline{a}b,aa^{}=aa^{}1),$$
where as usual $`\overline{A}=A`$ as a vector space. Then $`a\overline{a}`$ is a differential. Earlier, we showed that $`D(A)\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)T_A\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$.
Let us consider the commutative situation. Let $`X`$ be a smooth affine algebraic variety. Then we have the two following relations:
$`\mathrm{\Omega }^{\text{}}(X)=\mathrm{\Lambda }^{\text{}}\mathrm{\Omega }^1(X)\text{and}\mathrm{𝕜}[TX]=\mathrm{Sym}\mathrm{\Omega }^1(X),`$
where we use the convention $`\mathrm{\Omega }^1(X)=\mathrm{\Omega }_{\mathrm{com}}^1(\mathrm{𝕜}[X])`$. In the noncommutative case, there is no difference between the exterior and symmetric powers. Hence we see that the single differential graded algebra $`D(A)=\mathrm{\Omega }_{\mathrm{nc}}^{\text{}}(A)`$ can be simultaneously thought of as noncommutative differential forms and as functions on the “noncommutative tangent bundle.” When we wish to stress the latter interpretation, we will write $`\stackrel{ˇ}{D}(A)`$ instead of $`D(A)`$.
Fix a finite dimensional vector space $`E`$ and consider the variety $`\mathrm{Rep}_E^{\stackrel{ˇ}{D}(A)}`$ of algebra maps $`\stackrel{ˇ}{D}(A)\mathrm{End}_\mathrm{𝕜}E.`$
###### Proposition 19.4.1.
The varieties $`\mathrm{Rep}_E^{\stackrel{ˇ}{D}(A)}`$ and $`T\mathrm{Rep}_E^A`$ are isomorphic.
###### Proof.
An element of $`\mathrm{Rep}_E^{\stackrel{ˇ}{D}(A)}`$ is a homomorphism $`\stackrel{ˇ}{D}(A)\mathrm{End}E`$. But a homomorphism from $`\stackrel{ˇ}{D}(A)`$ can be specified by giving the image of $`a`$ and $`\overline{a}`$ for each $`aA`$. So, let $`\rho (a)`$ be the image of $`a`$, and $`\phi (a)`$ the image of $`\overline{a}`$. Then we can easily check that $`\rho :A\mathrm{End}E`$ must be a homomorphism, while $`\phi :A\mathrm{End}E`$ is a derivation. This is precisely a point of $`T\mathrm{Rep}_E^A`$. Given a point $`T\mathrm{Rep}_E^A`$, we can clearly reverse the arguments made above to construct a homomorphism $`\stackrel{ˇ}{D}(A)\mathrm{End}E`$. ∎
###### Theorem 19.4.2.
An associative algebra $`A`$ is formally smooth if and only if the natural map $`A\overline{A}`$ can be extended to a derivation of $`\stackrel{ˇ}{D}(A)`$ of degree $`+1`$.
Suppose we are in the commutative case, that is, $`A=\mathrm{𝕜}[X]`$ for some affine variety $`X`$. Then we view $`\stackrel{ˇ}{D}(A)`$ as the coordinate ring of the total space of the tangent bundle $`TX`$ of $`X`$. A derivation $``$ of $`\mathrm{𝕜}[TX]`$ is a vector field on $`TX`$. This gives a connection on $`TX`$. Indeed, since $``$ is a derivation, this connection has no torsion (that is, $`_{[\xi ,\eta ]}=_\xi \eta _\eta \xi `$).
###### Proof of Theorem.
We already know that $`A`$ is formally smooth if and only if every square zero extension
$$0MEA0$$
splits. To check that every square zero extension splits, it suffices to check this for the universal square zero extension, $`A\mathrm{}_c\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$, see 10.3.1. A splitting of the latter extension is provided by an algebra map $`\psi :AE`$ such that $`\psi j=\mathrm{id}`$, where $`j:AE`$ is the inclusion map. So, write $`\psi (a)=(a,\varphi (a))`$ for some function $`\varphi :A\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$ (we are not asserting any properties for $`\varphi `$ except linearity, which is obvious). Then we see that
$`\psi (a_1)\psi (a_2)=(a_1a_2,a_1\varphi (a_2)\varphi (a_1)a_2da_1da_2)`$ $`=(a_1a_2,\varphi (a_1a_2))`$
$`=\psi (a_1a_2).`$
So, we find that $`\varphi (a_1a_2)a_1\varphi (a_1)\varphi (a_1)a_2=da_1da_2.`$
Now we can define a derivation of $`\stackrel{ˇ}{D}(A)`$ by $`a\overline{a}`$ and $`\overline{a}\varphi (\overline{a})`$. The proof of the opposite implications follows from the universality of equation (10.3.1).
But the left hand side of the above equation is precisely $`\delta \varphi `$, where $`\delta `$ is the Hochschild differential. We write the right hand side as $`dd`$ to obtain the equation
$$\delta \varphi =dd.$$
Now, we claim that extending the map $`a\overline{a}`$ to a derivation $`+1`$ on $`\stackrel{ˇ}{D}(A)`$ is equivalent to giving a map $`\overline{A}\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$.
###### Lemma 19.4.3.
Giving a map $`\varphi `$ satisfying $`\delta \varphi =dd`$ is equivalent to giving an $`A`$-bimodule splitting of the sequence
$$0\mathrm{\Omega }_{\mathrm{nc}}^1(A)\stackrel{j}{}\mathrm{\Omega }_{\mathrm{nc}}^1(A)A\stackrel{m}{}\mathrm{\Omega }_{\mathrm{nc}}^1(A)\mathrm{\hspace{0.17em}0}$$
where $`m`$ is right multiplication, $`m(\omega a)=\omega a`$, and
$$j(\alpha da)=\alpha a1\alpha a,\text{for all}\alpha \mathrm{\Omega }_{\mathrm{nc}}^1(A),aA.$$
(By splitting, we mean a map $`p:\mathrm{\Omega }_{\mathrm{nc}}^1(A)`$ $`A\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$).
###### Proof.
Giving such a $`\varphi `$ is equivalent to giving a map $`\overline{A}\mathrm{\Omega }_{\mathrm{nc}}^2(A)`$, and we use this along with bilinearity to see that this is equivalent to giving a map
$$p:\mathrm{\Omega }_{\mathrm{nc}}^1(A)A=A\overline{A}A\mathrm{\Omega }_{\mathrm{nc}}^2(A).$$
Since $`\varphi `$ satisfies $`\delta \varphi =dd`$, we can check that $`p`$ is a splitting, i.e., $`pj=\mathrm{id}`$. Conversely, if $`pj=\mathrm{id}`$, we see that
$`pj(da_1da_2)`$ $`=p(da_1a_21da_1a_2)`$
$`=p(d(a_1a_2)1a_2da_21da_1a_2)=da_1da_2,`$
forces $`\delta p=dd`$. ∎
## 20. Serre functors and Duality
### 20.1.
We write $`\mathrm{𝖵𝖾𝖼𝗍}`$ for the category of finite dimensional vector spaces, and $`VV^{}=\mathrm{Hom}_\mathrm{𝕜}(V,\mathrm{𝕜}),`$ for the obvious duality functor on $`\mathrm{𝖵𝖾𝖼𝗍}`$.
In this section, we will freely use the language of derived categories. We write $`[n]`$ for the shift by $`n`$ in a triangulated category. A functor $`F:D_1D_2`$ between triangulated categories is said to be a triangulated functor if it takes distinguished triangles into distinguished triangles, and commutes with the shift functors.
A $`\mathrm{𝕜}`$-linear category $`D`$ is said to be Hom-finite if, for any two objects $`M,ND`$, the space $`\mathrm{Hom}_D(M,N)`$ has finite dimension over $`\mathrm{𝕜}`$.
###### Definition 20.1.1 (Bondal-Kapranov).
Let $`D`$ be a $`\mathrm{Hom}`$-finite triangulated category. An exact functor $`𝖲:DD`$ is called a Serre functor if there are functorial vector space isomorphisms
$$\mathrm{Hom}_D(M,N)^{}\mathrm{Hom}_D(N,𝖲(M)),\text{for any}M,ND.$$
Let $`D`$ be a $`\mathrm{Hom}`$-finite category. For any object $`MD`$, we consider the composite functor $`\mathrm{Hom}_D(M,)^{}:D\mathrm{𝖵𝖾𝖼𝗍},N\mathrm{Hom}_\mathrm{𝕜}(\mathrm{Hom}_D(M,N),\mathrm{𝕜})=\mathrm{Hom}_D(M,N)^{}`$. Observe that, if the category $`D`$ has a Serre functor $`𝖲`$, then the functor $`\mathrm{Hom}_D(M,)^{}`$ is, by definition, represented by the object $`𝖲(M)`$. Yoneda Lemma 2.1.2 insures that the object representing this functor is unique, if exists, up to (essentially unique) isomorphism. Using this, it is not difficult to deduce that any two Serre functors on a $`\mathrm{Hom}`$-finite category must be isomorphic to each other.
The definition of Serre functor is motivated by the following geometric example.
### 20.2. Serre duality.
Let $`X`$ be a smooth projective algebraic variety of dimension $`d`$. Write $`D_{\mathrm{coh}}^b(X)`$ for the bounded derived category of sheaves of $`𝒪_X`$-modules on $`X`$ with coherent cohomology sheaves. As a consequence of completeness of $`X`$, the category $`D_{\mathrm{coh}}^b(X)`$ is $`\mathrm{Hom}`$-finite. This follows from the well-known result, saying that $`dim\left(_iH^i(X,)\right)<\mathrm{}`$, for any coherent sheaf $``$ on $`X`$.
Let $`K_X:=\mathrm{\Omega }_X^d`$ be the canonical line bundle on $`X`$, the line bundle of top-degree differential forms on $`X`$, viewed as an invertible sheaf on $`X`$.
###### Proposition 20.2.1.
The functor $`MMK_X[d]`$ is a Serre functor on $`D_{\mathrm{coh}}^b(X)`$.
###### Remark 20.2.2.
$`(𝗂)`$The condition that $`X`$ is smooth is not very essential here. In the non-smooth case, one has replace $`K_X[d]`$ by the dualising complex $`\kappa __X`$, and to restrict oneself to the category $`D_{\mathrm{perf}}(X)D_{\mathrm{coh}}^b(X)`$ of perfect complexes, that is, the full triangulated subcategory in $`D_{\mathrm{coh}}^b(X)`$ formed by complexes which are quasi-isomorphic to bounded complexes of locally free sheaves. For example, if $`X`$ is a Cohen-Macaulay projective scheme, then the functor $`MM\kappa __X`$ gives a Serre functor on $`D_{\mathrm{perf}}(X)`$. In general, the Serre functor is an equivalence between the category $`D_{\mathrm{perf}}(X)`$ and the category of bounded coherent complexes with finite injective dimension.
$`(\mathrm{𝗂𝗂})`$According to Kontsevich, any $`\mathrm{Hom}`$-finite triangulated category with a Serre functor should be thought of as the category $`D_{\mathrm{coh}}^b(X)`$ for some complete ‘noncommutative space’ $`X`$, possibly singular. $`\mathrm{}`$
Below, we will interchangeably use the words “locally free sheaf” and “vector bundle” and, given such a vector bundle $`E`$, write $`E^{}`$ for the dual vector bundle.
The Proof of Proposition 20.2.1 uses the following important result
###### Theorem 20.2.3 (Grothendieck).
An algebraic variety $`X`$ is smooth if and only if every coherent sheaf on $`X`$ has a finite resolution by locally free sheaves, equivalently, if the (shifts of) vector bundles on $`X`$ generate the category $`D_{\mathrm{coh}}^b(X)`$.∎
###### Proof of Proposition 20.2.1.
The result is essentially a reformulation of the standard Serre duality. The latter says that, for any vector bundle $`E`$ on $`X`$, one has
$$H^i(X,E)^{}H^{di}(X,K_XE^{}),i(\text{Serre duality}).$$
Now, for any two vector bundles $`F_1,F_2,`$ on $`X`$, we have $`\mathrm{Ext}^{\text{}}(F_1,F_2)=H^{\text{}}(X,`$ $`om(F_1,F_2))=H^{\text{}}(X,F_1^{}F_2),`$ where $`\mathrm{Ext}^{\text{}}(,)`$ stands for the Ext-group in the abelian category $`\mathrm{Coh}(X)`$, and $`om(F_1,F_2)\mathrm{Coh}(X)`$ stands for the internal Hom-sheaf which, for vector bundles, is isomorphic to $`F_1^{}F_2`$. Using this, we compute
$`\mathrm{Hom}_{D_{\mathrm{coh}}^b(X)}(F_1,F_2[i])^{}`$ $`=\mathrm{Ext}^i(F_1,F_2)^{}=H^i(X,F_1^{}F_2)^{}`$
(by Serre duality) $`=H^{di}(X,K_X(F_1^{}F_2)^{})`$
$`=H^{di}(X,K_XF_1F_2^{})`$
$`=\mathrm{Ext}^{di}(F_2,K_XF_1)`$
$`=\mathrm{Hom}_{D_{\mathrm{coh}}^b(X)}(F_2[i],F_1K_X[d])`$
$`=\mathrm{Hom}_{D_{\mathrm{coh}}^b(X)}(F_2[i],𝖲(F_1)).`$
Thus, we have checked the defining property of Serre functor in the special case of vector bundles (more precisely, the chain of isomorphisms above may be refined to yield a morphism of functors $`\mathrm{Hom}(,N)^{}\mathrm{Hom}(N,𝖲())`$, which we have shown to be an isomorphism for locally free sheaves). The general case now follows from Proposition 20.2.3. ∎
We keep the above setup, and for any integer $`n1`$, set $`𝖲^n=𝖲{}_{^{^{}}}{}^{}𝖲{}_{^{^{}}}{}^{}\mathrm{}{}_{^{^{}}}{}^{}𝖲`$ ($`n`$ times). Clearly, we have $`𝖲^n:MMK_X^n[dn]`$. We see that any global section $`s\mathrm{\Gamma }(X,K_X^n)`$ gives, for each $`MD_{\mathrm{coh}}^b(X)`$, a morphism $`\mathrm{\Phi }_s:MMK_X^n=𝖲^n(M)[dn],mms`$. Thus, we have a morphism of functors $`\mathrm{\Phi }_s:\mathrm{Id}_{D_{\mathrm{coh}}^b(X)}𝖲^n[dn].`$ This way, we get a linear map of vector spaces
$$\mathrm{\Gamma }(X,K_X^n)\mathrm{Hom}(\mathrm{Id}_{D_{\mathrm{coh}}^b(X)},𝖲^n[dn]),s\mathrm{\Phi }_s,$$
(20.2.4)
which is easily seen to be an isomorphism.
We apply this to prove the following interesting result, first due to Bondal-Orlov.
###### Theorem 20.2.5.
Let $`X`$ and $`Y`$ be smooth projective varieties such that the canonical bundles $`K_X`$ and $`K_Y`$ are both ample line bundles on $`X`$ and $`Y`$, respectively. Then, any trianulated equivalence $`D_{\mathrm{coh}}^b(X)\stackrel{_{}}{}D_{\mathrm{coh}}^b(Y)`$ implies an isomorphism $`XY`$, of algebraic varieties.
###### Remark 20.2.6.
The Theorem says that a smooth projective variety with ample canonical class is completely determined by the corresponding triangulated category $`D_{\mathrm{coh}}^b(X)`$. In particular, with the assumptions above, one has $`D_{\mathrm{coh}}^b(X)D_{\mathrm{coh}}^b(Y)XY.`$ Such an implication is definitely false for varieties with non-ample, e.g. with trivial, canonical bundles.
On the other hand, it is not difficult to show that, for any algebraic varieties $`X`$ and $`Y`$, an equivalence $`\mathrm{Coh}(X)\mathrm{Coh}(Y),`$ of abelian categories, does imply an isomorphism $`XY`$ (Hint: the assignment sending a point $`xX`$ to the sky-scrapper sheaf at $`x`$ sets up a bijection between the set $`X`$ and the set of (isomorphism classes of) simple objects of the category $`\mathrm{Coh}(X)`$. This way, one recovers $`X`$ from $`\mathrm{Coh}(X)`$, as a set. A bit more efforts allow to recover $`X`$ as an algebraic variety, as well.) $`\mathrm{}`$
###### Proof of the Theorem (after Kontsevich).
Put $`D:=D_{\mathrm{coh}}^b(X)`$. We are going to give a canonical procedure of reconstructing the variety $`X`$ from the triangualted category $`D`$.
To this end, observe that, for each $`n,m0`$, the obvious sheaf morphism $`K_X^nK_X^mK_X^{(n+m)}`$ induces a linear map
$$\mathrm{\Gamma }(X,K_X^n)\mathrm{\Gamma }(X,K_X^m)\mathrm{\Gamma }(X,K_X^{(n+m)}).$$
Similarly, for each $`n,m0`$, there is a composition of morphisms of functors defined as follows
$`\mathrm{Hom}(\mathrm{Id}_D,𝖲^n[dn])\mathrm{Hom}(\mathrm{Id}_D,𝖲^m[dm])`$
$`\mathrm{Hom}(\mathrm{Id}_D,𝖲^{n+m}[d(n+m)]).`$
(here, we put $`𝖲^0:=\mathrm{Id}_D`$, by definition). This way, from (20.2.4) we deduce a graded algebra isomorphism
$$_{n0}\mathrm{\Gamma }(X,K_X^n)_{n0}\mathrm{Hom}(\mathrm{Id}_D,𝖲^n[dn]).$$
Now, if $`K_X`$ is very ample, then the variety $`X`$ may be obtained from the graded algebra on the LHS above via the standard $`\mathrm{Proj}`$-construction, that is, we have
$$X=\mathrm{Proj}\left(_{n0}\mathrm{\Gamma }(X,K_X^n)\right)\mathrm{Proj}\left(_{n0}\mathrm{Hom}(\mathrm{Id}_D,𝖲^n[dn])\right).$$
(20.2.7)
If $`K_X`$ is ample but not very ample, we replace $`K_X`$ in this formula by a sufficiently large power of $`K_X`$. Thus, formula (20.2.7) gives way to reconstruct the variety out of the corresponding derived category $`D_{\mathrm{coh}}^b(X)`$.
We observe next that the integer $`d=dimX`$ can be characterized as follows:
$``$ $`d`$ is the unique integer with the property that there exists an object $`MD_{\mathrm{coh}}^b(X)`$ such that $`𝖲(M)M[d]`$.
Now, let $`X,Y`$ be two smooth projective varieties such that $`D_{\mathrm{coh}}^b(X)D_{\mathrm{coh}}^b(Y)`$. Then, by the characterization above, one must have $`dimX=dimY`$. Hence, the integer $`d`$ in the RHS of (20.2.7) is equal to the one in a similar formula for $`Y`$. Further, the uniqueness of the Serre functor mentioned after Definition 20.1.1 implies that the Serre functor on $`D_{\mathrm{coh}}^b(X)`$ goes under the equivalence $`D_{\mathrm{coh}}^b(X)D_{\mathrm{coh}}^b(Y)`$ to the Serre functor on $`D_{\mathrm{coh}}^b(Y)`$. Hence, the equivalence yields a graded algebra isomorphism
$`{\displaystyle _{n0}}\mathrm{Hom}(\mathrm{Id}_{D_{\mathrm{coh}}^b(X)},`$ $`𝖲^{n+m}[d(n+m)])`$
$`{\displaystyle _{n0}}\mathrm{Hom}(\mathrm{Id}_{D_{\mathrm{coh}}^b(Y)},𝖲^{n+m}[d(n+m)]).`$
Therefore, the corresponding Proj-schemes are isomorphic, and we are done. ∎
###### Remark 20.2.8.
A similar result (with similar proof) holds in the case where the varieties have ample anti-canonical classes $`(K_X)^1`$ and $`(K_Y)^1`$. $`\mathrm{}`$
### 20.3. Calabi-Yau categories.
Recall that a smooth variaty $`X`$ is called a Calabi-Yau manifold if it has trivial canonical bundle, $`K_X𝒪_X`$. Assuming in addition that $`X`$ is projective and has dimension $`d`$, the Calabi-Yau property can be reformulated as an isomorphism of functors $`𝖲()()[d]`$. Motivated by this, one introduces the following
###### Definition 20.3.1.
A $`\mathrm{Hom}`$-finite triangulated category $`D`$, with Serre functor $`𝖲`$, is said to be a Calabi-Yau category of dimension $`d`$, if there is an isomorphism of functors $`𝖲()()[d]`$. In such a case, we write $`d=dimD`$.
We observe further that it makes sense to consider Calabi-Yau categories of fractional dimension. Specifically, we say that $`dimD=m/n,`$ provided there is an isomorphism of functors $`𝖲^n()()[m]`$.
###### Example 20.3.2 (Kontsevich).
Let $`A`$ be the associative algebra of upper-triangular $`n\times n`$-matrices (with zero diagonal entries). Then, one can show that the category $`D^b(A\text{-}\mathrm{𝗆𝗈𝖽})`$ has dimension $`\frac{n1}{n+1}`$. This category is, in effect, related to the category of coherent sheaves on the orbifold with $`𝐀_𝐧`$-type isolated singularity, i.e., with singularity of the form $`\mathrm{𝕜}^2/(/n)`$. $`\mathrm{}`$
The following important result is due to \[BK, Lemma 2.7\], cf. also \[BKR\].
###### Theorem 20.3.3.
Let $`D`$ and $`D^{}`$ be two triangulated categories with Serre functors, and let $`F:DD^{}`$ be an exact functor that intertwines the Serre functors on $`D`$ and $`D^{}`$. Assume in addition that
$``$ $`F`$ has a left adjoint $`F^{}:D^{}D`$ and the adjunction morphism $`\mathrm{Id}_DF{}_{^{^{}}}{}^{}F_{}^{}`$ is an isomorphism.
$``$ The category $`D`$ is indecomposable, i.e., there is no nontrivial decomposition $`D=C_1C_2`$, into triangulated subcategories.
$``$ $`D`$ is a Calabi-Yau category.
Then, the functor $`F`$ is an equivalence of triangulated categories.∎
### 20.4. Homological duality.
Given an associative algebra $`A`$, let $`D(A\text{-}\mathrm{𝗆𝗈𝖽})`$ be the derived category of all (unbounded) complexes of left $`A`$-modules.
###### Definition 20.4.1.
Let $`D_{\mathrm{perf}}(A)`$ be the full subcategory in $`D(A\text{-}\mathrm{𝗆𝗈𝖽})`$ formed by the complexes $`C`$, such that $`C`$ is quasi-isomorphic to a bounded complex of finite rank projective $`A`$-modules.
The following Lemma provides an interesting, purely category-theoretic interpretation of the category $`D_{\mathrm{perf}}(A)`$.
###### Lemma 20.4.2.
An object $`MD(A\text{-}\mathrm{𝗆𝗈𝖽})`$ belongs to $`D_{\mathrm{perf}}(A)`$ if and only if $`M`$ is compact, i.e., if the functor $`\mathrm{Hom}_{D(A\text{-}\mathrm{𝗆𝗈𝖽})}(M,)`$ commutes with arbitrary direct sums, cf. Definition 2.1.4.∎
Observe that, for any left $`A`$-module $`M`$, the space $`\mathrm{Hom}_A(M,A)`$ has a natural right $`A`$-module structure induced by right multiplication of $`A`$ on itself. This gives a functor $`\mathrm{Hom}_A(,A):A\text{-}\mathrm{𝗆𝗈𝖽}\mathrm{𝗆𝗈𝖽}\text{-}A=A^{\mathrm{op}}\text{-}\mathrm{𝗆𝗈𝖽}`$. Similarly, one has a functor $`\mathrm{Hom}_{A^{\mathrm{op}}}(,A):A^{\mathrm{op}}\text{-}\mathrm{𝗆𝗈𝖽}A\text{-}\mathrm{𝗆𝗈𝖽}.`$ Observe further that if $`M`$ is a projective left $`A`$-module, then $`\mathrm{Hom}_A(M,A)`$ is a projective right $`A`$-module, and vice versa. Therefore, the functor $`\mathrm{Hom}_A(,A)`$ gives rise to well-defined derived functors $`R\mathrm{Hom}_A(,A):D_{\mathrm{perf}}(A)D_{\mathrm{perf}}(A^{\mathrm{op}})`$.
Below, we will frequently use the following canonical isomorphisms
$`M\stackrel{_{}}{}R\mathrm{Hom}_{A^{\mathrm{op}}}(R\mathrm{Hom}_A(M,A),A),MD_{\mathrm{perf}}(A),`$
$`R\mathrm{Hom}_A(M,A)\stackrel{L}{}_AN\stackrel{_{}}{}R\mathrm{Hom}_A(M,N),M,ND_{\mathrm{perf}}(A).`$ (20.4.3)
$`A\stackrel{L}{}_{A^\mathrm{e}}(LR)R\stackrel{L}{}_AL,\text{for any}LD_{\mathrm{perf}}(A),RD_{\mathrm{perf}}(A^{\mathrm{op}}).`$
Here, each of the isomorphisms is clear for finite rank free modules, hence, holds for finite rank projective modules. This yields the result for arbitrary objects of $`D_{\mathrm{perf}}(A).`$
Recall next that an abelian category $`𝒞`$ is said to have homological dimension $`d`$ if, for any objects $`M,N𝒞`$, we have $`\mathrm{Ext}_𝒞^i(M,N)=0,`$ for all $`i>d`$. For a smooth variety of dimension $`d`$, the category $`\mathrm{Coh}(X)`$ is known to have dimension $`d`$ (and the inequality is in effect an equality).
###### Remark 20.4.4.
We recall that if $`A`$ and $`B`$ are both formally smooth algebras then $`AB`$ is not necessarily formally smooth. Similarly, if $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ and $`B\text{-}\mathrm{𝗆𝗈𝖽}`$ both have finite homological dimension then this is not necessarily so for $`(AB)\text{-}\mathrm{𝗆𝗈𝖽}`$, e.g., take $`A=B=𝕂`$, a field of infinite transcendence degree over $`\mathrm{𝕜}`$ (this example is due to Van den Bergh). $`\mathrm{}`$
Further, assume that the algebra $`A`$ is left Noetherian, and let $`D^b(A\text{-}\mathrm{𝗆𝗈𝖽})`$ be the full subcategory in $`D(A\text{-}\mathrm{𝗆𝗈𝖽})`$ formed by the complexes $`CD(A\text{-}\mathrm{𝗆𝗈𝖽})`$ such that
$``$ Each cohomology group $`H^i(C)`$ is a finitely generated $`A`$-module;
$``$ $`H^i(C)=0`$ for all but finitely many $`i`$.
It is easy to show that $`D^b(A\text{-}\mathrm{𝗆𝗈𝖽})`$ is a triangulated subcategory that contains $`D_{\mathrm{perf}}(A)`$. Furthermore, one proves
###### Lemma 20.4.5.
For a left Noetherian algebra $`A`$, the following conditions are equivalent:
$`(𝗂)`$The category $`A\text{-}\mathrm{𝗆𝗈𝖽}`$ has finite homological dimension;
$`(\mathrm{𝗂𝗂})`$The inclusion $`D_{\mathrm{perf}}(A){}_{}{}^{}_{}^{}D^b(A\text{-}\mathrm{𝗆𝗈𝖽})`$ is an equivalence;
$`(\mathrm{𝗂𝗂𝗂})`$Any finitely-generated left $`A`$-module has a finite resolution by finitely-generated projective $`A`$-modules. ∎
### 20.5. Auslander-Reiten functor.
Below, it will be helpful for us to observe that a left $`A^\mathrm{e}`$-module is the same thing as an $`A`$-bimodule, and also is the same thing as a right $`A^\mathrm{e}`$-module. This may be alternatively explained by the existence of the canonical algebra isomorphism $`(A^\mathrm{e})^{\mathrm{op}}A^\mathrm{e},`$ given by the flip.
Now, the object $`A^\mathrm{e}=AA`$ is clearly both a left and right $`A^\mathrm{e}`$-module. The left $`A^\mathrm{e}`$-action on $`A^\mathrm{e}`$ corresponds to the ‘outer’ $`A`$-bimodule structure on $`AA`$, explicitly given by $`(a^{},a^{\prime \prime }):xy(a^{}x)(ya^{\prime \prime }).`$ We will indicate this ‘outer action’ by writing $`A^\mathrm{e}={}_{_A}{}^{}AA__A.`$ More generally, given an $`A`$-bimodule $`M`$, we will use the notation $`{}_{_A}{}^{}M`$, resp., $`M__A`$, whenever we want to emphasize that $`M`$ is viewed as a left, resp. right, $`A`$-module. With these notations, the right $`A^\mathrm{e}`$-action on $`A^\mathrm{e}`$ corresponds to the ‘inner’ $`A`$-bimodule structure: $`A^\mathrm{e}=A__A{}_{_A}{}^{}A`$, explicitly given by $`(a^{},a^{\prime \prime }):xy(xa^{})(a^{\prime \prime }y).`$
View $`A`$ and $`A^\mathrm{e}`$ as a left $`A^\mathrm{e}`$-modules, and put
$$𝖴:=R\mathrm{Hom}_{D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})}(A,A^\mathrm{e}).$$
The right $`A^\mathrm{e}`$-module structure on $`A^\mathrm{e}`$ induces one on $`R\mathrm{Hom}_{D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})}(A,A^\mathrm{e})`$. This makes $`𝖴`$ a complex of right $`A^\mathrm{e}`$-modules. As we have explained, any right $`A^\mathrm{e}`$-module may be as well viewed as a left $`A^\mathrm{e}`$-module. Thus, we may (and will) regard $`𝖴`$ as an object of $`D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$.
###### Example 20.5.1.
Let $`V`$ be a finite dimensional vector space. The tensor algebra $`A=TV`$ is homologically smooth, and has a standard $`A`$-bimodule resolution:
$$0TVVTV\stackrel{\varkappa }{}TVTV\stackrel{\mathrm{𝚖𝚞𝚕𝚝}}{}TV0,$$
where the map $`\varkappa `$ is given by $`\varkappa :avb(av)ba(vb)`$ (this is a special case of Koszul bimodule-resolution for a general Koszul algebra, cf. e.g. \[BG\],\[VdB1\]). Therefore, we find that, for $`A=TV,`$ the object $`𝖴`$ is represented by the following two-term complex
$$TVTVTVV^{}TV,ab\underset{i=1}{\overset{r}{}}\left[(av_i)\stackrel{ˇ}{v}_iba\stackrel{ˇ}{v}_i(v_ib)\right],$$
where $`\{v_i\}`$ and $`\{\stackrel{ˇ}{v}_i\}`$ are dual bases of $`V`$ and $`V^{}`$, respectively. $`\mathrm{}`$
###### Definition 20.5.2.
The functor
$$D_{\mathrm{perf}}(A)D(A\text{-}\mathrm{𝗆𝗈𝖽}),\text{resp.,}D_{\mathrm{perf}}(A^\mathrm{e})D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽}),$$
given by $`M𝖴\stackrel{L}{}_AM,`$ will be called the Auslander-Reiten functor.
Auslander and Reiten considered a similar functor (on certain abelian categories) in their study of representation theory of finite dimensional algebras.
Next, we extend the notion of Hochschild homology and cohomology to objects of $`D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$ by the formulas
$$HH^i(A,M):=\mathrm{Hom}_{D_{\mathrm{perf}}(A^\mathrm{e})}(A,M[i])\text{and}HH_i(A,M):=H^i(A\stackrel{L}{}_{A^\mathrm{e}}M)$$
where the negative sign is chosen in order to make the present definition compatible with the standard definition of Hochschild homology of a bimodule, as given in §5.
With these definitions we have the following result, see \[VdB2\].
###### Proposition 20.5.3 (Duality).
For any $`MD_{\mathrm{perf}}(A^\mathrm{e})`$, and $`i`$, there is a natural isomorphism $`HH^i(A,M)HH_i(A,𝖴\stackrel{L}{}_AM).`$
###### Proof.
For any $`MD_{\mathrm{perf}}(A^\mathrm{e}),`$ we have $`HH^i(A,M)=\mathrm{Hom}_{D_{\mathrm{perf}}(A^\mathrm{e})}(A,M[i])`$. Therefore, by the second formula in (20.4), we find
$$HH^i(A,M)=H^i\left(\mathrm{Hom}_{D_{\mathrm{perf}}(A^\mathrm{e})}(A,A^\mathrm{e})\stackrel{L}{}_{A^\mathrm{e}}M\right)=H^i(𝖴\stackrel{L}{}_{A^\mathrm{e}}M).$$
Observe that, for any right $`A^\mathrm{e}`$-modules $`R`$ and left $`A^\mathrm{e}`$-module $`L`$, the object $`R\stackrel{L}{}_AL`$ carries an $`A`$-bimodule structure, equivalently, a left $`A^\mathrm{e}`$-module structure; furthermore, an analogue of the third isomorphism in (20.4) says $`R\stackrel{L}{}_{A^\mathrm{e}}LA\stackrel{L}{}_{A^\mathrm{e}}\left(R\stackrel{L}{}_AL\right).`$ Using this formula, and the previous calculation, we find
$`HH^i(A,M)=H^i(𝖴\stackrel{L}{}_{A^\mathrm{e}}M)=H^i(A\stackrel{L}{}_{A^\mathrm{e}}`$ $`(𝖴\stackrel{L}{}_AM))=HH_i(A,𝖴\stackrel{L}{}_AM).`$
We say that an associative (not necessarily commutative) algebra $`A`$ is Gorenstein of dimension $`d`$ if one has
$$\mathrm{Ext}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}^i(A,A^\mathrm{e})\{\begin{array}{cc}A\hfill & \text{if}i=d\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
Clearly, for a Gorenstein algebra $`A`$, in $`D_{\mathrm{perf}}(A^\mathrm{e})`$ one has $`𝖴A[d]`$. Hence, from Proposition 20.5.3 we deduce
###### Corollary 20.5.4.
Given a Gorenstein algebra $`A`$ of dimension $`d`$, for any $`MD_{\mathrm{perf}}(A^\mathrm{e})`$ there are canonical isomorphisms $`HH^{\text{}}(A,M)HH_d\text{}(A,M).`$
Assume next that $`A`$ is a finite-dimensional algebra of finite homological dimension. Then any finite rank projective $`A`$-module is finite dimensional, hence any object of $`D_{\mathrm{perf}}(A)`$, resp., $`D_{\mathrm{perf}}(A^{\mathrm{op}})`$, is quasi-isomorphic to a complex of $`\mathrm{𝕜}`$-finite dimensional $`A`$-modules. Thus, taking the $`\mathrm{𝕜}`$-linear dual of (a complex of) finite dimensional vector spaces induces a functor $`D_{\mathrm{perf}}(A^{\mathrm{op}})D^b(A\text{-}\mathrm{𝗆𝗈𝖽})\stackrel{_{}}{}D_{\mathrm{perf}}(A),M`$ $`M^{}:=\mathrm{Hom}_\mathrm{𝕜}(M,\mathrm{𝕜})`$.
###### Proposition 20.5.5.
Let $`A`$ be a finite-dimensional algebra of finite homological dimension. Then
$`(𝗂)`$The composite functor
$`MR\mathrm{Hom}_A(M,A)^{},`$ is a Serre functor on $`D_{\mathrm{perf}}(A)`$.
$`(\mathrm{𝗂𝗂})`$The functor $`M𝖴\stackrel{L}{}_AM`$ is an inverse to the Serre functor $`𝖲`$; furthermore, one has a functorial isomorphism $`𝖴\stackrel{L}{}_AMR\mathrm{Hom}_A(A^{},M).`$
Our proof of the Proposition given below exploits the following useful isomorphism that holds for any finite dimensional algebra $`A`$:
$$R\mathrm{Hom}_{A^\mathrm{e}}(A,E^{})(A\stackrel{L}{}_{A^\mathrm{e}}E)^{},\text{for any}ED_{\mathrm{perf}}(A^\mathrm{e}).\mathrm{}$$
(20.5.6)
###### Proof of Proposition 20.5.5..
Recall that, for any left $`A`$-modules $`M,N`$, the vector space $`\mathrm{Hom}_\mathrm{𝕜}(M,N)`$ has a natural $`A`$-bimodule structure, and we have a natural isomorphism
$$\mathrm{Hom}_A(M,N)=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,\mathrm{Hom}_\mathrm{𝕜}(M,N)).$$
(20.5.7)
Using this formula, for $`\mathrm{𝕜}`$-finite dimensional left $`A`$-modules $`M,N`$, we compute
$`\mathrm{Hom}_A(N,\mathrm{Hom}_A(M,A)^{})`$ $`=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,N^{}\mathrm{Hom}_A(M,A)^{})`$
$`=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,(\mathrm{Hom}_A(M,A)N)^{}).`$
Therefore, for any finite rank projective $`A`$-modules $`M,N`$, using the previous calculation, we find
$`\mathrm{Hom}_A(N,\mathrm{Hom}_A(M,A)^{})`$ $`=\mathrm{Hom}_{A\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}(A,(N\mathrm{Hom}_A(M,A))^{})`$
$`=\left(A\stackrel{L}{}_{A^\mathrm{e}}\left(N\mathrm{Hom}_A(M,A)\right)\right)^{}`$
by (20.5.6) $`=\left(\mathrm{Hom}_A(M,A)_AN\right)^{}`$
by (20.4) $`=\left(\mathrm{Hom}_A(M,N)\right)^{}.`$
Thus, we have established a natural isomorphism
$$\mathrm{Hom}_A(M,N)^{}\mathrm{Hom}_A(N,\mathrm{Hom}_A(M,A)^{}),$$
for any finite rank projective $`A`$-modules. It follows that a similar isomorphism holds for any objects $`M,ND_{\mathrm{perf}}(A)`$. Thus, we have verified the defining property of Serre functor, and part (i) is proved.
We claim next that the functor $`M𝖴\stackrel{L}{}_AM`$ is a left adjoint of the Serre functor, that is, one has a functorial isomorphism
$$\mathrm{Hom}_{D_{\mathrm{perf}}(A)}(N,𝖲(M))\mathrm{Hom}_{D_{\mathrm{perf}}(A)}(𝖴\stackrel{L}{}_AN,M).$$
(20.5.8)
To establish the isomorphism above, we first use the canonical adjunction isomorphism for tensor products. This says
$$R\mathrm{Hom}_A(𝖴\stackrel{L}{}_AN,M)R\mathrm{Hom}_{A^\mathrm{e}}(𝖴,M__\mathrm{𝕜}N^{})=R\mathrm{Hom}_{A^\mathrm{e}}(𝖴,\mathrm{Hom}_\mathrm{𝕜}(N,M)).$$
Now, using the second formula in (20.4) we compute
$`R\mathrm{Hom}_{A^\mathrm{e}}(𝖴,`$ $`\mathrm{Hom}_\mathrm{𝕜}(N,M))=R\mathrm{Hom}_{A^\mathrm{e}}(𝖴,A^\mathrm{e})\stackrel{L}{}_{A^\mathrm{e}}\mathrm{Hom}_\mathrm{𝕜}(N,M)`$
$`=R\mathrm{Hom}_{A^\mathrm{e}}(R\mathrm{Hom}_{A^\mathrm{e}}(A,A^\mathrm{e}),A^\mathrm{e})\stackrel{L}{}_{A^\mathrm{e}}\mathrm{Hom}_\mathrm{𝕜}(N,M)`$
by (20.4) $`=A\stackrel{L}{}_{A^\mathrm{e}}\mathrm{Hom}_\mathrm{𝕜}(N,M)`$
by (20.5.6) $`=\left(R\mathrm{Hom}_{A^\mathrm{e}}(A,\mathrm{Hom}_\mathrm{𝕜}(N,M))\right)^{}`$
$`=\left(R\mathrm{Hom}_A(M,N)\right)^{}=R\mathrm{Hom}_A(N,𝖲(M)),`$
where the last equality holds by the definition of Serre functor.
Thus we have proved our claim that, for the left adjoint functor $`{}_{}{}^{}𝖲`$, we have $`{}_{}{}^{}𝖲()=𝖴\stackrel{L}{}_A().`$ But the explicit form of the Serre functor provided by part (i) clearly shows that this functor is an equivalence. Hence, its left adjoint functor must be an inverse of $`𝖲`$, and the first statement of part (ii) follows.
To prove the last statement we compute, cf. \[CrB, §2\]:
$$R\mathrm{Hom}_A(A^{},M)R\mathrm{Hom}_A(A^{},A)\stackrel{L}{}_AMR\mathrm{Hom}_{A^\mathrm{e}}(A,AA)\stackrel{L}{}_AM,$$
where the first isomorphism is due to the second formula in (20.4) and the second isomorphism is due to (20.5.7). ∎
### 20.6. Homologically smooth algebras.
Recall that the category of $`A`$-bimodules has a natural monoidal structure $`M,NM_AN`$, where $`M_AN`$ is again an $`A`$-bimodule. This gives, at the level of derived categories, a monoidal structure $`D_{\mathrm{perf}}(A^\mathrm{e})D_{\mathrm{perf}}(A^\mathrm{e})D_{\mathrm{perf}}(A^\mathrm{e}),M,NM\stackrel{L}{}_AN`$, where we identify objects of $`D_{\mathrm{perf}}(A^\mathrm{e})`$ with complexes of $`A`$-bimodules. In a similar way, the category $`D_{\mathrm{perf}}(A)`$ is a module category over the monoidal category $`D_{\mathrm{perf}}(A^\mathrm{e})`$, with the module structure $`D_{\mathrm{perf}}(A^\mathrm{e})D_{\mathrm{perf}}(A)D_{\mathrm{perf}}(A)`$ given by the derived tensor product over $`A`$.
Fix an algebra $`A`$. We consider $`A`$ as an $`A`$-bimodule, that is, as an object of $`D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$.
Following Kontsevich, we introduce
###### Definition 20.6.1.
The algebra $`A`$ is called homologically smooth if $`A`$ is a *compact* object of $`D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$, equivalently, if $`AD_{\mathrm{perf}}(A^\mathrm{e})D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$, that is, if $`A`$ has a finite resolution by finitely-generated projective (left) $`A^\mathrm{e}`$-modules.
This definition is motivated by the following result
###### Lemma 20.6.2.
An affine algebraic variety $`X`$ is smooth if and only if its coordinate ring, $`\mathrm{𝕜}[X]`$, is a homologically smooth algebra (more generally, a scheme $`X`$ is smooth if and only if $`𝒪_{X_\mathrm{\Delta }}`$, the structure sheaf of the diagonal $`X_\mathrm{\Delta }X\times X,`$ is a compact object in $`D(𝒪_{X\times X}\text{-}\mathrm{𝗆𝗈𝖽})`$).
###### Proof.
If $`X`$ is smooth, then so is $`X\times X`$. Hence, by Grothendieck’s theorem 20.2.3, the structure sheaf $`𝒪_{X_\mathrm{\Delta }}`$ has a finite locally-free resolution. Conversely, let $`_d\mathrm{}_1_0{}_{}{}^{}_{}^{}𝒪_{X_\mathrm{\Delta }}`$ be a resolution of $`𝒪_{X_\mathrm{\Delta }}`$ by locally free sheaves on $`X\times X`$. Each term of the resolution, as well as the sheaf $`𝒪_{X_\mathrm{\Delta }}`$ itself, is flat relative to the first projection $`X\times XX`$. Hence, for any point $`xX`$, the resolution restricts to an exact sequence
$$_d|_{\{x\}\times X}\mathrm{}_1|_{\{x\}\times X}_0|_{\{x\}\times X}{}_{}{}^{}_{}^{}\mathrm{𝕜}_x,$$
where $`\mathrm{𝕜}_x`$ denotes the sky-scrapper sheaf at the point $`x`$. This exact sequence provides a bounded resolution of $`\mathrm{𝕜}_x`$ by locally free sheaves on $`X`$. Hence, by the standard regularity criterion, cf. e.g. \[Eis\], $`X`$ is smooth at the point $`x`$. Thus, $`X`$ is smooth at every point, and we are done. ∎
Thus, homologically smooth algebras should be thought of as coordinate rings of smooth ‘noncommutative spaces’. We remark that the property of being ”homologically smooth” is weaker than that of being ”formally smooth”: many algebras, such as polynomial algebras, universal enveloping algebras, etc., are homologically smooth, but not formally smooth. On the other hand, we have
###### Lemma 20.6.3.
$`(𝗂)`$Any formally smooth algebra is homologically smooth.
$`(\mathrm{𝗂𝗂})`$If $`A`$ and $`B`$ are homologically smooth, then so is $`AB`$, and $`A^{\mathrm{op}}`$.
###### Proof.
If $`A`$ is formally smooth, then we have a length two projective resolution $`\mathrm{\Omega }_{\mathrm{nc}}(A)A^\mathrm{e}{}_{}{}^{}_{}^{}A`$. Hence, $`AD_{\mathrm{perf}}(A^\mathrm{e})`$, and $`A`$ is homologically smooth. Part (ii) is straightforward, and is left to the reader. ∎
###### Lemma 20.6.4.
Let $`A`$ be an associative algebra such that $`A^\mathrm{e}`$ is Noetherian.
$`(𝗂)`$The following conditions on $`A`$ are equivalent
$``$ The algebra $`A`$ is homologically smooth;
$``$ The category $`A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽}`$ has finite homological dimension;
$``$ There exists $`d0`$ such that, for any $`A`$-bimodule $`M`$, one has$`HH^i(A,M)=0,`$ for all $`i>d`$.
$`(\mathrm{𝗂𝗂})`$If $`A`$ is homologically smooth, then the equivalent conditions (i)-(iii) of Lemma 20.4.5 hold for $`A`$.
###### Proof.
See \[VdB2\]. ∎
Assume now that $`A`$ is a homologically smooth algebra, so $`AD_{\mathrm{perf}}(A^\mathrm{e}).`$ Then $`A`$ clearly plays the role of unit for the monoidal structure $`()\stackrel{L}{}_A(),`$ on $`D_{\mathrm{perf}}(A^\mathrm{e}).`$ An object $`RD_{\mathrm{perf}}(A^\mathrm{e})`$ is said to be invertible if there exists another object $`R^{}D_{\mathrm{perf}}(A^\mathrm{e})`$, such that in $`D_{\mathrm{perf}}(A^\mathrm{e})`$ one has
$$R\stackrel{L}{}_AR^{}AR^{}\stackrel{L}{}_AR.$$
In this case, one calls $`R^{}`$ an inverse of $`R`$, which is uniquely defined up to isomorphism.
It is straightforward to see from the definitions that the object $`𝖴D(A^\mathrm{e}\text{-}\mathrm{𝗆𝗈𝖽})`$ belongs to $`D_{\mathrm{perf}}(A^\mathrm{e})`$ and, moreover, the derived tensor product with $`𝖴`$ preserves the category $`D_{\mathrm{perf}}(A)`$, that is, gives a functor
$$𝖴\stackrel{L}{}_A():D_{\mathrm{perf}}(A)D_{\mathrm{perf}}(A),M𝖴\stackrel{L}{}_AM.$$
(20.6.5)
###### Question 20.6.6 (Kontsevich).
Is it true that $`𝖴D_{\mathrm{perf}}(A^\mathrm{e})`$ is an invertible object, for any homologically smooth algebra $`A`$ ?
It is not known (to the author) whether the two-term complex representing the object $`𝖴`$ for the free algebra $`A=TV,dimV>1`$, see Example 20.5.1, is an invertible object in $`D_{\mathrm{perf}}\left((TV)^\mathrm{e}\right),`$ i.e., whether Question 20.6.6 has a positive answer for free associative algebras with more than one generator.
Assume that $`A`$ is a homologically smooth algebra such that $`𝖴`$ is invertible, and write $`𝖣D_{\mathrm{perf}}(A^\mathrm{e})`$ for the inverse of $`𝖴`$ (which is well-defined up to isomorphism). It is clear that the functors
$$D_{\mathrm{perf}}(A)D_{\mathrm{perf}}(A),M𝖴\stackrel{L}{}_AM,\text{resp.,}M𝖣\stackrel{L}{}_AM,$$
are mutually inverse auto-equivalences of $`D_{\mathrm{perf}}(A)`$. Furthermore, one can prove the following
$``$ There are algebra quasi-isomorphisms:
$$A^{\mathrm{op}}\stackrel{_{}}{}R\mathrm{Hom}_{D_{\mathrm{perf}}(A)}({}_{_A}{}^{}𝖣,{}_{_A}{}^{}𝖣),\text{and}A\stackrel{_{}}{}R\mathrm{Hom}_{D_{\mathrm{perf}}(A^{\mathrm{op}})}(𝖣__A,𝖣__A).$$
$``$ In $`D_{\mathrm{perf}}(A^\mathrm{e})`$, there is an isomorphism:
$$𝖣R\mathrm{Hom}_{D_{\mathrm{perf}}(A^\mathrm{e})}(A,{}_{_A}{}^{}𝖣𝖣__A).$$
###### Remark 20.6.7.
The above properties show that, if $`𝖴`$ is an invertible object, then its inverse, $`𝖣D_{\mathrm{perf}}(A^\mathrm{e})`$, is the rigid dualizing complex for $`A`$, as defined by Van den Bergh \[VdB1\]. This is known to be the case, for instance, if $`A=\mathrm{𝕜}[X]`$ is the coordinate ring of a smooth variety $`X`$. Then, $`D_{\mathrm{perf}}(\mathrm{𝕜}[X])=D_{\mathrm{coh}}^b(X)`$, and the functor $`𝖣\stackrel{L}{}_A()`$ reduces to $`MM\stackrel{L}{}_A\mathrm{\Omega }_{\mathrm{com}}(A)`$ (this is not a Serre functor because the category $`D_{\mathrm{perf}}(\mathrm{𝕜}[X])`$ is not $`\mathrm{Hom}`$-finite, since $`X`$ is not compact.) $`\mathrm{}`$
## 21. Geometry over an Operad
### 21.1.
We recall that, for any $`/(2)`$-graded vector space $`M=M_{\mathrm{ev}}M_{\mathrm{odd}}`$, one defines the parity reversal operator, $`\mathrm{\Pi }`$, such that $`\mathrm{\Pi }M=M_{\mathrm{odd}}M_{\mathrm{ev}}`$. Thus $`\mathrm{\Pi }M`$ is a $`/(2)`$-graded vector space again.
Let $`𝒫=\{𝒫(n),n=1,2,\mathrm{}\}`$ be a $`\mathrm{𝕜}`$-linear quadratic operad with $`𝒫(1)=\mathrm{𝕜}`$, see \[GiK\]. Let $`𝕊_n`$ denote the Symmetric group on $`n`$ letters. Given $`\mu 𝒫(n)`$ and a $`𝒫`$-algebra $`A`$, we will write: $`\mu _A(a_1,\mathrm{},a_n)`$ for the image of $`\mu a_1\mathrm{}a_n`$ under the structure map: $`𝒫(n)_{_{𝕊_n}}A^nA.`$ Following \[GiK, §1.6.4\], we introduce an enveloping algebra $`𝒰^𝒫A`$, the associative unital $`\mathrm{𝕜}`$-algebra such that the abelian category of (left) $`A`$-modules is equivalent to the category of left modules over $`𝒰^𝒫A`$, see \[GiK, Thm. 1.6.6\]. The algebra $`𝒰^𝒫A`$ is generated by the symbols: $`u(\mu ,a),\mu 𝒫(2),aA,`$ subject to certain relations, see \[Ba, §1.7\].
A $`𝒫`$-algebra in the monoidal category of $`/(2)`$-graded, (resp. $``$-graded) super-vector spaces, see \[GiK, §1.3.17-1.3.18\], will be referred to as a $`𝒫`$-super-algebra, (resp. graded super-algebra). Any $`𝒫`$-algebra may be regarded as a $`𝒫`$-superalgebra concentrated in degree zero. Following \[Gi, §5\], we define a free graded $`𝒫`$-algebra (resp. super-algebra) generated by $`V`$ by
$$𝖳__𝒫^{^{_{\text{}}}}V:=_{i1}𝒫(i)_{_{𝕊_i}}V^i\text{and}\stackrel{ˇ}{𝖳}__𝒫^{^{_{\text{}}}}{}_{}{}^{}V:=_{i1}𝒫(i)_{_{𝕊_i}}(\mathrm{\Pi }V)^i.$$
Fix a $`𝒫`$-algebra $`A`$. Following §3.4, we consider the category $`A\text{-algebras},`$ whose objects are pairs $`(B,f)`$, where $`B`$ is a $`𝒫`$-algebra and $`f:AB`$ is a $`𝒫`$-algebra morphism. Arguing as in §3.4, we get a functor $`A\text{-algebras}A\text{-modules}`$. By \[Gi, Lemma 5.2\], this functor has a right adjoint, i.e., we have
###### Lemma 21.1.1.
Given a $`𝒫`$-algebra $`A`$, there is a functor: $`MT_A^{^{_{\text{}}}}M,`$ (resp. $`M\stackrel{ˇ}{T}_A^{^{_{\text{}}}}M`$) assigning to a left $`A`$-module $`M`$ a graded $`𝒫`$-algebra $`T_A^{^{_{\text{}}}}M=_{i0}T_A^iM`$ (resp. graded $`𝒫`$-superalgebra $`\stackrel{ˇ}{T}_A^{^{_{\text{}}}}M=_{i0}T_A^i(\mathrm{\Pi }M)`$) equipped with a canonical $`𝒫`$-algebra isomorphism $`ı:A\stackrel{_{}}{}T_A^0M`$. Moreover, for any $`𝒫`$-algebra map: $`AB`$, one has a natural adjunction isomorphism:
$$\mathrm{Hom}_{_{A\text{-}\mathrm{𝗆𝗈𝖽𝗎𝗅𝖾𝗌}}}(M,B)\stackrel{_{}}{}\mathrm{Hom}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀𝖾𝖻𝗋𝖺𝗌}}}(T_A^{^{_{\text{}}}}M,B).\mathrm{}$$
An ideal $`I`$ in a $`𝒫`$-algebra $`A`$ will be called $`N`$-nilpotent if, for any $`nN,\mu 𝒫(n)`$, and $`a_1,\mathrm{},a_nA,`$ one has: $`\mu _A(a_1,\mathrm{},a_n)=0,`$ whenever at least $`N`$ among the elements $`a_1,\mathrm{},a_n`$ belong to $`I`$.
Given a left $`A`$-module $`M`$, one defines the square zero extension $`A\mathrm{}M`$, cf. \[Ba, Definition 3.2.6\], or \[Gi, Lemma 5.1\]. The following useful reformulation of the notion of a left $`A`$-module is due to \[Ba, 1.2\], \[Qu\]:
###### Lemma 21.1.2.
Giving a left $`A`$-module structure on a vector space $`M`$ is equivalent to giving a $`𝒫`$-algebra structure on $`A\mathrm{}M:=AM`$ such that the following conditions hold:
$`(𝗂)`$The imbedding: $`aa0`$ makes $`A`$ a $`𝒫`$-subalgebra in $`A\mathrm{}M`$.
$`(\mathrm{𝗂𝗂})`$$`M`$ is a 2-nilpotent ideal in $`A\mathrm{}M`$. ∎
###### Definition 21.1.3.
A $`\mathrm{𝕜}`$-linear map $`\theta :AM`$ is called a derivation if the map: $`ama\theta (a)+m,`$ is an automorphism of the $`𝒫`$-algebra $`A\mathrm{}M`$.
Let $`\mathrm{𝙳𝚎𝚛}__𝒫(A,M)`$ denote the $`\mathrm{𝕜}`$-vector space of all derivations from $`A`$ to $`M`$. It is straightforward to see that the ordinary commutator makes $`\mathrm{𝙳𝚎𝚛}__𝒫(A,A)`$ a Lie algebra.
Next we define, following \[Ba, Definition 4.5.2\], an $`A`$-module of Kähler differentials as the left $`𝒰^𝒫A`$-module, $`\mathrm{\Omega }__𝒫^1A`$, generated by the symbols $`da`$, for $`aA`$, subject to the relations:
$``$$`d(\lambda _1a_1+\lambda _2a_2)=\lambda _1da_1+\lambda _2da_2,\lambda _1,\lambda _2\mathrm{𝕜};`$
$``$$`d(\mu (a_1,a_2))=u(\mu ,a_1)da_2+u(\mu ^{(12)},a_2)da_1,\mu 𝒫(2),a_1,a_2A,`$
where $`u(\mu ,a)`$ denote the standard generators of $`𝒰^𝒫A`$, see \[Ba\].
By construction, $`\mathrm{\Omega }__𝒫^1A`$ is a left $`A`$-module, and the assignment $`ada`$ gives a derivation $`d\mathrm{𝙳𝚎𝚛}__𝒫(A,\mathrm{\Omega }__𝒫^1A)`$. Moreover, this derivation is universal in the following sense. Given any left $`A`$-module $`M`$ and a derivation $`\theta :AM`$, there exists an $`A`$-module morphism $`\mathrm{\Omega }^1\theta :\mathrm{\Omega }__𝒫^1AM`$, uniquely determined by the condition that $`(\mathrm{\Omega }^1\theta )(da)=\theta (a).`$ It follows that the $`A`$-module of Kähler differentials represents the functor $`\mathrm{𝙳𝚎𝚛}__𝒫(A,)`$, i.e., we have (see \[Ba, Remark 4.5.4\]):
###### Lemma 21.1.4.
For any left $`A`$-module $`M`$ there is a natural isomorphism:
$$\mathrm{𝙳𝚎𝚛}__𝒫(A,M)\mathrm{Hom}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^1A,M).\mathrm{}$$
###### Lemma 21.1.5.
There is a natural $`A`$-module morphism $`\delta :\mathrm{\Omega }^1(A\mathrm{}\mathrm{\Omega }__𝒫^1A)`$ $`\mathrm{\Omega }__𝒫^1A_{_{𝒰^𝒫A}}\mathrm{\Omega }__𝒫^1A`$.
###### Lemma 21.1.6.
Let $`A`$ be a $`𝒫`$-algebra, and $`M`$ a left $`A`$-module. Giving the structure of a left $`A\mathrm{}M`$-module on $`N`$ is equivalent to giving a $`𝒫`$-algebra structure on $`AMN`$ such that the following conditions hold:
$`(𝗂)`$The imbedding: $`aa00`$ makes $`A`$ a $`𝒫`$-subalgebra in $`AMN`$.
$`(\mathrm{𝗂𝗂})`$$`M`$ is a 2-nilpotent ideal in $`A\mathrm{}M`$. ∎
Next, we define $`\mathrm{\Omega }__𝒫^{\text{}}A`$, the differential envelope of a $`𝒫`$-algebra $`A`$, as the graded $`𝒫`$-super-algebra: $`\mathrm{\Omega }__𝒫^{\text{}}A=\stackrel{ˇ}{T}_A^{^{_{\text{}}}}(\mathrm{\Omega }__𝒫^1A).`$ The canonical derivation $`d:A\mathrm{\Omega }__𝒫^1A`$ extends, via the Leibniz rule, to a $`𝒫`$-superalgebra derivation $`d:\mathrm{\Omega }__𝒫^{^{_{\text{}}}}A\mathrm{\Omega }__𝒫^{+1}A`$, such that $`d^2=0`$. Thus, $`\mathrm{\Omega }__𝒫^{\text{}}A`$ is a differential graded $`𝒫`$-super-algebra. Further, there is a natural $`𝒫`$-superalgebra imbedding $`j:A=\mathrm{\Omega }__𝒫^0A{}_{}{}^{}_{}^{}\mathrm{\Omega }__𝒫^{\text{}}A`$. Lemmas 21.1.1 and 21.1.4 imply that this imbedding has the following universal property, cf. \[Ko2\] and \[KR\]: given a differential graded $`𝒫`$-superalgebra $`B`$ and a graded $`𝒫`$-algebra morphism $`f:AB`$, there exists a unique morphism $`\mathrm{\Omega }(f):\mathrm{\Omega }__𝒫^{\text{}}AB`$ of differential graded $`𝒫`$-superalgebras, such that: $`f=\mathrm{\Omega }(f){}_{^{^{}}}{}^{}j.`$
###### Proposition 21.1.7 (\[Gi\], Proposition 5.6).
$`(𝗂)`$There is a natural super-differential $`d:\mathrm{\Omega }__𝒫^{\text{}}A\mathrm{\Omega }__𝒫^{+1}A`$, $`d^2=0`$, such that its restriction $`A=\mathrm{\Omega }__𝒫^0A\mathrm{\Omega }__𝒫^1A`$ coincides with the canonical $`A`$-module derivation $`d:A\mathrm{\Omega }__𝒫^1A`$.
$`(\mathrm{𝗂𝗂})`$The differential graded $`𝒫`$-superalgebra $`(\mathrm{\Omega }__𝒫^{\text{}}A,d)`$ is the differential envelope of $`A`$.∎
We set $`\mathrm{\Theta }__𝒫^{\text{}}A=\mathrm{Hom}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^{^{_{\text{}}}}A,A)`$, and call the $`\mathrm{𝕜}`$-vector space $`\mathrm{\Theta }_pA:=`$$`\mathrm{Hom}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^pA,A)`$ the space of $`p`$-polyvectors on $`A`$.
###### Definition 21.1.8.
\[Ba\] A $`\mathrm{𝕜}`$-vector space is called a right $`A`$-module if it is a right $`𝒰^𝒫A`$-module.
Assume from now on that $`𝒫`$ is a Koszul operad, and write $`𝒫^!`$ for the quadratic dual operad, see \[GiK\]. Following \[GiK, §4.2\], \[KV\] and \[Ba, §4\], one defines Hochschild homology, $`H_{\text{}}^𝒫(A,N)`$, of $`A`$ with coefficients<sup>3</sup><sup>3</sup>3in \[GiK, §4.2\] only the case of trivial coefficients has been considered in a right $`A`$-module $`N`$. To this end, write $`\mathrm{𝚜𝚒𝚐𝚗}`$ for the 1-dimensional sign-representation, of $`𝕊_n`$, and given an $`𝕊_n`$-module $`E`$, set $`E^{}:=\mathrm{Hom}_\mathrm{𝕜}(E,\mathrm{𝕜})\mathrm{𝚜𝚒𝚐𝚗}`$. We define $`H_{\text{}}^𝒫(A,N)`$ as the homology groups of the differential graded vector space:
$$C_{\text{}}^𝒫(A,N)=\underset{n1}{}C_n^𝒫(A,N),C_n^𝒫(A,N)=N__\mathrm{𝕜}𝒫^!(n)^{}_{_{𝕊_n}}A^n,$$
(21.1.9)
As was pointed out in \[Ba, Prop. 4.5.3\], for any right $`A`$-module $`N`$, there is a natural isomorphism: $`H_1^𝒫(A,N)N_{𝒰^𝒫A}\mathrm{\Omega }__𝒫^1A`$. In particular, for $`N=𝒰^𝒫A`$, the enveloping algebra of $`A`$ regarded as a right $`A`$-module, one obtains an isomorphism of left $`A`$-modules:
$$\mathrm{\Omega }__𝒫^1AH_1^𝒫(A,𝒰^𝒫A),$$
(21.1.10)
where the left $`A`$-module structure on the RHS is induced by the left $`A`$-module structure on $`𝒰^𝒫A`$.
The right $`A`$-module structure on the enveloping algebra $`𝒰^𝒫A`$ enables us to form the chain complex, see (21.1.9) and \[Ba, §5\]:
$$𝖡_{}^𝒫A=C_{\text{}}^𝒫(A,𝒰^𝒫A)=\underset{n1}{}\left(𝒰^𝒫A__\mathrm{𝕜}𝒫^!(n)^{}_{_{𝕊_n}}A^n\right).$$
(21.1.11)
Further, the left $`A`$-module structure on $`𝒰^𝒫A`$ makes $`𝖡_{}^𝒫A`$ into the following augmented complex of left $`A`$-modules:
$$\mathrm{}𝖡_i^𝒫A𝖡_{i1}^𝒫A\mathrm{}𝖡_2^𝒫A𝖡_1^𝒫A{}_{}{}^{}_{}^{}\mathrm{\Omega }__𝒫^1A,$$
(21.1.12)
Here the augmentation: $`𝖡_1^𝒫A{}_{}{}^{}_{}^{}\mathrm{\Omega }__𝒫^1A`$ is induced by the tautological map: $`𝖡_1^𝒫A=𝒰^𝒫A__\mathrm{𝕜}𝒫^!(1)^{}__\mathrm{𝕜}A\stackrel{_{}}{}𝒰^𝒫AA`$, using the observation that the image of the morphism: $`𝖡_2^𝒫A𝖡_1^𝒫A`$ is the submodule of $`𝖡_1^𝒫A=𝒰^𝒫AA`$ generated by all the elements of the form: $`d(\mu (a_1,a_2))u(\mu ,a_1)da_2u(\mu ^{(12)},a_2)da_1`$, which is exactly the kernel of the projection: $`𝒰^𝒫AA{}_{}{}^{}_{}^{}\mathrm{\Omega }__𝒫^1A`$, see relations of condition (ii) in the definition of $`\mathrm{\Omega }__𝒫^1A`$.
In the associative case the complex $`𝖡_{}^𝒫A`$ is essentially the standard Bar-resolution of the algebra $`A`$ viewed as an $`A`$-bimodule. In particular, for an associative algebra $`A`$, the bimodule $`\mathrm{𝙲𝚘𝚔𝚎𝚛}(𝖡_2^𝒫A𝖡_1^𝒫A)`$ coincides with: $`\mathrm{𝙲𝚘𝚔𝚎𝚛}(A^4A^3),`$ which is equal, due to the exactness of the bar-resolution, to $`\mathrm{𝙺𝚎𝚛}(AAA).`$ For this reason, given an algebra $`A`$ over a general operad $`𝒫`$, we will refer to $`𝖡_{}^𝒫A`$, or to the complex (21.1.12), as the Bar-complex of $`A`$.
Similarly, there is a notion of co-homology, $`H_𝒫^{^{_{\text{}}}}(A,M),`$ with coefficients in a left $`A`$-module $`M`$, see \[KV\], and one has: $`\mathrm{𝙳𝚎𝚛}_{}^{}{}_{_𝒫}{}^{}(A,A)H_𝒫^1(A,A)`$.
###### Proposition 21.1.13.
For any $`𝒫`$-algebra $`A`$ and an $`A`$-module $`M`$ there is a canonical isomorphism:
$$H__𝒫^i(A,M)H^i\left(\mathrm{Hom}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(𝖡_{}^𝒫A,M)\right),i1.$$
###### Proof.
See \[Ba, Proposition 5.2\].∎∎
###### Conjecture 21.1.14.
If $`𝒫`$ is a Koszul operad, then $`𝖡_{}^𝒫A`$ is a resolution of the left $`A`$-module $`\mathrm{\Omega }__𝒫^1A`$; equivalently, and one has: $`H_i(𝖡_{}^𝒫A)=0,i>1.`$
The Conjecture would imply the following result.
###### Corollary 21.1.15.
We have: $`\mathrm{𝚃𝚘𝚛}_i^{^{𝒰^𝒫A}}(𝒰^𝒫A,A)=0,`$ for all $`i>0`$, and also:
$$H_{+1}^𝒫(A,N)\mathrm{𝚃𝚘𝚛}_{}^{^{𝒰^𝒫A}}(N,\mathrm{\Omega }__𝒫^1A),H__𝒫^{+1}(A,M)\mathrm{𝙴𝚡𝚝}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}^{^{_{\text{}}}}(\mathrm{\Omega }__𝒫^1A,M),$$
for any right $`A`$-module $`N`$ and left $`A`$-module $`M`$.∎
###### Conjecture 21.1.16.
There is a natural differential graded space morphism: $`\mathrm{\Omega }__𝒫^{^{_{\text{}}}}AC_{\text{}}^𝒫(A,A).`$
Following Grothendieck and Quillen, see \[CQ\], we introduce
###### Definition 21.1.17.
A $`𝒫`$-algebra $`A`$ is said to be formally-smooth if it satisfies the lifting property with respect to all nilpotent ideals, that is any $`𝒫`$-algebra homomorphism: $`AR/I`$ can be lifted to a $`𝒫`$-algebra homomorphism: $`AR`$, provided $`I`$ is a nilpotent ideal in $`B`$.
###### Proposition 21.1.18.
An algebra $`A`$ over a Koszul operad $`𝒫`$ is formally-smooth if and only if the following two conditions hold, cf. \[Gr\]:
$`(𝗂)`$$`\mathrm{\Omega }__𝒫^1A`$ is a projective left $`A`$-module;
$`(\mathrm{𝗂𝗂})`$For any presentation $`A=R/I`$, where $`R`$ is a free $`𝒫`$-algebra, the morphism $`ȷ`$ below, induced by the canonical map: $`I{}_{}{}^{}_{}^{}R\stackrel{d}{}\mathrm{\Omega }__𝒫^1R`$, is injective, i.e., the following canonical sequence is exact:
$$0I/I^2\stackrel{ȷ}{}𝒰^𝒫A_{_{𝒰^^_𝒫R}}\mathrm{\Omega }__𝒫^1R\mathrm{\Omega }__𝒫^1A\mathrm{\hspace{0.17em}0}.$$
###### Proof.
See \[Gr\].∎
Note that it has been shown in \[CQ\] that condition $`(\mathrm{𝗂𝗂})`$above automatically holds in the associative case (of course, it is not automatic in the associative commutative case, see \[Gr\]).
###### Proposition 21.1.19.
The following 3 properties of a $`𝒫`$-algebra $`A`$ are equivalent:
$`(𝗂)`$The left $`A`$-module $`\mathrm{\Omega }__𝒫^1A`$ is projective.
$`(\mathrm{𝗂𝗂})`$Any $`𝒫`$-algebra extension: $`I{}_{}{}^{}_{}^{}R{}_{}{}^{}_{}^{}A,`$ where $`I`$ is a 2-nilpotent ideal in
$`R`$ has a $`𝒫`$-algebra splitting $`A{}_{}{}^{}_{}^{}R`$, that is: $`RA\mathrm{}I`$.
$`(\mathrm{𝗂𝗂𝗂})`$For any left $`A`$-module $`M`$ we have: $`H__𝒫^2(A,M)=0`$.
###### Proof.
It has been shown in \[Ba, Theorem 3.4.2\] that 2-nilpotent $`𝒫`$-algebra extensions: $`I{}_{}{}^{}_{}^{}R{}_{}{}^{}_{}^{}A`$ are classified by the 2-d cohomology group: $`H__𝒫^2(A,I)`$. This proves that: (ii)$``$(iii). Further, by Corollary 21.1.15 we have: $`H__𝒫^2(A,I)=\mathrm{𝙴𝚡𝚝}_{_{𝖠\text{-}\mathrm{𝗆𝗈𝖽}}}^1(\mathrm{\Omega }__𝒫^1A,M).`$ Thus, (i)$``$(iii). ∎
###### Proposition 21.1.20.
If $`A`$ is a formally-smooth $`𝒫`$-algebra, and $`M`$ a projective $`A`$-module, then the $`𝒫`$-algebra $`T_AM`$ is formally-smooth.
###### Proof.
Copy the proof in \[CQ, Proposition 5.3(3)\]. ∎
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# Constraining Gravitino Dark Matter with the Cosmic Microwave Background
## I Introduction
Supersymmetry provides mainly two compelling candidates for cold dark matter: either the gravitino or the neutralino, depending on which one is the LSP. If we require that R-parity be conserved, the NLSP decays into the stable LSP and releases energy in standard model particles. At leading order, these late decays are two-body and the accompanying energy is mainly electromagnetic.
Part of the electromagnetic release is transferred to the cosmic microwave background radiation. The CMB then re-thermalizes through three relevant processes: Compton scattering $`(\gamma +e\gamma +e)`$, double-Compton scattering $`(\gamma +e\gamma +\gamma +e)`$ and bremsstrahlung $`(e+pe+p+\gamma )`$. The energy injection may distort the CMB, depending on the redshift at which it occurs and on the various time scales of the processes. Early NLSP decays can be fully thermalized, whereas distortions caused by injection from late decays cannot. Varying the NLSP and LSP masses, it is possible to control both the NLSP lifetime and the energy injected in the CMB.
The CMB is not only very isotropic, but it also has a very precise Planck spectrum. The FIRAS instrument aboard the COBE satellite constrains the deviation from a perfect blackbody spectrum in terms of a few numbers. Important for this work is the limit for the chemical potential fixsen:96
$$|\mu |9\times 10^5.$$
This bound has been used to derive limits for the energy released by NLSP decays as well as the NLSP lifetime. Observational constraints on the CMB spectrum can be translated into bounds for the stau NLSP and gravitino LSP masses. It should be mentioned that we explicitly assume that all gravitinos present in the universe are produced by stau decays. Some models suggest that gravitino LSPs may also be produced by scattering interactions after reheating, leading to less stringent constraints. Recently, several papers have employed an analytic approximation to determine these limits feng:0302 ; feng:0306 ; roszkowski:04 . The approximation used in feng:0302 ; feng:0306 ; roszkowski:04 is derived in Ref. hu:93:2 . This analytical result turns out to provide the most stringent limit on the gravitino dark matter model in some range for the NLSP and LSP masses. This prompted us to repeat the calculation numerically.
We find that the bounds for a chemical potential of $`\mu <9\times 10^5`$ given by hu:93:2 is a good approximation only for stau masses above 500 GeV. Below this mass, our bounds are less stringent and even disappear for staus lighter than 100 GeV. We also consider the limit $`\mu <10^5`$ and find that our results suggest lighter gravitinos or equivalently shorter stau lifetimes. Finally, we consider an upper bound on the chemical potential of $`2\times 10^6`$ as planned to be achieved in the DIMES experiment DIMES . We find that, if DIMES does not see a chemical potential, $`\mu <2\times 10^6`$, gravitinos cannot significantly contribute to the dark matter if supersymmetry breaking is gravity mediated.
After recalling some properties of the NLSP in the next section, we write down the full kinetic equation for the evolution of the photon one-particle distribution function in Section III. We express it as an evolution equation for a frequency-dependent chemical potential. This allows us to determine the chemical potential for a given point in the $`(m_{\mathrm{NLSP}},m_{\mathrm{LSP}})`$ and $`(m_{\mathrm{NLSP}},\tau _{\mathrm{NLSP}})`$ plane of the NLSP and to draw exclusion plots in Section IV. We summarize our conclusions in Section V.
## II NLSP properties
We investigate a super-gravity model with a gravitino LSP and a stau NLSP. We assume that the NLSP freezes out with thermal relic density and decays after a time determined by both its mass and the LSP mass. Due to the suppression of non-photonic decay channels, the branching ratio for decays to photons is set to be equal to one.
Gravitino LSPs are produced through NLSP decays $`\mathrm{NLSP}\mathrm{LSP}+\mathrm{SMP}`$, where SMP are standard model particles. Using the standard $`N=1`$ super-gravity Lagrangian, the rates for the various decay channels of the NLSP can be calculated.
### II.1 Slepton NLSP
We assume a gravity-mediated supersymmetry breaking model, where the gravitino LSP mass is of the order $`10^210^4`$ GeV and the stau NLSP lifetime of the order $`10^410^{10}`$ sec. In gauge-mediated supersymmetry breaking models the gravitino is also the LSP but is much lighter $`(m_{\stackrel{~}{G}}`$ keV), resulting in a much shorter NLSP lifetime. In such models, the present CMB constraints do therefore not apply.
The width for the decay of any sfermion $`\stackrel{~}{f}`$ to a gravitino $`\stackrel{~}{G}`$ for a negligible fermion mass is given by
$$\mathrm{\Gamma }(\stackrel{~}{f}f\stackrel{~}{G})=\frac{1}{48\pi M_{}^2}\frac{m_{\stackrel{~}{f}}^5}{m_{\stackrel{~}{G}}^2}\left[1\frac{m_{\stackrel{~}{G}}^2}{m_{\stackrel{~}{f}}^2}\right]^4,$$
(1)
where $`M_{}=(8\pi G_N)^{1/2}`$ is the reduced Planck mass.
Stau NLSPs decay to taus and gravitinos, which then decay to $`e`$, $`\mu `$, $`\pi ^0`$, $`\pi ^\pm `$ and $`\nu `$. As mentioned in Ref. feng:0306 , the electromagnetic energy produced in $`\tau `$ decays varies between $`ϵ_{\mathrm{EM}}^{\mathrm{min}}\frac{1}{3}E_\tau `$ and $`ϵ_{\mathrm{EM}}^{\mathrm{max}}=E_\tau `$. If not specified, a value of $`ϵ_{\mathrm{EM}}=0.8E_\tau `$ will be assumed throughout of this paper.
### II.2 Thermal relic density
Using the thermal relic density of the right-handed slepton NLSPs determined in Ref. scherrer:86 and the thermally-averaged cross section from Ref. asaka:00 , the stau relic abundance is given by
$$\mathrm{\Omega }_{\stackrel{~}{\tau }}^{\mathrm{th}}h^20.2\left[\frac{m_{\stackrel{~}{\tau }}}{\mathrm{TeV}}\right]^2.$$
(2)
As long as the staus do not decay, their time-dependent number density can be expressed as
$$n_{\stackrel{~}{\tau }}(t)0.26\mathrm{m}^3\left[\frac{\mathrm{GeV}}{m_{\stackrel{~}{\tau }}}\right]\left[\frac{T}{\mathrm{Kelvin}}\right]^3\mathrm{\Omega }_{\stackrel{~}{\tau }}.$$
(3)
The fact that the final gravitino density, $`\mathrm{\Omega }_{\stackrel{~}{G}}h^2=(m_{\stackrel{~}{G}}/m_{\stackrel{~}{\tau }})\mathrm{\Omega }_{\stackrel{~}{\tau }}h^2`$, is bounded by observations, $`\mathrm{\Omega }_{\stackrel{~}{G}}h^20.14`$ spergel:03 together with Eq. (2) implies an upper bound for the gravitino mass as a function of the stau mass:
$$\frac{m_{\stackrel{~}{G}}}{\mathrm{GeV}}<0.7\times 10^6\frac{\mathrm{GeV}}{m_{\stackrel{~}{\tau }}}.$$
(4)
## III Evolution equation for the photon number density
The decay of unstable particles into photons during the early stages of the universe can lead to distortions in the CMB. Depending on the redshift at which energy is injected, this may leave a measurable imprint of the early decays. This is the process which we now analyze in detail.
### III.1 Energy injection by particle decays
Energy injection resulting from NLSP decays heats the electrons, leading to a ratio $`T_e/T_{\mathrm{CMB}}`$ larger than one, where $`T_e`$ is the electron temperature and $`T_{\mathrm{CMB}}`$ is the temperature of the CMB. Due to the tight coupling between electrons and photons during the early stages of the universe, the energy surplus of the electrons is redistributed among the photons, distorting the CMB photon distribution from a blackbody spectrum. Assuming that the energy transfer between electrons and photons results in a Bose-Einstein spectrum with a frequency-independent chemical potential, it is possible to relate this resulting chemical potential to the number and energy density of the injected photons and electrons.
Following the analysis done in hu:93:2 , we can write the energy in a Bose-Einstein distribution as
$$\rho _{\mathrm{BE}}=4\sigma _{SB}T_e^4\left(1\frac{90\zeta (3)}{\pi ^4}\mu _{\mathrm{inj}}\right),$$
(5)
where we assume a small chemical potential. The number density is given by
$$n_{\mathrm{BE}}=\frac{2\zeta (3)}{\pi ^2}T_e^3\left(1\frac{\pi ^2}{6\zeta (3)}\mu _{\mathrm{inj}}\right).$$
(6)
Here $`\zeta `$ denotes the Riemann $`\zeta `$-function (see Ref. abramowitz:72 ). Furthermore, due to energy conservation, we know that the energy density may also be written as
$$\rho _{\mathrm{BE}}=\rho _\mathrm{P}+\rho _{\mathrm{decay}},$$
(7)
where $`\rho _{\mathrm{decay}}`$ is the energy density injected by the NLSP decays and $`\rho _P`$ is the density of the CMB photons. This equation is only valid if the injected photons are redistributed in a negligible amount of time compared to the time scales of double-Compton and bremsstrahlung. However, this does not hold for the low-frequency spectrum, where the photon-creating processes dominate.
More precisely, the injected energy density is given by the following differential equation
$$\frac{d\rho _{\mathrm{decay}}}{dt}=ϵ\frac{m_{\mathrm{NLSP}}^2m_{\mathrm{LSP}}^2}{2m_{\mathrm{NLSP}}}n_{\mathrm{NLSP}}(t)\frac{e^{t/\tau }}{\tau }4\frac{\dot{a}}{a}\rho _{\mathrm{decay}},$$
(8)
where $`\tau `$ is the lifetime of the NLSP. Due to the fact that tau decays also produce several neutrinos, the right hand side of Eq. (8) has been multiplied by a factor $`ϵ`$ describing the ratio of the injected energy to the total energy. As pointed out in Sec. II.1, $`ϵ`$ may have a value between 0.3 and 1; we set $`ϵ=0.8`$. In order to solve Eq. (8), we integrate both sides from $`t_{\mathrm{in}}`$ to $`t`$ and obtain
$`\rho _{\mathrm{decay}}`$ $`=`$ $`ϵ{\displaystyle \frac{m_{\mathrm{NLSP}}^2m_{\mathrm{LSP}}^2}{2m_{\mathrm{NLSP}}}}n_{\mathrm{NLSP}}(t_{\mathrm{in}})\left({\displaystyle \frac{a_{\mathrm{in}}}{a}}\right)^4`$ (9)
$`\times \{{\displaystyle \frac{1}{2}}\sqrt{\pi }[\mathrm{Erf}\left(\sqrt{t/\tau }\right)\mathrm{Erf}\left(\sqrt{t_{\mathrm{in}}/\tau }\right)]`$
$`[\sqrt{{\displaystyle \frac{t}{t_{\mathrm{in}}}}}e^{t/\tau }e^{t_{\mathrm{in}}/\tau }]\},`$
where $`\mathrm{Erf}`$ is the error function as defined in abramowitz:72 . Similarly, the photon number density is given by
$$n_{\mathrm{BE}}=n_\mathrm{P}+n_{\mathrm{decay}},$$
(10)
where $`n_{\mathrm{decay}}`$ is the injected photon number density given by the following differential equation
$$\frac{dn_{\mathrm{decay}}}{dt}=N_\gamma n_{\mathrm{NLSP}}(t)\frac{e^{t/\tau }}{\tau }3\frac{\dot{a}}{a}n_{\mathrm{decay}}.$$
(11)
Here $`N_\gamma `$ is the number of photons per stau decay injected into the spectrum. The solution of Eq. (11) is given by
$$n_{\mathrm{decay}}=N_\gamma n_{\mathrm{NLSP}}(t_{\mathrm{in}})\left(\frac{a_{\mathrm{in}}}{a}\right)^3\left(1e^{t/\tau }\right).$$
(12)
Inserting Eq. (9) and Eq. (7) into Eq. (5), as well inserting (12) and Eq. (10) into Eq. (6) we find the relations:
$`1+{\displaystyle \frac{n_{\mathrm{decay}}}{n_\mathrm{P}}}`$ $`=`$ $`\left({\displaystyle \frac{T_e}{T}}\right)^3\left(1{\displaystyle \frac{\pi ^2}{6\zeta (3)}}\mu _{\mathrm{inj}}\right)`$
$`1+{\displaystyle \frac{\rho _{\mathrm{decay}}}{\rho _\mathrm{P}}}`$ $`=`$ $`\left({\displaystyle \frac{T_e}{T}}\right)^4\left(1{\displaystyle \frac{90\zeta (3)}{\pi ^4}}\mu _{\mathrm{inj}}\right).`$ (13)
These equations cannot be solved simultaneously since there are three unknowns, $`\mu _{\mathrm{inj}}`$, $`T_e`$ and $`N_\gamma `$. However, it turns out that the chemical potential is independent of $`N_\gamma `$ up to an unreasonable photon number injection of $`10^7`$ photons per NLSP decay, as shown in the bottom panel of Fig. 1. The top panel shows the time evolution of the chemical potential for two different stau lifetimes.
Up to now, we have assumed that the energy injection caused by stau decays was instantaneously converted in a chemical potential through Compton scattering. However, there are also two other processes that do not conserve the photon numbers: bremsstrahlung and double-Compton scattering. The influence of both processes is discussed in the next section.
We shall use the value $`\mu _{\mathrm{inj}}`$ obtained in this section as initial condition for the numerical solution of the Boltzmann equation discussed below.
### III.2 Photon-matter interaction
When the universe is more than a few minutes old, the coupling of CMB photons and matter is basically due to three processes: Compton scattering, double Compton scattering and bremsstrahlung. If these processes are no longer very efficient, the spectrum can be distorted. Especially, if double Compton scattering and bremsstrahlung become weak, the photon number can no longer be changed and energy injection into the CMB leads to a chemical potential. With this in mind, we parameterize a general distorted spectrum by a Bose-Einstein distribution with a frequency-dependent dimensionless chemical potential $`\mu (x,t)`$,
$$n(x,t)=\frac{1}{\mathrm{exp}(x+\mu (x,t))1},$$
(14)
where $`x=h\nu /T_e`$ is the dimensionless photon frequency. The collision terms in the Boltzmann equation for the three relevant processes (Compton scattering, double-Compton scattering and bremsstrahlung) are studied in lightman:81 . The Boltzmann equation is given by
$`{\displaystyle \frac{n}{t}}=`$ $`{\displaystyle \frac{1}{t_{\gamma e}}}{\displaystyle \frac{T_e}{m_e}}{\displaystyle \frac{1}{x^2}}{\displaystyle \frac{}{x}}\left({\displaystyle \frac{n}{x}}+n+n^2\right)`$ (15)
$`+{\displaystyle \frac{Qg(x)}{t_{\gamma e}}}{\displaystyle \frac{1}{e^xx^3}}\left[1n\left(e^x1\right)\right]`$
$`+\left[1\theta (x1)\right]{\displaystyle \frac{1}{t_{\gamma e}}}{\displaystyle \frac{4\alpha }{3\pi }}\left({\displaystyle \frac{T_e}{m_e}}\right)^2{\displaystyle \frac{1}{x^3}}`$
$`\times \left[1n\left(e^x1\right)\right]{\displaystyle 𝑑xx^4(1+n)n}`$
$`{\displaystyle \frac{e^{x+\mu _{\mathrm{inj}}}}{\left(e^{x+\mu _{\mathrm{inj}}}1\right)^2}}{\displaystyle \frac{d\mu _{\mathrm{inj}}}{dt}}.`$
The first term describes Compton scattering, the second bremsstrahlung, the third term double Compton scattering and the fourth is the injection term given by the solution of Eq. (III.1). We have introduced the Heaviside function $`\theta `$ in the double-Compton scattering term, to take into account that it is active only for $`x<1`$. The constant $`t_{\gamma e}=(n_e\sigma _T)^1`$ is the Thomson scattering time, $`Q=2\sqrt{2\pi }(m_e/T_e)^{1/2}\alpha n_BT_e^31.7\times 10^{10}(\mathrm{MeV}/T)^{1/2}(T/T_e)^{7/2}\mathrm{\Omega }_Bh^2`$ and $`g(x)`$ is the Gaunt factor. More details on the collision terms can be found in Refs. lightman:81 ; hu:93:2 or padmanabhan:02:vol3 .
Since we expect a small value of the chemical potential, we can expand this equation to first order in $`\mu `$.
$`n(x,t)`$ $``$ $`n_0(x,t)+\mu (x,t){\displaystyle \frac{n_0}{\mu }}(x,t)`$ (16)
$`=`$ $`{\displaystyle \frac{1}{e^x1}}\mu (x,t){\displaystyle \frac{e^x}{(e^x1)^2}}.`$
The zeroth order is the equilibrium distribution, and the first order in $`\mu `$ describes the spectral distortion. The kinetic equations for the three relevant processes (Compton scattering, double Compton scattering and bremsstrahlung) lightman:81 , then becomes a linear equation for the evolution of the chemical potential $`(\mu ^{}=\mu /x)`$,
$`{\displaystyle \frac{e^x}{(e^x1)^2}}{\displaystyle \frac{}{t}}\mu (x,t)={\displaystyle \frac{2}{t_{\gamma e}}}{\displaystyle \frac{T_e}{m_e}}{\displaystyle \frac{xe^{2x}}{(e^x1)^4}}\times `$
$`\left[(44\mathrm{cosh}x+x\mathrm{sinh}x)\mu ^{}(x,t)x(\mathrm{cosh}x1)\mu ^{\prime \prime }(x,t)\right]`$
$`+{\displaystyle \frac{Qg(x)}{t_{\gamma e}}}{\displaystyle \frac{1}{x^3(e^x1)}}\mu (x,t)`$
$`+\left[1\theta (x1)\right]{\displaystyle \frac{1}{t_{\gamma e}}}{\displaystyle \frac{16\pi ^3\alpha }{45}}\left({\displaystyle \frac{T_e}{m_e}}\right)^2{\displaystyle \frac{e^x}{x^3(e^x1)}}\mu (x,t)`$
$`{\displaystyle \frac{e^x}{\left(e^x1\right)^2}}{\displaystyle \frac{d\mu _{\mathrm{inj}}}{dt}}.`$
(17)
We have solved both systems of equations numerically and find consistent results. When interested in values $`\mu 10^5`$, we start at $`t_{\mathrm{in}}=10^5`$ sec, but when we want to detect chemical potentials on the level of $`\mu 10^6`$ we have to start at $`t_{\mathrm{in}}=10^4`$ sec.
### III.3 Time evolution of the frequency-dependent chemical potential
As shown in Ref. padmanabhan:02:vol3 , energy injected at a redshift higher than $`z10^7`$, corresponding to a time $`t10^5`$ sec is fully thermalized. Furthermore, at decoupling time $`t_{\mathrm{dec}}10^{13}`$ sec, the CMB spectrum is frozen in and does not evolve anymore apart from redshifting the photon momenta. However, as shown in Fig. 2, the photon-creating processes are unable to reduce a chemical potential already as early as $`t10^8`$ sec, much before recombination.
We see from the top panel of Fig. 2 that, compared to the chemical potential of a Bose-Einstein spectrum (dashed blue line), double-Compton and bremsstrahlung significantly reduce the magnitude of the distortions from a blackbody spectrum: the chemical potential at late times has been reduced from $`7.4\times 10^5`$ to $`4.3\times 10^5`$. The bottom panel of Fig. 2 shows the frequency dependence of $`\mu `$ evaluated at different times. The high-frequency range is dominated by Compton scattering. The chemical potential is constant above $`x1`$, describing a true Bose-Einstein spectrum. The low-energy spectrum is dominated by the photon-creating processes which can destroy the chemical potential and lead to a Planck spectrum below $`x4\times 10^3`$ at recombination time.
It is clear that the later the energy injection, or equivalently the later the staus decay, the weaker are the photon-creating processes which would reduce the distortions. However, it should be kept in mind that our equations are valid only if Compton scattering can achieve a Bose-Einstein spectrum. By requiring that Compton scattering be well active during stau decays, we can put an upper limit on the stau lifetime. Following the analysis of padmanabhan:02:vol3 , a given spectrum can only relax to a Bose-Einstein spectrum before $`t_{\mathrm{BE}}10^9`$ sec. Therefore, the accuracy of a solution of Eq. (15) for a stau lifetime longer than $`\tau _{\mathrm{NLSP}}5\times 10^8`$ sec becomes questionable and untrustworthy for $`\tau _{\mathrm{NLSP}}10^9`$ sec.
On the other hand, after freeze-out of Compton scattering, the injected energy resulting from stau decays cannot be scattered downward in frequency. But there are several other processes that could leave an imprint on the measurable CMB spectrum. For example, photons produced during stau decays have an energy much greater than the electron mass and can create electron-positron pairs through the process $`\gamma +\gamma e^++e^{}`$. The rate of this process has a typical value of $`\mathrm{\Gamma }_{\mathrm{DP}}10^3`$ $`\mathrm{sec}^1`$, provided $`E_\gamma m_e^2/22T`$ kawasaki:95 . Compared to the Hubble rate $`H=1/2t`$, we see that this process plays an important role in heating up the electrons. However, the raise of the electron temperature does no longer affect the CMB spectrum when Compton scattering has already frozen out.
We have also analyzed the fact that the true electron temperature $`T_{e,\mathrm{true}}`$ is not the same as the electron temperature $`T_e`$ obtained by solving Eq. (III.1) due to the influence of bremsstrahlung and double-Compton scattering which reduce the chemical potential from the value $`\mu _{\mathrm{inj}}`$. Given the chemical potential $`\mu `$ at recombination time, we can calculate $`T_e`$ by inserting $`\mu `$ into one of the equations of the system (III.1). We have found that $`T_{e,\mathrm{true}}`$ only differs from $`T_e`$ by $`10^4`$, and that changing the electron temperature in Eq. (15) by such small amounts has no effect. (This is not so surprising, as the effect is of second order.)
## IV CMB Constraints on the stau and gravitino masses
The FIRAS instrument aboard the COBE satellite has measured a temperature $`T_0=2.725\pm 0.001`$ Kelvin mather:99 , and it was able to give an upper bound for the chemical potential fixsen:96 ; smoot:97 , $`|\mu |<9\times 10^5`$. This bound comes from measurements in the frequency range from 2 to 600 GHz, corresponding to $`x=h\nu /T_0[0.03,10]`$. There are also some measurements at lower frequencies, but their accuracy is worse, leading to a lower bound on $`\mu `$ which is by at least an order of magnitude higher. To obtain good accuracy in the measured interval, we numerically compute the chemical potential for $`x[10^4,15]`$. We require that the chemical potential be never higher than $`9\times 10^5`$ within the *experimental* range $`x[0.03,\mathrm{\hspace{0.17em}10}]`$. Outside that range, $`\mu `$ may be larger (experiments do not rule out deviations outside the frequency range $`[0.5\mathrm{GHz},600\mathrm{GHz}]`$).
A point in the $`(m_{\mathrm{NLSP}},m_{\mathrm{LSP}})`$-plane is considered to satisfy the CMB observational bound if the magnitude of the chemical potential never trespasses the limit $`9\times 10^5`$ within and only within the frequency range $`[0.03,\mathrm{\hspace{0.17em}10}]`$. Due to the limitations explained in the previous section, not every point in the $`(m_{\mathrm{NLSP}},m_{\mathrm{LSP}})`$-plane can be calculated, but we expect the chemical potential to be much smaller than the experimental limit for points where our calculation cannot be trusted.
An estimate of the chemical potential caused by an instantaneous energy injection is given in hu:93:2 . It is shown as a solid line in Figs. 3 to 6. This approximation is not very precise for small distortions. While it is in good agreement with our results for staus heavier than 500 GeV, it does not take account of the fact that light staus do not inject enough energy to significantly distort the spectrum. As shown in the upper panel of Fig. 3, the bound on the gravitino mass disappears for staus lighter than $`100`$ GeV. Moreover, due to the freeze-out of Compton scattering, we have introduced a limit on the gravitino masses corresponding to a stau lifetime $`10^9`$ sec. All gravitino masses leading to longer stau lifetimes are allowed.
Our results match quite well with the approximation given in hu:93:2 for $`\mu <10^5`$, as shown in the bottom panel of Fig. 3. We find the same limit for light staus and a somewhat more stringent limit for heavy staus. However, when compared to the stau lifetime (see the bottom panel of Fig.4), our numerical results give a limit on the lifetime that is up to five times shorter than the one obtained by using the approximation.
As mentioned in Sec. III.1, the injected energy depends on the energy going into neutrinos. This is described by the parameter $`ϵ`$. Contrary to what is claimed in feng:0306 , our results depend significantly on this parameter. We considered the two extreme cases $`ϵ=0.3`$ and $`ϵ=1`$ and the usual case $`ϵ=0.8`$ (Fig. 5). The case $`ϵ=0.3`$ is much less constraining than $`ϵ=0.8`$, even completely disappearing for staus lighter than 140 GeV. On the other hand, the cases $`ϵ=1`$ and $`ϵ=0.8`$ match well down to $`m_{\stackrel{~}{\tau }}200`$GeV below which the former becomes more stringent.
Future missions like the Absolute Radiometer for Cosmology, Astrophysics Diffuse Emission (ARCADE) arcade ; kogut:04 or the Diffuse Microwave Emission Survey (DIMES) DIMES experiments may improve sensitivities in the poorly studied centimeter-wavelength band, improving the limit on the chemical potential to about $`|\mu |<2\times 10^6`$. In our model, if neither DIMES nor ARCADE is able to measure distortions of the CMB, gravitinos could only contribute to the missing dark matter if $`\tau _{\mathrm{NLSP}}10^9`$ sec or $`\tau _{\mathrm{NLSP}}2\times 10^5`$ sec (see Fig. 6). However, combining our results with other constraints feng:231 , we find that gravitinos could not significantly contribute to the dark matter for such a bound.
## V Summary
We have studied the effect on the CMB from stau NLSP decays into gravitino LSPs, assuming that the staus freeze out with their thermal relic density. We have numerically solved the kinetic equation for the photon number density with non-instantaneous energy injection. We have found that our numerical results are in good accordance with the analytical approximation hu:93:2 for the induced chemical potential $`\mu <9\times 10^5`$ if the stau is heavier than 500GeV, but differs considerably for lighter stau masses. For light staus the constraints are weaker and even disappear for $`m_{\stackrel{~}{\tau }}<100`$GeV. On the other hand, the approximation underestimates the limits for stronger constraints given by $`\mu <10^5`$ or even more for $`\mu <2\times 10^6`$. This limit, which could be achieved in planned experiments DIMES ; arcade , together with other contraints feng:231 would completely exclude the gravitino as dark matter candidate in models with gravity-mediated supersymmetry breaking. However, allowing a gravitino production after reheating leads to less stringent constraints than our results. We also found that the results depend sensitively on the energy injection parameter $`ϵ`$.
###### Acknowledgements.
We are grateful to Leszek Roszkowski for many helpful remarks and to Jürg Fröhlich for discussions. We also thank Choi Ki-Young for pointing out an important error in the first version of this paper. This work is supported by the Fonds National Suisse.
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# Nonlinear Schrödinger equations with symmetric multi-polar potentials
## 1. Introduction
Schrödinger equations with Hardy-type singular potentials have been the object of a quite large interest in the recent literature, see e.g. . The singularity of inverse square potentials $`V(x)\lambda |x|^2`$ is critical both from the mathematical and the physical point of view. As it does not belong to the Kato’s class, it cannot be regarded as a lower order perturbation of the laplacian but strongly influences the properties of the associated Schrödinger operator. Moreover, from the point of view of nonrelativistic quantum mechanics, among potentials of type $`V(x)\lambda |x|^\alpha `$, the inverse square case represent a transition threshold: for $`\lambda <0`$ and $`\alpha >2`$ (attractively singular potential), the energy is not lower-bounded and a particle near the origin in the presence a potential of this type “falls” to the center, whereas if $`\alpha <2`$ the discrete spectrum has a lower bound (see ). Moreover inverse square singular potentials arise in many fields, such as quantum mechanics, nuclear physics, molecular physics, and quantum cosmology; we refer to for further discussion and motivation.
The case of multi-polar Hardy-type potentials was considered in . In particular in the authors studied the ground states of the following class of nonlinear elliptic equations with a critical power-nonlinearity and a potential exhibiting multiple inverse square singularities:
(1)
$$\{\begin{array}{cc}\mathrm{\Delta }v\underset{i=1}{\overset{k}{}}\frac{\lambda _i}{|xa_i|^2}v=v^{2^{}1},\hfill & \\ v>0\text{in }^N\{a_1,\mathrm{},a_k\}.\hfill & \end{array}$$
For Schrödinger operators $`\mathrm{\Delta }+V`$ the potential term $`V`$ describes the interactions of the quantum particles with the environment. Hence, multi-singular inverse-square potentials are associated with the interaction of particles with a finite number of electric dipoles. The mathematical interest in this problem rests in its criticality, for the exponent of the nonlinearity as well as the singularities share the same order of homogeneity with the laplacian.
The analysis carried out in highlighted how the existence of solutions to (1) heavily depends on the strength and the location of the singularities. For the scaling properties of the problem, the mutual interaction among the poles actually depends only on the shape of their configuration. When the poles form a symmetric structure, it is natural to wonder how the symmetry affects such mutual interaction. The present paper means to study this aspect from the point of view of the existence of solutions inheriting the same symmetry properties as the set of singularities. More precisely we deal with a class of nonlinear elliptic equations on $`^N`$, $`N3`$, involving a critical power-nonlinearity as well as a Hardy-type potential which is singular on sets exhibiting some simple kind of symmetry, as depicted in figures 1–3.
Let us start by considering a potential featuring multiple inverse square singularities located on the vertices of $`k`$-side regular concentric polygons. Let us write $`^N=^2\times ^{N2}`$. For $`k`$, we consider the group $`_k\times 𝕊𝕆(N2)`$ acting on $`𝒟^{1,2}(^N)`$ as $`u(y^{},z)v(y^{},z)=u(e^{2\pi \sqrt{1}/k}y^{},Tz)`$, $`T`$ being any rotation of $`^{N2}`$.
Let us consider $`m`$ regular polygons (with $`k`$ sides) centered at the origin and lying on the plane $`^2\times \{0\}^N`$. Let us denote as $`a_i^{\mathrm{}}`$, $`i=1,2,\mathrm{},k`$, the vertices of the $`\mathrm{}`$-th polygon, $`\mathrm{}=1,2,\mathrm{},m`$. Since the polygons are regular, we have $`a_i^{\mathrm{}}=e^{2\pi \sqrt{1}/k}a_{i1}^{\mathrm{}}`$. We look for $`_k\times 𝕊𝕆(N2)`$-invariant solutions to the following equation
(2)
$$\{\begin{array}{cc}\mathrm{\Delta }v\frac{\lambda _0}{|x|^2}v\underset{\mathrm{}=1}{\overset{m}{}}\underset{i=1}{\overset{k}{}}\frac{\lambda _{\mathrm{}}}{|xa_i^{\mathrm{}}|^2}v=v^{2^{}1},\hfill & \\ v>0\text{in }^N\{0,a_i^{\mathrm{}}:1\mathrm{}m,1ik\},\hfill & \end{array}$$
where $`2^{}=\frac{2N}{N2}`$, i.e. for solutions belonging to the space
$$𝒟_k^{1,2}(^N)=\{u(z,y)𝒟^{1,2}(^N):u(e^{2\pi \sqrt{1}/k}z,y)=u(z,|y|)\},$$
see figures 1–3. Here $`𝒟^{1,2}(^N)`$ denotes the closure of the space $`𝒟(^N)`$ of smooth functions with compact support with respect to the Dirichlet norm
$$u_{𝒟^{1,2}(^N)}:=\left(_^N|u|^2𝑑x\right)^{1/2}.$$
Let us denote as
$$r_{\mathrm{}}=|a_1^{\mathrm{}}|=|a_2^{\mathrm{}}|=\mathrm{}=|a_k^{\mathrm{}}|,\text{for }\mathrm{}=1,2,\mathrm{},m,$$
the radius of the $`\mathrm{}`$-polygon and as
$$\mathrm{\Lambda }_{\mathrm{}}=k\lambda _{\mathrm{}},\mathrm{}=1,\mathrm{},m,$$
the total mass of poles located on the $`\mathrm{}`$-th polygon. The Rayleigh quotient associated with problem (2) in $`𝒟_k^{1,2}(^N)`$ is
(3) $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)=\underset{\begin{array}{c}u𝒟_k^{1,2}(^N)\\ u0\end{array}}{inf}{\displaystyle \frac{{\displaystyle _^N}|u|^2𝑑x{\displaystyle _^N}\left({\displaystyle \frac{\lambda _0}{|x|^2}}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\mathrm{\Lambda }_{\mathrm{}}}{k|xa_i^{\mathrm{}}|^2}}\right)u^2(x)𝑑x}{\left({\displaystyle _^N}|u|^2^{}𝑑x\right)^{2/2^{}}}}.`$
It is well known that minimizers of (3) solve equation (2) up to a Lagrange multiplier. Theorems 1.4 and 1.5 give sufficient conditions for the attainability of $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ for large values of $`k`$. Letting $`k\mathrm{}`$, the Schrödinger operator converges, in the sense of distributions, to the operator associated with a continuous distribution of mass on concentric circles. We stress that the convergence of the potentials does not hold in the natural way, i.e. in $`L_{\mathrm{loc}}^p(^N)`$ for any $`p\frac{N}{2}`$, because of the singularity. To formulate the limiting problem, for any $`r>0`$, we denote as
$$S_r:=\{(x,0)^2\times ^{N2}:|x|=r\}$$
the circle of radius $`r`$ lying on the plane $`^2\times \{0\}^N`$ and consider the distribution $`\delta _{S_r}𝒟^{}(^N)`$ supported in $`S_r`$ and defined by
$$_{𝒟^{}(^N)}\delta _{S_r},\phi _{𝒟(^N)}:=\text{ }_{S_r}\phi (x)d\sigma (x)=\frac{1}{2\pi r}_{S_r}\phi (x)d\sigma (x)\text{for any }\phi 𝒟(^N),$$
where $`d\sigma `$ is the line element on $`S_r`$. We look for solutions to the following equation
(4)
$$\{\begin{array}{cc}\mathrm{\Delta }v\frac{\lambda _0}{|x|^2}v\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}\left(\delta _{S_r_{\mathrm{}}}\frac{1}{|x|^2}\right)v=v^{2^{}1},\hfill & \\ v>0\text{in }^N\{0,_{1\mathrm{}m}S_r_{\mathrm{}}\},\hfill & \end{array}$$
which are invariant by the action of the group $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$. To this purpose, the natural space to set the problem is
$$𝒟_{\mathrm{circ}}^{1,2}(^N)=\{u(z,y)𝒟^{1,2}(^N):u(z,y)=u(|z|,|y|)\},$$
and to consider the associated Rayleigh quotient
(5) $`S_{\mathrm{circ}}`$ $`(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$
$`=\underset{\begin{array}{c}u𝒟_{\mathrm{circ}}^{1,2}(^N)\\ u0\end{array}}{inf}{\displaystyle \frac{{\displaystyle _^N}|u|^2𝑑x{\displaystyle _^N}\left({\displaystyle \frac{\lambda _0}{|y|^2}}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}\text{ }{\displaystyle _{S_r_{\mathrm{}}}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)u^2(y)𝑑y}{\left({\displaystyle _^N}|u|^2^{}𝑑x\right)^{2/2^{}}}}.`$
The following theorem contains a Hardy type inequality for potentials which are singular at circles.
###### Theorem 1.1.
Let $`N3`$ and $`r>0`$. For any $`u𝒟^{1,2}(^N)`$ the map $`yu(y)_{S_r}\frac{d\sigma (x)}{|xy|^2}`$ belongs to $`L^2(^N)`$ and
(6)
$$\left(\frac{N2}{2}\right)^2_^N|u(y)|^2\left(\text{ }_{S_r}\frac{d\sigma (x)}{|xy|^2}\right)𝑑y_^N|u(y)|^2𝑑y.$$
Moreover the constant $`\left(\frac{N2}{2}\right)^2`$ is optimal and not attained.
Hardy type inequalities involving singularities at smooth compact boundaryless manifolds have been considered by several authors, see and references therein. In the aforementioned papers, the potentials taken into account are of the type $`|dist(x,\mathrm{\Sigma })|^2`$, where $`dist(x,\mathrm{\Sigma })`$ denotes the distance from a smooth compact manifold $`\mathrm{\Sigma }`$. We point out that such kind of potentials are quite different from the ones we are considering. Indeed an explicit computation yields
(7)
$$V^r(y):=\text{ }_{S_r}\frac{d\sigma (x)}{|xy|^2}=\frac{1}{\sqrt{(r^2+|y|^2)^24r^2|y^{}|^2}}\text{for all }y=(y^{},z)^N=^2\times ^{N2}.$$
Hence
$`V^r(y){\displaystyle \frac{1}{|y|^2}}\text{as}|y|+\mathrm{}`$
whereas
$`V^r(y)={\displaystyle \frac{1}{\sqrt{r^2+|y|^22r|y^{}|}}}{\displaystyle \frac{1}{\sqrt{r^2+|y|^2+2r|y^{}|}}}{\displaystyle \frac{1}{2r\left||y|r\right|}}\text{as}dist(y,S_r)0.`$
Hence the singularity at the circle of $`V^r`$ is weaker than the inverse square distance potential considered in , but has the same behavior at $`\mathrm{}`$. We also remark that $`V^r`$ is “regular” in the sense of the classification of singular potentials given in .
Arguing as in \[14, Proposition 1.1\], it is easy to verify that solvability of equations (2) and (4) requires the positivity of the associated quadratic forms. Let us consider for example the quadratic form associated with potentials singular on circles, i.e.
$$Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}=_^N|u|^2𝑑x_^N\left(\frac{\lambda _0}{|y|^2}+\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}\text{ }_{S_r_{\mathrm{}}}\frac{d\sigma (x)}{|xy|^2}\right)u^2(y)𝑑y.$$
From (6) and Sobolev’s inequality, it follows that
$$Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}\left[1\frac{4}{(N2)^2}\left(\lambda _0^++_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}}^+\right)\right]_^N|u|^2𝑑x,$$
where $`t^+:=\mathrm{max}\{t,0\}`$ denotes the positive part. Hence $`Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}`$ is positive definite whenever
(8)
$$\lambda _0^++\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}^+<\frac{(N2)^2}{4},$$
see \[14, Proposition 1.2\] for further discussion on the positivity of the quadratic form. Condition (8) also ensures the positivity of the quadratic form associated to problem (2).
The attainability of $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ and $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ requires a delicate balance between the contribution of positive and negative masses. In particular, if $`N4`$ and all the masses have the same sign, $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ and $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ are never achieved.
###### Theorem 1.2.
Let $`N4`$, $`\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m`$, $`r_1,r_2,\mathrm{},r_m^+`$ satisfy (8). If
$`\text{either}(\mathrm{i})`$ $`\mathrm{\Lambda }_{\mathrm{}}<0\text{for all }\mathrm{}=1,\mathrm{},m`$
$`\text{or}(\mathrm{ii})`$ $`\lambda _0>0,\mathrm{\Lambda }_{\mathrm{}}>0\text{for all }\mathrm{}=1,\mathrm{},m`$
then neither $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ nor $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ are attained.
The analysis we are going to carry out in the present paper will highlight that, from the point of view of minimization of the Rayleigh quotient, spreading mass all over a continuum is more convenient than localization of mass at isolated points.
###### Theorem 1.3.
Let $`N4`$, $`\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m`$, $`r_1,r_2,\mathrm{},r_m^+`$ satisfy (8) and
(9) $`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}0,\lambda _0<{\displaystyle \frac{(N2)^2}{4}},`$
and
(10) $`\{\begin{array}{cc}{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\mathrm{\Lambda }_{\mathrm{}}}{|r_{\mathrm{}}|^2}}>0,\hfill & \text{if}\lambda _0{\displaystyle \frac{N(N4)}{4}},\hfill \\ {\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\mathrm{\Lambda }_{\mathrm{}}}{|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}}}>0,\hfill & \text{if}{\displaystyle \frac{N(N4)}{4}}<\lambda _0<{\displaystyle \frac{(N2)^2}{4}}.\hfill \end{array}`$
Then the infimum in (5) is achieved. In particular equation (4) admits a solution which is $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$-invariant.
As problem (5) is the limit of (3), when $`k\mathrm{}`$, we expect the assumptions of Theorem 1.3 to ensure the existence of solutions to (3) provided $`k`$ is sufficiently large. Indeed the theorem below states that (9) and (10) are sufficient conditions on radii and masses of the polygons for the infimum in (3) to be achieved when $`k`$ is large.
###### Theorem 1.4.
Let $`N4`$, $`\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m`$, $`r_1,r_2,\mathrm{},r_m^+`$ satisfy (8), (9) and (10). For any $`\mathrm{}=1,\mathrm{},m`$ and $`k`$ let $`\{a_i^{\mathrm{}}\}_{i=1,\mathrm{},k}`$ be the vertices of a regular $`k`$-side polygon centered at $`0`$ of radius $`r_{\mathrm{}}`$ and let $`\lambda _{\mathrm{}}=\mathrm{\Lambda }_{\mathrm{}}/k`$. Then if $`k`$ is sufficiently large, the infimum in (3) is achieved. In particular equation (2) admits a solution which is $`_k\times 𝕊𝕆(N2)`$-invariant.
When $`N>4`$, it is possible to estimate how large $`k`$ must be in order to obtain the above existence result. This is the content of the following theorem.
###### Theorem 1.5.
Assume that $`N>4`$, $`\lambda _0^++k_{j=1}^m\lambda _j^+<\frac{(N2)^2}{4}`$,
(11) $`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\lambda _{\mathrm{}}0,`$
(12) $`\lambda _1\lambda _2\mathrm{}\lambda _m{\displaystyle \frac{N(N4)}{4}},`$
(13) $`\lambda _0<{\displaystyle \frac{(N2)^2}{4}},`$
(14) $`\{\begin{array}{cc}{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^2}}>0,\hfill & \text{if}\lambda _0{\displaystyle \frac{N(N4)}{4}},\hfill \\ {\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}}}>0,\hfill & \text{if}{\displaystyle \frac{N(N4)}{4}}<\lambda _0<{\displaystyle \frac{(N2)^2}{4}},\hfill \end{array}`$
(15) $`{\displaystyle \frac{\lambda _0}{|r_m|^2}}+\lambda _m{\displaystyle \underset{i=1}{\overset{k1}{}}}{\displaystyle \frac{1}{4r_m^2\left|\mathrm{sin}\frac{i\pi }{k}\right|^2}}+{\displaystyle \underset{\mathrm{}=1}{\overset{m1}{}}}\lambda _{\mathrm{}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{r_m^2+r_{\mathrm{}}^22r_mr_{\mathrm{}}\mathrm{cos}\left(\frac{2\pi i}{k}+\mathrm{\Theta }_j\mathrm{}\right)}}>0,`$
where $`\mathrm{\Theta }_j\mathrm{}`$ denoted the minimum angle formed by vectors $`a_i^j`$ and $`a_s^{\mathrm{}}`$, (see figure 4). Then the infimum in (3) is achieved. In particular equation (2) admits a solution.
This paper is organized as follows. In section 2 we recall some known facts about the single-polar problem and study the behavior of any solution (radial and non radial) to the one pole-equation near the singularities $`0`$ and $`\mathrm{}`$. In section 3 we prove the Hardy type inequality for potentials which are singular at circles stated in Theorem 1.1. Section 4 contains an analysis of possible reasons for lack of compactness of minimizing sequences of problems (3) and (5) and a local Palais-Smale condition below some critical thresholds. In section 5 we provide some interaction estimates which are needed in section 6 to compare the concentration levels of minimization sequences and consequently to prove Theorem 1.3. Section 7 contains the study of behavior of energy levels of minimizing sequences as $`k\mathrm{}`$ which is needed in section 8 to prove Theorem 1.4. Last section is devoted to the proof of Theorem 1.5.
Figures 1–3 are the plot of the test functions used to estimate the energy levels of Palais-Smale sequences and represent the expected shape of solutions found in Theorem 1.4 (assuming the last $`N2`$ variables to be $`0`$), which is based on the knowledge of their behavior at singularities (see and Theorem 2.2), which is known to be be singular at poles with positive mass and vanishing at poles with negative mass.
1. $`k=10`$, $`m=2`$, $`\lambda _0,\lambda _1>0`$, $`\lambda _2<0`$.
Notation. We list below some notation used throughout the paper.
* $`B(a,r)`$ denotes the ball $`\{x^N:|x|<r\}`$ in $`^N`$ with center at $`a`$ and radius $`r`$.
* $`S_r=\{(x,0)^2\times ^{N2}:|x|=r\}`$ denotes the circle of radius $`r`$ in the plane $`^2\times \{0\}`$.
* $`\delta _x`$ denotes the Dirac mass located at point $`x^N`$.
* $`𝒟(^N)`$ is the space of smooth functions with compact support in $`^N`$.
* $`𝒟^{1,2}(^N)`$ is the closure of $`𝒟(^N)`$ with respect to the Dirichlet norm $`(_^N|u|^2𝑑x)^{1/2}`$.
* $`C_0(^N)`$ denotes the closure of continuous functions with compact support in $`^N`$ with respect to the uniform norm.
* $`dist(x,A)`$ denotes the distance of the point $`x^N`$ from the set $`A^N`$.
* $`_p`$ denotes the norm in the Lebesgue space $`L^p(^N)`$.
* $`A\mathrm{}B`$ denotes the symmetric difference of sets $`A`$ and $`B`$, i.e. $`A\mathrm{}B=(AB)(BA)`$.
| |
| --- |
| 2. $`k=20`$, $`m=2`$, $`\lambda _0,\lambda _1>0`$, $`\lambda _2<0`$. |
| |
| 3. $`k=6`$, $`m=2`$, $`\lambda _1>0`$, $`\lambda _0,\lambda _2<0`$. |
## 2. The problem with one singularity
For any $`\lambda <(N2)^2/4`$, the problem with one singularity
(16)
$$\{\begin{array}{cc}\mathrm{\Delta }u=\frac{\lambda }{|x|^2}u+u^{2^{}1},x^N,\hfill & \\ u>0\text{ in }^N\{0\},\text{ and }u𝒟^{1,2}(^N),\hfill & \end{array}$$
admits a family of positive solutions given by
(17)
$$w_\mu ^\lambda (x)=\mu ^{\frac{N2}{2}}w_1^\lambda \left(\frac{x}{\mu }\right),\mu >0,$$
where we denote
(18)
$$w_1^\lambda (x)=\frac{\left(N(N2)\nu _\lambda ^2\right)^{\frac{N2}{4}}}{\left(|x|^{1\nu _\lambda }(1+|x|^{2\nu _\lambda })\right)^{\frac{N2}{2}}},\text{and}\nu _\lambda =\left(1\frac{4\lambda }{(N2)^2}\right)^{1/2}.$$
Moreover, when $`0\lambda <(N2)^2/4`$, all $`w_\mu ^\lambda (x)`$ minimize the associated Rayleigh quotient and the minimum can be computed as:
(19)
$$S(\lambda ):=\underset{u𝒟^{1,2}(^N)\{0\}}{inf}\frac{Q_\lambda (u)}{\left(_^N|u|^2^{}𝑑x\right)^{2/2^{}}}=\frac{Q_\lambda (w_\mu ^\lambda )}{\left(_^N|w_\mu ^\lambda |^2^{}𝑑x\right)^{2/2^{}}}=\left(1\frac{4\lambda }{(N2)^2}\right)^{\frac{N1}{N}}S,$$
where we denoted the quadratic form $`Q_\lambda (u)=_^N|u|^2𝑑x\lambda _^N\frac{u^2}{|x|^2}𝑑x`$, see , and $`S`$ is the best constant in the Sobolev inequality
$$Su_{L^2^{}(^N)}^2u_{𝒟^{1,2}(^N)}^2.$$
As minimizers of problem (19), we consider
(20)
$$z_\mu ^\lambda (x)=\frac{w_\mu ^\lambda (x)}{\left(_^N|w_\mu ^\lambda |^2^{}𝑑x\right)^{1/2^{}}}=\alpha _{\lambda ,N}\mu ^{\frac{N2}{2}}\left(\left|\frac{x}{\mu }\right|^{1\nu _\lambda }+\left|\frac{x}{\mu }\right|^{1+\nu _\lambda }\right)^{\frac{N2}{2}}$$
where $`\alpha _{\lambda ,N}=\left(N(N2)\nu _\lambda ^2\right)^{\frac{N2}{4}}w_1^\lambda _{L^2^{}}^1`$ is a positive constant depending only on $`\lambda `$ and $`N`$, so that for $`0\lambda <(N2)^2/4`$
(21)
$$S(\lambda )=Q_\lambda (z_\mu ^\lambda )\text{for all }\mu >0.$$
For $`\mathrm{}<\lambda <(N2)^2/4`$, we also set
(22)
$$S_k(\lambda ):=\underset{u𝒟_k^{1,2}(^N)\{0\}}{inf}\frac{Q_\lambda (u)}{\left(_^N|u|^2^{}𝑑x\right)^{2/2^{}}}.$$
We note that $`S(\lambda )S_k(\lambda )`$, and equality holds whenever $`\lambda 0`$. Moreover the following result has been proved in .
###### Lemma 2.1 (see , Lemma 6.1).
Let $`N4`$. If $`S_k(\lambda )<k^{2/N}S`$ then $`S_k(\lambda )`$ is achieved.
In it is proved that if $`\lambda (0,(N2)^2/4)`$ then all solutions to (16) are of the form (17) while if $`\lambda <<0`$ then also non radial solutions to (16) can exist. The behavior of any solution (radial and non radial) to problem (16) near the singularities $`0`$ and $`\mathrm{}`$ is described by the following theorem.
###### Theorem 2.2.
If $`\lambda <(N2)^2/4`$ and $`u𝒟^{1,2}(^N)`$ is a solution to problem (16), then there exist positive constant $`\kappa _0(u)`$ and $`\kappa _{\mathrm{}}(u)`$ depending on $`u`$ such that
(23) $`u(x)=|x|^{\frac{N2}{2}(1\nu _\lambda )}\left[\kappa _0(u)+O(|x|^\alpha )\right],\text{as }x0,`$
(24) $`u(x)=|x|^{\frac{N2}{2}(1+\nu _\lambda )}\left[\kappa _{\mathrm{}}(u)+O(|x|^\alpha )\right],\text{as }|x|+\mathrm{},`$
for some $`\alpha (0,1)`$.
###### Remark 2.3.
Putting together (2324) we deduce that there exists a positive constant $`\kappa (u)`$ depending on $`u`$ such that
(25)
$$\frac{1}{\kappa (u)}w_1^\lambda (x)u(x)\kappa (u)w_1^\lambda (x).$$
Proof of Theorem 2.2. Set
$$a_\lambda =\frac{(N2)\left(1\nu _\lambda \right)}{2},v(x)=|x|^{a_\lambda }u(x).$$
Then the function $`v`$ belongs to $`𝒟_{a_\lambda }^{1,2}(^N)`$ where $`𝒟_{a_\lambda }^{1,2}(^N)`$ denotes the space obtained by completion of $`𝒟(^N)`$ with respect to the weighted Dirichlet norm
$$v_{𝒟_{a_\lambda }^{1,2}(^N)}:=\left(_^N|x|^{2a_\lambda }|v|^2𝑑x\right)^{1/2}.$$
Moreover $`v`$ solves equation
(26)
$$div\left(|x|^{2a_\lambda }v\right)=\frac{v^{2^{}1}}{|x|^{2^{}a_\lambda }}.$$
From \[13, Theorem 1.2\], it follows that $`v`$ is Hölder continuous; in particular expansion (23) holds for $`\kappa _0(u)=v(0)`$ and some $`\alpha (0,1)`$. Moreover $`v(0)`$ is strictly positive in view of Harnack’s inequality for degenerate operators proved in , see also ; we mention that weights of type $`|x|^{2a}`$ with $`a<\frac{N2}{2}`$ belong to the class of quasi-conformal weights considered in .
To deduce (24), we perform the change of variable
$$\stackrel{~}{v}(x)=|x|^{a_\lambda (N2)}u\left(\frac{x}{|x|^2}\right),$$
and observe that the transformed function $`\stackrel{~}{v}`$ solves equation (26). Hence \[13, Theorem 1.2\] yields that $`\stackrel{~}{v}`$ is Hölder continuous and admits the following expansion
$$\stackrel{~}{v}(x)=\stackrel{~}{v}(0)+O(|x|^\alpha ),\text{as }x0,$$
for some $`\alpha (0,1)`$, where $`\stackrel{~}{v}(0)>0`$ in view of Harnack’s inequality in . Coming back to $`u`$ we obtain that $`u`$ satisfies (24) with $`\kappa _{\mathrm{}}(u)=\stackrel{~}{v}(0)>0`$.
## 3. Hardy’s inequality with singularity on a circle
We prove now the Hardy type inequality for potentials which are singular at circles stated in Theorem 1.1.
Proof of Theorem 1.1. Let us consider the minimization problem
$$I(S_r):=\underset{\begin{array}{c}u𝒟^{1,2}(^N)\\ u0\end{array}}{inf}\frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|u(y)|^2\left(\text{ }{\displaystyle _{S_r}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y}=\underset{\begin{array}{c}u𝒟(^N\{0\})\\ u0\end{array}}{inf}\frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|u(y)|^2\left(\text{ }{\displaystyle _{S_r}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y},$$
where the last equality is due to density of $`𝒟(^N\{0\})`$ in $`𝒟^{1,2}(^N)`$ (see e.g. \[5, Lemma 2.1\]). An easy calculation shows that for any $`u𝒟^{1,2}(^N)`$
$$\frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|u(y)|^2\left(\text{ }{\displaystyle _{S_r}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y}=\frac{{\displaystyle _^N}|v(y)|^2𝑑y}{{\displaystyle _^N}|v(y)|^2\left(\text{ }{\displaystyle _{S_1}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y}$$
where $`v(y)=u(ry)`$. Therefore
(27)
$$I(S_r)=I(S_1)\text{for any }r>0.$$
In view of (27), it is enough to prove the theorem for $`r=1`$. The proof consists in three steps.
Step 1: Inequality (6) i.e. $`I\mathbf{(}S_\mathrm{𝟏}\mathbf{)}\mathbf{}\mathbf{\left(}\frac{N\mathbf{}\mathrm{𝟐}}{\mathrm{𝟐}}\mathbf{\right)}^\mathrm{𝟐}`$.
For any $`u𝒟^{1,2}(^N)`$, $`u0`$ a.e., we consider the Schwarz symmetrization $`u^{}`$ of $`u`$ defined as
(28)
$$u^{}(x):=inf\{t>0:\left|\{y^N:u(y)>t\}\right|\omega _N|x|^N\}$$
where $`||`$ denotes the Lebesgue measure of $`^N`$ and $`\omega _N`$ is the volume of the standard unit $`N`$-ball. From \[31, Theorem 21.8\], it follows that for any $`xS_1`$
$$_^N\frac{|u(y)|^2}{|xy|^2}𝑑y_^N|u^{}(y)|^2\left[\left(\frac{1}{|xy|}\right)^{}\right]^2.$$
Since $`\left(\frac{1}{|xy|}\right)^{}=\frac{1}{|y|}`$, we deduce
(29)
$$_^N\frac{|u(y)|^2}{|xy|^2}𝑑y_^N\frac{|u^{}(y)|^2}{|y|^2}𝑑y.$$
Moreover by Polya-Szego inequality
(30)
$$_^N|u^{}|^2_^N|u|^2.$$
From (2930) and the classical Hardy’s inequality, it follows that, for any $`u𝒟^{1,2}(^N)\{0\}`$, $`u0`$ a.e.,
$`{\displaystyle \frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|u(y)|^2\left(\text{ }{\displaystyle _{S_1}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y}}`$ $`={\displaystyle \frac{{\displaystyle _^N}|u(y)|^2𝑑y}{\text{ }{\displaystyle _{S_1}}\left({\displaystyle _^N}|u(y)|^2{\displaystyle \frac{dy}{|xy|^2}}\right)𝑑\sigma (x)}}`$
$`{\displaystyle \frac{{\displaystyle _^N}|u^{}(y)|^2𝑑y}{{\displaystyle _^N}{\displaystyle \frac{|u^{}(y)|^2}{|y|^2}}𝑑y}}\left({\displaystyle \frac{N2}{2}}\right)^2.`$
Due to evenness of the quotient we are minimizing, to compute $`I(S_1)`$ it is enough to take the infimum over positive functions. Hence passing to the infimum in the above inequality, we obtain $`I(S_1)\left(\frac{N2}{2}\right)^2`$.
Step 2: Optimality of the constant, i.e. $`I\mathbf{(}S_\mathrm{𝟏}\mathbf{)}\mathbf{=}\mathbf{\left(}\frac{N\mathbf{}\mathrm{𝟐}}{\mathrm{𝟐}}\mathbf{\right)}^\mathrm{𝟐}`$.
We fix $`u𝒟(^N\{0\})`$ and let $`0<r<R`$ be such that $`suppu\{x^N:r<|x|<R\}`$. For any $`0<\lambda <<1`$, we set $`\stackrel{~}{u}_\lambda (x)=u(\lambda x)`$. Hence we have
(31) $`I(S_1)`$ $`{\displaystyle \frac{{\displaystyle _^N}|\stackrel{~}{u}_\lambda (y)|^2𝑑y}{{\displaystyle _^N}|\stackrel{~}{u}_\lambda (y)|^2\left(\text{ }{\displaystyle _{S_1}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)𝑑y}}={\displaystyle \frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|u(y)|^2\left(\text{ }{\displaystyle _{S_1}}{\displaystyle \frac{d\sigma (x)}{|\lambda xy|^2}}\right)𝑑y}}.`$
Since $`\frac{1}{|\lambda xy|^2}\frac{4}{r^2}`$ for all $`ysuppu`$, $`xS_1`$, and $`0<\lambda <\frac{r}{2}`$, by Dominated Convergence Theorem we deduce that $`_^N|u(y)|^2\left(\text{ }_{S_1}\frac{d\sigma (x)}{|\lambda xy|^2}\right)𝑑y`$ converges to $`_^N\frac{|u(y)|^2}{|y|^2}𝑑y`$ as $`\lambda 0`$, hence passing to the limit in (31) we obtain
$$I(S_1)\frac{_^N|u(y)|^2𝑑y}{_^N|y|^2|u(y)|^2𝑑y}\text{for any }uu𝒟(^N\{0\}).$$
By density of $`𝒟(^N\{0\})`$ in $`𝒟^{1,2}(^N)`$ we deduce
$$I(S_1)\underset{𝒟^{1,2}(^N)\{0\}}{inf}\frac{_^N|u(y)|^2𝑑y}{_^N|y|^2|u(y)|^2𝑑y}=\left(\frac{N2}{2}\right)^2$$
where the last equality follows from the optimality of the constant $`\left(\frac{N2}{2}\right)^2`$ in the classical Hardy inequality (see \[18, Lemma 2.1\]. Collecting the above inequality with the one proved in Step 1, we find $`I(S_1)=\left(\frac{N2}{2}\right)^2`$.
Step 3: The infimum $`I\mathbf{(}S_\mathrm{𝟏}\mathbf{)}`$ is not attained.
Arguing by contradiction, assume that the infimum $`I(S_1)`$ is achieved by some $`\overline{u}𝒟^{1,2}(^N)`$. We can assume that $`\overline{u}0`$ (otherwise we consider $`|\overline{u}|`$ which is also a minimizer by evenness of the quotient). Hence from (29) and (30)
$$\left(\frac{N2}{2}\right)^2=\frac{{\displaystyle _^N}|\overline{u}(y)|^2𝑑y}{\text{ }{\displaystyle _{S_1}}\left({\displaystyle _^N}|\overline{u}(y)|^2{\displaystyle \frac{dy}{|xy|^2}}\right)𝑑\sigma (x)}\frac{{\displaystyle _^N}|\overline{u}^{}(y)|^2𝑑y}{{\displaystyle _^N}{\displaystyle \frac{|\overline{u}^{}(y)|^2}{|y|^2}}𝑑y}\left(\frac{N2}{2}\right)^2.$$
Therefore the above inequalities are indeed equalities and this implies that the infimum
(32)
$$\left(\frac{N2}{2}\right)^2=\underset{u𝒟^{1,2}(^N)\{0\}}{inf}\frac{{\displaystyle _^N}|u(y)|^2𝑑y}{{\displaystyle _^N}|y|^2|u(y)|^2𝑑y}$$
which yields the best constant in the classical Hardy inequality, is achieved by $`\overline{u}^{}`$. Since it is known that the infimum in (32) cannot be attained (see \[28, Remark 1.2\]), we reach a contradiction.
## 4. The Palais-Smale condition under $`_k\times 𝕊𝕆(N2)`$ and $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$-invariance
Let us define the functional $`J_k:𝒟^{1,2}(^N)`$ associated to equation (2) as
(33) $`J_k(u)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^N}|u|^2𝑑x{\displaystyle \frac{\lambda _0}{2}}{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|x|^2}}𝑑x`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{2}}{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|xa_i^{\mathrm{}}|^2}}𝑑x{\displaystyle \frac{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}{2^{}}}{\displaystyle _^N}|u|^2^{}𝑑x.`$
The choice of location of the singularities ensures that $`J_k`$ is $`_k\times 𝕊𝕆(N2)`$-invariant. Since $`_k\times 𝕊𝕆(N2)`$ acts by isometries on $`𝒟^{1,2}(^N)`$, we can apply the Principle of Symmetric Criticality by Palais to deduce that the critical points of $`J_k`$ restricted to $`𝒟_k^{1,2}(^N)`$ are also critical points of $`J_k`$ in $`𝒟^{1,2}(^N)`$. Therefore, if $`u`$ is a critical point of $`J_k`$ in $`𝒟_k^{1,2}(^N)`$, $`u>0`$ outside singularities, then $`v=S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_{\mathrm{}})^{1/(2^{}2)}u`$ is a solution to equation (2).
The following theorem provides a local Palais-Smale condition for $`J_k`$ restricted to $`𝒟_k^{1,2}(^N)`$ below some critical threshold. We emphasize that the invariance of the problem by the action of a subgroup of orthogonal transformation allows to recover some compactness, in the sense that concentration points of invariant functions must be located in some symmetric way, thus reducing the possibility of loss of compactness. The restriction on dimension $`N4`$ is required to avoid the presence of possible concentration points on $`\{0\}\times ^{N2}`$. Indeed when $`N=3`$, $`𝕊𝕆(N2)=𝕊𝕆(1)`$ is a discrete group, making thus possible concentration at points on the axis $`\{0\}\times `$.
We mention that the Concentration-Compactness method under the action of $`_k\times 𝕊𝕆(N2)`$ was used by several authors to find $`k`$-bump solutions with prescribed symmetry for different classes of nonlinear elliptic equations: nonlinear Schrödinger equation in , nonlinear elliptic equations in symmetric domains in , nonlinear elliptic equations of Caffarelli-Kohn-Nirenberg type in , and elliptic equations with Hardy potential and critical growth in .
###### Theorem 4.1.
Assume $`N4`$ and $`\lambda _0^++k_{j=1}^m\lambda _j^+<\frac{(N2)^2}{4}`$. Let $`\{u_n\}_n𝒟_k^{1,2}(^N)`$ be a Palais-Smale sequence for $`J_k`$ restricted to $`𝒟_k^{1,2}(^N)`$, namely
$$\underset{n\mathrm{}}{lim}J_k(u_n)=c<\mathrm{}\text{ in }\text{and}\underset{n\mathrm{}}{lim}J_k^{}(u_n)=0\text{ in the dual space }(𝒟_k^{1,2}(^N))^{}.$$
If
(34) $`c<{\displaystyle \frac{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{1\frac{N}{2}}}{N}}\mathrm{min}\{k^{\frac{2}{N}}S,k^{\frac{2}{N}}S(\lambda _1),\mathrm{},k^{\frac{2}{N}}S(\lambda _m),S_k(\lambda _0),S_k(\lambda _0+k{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\lambda _{\mathrm{}})\}^{\frac{N}{2}},`$
then $`\{u_n\}_n`$ has a subsequence strongly converging in $`𝒟_k^{1,2}(^N)`$.
Proof. Let $`\{u_n\}_n`$ be a Palais-Smale sequence for $`J_k`$ in $`𝒟_k^{1,2}(^N)`$, then from Hardy’s and Sobolev’s inequalities it is easy to prove that $`\{u_n\}_n`$ is a bounded sequence in $`𝒟^{1,2}(^N)`$. Hence, up to a subsequence, $`u_nu_0\text{ in }𝒟^{1,2}(^N)`$ and $`u_nu_0`$ almost everywhere. Therefore, from the Concentration Compactness Principle by P. L. Lions (see ), we deduce the existence of a subsequence, still denoted by $`\{u_n\}`$, for which there exist an at most countable set $`𝒥`$, points $`x_j^N\{0,a_i^{\mathrm{}},1ik,1\mathrm{}m\}`$, real numbers $`\mu _{x_j},\nu _{x_j}`$, $`j𝒥`$, and $`\mu _0`$, $`\nu _0`$, $`\mu _{a_i^{\mathrm{}}}`$, $`\nu _{a_i^{\mathrm{}}}`$, $`\gamma _{a_i^{\mathrm{}}}`$, $`i=1,\mathrm{},k`$, $`\mathrm{}=1,\mathrm{},m`$, such that the following convergences hold in the sense of measures
(35) $`|u_n|^2d\mu |u_0|^2+\mu _0\delta _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}\mu _{a_i^{\mathrm{}}}\delta _{a_i^{\mathrm{}}}+{\displaystyle \underset{j𝒥}{}}\mu _{x_j}\delta _{x_j},`$
(36) $`|u_n|^2^{}d\nu =|u_0|^2^{}+\nu _0\delta _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}\nu _{a_i^{\mathrm{}}}\delta _{a_i^{\mathrm{}}}+{\displaystyle \underset{j𝒥}{}}\nu _{x_j}\delta _{x_j},`$
(37) $`\lambda _0{\displaystyle \frac{u_n^2}{|x|^2}}d\gamma _0=\lambda _0{\displaystyle \frac{u_0^2}{|x|^2}}+\gamma _0\delta _0,`$
(38) $`\lambda _{\mathrm{}}{\displaystyle \frac{u_n^2}{|xa_i^{\mathrm{}}|^2}}d\gamma _{a_i^{\mathrm{}}}=\lambda _{\mathrm{}}{\displaystyle \frac{u_0^2}{|xa_i^{\mathrm{}}|^2}}+\gamma _{a_i^{\mathrm{}}}\delta _{a_i^{\mathrm{}}},\text{for any}i=1,\mathrm{},k,\mathrm{}=1,\mathrm{},m.`$
From Sobolev’s inequality it follows that
(39)
$$S\nu _{x_j}^{\frac{2}{2^{}}}\mu _{x_j}\text{ for all }j𝒥\text{and}S\nu _{a_i}^{\frac{2}{2^{}}}\mu _{a_i}\text{ for all }i=1,\mathrm{},k.$$
Possible concentration at infinity of the sequence can be quantified by the following numbers
(40)
$$\nu _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}|u_n|^2^{}𝑑x,\mu _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}|u_n|^2𝑑x,$$
and
$$\gamma _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}\left(\lambda _0+k\underset{\mathrm{}=1}{\overset{m}{}}\lambda _{\mathrm{}}\right)\frac{u_n^2}{|x|^2}𝑑x.$$
Let us first prove that possible concentration points are located in a symmetric fashion. Pointwise convergence of $`u_n𝒟_k^{1,2}(^N)`$ to $`u_0`$ implies that $`u_0`$ is invariant by the $`_k\times 𝕊𝕆(N2)`$-action, hence $`u_0𝒟_k^{1,2}(^N)`$. Moreover, for any $`\varphi C_0(^N)`$ and for any $`\tau _k\times 𝕊𝕆(N2)`$ we have
$$_^N|u_n|^2^{}\varphi =_^N|u_n|^2^{}(\varphi \tau ^1)$$
that, passing to the limit as $`n\mathrm{}`$, yields
$$\underset{\mathrm{}=1}{\overset{m}{}}\underset{i=1}{\overset{k}{}}\nu _{a_i^{\mathrm{}}}\varphi (a_i^{\mathrm{}})+\underset{j𝒥}{}\nu _{x_j}\varphi (x_j)=\underset{\mathrm{}=1}{\overset{m}{}}\underset{i=1}{\overset{k}{}}\nu _{a_i^{\mathrm{}}}\varphi (\tau ^1(a_i^{\mathrm{}}))+\underset{j𝒥}{}\nu _{x_j}\varphi (\tau ^1(x_j)).$$
Choosing $`\varphi =\varphi _{i,\epsilon }^{\mathrm{}}`$ such that $`\varphi _{i,\epsilon }^{\mathrm{}}1`$ on $`B(a_i^{\mathrm{}},\epsilon /2)`$, $`\varphi _{i,\epsilon }^{\mathrm{}}0`$ on $`^NB(a_i^{\mathrm{}},\epsilon )`$, $`0\varphi _{i,\epsilon }^{\mathrm{}}1`$, and letting $`\epsilon 0`$, we find that $`\nu _{a_i^{\mathrm{}}}=\nu _{\tau ^1(a_i^{\mathrm{}})}`$. Hence we deduce that for any $`\mathrm{}=1,\mathrm{},m`$ there exists $`\nu _a^{\mathrm{}}`$ such that
$$\nu _{a_i^{\mathrm{}}}=\nu _a^{\mathrm{}}\text{for any }i=1,\mathrm{},k.$$
Fix $`j𝒥`$ and let $`\varphi =\varphi _{j,\epsilon }`$ such that $`\varphi _{j,\epsilon }1`$ on $`B(x_j,\epsilon /2)`$, $`\varphi _{j,\epsilon }0`$ on $`^NB(x_j,\epsilon )`$, $`0\varphi _{j,\epsilon }1`$. Letting $`\epsilon 0`$, we find that
$`(\mathrm{i})`$ $`\text{either }\nu _{x_j}=0`$
$`(\mathrm{ii})`$ $`\text{or for any }\tau Z_k\times 𝕊𝕆(N2)\text{ there exists }i𝒥\text{ such that }\tau (x_j)=x_i.`$
When $`N4`$, $`𝕊𝕆(N2)`$ is a continuous group, hence for any $`x^2\times \{0\}^N`$, the set $`\{\tau (x):\tau Z_k\times 𝕊𝕆(N2)\}`$ is more than countable. If alternative (ii) holds, from at most countability of $`𝒥`$ we deduce that $`x_j^2\times \{0\}^N`$. Moreover, arguing as above we can prove that if $`\tau (x_j)=x_i`$ for some $`\tau Z_k\times 𝕊𝕆(N2)`$ then $`\nu _{x_i}=\nu _{x_j}`$. Hence we can rewrite (36) as
$$|u_n|^2^{}d\nu =|u_0|^2^{}+\nu _0\delta _0+\underset{\mathrm{}=1}{\overset{m}{}}\nu _a^{\mathrm{}}\underset{i=1}{\overset{k}{}}\delta _{a_i^{\mathrm{}}}+\underset{\mathrm{}}{}\nu _y^{\mathrm{}}\underset{i=1}{\overset{k}{}}\delta _{y_i^{\mathrm{}}}$$
where $``$ is an at most countable set,
$$\{y_i^{\mathrm{}}:i=1,\mathrm{},k,\mathrm{}\}\{x_j:j𝒥\}\text{and}y_i^{\mathrm{}}=e^{2\pi \sqrt{1}/k}y_{i1}^{\mathrm{}}.$$
Concentration at non singular points. We claim that
(41)
$$𝒥\text{is finite and for }j𝒥\text{ either }\nu _{x_j}=0\text{ or }\nu _{x_j}\left(\frac{S}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
Indeed, if $`\nu _{x_j}>0`$, we have $`x_j=y_s^{\mathrm{}}`$ and $`\nu _{x_j}=\nu _y^{\mathrm{}}`$ for some $`1sk`$, $`\mathrm{}`$. For $`\epsilon >0`$, let $`\varphi _{\mathrm{}}^\epsilon `$ be a smooth cut-off function in $`𝒟_k^{1,2}(^N)`$, $`0\varphi _{\mathrm{}}^\epsilon (x)1`$ such that
$$\varphi _{\mathrm{}}^\epsilon (x)=1\text{ if }x\underset{i=1}{\overset{k}{}}B(y_i^{\mathrm{}},\frac{\epsilon }{2}),\varphi _{\mathrm{}}^\epsilon (x)=0\text{ if }x\underset{i=1}{\overset{k}{}}B(y_i^{\mathrm{}},\epsilon ),\text{and}|\varphi _{\mathrm{}}^\epsilon |\frac{4}{\epsilon }.$$
Testing $`J_k^{}(u_n)`$ with $`u_n\varphi _{\mathrm{}}^\epsilon `$ we obtain
$`0`$ $`=\underset{\epsilon 0}{lim}\underset{n\mathrm{}}{lim}J_k^{}(u_n),u_n\varphi _{\mathrm{}}^\epsilon {\displaystyle \underset{i=1}{\overset{k}{}}}\mu _{y_i^{\mathrm{}}}kS_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _y^{\mathrm{}}.`$
By (39) we have that $`S(\nu _y^{\mathrm{}})^{\frac{2}{2^{}}}\mu _{y_i^{\mathrm{}}}`$, then we obtain that $`\nu _y^{\mathrm{}}\left(\frac{S}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}`$. Hence either $`\nu _{x_j}=0`$ or $`\nu _{x_j}\left(\frac{S}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}`$, which implies that $`𝒥`$ is finite. In particular also $``$ is finite. The claim is proved.
Concentration at vertices of polygons. We claim that
(42)
$$\text{ for each}\mathrm{}=1,2,\mathrm{},m\text{either}\nu _a^{\mathrm{}}=0\text{or}\nu _a^{\mathrm{}}\left(\frac{S(\lambda _{\mathrm{}})}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
In order to prove claim (42), for each $`i=1,2,\mathrm{},k`$ and $`\mathrm{}=1,2,\mathrm{},m`$ we consider the smooth cut-off function $`\varphi _{i,\epsilon }^{\mathrm{}}`$ satisfying $`0\varphi _{i,\epsilon }^{\mathrm{}}(x)1`$,
$$\varphi _{i,\epsilon }^{\mathrm{}}(x)=1\text{ if }|xa_i^{\mathrm{}}|\frac{\epsilon }{2},\varphi _{i,\epsilon }^{\mathrm{}}(x)=0\text{ if }|xa_i^{\mathrm{}}|\epsilon ,\text{and}|\varphi _{i,\epsilon }^{\mathrm{}}|\frac{4}{\epsilon }.$$
From (19) we obtain that
$$\frac{_^N|(u_n\varphi _{i,\epsilon }^{\mathrm{}})|^2𝑑x\lambda _{\mathrm{}}_^N|xa_i^{\mathrm{}}|^2|\varphi _{i,\epsilon }^{\mathrm{}}|^2u_n^2𝑑x}{\left(_^N|\varphi _{i,\epsilon }^{\mathrm{}}u_n|^2^{}\right)^{2/2^{}}}S(\lambda _{\mathrm{}})$$
hence passing to limit as $`n\mathrm{}`$ and $`\epsilon 0`$ we obtain
(43)
$$\mu _{a_i^{\mathrm{}}}\gamma _{a_i^{\mathrm{}}}+S(\lambda _{\mathrm{}})\left(\nu _a^{\mathrm{}}\right)^{2/2^{}}.$$
For $`\epsilon >0`$, let $`\psi _{\mathrm{}}^\epsilon `$ be a smooth cut-off function in $`𝒟_k^{1,2}(^N)`$, $`0\psi _{\mathrm{}}^\epsilon (x)1`$ such that
$$\psi _{\mathrm{}}^\epsilon (x)=1\text{ if }x\underset{i=1}{\overset{k}{}}B(a_i^{\mathrm{}},\frac{\epsilon }{2}),\psi _{\mathrm{}}^\epsilon (x)=0\text{ if }x\underset{i=1}{\overset{k}{}}B(a_i^{\mathrm{}},\epsilon ),\text{and}|\psi _{\mathrm{}}^\epsilon |\frac{4}{\epsilon }.$$
Testing $`J_k^{}(u_n)`$ with $`u_n\psi _{\mathrm{}}^\epsilon `$ and letting $`n\mathrm{}`$ and $`\epsilon 0`$ we infer that
(44)
$$\underset{i=1}{\overset{k}{}}\mu _{a_i^{\mathrm{}}}\underset{i=1}{\overset{k}{}}\gamma _{a_i^{\mathrm{}}}kS_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _a^{\mathrm{}}.$$
From (43) and (44) we derive (42).
Concentration at the origin. We claim that
(45)
$$\text{either}\nu _0=0\text{or}\nu _0\left(\frac{S_k(\lambda _0)}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
In order to prove claim (45), we consider a smooth cut-off function $`\psi _0^\epsilon 𝒟_k^{1,2}(^N)`$ satisfying $`0\psi _0^\epsilon (x)1`$,
$$\psi _0^\epsilon (x)=1\text{ if }|x|\frac{\epsilon }{2},\psi _0^\epsilon (x)=0\text{ if }|x|\epsilon ,\text{and}|\psi _0^\epsilon |\frac{4}{\epsilon }.$$
From (22) we obtain that
$$\frac{_^N|(u_n\psi _0^\epsilon )|^2𝑑x\lambda _0_^N|x|^2|\psi _0^\epsilon |^2u_n^2𝑑x}{\left(_^N|\psi _0^\epsilon u_n|^2^{}\right)^{2/2^{}}}S_k(\lambda _0)$$
hence passing to limit as $`n\mathrm{}`$ and $`\epsilon 0`$ we obtain
(46)
$$\mu _0\gamma _0+S_k(\lambda _0)\left(\nu _0\right)^{2/2^{}}.$$
On the other hand, testing $`J_k^{}(u_n)`$ with $`u_n\psi _0^\epsilon `$ and letting $`n\mathrm{}`$ and $`\epsilon 0`$ we infer that
(47)
$$\mu _0\gamma _0S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _0.$$
From (46) and (47) we deduce (45).
Concentration at infinity. We claim that
(48)
$$\text{either}\nu _{\mathrm{}}=0\text{or}\nu _{\mathrm{}}\left(\frac{S_k(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}})}{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
In order to prove (48), we study the possibility of concentration at $`\mathrm{}`$. Let $`\psi _R`$ be a regular radial cut-off function such that
$$0\psi _R(x)1,\psi _R(x)=\{\begin{array}{c}1,\text{ if }|x|>2R\hfill \\ 0,\text{ if }|x|<R,\hfill \end{array}\text{and}|\psi _R|\frac{2}{R}.$$
From (22) we obtain that
(49)
$$\frac{_^N|(u_n\psi _R)|^2𝑑x\left(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}}\right)_^N\frac{\psi _R^2u_n^2}{|x|^2}𝑑x}{\left(_^N|\psi _Ru_n|^2^{}\right)^{2/2^{}}}S_k\left(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}}\right).$$
Taking $`lim\; sup`$ as $`n\mathrm{}`$ and limit as $`R+\mathrm{}`$, standard calculations yield
(50)
$$\mu _{\mathrm{}}\gamma _{\mathrm{}}S_k\left(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}}\right)\nu _{\mathrm{}}^{2/2^{}}.$$
Testing $`J_k^{}(u_n)`$ with $`u_n\psi _R`$ and letting $`n\mathrm{}`$ and $`R+\mathrm{}`$, we obtain
(51)
$$\mu _{\mathrm{}}\gamma _{\mathrm{}}S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _{\mathrm{}}.$$
Claim (48) follows from (50) and (51).
As a conclusion we obtain
(52) $`c`$ $`=J_k(u_n){\displaystyle \frac{1}{2}}J_k^{}(u_n),u_n+o(1)={\displaystyle \frac{1}{N}}S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m){\displaystyle _^N}|u_n|^2^{}𝑑x+o(1)`$
$`={\displaystyle \frac{S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}{N}}\left\{{\displaystyle _^N}|u_0|^2^{}𝑑x+\nu _0+k{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\nu _a^{\mathrm{}}+k{\displaystyle \underset{\mathrm{}}{}}\nu _y^{\mathrm{}}+\nu _{\mathrm{}}\right\}.`$
From (34), (41), (42), (45), (48), and (52), we deduce that $`\nu _0=0`$, $`\nu _y^{\mathrm{}}=0`$ for any $`\mathrm{}`$, $`\nu _a^{\mathrm{}}=0`$ for any $`\mathrm{}=1,\mathrm{},m`$, and $`\nu _{\mathrm{}}=0`$. Then, up to a subsequence, $`u_nu_0`$ in $`𝒟_k^{1,2}(^N)`$.
The functional $`J:𝒟^{1,2}(^N)`$ associated to equation (4) is
(53) $`J(u)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^N}|u|^2𝑑x{\displaystyle \frac{\lambda _0}{2}}{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|x|^2}}𝑑x`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\mathrm{\Lambda }_{\mathrm{}}}{2}}{\displaystyle _^N}\left(\text{ }{\displaystyle _{S_r_{\mathrm{}}}}{\displaystyle \frac{u^2(y)}{|xy|^2}}𝑑\sigma (x)\right)𝑑y{\displaystyle \frac{S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}{2^{}}}{\displaystyle _^N}|u|^2^{}𝑑x.`$
The functional $`J`$ is $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$-invariant. Since $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$ acts by isometries on $`𝒟^{1,2}(^N)`$, the Principle of Symmetric Criticality by Palais implies that the critical points of $`J`$ restricted to $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$ are also critical points of $`J`$ in $`𝒟^{1,2}(^N)`$. Therefore, any critical point of $`J`$ in $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$ provides a solution to equation (4).
The following theorem is the analogous of Theorem 4.1 for $`J`$ restricted to $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$. However, the fact that the singularities are spread over circles instead of being concentrated at atoms reduces the possibility of lack of compactness. Indeed, according to P.L. Lions Concentration-Compactness Principle, possible forms of “non compactness” which can cause failure of Palais-Smale condition are loss of mass at infinity and concentration at an most countable set of points. When considering $`J`$ restricted to $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$, it turns out that if $`\overline{x}`$ is a concentration point of a Palais-Smale sequence, then all points of the orbit $`𝒪(\overline{x})=\{\tau \overline{x}:\tau 𝕊𝕆(2)\times 𝕊𝕆(N2)\}`$ must be concentration points. On the other hand, when $`N4`$ both groups $`𝕊𝕆(2)`$ and $`𝕊𝕆(N2)`$ are continuous, hence the only point $`\overline{x}`$ for which $`𝒪(\overline{x})`$ is at most countable is the origin. Hence concentration can occur only at $`0`$ and at $`\mathrm{}`$.
We mention that action of this type of groups was considered in to find nonradial solutions to a Euclidean scalar field equation. We refer to \[30, §1.5\] for a discussion on the relation between symmetry and compactness in variational problems.
Let us define
(54)
$$S_{\mathrm{circ}}(\lambda _0)=\underset{\begin{array}{c}u𝒟_{\mathrm{circ}}^{1,2}(^N)\\ u0\end{array}}{inf}\frac{{\displaystyle _^N}|u|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|x|^2}}𝑑x}{\left({\displaystyle _^N}|u|^2^{}𝑑x\right)^{2/2^{}}}.$$
The following theorem provides a threshold up to which $`J`$ satisfies Palais-Smale condition.
###### Theorem 4.2.
Assume $`N4`$ and $`\lambda _0^++_{j=1}^m\mathrm{\Lambda }_j^+<\frac{(N2)^2}{4}`$. Let $`\{u_n\}_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ be a Palais-Smale sequence for $`J`$ restricted to $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$, namely
$$\underset{n\mathrm{}}{lim}J(u_n)=c<\mathrm{}\text{ in }\text{and}\underset{n\mathrm{}}{lim}J^{}(u_n)=0\text{ in the dual space }(𝒟_{\mathrm{circ}}^{1,2}(^N))^{}.$$
If
(55) $`c<{\displaystyle \frac{S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{1\frac{N}{2}}}{N}}\mathrm{min}\{S_{\mathrm{circ}}(\lambda _0),S_{\mathrm{circ}}(\lambda _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}})\}^{\frac{N}{2}},`$
then $`\{u_n\}_n`$ has a converging subsequence in $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$.
Proof. Let $`\{u_n\}`$ be a Palais-Smale sequence for $`J`$ in $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$, then from Sobolev’s inequality and (6), it is easy to prove that $`\{u_n\}`$ is a bounded sequence in $`𝒟^{1,2}(^N)`$. Hence, up to a subsequence, $`u_nu_0\text{ in }𝒟^{1,2}(^N)`$, $`u_nu_0`$ almost everywhere, and, from the Concentration Compactness Principle by P. L. Lions
(56) $`|u_n|^2d\mu |u_0|^2+\mu _0\delta _0+{\displaystyle \underset{j𝒥}{}}\mu _{x_j}\delta _{x_j},`$
(57) $`|u_n|^2^{}d\nu =|u_0|^2^{}+\nu _0\delta _0++{\displaystyle \underset{j𝒥}{}}\nu _{x_j}\delta _{x_j},`$
(58) $`\lambda _0{\displaystyle \frac{u_n^2}{|x|^2}}d\gamma _0=\lambda _0{\displaystyle \frac{u_0^2}{|x|^2}}+\gamma _0\delta _0`$
where $`𝒥`$ is an at most countable set, $`x_j^N\{0\}`$, $`\mu _{x_j},\nu _{x_j}`$, $`j𝒥`$, $`\mu _0,\nu _0,\gamma _0`$, and the above convergences hold in the sense of measures. We quantify how much the sequence concentrates at infinity by the quantities $`\nu _{\mathrm{}}`$, $`\mu _{\mathrm{}}`$ defined as in (40) and
$$\gamma _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}\left(\lambda _0+\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}\right)\frac{u_n^2}{|x|^2}𝑑x.$$
From pointwise convergence of $`u_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ to $`u_0`$ we deduce that $`u_0`$ is invariant by the $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$-action, hence $`u_0𝒟_{\mathrm{circ}}^{1,2}(^N)`$. Moreover, for any $`\varphi C_0(^N)`$ and for any $`\tau 𝕊𝕆(2)\times 𝕊𝕆(N2)`$ we have
$$_^N|u_n|^2^{}\varphi =_^N|u_n|^2^{}(\varphi \tau ^1)$$
that, passing to the limit as $`n\mathrm{}`$, yields
$$\underset{j𝒥}{}\nu _{x_j}\varphi (x_j)=\underset{j𝒥}{}\nu _{x_j}\varphi (\tau ^1(x_j)).$$
Arguing as in the proof of Theorem 4.1, we deduce that for any $`j𝒥`$, either $`\nu _{x_j}=0`$ or for any $`\tau 𝕊𝕆(2)\times 𝕊𝕆(N2)`$ there exists $`i𝒥`$ such that $`\tau (x_j)=x_i`$. Namely, if for some $`j𝒥`$, $`\nu _{x_j}0`$, then $`𝒪(x_j)\{x_i:i𝒥\}`$. When $`N4`$, this is not possible since $`𝒥`$ is at most countable whereas $`𝒪(x_j)`$ is more than countable. Therefore $`\nu _{x_j}=0`$ for all $`j𝒥`$ and we can rewrite (57) as
$$|u_n|^2^{}d\nu =|u_0|^2^{}+\nu _0\delta _0.$$
Concentration at the origin. We claim that
(59)
$$\text{either}\nu _0=0\text{or}\nu _0\left(\frac{S_{\mathrm{circ}}(\lambda _0)}{S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
In order to prove claim (59), we consider a smooth cut-off function $`\psi _0^\epsilon 𝒟_{\mathrm{circ}}^{1,2}(^N)`$ satisfying $`0\psi _0^\epsilon (x)1`$,
$$\psi _0^\epsilon (x)=1\text{ if }|x|\frac{\epsilon }{2},\psi _0^\epsilon (x)=0\text{ if }|x|\epsilon ,\text{and}|\psi _0^\epsilon |\frac{4}{\epsilon }.$$
From (54) we obtain that
$$\frac{_^N|(u_n\psi _0^\epsilon )|^2𝑑x\lambda _0_^N|x|^2|\psi _0^\epsilon |^2u_n^2𝑑x}{\left(_^N|\psi _0^\epsilon u_n|^2^{}\right)^{2/2^{}}}S_{\mathrm{circ}}(\lambda _0)$$
hence passing to limit as $`n\mathrm{}`$ and $`\epsilon 0`$ we obtain
(60)
$$\mu _0\gamma _0+S_{\mathrm{circ}}(\lambda _0)\left(\nu _0\right)^{2/2^{}}.$$
On the other hand, testing $`J^{}(u_n)`$ with $`u_n\psi _0^\epsilon `$ we obtain
$$\begin{array}{c}_^N|u_n|^2\psi _0^\epsilon +_^Nu_nu_n\psi _0^\epsilon \lambda _0_^N\frac{u_n^2\psi _0^\epsilon }{|x|^2}\hfill \\ \hfill \underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_m_^N\left(\text{ }_{S_r_{\mathrm{}}}\frac{u_n^2(y)\psi _0^\epsilon (y)}{|xy|^2}𝑑\sigma (x)\right)𝑑yS_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)_^N\psi _0^\epsilon |u_n|^2^{}=o(1).\end{array}$$
In view of (7), for $`\epsilon `$ small we have
$$\begin{array}{c}_^N\left(\text{ }_{S_r}\frac{u_n^2(y)\psi _0^\epsilon (y)}{|xy|^2}𝑑\sigma (x)\right)𝑑y=_^N\frac{u_n^2(y)\psi _0^\epsilon (y)}{\sqrt{(r^2+|y|^2)^24r^2|y^{}|^2}}𝑑x\hfill \\ \hfill _^N\frac{u_n^2(y)\psi _0^\epsilon (y)}{|r^2|y|^2|}𝑑x\frac{1}{r^2\epsilon ^2}_{|y|<\epsilon }u_n^2(y)𝑑y\mathrm{const}\left(_^N|u_n|^2^{}\right)^{2/2^{}}\frac{\epsilon ^2}{r^2\epsilon ^2}\underset{\epsilon 0}{}0,\end{array}$$
therefore letting $`n\mathrm{}`$ and $`\epsilon 0`$ we infer that
(61)
$$\mu _0\gamma _0S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _0.$$
From (60) and (61) we deduce (59).
Concentration at infinity. We claim that
(62)
$$\text{either}\nu _{\mathrm{}}=0\text{or}\nu _{\mathrm{}}\left(\frac{S_{\mathrm{circ}}(\lambda _0+_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}})}{S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2,\mathrm{},\mathrm{\Lambda }_m)}\right)^{N/2}.$$
Indeed, let $`\psi _R`$ be a smooth radial cut-off function such that $`0\psi _R(x)1`$, $`\psi _R(x)=1`$ if $`|x|>2R`$, $`\psi _R(x)=0`$ if $`|x|<R`$, and $`|\psi _R|2/R`$. Using (5) and taking $`lim\; sup`$ as $`n\mathrm{}`$ and limit as $`R+\mathrm{}`$, it is easy to show that
(63)
$$\mu _{\mathrm{}}\gamma _{\mathrm{}}S_{\mathrm{circ}}\left(\lambda _0+_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}}\right)\nu _{\mathrm{}}^{2/2^{}}.$$
On the other hand, testing $`J^{}(u_n)`$ with $`u_n\psi _R`$ we obtain
(64) $`{\displaystyle _^N}|u_n|^2\psi _R+{\displaystyle _^N}u_nu_n\psi _R\lambda _0{\displaystyle _^N}{\displaystyle \frac{u_n^2\psi _R}{|x|^2}}`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_m{\displaystyle _^N}\left(\text{ }{\displaystyle _{S_r_{\mathrm{}}}}{\displaystyle \frac{u_n^2(y)\psi _R(y)}{|xy|^2}}𝑑\sigma (x)\right)𝑑yS_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m){\displaystyle _^N}\psi _R|u_n|^2^{}=o(1).`$
Let $`\overline{R}>\mathrm{max}\{r_{\mathrm{}}:\mathrm{}=1,\mathrm{},m\}`$. If $`R\overline{R}`$, in view of (7) we have for all $`\mathrm{}=1.\mathrm{},m`$
$$\left|\text{ }_{S_r_{\mathrm{}}}\frac{u_n^2(y)\psi _R(y)}{|xy|^2}𝑑\sigma (x)\frac{u_n^2(y)\psi _R(y)}{|y|^2}\right|=\frac{u_n^2(y)\psi _R(y)}{|y|^2}\left|\frac{|y|^2\sqrt{(r_{\mathrm{}}^2+|y|^2)^24r_{\mathrm{}}^2|y^{}|^2}}{\sqrt{(r_{\mathrm{}}^2+|y|^2)^24r_{\mathrm{}}^2|y^{}|^2}}\right|$$
$$\frac{u_n^2(y)\psi _R(y)}{|y|^2}\frac{r_{\mathrm{}}^2}{|y|^2}\frac{6|y|^2+r_{\mathrm{}}^2}{|y|^2r_{\mathrm{}}^2}\frac{u_n^2(y)\psi _R(y)}{|y|^4}\frac{7\overline{R}^4}{\overline{R}^2r_{\mathrm{}}^2}.$$
Since
$$_^N\frac{u_n^2(y)\psi _R(y)}{|y|^4}𝑑y\frac{1}{R^2}_{|y|>R}\frac{u_n^2(y)}{|y|^2}𝑑y\frac{\mathrm{const}}{R^2}\underset{R+\mathrm{}}{}0$$
we deduce that
$$_^N\text{ }_{S_r_{\mathrm{}}}\frac{u_n^2(y)\psi _R(y)}{|xy|^2}𝑑\sigma (x)𝑑y=_^N\frac{u_n^2(y)\psi _R(y)}{|y|^2}𝑑y+o(1)\text{as }R+\mathrm{}.$$
Therefore, letting $`n\mathrm{}`$ and $`R+\mathrm{}`$ in (64), we obtain
(65)
$$\mu _{\mathrm{}}\gamma _{\mathrm{}}S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\nu _{\mathrm{}}.$$
Claim (62) follows from (63) and (65).
As a conclusion we obtain
(66) $`c`$ $`=J(u_n){\displaystyle \frac{1}{2}}J^{}(u_n),u_n+o(1)={\displaystyle \frac{S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)}{N}}\left\{{\displaystyle _^N}|u_0|^2𝑑x+\nu _0+\nu _{\mathrm{}}\right\}.`$
From (55), (59), (62), and (66), we deduce that $`\nu _0=0`$ and $`\nu _{\mathrm{}}=0`$. Then, up to a subsequence, $`u_nu_0`$ in $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$.
## 5. Interaction estimates
For any $`u𝒟^{1,2}(^N)`$, let us consider the family of functions obtained from $`u`$ by dilation, i.e.
(67)
$$u_\mu (x)=\mu ^{\frac{N2}{2}}u(x/\mu ),\mu >0.$$
The following lemma describes the behavior of $`|x+\xi |^2|u_\mu ^\lambda |^2`$ as $`\mu 0`$ for any solution $`u^\lambda `$ of equation (16). We mention that estimates below were obtained in for radial solutions to (16) (i.e. for functions $`w_\mu ^{(\lambda )}`$ in (17)).
###### Lemma 5.1.
Let $`u^\lambda 𝒟^{1,2}(^N)`$ be a solution to (16). For any $`\xi ^N`$ there holds
$$_^N\frac{|u_\mu ^\lambda |^2}{|x+\xi |^2}𝑑x=\{\begin{array}{cc}\frac{\mu ^2}{|\xi |^2}_^N|u^\lambda |^2𝑑x+o\left(\mu ^2\right)\hfill & \text{if }\lambda <\frac{N(N4)}{4},\hfill \\ \kappa _{\mathrm{}}(u^\lambda )^2\frac{\mu ^2|\mathrm{ln}\mu |}{|\xi |^2}+o(\mu ^2|\mathrm{ln}\mu |)\hfill & \text{if }\lambda =\frac{N(N4)}{4},\hfill \\ \kappa _{\mathrm{}}(u^\lambda )^2\beta _{\lambda ,N}\frac{\mu ^{\sqrt{(N2)^24\lambda }}}{|\xi |^{\sqrt{(N2)^24\lambda }}}+o\left(\mu ^{\sqrt{(N2)^24\lambda }}\right)\hfill & \text{if }\lambda >\frac{N(N4)}{4},\hfill \end{array}$$
as $`\mu 0`$, where
$$\beta _{\lambda ,N}=_^N\frac{dx}{|x|^2|xe_1|^{N2+\sqrt{(N2)^24\lambda }}},e_1=(1,0,\mathrm{},0)^N.$$
Proof. We have that
(68) $`{\displaystyle _^N}{\displaystyle \frac{|u_\mu ^\lambda |^2}{|x+\xi |^2}}𝑑x`$ $`=\mu ^2{\displaystyle _^N}{\displaystyle \frac{|u^\lambda |^2}{|\mu x+\xi |^2}}𝑑x=\mu ^2{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}{\displaystyle \frac{|u^\lambda |^2}{|\mu x+\xi |^2}}𝑑x+\mu ^2{\displaystyle _{|x|>\frac{|\xi |}{2\mu }}}{\displaystyle \frac{|u^\lambda |^2}{|\mu x+\xi |^2}}𝑑x.`$
For $`\lambda <\frac{N(N4)}{4}`$, from (25) we have that $`u^\lambda L^2(^N)`$. From (25) and \[14, Proof of Lemma 3.4\] we deduce
$`|{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}`$ $`|u^\lambda (x)|^2[{\displaystyle \frac{1}{|\mu x+\xi |^2}}{\displaystyle \frac{1}{|\xi |^2}}]dx|`$
$`\kappa (u^\lambda )^2{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|w_1^\lambda (x)|^2\left|{\displaystyle \frac{1}{|\mu x+\xi |^2}}{\displaystyle \frac{1}{|\xi |^2}}\right|𝑑x=o(1)`$
as $`\mu 0`$ and hence, since $`u^\lambda L^2(^N)`$,
(69) $`{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2{\displaystyle \frac{dx}{|\mu x+\xi |^2}}={\displaystyle \frac{1}{|\xi |^2}}{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2𝑑x+o(1)={\displaystyle \frac{1}{|\xi |^2}}{\displaystyle _^N}|u^\lambda (x)|^2𝑑x+o(1).`$
On the other hand, from \[14, Proof of Lemma 3.4\] we have
(70) $`\mu ^2{\displaystyle _{|x|>\frac{|\xi |}{2\mu }}}{\displaystyle \frac{|u^\lambda |^2}{|\mu x+\xi |^2}}𝑑x`$ $`\kappa (u^\lambda )^2\mu ^{2N}{\displaystyle _{|x\xi |\frac{|\xi |}{2}}}\left|w_1^\lambda \left({\displaystyle \frac{x\xi }{\mu }}\right)\right|^2{\displaystyle \frac{dx}{|x|^2}}`$
$`=O\left(\mu ^{\nu _\lambda (N2)}\right)=o(\mu ^2).`$
From (68), (69), and (70) we deduce that
$$_^N\frac{|u^\lambda |^2}{|x+\xi |^2}𝑑x=\frac{\mu ^2}{|\xi |^2}_^N|u^\lambda (x)|^2𝑑x+o(\mu ^2).$$
For $`\lambda =\frac{N(N4)}{4}`$, from (25) and \[14, Proof of Lemma 3.4\] we deduce that
(71) $`{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2{\displaystyle \frac{dx}{|\mu x+\xi |^2}}={\displaystyle \frac{1}{|\xi |^2}}{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2𝑑x+O(1).`$
On the other hand, from (24) we obtain
(72) $`{\displaystyle _{|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2𝑑x`$ $`={\displaystyle _{1<|x|<\frac{|\xi |}{2\mu }}}|u^\lambda (x)|^2𝑑x+{\displaystyle _{|x|<1}}|u^\lambda (x)|^2𝑑x`$
$`=\kappa _{\mathrm{}}(u^\lambda )^2{\displaystyle _{1<|x|<\frac{|\xi |}{2\mu }}}|x|^N𝑑x+O\left({\displaystyle _{1<|x|<\frac{|\xi |}{2\mu }}}|x|^{N2\alpha }𝑑x\right)+O(1)`$
$`=\kappa _{\mathrm{}}(u^\lambda )^2|\mathrm{ln}\mu |+O(1).`$
Arguing as above (see (70)), we obtain
(73) $`\mu ^2{\displaystyle _{|x|>\frac{|\xi |}{2\mu }}}{\displaystyle \frac{|u^\lambda |^2}{|\mu x+\xi |^2}}𝑑x=O(\mu ^2).`$
Gathering (68), (71), (72) and (73) we deduce that
$$_^N\frac{|u_\mu ^\lambda |^2}{|x+\xi |^2}𝑑x=\kappa _{\mathrm{}}(u^\lambda )^2\frac{\mu ^2|\mathrm{ln}\mu |}{|\xi |^2}+o(\mu ^2|\mathrm{ln}\mu |).$$
For $`\lambda >\frac{N(N4)}{4}`$, in view of (24) we have that
(74) $`{\displaystyle _^N}`$ $`{\displaystyle \frac{|u_\mu ^\lambda |^2}{|x+\xi |^2}}dx=\mu ^{\nu _\lambda (N2)}\left[\kappa _{\mathrm{}}^2(u^\lambda ){\displaystyle _^N}{\displaystyle \frac{1}{|x|^2|x\xi |^{(N2)(1+\nu _\lambda )}}}+o(1)\right]dx`$
As observed in \[14, Proof of Lemma 3.4\], the function
$$\phi (\xi ):=_^N\frac{dx}{|x|^2|x\xi |^{(N2)(1+\nu _\lambda )}}$$
can be written as
(75)
$$\phi (\xi )=|\xi |^{\sqrt{(N2)^24\lambda }}\phi (\xi /|\xi |)=|\xi |^{\sqrt{(N2)^24\lambda }}\phi (e_1).$$
(74) and (75) yield the required estimate for $`\lambda >\frac{N(N4)}{4}`$.
Let us now study the interaction between two minimizers of (19), i.e. functions $`z_\mu ^\lambda `$ in (20), centered at different points as $`\mu 0`$. To this aim we note that a direct calculation yields
(76) $`z_1^\lambda (x)=|x|^{\frac{N2}{2}(1+\nu _\lambda )}\left[\alpha _{\lambda ,N}+O(|x|^\alpha )\right],`$
(77) $`z_1^\lambda (x)=|x|^{\frac{N+2}{2}\nu _\lambda \frac{N2}{2}}x\left[\alpha _{\lambda ,N}\frac{N2}{2}(1+\nu _\lambda )+O(|x|^\alpha )\right],`$
for all $`0<\alpha 2\nu _\lambda `$. From (20), (76), and (77), it is easy to deduce the following result.
###### Lemma 5.2.
For any $`\lambda (0,(N2)^2/4)`$ and $`\xi ,\zeta ^N`$, $`\xi 0`$, there holds
$$_^N\frac{z_\mu ^\lambda (x)z_\mu ^\lambda (x+\xi )}{|x+\zeta |^2}𝑑x=\mu ^{\sqrt{(N2)^24\lambda }}\left[\alpha _{\lambda ,N}^2_^N\frac{dx}{|x|^{(1+\nu _\lambda )\frac{N2}{2}}|x+\xi |^{(1+\nu _\lambda )\frac{N2}{2}}|x+\zeta |^2}+o(1)\right]$$
and
$`{\displaystyle _^N}z_\mu ^\lambda (x)z_\mu ^\lambda (x+\xi )𝑑x=\mu ^{\sqrt{(N2)^24\lambda }}\left[\alpha _{\lambda ,N}^2\frac{(N2)^2}{4}(1+\nu _\lambda )^2\gamma _{\lambda ,N}|\xi |^{\sqrt{(N2)^24\lambda }}+o(1)\right]`$
as $`\mu 0`$, where
$$\gamma _{\lambda ,N}=_^N\frac{x(x+e_1)dx}{|x|^{\frac{N+2}{2}+\nu _\lambda \frac{N2}{2}}|xe_1|^{\frac{N+2}{2}+\nu _\lambda \frac{N2}{2}}},e_1=(1,0,\mathrm{},0)^N.$$
## 6. Comparison between concentration levels for $`J`$ and proof of Theorem 1.3
In order to compare the level $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ with the level $`S_{\mathrm{circ}}(\lambda _0)`$ of possible concentration at $`0`$, we need the following lemma, which states that the infimum in (54) is achieved if $`N4`$. Such a result does not come unexpected, since it can be seen as the analogue of Lemma 2.1 when $`k=\mathrm{}`$; indeed when $`k`$ becomes larger and larger, assumption $`S_k(\lambda )<k^{2/N}S`$ of Lemma 2.1 is weakened till it is no more needed in the limiting problem corresponding to singularities spread over circles.
###### Lemma 6.1.
For any $`\lambda _0(\mathrm{},(N2)^2/4)`$ and $`N4`$, the infimum in (54) is achieved.
Proof. Hardy’s and Sobolev’s inequalities imply that $`S_{\mathrm{circ}}(\lambda _0)\left(1\frac{4\lambda _0}{(N2)^2}\right)S>0`$. Let $`\{u_n\}_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ be a minimizing sequence such that $`_^N|u_n|^2^{}=1`$. By virtue of the Ekeland’s variational principle we can assume that $`\{u_n\}_n`$ is a Palais-Smale sequence for the functional
$$F(u)=\frac{1}{2}_^N|u|^2𝑑x\frac{\lambda _0}{2}_^N\frac{|u(x)|^2}{|x|^2}𝑑x\frac{S_{\mathrm{circ}}(\lambda _0)}{2^{}}_^N|u(x)|^2^{}𝑑x,u𝒟_{\mathrm{circ}}^{1,2}(^N),$$
i.e.
$$\underset{n\mathrm{}}{lim}F(u_n)=\frac{S_{\mathrm{circ}}(\lambda _0)}{N}\text{ in }\text{and}\underset{n\mathrm{}}{lim}F^{}(u_n)=0\text{ in the dual space }(𝒟_{\mathrm{circ}}^{1,2}(^N))^{}.$$
Let
$$w_n(x)=\sigma _n^{\frac{N2}{2}}u_n(\sigma _n^1x)$$
where $`\sigma _n`$ is chosen in such a way that
(78)
$$_{B(0,1)}|w_n(x)|^2^{}𝑑x=_{B(0,\sigma _n^1)}|u_n(x)|^2^{}𝑑x=\frac{1}{2}.$$
Scaling invariance ensures that $`\{w_n\}_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ is also a minimizing sequence for (54),
(79)
$$_^N|w_n(x)|^2𝑑x=1,\underset{n+\mathrm{}}{lim}_^N\left(|w_n(x)|^2\lambda _0\frac{|w_n(x)|^2}{|x|^2}\right)𝑑x=S_{\mathrm{circ}}(\lambda _0),$$
and
(80)
$$F^{}(w_n)0\text{ in the dual space }(𝒟_{\mathrm{circ}}^{1,2}(^N))^{}.$$
Since $`\{w_n\}_n`$ is bounded in $`𝒟^{1,2}(^N)`$, up to a subsequence $`w_n`$ converges to $`w`$ weakly in $`𝒟^{1,2}(^N)`$ and almost everywhere. Pointwise convergence implies that $`w𝒟_{\mathrm{circ}}^{1,2}(^N)`$. From the Concentration Compactness Principle by P. L. Lions and taking into account that, when $`N4`$, $`𝕊𝕆(2)\times 𝕊𝕆(N2)`$-invariant functions can concentrate only at $`0`$ and at $`\mathrm{}`$, as already pointed out in the proof of Theorem 4.2, we have
(81) $`|w_n|^2d\mu |w|^2+\mu _0\delta _0,|w_n|^2^{}d\nu =|w|^2^{}+\nu _0\delta _0,\text{and}{\displaystyle \frac{w_n^2}{|x|^2}}d\gamma ={\displaystyle \frac{w^2}{|x|^2}}+\gamma _0\delta _0.`$
The amount of concentration at infinity is quantified by the following numbers
$$\nu _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}|w_n|^2^{}𝑑x,\mu _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}|w_n|^2𝑑x,$$
and
$$\gamma _{\mathrm{}}=\underset{R\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; sup}_{|x|>R}\frac{w_n^2}{|x|^2}𝑑x.$$
Arguing as we did to prove (60) and (63), we can easily obtain
(82)
$$\nu _0^{\frac{2}{N}}(\mu _0\lambda _0\gamma _0)S_{\mathrm{circ}}(\lambda _0)\nu _0$$
and
(83)
$$(\mu _{\mathrm{}}\lambda _0\gamma _{\mathrm{}})S_{\mathrm{circ}}(\lambda _0)\nu _{\mathrm{}}^{\frac{2}{2^{}}}.$$
Step 1: we prove that $`w0`$. By contradiction, assume that $`w0`$. Then from (7981) we deduce
(84)
$$1=\nu _0+\nu _{\mathrm{}}$$
and
(85)
$$S_{\mathrm{circ}}(\lambda _0)=_^N𝑑\mu \lambda _0_^N𝑑\gamma +\mu _{\mathrm{}}\lambda _0\gamma _{\mathrm{}}.$$
Let $`\psi _R`$ be a smooth radial cut-off function such that $`0\psi _R(x)1`$, $`\psi _R(x)=1`$ if $`|x|>2R`$, $`\psi _R(x)=0`$ if $`|x|<R`$, and $`|\psi _R|2/R`$. From (80), we have
$`o(1)`$ $`=F^{}(w_n),w_n\psi _R`$
$`={\displaystyle _^N}|w_n|^2\psi _R+{\displaystyle _^N}w_nw_n\psi _R\lambda _0{\displaystyle _^N}{\displaystyle \frac{w_n\psi _R}{|x|^2}}S_{\mathrm{circ}}(\lambda _0){\displaystyle _^N}|w_n|^2^{}\psi _R.`$
Taking $`lim\; sup`$ as $`n\mathrm{}`$ and limit as $`R+\mathrm{}`$, we find
(86)
$$\mu _{\mathrm{}}\lambda _0\gamma _{\mathrm{}}=S_{\mathrm{circ}}(\lambda _0)\nu _{\mathrm{}}.$$
From (8586) it follows that
(87)
$$_^N𝑑\mu \lambda _0_^N𝑑\gamma =S_{\mathrm{circ}}(\lambda _0)(1\nu _{\mathrm{}}).$$
From (84), (82), and (87), it follows that
(88) $`1\nu _{\mathrm{}}`$ $`=\nu _0S_{\mathrm{circ}}(\lambda _0)^1\nu _0^{\frac{2}{N}}(\mu _0\lambda _0\gamma _0)S_{\mathrm{circ}}(\lambda _0)^1\nu _0^{\frac{2}{N}}\left({\displaystyle _^N}𝑑\mu \lambda _0{\displaystyle _^N}𝑑\gamma \right)`$
$`=\nu _0^{\frac{2}{N}}(1\nu _{\mathrm{}})=(1\nu _{\mathrm{}})^{1+\frac{2}{N}}.`$
On the other hand, from (78) we have
$$_{^N\{B(0,R)\}}|w_n|^2^{}𝑑x=1_{B(0,R)}|w_n|^2^{}𝑑x1_{B(0,1)}|w_n|^2^{}𝑑x=\frac{1}{2}\text{for all }R>1,$$
hence
(89)
$$\nu _{\mathrm{}}\frac{1}{2}.$$
From (8889) we deduce that $`\nu _{\mathrm{}}=0`$. From (84) it follows that $`\nu _0=1`$. Therefore (78) implies
$$\frac{1}{2}=_{B(0,1)}|w_n|^2^{}\underset{n\mathrm{}}{}d\nu (B(0,1))=1,$$
thus giving rise to a contradiction.
Step 2: we prove that
$$_^N|w|^2^{}=1\text{and}_^N|w|^2𝑑x\lambda _0_^N\frac{|w|^2}{|x|^2}𝑑x=S_{\mathrm{circ}}(\lambda _0).$$
Let $`\rho =_^N|w|^2^{}`$. From (81), we have that $`1=\rho +\nu _0+\nu _{\mathrm{}}`$ and, in view of step 1, $`\rho (0,1]`$, i.e. $`\nu _0+\nu _{\mathrm{}}[0,1)`$. Since
$`S_{\mathrm{circ}}(\lambda _0)`$ $`={\displaystyle _^N}𝑑\mu \lambda _0{\displaystyle _^N}𝑑\gamma +\mu _{\mathrm{}}\lambda _0\gamma _{\mathrm{}}`$
$`{\displaystyle _^N}|w|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{|w|^2}{|x|^2}}𝑑x+\mu _0\lambda _0\gamma _0+\mu _{\mathrm{}}\lambda _0\gamma _{\mathrm{}},`$
from (82), (83), and concavity of the function $`tt^{2/2^{}}`$ we deduce
$`{\displaystyle _^N}|w|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{|w|^2}{|x|^2}}𝑑x`$ $`S_{\mathrm{circ}}(\lambda _0)\left(1\nu _0^{\frac{2}{2^{}}}\nu _{\mathrm{}}^{\frac{2}{2^{}}}\right)S_{\mathrm{circ}}(\lambda _0)\left(1\nu _0\nu _{\mathrm{}}\right)^{\frac{2}{2^{}}}`$
$`=S_{\mathrm{circ}}(\lambda _0)\rho ^{2/2^{}}{\displaystyle _^N}|w|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{|w|^2}{|x|^2}}𝑑x.`$
Hence all the above inequalities are indeed equalities; in particular $`\left(1\nu _0^{\frac{2}{2^{}}}\nu _{\mathrm{}}^{\frac{2}{2^{}}}\right)=\left(1\nu _0\nu _{\mathrm{}}\right)^{\frac{2}{2^{}}}`$ which is possible only when $`\nu _0=\nu _{\mathrm{}}=0`$. Therefore $`\rho =1`$ and
$$_^N|w|^2𝑑x\lambda _0_^N\frac{|w|^2}{|x|^2}𝑑x=S_{\mathrm{circ}}(\lambda _0),$$
i.e. $`w`$ attains the infimum.
We now provide a sufficient condition for the infimum in (5) to stay below the level $`S_{\mathrm{circ}}(\lambda _0)`$, at which possible concentration at $`0`$ can occur.
###### Lemma 6.2.
Let $`\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m`$, $`r_1,r_2,\mathrm{},r_m^+`$ satisfy (10). Then
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_{\mathrm{circ}}(\lambda _0).$$
Proof. From Lemma 6.1, we have that $`S_{\mathrm{circ}}(\lambda _0)`$ is attained by some $`u^{\lambda _0}𝒟_{\mathrm{circ}}^{1,2}(^N)`$. By homogeneity of the Rayleigh quotient, we can assume $`|u^{\lambda _0}|^2^{}=1`$. Moreover, the function $`v^{\lambda _0}=S_{\mathrm{circ}}(\lambda _0)^{1/(2^{}2)}|u^{\lambda _0}|`$ is a nonnegative solution to (16), hence we can apply Lemma 5.1 to study the behavior of $`_^N\frac{|u_\mu ^{\lambda _0}|^2}{|x+\xi |^2}𝑑x`$ as $`\mu 0`$, where $`u_\mu ^{\lambda _0}`$ are defined in (67). Hence for some positive constant $`\stackrel{~}{\kappa }`$
$$\begin{array}{c}_^N\left(\text{ }_{S_r_{\mathrm{}}}\frac{|u_\mu ^{\lambda _0}(y)|^2}{|xy|^2}𝑑\sigma (x)\right)𝑑y\hfill \\ \hfill =\{\begin{array}{cc}\frac{\mu ^2}{r_{\mathrm{}}^2}_^N|u_1^{\lambda _0}|^2𝑑x+o\left(\mu ^2\right)\hfill & \text{if }\lambda _0<\frac{N(N4)}{4},\hfill \\ \stackrel{~}{\kappa }^2\frac{\mu ^2|\mathrm{ln}\mu |}{r_{\mathrm{}}^2}+o(\mu ^2|\mathrm{ln}\mu |)\hfill & \text{if }\lambda _0=\frac{N(N4)}{4},\hfill \\ \stackrel{~}{\kappa }^2\beta _{\lambda _0,N}\mu ^{\sqrt{(N2)^24\lambda _0}}|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}+o\left(\mu ^{\sqrt{(N2)^24\lambda _0}}\right)\hfill & \text{if }\lambda _0>\frac{N(N4)}{4},\hfill \end{array}\end{array}$$
as $`\mu 0`$, where $`\beta _{\lambda _0,N}`$ is defined in Lemma 5.1. Therefore
(90)
$$\begin{array}{c}S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\hfill \\ \hfill _^N|u_\mu ^{\lambda _0}|^2𝑑y\lambda _0_^N\frac{|u_\mu ^{\lambda _0}(y)|^2}{|y|^2}𝑑y\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}_^N\left(\text{ }_{S_r_{\mathrm{}}}\frac{|u_\mu ^{\lambda _0}(y)|^2}{|xy|^2}𝑑\sigma (x)\right)𝑑y\\ \hfill =S_{\mathrm{circ}}(\lambda _0)\{\begin{array}{cc}\mu ^2\left(_^N|u_1^{\lambda _0}|^2\right)\left(\underset{\mathrm{}=1}{\overset{m}{}}\frac{\mathrm{\Lambda }_{\mathrm{}}}{r_{\mathrm{}}^2}+o(1)\right)\hfill & \text{if }\lambda _0<\frac{N(N4)}{4},\hfill \\ \mu ^2|\mathrm{ln}\mu |\stackrel{~}{\kappa }^2\left(\underset{\mathrm{}=1}{\overset{m}{}}\frac{\mathrm{\Lambda }_{\mathrm{}}}{r_{\mathrm{}}^2}+o(1)\right)\hfill & \text{if }\lambda _0=\frac{N(N4)}{4},\hfill \\ \mu ^{\sqrt{(N2)^24\lambda _0}}\stackrel{~}{\kappa }^2\beta _{\lambda _0,N}\left(\underset{\mathrm{}=1}{\overset{m}{}}\frac{\mathrm{\Lambda }_{\mathrm{}}}{|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}}+o(1)\right)\hfill & \text{if }\lambda _0>\frac{N(N4)}{4},\hfill \end{array}\end{array}$$
as $`\mu 0`$. Taking $`\mu `$ sufficiently small, assumption (10) yields $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_{\mathrm{circ}}(\lambda _0)`$.
Proof of Theorem 1.3. Let $`\{u_n\}_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ be a minimizing sequence for (5). From the homogeneity of the quotient there is no restriction requiring $`u_n_{L^2^{}(^N)}=1`$. Moreover from Ekeland’s variational principle we can assume that $`\{u_n\}_n𝒟_{\mathrm{circ}}^{1,2}(^N)`$ is a Palais-Smale sequence, more precisely $`J^{}(u_n)0`$ in $`(𝒟_{\mathrm{circ}}^{1,2}(^N))^{}`$ and $`J(u_n)\frac{1}{N}S_{\mathrm{circ}}(\lambda _0\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. From assumption (9) we deduce that
(91)
$$S_{\mathrm{circ}}(\lambda _0)S_{\mathrm{circ}}\left(\lambda _0+\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}\right).$$
From Lemma 6.2 and (91), it follows that the level of the minimizing Palais-Smale sequence satisfies assumption (55). Hence from Theorem 4.2, $`\{u_n\}_n`$ has a subsequence strongly converging to some $`u_0𝒟_{\mathrm{circ}}^{1,2}(^N)`$ such that $`J(u_0)=\frac{1}{N}S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. Hence $`u_0`$ achieves the infimum in (5). Since $`J`$ is even, also $`|u_0|`$ is a minimizer in (5) and then $`v_0=S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{1/(2^{}2)}|u_0|`$ is a nonnegative solution to equation (4). The maximum principle implies the positivity outside singular circles of such a solution.
## 7. Limit of $`S_k(\lambda )`$ as $`k\mathrm{}`$.
Since $`𝒟_{\mathrm{circ}}^{1,2}(^N)𝒟_k^{1,2}(^N)`$, for any $`\lambda (\mathrm{},(n2)^2/4)`$, there holds
(92)
$$0<S(\lambda )S_k(\lambda )S_{\mathrm{circ}}(\lambda ).$$
From (92) and Lemma 2.1, it follows easily the following result.
###### Lemma 7.1.
Let $`\lambda (\mathrm{},(n2)^2/4)`$ and $`N4`$. Then there exists $`\overline{k}=\overline{k}(\lambda ,N)`$ such that $`S_k(\lambda )`$ is achieved for all $`k\overline{k}`$.
Let us now study the limit of $`S_k(\lambda )`$ as $`k\mathrm{}`$. Theorem 7.3 provides convergence of $`S_k(\lambda )`$ to $`S_{\mathrm{circ}}(\lambda )`$. To prove it we will need the following proposition.
###### Proposition 7.2.
Let $`\lambda (\mathrm{},(n2)^2/4)`$ and let $`\{w_k\}_k`$ be a sequence in $`𝒟^{1,2}(^N)`$ such that $`w_k𝒟_k^{1,2}(^N)`$,
(93)
$$_^N|w_k|^2^{}=1,Q_\lambda (w_k)=S_k(\lambda ),$$
and $`w_k`$ converges weakly to $`0`$ in $`𝒟^{1,2}(^N)`$ (at least along a subsequence). Then, for any $`r>0`$ and $`\epsilon (r,r)`$, there exists $`\rho `$ such that $`0<|\rho |<|\epsilon |`$ and , for a subsequence,
$`\text{either}{\displaystyle _{B(0,r+\rho )}}|w_k|^20,{\displaystyle _{B(0,r+\rho )}}|w_k|^2^{}0,\text{and }{\displaystyle _{B(0,r+\rho )}}{\displaystyle \frac{|w_k|^2}{|x|^2}}0,`$
$`\text{or}{\displaystyle _{^NB(0,r+\rho )}}|w_k|^20,{\displaystyle _{^NB(0,r+\rho )}}|w_k|^2^{}0,\text{and }{\displaystyle _{^NB(0,r+\rho )}}{\displaystyle \frac{|w_k|^2}{|x|^2}}0.`$
Proof. An analogous result is proved in for minimizing sequences of quotient (19). Since the proof of Proposition 7.2 is similar, we will be sketchy.
Let $`\epsilon (0,r)`$ (the proof for $`\epsilon `$ negative is similar). Since
$$_r^{r+\epsilon }𝑑\rho _{B(0,\rho )}|w_k|^2=_{B(0,r+\epsilon )B(0,r)}|w_k|^2,$$
we can choose $`\rho (0,\epsilon )`$ such that, for infinitely many $`k`$’s (i.e. along a subsequence still denoted as $`\{w_k\}_k`$)
(94)
$$_{B(0,r+\rho )}|w_k|^2\frac{2}{\epsilon }_{B(0,r+\epsilon )B(0,r)}|w_k|^2.$$
From the uniform bound of $`S_k(\lambda )`$ (see (92)) and equivalence of $`Q_\lambda `$ to $`𝒟^{1,2}(^N)`$-norm, it follows that
(95)
$$_{B(0,r+\epsilon )B(0,r)}|w_k|^2\mathrm{const}Q_\lambda (w_k)=\mathrm{const}S_k(\lambda )\mathrm{const}.$$
From (9495), it follows that $`\left\{w_k|_{B(0,r+\rho )}\right\}_k`$ is bounded in $`H^1\left(B(0,r+\rho )\right)`$. By compactness of the embedding $`H^1\left(B(0,r+\rho )\right)H^{1/2}\left(B(0,r+\rho )\right)`$ and weak convergence to $`0`$, we conclude that, up to subsequence, $`\left\{w_k|_{B(0,r+\rho )}\right\}_k`$ converges strongly to $`0`$ in $`H^{1/2}\left(B(0,r+\rho )\right)`$.
Let $`w_k^1`$ (respectively $`w_k^2`$) be the harmonic functions in $`B(0,r+\epsilon )B(0,r+\rho )`$ (respectively $`B(0,r+\rho )B(0,r\epsilon )`$) such that
$$\{\begin{array}{cc}w_k^1=w_k\hfill & \text{on }B(0,r+\rho ),\hfill \\ w_k^1=0\hfill & \text{on }B(0,r+\epsilon ),\hfill \end{array}\{\begin{array}{cc}w_k^2=w_k\hfill & \text{on }B(0,r+\rho ),\hfill \\ w_k^2=0\hfill & \text{on }B(0,r\epsilon ).\hfill \end{array}$$
By continuity of $`\mathrm{\Delta }^1`$ we have that $`w_k^10`$ strongly in $`H^1(B(0,r+\epsilon )B(0,r+\rho ))`$ and $`w_k^20`$ strongly in $`H^1(B(0,r+\rho )B(0,r\epsilon ))`$. Moreover symmetry properties of boundary data ensure that $`w_k^1,w_k^2𝒟_k^{1,2}(^N)`$. Let us now set
$$\{\begin{array}{cc}u_k^1=w_k\hfill & \text{in }B(0,r+\rho ),\hfill \\ u_k^1=w_k^1\hfill & \text{in }B(0,r+\epsilon )B(0,r+\rho ),\hfill \\ u_k^1=0\hfill & \text{in }^NB(0,r+\epsilon ),\hfill \end{array}\{\begin{array}{cc}u_k^2=w_k\hfill & \text{in }^NB(0,r+\rho ),\hfill \\ u_k^2=w_k^2\hfill & \text{in }B(0,r+\rho )B(0,r\epsilon ),\hfill \\ u_k^2=0\hfill & \text{in }B(0,r\epsilon ).\hfill \end{array}$$
Direct computations yield
$`Q_\lambda (u_k^1)=S_k(\lambda ){\displaystyle _{B(0,r+\rho )}}|w_k|^2^{}𝑑x+o(1),`$
$`Q_\lambda (u_k^2)=S_k(\lambda ){\displaystyle _{^NB(0,r+\rho )}}|w_k|^2^{}𝑑x+o(1),`$
$`Q_\lambda (u_k^1)+Q_\lambda (u_k^2)=S_k(\lambda ){\displaystyle _^N}|w_k|^2^{}𝑑x+o(1)=S_k(\lambda )+o(1)=Q_\lambda (w_k)+o(1),`$
$`{\displaystyle _^N}|u_k^1|^2^{}𝑑x+{\displaystyle _^N}|u_k^2|^2^{}𝑑x={\displaystyle _^N}|w_k|^2^{}𝑑x+o(1).`$
We claim that either $`Q_\lambda (u_k^1)0`$ or $`Q_\lambda (u_k^2)0`$ along some subsequence. Indeed, assume that $`Q_\lambda (u_k^1)\to ̸0`$ along any subsequence, i.e. $`Q_\lambda (u_k^1)`$ stays bounded below away from $`0`$. From above, (93) and (22), it follows
$`{\displaystyle \frac{Q_\lambda (u_k^1)}{u_k^1_2^{}^2}}`$ $`={\displaystyle \frac{Q_\lambda (w_k)Q_\lambda (u_k^2)+o(1)}{\left(w_k_2^{}^2^{}u_k^2_2^{}^2^{}+o(1)\right)^{2/2^{}}}}S_k(\lambda ){\displaystyle \frac{Q_\lambda (w_k)Q_\lambda (u_k^2)+o(1)}{\left(Q_\lambda (w_k)^{2^{}/2}Q_\lambda (u_k^2)^{2^{}/2}+o(1)\right)^{2/2^{}}}}<S_k(\lambda )`$
in contradiction with (22), unless $`Q_\lambda (u_k^2)0`$ along some subsequence. The claim is thereby proved. The statement of the proposition follows from equivalence to norm of $`Q_\lambda `$, Hardy’s and Sobolev’s inequalities.
###### Theorem 7.3.
Let $`\lambda (\mathrm{},(n2)^2/4)`$ and $`N4`$. Then $`lim_{k+\mathrm{}}S_k(\lambda )=S_{\mathrm{circ}}(\lambda )`$.
Proof. From Lemma 7.1 we know that, for $`k`$ sufficiently large, $`S_k(\lambda )`$ is achieved, hence there exists some $`u_k𝒟_k^{1,2}(^N)`$ such that
$$_^N|u_k|^2^{}=1\text{and}Q_\lambda (u_k)=S_k(\lambda ).$$
From the uniform bound of $`S_k(\lambda )`$ (see (92)) and equivalence of $`Q_\lambda `$ to $`𝒟^{1,2}(^N)`$-norm, it follows that $`\{u_k\}_k`$ is bounded in $`𝒟^{1,2}(^N)`$. Let us set
(96)
$$\stackrel{~}{u}_k(x)=R_k^{\frac{N2}{2}}u_k\left(\frac{x}{R_k}\right)\text{and}v_k(x)=(S_k(\lambda ))^{\frac{1}{2^{}2}}\stackrel{~}{u}_k(x)$$
where $`R_k`$ is chosen such that
$$_{B(0,R_k)}\left[|u_k(x)|^2\lambda \frac{|u_k(x)|^2}{|x|^2}\right]𝑑x=_{^NB(0,R_k)}\left[|u_k(x)|^2\lambda \frac{|u_k(x)|^2}{|x|^2}\right]𝑑x=\frac{1}{2}S_k(\lambda ).$$
Invariance by scaling yields
(97) $`{\displaystyle _^N}|\stackrel{~}{u}_k|^2^{}=1,`$ $`Q_\lambda (\stackrel{~}{u}_k)=S_k(\lambda ),`$
(98) $`{\displaystyle _^N}|v_k|^2^{}=(S_k(\lambda ))^{\frac{N}{2}},`$ $`Q_\lambda (v_k)=(S_k(\lambda ))^{\frac{N}{2}},`$
and
(99)
$$_{B(0,1)}[|v_k(x)|^2\lambda \frac{|v_k(x)|^2}{|x|^2}]dx=_{^NB(0,1)}[v_k(x)|^2\lambda \frac{|v_k(x)|^2}{|x|^2}],dx=\frac{1}{2}(S_k(\lambda ))^{\frac{N}{2}}.$$
Invariance by scaling also implies that $`\{\stackrel{~}{u}_k\}_k`$ is bounded in $`𝒟^{1,2}(^N)`$, hence there exists a subsequence (still denoted as $`\{\stackrel{~}{u}_k\}_k`$) weakly converging to some $`\stackrel{~}{u}_0`$ in $`𝒟^{1,2}(^N)`$.
Claim 1. We claim that $`\stackrel{~}{u}_00`$. Assume by contradiction that $`\stackrel{~}{u}_00`$. Using Proposition 7.2 for sequence $`\stackrel{~}{u}_k`$ with $`r=1`$ and $`\epsilon =\pm \frac{1}{4}`$ and taking into account (99), (92), and (96), we deduce that there exist $`\rho ^+(0,1/4)`$ and $`\rho ^{}(1/4,0)`$ such that, up to a subsequence,
(100) $`{\displaystyle _{B(0,1+\rho ^{})}}|v_k|^20,{\displaystyle _{B(0,1+\rho ^{})}}|v_k|^2^{}0,\text{and }{\displaystyle _{B(0,1+\rho ^{})}}{\displaystyle \frac{|v_k|^2}{|x|^2}}0,`$
(101) $`{\displaystyle _{^NB(0,1+\rho ^+)}}|v_k|^20,{\displaystyle _{^NB(0,1+\rho ^+)}}|v_k|^2^{}0,\text{and }{\displaystyle _{^NB(0,1+\rho ^+)}}{\displaystyle \frac{|v_k|^2}{|x|^2}}0.`$
Note that weak convergence of $`\stackrel{~}{u}_k0`$ in $`𝒟^{1,2}(^N)`$, (92), and (96), imply weak convergence of $`v_k0`$ in $`𝒟^{1,2}(^N)`$. Let $`\eta `$ be a smooth radial cut off function such that $`0\eta 1`$, $`\eta (x)1`$ for $`1+\rho ^{}|x|1+\rho ^+`$ and $`\eta (x)0`$ for $`|x|[3/4,5/4]`$. Set $`\stackrel{~}{v}_k:=\eta v_k`$. Clearly $`\stackrel{~}{v}_k𝒟_k^{1,2}(^N)`$. By choice of $`\eta `$ and (100101) we have
$`Q_\lambda (\stackrel{~}{v}_k)=Q_\lambda (v_k)+o(1),\stackrel{~}{v}_kv_k_{𝒟^{1,2}(^N)}=o(1),{\displaystyle _^N}|\stackrel{~}{v}_k|^2^{}={\displaystyle _^N}|v_k|^2^{}+o(1).`$
Let us define
$$f(u):=\frac{1}{2}_{B(0,5/4)B(0,3/4)}|u|^2\frac{1}{2^{}}_{B(0,5/4)B(0,3/4)}|u|^2^{},uH_0^1(B(0,5/4)B(0,3/4)).$$
From (100101), (92), and (98), it is easy to verify that
$$f^{}(\stackrel{~}{v}_k)0\text{in }(𝒟^{1,2}(^N))^{}\text{and}f(\stackrel{~}{v}_k)=\frac{1}{N}(S_k(\lambda ))^{N/2}+o(1)\frac{1}{N}(S_{\mathrm{circ}}(\lambda ))^{N/2}+o(1),$$
i.e. $`\stackrel{~}{v}_k`$ is a Palais-Smale sequence for $`f`$ in $`H_0^1(B(0,5/4)B(0,3/4))`$. From Struwe’s representation lemma for diverging Palais-Smale sequences \[27, Theorem III.3.1\], we deduce the existence of an integer $`M`$, $`M`$ sequences of points $`\{x_k^i\}_kB(0,5/4)B(0,3/4)`$ and $`M`$ sequences of radii $`\{R_k^i\}_k`$, $`i=1,\mathrm{},M`$, such that $`lim_kR_k^i=+\mathrm{}`$ and
(102)
$$\stackrel{~}{v}_k(x)=\underset{i=1}{\overset{M}{}}(R_k^i)^{\frac{N2}{2}}\stackrel{~}{v}_0(R_k^i(xx_k^i))+_k(x)\text{where }_k0\text{ in }𝒟^{1,2}(^N),$$
$`\stackrel{~}{v}_0=w_1^0`$ and $`w_1^0`$ is the Talenti-Aubin function in (17). Let us consider the sequence $`\{x_k^1\}_k`$; up to subsequence we can assume that it converges to some point $`x^1\overline{B(0,5/4)B(0,3/4)}`$. Let us write $`x^1=(z_1,y_1)^2\times ^{N2}`$ where $`z_1=|z_1|e^{\overline{\theta }\sqrt{1}}`$.
Let us first assume that $`|z_1|0`$. We fix $`J`$ and for any $`i=1,2,\mathrm{},J`$ we set
$$S_i=\{(z,y)=(|z|e^{\theta \sqrt{1}},y)^2\times ^{N2}:\overline{\theta }\frac{(2i1)\pi }{J}<\theta <\overline{\theta }+\frac{(2i+1)\pi }{J}\}.$$
Note that $`x^1S_1`$ and there exists $`\delta =\delta (J)>0`$ such that
$$B(x^1,\delta )\{(|z|e^{\theta \sqrt{1}},y)^2\times ^{N2}:\overline{\theta }\frac{(2i1)\pi }{2J}<\theta <\overline{\theta }+\frac{(2i+1)\pi }{2J}\}.$$
Choose $`\overline{k}=\overline{k}(\delta )`$ such that for all $`k\overline{k}`$
$$x_k^1=(z_k^1,y_k^1)B(x^1,\frac{\delta }{2})\text{and}(R_k^1)^1<\frac{\delta }{2}.$$
Moreover, if $`\overline{k}`$ is chosen sufficiently large, for each $`i=1,2,\mathrm{},J`$ it is possible to find $`\tau _k^i_k`$ such that $`(\tau _k^iz_1,y_1)`$ stays in the middle half of $`S_i`$, i.e.
$$(\tau _k^iz_1,y_1)\{(z,y)=(|z|e^{\theta \sqrt{1}},y)^2\times ^{N2}:\overline{\theta }\frac{(2i1)\pi }{2J}<\theta <\overline{\theta }+\frac{(2i+1)\pi }{2J}\}.$$
Hence
$$B((\tau _k^iz_1,y_1),\delta )S_i,(\tau _k^iz_k^1,y_k^1)B((\tau _k^iz_1,y_1),\frac{\delta }{2})$$
and consequently
$$B((\tau _k^iz_k^1,y_k^1),\frac{\delta }{2})S_i$$
which yields
$$B((\tau _k^iz_k^1,y_k^1),(R_k^i)^1)S_i.$$
In particular the $`J`$ balls $`B((\tau _k^iz_k^1,y_k^1),(R_k^i)^1)`$ are disjoint, hence, by symmetry properties of $`\stackrel{~}{v}_k`$ we have that
$`(S_k(\lambda ))^{N/2}+o(1)={\displaystyle _^N}|\stackrel{~}{v}_k|^2^{}`$ $`{\displaystyle \underset{i=1}{\overset{J}{}}}{\displaystyle _{B((\tau _k^iz_k^1,y_k^1),(R_k^i)^1)}}|\stackrel{~}{v}_k|^2^{}={\displaystyle \underset{i=1}{\overset{J}{}}}{\displaystyle _{B(x_k^1,(R_k^i)^1)}}|\stackrel{~}{v}_k|^2^{}.`$
On the other hand from (102) we have, for $`k`$ large,
$`(`$ $`{\displaystyle _{B(x_k^1,(R_k^i)^1)}}|\stackrel{~}{v}_k|^2^{})^{\frac{1}{2^{}}}`$
$`\left({\displaystyle _{B(x_k^1,(R_k^i)^1)}}(R_k^1)^N\left|\stackrel{~}{v}_0(R_k^1(xx_k^1))\right|^2^{}\right)^{\frac{1}{2^{}}}\left({\displaystyle _{B(x_k^1,(R_k^i)^1)}}|_k|^2^{}\right)^{\frac{1}{2^{}}}{\displaystyle \frac{1}{2}}\left({\displaystyle _{B(0,1)}}\stackrel{~}{v}_0^2^{}\right)^{\frac{1}{2^{}}}.`$
Therefore
$$(S_k(\lambda ))^{N/2}+o(1)\frac{J}{2^2^{}}_{B(0,1)}\stackrel{~}{v}_0^2^{}$$
and, in view of (92)
$$(S_{\mathrm{circ}}(\lambda ))^{N/2}\frac{J}{2^2^{}}_{B(0,1)}\stackrel{~}{v}_0^2^{}.$$
Letting $`J+\mathrm{}`$, we find a contradiction. Claim 1 is thereby proved in the case $`|z_1|0`$. The case $`|z_1|=0`$ can be treated exploiting the radial symmetry of functions $`\stackrel{~}{u}_k`$ in the last $`N2`$ variables with a similar argument (even simpler due to the stronger symmetry).
Claim 2. We claim that $`\stackrel{~}{u}_0𝒟_{\mathrm{circ}}^{1,2}(^N)`$. We first note that $`\stackrel{~}{u}_k`$ satisfy the equation $`\mathrm{\Delta }\stackrel{~}{u}_k\lambda \frac{\stackrel{~}{u}_k}{|x|^2}=S_k(\lambda )\stackrel{~}{u}_k^{2^{}1}`$. From (92), we can assume that $`S_k(\lambda )L(0,+\mathrm{})`$ at least for a subsequence. Hence, due to weak convergence of $`\stackrel{~}{u}_k\stackrel{~}{u}_0`$, we can pass to the limit in the equation to find that $`\stackrel{~}{u}_0`$ satisfies the equation $`\mathrm{\Delta }\stackrel{~}{u}_0\lambda \frac{\stackrel{~}{u}_0}{|x|^2}=L\stackrel{~}{u}_0^{2^{}1}`$. By classical regularity theory for elliptic equations, we deduce that $`\stackrel{~}{u}_0`$ is a smooth function outside the origin.
Let $`R>0`$. Assume that there exist $`(z_1,y),(z_2,y)B(0,R)(^2\times ^{N2})`$, $`|z_1|=|z_2|`$, such that $`\stackrel{~}{u}_0(z_1,y)\stackrel{~}{u}_0(z_2,y)`$. Then there exist $`\delta >0`$ such that $`\stackrel{~}{u}_0(x)\stackrel{~}{u}_0(y)`$ for any $`xB((z_1,y),\delta )`$, $`yB((z_2,y),\delta )`$. Let $`0<\epsilon <\frac{1}{2}|B(0,\delta )|`$. Since, up to a subsequence, $`\stackrel{~}{u}_k\stackrel{~}{u}_0`$ a.e. in $`B(0,R)`$, by the Severini-Egorov Theorem, there exists a measurable set $`\mathrm{\Omega }B(0,R)`$ such that $`|\mathrm{\Omega }|<\epsilon `$ and $`\stackrel{~}{u}_k\stackrel{~}{u}_0`$ uniformly in $`B(0,R)\mathrm{\Omega }`$. Hence for $`k`$ large, $`\stackrel{~}{u}_k(x)\stackrel{~}{u}_k(y)`$ for any $`xB((z_1,y),\delta )\mathrm{\Omega }`$, $`yB((z_2,y),\delta )\mathrm{\Omega }`$. On the other hand, if $`k`$ is large enough, there exists $`\tau _k_k`$ such that
$$\left|\tau _k(B((z_1,y),\delta ))\mathrm{}B((z_2,y),\delta )\right|<\epsilon ,$$
where $`\mathrm{}`$ denotes the symmetric difference of sets. Hence
$$\left|\left(\tau _k(B((z_1,y),\delta ))B((z_2,y),\delta )\right)\mathrm{\Omega }\right|>|B(0,\delta )|2\epsilon >0.$$
In particular the set $`\left(\tau _k(B((z_1,y),\delta ))B((z_2,y),\delta )\right)\mathrm{\Omega }`$ has non-zero measure. If $`(z,y)\left(\tau _k(B((z_1,y),\delta ))B((z_2,y),\delta )\right)\mathrm{\Omega }`$, then $`z=\tau _k\stackrel{~}{z}`$ with $`(\stackrel{~}{z},y)B((z_1,y),\delta )`$ and by symmetry of $`\stackrel{~}{u}_k`$, $`\stackrel{~}{u}_k(z,y)=\stackrel{~}{u}_k(\stackrel{~}{z},y)`$, thus giving a contradiction. Hence $`\stackrel{~}{u}_0`$ is invariant by the $`𝕊𝕆(2)`$-action on the first two variables on $`B(0,R)`$ for any $`R`$. Invariance by the $`𝕊𝕆(N2)`$-action on the last $`(N2)`$ variables follows easily from pointwise convergence. Then we conclude that $`\stackrel{~}{u}_0𝒟_{\mathrm{circ}}^{1,2}(^N)`$.
Hence we have proved that, up to a subsequence, $`\stackrel{~}{u}_k\stackrel{~}{u}_0`$ in $`𝒟^{1,2}(^N)`$, with $`\stackrel{~}{u}_0𝒟_{\mathrm{circ}}^{1,2}(^N)\{0\}`$. Weak convegence yields
(103)
$$Q_\lambda (\stackrel{~}{u}_k)=Q_\lambda (\stackrel{~}{u}_0)+Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)+o(1)$$
while Brezis-Lieb Lemma implies
(104)
$$\stackrel{~}{u}_k_2^{}^2^{}=\stackrel{~}{u}_0_2^{}^2^{}+\stackrel{~}{u}_k\stackrel{~}{u}_0_2^{}^2^{}+o(1).$$
From (103), (104), (97), and (22) we have
$$S_k(\lambda )\frac{Q_\lambda (\stackrel{~}{u}_0)}{\stackrel{~}{u}_0_2^{}^2}S_k(\lambda )\frac{Q_\lambda (\stackrel{~}{u}_k)Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)+o(1)}{\left(Q_\lambda (\stackrel{~}{u}_k)^{2^{}/2}Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)^{2^{}/2}+o(1)\right)^{2/2^{}}}.$$
Hence
$$\frac{Q_\lambda (\stackrel{~}{u}_k)Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)+o(1)}{\left(Q_\lambda (\stackrel{~}{u}_k)^{2^{}/2}Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)^{2^{}/2}+o(1)\right)^{2/2^{}}}1.$$
Since $`Q_\lambda (\stackrel{~}{u}_k)`$ stay bounded away from $`0`$, this is possible only when $`Q_\lambda (\stackrel{~}{u}_k\stackrel{~}{u}_0)0`$. Since $`Q_\lambda ^{1/2}`$ is an equivalent norm, we deduce that $`\stackrel{~}{u}_k\stackrel{~}{u}_0`$ in $`𝒟^{1,2}(^N)`$. In particular $`\stackrel{~}{u}_0_2^{}=lim_k\stackrel{~}{u}_k_2^{}=1`$. Hence, by weakly lower semi-continuity of $`Q_\lambda `$, (97), and (92)
$$S_{\mathrm{circ}}(\lambda )\frac{Q_\lambda (\stackrel{~}{u}_0)}{\stackrel{~}{u}_0_2^{}^2}\underset{k}{lim\; inf}Q_\lambda (\stackrel{~}{u}_k)=\underset{k}{lim\; inf}S_k(\lambda )\underset{k}{lim\; sup}S_k(\lambda )S_{\mathrm{circ}}(\lambda ).$$
Therefore all the above inequalities are indeed equalities. We have thus proved that along a subsequence, $`S_k(\lambda )`$ converges to $`S_{\mathrm{circ}}(\lambda )`$. The Uryson’s property yields convergence of the entire sequence.
## 8. Proof of Theorem 1.4
The proof of Theorem 1.4 is based on Theorem 7.3 and the following lemma.
###### Lemma 8.1.
$`lim\; sup_{k+\mathrm{}}S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$.
Proof. Let $`\epsilon >0`$. Then from (5) and density of $`𝒟(^N\{0\})𝒟_{\mathrm{circ}}^{1,2}(^N)`$ in $`𝒟_{\mathrm{circ}}^{1,2}(^N)`$, there exists $`u𝒟(^N\{0\})𝒟_{\mathrm{circ}}^{1,2}(^N)`$ such that $`_^N|u|^2^{}=1`$ and
(105)
$$_^N|u|^2𝑑x\lambda _0_^N\frac{u^2(x)}{|x|^2}𝑑x\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}_^N\left(\text{ }_{S_r_{\mathrm{}}}\frac{u^2(y)}{|xy|^2}𝑑\sigma (x)\right)𝑑y<S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)+\epsilon .$$
For any $`\mathrm{}=1,\mathrm{},m`$, set
$$f_{\mathrm{}}(x):=_^N\frac{|u(y+x)|^2}{|y|^2}𝑑y,xS_r_{\mathrm{}}.$$
It is easy to check that $`f_{\mathrm{}}C^0(S_r_{\mathrm{}})`$; indeed if $`x_nS_r_{\mathrm{}}`$ converge to $`xS_r_{\mathrm{}}`$, by the Dominated Convergence Theorem we conclude that $`lim_nf_{\mathrm{}}(x_n)=f_{\mathrm{}}(x)`$. Hence the Riemann sum
$$\frac{1}{k}\underset{i=1}{\overset{k}{}}f_{\mathrm{}}(a_i^{\mathrm{}})=\frac{1}{k}\underset{i=1}{\overset{k}{}}_^N\frac{|u(y)|^2}{|ya_i^{\mathrm{}}|^2}𝑑y$$
converges to the integral
$$\frac{1}{2\pi r_{\mathrm{}}}_{S_r_{\mathrm{}}}f_{\mathrm{}}(x)𝑑\sigma (x)=_^N\left(\text{ }_{S_r_{\mathrm{}}}\frac{u^2(y)}{|xy|^2}𝑑\sigma (x)\right)𝑑y.$$
Hence there exists $`\overline{k}`$ such that for all $`k\overline{k}`$
(106) $`{\displaystyle _^N}`$ $`|u|^2dx\lambda _0{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|x|^2}}𝑑x{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\mathrm{\Lambda }_{\mathrm{}}}{k}}{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|xa_i^{\mathrm{}}|^2}}𝑑x\epsilon `$
$`{\displaystyle _^N}|u|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{u^2(x)}{|x|^2}}𝑑x{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}{\displaystyle _^N}\left(\text{ }{\displaystyle _{S_r_{\mathrm{}}}}{\displaystyle \frac{u^2(y)}{|xy|^2}}𝑑\sigma (x)\right)𝑑y.`$
From (105106), we deduce that
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)\epsilon <S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)+\epsilon .$$
Taking $`lim\; sup`$ as $`k+\mathrm{}`$, since $`\epsilon `$ is arbitrary we reach the conclusion.
Proof of Theorem 1.4. As in the proof of Theorem 1.3, we can find a minimizing sequence $`\{u_n\}_n𝒟_k^{1,2}(^N)`$ for (3) with the Palais-Smale property. Under assumption (10), Lemma 6.2 yields
(107)
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_{\mathrm{circ}}(\lambda _0),$$
while (9) implies
(108)
$$S_k(\lambda _0)S_k\left(\lambda _0+k\underset{\mathrm{}=1}{\overset{k}{}}\lambda _{\mathrm{}}\right).$$
Let $`0<\epsilon <S_{\mathrm{circ}}(\lambda _0)S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. From Lemma 8.1, there exists $`k_1=k_1(\epsilon )`$ such that for all $`kk_1`$
(109)
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)+\epsilon .$$
From Theorem 7.3 (107), there exists $`k_2`$ such that for all $`kk_2`$
(110)
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)+\epsilon <S_k(\lambda _0)S_{\mathrm{circ}}(\lambda _0).$$
Let $`k_3`$ be such that for all $`kk_3`$
(111)
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)+\epsilon \mathrm{min}\{k^{\frac{2}{N}}S,k^{\frac{2}{N}}S(\lambda _1),\mathrm{},k^{\frac{2}{N}}S(\lambda _m)\}.$$
From (108111), we conclude that for all $`k\mathrm{max}\{k_1,k_2,k_3\}`$
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<\mathrm{min}\{k^{\frac{2}{N}}S,k^{\frac{2}{N}}S(\lambda _1),\mathrm{},k^{\frac{2}{N}}S(\lambda _m),S_k(\lambda _0),S_k\left(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}}\right)\}.$$
From above and the Palais-Smale condition proved in Theorem 4.1, we deduce that $`\{u_n\}_n`$ has a subsequence strongly converging to some $`u_0𝒟_k^{1,2}(^N)`$ such that $`J_k(u_0)=\frac{1}{N}S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. Hence $`u_0`$ achieves the infimum in (3). Since $`J_k`$ is even, also $`|u_0|`$ is a minimizer in (3) and then $`v_0=S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{1/(2^{}2)}|u_0|`$ is a nonnegative solution to equation (2). The maximum principle implies the positivity outside singularities of such a solution.
## 9. Proof of Theorem 1.5
We now provide a sufficient condition for the infimum in (3) to stay below the level $`k^{2/N}S(\lambda _j)`$, in correspondence of which possible concentration at singular points located at the $`j`$-th polygon can occur. We denote by $`\mathrm{\Theta }_j\mathrm{}`$ the minimum angle formed by vectors $`a_i^j`$ and $`a_s^{\mathrm{}}`$, see figure below.
Figure 4 (The angle $`\mathrm{\Theta }_j\mathrm{}`$.)
The following lemma can be proved by standard trigonometry calculus.
###### Lemma 9.1.
For any $`i,s=1,2,\mathrm{},k`$, and $`j,\mathrm{}=1,2,\mathrm{},m`$, there holds
$`|a_i^ja_s^j|=2r_j\left|\mathrm{sin}{\displaystyle \frac{(si)\pi }{k}}\right|,`$
$`|a_i^ja_s^{\mathrm{}}|^2=r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}\left({\displaystyle \frac{2\pi (is)}{k}}+\mathrm{\Theta }_j\mathrm{}\right).`$
###### Lemma 9.2.
Let $`j\{1,2,\mathrm{},m\}`$. If
(112) $`0<\lambda _j{\displaystyle \frac{N(N4)}{4}}`$
and
(113) $`{\displaystyle \frac{\lambda _0}{|r_j|^2}}+\lambda _j{\displaystyle \underset{i=1}{\overset{k1}{}}}{\displaystyle \frac{1}{4r_j^2\left|\mathrm{sin}\frac{i\pi }{k}\right|^2}}+{\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{m}{}}}\lambda _{\mathrm{}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}\left(\frac{2\pi i}{k}+\mathrm{\Theta }_j\mathrm{}\right)}}>0,`$
then
(114) $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<k^{2/N}S(\lambda _j).`$
Proof. Let $`z(x)=_{i=1}^kz_\mu ^{\lambda _j}(xa_i^j)𝒟_k^{1,2}(^N)`$. Then
$`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ $`{\displaystyle \frac{Q_{\lambda _0,\lambda _1,\mathrm{},\lambda _m}(z)}{\left({\displaystyle _^N}|z|^2^{}𝑑x\right)^{2/2^{}}}},`$
where $`Q_{\lambda _0,\lambda _1,\mathrm{},\lambda _m}`$ denotes the quadratic form defined by
$$Q_{\lambda _0,\lambda _1,\mathrm{},\lambda _m}(u)=_^N|u|^2𝑑x_^N\frac{\lambda _0}{|x|^2}u^2(x)𝑑x\underset{\mathrm{}=1}{\overset{m}{}}\underset{i=1}{\overset{k}{}}\lambda _{\mathrm{}}_^N\frac{u^2(x)}{|xa_i^{\mathrm{}}|^2}𝑑x.$$
Note that
$$\left(_^N|z|^2^{}𝑑x\right)^{2/2^{}}k^{2/2^{}}.$$
Moreover, from (19), Lemmas 5.1 and 5.2 we find that
$`Q_{\lambda _0,\lambda _1,\mathrm{},\lambda _m}(z)=k{\displaystyle _^N}|z_\mu ^{\lambda _j}|^2𝑑xk\lambda _j{\displaystyle _^N}{\displaystyle \frac{|z_1^{\lambda _j}|^2}{|x|^2}}𝑑x\lambda _0{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle _^N}{\displaystyle \frac{|z_\mu ^{\lambda _j}|^2}{|x+a_i^j|^2}}𝑑x`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ (i,\mathrm{})(s,j)\end{array}}{\overset{k}{}}}{\displaystyle _^N}{\displaystyle \frac{\lambda _{\mathrm{}}|z_\mu ^{\lambda _j}|^2}{|x+a_s^ja_i^{\mathrm{}}|^2}}𝑑x+{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ is\end{array}}{\overset{k}{}}}{\displaystyle _^N}z_\mu ^{\lambda _j}(xa_i^j)z_\mu ^{\lambda _j}(xa_s^j)𝑑x`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}i=1,s=1,t=1\\ st\end{array}}{\overset{k}{}}}{\displaystyle _^N}{\displaystyle \frac{\lambda _{\mathrm{}}z_\mu ^{\lambda _j}(xa_s^j)z_\mu ^{\lambda _j}(xa_t^j)}{|xa_i^{\mathrm{}}|^2}}𝑑x\lambda _0{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ is\end{array}}{\overset{k}{}}}{\displaystyle _^N}{\displaystyle \frac{z_\mu ^{\lambda _j}(xa_i^j)z_\mu ^{\lambda _j}(xa_s^j)}{|x|^2}}𝑑x`$
$`=\{\begin{array}{cc}kS(\lambda _j)\mu ^2\left(_^N|z_1^{\lambda _j}|^2\right)\left[\frac{k\lambda _0}{|r_j|^2}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ (i,\mathrm{})(s,j)\end{array}}{\overset{k}{}}}\frac{\lambda _{\mathrm{}}}{|a_i^{\mathrm{}}a_s^j|^2}+o(1)\right]\text{if }\lambda _j<\frac{N(N4)}{4},\hfill & \\ kS(\lambda _j)\mu ^2|\mathrm{ln}\mu |\alpha _{\lambda _j,N}^2\left[\frac{k\lambda _0}{|r_j|^2}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ (i,\mathrm{})(s,j)\end{array}}{\overset{k}{}}}\frac{\lambda _{\mathrm{}}}{|a_i^{\mathrm{}}a_s^j|^2}+o(1)\right]\text{if }\lambda _j=\frac{N(N4)}{4}.\hfill & \end{array}`$
Therefore (114) holds provided
(115) $`{\displaystyle \frac{k\lambda _0}{|r_j|^2}}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}i=1,s=1\\ (i,\mathrm{})(s,j)\end{array}}{\overset{k}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|a_i^{\mathrm{}}a_s^j|^2}}>0.`$
It is easy to verify that assumption (113) and Lemma 9.1 imply (115).
###### Lemma 9.3.
For any $`j,\mathrm{}=1,\mathrm{},m`$, $`j\mathrm{}`$, there holds
(116) $`\underset{k\mathrm{}}{lim}{\displaystyle \frac{1}{\mathrm{\Lambda }_j}}\left[\lambda _j{\displaystyle \underset{i=1}{\overset{k1}{}}}{\displaystyle \frac{1}{4r_j^2\left|\mathrm{sin}\frac{i\pi }{k}\right|^2}}\right]=+\mathrm{},`$
(117) $`\lambda _{\mathrm{}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}\left(\frac{2\pi i}{k}+\mathrm{\Theta }_j\mathrm{}\right)}}=\mathrm{\Lambda }_{\mathrm{}}O(1)\text{as }k+\mathrm{}.`$
Proof. A direct calculation yields
$`{\displaystyle \frac{\lambda _j}{\mathrm{\Lambda }_j}}{\displaystyle \underset{i=1}{\overset{k1}{}}}{\displaystyle \frac{1}{4r_j^2\left|\mathrm{sin}\frac{i\pi }{k}\right|^2}}`$ $`{\displaystyle \frac{1}{k}}{\displaystyle _1^{k/2}}{\displaystyle \frac{ds}{4r_j^2\left|\mathrm{sin}\frac{s\pi }{k}\right|^2}}{\displaystyle \frac{1}{\pi }}{\displaystyle _{\pi /k}^{\pi /2}}{\displaystyle \frac{dt}{4r_j^2|\mathrm{sin}t|^2}}\underset{k+\mathrm{}}{}+\mathrm{}.`$
On the other hand
$`{\displaystyle \frac{\lambda _{\mathrm{}}}{\mathrm{\Lambda }_{\mathrm{}}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}\left(\frac{2\pi i}{k}+\mathrm{\Theta }_j\mathrm{}\right)}}`$ $`{\displaystyle \frac{1}{k}}{\displaystyle _0^k}{\displaystyle \frac{ds}{r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}\left(\frac{2\pi s}{k}+\mathrm{\Theta }_j\mathrm{}\right)}}`$
$`={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{dt}{r_j^2+r_{\mathrm{}}^22r_jr_{\mathrm{}}\mathrm{cos}(t+\mathrm{\Theta }_j\mathrm{})}}^+,`$
thus proving (117).
###### Remark 9.4.
Lemma 9.3 implies that if we fix $`\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m`$ and let $`k+\mathrm{}`$, then the quantity in formula (113) tends to $`+\mathrm{}`$. Hence condition (113) is satisfied for $`k`$ sufficiently large.
Let us now compare levels $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ and $`S_k(\lambda _0)`$ which is related to concentration at the origin. Two cases can occur:
$`(i)`$ $`S_k(\lambda _0)k^{2/N}S,`$
$`(\mathrm{𝑖𝑖})`$ $`S_k(\lambda _0)<k^{2/N}S.`$
In case (i), since $`S=S(0)`$ and $`\lambda S(\lambda )`$ is a nonincreasing function, to exclude that the infimum in (22) stays above $`S_k(\lambda _0)`$ it is enough to compare $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ with $`k^{2/N}S(\lambda _j)`$ where $`\lambda _j=\mathrm{max}\{\lambda _{\mathrm{}}\}_{1\mathrm{}m}`$, as we have done in Lemma 9.2.
The study of case (ii) is based on Lemma 2.1. Indeed, using Lemma 2.1 and estimates of Lemma 5.1, we can prove the following lemma.
###### Lemma 9.5.
If $`N4`$, $`S_k(\lambda _0)<k^{2/N}S`$, and one of the following assumptions is satisfied
(118) $`\lambda _0{\displaystyle \frac{N(N4)}{4}}`$ $`\text{and}{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^2}}>0,`$
(119) $`{\displaystyle \frac{N(N4)}{4}}<\lambda _0<{\displaystyle \frac{(N2)^2}{4}}`$ $`\text{and}{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}}}>0,`$
then
(120)
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_k(\lambda _0).$$
Proof. From Lemma 2.1, we have that $`S_k(\lambda _0)`$ is attained by some $`u^{\lambda _0}𝒟_k^{1,2}(^N)`$. By homogeneity of the Rayleigh quotient, we can assume $`|u^{\lambda _0}|^2^{}=1`$. Furthermore, the function $`v^{\lambda _0}=S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{1/(2^{}2)}|u^{\lambda _0}|`$ is a nonnegative solution to (16), hence we can apply Lemma 5.1 to study the behavior of $`_^N\frac{|u_\mu ^{\lambda _0}|^2}{|x+\xi |^2}𝑑x`$ as $`\mu 0`$. Hence we obtain that there exists some positive constant $`\kappa _0`$ such that
(121) $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$
$`{\displaystyle _^N}|u_\mu ^{\lambda _0}(x)|^2𝑑x\lambda _0{\displaystyle _^N}{\displaystyle \frac{|u_\mu ^{\lambda _0}(x)|^2}{|x|^2}}𝑑x{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}\lambda _{\mathrm{}}{\displaystyle _^N}{\displaystyle \frac{|u_\mu ^{\lambda _0}(x)|^2}{|xa_i^{\mathrm{}}|^2}}𝑑x`$
$`=S_k(\lambda _0)\{\begin{array}{cc}\mu ^2_^N|u^{\lambda _0}|^2\left({\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^2}}+o(1)\right)\hfill & \text{if }\lambda _0<\frac{N(N4)}{4}\hfill \\ \kappa _0^2\mu ^2|\mathrm{ln}\mu |\left({\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^2}}+o(1)\right)\hfill & \text{if }\lambda _0=\frac{N(N4)}{4}\hfill \\ \kappa _0^2\beta _{\lambda ,N}\mu ^{\sqrt{(N2)^24\lambda }}\left({\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{\lambda _{\mathrm{}}}{|r_{\mathrm{}}|^{\sqrt{(N2)^24\lambda _0}}}}+o(1)\right)\hfill & \text{if }\lambda _0>\frac{N(N4)}{4}.\hfill \end{array}`$
Taking $`\mu `$ sufficiently small we obtain that either assumption (118) or (119) yield (120).
Proof of Theorem 1.5. As in the proof of Theorem 1.3, we can find a minimizing sequence $`\{u_n\}_n`$ which has the Palais-Smale property, more precisely $`J_k^{}(u_n)0`$ in $`(𝒟_k^{1,2}(^N))^{}`$ and $`J_k(u_n)\frac{1}{N}S_k(\lambda _0\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. Assumption (11) yields
(122)
$$S_k\left(\lambda _0+k\underset{\mathrm{}=1}{\overset{k}{}}\lambda _{\mathrm{}}\right)S_k(\lambda _0).$$
Note also that (12) and (14) imply that $`N>4`$ and $`\lambda _m>0`$. Two cases can occur. If $`S_k(\lambda _0)<k^{\frac{2}{N}}S`$, then Lemma 9.5 and assumption (14) yield
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<S_k(\lambda _0).$$
If $`S_k(\lambda _0)k^{2/N}S`$, from monotonicity we have $`S_k(\lambda _0)k^{2/N}S>k^{2/N}S(\lambda _m)`$. In both cases from Lemma 9.2 and (15) we deduce $`k^{2/N}S(\lambda _m)>S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$. Hence we obtain that
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)<\mathrm{min}\{k^{\frac{2}{N}}S,k^{\frac{2}{N}}S(\lambda _1),\mathrm{},k^{\frac{2}{N}}S(\lambda _m),S_k(\lambda _0),S_k\left(\lambda _0+k_{\mathrm{}=1}^m\lambda _{\mathrm{}}\right)\}.$$
From above and the Palais-Smale condition proved in Theorem 4.1, we deduce that $`\{u_n\}_n`$ has a subsequence strongly converging to some $`u_0𝒟_k^{1,2}(^N)`$ which achieves the infimum in (3). Moreover $`v_0=S_k(\lambda _0\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)^{\frac{1}{2^{}2}}|u_0|`$ is a solution to (2).
###### Remark 9.6.
Theorem 1.5 contains an alternative proof to Theorem 1.4 in the case $`N>4`$, as it follows easily gathering Theorem 1.5 and Remark 9.4. Note that the assumption $`N>4`$ is needed to ensure that (12) and (14) hold. However, with respect to Theorem 1.4, it contains a more precise information on how $`k`$ must be large in order to solve the problem.
## 10. Proof of Theorem 1.2
Proof of Theorem 1.2. Assume first (i). Let $`\epsilon >0`$. Then from (3) and density of $`𝒟(^N\{0\})𝒟_k^{1,2}(^N)`$ in $`𝒟_k^{1,2}(^N)`$, there exists $`u𝒟(^N\{0\})𝒟_k^{1,2}(^N)`$ such that $`Q_{\lambda _0}(u)S_k(\lambda _0)+\epsilon `$. Let
$$u_\mu (x)=\mu ^{\frac{N2}{2}}u(x/\mu ),\mu >0.$$
By Dominated Convergence Theorem it is easy to verify that
$$\underset{\mu 0}{lim}_^N\frac{|u_\mu (x)|^2}{|xa_i^{\mathrm{}}|^2}=0,$$
hence
$`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ $`{\displaystyle \frac{_^N|u_\mu |^2𝑑x\lambda _0_^N\frac{u_\mu ^2(x)}{|x|^2}_^N_{\mathrm{}=1}^m\frac{\mathrm{\Lambda }_{\mathrm{}}}{k}_{i=1}^k\frac{u_\mu ^2(x)}{|xa_i^{\mathrm{}}|^2}dx}{\left(_^N|u_\mu |^2^{}𝑑x\right)^{2/2^{}}}}`$
$`=Q_{\lambda _0}(u)+o(1)S_k(\lambda _0)+\epsilon +o(1)\text{as }\mu 0.`$
Letting $`\mu 0`$ we have that $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)S_k(\lambda _0)+\epsilon `$ for all $`\epsilon >0`$. Hence
$$S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)S_k(\lambda _0).$$
Assume by contradiction that $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ is attained by $`\overline{u}𝒟_k^{1,2}(^N)\{0\}`$, then
$`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ $`={\displaystyle \frac{_^N|\overline{u}|^2𝑑x\lambda _0_^N\frac{\overline{u}^2(x)}{|x|^2}_^N_{\mathrm{}=1}^m\frac{\mathrm{\Lambda }_{\mathrm{}}}{k}_{i=1}^k\frac{\overline{u}^2(x)}{|xa_i^{\mathrm{}}|^2}dx}{\left(_^N|\overline{u}|^2^{}𝑑x\right)^{2/2^{}}}}`$
$`>{\displaystyle \frac{_^N|\overline{u}|^2𝑑x_^N\frac{\lambda _0}{|y|^2}\overline{u}^2(y)𝑑y}{\left(_^N|\overline{u}|^2^{}𝑑x\right)^{2/2^{}}}}S_k(\lambda _0),`$
giving rise to a contradiction. The proof of non-attainability of $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ is analogous and is based on (90).
Assume now that (ii) holds. Then for all $`u𝒟_{\mathrm{circ}}^{1,2}(^N)`$, $`u0`$, denoting by $`u^{}`$ the Schwarz symmetrization of $`u`$ (see (28)), from (29) and (30) it follows that
$`{\displaystyle \frac{Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}(u)}{\left(_^N|u|^2^{}𝑑x\right)^{2/2^{}}}}`$ $`{\displaystyle \frac{_^N|u^{}|^2𝑑x\left(\lambda _0+_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}}\right)_^N\frac{|u^{}(y)|^2}{|y|^2}𝑑y}{\left(_^N|u^{}|^2^{}𝑑x\right)^{2/2^{}}}}`$
$`S_{\mathrm{circ}}\left(\lambda _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}\right)=S\left(\lambda _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}\right).`$
Hence $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)S_{\mathrm{circ}}\left(\lambda _0+_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}}\right)`$. On the other hand, setting $`\mathrm{\Lambda }=\lambda _0+_{\mathrm{}=1}^m\mathrm{\Lambda }_{\mathrm{}}`$, and using \[14, Corollary 3.2\] we obtain
$`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ $`Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}(z_\mu ^\mathrm{\Lambda })=Q_\mathrm{\Lambda }(z_\mu ^\mathrm{\Lambda })+o(1)=S(\mathrm{\Lambda })+o(1)`$
$`=S_{\mathrm{circ}}\left(\lambda _0+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}\right)+o(1)\text{as }\mu \mathrm{}.`$
Then
(123)
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)=S_{\mathrm{circ}}\left(\lambda _0+\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{\Lambda }_{\mathrm{}}\right).$$
If $`S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ was attained by some $`\overline{u}𝒟_{\mathrm{circ}}^{1,2}(^N)\{0\}`$, then
$$S_{\mathrm{circ}}(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)=\frac{Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}(\overline{u})}{\left(_^N|\overline{u}|^2^{}𝑑x\right)^{2/2^{}}}\frac{Q_\mathrm{\Lambda }(\overline{u}^{})}{\left(_^N|\overline{u}^{}|^2^{}𝑑x\right)^{2/2^{}}}S_{\mathrm{circ}}(\mathrm{\Lambda }).$$
Due to (123), all above inequalities are indeed equalities; in particular $`Q_{\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m}^{\mathrm{circ}}(\overline{u})=Q_\mathrm{\Lambda }(\overline{u}^{})`$ which (taking into account Polya-Szego inequality and (29)) yields
$`0`$ $`{\displaystyle _^N}|\overline{u}|^2𝑑x{\displaystyle _^N}|\overline{u}^{}|^2𝑑x`$
$`={\displaystyle _^N}\left({\displaystyle \frac{\lambda _0}{|y|^2}}+{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}\mathrm{\Lambda }_{\mathrm{}}\text{ }{\displaystyle _{S_r_{\mathrm{}}}}{\displaystyle \frac{d\sigma (x)}{|xy|^2}}\right)\overline{u}^2(y)𝑑y\mathrm{\Lambda }{\displaystyle _^N}{\displaystyle \frac{|\overline{u}^{}(y)|^2}{|y|^2}}𝑑y0.`$
Then $`_^N|\overline{u}|^2𝑑x=_^N|\overline{u}^{}|^2𝑑x`$. From , it follows that $`\overline{u}`$ must be spherically symmetric with respect to some point. Since $`\overline{u}`$ is a solution to equation (4) (up to some Lagrange’s multiplier), the potential in equation (4) must be spherically symmetric, thus giving rise to a contradiction. The proof of non-attainability of $`S_k(\lambda _0,\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m)`$ is contained in \[14, Theorem 1.3\].
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# Relaxation time scales in collective dynamics of liquid alkali metals
## I Introduction
The dynamic structure factor $`S(k,\omega )`$ is an experimentally measured term, containing information about the processes in a liquid with long- and short-time scales. It can be used to judge on the microscopic behavior in a system on the basis of its spectra, obtained by means of Inelastic Neutron Scattering (INS) Balucani\_Zoppi ; Copley or Inelastic X-ray Scattering (IXS) Burkel . As for simple liquids, at present a great amount of experimental data of $`S(k,\omega )`$ has been accumulated, in particular, for liquid alkali metals. These data indicate legibly the presence of the collective propagation excitations beyond the hydrodynamic region. The characteristic feature of liquid alkali metals is a triple-peak shape of $`S(k,\omega )`$ lasted to $`k0.8k_m`$, where $`k_m`$ corresponds to the first maximum of the static structure factor $`S(k)`$. Moreover, the frequency of the side peak achieves its maximum at $`k0.55k_m`$. The propagation of these high-frequency waves cannot be obtained within a hydrodynamic treatment, therefore, they are related in some works to the so-called kinetic collective excitations. The impossibility to describe these microscopic phenomena and, therefore, to reproduce qualitatively the experimental $`S(k,\omega )`$ by means of ordinary hydrodynamic equations led to the development of other theoretical models and approaches.
One of the simplest and perhaps the earliest modeling approaches is the so-called viscoelastic theory. It allows one to obtain the central quasi-elastic line as well as two inelastic peaks symmetrically located around $`\omega =0`$ for mesoscopic space-frequency region. However, as shown in Refs. Bodensteiner ; Tullio0 , this model can not be used for the exact reproduction of the experimental spectral shapes of $`S(k,\omega )`$ (see, for instance, the cases of liquid cesium and lithium in Refs. Bodensteiner ; Tullio0 ). Therefore, in Ref. Tullio0 the double-scale model for the viscous relaxation process with fast and slow time scales was tested, and as a result a good agreement with the IXS experimental data for the dynamic structure factor was received. Recently the similar approach was also applied for the description of relaxation processes in H-bonded liquids Angelini ; Balucani1 . The existence of two time scales in this model reflects the presence of physically different decay mechanisms. A faster process is hypothetically associated with interactions between an atom and the “cage” of its nearest neighbors, and a slower one is identified with the well-known structural $`(\alpha )`$ process. However, relaxations of both processes are approximated by exponential dependencies. In recent works the viscoelastic model has been improved by means of the Markovian closure on the next relaxation level of Zwanzig-Mori hierarchy Cabrillo , it is equivalent to the exponential relaxation on this level. It is worth mentioning two others methods, one of which is related to the extension of the usual hydrodynamic analytical expressions by modification of hydrodynamic modes to $`k`$-dependence (see, for instance, Ref. McGreevy ). This method assumes the existence of non-hydrodynamical additional modes. The second approach is related to the so-called concept of generalized collective modes, which was proposed for the investigation of the time correlation functions (TCF’s) beyond the hydrodynamic region Bruin . The key idea of this method consists in the correct choice of the basic set of dynamical variables.
All these methods are more or less successfully used for the description of collective dynamics in liquids. They have common property. Namely, they are actually constructed on heuristic assumption about the presence of exponential decay (or combination of exponential decay contributions) in some relaxation processes. Nevertheless, the transition and imposition of different relaxation modes in disorder systems can occur even in case of a concrete relaxation process, that complicates the selection of the analytical time dependence for the corresponding TCF. This fact is proved by the successful application of different mode-coupling theories. On the other hand, this difficulty can be resolved by means of analysis and comparison of the resulting time scales of relaxation processes. Therefore, in the present work we suggest the approach, which allows us to avoid the immediate approximation of relaxation processes by analytical functions. It is based on the development of Bogoliubov’s ideas about the hierarchy of relaxation times in liquids Bogoliubov , adapted to the formalism of time correlation functions.
One of the open problems in studying of liquid state (in particular, of the microdynamics of simple liquids) is to describe and understand on a general ground the common features of different relaxation processes Angelini . It is well known, that the dispersion of the side (high-frequency) peak of dynamic structure factor is the same for all alkali metals. Moreover, it is also valid in case of more complex systems, for example, for liquid alloys Bove . Then the following questions arise: Is the origin of relaxation processes the same for liquid systems with the similar features? Can the unified description be applied to these systems? As for the group of melting alkalis, it has been indicated in Ref. Balucani\_PRB that both the equilibrium and the time-dependent correlations can be cast in a properly scaled form for all the the alkali metals. Further, it was justified by ab initio molecular dynamic studying in Ref. Balucani\_PRB too. Experimental confirmation of this result was impossible over a long period particularly because of the difficulties related to the technique of INS due to the deficient precision of the experimental data. Recently, due to progress in IXS technique this issue was considered again Tullio1 . In this work we present investigations related to the determination of corresponding scale transitions for liquid systems.
The organization of the paper is as follows. In the next Section, we describe the theoretical formalism, and the comparison with the experimental data and other theories is carried out. The possibility of scale uniformity of dynamical processes in the group of liquid alkali metals is analyzed and discussed in Section III. The scale-crossing relations are also presented here. Finally, we come up with some concluding remarks in Section IV.
## II Theoretical formalism
### II.1 Basic notions
Let us consider the liquid system of $`N`$ identical classical particles of the mass $`m`$ in the volume $`V`$ and take the density fluctuations
$$W_0(\text{k})=\frac{1}{\sqrt{N}}\underset{j=1}{\overset{N}{}}e^{i\text{k}\text{r}_j}$$
(1)
as an initial dynamical variable. To construct a some set of dynamical variables necessary for the description of the evolution of the system we use the technique of projection operators of Zwanzig-Mori Zwanzig ; Mori . It is a formal version of the Gram-Schmidt orthogonalization process, which allows one to obtain the set of orthogonal variables
$`\text{W}(k)=\{W_0(k),W_1(k),W_2(k),\mathrm{},W_j(k),\mathrm{}\}.`$ (2)
They satisfy the condition $`W_j^{}W_l=\delta _{j,l}|W_j|^2`$ and are connected by the following recurrent relation Yulm\_Khus :
$`W_{j+1}(k)=W_j(k)\mathrm{\Omega }_j^2(k)W_{j1}(k),`$
$`j=0,1,2,\mathrm{};W_1(k)=0.`$ (3)
Here the characteristic of the corresponding $`j`$th relaxation process, the so-called frequency parameter $`\mathrm{\Omega }_j^2(k)`$, appears, $``$ is the Liouville operator
$$=i\left\{\underset{j=1}{\overset{N}{}}\frac{\text{p}_j_j}{m}\underset{i>j=1}{\overset{N}{}}_ju(j,i)(_p^j_p^i)\right\}$$
(4)
with the momentum of the $`j`$th particle $`\text{p}_j`$ and the pair potential $`u(j,i)`$. So, if $`W_0(k)`$ is the density fluctuations, then $`W_1(k)`$ is the longitudinal component of the momentum density and so it goes on.
The TCF’s for the corresponding dynamical variables are given by
$$M_{jl}(k,t)=\frac{W_j^{}(k)e^{i_{22}^{(l)}t}W_l(k)}{W_j^{}(k)W_l(k)},j,l=1,2,\mathrm{}$$
(5)
For convenience normalized time correlation functions are used here. The time-evolution operator of Eq. (5) contains the reduced Liouville operator
$$_{22}^{(l)}=\left(1\underset{j=1}{\overset{l}{}}\mathrm{\Pi }_j\right)\left(1\underset{j=1}{\overset{l}{}}\mathrm{\Pi }_j\right),$$
(6)
defined by the following projection operators
$$\mathrm{\Pi }_j=\frac{W_j(k)W_j^{}(k)}{|W_j(k)|^2},\mathrm{\Pi }_j\mathrm{\Pi }_l=\delta _{j,l}\mathrm{\Pi }_j.$$
(7)
From the condition of orthogonalization of the dynamical variables we obtain the initial values for the TCF’s of Eq. (5):
$$M_{jl}(k,t=0)=\{\begin{array}{ccc}\hfill 0,& \text{if}jl,& \text{cross-correlations}\hfill \\ \hfill 1,& \text{if}j=l,& \text{autocorrelations}\hfill \end{array}$$
(8)
These correlation functions $`M_{jl}(k,t)`$ are symmetrical in $`l`$ and $`j`$, i.e.,
$$M_{jl}(k,t)=M_{lj}(k,t).$$
(9)
Autocorrelation functions of Zwanzig-Mori formalism have the following property: every autocorrelation function of the higher order $`M_j(k,t)=M_{jj}(k,t)`$ is a memory function for the previous one, i.e. $`M_{j1}(k,t)`$ (autocorrelation functions will be marked by one index only in accordance with the used variable), and they are interrelated by integro-differential non-Markovian equations of the form:
$$\frac{dM_{j1}(k,t)}{dt}+\mathrm{\Omega }_j^2(k)_0^t𝑑\tau M_j(k,\tau )M_{j1}(k,t\tau )=0.$$
(10)
Differentiating the first equation of the chain (10), i.e. $`j=1`$, one obtain the generalized Langevin equation:
$$\frac{d^2M_0(k,t)}{dt^2}+\mathrm{\Omega }_1^2(k)M_0(k,t)\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_2^2(k)_0^t𝑑\tau _0^\tau 𝑑\tau ^{^{}}M_2(k,t\tau )M_1(k,t\tau ^{^{}})M_0(k,\tau ^{^{}})=0.$$
(11)
One the other hand, these functions describe concrete relaxation processes, the physical meaning of which may be established from direct definitions of TCF’s. For instance, $`M_0(k,t)`$ describes the dynamics of fluctuations of density correlations in the system, $`M_1(k,t)`$ is the TCF of the fluctuations of the longitudinal component of the momentum density, $`M_2(k,t)`$ contains the TCF of fluctuations of energy density. So, these quantities are associated with the TCF’s of the well-known hydrodynamic “slow” variables. These TCF’s have characteristic time scales, which can be found from
$$\tau _j(k)=\text{Re}_0^{\mathrm{}}𝑑tM_j(k,t)=\text{Re}\stackrel{~}{M}_j(k,s=0),$$
(12)
where $`\stackrel{~}{M}_j(k,s)`$ is the Laplace transform of the corresponding TCF, i.e. $`\stackrel{~}{M}_j(k,s)=_0^{\mathrm{}}𝑑te^{st}M_j(k,t)`$ time ; Egelstaff ; Costa .
So, the memory function approach with single initial dynamical variable extracts the whole set, which describes the relaxation processes of the corresponding relaxation levels. In fact, the well-known problem of the choice of a set of variables required for the correct description of the system dynamics here is reduced (i) to the search of the number of variables for a priori known succession $`\text{W}(k)`$, that was excellently shown by the recurrent relation approach in a works of Lee Lee ; Omega ; and/or (ii) to the finding the correct closure of the chain (10).
The ratio between $`\tau _0(k)`$, $`\tau _1(k)`$ and $`\tau _2(k)`$ may be quite arbitrary. In the hydrodynamic region ($`k0`$, $`\omega 0`$) they take large values due to the slow changes of the correspondent variables: densities of mass, momentum and energy. Further, one can suggest that the relaxation times of the following TCF’s, in comparison with the scales of these three variables, are comparable, i.e. $`\tau _3(k)\tau _4(k)`$. We emphasize here that this assumption does not contradict the viscoelastic model, which presupposes that $`\tau _2(k)\tau _3(k)`$. Obviously, this key condition of the viscoelastic theory is just a special case in our approach. Simultaneously, our approach does not deny the presence of the long-lasting time tail of $`M_2(k,t)`$, which may be adequately taken into account by the mode-coupling theory Gotze .
Then, taking into account Eq. (12) one can find
$$M_4(k,t)=M_3(k,t)+h(k,t),$$
(13)
where the “tail” function $`h(k,t)`$ appears. From the short-time asymptotic of the time autocorrelation functions and the condition of the long time attenuation of correlation Eq. (13) yields the following properties of $`h(k,t)`$:
$$\underset{t0}{lim}h(k,t)=\underset{t\mathrm{}}{lim}h(k,t)=0,$$
(14)
this function must have at least one crossing with the time axis at the intermediate region tail . Eq. (13) allows us to obtain the closure of hierarchy of equations of the form (10) at the fourth level ($`j=4`$) and by means of Laplace transformation to find its exact solution for $`\stackrel{~}{M}_0(k,i\omega )`$, in particular, which is directly related to the experimentally available term, the dynamic structure factor, $`S(k,\omega )`$. The expression for the resulting $`S(k,\omega )`$ is given in work our\_JCP in terms of the first four frequency parameters $`\mathrm{\Omega }_1^2(k)`$, $`\mathrm{\Omega }_2^2(k)`$, $`\mathrm{\Omega }_3^2(k)`$, $`\mathrm{\Omega }_4^2(k)`$ and the Laplace transform of tail function, i.e., $`\stackrel{~}{h}(k,i\omega )`$. In some cases, the regime with $`h(k,t)0`$ may be realized. It can be observed in some parts of time (frequency) scale. In this case we find the following expression for the dynamic structure factor:
$`S(k,\omega )`$ $`=`$ $`{\displaystyle \frac{S(k)}{2\pi }}\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_2^2(k)\mathrm{\Omega }_3^2(k)[4\mathrm{\Omega }_4^2(k)\omega ^2]^{\frac{1}{2}}\{\mathrm{\Omega }_1^4(k)\mathrm{\Omega }_3^4(k)`$ (15)
$`+`$ $`\omega ^2[\mathrm{\Omega }_1^4(k)\mathrm{\Omega }_4^2(k)2\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_3^4(k)\mathrm{\Omega }_1^4(k)\mathrm{\Omega }_3^2(k)`$
$`+2\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_2^2(k)\mathrm{\Omega }_4^2(k)\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_2^2(k)\mathrm{\Omega }_3^2(k)+\mathrm{\Omega }_2^4(k)\mathrm{\Omega }_4^2(k)]`$
$`+`$ $`\omega ^4[\mathrm{\Omega }_3^4(k)2\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_4^2(k)+2\mathrm{\Omega }_1^2(k)\mathrm{\Omega }_3^2(k)2\mathrm{\Omega }_2^2(k)\mathrm{\Omega }_4^2(k)+\mathrm{\Omega }_2^2(k)\mathrm{\Omega }_3^2(k)]`$
$`+`$ $`\omega ^6[\mathrm{\Omega }_4^2(k)\mathrm{\Omega }_3^2(k)]\}^1.`$
This equation is also expressed through the first four frequency parameters, which are directly related to the first five even frequency moments of dynamics structure factor. It is necessary to note that this expression is obtained in the way completely different from the theory of moments Orthner .
### II.2 Comparison with IXS experiment and relationship with other theoretical approaches
In Fig. $`1`$ we report the dynamic structure factor $`S(k,\omega )`$ of liquid lithium ($`T=475`$K) for some wave numbers calculated from Eq. (15) (solid line) and obtained from IXS experiment (circles) Tullio1 . Being used in theoretical computations the static structure factor $`S(k)`$ for both cases was taken from Ref. Waseda . The first frequency parameter was directly defined from its definition $`\mathrm{\Omega }_1^2(k)=K_BTk^2/mS(k)`$. The second frequency parameter $`\mathrm{\Omega }_2^2(k)`$ is related to the fourth frequency moment. We found this parameter from the values of the infinite frequency sound velocity $`c_{\mathrm{}}(k)`$ Tullio0 ; Tullio1 by means of relation $`c_{\mathrm{}}(k)=\sqrt{\mathrm{\Omega }_1^2(k)+\mathrm{\Omega }_2^2(k)}/k`$. The high-order parameters were found by comparison with the experiment. Eventually, we have revealed that all the frequency parameters have the similar dispersion. In particular, they have the first principal maximum at the same wave numbers such as the side peak of $`S(k,\omega )`$, i.e. at $`k0.55k_m`$, and any low order parameter is less than the high order one.
We would like to emphasize that the theoretical $`S(k,\omega )`$ and, in particular, the position of the side peak, is very sensitive to the magnitude of $`\mathrm{\Omega }_2^2(k)`$. The magnitudes of $`\mathrm{\Omega }_3^2(k)`$ and $`\mathrm{\Omega }_4^2(k)`$ influence the form of $`S(k,\omega )`$. However, it is not so important to know these parameters separately as their ratio, i.e. $`\mathrm{\Omega }_4^2(k)/\mathrm{\Omega }_3^2(k)`$.
To compare the theoretical outcome with the experiment we modified it to account for the quantum mechanical detailed balance condition according to
$$S_q(k,\omega )\frac{\mathrm{}\omega /K_BT}{1e^{\mathrm{}\omega /K_BT}}S(k,\omega ),$$
(16)
and then broadened it for the finite experimental resolution effects $`R(k,\omega )`$ Tullio0 :
$$R(k,\omega \omega ^{})S_q(k,\omega ^{})𝑑\omega .$$
(17)
From Fig. $`1`$ one can see that the above described theoretical approach yields a good agreement with IXS data of both systems.
Now we can execute a more detailed study of the obtained results and compare them with other approaches: the usual viscoelastic model, the double-scale viscous model and the generalized mode approach. The common feature of these theories is the use of the time autocorrelation function $`M_{jl}(k,t)`$ of Eq. (5) at $`j=l=2`$. So, the viscoelastic and the double-viscosity models are based on approximations to this term, and $`M_2(k,t)`$ plays a key role in these theories. As for our approach, it gives the following form for Laplace transform of $`M_2(k,t)`$
$`\stackrel{~}{M}_2(k,s)`$ $`=`$ $`[s+\mathrm{\Omega }_3^2(k)\stackrel{~}{M}_3(k,s)]^1`$ (18a)
$`=`$ $`{\displaystyle \frac{s+\mathrm{\Omega }_4^2(k)\stackrel{~}{M}_3(k,s)}{s^2+\mathrm{\Omega }_4^2(k)\stackrel{~}{M}_3(k,s)s+\mathrm{\Omega }_3^2(k)}},`$
$$\stackrel{~}{M}_3(k,s)=\frac{s+\sqrt{s^2+4\mathrm{\Omega }_4^2(k)}}{2\mathrm{\Omega }_4^2(k)},$$
(18b)
which are obtained by Laplace transform of the third and fourth ($`j=3,4`$) equations of the chain (10).
To pass from the frequency dependence of $`\stackrel{~}{M}_2(k,i\omega )`$ to the time one, let us consider the low frequency region restricted by the value $`2\mathrm{\Omega }_4(k)`$. For convenience we introduce here a small parameter (at the fixed wave number $`k`$):
$$\xi =\frac{s^2}{4\mathrm{\Omega }_4^2},|\xi |1.$$
(19)
Taking into account the fact that the found values of $`\mathrm{\Omega }_4^2(k)`$ for liquid sodium and lithium achieve $`10^{29}10^{30}s^2`$ for the low-$`k`$ region, we span by introducing parameter $`\xi `$ the frequency (time) range $`\omega <10^{15}s^1`$ ($`t>10^{15}s`$), which is important for us and is available experimentally.
Expanding the radicand in Eq. (18b) as a series in the parameter $`\xi `$
$$\sqrt{1+\xi }=1+\frac{\xi }{2}\frac{\xi ^2}{8}+\mathrm{},$$
(20)
we can rewrite it in the following way
$$\stackrel{~}{M}_3(s)=\frac{s}{2\mathrm{\Omega }_4^2}+\frac{1}{\mathrm{\Omega }_4}+\frac{s^2}{8\mathrm{\Omega }_4^3}\frac{s^4}{32\mathrm{\Omega }_4^5}+\mathrm{}.$$
(21)
By restricting the number of terms in the series (20) \[and, accordingly, in Eq. (21)\] we receive from Eq. (18a) the linear combination of the Lorentz functions
$$\stackrel{~}{M}_2(k,s)=\underset{j}{}\frac{A_j(k)}{s+\tau _j^1(k)},j=1,2,3,5,\mathrm{},$$
(22)
the number of which will be increased at the increase of the number of terms in the series (20). The quantities $`A_j(k)`$ and $`\tau _j(k)`$ are expressed by the relaxation frequencies $`\mathrm{\Omega }_3^2(k)`$ and $`\mathrm{\Omega }_4^2(k)`$. Going over to the time scale by the inverse Laplace transform Abramov we obtain
$$M_2(k,t)=\underset{j}{}A_j(k)e^{t/\tau _j(k)}.$$
(23)
By restricting the first term of the series (20) only we receive the simplest model from the first equality of Eq. (18a) with Eq. (21)
$$M_2(k,t)=e^{t/\tau (k)},$$
(24)
which corresponds to the viscoelastic model with the relaxation time $`\tau (k)=\mathrm{\Omega }_4(k)/\mathrm{\Omega }_3^2(k)`$, and from the second equality of Eq. (18a) the double exponential model, i.e Eq. (23) at $`j=2`$, with the following time relaxation parameters
$`\tau _{1,2}(k)=\left[\mathrm{\Omega }_4(k)\pm \sqrt{\mathrm{\Omega }_4^2(k)\mathrm{\Omega }_3^2(k)}\right]^1`$ (25)
and the weight factor
$$A(k)=\frac{\mathrm{\Omega }_4(k)+\sqrt{\mathrm{\Omega }_4^2(k)\mathrm{\Omega }_3^2(k)}}{2\sqrt{\mathrm{\Omega }_4^2(k)\mathrm{\Omega }_3^2(k)}}.$$
(26)
This case may be related to the double-time viscous model Tullio1 ; Tullio0 , two-time exponential ansatz Levesque ; Egelstaff . In the general form Eq. (23) corresponds to the framework of generalized collective mode approach Schepper with the sum of the weighed exponents for the TCF $`M_2(k,t)`$, where $`\tau _j^1(k)`$ denote eigenvalues of a generalized dynamic matrix with the elements consisting of static correlation functions, and the weight factors $`A_j(k)`$ are the amplitudes describing the contribution of the corresponding modes.
So, it is obvious that the theory underlying Eq. (15) prescribes such behavior of the second order memory function $`M_2(k,t)`$, which may be represented in the form of Eq. (23) and can be reduced to the above-mentioned models. Eq. (23) is in fact an expansion of $`M_2(k,t)`$ into decay channels embedded in this function.
## III Scale uniformity of dynamics processes in liquid alkali metals
The determination of the scale uniformity of structural and dynamical features for different groups of liquids is very important for the physics of liquid state. On the one hand, it allows one to apply the unified theoretical description to the whole group. On the other hand, it allows one to remove the difficulties related to obtaining the experimental data. The fact is that until recently the microscopic dynamics of liquids could be experimentally probed by INS only. However, there were often different problems related to, firstly, separation of collective and one-particle contributions, and, secondly, gross experimental errors (and even with impossibility to obtain data) for different ($`k,\omega )`$-regions. Recent progress in the technique of IXS has allowed one to clear some of the obstacles Burkel . Ten years ago the possibility of the unified description of the structural and dynamical properties of different liquid alkali metals near the melting point was found by the comprehensive molecular dynamics simulation study Balucani\_PRB , where the adopted potential model of Price, Singwi and Tosi was used, and the scale passage was executed on the basis of the potential parameters. The recent sketchy attempt of testing this outcome experimentally has shown its inconsistency Tullio1 .
In present work we also execute the comparison of the dynamic structure factor spectra of liquid lithium and sodium. As known from the experimental results, the dynamic structure factor $`S(k,\omega )`$ depends strongly on the temperature $`T`$ and the wave number $`k`$. So, one can define the reduced forms of these terms as $`T/T_m`$ and $`k/k_m`$, where $`T_m`$ is the melting temperature and $`k_m`$ is the main peak position in the static structure factor $`S(k)`$ for the corresponding system. The scale time interval $`t^{}`$ can be expressed as $`t^{}=k^1\sqrt{m/K_BT}`$. Thus defined time unit $`t^{}`$ is different from the one introduced in Ref. Tullio1 , because the present term varies with the change of space and temperature characteristics. Although we do not exclude the possibility, that this scale unit may be independent of the temperature and the wave number for other systems (for instance, semi-conductors, or H-bonded liquids).
In Fig. $`2`$ we report the comparison of $`S(k,\omega )`$ spectra for liquid lithium and sodium Tullio0 ; Tullio1 at approximately the same reduced temperatures $`T/T_m`$ and wave numbers $`k/k_m`$. Namely, $`T/T_m=1.049`$ for liquid lithium and $`1.051`$ in case of sodium. From this figure one can see that dynamic structure factor practically coincides in the first two higher cases. From the lower plot of Fig. $`2`$ one can see, that the position of inelastic and central peaks for both systems is the same. However, though the overall coincidence of spectra is observed at intermediate frequencies only, the peak altitudes are a little different. Such deviation can easily be explained by the fact that the plot for liquid lithium is presented for a higher value of the reduced wave number, $`0.75`$, whereas in case of sodium $`k/k_m=0.73`$. As known, these wave numbers correspond to the so-called the de Gennes narrowing region characterized by a strong $`k`$dependence. In other words, a higher section at $`k`$ of flat $`S(k,\omega )/t^{}`$ is presented for lithium than for sodium. It is necessary to take into account, that the values of the reduced temperatures for both systems are also slightly different.
Notice that the time unit $`t^{}`$ depends on the system features ($`m`$), on the probed spatial region ($`k`$) and the temperature regime ($`T`$) in contrast to scale units $`k_m`$ and $`T_m`$, which remain unchanged the spatial region and the temperature of the system is revised.
At result, the experimental or theoretical $`S(k,\omega )`$ for any single metal allows one easily to restore this term for the whole group of alkali metals at same reduced conditions, $`k/k_m`$ and $`T/T_m`$. Moreover, the theory developed for the concrete separate alkali metal may be simply extended to the whole group.
As an example, in Fig. $`3`$ we report the dynamic structure factor of liquid potassium $`S^K(k,\omega )`$ obtained on the basis IXS data for liquid sodium $`S^{Na}(k,\omega )`$. The transition $`S^{Na}(k,\omega )S^K(k,\omega )`$ has been executed by means of the following scale reductions:
$$S^K(k,\omega )=S^{Na}(k,\omega )\frac{k^{Na}}{k^K}\sqrt{\frac{m^KT^{Na}}{m^{Na}T^K}},$$
(27a)
$$\omega ^K=\omega ^{Na}\frac{k^K}{k^{Na}}\sqrt{\frac{m^{Na}T^K}{m^KT^{Na}}},$$
(27b)
$$T^K=\frac{T_m^KT^{Na}}{T_m^{Na}},k^K=\frac{k_m^Kk^{Na}}{k_m^{Na}}.$$
(27c)
By the top subscript we note the corresponding system ($`K`$ or $`Na`$).
## IV Concluding remarks
The following results are presented in this work.
(i) The theory, developed on the basis of Bogoliubov ideas about the hierarchy of relaxation times, allows one to obtain dynamic structure factor, reproducing adequately experimental IXS spectra for liquid alkali metals (in particular, for liquid lithium and sodium) in the region of low values of wave number.
(ii) The expansion of the second order memory function into exponential decay channels, used (sometimes intuitively) in others theories, may be easily obtained within the framework of the presented approach. This is the evidence of the multi-mode character of decay of the observed relaxation process.
(iii) An important result of this work is the confirmation of the proposition about the unitary description of the dynamical features of liquid alkali metals, and finding of corresponding scale transition relations.
## V Acknowledgments
The authors are grateful to T. Scopigno for providing IXS data, and to M. H. Lee for fruitful discussions. A. V. M. acknowledges G. Garberoglio and F. A. Oliveira for the useful correspondence. This work was supported by the Russian Ministry of Education and Science (Grants No. 03-06-00218a, A03-2.9-336) and RFBR (No. 02-02-16146).
## VI Figure captions
Fig. 1. Dynamic structure factor of liquid lithium at the temperature $`T=475`$K. The solid lines are the results of the theoretic model (15), whereas the open circles are the IXS data Tullio0 . The theoretical lineshapes have been modified to account for the quantum mechanical detailed balance condition and broadened for the finite experimental resolution effects as described in the text. The wave numbers $`k`$ are given in a reduced form, where $`k_m`$ is the main peak position in the static structure factor $`S(k)`$.
Fig. 2. IXS spectra of liquid lithium at $`T=475`$K ($`\mathrm{}\mathrm{}\mathrm{}`$) and liquid sodium at $`T=390`$K ($``$) Tullio0 ; Tullio1 in the reduced units. The scale frequency $`\omega ^{}`$ is chosen as the term inverse proportional to $`t^{}`$.
Fig. 3. The dynamic structure factor $`S(k,\omega )`$ of liquid potassium at $`T=354.1`$K calculated from IXS data of liquid sodium at $`T=390`$K by the scale reduction described in the text.
Fig. 1
Fig. 2
Fig. 3
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# Bounds for the 𝑏-chromatic number of some families of graphs
## 1 Introduction and related results
All graphs in this paper are finite, simple and undirected graphs. By clique number of a graph $`G`$ we mean the largest order of a complete subgraph in $`G`$ and denote it by $`\omega (G)`$. Also $`\alpha (G)`$ stands for the largest number of independent vertices in $`G`$. For other notations which are not defined here we refer the reader to .
An antimatching of a graph $`G`$ is a matching of its complement.
A proper coloring of $`G`$ is a coloring of the vertices such that any two adjacent vertices have different colors. Given a proper coloring of $`G`$, a $`t`$-dominating set $`T=\{x_1,\mathrm{},x_t\}`$ is a set of vertices such that $`T`$ is colored by $`t`$ colors and each $`x_i`$ is adjacent to $`t1`$ vertices of different colors. In that case, and if $`G`$ is colored by exactly $`t`$ colors, we say we have a $`t`$-dominating coloring (or $`b`$-coloring with $`t`$ colors). We denote by $`\phi (G)`$ the maximum number $`t`$ for which there exists a $`t`$-dominating set in a coloring of $`V(G)`$ by $`t`$ colors. This parameter has been defined by Irving and Manlove , and is called the $`b`$-chromatic number of $`G`$. In a $`b`$-coloring of a graph $`G`$ with $`b`$ colors, any vertex $`v`$ which has at least $`b1`$ neighbors with different colors is called a representative. We note that in any $`b`$-coloring of $`G`$ with $`b`$ colors there should be at least $`b`$ representatives with $`b`$ different colors.
It is known that $`\chi (G)\phi (G)\mathrm{\Delta }+1`$. Let $`G`$ be a graph with decreasing degree sequence $`d(x_1)d(x_2)\mathrm{}d(x_n)`$ and let $`m(G)=max\{i:d(x_i)(i1)\}`$. In , the authors proved that for any graph $`G`$, $`\phi (G)m(G)`$ and they show that for tree $`T`$ the inequality $`m(T)1\phi (G)m(T)`$ is satisfied.
Also in it is shown that determining $`\phi `$ is NP-hard for general graphs, but polynomial for trees.
Some authors have obtained upper or lower bounds for $`\phi (G)`$ when $`G`$ belongs to some special families of graphs. In , $`b`$-chromatic number of graphs with girths five and six has been studied. Let $`G`$ be a graph of girth at least $`5`$, of minimum degree $`\delta `$ and of diameter $`D`$, it is shown in that $`\phi (G)>min\{\delta ,D/6\}`$ and that if $`G`$ is $`d`$-regular, of girth at least six, then $`\phi (G)=d+1`$. In this last case the construction of a $`b`$-dominating coloring is done in a polynomial time.
Kratochvil et al. in showed that for a $`d`$-regular graph $`G`$ with at least $`d^4`$ vertices, $`\phi (G)=d+1`$.
In , Kouider and Mahéo discuss on the $`b`$-chromatic number of the cartesian product $`G\mathrm{}H`$ of two graphs $`G`$ and $`H`$. They prove that $`\phi (G\mathrm{}H)\phi (G)+\phi (H)1`$ when $`G`$ (resp. $`H`$) admits $`\phi (G)`$ (resp. $`\phi (H)`$) dominating set which is stable set.
We also recall the following result of Klein and Kouider . Let $`𝒟`$ be $`K_4e`$. Let $`G`$ be a $`P_4`$-free graph, then $`\phi (G)=\omega (G)`$, for any induced subgraph of $`G`$ if and only if $`G`$ is $`2𝒟`$-free and $`3P_3`$-free.
The aim of this paper is to obtain an upper bound for $`b`$-chromatic number of a graph $`G`$ when $`G`$ is restricted to be in special families of graphs. In section $`2`$ we consider $`K_{1,t}`$ -free graphs. In section $`3`$ we give an upper bound in terms of clique number and minimum clique partition of a graph. Finally in section $`4`$ bipartite graphs will be considered. We also show that all the bounds obtained in this paper are tight.
## 2 $`K_{1,t}`$ -free graphs
In this section we give an upper bound for the $`b`$-chromatic number of $`K_{1,t}`$ -free graphs, when $`t3`$. If $`t=2`$ then the graph should be a complete graph for which the $`b`$-chromatic number is the same as chromatic number.
###### Theorem 1
. Let $`G`$ be a $`K_{1,t}`$ -free graph where $`t3`$, then $`\phi (G)(t1)(\chi (G)1)+1`$.
Proof. Suppose $`\phi (G)=b`$. Let $`C`$ be a color class in a $`b`$-coloring of $`G`$ with $`b`$ colors, and let $`x`$ be any representative of the class $`C`$. Among the neighbors of the vertex $`x`$ there exist a set say $`S`$ of $`b1`$ vertices with distinct colors. Let $`H`$ be the subgraph induced by $`S`$. By the hypothesis on the graph $`G`$ we have $`\alpha (H)t1`$ and also $`\chi (H)\chi (G)1`$. So
$$b1=|V(H)|\alpha (H).\chi (H)(t1)(\chi (G)1).$$
Therefore $`b(t1)(\chi (G)1)+1`$. $`\mathrm{}`$
In the following we show that the bound of the theorem can be achieved for each $`t`$.
###### Proposition 1
. For any integer $`t3`$ and $`k`$, there exists a $`K_{1,t}`$ -free graph $`G`$ such that $`\chi (G)=k`$ and $`\phi (G)=(t1)(k1)+1`$.
Proof. Suppose the graph $`H`$ is defined as a vertex $`v`$ such that its neighbors form $`t1`$ mutually disjoint cliques with $`k1`$ vertices. Now we take $`(t1)(k1)+1`$ disjoint copies of $`H`$ and connect them sequentially by exactly one edge between any two consecutive copies. These edges can be incident to any vertex other than $`v`$ and its copies in other copies of $`H`$. We denote the resulting graph by $`G`$. It is easily seen that $`G`$ satisfies the conditions of theorem. $`\mathrm{}`$
We have now the following immediate corollary of theorem 1.
###### Corollary
. If $`G`$ is a claw-free graph, then $`\phi (G)2\chi (G)1`$.
In the important fact $`\chi (G)2\omega (G)`$ is proved for a claw-free graph $`G`$ satisfying $`\alpha (G)3`$, therefore using this result we obtain $`\phi (G)4\omega (G)1`$.
## 3 $`b`$-coloring and minimum clique partition
In this section we give a bound for the $`b`$-chromatic number of a graph $`G`$ in terms of its minimum clique partition. A clique partition for a graph $`G`$ is any partition of $`V(G)`$ into subsets say $`C_1,C_2,\mathrm{},C_k`$ in such a way that the subgraph of $`G`$ induced by $`C_i`$ is a clique, for each $`i`$. We denote by $`\theta (G)`$ the minimum number of subsets in a clique partition of the graph $`G`$. We note that for any graph $`G`$, $`\chi (\overline{G})=\theta (G)`$; also, if $`\theta (G)=k`$ then $`G`$ is the complement of a $`k`$-partite graph. Therefore the following result applies for all graphs.
###### Theorem 2
. Let $`G`$ be a graph with clique partition number $`\theta (G)=k`$ and clique number $`\omega `$, then $`\phi (G){\displaystyle \frac{k^2\omega }{2k1}}`$.
Proof. If $`k=1`$ then $`G`$ is complete and equality holds in the inequality of theorem. We suppose now $`k2`$. As $`\theta (G)=k`$, therefore $`\alpha (G)k`$. Let us consider a $`b`$-coloring of $`G`$ with $`\phi (G)=b`$ colors. Let $`i_j`$ be the number of color classes with exactly $`j`$ elements. As $`\alpha (G)k`$, we know that $`i_j=0`$ for $`jk+1`$. So we have
$$b=\underset{j=1}{\overset{k}{}}i_j.$$
By hypothesis, there exists a partition of $`V(G)`$ into $`k`$ complete subgraphs, therefore if $`n`$ is the order of $`G`$,
$$n=\underset{j=1}{\overset{k}{}}j.i_j=b+\underset{j=2}{\overset{k}{}}(j1)i_jk\omega .(\mathrm{𝟏})$$
Suppose first that $`i_1=0`$. Then any color class in the $`b`$-coloring of $`G`$ with $`b`$ colors contains at least two vertices. This shows that $`bn/2`$ and so $`bk\omega /2`$. Finally $`b{\displaystyle \frac{k^2}{2k1}}\omega `$, because $`{\displaystyle \frac{k}{2}}{\displaystyle \frac{k^2}{2k1}}`$.
Suppose now $`i_11`$ and let $`C_i=\{x_i\}`$ for $`i=1,\mathrm{},i_1`$. Then any representative of any color $`j`$ is adjacent to any $`x_i`$, where $`i,ji_1`$ and $`ij`$.
It follows that $`\{x_1,\mathrm{},x_{i_1}\}`$ induces a complete subgraph of $`G`$. On the other hand, by the fact that there exists a partition of $`V(G)`$ into $`k`$ cliques and the pigeonhole principle, at least $`{\displaystyle \frac{_{j=2}^ki_j}{k}}`$ of representative vertices form a complete graph. We know from above that any representative of any color $`j`$ is adjacent to any $`x_i,ij,ii_1`$, consequently there is a complete subgraph of at least $`i_1+{\displaystyle \frac{_{j=2}^ki_j}{k}}`$ vertices. We get the following inequality
$$i_1+\frac{_{j=2}^ki_j}{k}\omega $$
in other words,
$$ki_1+\underset{j=2}{\overset{k}{}}i_jk\omega .(\mathrm{𝟐})$$
Now we have
$$(2k1)b=\underset{j=1}{\overset{k}{}}(2k1)i_j=(k1)(\underset{j=1}{\overset{k}{}}ji_j)+ki_1+i_2\underset{j=3}{\overset{k}{}}((k1)j2k+1)i_jfork3,$$
or
$$(2k1)b=(k1)(\underset{j=1}{\overset{k}{}}ji_j)+ki_1+i_2fork=2.$$
So we have
$$(2k1)b(k1)n+ki_1+i_2$$
and by inequality (1),
$$(2k1)b(k1)k\omega +ki_1+i_2k^2\omega k(\omega i_1\frac{i_2}{k}).$$
By inequality (2),
$$(2k1)bk^2\omega .$$
The theorem is proved. $`\mathrm{}`$
###### Proposition 2
For any positive integers $`k2`$ and $`\omega `$ divisible by $`2k1`$, there exists a graph $`G`$ with $`\theta (G)=k`$ and with clique number $`\omega `$, such that $`\phi (G)={\displaystyle \frac{k^2\omega }{2k1}}`$.
Proof. In order to construct our graph we first consider three sets of mutually disjoint cliques $`\{A_1,\mathrm{},A_k\}`$, $`\{B_1,\mathrm{},B_k\}`$ and $`\{C_1,\mathrm{},C_k\}`$ where $`|A_i|={\displaystyle \frac{\omega }{2k1}}`$, $`|B_i|=|C_i|={\displaystyle \frac{(k1)\omega }{2k1}}`$, for each $`i=1,\mathrm{},k`$. We put an edge between any two vertices $`u`$ and $`v`$ in $`A_i`$ and $`A_j`$ for each $`i`$ and $`j`$, therefore $`{\displaystyle \underset{i}{}}A_i`$ forms a clique with $`{\displaystyle \frac{k\omega }{2k1}}`$ vertices. Then we join any vertex in $`A_i`$ to any vertex in $`B_j`$ for each $`i`$ and $`j`$, and also we join the vertices of $`A_i`$ to all the vertices of $`C_i`$, for each $`i`$. We don’t have any edge between any two vertices of $`B_i`$ and $`B_j`$ when $`ij`$ and the same holds for $`C_i`$’s. Finally we put an edge between any two vertices $`vB_i`$ and $`uC_j`$ if $`ij`$.
We color the vertices in $`{\displaystyle \underset{i}{}}A_i`$ with $`1,2,\mathrm{},{\displaystyle \frac{k\omega }{2k1}}`$ and the vertices of $`{\displaystyle \underset{i}{}}B_i`$ with distinct colors $`{\displaystyle \frac{k\omega }{2k1}}+1,\mathrm{},{\displaystyle \frac{k^2\omega }{2k1}}`$. The colors in $`C_i`$ will be the same as $`B_i`$ for each $`i`$. All the vertices of $`A={\displaystyle \underset{i}{}}A_i`$ are representatives and the same holds for $`B={\displaystyle \underset{i}{}}B_i`$.
Now it is enough to show that the constructed graph $`G`$ has the clique number $`\omega `$. We first observe that if we identify each of cliques $`A_i`$’s, $`B_i`$’s and $`C_i`$’s with single vertices $`a_i`$’s, $`b_i`$’s and $`c_i`$’s, respectively, then we may define a graph $`H`$ with $`3k`$ vertices with vertex set $`\{a_1,\mathrm{},a_k,b_1,\mathrm{},b_k,c_1,\mathrm{},c_k\}`$ where there is an edge between two vertices $`u`$ and $`v`$ if and only if their corresponding cliques are jointed in the graph $`G`$. Therefore to find the maximum number of vertices in a clique of the graph $`G`$, it is enough to check all cliques in $`H`$. Let us first set $`A=\{a_1,\mathrm{},a_k\}`$, $`B=\{b_1,\mathrm{},b_k\}`$ and $`C=\{c_1,\mathrm{},c_k\}`$. Let $`K`$ be a clique in $`H`$. There are two possibilities:
1. There is no vertex from $`C`$ in $`K`$. In this case $`K`$ may contain all vertices in $`A`$ and at most one from $`B`$, i.e. with at most $`k+1`$ vertices. This clique results in a clique in $`G`$ with $`{\displaystyle \frac{k\omega }{2k1}}+{\displaystyle \frac{(k1)\omega }{2k1}}=\omega `$ vertices.
2. There is one vertex from $`C`$ in $`K`$. In this case $`K`$ contains only one vertex from $`C`$ and at most one vertex from $`A`$ and one from $`B`$. And this may happen when we consider for example $`a_1`$ and its neighbor in $`C`$ and a suitable vertex in $`B`$. This clique of order three results in a clique in $`G`$ with $`{\displaystyle \frac{(k1)\omega }{2k1}}+{\displaystyle \frac{\omega }{2k1}}=\omega `$ vertices. $`\mathrm{}`$
The following result is an immediate corollary of theorem 2.
###### Corollary
. For any graph $`G`$, with clique-number $`\omega (G)`$,
$$\phi (G)\frac{\chi ^2(\overline{G})}{2\chi (\overline{G})1}\omega (G).$$
In the case that $`G`$ is the complement of a bipartite graph we have more knowledge on its $`b`$-colorings. We first introduce some special graphs which play an important role in $`b`$-colorings of the complement of bipartite graphs. Before we begin let us mention that when we say there is an anti-matching between two subsets $`X`$ and $`Y`$ in a graph $`G`$, it means that there exists a matching between $`X`$ and $`Y`$ in the complement of $`G`$.
Let $`G`$ be the complement of a bipartite graph with a bipartition $`(X,Y)`$ in such a way that there are partitions of $`X`$ and $`Y`$ into three subsets as $`X=A_1B_1C_1`$ and $`Y=A_2B_2C_2`$ such that the following properties hold:
1. Any vertex in $`A_1`$ is adjacent to any vertex in $`A_2B_2`$, hence the subgraph induced by $`A_1A_2B_2`$ in $`G`$ is a clique. Also any vertex in $`A_2`$ is adjacent to any vertex in $`C_1`$.
2. $`|B_1|=|B_2|`$ and there is a perfect anti-matching between $`B_1`$ and $`B_2`$.
3. $`|C_1|=|C_2|`$ and there is a perfect anti-matching between $`C_1`$ and $`C_2`$.
In this case by letting $`b=|A_1A_2|+|B_1|+|C_1|=|X|+|A_2|`$, we say $`G`$ belongs to the family $`𝒜_b`$. In fact $`𝒜_b`$ consists of all the complement of bipartite graphs $`G`$ which admits the above-mentioned properties. Let us remark that $`\phi (G)b`$ for any graph $`G`$ belonging to $`𝒜_b`$. In fact, we color $`XA_2`$ with different colors; using the antimatchings, we give to $`B_2`$ the same colors as $`B_1`$, and to $`C_2`$ the same colors as $`C_1`$.
###### Theorem 3
. Let $`G`$ be the complement of a bipartite graph, then $`\phi (G){\displaystyle \frac{4\omega }{3}}`$. Furthermore, there is a $`b`$-coloring for $`G`$ with $`b`$ colors if and only if $`G`$ is in $`𝒜_b`$.
Proof. The inequality $`\phi (G){\displaystyle \frac{4\omega }{3}}`$ follows from theorem 2 where we put $`k=2`$.
If $`G`$ is in $`𝒜_b`$ then by the comment before theorem 3 there is a $`b`$-coloring for $`G`$ with $`b`$ colors.
Suppose now we have a $`b`$-coloring for $`G=(X,Y;E)`$ with $`b`$ colors $`\{1,2,\mathrm{},b\}`$. Let the color classes be $`U_1`$, $`U_2`$, …, $`U_b`$ and without loss of generality we may suppose that $`|U_i|=1`$ for $`i=1,\mathrm{},t`$. Therefore $`|U_i|=2`$ for $`i>t`$. Set $`A_1=X_{i=1}^tU_i`$ and $`A_2=Y_{i=1}^tU_i`$.
Let $`u_i,i=t+1,\mathrm{},s`$ be the representatives contained in $`X`$ and they form a set $`B_1`$; let $`u_i,i=s+1,\mathrm{},b`$ be the remaining representatives, these are by definition in $`YA_2`$ and they form a set $`C_2`$. As each color class for $`it+1`$, has exactly $`2`$ elements, there exists a set $`B_2`$ in $`Y`$ with $`|B_2|=|B_1|`$ with the same colors as $`B_1`$. Similarly there exists a set $`C_1`$ in $`Y`$ with $`|C_2|=|C_1|`$ with the same colors as $`C_2`$.
There are perfect anti-matchings, one between $`B_1`$ and $`B_2`$ and another between $`C_1`$ and $`C_2`$. By the property of being representative for each element of $`B_1C_2`$, and by the unicity of the elements colored by the colors of $`A_1A_2`$, $`A_1C_2`$ is a clique, $`A_2B_1`$ is also a clique. Considering now the partitions $`X=A_1B_1C_1`$ and $`Y=A_2B_2C_2`$ we conclude that $`G`$ belongs to $`𝒜_b`$. $`\mathrm{}`$
We get easily the following consequence.
###### Corollary
. Let $`G`$ be the complement of a bipartite graph. Then
$$\phi (G)=b$$
if and only if $`\mathrm{max}\{k:G𝒜_k\}=b`$.
Let us remark that for the larger class of graphs $`G`$ with $`\alpha (G)=2`$, there is no linear bound for $`b`$-chromatic number (even for chromatic number) in terms of $`\omega (G)`$ because, as pointed out in , or each $`k`$ there is a graph $`G`$ with $`\alpha (G)=2`$ such that $`\chi (G)k/2`$ and $`\omega (G)=o(k)`$.
## 4 Bipartite graphs
In this chapter we suppose $`G`$ is a bipartite graph. In the following by the biclique number of $`G`$ we mean the minimum number of mutually disjoint complete bipartite graphs which cover the vertices of $`G`$. Any subgraph of $`G`$ which is complete bipartite graph is called a biclique of $`G`$.
###### Theorem 4
. Let $`G`$ be a bipartite graph with bipartition $`(X,Y)`$, on $`n`$ vertices and biclique number $`t`$. Then
$$\phi (G)\frac{nt+4}{2}.$$
Proof. We first prove the theorem for graphs $`G=(X,Y)`$ which admits a $`b`$-coloring with $`b=\phi (G)`$ colors where there is at least one representative in $`X`$ and also one in $`Y`$. Let these representatives be $`vX`$ and $`uY`$. Then $`v`$ has at least $`b1`$ neighbors in $`Y`$ and also $`u`$ has at least $`b1`$ neighbors in $`X`$. These give us two bicliques with cardinality at least $`2b2`$ and at most $`2b`$. As $`t`$ is the biclique number of $`G`$ there should be at least $`t2`$ vertices in $`G`$. Therefore $`n2b2+t2`$ and so $`b{\displaystyle \frac{nt+4}{2}}`$.
Now we may suppose that in a $`b`$-coloring of $`G`$, all the representatives are in a same part say $`X`$. Let $`i_j`$ be the number of color classes in the $`b`$-coloring with exactly $`j`$ colors in part $`Y`$. There are two possibilities.
Suppose first that $`i_11`$. Let $`w`$ be the vertex in any color class with cardinality one in the part $`Y`$. Then $`w`$ belongs to $`Y`$ and has $`b1`$ neighbors which are representatives of different colors. So $`w`$ is representative. This is a contradiction with the hypothesis on $`X`$.
Now let $`p`$ be the minimum number with $`i_p0`$. So $`p2`$. We have $`nb+bp=b(p+1)`$. We may suppose at this stage that all vertices in $`X`$ are representatives and of different colors, and, also any vertex $`y`$ in $`Y`$ is adjacent to, at least, some representative and is the unique vertex of color c(y) of this representative. Otherwise, if we delete those vertices in $`X`$ which are not representatives and also vertices in $`Y`$ without the previous property, we prove the inequality of the theorem, for the resulting graph $`G^{}`$. Let $`nl`$ be its order. We have
$$\phi (G^{})\frac{nlt^{}+4}{2}.$$
As the inequality $`tt^{}+l`$ holds, it can be seen easily that we get
$$\phi (G^{})\frac{nt+4}{2}.$$
We also have, by construction of $`G^{}`$, $`\phi (G)\phi (G^{})`$.
So it is enough to prove the theorem for the case where all the vertices in $`X`$ are representative and any vertex in $`Y`$ is adjacent to some representative. By these hypothesis, and as the coloring is proper, we have $`tb`$. Finally since $`nb(p+1)`$ and $`p2`$ then $`2bn(p1)bnbnt`$. Therefore $`b{\displaystyle \frac{nt}{2}}`$. $`\mathrm{}`$
###### Proposition 3
. For any integer $`p3`$, there is a bipartite graph $`G`$ with $`n=3p4`$ vertices and biclique number $`t=p1`$ such that $`b(G)=p={\displaystyle \frac{nt+4}{2}}`$.
Proof. We first consider a complete bipartite graph $`K_{p1,p1}`$ minus a matching with size $`p2`$. We color one part say $`X`$ of this graph with $`1,3,4,\mathrm{},p`$ and other part say $`Y`$ with $`2,3,4,\mathrm{},p`$ so that vertices with colors $`1`$ and $`2`$ are adjacent. Then we add $`p2`$ extra vertices to the part $`X`$ and color all of them with $`2`$. Now put a matching with size $`p2`$ between these extra vertices in $`X`$ and all the vertices in $`Y`$ except the one colored by $`2`$. The resulting graph $`G`$ is a graph of order $`n=3p4`$ with a $`b`$-coloring with $`p`$ colors. In fact, $`b(G)`$ is exactly equal to $`p`$ because $`\mathrm{\Delta }(G)=p1`$. By the precedent theorem,
$$b(G)\frac{nt+4}{2}.$$
It is then enough to show that $`tp1`$ to get the reverse inequality. Because there are $`p2`$ vertices with degree one, at least $`p2`$ cliques are required to cover these vertices. We observe that we need an extra clique to cover the vertex colored by $`2`$ in $`Y`$. Now we get the equality $`b(G)={\displaystyle \frac{nt+4}{2}}`$. $`\mathrm{}`$
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# Cavity Approach to the Random Solid State
(June 7, 2005)
## Abstract
The cavity approach is used to address the physical properties of random solids in equilibrium. Particular attention is paid to the fraction of localized particles and the distribution of localization lengths characterizing their thermal motion. This approach is of relevance to a wide class of random solids, including rubbery media (formed via the vulcanization of polymer fluids) and chemical gels (formed by the random covalent bonding of fluids of atoms or small molecules). The cavity approach confirms results that have been obtained previously via replica mean-field theory, doing so in a way that sheds new light on their physical origin.
Introduction—Permanent random chemical bonds, when introduced in sufficient number between the (atomic, molecular or macromolecular) constituents of a fluid, cause a phase transition—the vulcanization transition—to a new equilibrium state of matter: the random solid state. In this state, some fraction of the molecules are spontaneously localized, and thus undergo thermal fluctuations about mean positions, the collection of which positions are random (i.e. exhibit no long-range regularity). As this is a problem of statistical mechanics (the constituents are undergoing thermal motion) in the presence of quenched randomness (the constraints imposed by the random bonding), it has been addressed via the replica technique. Indeed, the example of cross-linked macromolecular matter served as early motivation for Edwards’ development of the replica technique (see, e.g., Ref. REF:GGinSTG ).
Our purpose in this Letter is to consider two central diagnostics of the random solid state: the fraction $`Q`$ of localized particles and the statistical distribution $`𝒩(\xi ^2)`$ of squared localization lengths $`\xi ^2`$ of these localized particles. Specifically, we show how results for these quantities can be obtained in an elementary and physically transparent way, via the cavity method. The cavity method has proven flexible and powerful in the analysis of a variety of other disordered systems, e.g., spin glasses \[REF:MPVbook, \]. The present work is based on the version used to address spin glasses having finite connectivity \[REF:MP2001-finite-conn, \].
The results that we shall obtain via the cavity method are amongst those already known via a (less elementary and less physically transparent) application of the replica technique, together with a mean-field approximation; for reviews see Refs. \[CGZ-AdvPhy, ,G-Trieste, \]. Thus, it was already known \[Castillo, 1994\] that (as is typical for mean-field theories) $`Q`$ obeys a transcendental equation, in this case
$$1Q=\mathrm{exp}(\mu ^2Q),$$
(1)
where $`\mu ^2`$ is a parameter that controls the density of cross-links. The instability of the fluid state and its replacement by the random solid state are signaled by the emergence, as $`\mu `$ is increased beyond a critical value (here unity), of a positive solution to Eq. (1), although the formulation is not restricted to the critical regime. This result for $`Q`$ was, in essence, found by Erdös and Rényi, in their classic work on the statistical properties of random graphs \[REF:ER, \]. As for the distribution $`𝒩(\xi ^2)`$, at the mean-field level and for near-critical values of the cross-link density, the replica approach yields the scaling form
$$𝒩(\xi ^2)=\frac{2}{ϵ\xi ^4}\pi \left(\frac{2}{ϵ\xi ^2}\right),$$
(2)
where $`ϵ`$ \[$`(\mu ^21)`$\] measures the distance from the critical point, and in which the scaling function $`\pi `$ obeys the nonlinear integro-differential equation
$$\frac{d}{d\theta }\left({\scriptscriptstyle \frac{1}{2}}\theta ^2\pi (\theta )\right)=\pi (\theta )\left(\pi \pi \right)(\theta ),$$
(3)
where $`\pi \pi `$ indicates a Laplace convolution REF:Normal . This result was obtained from a semi-microscopic approach by Castillo et al. \[Castillo, 1994\] and re-derived via a Landau-type theory by Peng et al. \[Peng, 1998\]. Some support for the results for $`Q`$ and $`\pi `$ has been obtained numerically \[Barsky, 1996\] and experimentally \[REF:Dinsmore, \].
More recently, the mean-field level results have been improved in two directions: via renormalization-group analysis in the vicinity of the upper critical dimension (which is six for this random solidification transition) \[Peng, 2000; REF:Janssen+Stenull, \], and via analysis of the Goldstone fluctuations \[REF:elast-swagatam, \] (notably in two dimensions, where their consequences are—not surprisingly—dramatic). We pause to mention that we view the cavity method as complementing rather than supplanting the replica method: the latter provides access to the powerful array of field-theoretic tools, the former opens the way to a perturbative treatment of correlations beyond mean-field.
Cavity method for randomly cross-linked macromolecules, especially in the vicinity of the random solidification transition—We begin by considering a system of vulcanized macromolecules, as depicted in Fig. 1a. We characterize the system by the fraction $`Q`$ of localized chain segments and the statistical distribution $`𝒩`$ of mean squared localization lengths $`\xi ^2`$ of the localized particles. We then envisage adding a further macromolecule to the system, as shown in Fig. 1b. Of all the segments on this chain, we suppose that a certain number $`N_\mathrm{c}`$ are sufficiently close to segments of the original system to have a chance of becoming cross-linked to them. We suppose that fluctuations in this number are sufficiently small that we may neglect them. Next, we consider a random cross-linking process that results, with independent probability $`p`$, in cross-links actually being introduced between each of the $`N_\mathrm{c}`$ close pairs. Within this framework, the probability that exactly $`k`$ cross-links are introduced is then given by the binomial formula: $`\left(\genfrac{}{}{0pt}{}{N_\mathrm{c}}{k}\right)p^k(1p)^{N_\mathrm{c}k}`$. We now ask the question: What is the probability $`𝒫_k`$ that exactly $`k`$ of these cross-links are made to localized segments? To answer it, we make the approximation that the probabilities for the segments of the original system to be localized describe independent random variables, in which case the probability of any one such segment being localized is $`Q`$. Then, collecting together the contributions to this probability, which arise from $`k^{}`$ (with $`kk^{}N_\mathrm{c}`$) cross-links being formed, of which exactly $`k`$ are to localized segments, we arrive at the formula
$$𝒫_k=\underset{k^{}=k}{\overset{N_\mathrm{c}}{}}\left(\genfrac{}{}{0pt}{}{N_\mathrm{c}}{k^{}}\right)p^k^{}(1p)^{N_\mathrm{c}k^{}}\left(\genfrac{}{}{0pt}{}{k^{}}{k}\right)Q^k(1Q)^{k^{}k}.$$
(4)
Via a straightforward application of the binomial theorem one can perform this summation, and hence arrive at the result
$$𝒫_k=\left(\genfrac{}{}{0pt}{}{N_\mathrm{c}}{k}\right)(pQ)^k(1pQ)^{N_\mathrm{c}k}.$$
(5)
Let us evaluate these probabilities for the three cases of relevance, viz.,
$`𝒫_0`$ $`=`$ $`(1pQ)^{N_\mathrm{c}},`$ (6a)
$`𝒫_1`$ $`=`$ $`N_\mathrm{c}pQ(1pQ)^{N_\mathrm{c}1},`$ (6b)
$`𝒫_2`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}N_\mathrm{c}(N_\mathrm{c}1)(pQ)^2(1pQ)^{N_\mathrm{c}2}.`$ (6c)
In the limit of main interest, viz. $`pQ1`$, these probabilities simplify to
$$(𝒫_0,𝒫_1,𝒫_2)\mathrm{e}^{N_\mathrm{c}pQ}(1,N_\mathrm{c}pQ,{\scriptscriptstyle \frac{1}{2}}N_\mathrm{c}(N_\mathrm{c}1)(pQ)^2).$$
(7)
To arrive at a self-consistent equation for $`Q`$ (as a function of $`p`$ and $`N_\mathrm{c}`$) we require that the probability of the added macromolecule being cross-linked to exactly zero localized segments be $`1Q`$, which gives
$$1Q=\left(1pQ\right)^{N_\mathrm{c}}.$$
(8)
In the limit $`pQ1`$ this becomes
$$1Q=\mathrm{exp}\left(N_\mathrm{c}pQ\right),$$
(9)
i.e., Eq. (1), provided we make the (physically sensible) identification $`pN_\mathrm{c}\mu ^2`$.
To arrive at a self-consistent equation for $`𝒩`$ (as a function of $`p`$, $`N_\mathrm{c}`$ and $`\xi ^2`$) we address the probability that the segment at arc-length $`\sigma `$ (with $`0\sigma 1`$) on the added macromolecule has squared localization length $`\xi ^2`$. As the segments of the added macromolecule will only be localized if they are attached to at least one localized segment of the original system, we should replace the probabilities $`\{𝒫_k\}_{k=0}^2`$ by the probabilities conditioned on at least one attachment being to a localized segment of the original system (i.e. $`k=1,2`$). Thus we arrive at the probabilities $`(\widehat{𝒫}_1,\widehat{𝒫}_2)`$ which we write in the form $`(1a,a)`$. Consider the situation in which the added macromolecule is attached at its arc-length $`\sigma _1`$ to a single segment of the original system, that segment having squared localization length $`\xi _1^2`$, as depicted in Fig. 2a. Furthermore, suppose that the added chain is Gaussian. Then, by the elementary properties of random walks, the mean squared spatial separation of segments separated by arc-length $`|\sigma \sigma _1|`$ is given by $`|\sigma \sigma _1|_\mathrm{g}^2`$, where $`_\mathrm{g}^2`$ is the mean squared end-to-end distance of each chain. Thus, if $`𝒩`$ is the distribution for the squared localization length of the segment of the original system to which the new chain is attached at arc-length $`\sigma _1`$, the distribution of the squared localization lengths $`\xi ^2`$ for the segment at arc-length $`\sigma `$ on the new chain will be given by
$$𝑑\xi _1^2𝒩(\xi _1^2)\delta \left(\xi ^2\left(\xi _1^2+|\sigma \sigma _1|_\mathrm{g}^2\right)\right).$$
(10)
Now, supposing that the addition to the squared localization length, $`|\sigma \sigma _1|_\mathrm{g}^2`$, is small compared with the localization lengths that feature with appreciable weight in $`𝒩`$, this approximates to
$$𝒩\left(\xi ^2|\sigma \sigma _1|_\mathrm{g}^2\right)𝒩(\xi ^2)|\sigma \sigma _1|_\mathrm{g}^2𝒩^{}(\xi ^2).$$
(11)
Next consider the situation in which the added macromolecule is attached at its arc-lengths $`\sigma _1`$ and $`\sigma _2`$ to two segments of the original system, these segments having respective squared localization length $`\xi _1^2`$ and $`\xi _2^2`$, as shown in Fig. 2b. In fact, as the probability that the added chain has two cross-links to localized segments of the original system is (in the limit of interest) much smaller than the probability that it has one cross-link, we shall not need to keep track of the arc-length locations of the cross-links; in this situation it will be adequate to treat the added chain as a point object. Then, as this object is attached to two localized objects, it is pinned more sharply than either, this parallel form of pinning giving rise to a smaller squared localization length $`\xi ^2`$, via the formula
$$\xi ^2=\xi _1^2+\xi _2^2.$$
(12)
So, assuming that $`\xi _1^2`$ and $`\xi _2^2`$ are independent \[and thus governed by the joint distribution $`𝒩(\xi _1^2)𝒩(\xi _2^2)`$\], the distribution of $`\xi ^2`$ is given by
$$𝑑\xi _1^2𝒩(\xi _1^2)𝑑\xi _2^2𝒩(\xi _2^2)\delta \left(\xi ^2\left(\xi _1^2+\xi _2^2\right)^1\right).$$
(13)
We now put these results together to construct the distribution of squared localization lengths for segment $`\sigma `$ of the added chain, arriving at
$`(1a)\left(𝒩(\xi ^2)|\sigma \sigma _1|_\mathrm{g}^2𝒩^{}(\xi ^2)\right)`$
$`+a`$ $`{\displaystyle 𝑑\xi _1^2𝒩(\xi _1^2)𝑑\xi _2^2𝒩(\xi _2^2)\delta \left(\xi ^2\left(\xi _1^2+\xi _2^2\right)^1\right)}`$ (14)
Finally, we average over the segment $`\sigma `$ of the added chain, as well as the location $`\sigma _1`$ of the cross-link, using $`𝑑\sigma 𝑑\sigma ^{}|\sigma \sigma ^{}|=1/3`$, thus arriving at the self-consistent equation obeyed by the distribution of squared localization lengths:
$`𝒩(\xi ^2)=(1a)\left(𝒩(\xi ^2){\scriptscriptstyle \frac{1}{3}}_\mathrm{g}^2𝒩^{}(\xi ^2)\right)`$
$`+a`$ $`{\displaystyle 𝑑\xi _1^2𝒩(\xi _1^2)𝑑\xi _2^2𝒩(\xi _2^2)\delta \left(\xi ^2\left(\xi _1^2+\xi _2^2\right)^1\right)}`$ (15)
Observe that by integrating both sides over $`\xi ^2`$ and invoking the property that $`𝒩`$ vanishes at the limits $`\xi ^2=0`$ and $`\mathrm{}`$, we recover the normalization condition that $`𝑑\xi ^2𝒩(\xi ^2)=1`$. The scaling property of $`𝒩`$ shows up via the following change of dependent and independent variables:
$`\xi ^2`$ $``$ $`\theta {\displaystyle \frac{2}{3}}{\displaystyle \frac{1a}{a}}{\displaystyle \frac{_\mathrm{g}^2}{\xi ^2}},`$ (16a)
$`𝒩(\xi ^2)`$ $``$ $`\pi (\theta ){\displaystyle \frac{3}{2}}{\displaystyle \frac{a}{1a}}{\displaystyle \frac{\xi ^2}{_\mathrm{g}^2}}\xi ^2𝒩(\xi ^2),`$ (16b)
under which Eq. (Cavity Approach to the Random Solid State) becomes the sought integro-differential equation (3). We identify the parameter $`ϵ`$ in Eq. (2) with $`\frac{3a}{(1a)_\mathrm{g}}`$. Notice that the cavity method allows us to compute corrections to Eq. (3) perturbatively in $`1/_\mathrm{g}`$.
Cavity method for randomly bonded Brownian particles at arbitrary bonding densities—The cavity approach can be extended to address the chemical gelation transition, i.e., the transition triggered by the introduction of random covalent bonds between atoms or small molecules (rather than macromolecules) in the liquid state. We shall address the model studied previously by Broderix et al. BWZ-gelation , which consists of a collection of point particles undergoing Brownian motion at a certain temperature. Permanent bonds are then introduced at random between nearby particles, so that pairs of bonded particles become constrained softly (i.e. by a spring-like, harmonic potential), so that the probability distribution $`\varphi `$ of their separations $`𝐫_j𝐫_k`$ is Gaussian and characterized by a length-scale $`l`$:
$$\varphi (𝐫_j𝐫_k)\mathrm{exp}\left(|𝐫_j𝐫_k|^2/2l^2\right),$$
(17)
where $`𝐫_{j,k}`$ are the position vectors of particles $`j`$ and $`k`$. This model is sufficiently simple that the analysis of it need not be restricted to the critical region.
To approach the statistics of this system of randomly bonded Brownian particles using the cavity method, we consider the process of adding a new particle. The combinatorics of the bonding follows the form taken for the system of cross-linked macromolecules; we simply need to convert the notion of contact points into a spherical region of a certain radius in which the likelihood of particles being bonded to one another is concentrated. This sphere is centered on the new particle and, on average, includes $`N_\mathrm{c}`$ of the existing particles. Then, bonds are randomly introduced, with probability $`p`$, between the new particle and some of the $`N_\mathrm{c}`$ existing particles that are nearby. Thus, the foregoing combinatorics continues to apply, and we arrive at the formula for the probability of having exactly $`k`$ bonds with the infinite cluster given in Eq. (5). As a consequence, we obtain the foregoing result for the fraction of the infinite cluster $`Q`$, Eq. (9).
The physics of the localization lengths is, in fact, simpler for bonded Brownian particles. When the new particle is connected via a spring of length-scale $`l`$ to one localized particle, its localization length $`\mathrm{\Xi }_1`$ is given by $`\mathrm{\Xi }_1^2=(\xi _1^2+l^2)^1`$. When it is connected in parallel via identical springs to $`k`$ localized particles, its localization length $`\mathrm{\Xi }_k`$ is given by
$$\frac{1}{\mathrm{\Xi }_k^2}\underset{j=1}{\overset{k}{}}\frac{1}{\xi _j^2+l^2}.$$
(18)
To construct the distribution of the squared localization length of the new particle, we shall average over all possible numbers of bonds, weighted with their corresponding probabilities. These probabilities follow from the probabilities $`𝒫_k`$ given by Eq. (5), but normalized by a factor $`Q^1`$ because the new particle will only be localized if it is bonded to at least one particle in the infinite cluster. Hence, we arrive at a self-consistency equation for localization-length distribution for the randomly bonded Brownian particle model:
$$𝒩(\xi ^2)=\underset{k=1}{\overset{\mathrm{}}{}}\widehat{𝒫}_k_0^{\mathrm{}}𝑑\xi _1^2𝒩(\xi _1^2)\mathrm{}𝑑\xi _k^2𝒩(\xi _k^2)\delta (\xi ^2\mathrm{\Xi }_k^2),$$
(19)
where the conditional probabilities $`\widehat{𝒫}_k`$ are given by $`\widehat{𝒫}_k=𝒫_k/Q`$ (for $`k=1,2,3,\mathrm{}`$). The distribution of localization lengths for the randomly bonded Brownian particle model was studied previously by Broderix et al. BWZ-gelation using the replica method and a Mayer cluster expansion. To see that the cavity approach result (19) recovers their result, we take the limit $`N_\mathrm{c}\mathrm{}`$ whilst keeping finite the mean number of bonds from the new particle (either to the infinite cluster or to delocalized particles), i.e. $`pN_\mathrm{c}`$. In this limit, the binomial distribution tends to a Poisson distribution: $`\widehat{𝒫}_k\left((pN_\mathrm{c})^kQ^{k1}/k!\right)\text{e}^{pN_\mathrm{c}Q}.`$ The result of Broderix et al. BWZ-gelation then follows from Eq. (19) by (i) transforming to a distribution for $`\tau 1/\xi ^2`$, and (ii) making the identifications $`\kappa =l^2`$ and $`c=pN_\mathrm{c}`$ (i.e. the mean of number of bonds associated with a single particle).
Acknowledgments—We thank H. Castillo and A. Zippelius for valuable discussions. PG, MM and MW thank for its hospitality the Kavli Institute for Theoretical Physics at the University of California–Santa Barbara, where this work was begun. PG also thanks for its hospitality the Laboratoire de Physique Théorique et Modèles Statistiques at l’Université Paris Sud, where it was completed. This work was supported by NSF DMR02-05858 (XM, PMG) NSF PH99-07949 (PG, MM, MW).
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# 1 Introduction
## 1 Introduction
The nuclear reaction mechanism of fragmentation has been investigated for more than 60 years. A turning point in this study was marked by Jacobsson et al. in 1982, who measured multiple fragment production in nuclear emulsions containing Ag and Br, irradiated with $`{}_{}{}^{12}\mathrm{C}`$ at 55$`A`$ and 110$`A`$ MeV/nucleon . These data stimulated development of new models to explain the formation of multiple fragments by a “liquid-gas” phase transition in hot nuclear matter (see, e.g., ).
The isospin dependence in the equation of state of nuclear matter is very important, being at the same time poorly known property of neutron-rich nuclear matter . In recent years, much attention has been paid to the isospin dependence both in nucleus-nucleus experiments with an excess of neutrons in the bombarding and/or target nuclei and in experiments with different types of light projectiles on targets with different neutron/proton ratios . Such investigations may help obtain information about the equation of state of the asymmetric nuclear matter.
Many experiments have been devoted to the study of nuclear multifragmentation, where several fragments in the mass region $`3Z20`$ are formed from hot nuclear matter . Observation of isoscaling, that is, the dependence of fragment formation probabilities on the third component of their isotopic spins, has increased the possibility of obtaining information on the formation mechanisms of these fragments . Early work in this field was done by Bogatin et al. and has continued .
Recently experimentalists and theorists have focussed on the investigation of formation mechanisms of heavy fragments ($`Z20`$) by different projectiles ($`\gamma `$ \- rays, $`\pi ^{}`$\- meson etc.) . Experiments with a direct registration of heavy fragments do not provide a comprehensive understanding of fragments in this mass region. The induced-radioactivity method adds more possibilities and the investigation of mechanisms of heavy-fragment production becomes more realistic .
The aim of the present work is to investigate the formation of product nuclei on separated tin isotopes by proton and deuteron beams of 3.65 GeV/nucleon in three mass regions of product nuclides: $`7A30`$, $`40A80`$, and $`81A96`$.
## 2 Experimental results and discussion
Targets of enriched tin isotopes $`{}_{}{}^{112}\mathrm{Sn},^{118}\mathrm{Sn},^{120}\mathrm{Sn}`$, and $`{}_{}{}^{124}\mathrm{Sn}`$ are irradiated at the Nuclotron and Synchrophasotron of the LHE JINR in Dubna by proton and deuteron beams with energy of 3.65 GeV/nucleon. The description of the experiment is given in . The measurement of products in the mass range $`7A80`$ is performed by studying the induced activity. The radioactive nuclei obtained are identified by the characteristic $`\gamma `$-rays and their half-lives. For beam monitoring, we employ the reactions <sup>27</sup>Al$`(d,3p2n)^{24}`$Na and <sup>27</sup>Al$`(p,3pn)^{24}`$Na, whose cross sections are taken as of $`14.2\pm 0.2`$ mb and $`10.6\pm 0.8`$ mb , respectively. Cross sections for about 70 products from each target for the proton and deuteron beams are obtained. The measured cross sections of all products are shown in Tables 1 and 2: index “I” denotes independent yields while “C” indicates cumulative ones.
Table 1. Measured product cross sections for p + <sup>112,118,120,124</sup>Sn
| Product | Type | Cross section (mb) | | | |
| --- | --- | --- | --- | --- | --- |
| | | $`{}_{}{}^{112}\mathrm{Sn}`$ | $`{}_{}{}^{118}\mathrm{Sn}`$ | $`{}_{}{}^{120}\mathrm{Sn}`$ | $`{}_{}{}^{124}\mathrm{Sn}`$ |
| <sup>7</sup>Be | I | 13.9$`\pm `$1.5 | 9.4$`\pm `$0.3 | 8.2$`\pm `$ 1.4 | 7.5$`\pm `$ 0.8 |
| <sup>22</sup>Na | C | 2.3$`\pm `$0.3 | 2.4$`\pm `$0.4 | 2.1$`\pm `$ 0.4 | 1.7$`\pm `$0.2 |
| <sup>24</sup>Na | C | 3.25$`\pm `$0.3 | 3.23$`\pm `$0.2 | 3.69$`\pm `$0.3 | 3.97$`\pm `$0.3 |
| <sup>28</sup>Mg | C | 0.39$`\pm `$0.06 | 0.53$`\pm `$0.05 | 0.75$`\pm `$0.08 | 0.89$`\pm `$0.07 |
| <sup>38</sup>Cl | I | | 1.67$`\pm `$ 0.2 | 1.5$`\pm `$ 0.2 | |
| <sup>39</sup>Cl | C | | 0.57$`\pm `$0.02 | 0.34$`\pm `$0.07 | |
| <sup>42</sup>K | C | 1.76$`\pm `$0.11 | 1.85$`\pm `$0.25 | 1.95$`\pm `$0.2 | 2.1$`\pm `$0.2 |
| <sup>43</sup>K | C | 0.74$`\pm `$0.06 | 0.85$`\pm `$ 0.08 | 1.04$`\pm `$ 0.1 | 1.32$`\pm `$ 0.1 |
| <sup>43</sup>Sc | C | 0.72$`\pm `$0.18 | 0.6$`\pm `$0.2 | 0.45$`\pm `$0.2 | 0.2$`\pm `$0.03 |
| <sup>44g</sup>Sc | I | 0.97$`\pm `$ 0.09 | 0.54$`\pm `$0.15 | 0.58$`\pm `$0.04 | 0.36$`\pm `$0.09 |
| <sup>44m</sup>Sc | I | 2.28$`\pm `$0.1 | 1.45$`\pm `$0.06 | 1.4$`\pm `$0.07 | 1.7$`\pm `$0.1 |
| <sup>46</sup>Sc | I | 2.2$`\pm `$ 0.2 | 2.35$`\pm `$0.2 | 2.6$`\pm `$0.4 | 2.4$`\pm `$0.2 |
| <sup>48</sup>Sc | I | 0.3$`\pm `$ 0.05 | 0.38$`\pm `$0.07 | 0.42$`\pm `$0.05 | 0.7$`\pm `$0.09 |
| <sup>48</sup>Cr | C | 0.19$`\pm `$0.07 | 0.1$`\pm `$ 0.01 | 0.13$`\pm `$ 0.04 | |
| <sup>51</sup>Cr | C | | 3.7$`\pm `$0.6 | 3.3$`\pm `$0.6 | |
| <sup>48</sup>V | I | 2.7$`\pm `$ 0.15 | 1.68$`\pm `$ 0.1 | 1.76$`\pm `$ 0.1 | 1.06$`\pm `$ 0.1 |
| <sup>52</sup>Mn | C | 1.9$`\pm `$0.04 | 1.12$`\pm `$0.03 | 1.03$`\pm `$0.04 | 0.7$`\pm `$0.08 |
| <sup>54</sup>Mn | I | 6.1$`\pm `$0.3 | 5.1$`\pm `$0.25 | 4.8$`\pm `$0.3 | 4.2$`\pm `$0.3 |
| <sup>56</sup>Mn | C | 0.8$`\pm `$ 0.03 | 1.08$`\pm `$0.07 | 1.33$`\pm `$0.1 | 1.86$`\pm `$0.08 |
| <sup>59</sup>Fe | C | 0.37$`\pm `$0.05 | 0.83$`\pm `$ 0.09 | 0.85$`\pm `$ 0.07 | 1.17$`\pm `$ 0.1 |
Table 1 (continued)
| Product | Type | Cross section (mb) | | | |
| --- | --- | --- | --- | --- | --- |
| | | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| <sup>56</sup>Co | C | 1.8$`\pm `$0.1 | 1.7$`\pm `$ 0.2 | 1.75$`\pm `$ 0.2 | 1.54$`\pm `$0.2 |
| <sup>58</sup>Co | C | 5.7$`\pm `$ 0.4 | 4.8$`\pm `$ 0.5 | 4.8$`\pm `$ 0.4 | 3.6$`\pm `$0.4 |
| <sup>67</sup>Cu | C | 0.1$`\pm `$0.02 | 0.29$`\pm `$0.05 | 0.44$`\pm `$0.04 | 0.52$`\pm `$0.05 |
| <sup>65</sup>Zn | C | 9.1$`\pm `$0.3 | 6.9$`\pm `$0.3 | 6.9$`\pm `$0.3 | 5.1$`\pm `$0.3 |
| <sup>66</sup>Ga | C | 4.8$`\pm `$0.4 | 3.1$`\pm `$0.2 | 2.9$`\pm `$0.3 | 2.5$`\pm `$0.3 |
| <sup>67</sup>Ga | C | 8.9$`\pm `$0.07 | 6.7$`\pm `$ 0.4 | 6.4$`\pm `$ 0.4 | 5.3$`\pm `$ 0.5 |
| <sup>69</sup>Ge | C | 7.5$`\pm `$ 0.7 | 5.3$`\pm `$ 0.5 | 5.0$`\pm `$ 0.3 | 3.9$`\pm `$ 0.5 |
| <sup>77</sup>Ge | C | | 0.2$`\pm `$ 0.05 | 0.16$`\pm `$ 0.05 | |
| <sup>70</sup>As | C | 3.0$`\pm `$ 0.5 | 1.9$`\pm `$ 0.5 | 1.42$`\pm `$ 0.2 | 2.1$`\pm `$ 0.6 |
| <sup>71</sup>As | C | 8.01$`\pm `$0.8 | 5.54$`\pm `$ 0.06 | 5.3$`\pm `$ 0.4 | 4.1$`\pm `$ 0.5 |
| <sup>72</sup>As | C | 2.8$`\pm `$ 0.4 | 2.1$`\pm `$ 0.3 | 1.9$`\pm `$ 0.4 | 2.1$`\pm `$ 0.6 |
| <sup>74</sup>As | I | 1.92$`\pm `$0.15 | 2.6$`\pm `$ 0.3 | 3.07$`\pm `$ 0.4 | 3.5$`\pm `$ 0.25 |
| <sup>76</sup>As | I | 4.5$`\pm `$ 0.4 | 5.1$`\pm `$ 0.4 | 5.3$`\pm `$ 0.5 | 6.3$`\pm `$ 0.4 |
| <sup>73</sup>Se | C | 5.7$`\pm `$ 0.15 | 3.8$`\pm `$ 0.15 | 3.5$`\pm `$ 0.15 | 2.4$`\pm `$ 0.2 |
| <sup>75</sup>Se | C | 13.2$`\pm `$0.5 | 10.3$`\pm `$ 0.7 | 10.1$`\pm `$ 1.0 | 8.8$`\pm `$ 0.7 |
| <sup>76</sup>Br | C | 10.5$`\pm `$1 | 7.3$`\pm `$0.5 | 6.6$`\pm `$0.7 | 5.2$`\pm `$0.4 |
| <sup>77</sup>Br | I | 10.4$`\pm `$0.4 | 8.4$`\pm `$0.2 | 8.4$`\pm `$0.2 | 7.8$`\pm `$0.2 |
| <sup>82</sup>Br | I | | 0.26$`\pm `$0.04 | 0.25$`\pm `$0.02 | 0.49$`\pm `$0.03 |
| <sup>76</sup>Kr | C | 1.6$`\pm `$ 0.2 | 1.03$`\pm `$ 0.05 | 0.9$`\pm `$ 0.07 | 0.55$`\pm `$0.05 |
| <sup>77</sup>Kr | C | | 3.2$`\pm `$ 0.3 | 2.7$`\pm `$ 0.3 | 1.76$`\pm `$0.16 |
| <sup>81</sup>Rb | C | 16.4$`\pm `$0.6 | 11.7$`\pm `$0.2 | 11.5$`\pm `$0.4 | 8.8$`\pm `$0.6 |
| <sup>82m</sup>Rb | I | 5.2$`\pm `$0.05 | 5.3$`\pm `$0.4 | 5.9$`\pm `$0.4 | 5.7$`\pm `$0.3 |
| <sup>83</sup>Rb | C | 18.6$`\pm `$ 0.7 | 15.3$`\pm `$ 0.5 | 16.0$`\pm `$ 0.6 | 13.8$`\pm `$ 0.3 |
| <sup>84g</sup>Rb | I | 1.4$`\pm `$ 0.3 | 2.3$`\pm `$ 0.2 | 2.8$`\pm `$ 0.2 | 4.4$`\pm `$ 0.6 |
| <sup>86</sup>Rb | I | 0.34$`\pm `$0.09 | | 1.09$`\pm `$0.14 | 1.85$`\pm `$0.24 |
| <sup>83</sup>Sr | C | 14.6$`\pm `$0.1 | 10.3$`\pm `$0.3 | 9.8$`\pm `$0.2 | 7.4$`\pm `$0.5 |
| <sup>85</sup>Sr | C | 21$`\pm `$1.6 | 17.5$`\pm `$1.7 | 17.3$`\pm `$1.5 | 15.2$`\pm `$0.9 |
| <sup>84m</sup>Y | I | 5.1$`\pm `$0.5 | 3.4$`\pm `$0.3 | 2.6$`\pm `$0.4 | 2.1$`\pm `$0.4 |
| <sup>86m</sup>Y | I | 4.9$`\pm `$ 0.4 | 6.4$`\pm `$ 0.1 | 6.7$`\pm `$ 0.4 | 5.5$`\pm `$ 0.4 |
| <sup>87m</sup>Y | C | 18.6$`\pm `$ 0.7 | 15.5$`\pm `$ 0.9 | 15.8$`\pm `$ 0.3 | 12.8$`\pm `$ 0.4 |
| <sup>87g</sup>Y | I | 4.3$`\pm `$ 0.4 | 3.5$`\pm `$ 0.3 | 4.0$`\pm `$ 0.3 | 2.8$`\pm `$ 0.2 |
| <sup>86</sup>Zr | C | 8.8$`\pm `$ 0.5 | 4.8$`\pm `$ 0.15 | 3.5$`\pm `$ 0.1 | 2.3$`\pm `$ 0.25 |
| <sup>88</sup>Zr | C | 20.2$`\pm `$2.0 | 13.8$`\pm `$0.8 | 14.4$`\pm `$1.0 | 10.2$`\pm `$0.9 |
| <sup>89</sup>Zr | C | 20.2$`\pm `$0.5 | 16.35$`\pm `$0.3 | 16.4$`\pm `$0.3 | 13.1$`\pm `$0.7 |
| <sup>90</sup>Nb | C | 18.2$`\pm `$ 1.0 | 12.4$`\pm `$ 0.4 | 11.6$`\pm `$ 1.2 | 8.6$`\pm `$ 0.4 |
| <sup>95g</sup>Nb | C | | 0.8$`\pm `$ 0.03 | 1.75$`\pm `$ 0.07 | 2.40$`\pm `$ 0.25 |
| <sup>95m</sup>Nb | I | | 0.17$`\pm `$0.08 | 0.35$`\pm `$ 0.06 | |
| <sup>96</sup>Nb | I | 0.33$`\pm `$0.08 | 0.42$`\pm `$ 0.07 | 0.65$`\pm `$0.06 | 1.14$`\pm `$ 0.16 |
| <sup>90</sup>Mo | C | 5.9$`\pm `$0.2 | 2.6$`\pm `$0.3 | 2.1$`\pm `$ 0.3 | 1.1$`\pm `$ 0.1 |
| <sup>93m</sup>Mo | I | 3.5$`\pm `$ 0.2 | 4.1$`\pm `$ 0.4 | 4.4$`\pm `$ 0.3 | 3.8$`\pm `$ 0.2 |
| <sup>99</sup>Mo | C | 0.19$`\pm `$0.02 | 0.26$`\pm `$ 0.02 | 0.62$`\pm `$0.13 | 1.65$`\pm `$0.25 |
| <sup>93</sup>Tc | C | 12.35$`\pm `$0.8 | 6.95$`\pm `$0.3 | 5.7$`\pm `$ 0.5 | 3.4$`\pm `$ 0.3 |
Table 1 (continued)
| Product | Type | Cross section (mb) | | | |
| --- | --- | --- | --- | --- | --- |
| | | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| <sup>94</sup>Tc | I | 9.8$`\pm `$ 0.2 | 6.5$`\pm `$ 0.1 | 6.7$`\pm `$ 0.2 | 4.4$`\pm `$ 0.3 |
| <sup>95g</sup>Tc | I | 12.4$`\pm `$ 0.6 | 9.8$`\pm `$ 0.4 | 8.3$`\pm `$ 0.3 | 7.5$`\pm `$ 0.3 |
| <sup>95m</sup>Tc | I | 1.0$`\pm `$0.1 | 1.1$`\pm `$0.1 | 0.8$`\pm `$0.08 | 0.56$`\pm `$0.10 |
| <sup>96</sup>Tc | I | 4.4$`\pm `$ 0.1 | 6.7$`\pm `$ 0.2 | 7.4$`\pm `$ 0.25 | 6.4$`\pm `$0.1 |
Table 2. Measured product cross sections for d + <sup>112,118,120,124</sup>Sn
| Product | Type | Cross section (mb) | | | |
| --- | --- | --- | --- | --- | --- |
| | | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| <sup>7</sup>Be | I | 31.1$`\pm `$ 2.7 | | 25.5$`\pm `$ 2.5 | 23.1$`\pm `$4.0 |
| <sup>22</sup>Na | C | 23.3$`\pm `$ 0.4 | 8.1$`\pm `$ 1.5 | 4.1$`\pm `$ 0.1 | 3.5$`\pm `$0.9 |
| <sup>24</sup>Na | C | 6.2$`\pm `$ 0.4 | 10.0$`\pm `$ 1.3 | 9.9$`\pm `$0.8 | 12.2$`\pm `$ 1.1 |
| <sup>28</sup>Mg | C | 1.0$`\pm `$ 0.1 | 1.8$`\pm `$0.1 | 1.6$`\pm `$ 0.2 | 2.9$`\pm `$ 0.8 |
| <sup>38</sup>S | C | | | 0.37$`\pm `$0.04 | 0.56$`\pm `$0.06 |
| <sup>38</sup>Cl | I | | | 3.0$`\pm `$ 0.2 | 3.5$`\pm `$ 0.5 |
| <sup>39</sup>Cl | C | | | 1.1$`\pm `$ 0.1 | 1.8$`\pm `$ 0.6 |
| <sup>42</sup>K | C | 2.4$`\pm `$ 0.2 | 4.5$`\pm `$ 1.2 | 3.9$`\pm `$ 0.5 | 5.2$`\pm `$ 0.5 |
| <sup>43</sup>K | C | 1.2$`\pm `$ 0.1 | 2.4$`\pm `$ 0.4 | 2.1$`\pm `$0.3 | 3.1$`\pm `$ 0.3 |
| <sup>43</sup>Sc | C | 1.2$`\pm `$ 0.1 | 1.4$`\pm `$0.02 | 1.4$`\pm `$ 0.1 | 2.1$`\pm `$ 0.1 |
| <sup>44m</sup>Sc | I | 2.9$`\pm `$0.1 | 3.2$`\pm `$ 0.4 | 2.7$`\pm `$ 0.7 | 2.0$`\pm `$ 0.3 |
| <sup>44g</sup>Sc | I | 1.7$`\pm `$ 0.4 | 2.0$`\pm `$ 0.3 | 1.5$`\pm `$ 0.2 | 1.5$`\pm `$0.2 |
| <sup>46</sup>Sc | I | 3.3$`\pm `$ 0.8 | 6.1$`\pm `$ 0.8 | 6.1$`\pm `$0.3 | 6.6$`\pm `$ 0.3 |
| <sup>47</sup>Sc | C | 3.5$`\pm `$ 0.2 | | | |
| <sup>48</sup>Sc | I | 0.5$`\pm `$ 0.1 | 1.0$`\pm `$ 0.09 | 1.1$`\pm `$ 0.2 | 1.6$`\pm `$ 0.3 |
| <sup>48</sup>V | C | 3.5$`\pm `$ 0.3 | 3.5$`\pm `$ 0.4 | 3.2$`\pm `$ 0.1 | 2.9$`\pm `$ 0.5 |
| <sup>51</sup>Cr | C | 14.1$`\pm `$ 1.4 | 6.1$`\pm `$ 0.4 | 7.4$`\pm `$0.8 | 5.7$`\pm `$ 0.5 |
| <sup>52g</sup>Mn | C | 2.3$`\pm `$ 0.4 | 2.1$`\pm `$ 0.2 | 2.0$`\pm `$ 0.4 | 1.5$`\pm `$ 0.3 |
| <sup>56</sup>Mn | C | 2.4$`\pm `$ 0.6 | 3.1$`\pm `$ 0.4 | 2.9$`\pm `$ 0.1 | 4.3$`\pm `$ 0.3 |
| <sup>59</sup>Fe | C | 0.68$`\pm `$0.03 | 1.5$`\pm `$ 0.2 | 1.7$`\pm `$ 0.1 | 2.7$`\pm `$0.2 |
| <sup>55</sup>Co | C | 0.35$`\pm `$0.05 | | | |
| <sup>56</sup>Co | C | 1.9$`\pm `$ 0.1 | 1.9$`\pm `$ 0.3 | 1.4$`\pm `$ 0.1 | 1.1$`\pm `$ 0.1 |
| <sup>57</sup>Co | C | 6.9$`\pm `$ 0.2 | 11.8$`\pm `$0.3 | 5.9$`\pm `$ 0.2 | 4.8$`\pm `$ 0.1 |
| <sup>58</sup>Co | I | 8.3$`\pm `$ 0.3 | 9.9$`\pm `$ 0.2 | 7.6$`\pm `$ 1.0 | 7.8$`\pm `$ 0.5 |
| <sup>60</sup>Cu | C | 1.9$`\pm `$ 0.4 | 1.04$`\pm `$0.19 | 1.7$`\pm `$ 0.1 | 0.5$`\pm `$ 0.09 |
| <sup>67</sup>Cu | C | | | 0.34$`\pm `$ 0.04 | 0.34$`\pm `$ 0.01 |
| <sup>62</sup>Zn | C | 1.25$`\pm `$0.05 | | | 0.3$`\pm `$0.03 |
| <sup>65</sup>Zn | C | 16.5$`\pm `$1.0 | | 10.7$`\pm `$ 0.2 | 10.1$`\pm `$ 0.4 |
| <sup>69m</sup>Zn | I | 0.39$`\pm `$ 0.05 | 0.84$`\pm `$0.05 | 1.1$`\pm `$ 0.1 | 1.4$`\pm `$ 0.1 |
| <sup>66</sup>Ga | C | 5$`\pm `$ 0.5 | | 3.3$`\pm `$0.3 | 3.4$`\pm `$ 0.3 |
| <sup>67</sup>Ga | C | 8.6$`\pm `$ 0.3 | 9.7$`\pm `$ 1.3 | 9.3$`\pm `$ 0.7 | 8.7$`\pm `$ 0.3 |
| <sup>73</sup>Ga | C | 0.48$`\pm `$0.12 | | | |
Table 2 (continued)
| Product | Type | Cross section (mb) | | | |
| --- | --- | --- | --- | --- | --- |
| | | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| <sup>69</sup>Ge | C | 6.2$`\pm `$ 0.6 | 8.8$`\pm `$ 1.6 | 7.2$`\pm `$ 0.5 | 6.7$`\pm `$ 0.9 |
| <sup>77</sup>Ge | C | | | 1.3$`\pm `$0.1 | 2.5$`\pm `$ 0.3 |
| <sup>71</sup>As | C | 7.4$`\pm `$ 0.2 | 7.73$`\pm `$0.8 | 6.7$`\pm `$ 0.2 | 5.4$`\pm `$ 0.4 |
| <sup>74</sup>As | I | 2.1$`\pm `$ 0.1 | 2.65$`\pm `$ 0.5 | 4.3$`\pm `$ 0.1 | 6.9$`\pm `$ 0.8 |
| <sup>78</sup>As | C | 1.35$`\pm `$0.05 | | | 0.9$`\pm `$ 0.2 |
| <sup>73g</sup>Se | C | 6.6$`\pm `$ 0.1 | 6.4$`\pm `$ 0.4 | 5.2$`\pm `$ 0.6 | 4.3$`\pm `$ 0.3 |
| <sup>75</sup>Se | I | 7.0$`\pm `$ 0.77 | 9.4$`\pm `$ 1.49 | 10.8$`\pm `$0.93 | 10.3$`\pm `$ 1.06 |
| <sup>75</sup>Br | C | 6.8$`\pm `$ 0.6 | 5.4$`\pm `$ 0.5 | 4.4$`\pm `$ 0.3 | 3.5$`\pm `$ 0.3 |
| <sup>77</sup>Br | I | 4.5$`\pm `$ 0.55 | 6.9$`\pm `$ 1.18 | 6.0$`\pm `$ 0.47 | 6.4$`\pm `$ 0.74 |
| <sup>77</sup>Kr | C | 5.2$`\pm `$ 0.6 | 3.5$`\pm `$ 0.5 | 3.5$`\pm `$ 0.2 | 2.4$`\pm `$ 0.2 |
| <sup>79</sup>Kr | I | 1.2$`\pm `$ 0.13 | | 8.9$`\pm `$1.22 | 8.0$`\pm `$ 1.01 |
| <sup>79</sup>Rb | C | 10.1$`\pm `$1.0 | | 2.8$`\pm `$0.3 | 1.9$`\pm `$0.2 |
| <sup>81</sup>Rb | C | 14.1$`\pm `$ 0.7 | 15.7$`\pm `$ 0.6 | 12.8$`\pm `$ 0.4 | 10.9$`\pm `$ 0.7 |
| <sup>83</sup>Rb | C | 20.3$`\pm `$ 0.8 | 24.15$`\pm `$ 0.95 | 21.7$`\pm `$ 1.0 | 21.6$`\pm `$ 0.6 |
| <sup>84</sup>Rb | I | 1.2$`\pm `$ 0.1 | 3.2$`\pm `$ 0.3 | 4.1$`\pm `$ 0.5 | 6.4$`\pm `$ 0.3 |
| <sup>82</sup>Sr | C | 10.1$`\pm `$1.5 | 8.8$`\pm `$ 1.5 | 6.0$`\pm `$0.6 | 4.6$`\pm `$0.5 |
| <sup>83</sup>Sr | C | 13.1$`\pm `$ 0.3 | 13.7$`\pm `$ 1.2 | 10.3$`\pm `$1.3 | 9.0$`\pm `$ 0.6 |
| <sup>85</sup>Sr | I | | | 13.3$`\pm `$1.55 | 15.9$`\pm `$ 1.4 |
| <sup>85m</sup>Y | C | | | 7.7$`\pm `$ 0.9 | 3.5$`\pm `$ 0.3 |
| <sup>85g</sup>Y | C | | | | 2.9$`\pm `$ 0.7 |
| <sup>86m</sup>Y | I | 5.0$`\pm `$ 0.3 | 8.3$`\pm `$ 0.2 | 7.6$`\pm `$ 0.1 | 6.6$`\pm `$ 0.2 |
| <sup>86g</sup>Y | I | 6.5$`\pm `$ 0.5 | | 8.9$`\pm `$ 0.1 | 8.8$`\pm `$ 0.8 |
| <sup>87g</sup>Y | C | 20.5$`\pm `$ 2.1 | | 20.0$`\pm `$ 2.0 | 19.3$`\pm `$ 2.0 |
| <sup>88</sup>Y | I | | 10.2$`\pm `$1.7 | 5.6$`\pm `$ 0.5 | 7.8$`\pm `$ 0.8 |
| <sup>86</sup>Zr | C | 7.4$`\pm `$0.2 | 5.4$`\pm `$ 0.1 | 3.7$`\pm `$ 0.3 | 2.8$`\pm `$ 0.1 |
| <sup>88</sup>Zr | C | 18.4$`\pm `$ 2.2 | | 16.8$`\pm `$ 0.2 | 14.6$`\pm `$ 0.2 |
| <sup>89</sup>Zr | C | 18.4$`\pm `$ 0.4 | | 17.8$`\pm `$ 0.5 | 16.9$`\pm `$ 0.7 |
| <sup>93m</sup>Tc | I | 10.1$`\pm `$ 0.6 | | 6.6$`\pm `$ 0.5 | 2.0$`\pm `$ 0.2 |
| <sup>94m</sup>Tc | I | 9.5$`\pm `$ 1.0 | | | 5.6$`\pm `$ 0.2 |
| <sup>94g</sup>Tc | I | 4.3$`\pm `$ 0.3 | 1.9$`\pm `$0.46 | | 2.4$`\pm `$ 0.8 |
| <sup>95m</sup>Tc | I | | 15.1$`\pm `$0.7 | | |
| <sup>95g</sup>Tc | I | 1.23$`\pm `$0.05 | 1.23$`\pm `$0.07 | 1.2$`\pm `$0.2 | 1.0$`\pm `$ 0.1 |
| <sup>96g</sup>Tc | I | 4.4$`\pm `$ 0.3 | 9.1$`\pm `$0.6 | | 8.4$`\pm `$ 1.3 |
To reveal the production mechanisms of light nuclei, the experimental results are analyzed from the viewpoint of:
1) exponential dependence of cross sections on the mass and charge numbers;
2) including isospin dependence.
Investigations by many authors have showed that the yields of fragments from various nuclear reactions can be represented as $`\sigma (A_f)A_f^\tau `$ and $`\sigma (Z_f)Z_f^\tau `$, where $`\tau `$ has values of about $`1.5`$$`2`$ depending on the reactions, where $`A_f`$ and $`Z_f`$ are the mass and charge numbers of the fragments. Note that calculations by the Statistical Multifragmentation Model (SMM) for the mass region of fragments discussed here provide an exponential dependence with $`\tau =2.2`$.
The isospin dependence of the available experimental yields points to an isoscaling behavior. In the case of multifragmentation, the ratio of the yields of fragments produced from different targets has an exponential dependence on the number of protons and neutrons of the product isotopes described by the formula :
$$R_{21}(t_3)=Y_2(N,Z)/Y_1(N,Z)=C\mathrm{exp}(\alpha N+\beta Z),$$
(1)
where $`Y(N,Z)`$ is the yield of fragment with $`Z`$ protons and $`N`$ neutrons, and $`t_3=(NZ)/2`$ is the third projection of the fragment isospin. Indices 1 and 2 correspond to different targets with different isotopic compositions, with 2 corresponding to the more neutron-rich target and where $`C`$ is a normalization parameter. In Ref. , the parameters $`\alpha `$ and $`\beta `$ were expressed using the difference of chemical potentials of the two systems as following: $`\alpha =\mathrm{\Delta }\mu _n/T`$, $`\beta =\mathrm{\Delta }\mu _p/T`$, where $`T`$ is the temperature of the excited nucleus.
Since in our measurements we use targets of different isotopes of the same element, we analyze our data with the following formula:
$$R_{21}(t_3)=Y_2(N,Z)/Y_1(N,Z)=\mathrm{exp}(C+Bt_3),$$
(2)
where $`C`$ and $`B`$ are fitting parameters . The parameter $`B`$ is related to the difference of the chemical potentials of protons and neutrons in the fragment and depends on the temperature of the excited nucleus; therefore it may reveal information about the formation mechanism of the corresponding product.
Figure 1 shows the dependence of $`Y_2/Y_1`$ on $`t_3`$ for the entire mass region of product nuclei from proton-induced reactions for different values of the difference in the neutron numbers of considered pairs of targets $`\mathrm{\Delta }N`$. Similar dependences for deuteron-induced reactions are shown in Figure 2. In both these figures, symbols show the measured data while lines show their fit with formula (2).
Tables 3 and 4 present the values of the fitting parameter $`B`$ for different combinations of targets pairs and for different mass regions of product nuclei for proton- and deuteron-induced reactions, respectively. Figure 3 shows the dependence of the parameter $`B`$ on the difference of neutron numbers in a pair of targets, $`\mathrm{\Delta }N`$, for different mass regions of products from proton-induced reactions.
The value of the parameter $`B`$ increases linearly with increasing $`\mathrm{\Delta }N`$. $`B`$ also increases with increasing mass of the product nuclei. The dependence of parameter $`B`$ on the difference of the neutron numbers in a pair of targets, $`\mathrm{\Delta }N`$, is fitted using the following formula:
$$B=k+d\mathrm{\Delta }N,$$
(3)
where $`k=0.036\pm 0.01`$ and $`d=0.094\pm 0.016`$ for the mass region $`7A30`$, $`k=0.0008\pm 0.0001`$ and $`d=0.071\pm 0.005`$ for the mass region $`40A80`$, and $`k=0.113\pm 0.060`$ and $`d=0.033\pm 0.008`$ for the mass region $`A80`$. The value of the parameter $`d`$ changes with the mass number of the products, and could be a factor in understanding the formation mechanism of the final nuclides.
From Tables 3 and 4, we see that for the production of <sup>93-96</sup>Tc and <sup>81-86</sup>Rb on the pair of targets $`{}_{}{}^{124}\mathrm{Sn}/^{112}\mathrm{Sn}`$, the parameter $`B`$ has values of $`1.07\pm 0.32`$ and $`0.94\pm 0.20`$ for proton-induced reactions and $`1.10\pm 0.40`$ and $`1.17\pm 0.29`$ for deuteron-induced reactions, respectively. This agrees with similar values of $`B`$ of $`1.22\pm 0.12`$ and $`1.23\pm 0.13`$ found in the literature for such products at a higher energy of 8.1 GeV . This allows us to conclude that residual products in this mass region are produced via spallation processes of successive particle evaporation.
Table 3. Mean values of the fitting parameter $`B`$ for different target pairs (with the difference in the excess neutron number of $`\mathrm{\Delta }N`$) bombarded by protons
| Product nuclei | $`\mathrm{\Delta }N=2`$ | $`\mathrm{\Delta }N=4`$ | $`\mathrm{\Delta }N=6`$ | $`\mathrm{\Delta }N=8`$ | $`\mathrm{\Delta }N=12`$ |
| --- | --- | --- | --- | --- | --- |
| $`7A30`$ | 0.19$`\pm `$0.03 | 0.27$`\pm `$0.06 | 0.23$`\pm `$0.03 | 0.44$`\pm `$0.10 | 0.51$`\pm `$0.04 |
| $`40A60`$ | 0.11$`\pm `$0.03 | 0.34$`\pm `$0.05 | 0.41$`\pm `$0.04 | 0.56$`\pm `$0.04 | 0.85$`\pm `$0.04 |
| $`70A80`$ | 0.18$`\pm `$0.05 | 0.25$`\pm `$0.13 | 0.36$`\pm `$0.09 | 0.51$`\pm `$0.10 | 0.78$`\pm `$0.21 |
| <sup>81-86</sup>Rb | 0.25$`\pm `$0.02 | 0.32$`\pm `$0.04 | 0.66$`\pm `$0.02 | 0.62$`\pm `$0.15 | 0.94$`\pm `$0.20 |
| <sup>93-96</sup>Tc | 0.24$`\pm `$0.07 | 0.22$`\pm `$0.08 | 0.46$`\pm `$0.15 | 0.85$`\pm `$0.25 | 1.07$`\pm `$0.32 |
Table 4. The same as in Table 3, but for deuteron-induced reactions
| Product nuclei | $`\mathrm{\Delta }N=4`$ | $`\mathrm{\Delta }N=6`$ | $`\mathrm{\Delta }N=8`$ | $`\mathrm{\Delta }N=12`$ |
| --- | --- | --- | --- | --- |
| $`7A30`$ | 0.27$`\pm `$0.05 | 0.26$`\pm `$0.05 | 0.28$`\pm `$0.12 | 0.55$`\pm `$0.07 |
| $`40A60`$ | 0.38$`\pm `$0.1 | 0.41$`\pm `$0.09 | 0.41$`\pm `$0.10 | 0.78$`\pm `$0.14 |
| $`70A80`$ | 0.19$`\pm `$0.07 | 0.5$`\pm `$0.1 | 0.57$`\pm `$0.14 | 0.77$`\pm `$0.19 |
| <sup>81-86</sup>Rb | 0.30$`\pm `$0.08 | | 0.87$`\pm `$0.22 | 1.17$`\pm `$0.29 |
| $`{}_{}{}^{9396}T`$c | 0.15$`\pm `$0.08 | 0.48$`\pm `$0.30 | 0.51$`\pm `$0.07 | 1.10$`\pm `$0.40 |
On the other hand, much smaller values of the fitting parameter $`B`$ in the mass region $`7A30`$ may point to a possible multifragmentation mechanism in the formation of these light fragments .
A different situation may be seen in the mass region $`40A60`$, both for proton- and deuteron-induced reactions. The values of $`B`$ in this mass region is generally lower than for the heavy products <sup>81-86</sup>Rb and <sup>93-96</sup>Tc, but higher than for light fragments with $`7A30`$. This may be understood if we assume that intermediate-mass nuclei are produced not only via evaporation of particles (the spallation mechanism) but also include a contribution from multifragmentation processes. This assumption is in agreement with results of our earlier studies at bombarding proton energies of 0.66, 1.0, and 8.1 GeV: We found that an observed increase in the measured yields of intermediate-mass products can be described in the frameworks of the Intra-Nuclear Cascade (INC) model merged with SMM , i.e., by the INC+SMM model, which considers a contribution of multifragmentation to the formation of such intermediate-mass nuclei.
In the present work, we compare the measured cross sections with predictions by the FLUKA , LAHET , CEM03 , and LAQGSM03 codes (none of them considers the multifragmentation mechanism of fragment production). The first three codes are only applied to the proton-induced reactions, while LAGQSM03 is used for both protons and deuterons. In order to compare the measured cumulative cross sections with calculations, the corresponding theoretical cumulative yields were estimated from the calculated independent cross sections.
Figures 4, 5, and 6 show dependencies of ratios of theoretical to experimental cross sections as functions on the product mass numbers for deuteron- and proton-induced reaction, respectively. We see that, as a rule, all models describe most of the measured cross sections of heavy and medium products within a factor of two. Except for the CEM03 code, the agreement with the measured yields of light fragments is much worse, where the other codes underestimate some measured cross sections by up to two orders of magnitude and more. This could be related to the fact that all the models used here do not consider multifragmentation. But it is also true that they do not include simpler fission/fragmentation production mechanisms, either.
To have a better overall quantitative comparison of experimental data with calculations, we have analyzed our data using the mean deviation factor method suggested first by R. Michel :
$$F=10^{\sqrt{(\mathrm{log}[\sigma ^{cal}/\sigma ^{exp}])^2}},$$
(4)
with its standard deviation
$$S\left(F\right)=10^{\sqrt{\left(\left|\mathrm{log}\left(\sigma ^{cal}/\sigma ^{exp}\right)\right|\mathrm{log}(F)\right)^2}},$$
(5)
where $`<>`$ stands for averaging over all the products included in the comparison.
Values of the average deviation factor $`F`$ and its standard deviation $`S\left(F\right)`$ are listed in Table 5 for deuteron-induced reactions and in Table 6 for reactions with protons, respectively.
Table 5. Mean deviations of product yields calculated by LAQGSM03 from the measured data (parameters $`<F>\pm S(<F>)`$) for deuteron-induced reactions averaged over all compared cross sections
| | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| --- | --- | --- | --- | --- |
| $`<F>\pm S(<F>)`$ | 3.305$`\pm `$3.08 | 2.04$`\pm `$1.69 | 2.96$`\pm `$2.75 | 2.41$`\pm `$2.75 |
Table 6. Mean deviations of theoretical product yields from the measured data (parameters $`<F>\pm S(<F>)`$) for proton-induced reactions averaged over all compared cross sections
| Models used | <sup>112</sup>Sn | <sup>118</sup>Sn | <sup>120</sup>Sn | <sup>124</sup>Sn |
| --- | --- | --- | --- | --- |
| LAHET | 4.07$`\pm `$2.77 | 3.49$`\pm `$2.40 | 3.37$`\pm `$2.31 | 3.61$`\pm `$2.56 |
| FLUKA | 5.92$`\pm `$4.18 | 7.84$`\pm `$5.19 | 8.87$`\pm `$5.42 | 6.97$`\pm `$4.29 |
| LAQGSM03 | 5.10$`\pm `$3.86 | 3.44$`\pm `$2.62 | 3.09$`\pm `$2.22 | 3.16$`\pm `$2.14 |
| CEM03 | 3.66$`\pm `$3.02 | 3.26$`\pm `$2.79 | 4.04$`\pm `$3.29 | 3.60$`\pm `$3.01 |
The present analysis points to a possible formation of light nuclides via multifragmentation, which would suggest a “liquid-gas” phase transition taking place in hot nuclear matter formed by irradiation of target nuclei with high-energy particles. The intermediate-mass products are probably formed mainly via evaporation, but some contribution from multifragmentation is also possible, according to our study.
Acknowledgments
We thank Dr. A. J. Sierk for useful discussions and help. The work was partially supported by the Advanced Simulation Computing (ASC) Program at the Los Alamos National Laboratory operated by the University of California for the U. S. Department of Energy and by the Moldovan-US Bilateral Grants Program, CRDF Project MP2-3045-CH-02 and the NASA ATP01 Grant NRA-01-01-ATP-066.
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# B PHYSICS PROSPECTS AT LHC
## 1 Introduction
First pp collisions are foreseen in the Large Hadron Collider at Cern in summer 2007. At that time several precise measurements will be available from B-Factories and Tevatron experiments to test the CKM paradigm of flavour structure and CP violation. However New Physics could be hidden in B decays, specially those involving box and penguin diagrams. On the other hand, if New Physics will be found at LHC in direct searches, B physics measurements will have to sort out the corresponding flavour structure.
The B Physics program at LHC is rich. It will include precise measurement of $`B_s^0\overline{B_s^0}`$ mixing (mass difference, width difference and phase), precise measurements of the angle $`\gamma `$ (including from processes only at tree level), several measurements of other CP phases in different channels for over-constraining the Unitarity Triangles, search for New Physics effects appearing in rare exclusive and inclusive B decays, studies of b-baryons and B<sub>c</sub> physics and also studies of b production.
B physics at LHC has the great advantage of a high $`b\overline{b}`$ cross section ($`\sigma _{bb}500\mu `$b), several orders of magnitude higher than at the $`\mathrm{{\rm Y}}(4S)`$, and of the production of all species of $`b`$-hadrons, including $`B_s`$, $`b`$-baryons and $`B_c`$. The challenge in the analysis is related to the presence of the underlying event, to the high particle multiplicity and to the high rate of background events (the inelastic cross section is $`80`$ mb ). These features demand to the experiments an excellent trigger capability, with good efficiency also on fully hadronic decay modes of $`b`$-hadrons, excellent tracking and vertexing performance, allowing for high mass resolution and proper time resolution, and excellent particle identification to separate exclusive decays.
## 2 LHC experiments
The LHC will collide protons at 14 TeV with a bunch crossing rate of 40MHz.
Two experiments, ATLAS and CMS, are omni-purpose and optimized for discovery physics. Their B physics program is mainly pursued in the first years of running, when the LHC luminosity is expected to be 1-2$`\times 10^{33}`$cm<sup>-2</sup>s<sup>-1</sup>. In the subsequent years at high luminosity ($`10^{34}`$cm<sup>-2</sup>s<sup>-1</sup>), when several pp collisions per bunch crossing will pile up, only search for very rare B decays with clear signatures will be performed. Reaches in B Physics will depend on the chosen trigger strategy and allocated bandwidth. B events will be mainly triggered by high $`p_T`$ muons or di-muon triggers. CMS also exploits on-line tracking for the selection of exclusive B events at High Level Triggel ( $`^\mathrm{?}`$), while ATLAS foresees a flexible trigger strategy with the progressive addition of other triggers ( $`^\mathrm{?}`$).
LHCb is the LHC experiment dedicated to B physics. It will locally tune the luminosity, by de-focusing the beams, to 2$`\times 10^{32}`$cm<sup>-2</sup>s<sup>-1</sup>, in order to limit pile up of pp interactions. Taking the nominal year period as $`10^7`$s, an integrated luminosity of 2 fb<sup>-1</sup> per year is expected, corresponding to $`10^{12}`$ $`b\overline{b}`$ events/year. LHCb is a single-arm forward detector in the polar region 10-300 mrad, with good acceptance for $`b`$ events due to the forward peaked production of $`b`$-hadrons at LHC. A schematic view of the LHCb detector is shown in Figure 1. A description of the detector and its performances can be found in ( $`^\mathrm{?}`$).
The LHCb trigger is operating in three stages. The Level-0 reduces the rate to 1 MHz requiring the presence of leptons or photons or hadrons with high p<sub>T</sub> while the Level-1 selects on high impact parameter, high p<sub>T</sub> tracks. The High Level Trigger is a software trigger using the full information on the event. Its output contains 200 Hz of exclusive B candidates and about 1.8 KHz of inclusive channels to be used also for calibration purposes and systematic studies.
## 3 Prospects on $`B_s^0\overline{B_s^0}`$ mixing measurements
### 3.1 Measurements of $`\mathrm{\Delta }m_s`$
The $`B_s^0\overline{B_s^0}`$ oscillation has been proven to be too fast to be resolved at LEP and SLC experiments. The best limit today is $`\mathrm{\Delta }m_s>14.5`$ ps<sup>-1</sup> at 95%CL. The Tevatron is at present the only available source of $`B_s^0`$ mesons and CDF and DØ have the chance to find a mixing signal in the coming years. The measurement requires best performances in the event reconstruction and purity, proper time resolution and flavour tagging.
But the definitive answer on $`B_s^0\overline{B_s^0}`$ mixing may come from LHC. The best channel for these studies is $`B_s^0D_s^{}\pi ^+`$. Results of LHCb full simulation indicate a proper time resolution of $`\sigma _\tau 40`$ fs and an annual yield of 80.000 events with a signal over background ratio of about 3. The effective efficiency for flavour tagging is estimated to be about 7%. The expected proper time distribution of tagged events is shown in Figure 2 for two different values of $`\mathrm{\Delta }m_s`$. In one year of data-taking a 5$`\sigma `$ observation of oscillation is expected if $`\mathrm{\Delta }m_s<68`$ ps<sup>-1</sup>. Once observed, the precision to measure $`\mathrm{\Delta }m_s`$ is $`0.01`$ ps<sup>-1</sup> .
ATLAS will also make a $`5\sigma `$ observation of oscillations if $`\mathrm{\Delta }m_s<22`$ ps<sup>-1</sup>, in 10 fb<sup>-1</sup>. Most recent expectation of CMS is lower, due to restriction to the trigger bandwidth allocated to this channel.
### 3.2 $`\varphi _s`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ measurements
The phase $`\varphi _s`$ of $`B_s^0\overline{B_s^0}`$ mixing is expected to be very small in the Standard Model $`\varphi _s=2\chi =2\lambda ^2\eta 0.04`$ resulting in a high sensitivity to possible New Physics contributions in $`bs`$ transitions. Hints of New Physics could also be found in the measurement of the decay width difference between the two CP eigenstates $`\mathrm{\Delta }\mathrm{\Gamma }_s=\mathrm{\Gamma }(B_L)\mathrm{\Gamma }(B_H)`$. In the Standard Model $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ is expected to be of the order of 10%. Both quantities can be measured using $`B_s^0J/\psi \varphi `$ decays $`(J/\psi \mu \mu ,\varphi KK)`$. In a decay to two vector mesons three distinct amplitudes contribute: two CP even and one CP odd. The CP components can be disentangled on a statistical basis taking into account the distribution of the so-called transversity angle $`\theta _{tr}`$, defined as the angle between the positive lepton and the $`\varphi `$ decay plane, in the $`J/\psi `$ rest frame. The physics parameters can be extracted from a simultaneous fit to the proper time, $`cos(\theta _{tr})`$ and $`\mathrm{\Delta }m_s`$ distributions of tagged events. In one year of data-taking LHCb expects to collect 100.000 $`J/\psi (\mu \mu )\varphi `$ decays and to obtain a precision on $`sin(\varphi _s)`$ of about 0.06 and precision on $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ of about 0.018 (for $`\mathrm{\Delta }m_s`$=20 ps<sup>-1</sup>). The sensitivity will be increased adding $`B_s^0J/\psi \eta `$ events, which are pure CP eigenstates. About 7000 events per year are expected in this channel.
CMS and ATLAS, with 30 fb<sup>-1</sup>, expect a sensitivity on $`sin(\varphi _s)`$ of 0.03-0.04 and a sensitivity on $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ of 0.015-0.012, respectively.
## 4 Measurements of $`\gamma `$
### 4.1 $`\gamma `$ measurements from $`B_sD_sK`$ decays
$`B_s^0D_s^\pm K^{}`$ and $`\overline{B_s^0}D_s^{}K^\pm `$ decays can proceed through the tree diagrams shown in Figure 3, which can interfere via mixing. From the measurement of the time-dependent decay asymmetries the phase $`\gamma +\varphi _s`$ can be extracted, together with a strong phase. If $`\varphi _s`$ has been determined otherwise, $`\gamma `$ can be extracted, with little theoretical uncertainty, and insensitive to New Physics.
Strong particle identification capabilities are required to separate $`B_sD_sK`$ decays from $`B_sD_s\pi `$ background having $``$12 times larger branching fraction. The performance of the LHCb RICH detectors will be fully adequate, as it is shown in Figure 4. Monte Carlo studies have shown that 5400 $`D_s^{}K^\pm `$ events will be collected in one year of data-taking, with S/B ratio, estimated from $`b\overline{b}`$ events, larger than 1. Hereafter only limits on S/B are quoted, due to the limited statistics with which the value of B has been determined. The $`D_s^{}K^\pm `$ asymmetries are shown in Figure 3. A sensitivity of $`\sigma _\gamma =14`$ degrees can be obtained if $`\mathrm{\Delta }m_s`$=20 ps<sup>-1</sup>.
### 4.2 $`\gamma `$ measurements from $`B^0D^0K^0`$ decays
A theoretically clean determination of the angle $`\gamma `$ can be performed using $`B^0D^0K^0`$ decays. The tree diagrams of the decay are shown in Figure 5. The method, described in ( $`^\mathrm{?}`$), is based on the measurement of six time-integrated decay rates for $`B_d^0D^0K^0,\overline{D^0}K^0,D_{CP}K^0`$ and their CP conjugates; the decays are self-tagged through $`K^0K^+\pi ^{}`$, while the CP auto-states $`D_{CP}`$ can be reconstructed in $`K^+K^{}`$ and $`\pi ^+\pi ^{}`$ modes. This method is similar to the analysis of $`B^\pm D^0K^\pm `$ decays, already performed at the B-Factories, but has the advantage of using two colour-suppressed diagrams with an expected ratio of amplitudes $`r=|A(B^0D^0K^0)|/|A(B^0\overline{D^0}K^0)|0.4`$.
LHCb in one year of data taking expects to collect a total of about 4.000 signal events leading to a sensitivity on $`\gamma `$ of $`\sigma _\gamma 8`$ degrees.
### 4.3 $`\gamma `$ measurements from $`B_d^0\pi ^+\pi ^{}`$ and $`B_s^0K^+K^{}`$ decays
Several strategies have been proposed ( $`^\mathrm{?}`$) to extract informations on the angle $`\gamma `$ from two body charmless decays of $`B`$ mesons, some of them make use of assumptions on dynamics or on U-spin flavour symmetry.
RICH detectors allow to separate the $`K/\pi `$ channels with high efficiency and purity, as shown in Figure 4. LHCb in one year of data taking expects to collect 26.000 $`B_d^0\pi ^+\pi ^{}`$, 37.000 $`B_s^0K^+K^{}`$ and 135.000 $`B_d^0K^+\pi ^{}`$ decays, with mass resolution $`\sigma (M_B)17`$ MeV and a proper time resolution of $`\sigma _\tau 30`$ fs. The two time dependent CP asymmetries
$$A_{CP}(B_d^0\pi ^+\pi ^{})(t)=A_{CP}^{dir,\pi \pi }cos(\mathrm{\Delta }m_dt)+A_{CP}^{mix,\pi \pi }sin(\mathrm{\Delta }m_dt)$$
$$A_{CP}(B_s^0K^+K^{})(t)=A_{CP}^{dir,KK}cos(\mathrm{\Delta }m_st)+A_{CP}^{mix,KK}sin(\mathrm{\Delta }m_st)$$
will be used to fit the four CP asymmetries, which will be extracted with a precision of about 6%. Following the method suggested in ( $`^\mathrm{?}`$), U-spin symmetry can be exploited to constrain the penguin to tree ratios in the two decays to be the same. Assuming the knowledge of the mixing phases $`\varphi _d`$ and $`\varphi _s`$ from previous measurements, the $`\gamma `$ angle can be extracted with a precision $`\sigma _\gamma 5`$ degrees in one year of data-taking. Additional measurements can be used to test the uncertainty related to the U-spin assumptions.
## 5 $`\alpha `$ measurements from $`B^0\rho \pi `$ decays
A time dependent Dalitz plot analysis of the three body decay $`B^0\rho \pi \pi ^+\pi ^{}\pi ^0`$ allows a clean extraction of the angle $`\alpha =\pi \beta \gamma `$, as suggested in ( $`^\mathrm{?}`$). LHCb expects to reconstruct 14000 decays per year (with S/B$`>1.3`$) in this channel. The Dalitz distribution is shown in Figure 6. An 11-parameter fit has been used to get an independent measurement of tree and penguin parameters, taking into account resonant and non resonant background sources. A sensitivity $`\sigma _\alpha <`$10 degrees can be obtained in one year of data-taking.
## 6 $`B_d^0K^0\mu ^+\mu ^{}`$
$`B_d^0K^0\mu ^+\mu ^{}`$ is a rare decay with branching fraction of the order of $`10^6`$ which has a clear experimental signature. The forward-backward asymmetry is defined as
$$A_{FB}(\text{ŝ})=(_0^1𝑑cos\theta _1^0𝑑cos\theta )\frac{d\mathrm{\Gamma }^2}{d\text{ŝ}dcos\theta }$$
where $`\theta `$ is the angle between the $`\mu ^+`$ and the $`K^0`$ in the di-muon rest frame, and ŝ=$`(m_{\mu ^+\mu ^{}}/m_B)^2`$. The forward-backward asymmetry is a sensitive probe of New Physics. In the Standard Model the value of ŝ for which $`A_{FB}`$(ŝ) is zero can be calculated with a 5% precision. Models with non-standard values of Wilson coefficients $`C_7,C_9,C_{10}`$ predict $`A_{FB}`$(ŝ) of opposite sign or without zero point.
LHCb will select 4400 decays per year (with S/B$`>0.4`$), this allows a determination of branching fractions and CP asymmetries with a precision of few percent. Using a toy Monte Carlo to determine the sensitivity in the forward-backward asymmetry measurement, including background subtraction, an uncertainty of 0.06 on the location of ŝ<sub>0</sub> is found, in 1 year of data-taking.
ATLAS will also collect about 2000 events of $`B_d^0K^0\mu ^+\mu ^{}`$, with S/B$`>7`$ in 30 fb<sup>-1</sup>.
## 7 $`B_s\mu ^+\mu ^{}`$
$`B_s\mu ^+\mu ^{}`$ is a rare decay involving flavour changing neutral currents whose branching ratio is estimated to be $`BR(B_s\mu ^+\mu ^{})=(3.5\pm 0.1)\times 10^9`$ in the Standard Model ( $`^\mathrm{?}`$). In various supersymmetric extensions of the Standard model it can be enhanced by one to three orders of magnitude, being $`BR(tan\beta )^6`$, for large $`tan\beta `$. The best upper limit on the branching ratio at present come from experiments at Tevatron, and reachs few$`\times 10^7`$ at 95% CL.
In the SM context, LHCb expects to select 17 events per year, with a resolution on the B<sub>s</sub> mass of 18 MeV/c<sup>2</sup>. The background determination is still incomplete and require additional Monte Carlo statistics. No events were selected in the $`10^7`$ $`b\overline{b}`$ event sample used so far.
CMS has studied a selection at the High Level Trigger, giving $`B_s\mu ^+\mu ^{}`$ candidates with a rate smaller than 1.7 Hz, and a resolution on the B<sub>s</sub> mass of 74 MeV/c<sup>2</sup>. In 10 fb<sup>-1</sup> 47 signal events are selected. With a refined selection at the offline level 7 signal events are expected to be retained, with less than 1 background event.
For this search ATLAS and CMS will also exploit the high luminosity runs. In 100 fb<sup>-1</sup> (1 year at 10<sup>34</sup>) 92 signal events are expected (with 660 of background) and 26 signal events (with $`6`$ of background) respectively. The different levels of background can be attributed to different vertex reconstruction and selection, however an update on these estimations is expected.
In conclusion there are good prospects of significant measurement in this channel, even for the SM value of the branching ratios.
## 8 Conclusion
In the coming years CP asymmetries will be measured at LHC using several $`B_d^0`$ and $`B_s^0`$ mesons and $`b`$-baryons decay channels. Very rare decays will also be studied, thanks to the high $`b\overline{b}`$ cross section available. This program is complementary to the B-Factories one and will allow to complete and improve the available results and possibly to reveal first signals of new Physics.
## References
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# 1 Introduction
## 1 Introduction
The possibility that Lorentz invariance(LI) may be violated at high energies in 4d with testable consequences has become a subject of much interest in the past few years. If one considers rotational invariance to be sacrosanct due to the strength of existing experimental constraints, Lorentz violation(LV) must take the form of a breakdown of invariance under boosts. Such a scenario is suggestive of spontaneous LV where in some preferred frame, e.g., a time-like 4-vector takes on a vacuum expectation value with components $`(1,0,0,0)`$.
While LV of the type described above have yet to be observed in 4d it is clear that if $`n`$ extra (flat) compact dimensions exist, perhaps near the TeV scale, they obviously behave differently than do the large dimensions with which we are familiar. Clearly the LI in the $`D=4+n`$ dimensional space has been lost through the compactification process; LI suffers further damage if the compactified manifold (here considered to be a torus) is orbifolded. In the usual bottom-up analysis, the $`D`$-dimensional action is conventionally written in a completely LI manner with all LV arising from the compactification/orbifolding process in the IR. Thus, in the UV, such a theory apparently maintains LI. While such an approach may be the simplest one to analyze it certainly does not address the larger issue as to how or why the $`n`$ dimensions have become compact. Some unknown UV physics must have triggered this compactification process making the 4 large dimensions (or, depending on one’s viewpoint, the $`n`$ compactified dimensions) ‘special’ otherwise we would be living in $`4+n`$ dimensions. Thus the true UV physics cannot be completely LI in $`D`$-dimensions. It would be interesting to ask if this UV physics has left any signatures for us to find at accelerators that are beginning to probe the TeV energy regime near the compactification scale, $`R^1`$. If present, how might such effects be parameterized?
For simplicity let us consider the case of one extra dimension, i.e., a theory in 5d with the extra dimension compactified as usual on $`S^1/Z_2`$ with $`R^1<`$ few TeV. If the UV breaking of LI is spontaneous, as in the 4d case discussed above, we can imagine that it was triggered by a 5-vector taking on the vev $`(0,0,0,0,1)`$ in some frame thus leaving us with 4d LI with the fifth dimension becoming ‘special’. Although our approach will not capture any of the detailed dynamics associated with how such a vev was generated, it may be able to probe some of the residual LV physics which could remain accessible at the TeV scale. A framework for this type of analysis already exists for the Standard Model(SM) in 4d, the so-called SM Extension. This framework is particularly appealing for a number of reasons one of which is that it isolates the dominant effects of explicit LV in operators of the lowest possible dimension. In order to have a specific model in which to work in 5d that is closely analogous to this 4d case, we here adapt this particular framework to the 5d version of SM where all fields are in the bulk, i.e., the Universal Extra Dimensions(UED) scenario, but allowing LV to occur only in the fifth dimension.
As we will see in detail below the existence of LV in 5d induces a number of new effects within the 5d UED framework such as: ($`i`$) The KK spectra of the gauge and Higgs bosons as well as those for all of the left- and right-handed fermions can each be rescaled arbitrarily. This can increase the possible confusion between UED and SUSY scenarios at the LHC. ($`ii`$) New loop-induced parity violating effects are generated within previously parity conserving sectors of the model, e.g., in 5d QED/QCD. One signature of this is the generation of anapole moment-type couplings between fermions and gauge bosons. ($`iii`$) KK-parity, which is exact even at loop order in ordinary UED, becomes broken through mixings in the fermion KK mass matrices. This results in loop-induced mixing among gauge KK states and a finite lifetime for the Lightest KK Particle(LKKP) which usually considered as a stable dark matter candidate in the UED scenario. ($`iv`$) It is possible that the 5d LV operators may be the source of electroweak symmetry breaking.
The outline of this paper is as follows: In Section 2 we discuss the adaptation of the LV SM Extension operators to 5d under the assumption that LI is broken only in 5d. Here we also show how the 5d analogs of the 4d CPT violating (as well as LV) operators can be (almost) removed from the action by suitable field redefinitions. We will show that through these redefinitions these operators may induce spontaneous symmetry breaking by generating a negative mass squared for scalar fields. This leads to a potential correlation between the SM Higgs vev and the size of the extra dimension, $`R^1`$. In Section 3, we analyze how the remaining operators lead to modifications in the KK spectra of the SM gauge, scalar and fermion fields. We show that having a different KK spectrum for left- and right-handed fermions, which is possible now that 5d LI is broken, yields a potential source for parity violation in QED/QCD in 4d. In Section 4 we demonstrate that KK-parity violation results from the Yukawa sector of the theory that normally generates masses for the would be fermion zero-modes. This again arises from the field redefinitions employed earlier to remove the analogs of the CPT violation operators. KK-parity violation is shown to lead to a number of new effects such as the instability of the lightest KK-parity odd particle as well as general mixing among the KK-even and KK-odd excitations. A discussion and our conclusions can be found in Section 5. The Appendix contains a brief discussion of LV in 6d for scalar fields.
## 2 Lorentz Violating Operators
Given the field content of the SM, the SM Extension provides for us a relatively short list of the lowest-dimension gauge invariant 4d LV operators which may also be CPT violating. We can easily adapt this list to our purposes and restrict ourselves to those cases where only 5d LI is lost while 4d LI remains. This requirement turns out to be highly restrictive as many of the 4d LV operators do not have 5d analogs which can lead to loss of LI in only 5d. In principle, one can extend this list by including new operators generated by gravity modifications in 5-d as has been done in Ref.. However, since we are restricting ourselves to the UED framework with one additional TeV scale extra dimension wherein gravitational effects we will not discuss this possibility here.
Systematically going through the list 4d operators we find a number whose generalizations to 5d cannot satisfy our constraints. Consider for example the LV set of 4d operators
$$\mathrm{\Phi }^{}(\alpha _{\mu \nu }B^{\mu \nu }+\beta _{\mu \nu }T^aW_a^{\mu \nu })\mathrm{\Phi },$$
(1)
where $`\mathrm{\Phi }`$ is the Higgs scalar and $`B^{\mu \nu }`$ and $`W_a^{\mu \nu }`$ are the $`U(1)`$ and $`SU(2)`$ field strength tensors; $`\alpha _{\mu \nu }`$ and $`\beta _{\mu \nu }`$ are sets of numerical coefficients. Generalizing to 5d we immediately obtain
$$\mathrm{\Phi }^{}(\alpha _{MN}B^{MN}+\beta _{MN}T^aW_a^{MN})\mathrm{\Phi }.$$
(2)
We now ask what values of $`A,B`$ are allowed for the coefficients above without violating 4d LI: if $`A,B`$ both take on 4d indices then 4d LI will be broken. Similarly, if, e.g., $`A=\mu `$ and $`B=5`$ then 4d LI is again lost; the last possibility, i.e., $`A=B=5`$, yields zero due to the index antisymmetry. Thus the generalization of operators such as this in 4d to 5d does not yield anything interesting given the assumptions of our analysis. As another example of this, consider the 4d operator
$$\kappa _{\mu \nu }\overline{D}\sigma ^{\mu \nu }S\mathrm{\Phi }+h.c.,$$
(3)
where $`D(S)`$ is an $`SU(2)_L`$ doublet(singlet) fermion field and $`\kappa _{\mu \nu }`$ a set of numerical coefficients. In 5d this trivially generalizes to
$$\kappa _{AB}\overline{D}\sigma ^{AB}S\mathrm{\Phi }+h.c..$$
(4)
As in the previous example the various possible choices of $`A,B`$ yield either LV in 4d or are zero by the antisymmetry of the indices.
Let us now turn to the set of surviving operators. As an example, in 4d, suppressing flavor indices one has the following possible ‘kinetic’ LV term, e.g., for an $`SU(2)_L`$ singlet fermion field, $`S`$:
$$\frac{i}{2}(c_S)_{\mu \nu }\overline{S}\gamma ^\mu \overline{D}^\nu S,$$
(5)
where $`\overline{D}^\nu `$ is a covariant derivative acting in both directions and the $`c_S`$ are a set of dimensionless numerical coefficients; we expect these coefficients to be very small in 4d. Here we wish to generalize this operator to 5d, i.e.,
$$\frac{i}{2}(c_S)_{AB}\overline{S}\mathrm{\Gamma }^A\overline{D}^BS,$$
(6)
where $`\mathrm{\Gamma }^\mu =\gamma ^\mu `$, $`\mathrm{\Gamma }^5=i\gamma _5`$ and only keep terms which are LV in 5d but not in 4d. Clearly, given the experience of the examples above, we can only choose $`A=B=5`$ and taking $`k_S=(c_S)_{55}`$ this term becomes
$$\frac{k_S}{2}[\overline{S}\gamma _5D_5S(D_5\overline{S})\gamma _5S].$$
(7)
Perhaps, more naturally, in 5d we might anticipate that $`k_S=O(1)`$ since LI in 4d remains unbroken. Of course we may expect a similar term to be present for an $`SU(2)_L`$ doublet as in the singlet case described above, with $`k_Sk_D`$, but which need not be the same value. We will assume for simplicity that these 5d fermion terms are flavor-diagonal in what follows and denote their set of coefficients more generically by $`k_\mathrm{\Psi }`$. Going through and attempting to generalize the remaining set of SM Extension 4d operators we find that only very few of them satisfy our 5d requirements; in addition to the ‘kinetic’-like operator above for fermions we find the following possibilities:
$`{\displaystyle \frac{1}{4}}k_{\kappa \lambda \mu \nu }F^{\kappa \lambda }F^{\mu \nu }`$ $``$ $`{\displaystyle \frac{\lambda }{4}}\left(F_{\mu 5}F^{\mu 5}+F_{5\mu }F^{5\mu }\right)`$
$`k_{\mu \nu }^{}(D^\mu \mathrm{\Phi })^{}(D^\nu \mathrm{\Phi })`$ $``$ $`k_\mathrm{\Phi }(D_5\mathrm{\Phi })^{}(D_5\mathrm{\Phi })`$
$`a_\mu \overline{f}\gamma ^\mu f`$ $``$ $`i\alpha \overline{f}\gamma _5f`$
$`i(k_\varphi )^\mu \mathrm{\Phi }^{}D_\mu \mathrm{\Phi }+h.c.`$ $``$ $`ih_\mathrm{\Phi }\mathrm{\Phi }^{}D_5\mathrm{\Phi }+h.c.,`$ (8)
where the ‘mapping’ from 4d to 5d is shown explicitly. In the equation above, $`\mathrm{\Phi }`$ represents the Higgs doublet as before and, correspondingly, $`F`$ represents any of the SM field strength tensors. Likewise, $`f`$ is either an $`SU(2)_L`$ doublet, $`D`$, or singlet, $`S`$, fermion. Note that the first two operators above lead to modifications of gauge and Higgs kinetic terms as was the case for the fermionic operator discussed previously. Also note that in 4d the last two operators in the equation above are CPT violating; we note that in 5d their coefficients must be $`Z_2`$-odd in a manner similar to that of any 5d bulk fermion mass term. The parameters $`\lambda ,k_{\mathrm{\Phi },\mathrm{\Psi }}`$ are dimensionless quantities which we might expect to be of order unity; they must be highly suppressed quantities in the usual 4d SM Extension. On the otherhand the coefficients $`\alpha `$ and $`h_\mathrm{\Phi }`$ have the dimensions of mass and might most naturally be of order $`R^1`$. While it is possible that higher dimensional operators may also be present in addition to the ones considered above these are likely to be Planck suppressed and can be safely ignored in the analysis below.
The LV operators that we have found above are for a 5d scenario. It would be interesting to consider how this operator set would be modified by going to even higher dimensions, e.g., 6d. Here we could imagine that not only is LI of the type that we have been discussing violated in the higher dimensional action but also rotational invariance in the extra dimensions may be lost leading to very interesting new physics. Such possibilities will be considered briefly in the Appendix and in detail elsewhere.
It was noted in Ref that some of these CPT violating operators can be eliminated from the action in 4d by suitable field redefinitions. This remains especially true here in 5d (since we are only considering LV in the one extra dimension) but with interesting consequences since these field redefinitions will now depend on the co-ordinate of the extra fifth dimension. In a way, this field redefinition resums the effects of these operators in a non-perturbative way into the fields themselvess. Leaving these operators in the action, one would obtain similar effects order by order in the perturbation theory in the new LV parameters. The field redefinition simplifies the action and allows us to provide a non-perturbative treatment of these operators. As an example, consider the scalar part of the action including the two relevant LV terms above:
$$d^4x𝑑y\left[(D_A\mathrm{\Phi })^{}(D^A\mathrm{\Phi })V(\mathrm{\Phi }^{}\mathrm{\Phi })k_\mathrm{\Phi }(D_5\mathrm{\Phi })^{}(D_5\mathrm{\Phi })+ih_\mathrm{\Phi }(\mathrm{\Phi }^{}D_5\mathrm{\Phi }\mathrm{\Phi }D_5\mathrm{\Phi }^{})\right],$$
(9)
where we use $`y`$ as the co-ordinate for the extra dimension and $`V`$ is the usual scalar potential. If we now make a field redefinition
$$\mathrm{\Phi }e^{i\mathrm{\Sigma }_\mathrm{\Phi }y}\mathrm{\Phi },$$
(10)
where $`\mathrm{\Sigma }`$ is a parameter whose value we choose to be (recall $`h_\mathrm{\Phi }`$ is $`Z_2`$-odd)
$$\mathrm{\Sigma }_\mathrm{\Phi }=\frac{h_\mathrm{\Phi }}{1+k_\mathrm{\Phi }},$$
(11)
then the action becomes
$$d^4x𝑑y\left[(D_A\mathrm{\Phi })^{}(D^A\mathrm{\Phi })V(\mathrm{\Phi }^{}\mathrm{\Phi })+\mathrm{\Sigma }_\mathrm{\Phi }^2(1+k_\mathrm{\Phi })\mathrm{\Phi }^{}\mathrm{\Phi }k_\mathrm{\Phi }(D_5\mathrm{\Phi })^{}(D_5\mathrm{\Phi })\right],$$
(12)
thus eliminating the ‘CPT violating’ term but now introducing a new contribution to the scalar potential. Note that although the parameter $`\mathrm{\Sigma }_\mathrm{\Phi }`$ must be $`Z_2`$-odd to maintain the original symmetry only its square now appears in the action. Though the ‘CPT’ violating operator no longer appears its effects will remain important as we will see below.
Interestingly it is possible that this new quadratic term could produce a negative mass squared for the scalar Higgs field thus generating spontaneous symmetry breaking in the potential. Since we imagine that most naturally $`\mathrm{\Sigma }_\mathrm{\Phi }R^1`$ in magnitude this would tell us that the weak scale is linked to the size of the compactification scale up to order one corrections. Note that our field redfinition is consistent with the original $`Z_2`$ symmetry. Also note that additional kinetic term proportional to $`k_\varphi `$ can induce tachyonic states, i.e., instabilities and /or causality violations, unless the value of this parameter is restricted to be $`1`$. Such effects were observed in the 4d case when LV was present\[stab\].
A similar field redefinition trick can also be applied in the fermion sector to rid ourselves of the ‘CPT violating’ term. Let $`\mathrm{\Psi }`$ be any 5d fermion field; the relevant action is then
$$d^4x𝑑y\left[\frac{i}{2}\overline{\mathrm{\Psi }}\mathrm{\Gamma }^A\overline{D}_A\mathrm{\Psi }\frac{1}{2}k_\mathrm{\Psi }\overline{\mathrm{\Psi }}\gamma _5\overline{D}_5\mathrm{\Psi }i\alpha \overline{\mathrm{\Psi }}\gamma _5\mathrm{\Psi }\right],$$
(13)
where we have neglected any bulk mass terms as is standard in UED. Now we make the field redefinition (which maintains the original $`Z_2`$ symmetry)
$$\mathrm{\Psi }e^{i\mathrm{\Sigma }_\mathrm{\Psi }y}\mathrm{\Psi },$$
(14)
with
$$\mathrm{\Sigma }_\mathrm{\Psi }=\frac{\alpha }{1+k_\mathrm{\Psi }},$$
(15)
and the ‘CPT violating’ term is eliminated leaving the action
$$d^4x𝑑y\left[\frac{i}{2}\overline{\mathrm{\Psi }}\mathrm{\Gamma }^A\overline{D}_A\mathrm{\Psi }\frac{1}{2}k_\mathrm{\Psi }\overline{\mathrm{\Psi }}\gamma _5\overline{D}_5\mathrm{\Psi }\right],$$
(16)
this time with no additional terms. As in the case above we note that the coefficient $`\mathrm{\Sigma }_\mathrm{\Psi }`$ must be $`Z_2`$-odd. Thus out of the five possible LV structures in 5d we can eliminate two of them by field redefinitions; as we will see these redefinitions will return to haunt us later on. Note that the remaining LV terms are all modifications to kinetic terms and are all dimension-5, i.e., they are of the same dimension as are the usual SM-like terms in the 5d action.
## 3 Influence of LV Terms: KK Spectrum
The three remaining LV terms have a common feature: they are modifications of the 5d kinetic terms for fermions, Higgs bosons and gauge fields. They are analogous to (but not the same as) the addition of brane kinetic terms in the action and will produce similar effects as we will now see. We remind the reader that the LV contributions discussed below occur at the tree level while the somewhat analogous effects observed in the UED occur at loop order.
For simplicity let us first examine the case of the free scalar(Higgs) field; the action is then
$$d^4x𝑑y\left[(_A\mathrm{\Phi })^{}(^A\mathrm{\Phi })\mu ^2\mathrm{\Phi }^{}\mathrm{\Phi }k_\mathrm{\Phi }(_5\mathrm{\Phi })^{}(_5\mathrm{\Phi })\right],$$
(17)
where we have allowed a standard bulk mass term only for this discussion. Performing the Kaluza-Klein(KK) decomposition as usual
$$\mathrm{\Phi }=\underset{n}{}\varphi _n(x)\chi _n(y),$$
(18)
and imposing the orbifold boundary conditions one obtains the usual eigenfunctions $`\chi _n(y)\mathrm{cos}q_ny`$ for $`Z_2`$-even fields where $`q_n^2=m_n^2\mu ^2`$, with $`m_n`$ being the KK mass. In addition, these wavefunctions also have the standard normalization $`𝑑y\chi _n(y)\chi _m(y)=\delta _{nm}`$. However, the KK spectrum is now somewhat different than usual:
$$m_n^2=\mu ^2+\frac{n^2}{R^2}(1+k_\mathrm{\Phi }),$$
(19)
where $`R`$ is the compactification radius and $`n`$ is an integer. The effect of the LV term is to rescale the KK excitation spectrum by some arbitrary amount. (Recall the we expect the dimensionless quantity $`k_\mathrm{\Phi }`$ to be as large as order unity.) This is quite similar to the loop-induced radiative Higgs mass shift found in the case of UED induced by brane kinetic terms. In that model the size of the contribution was logarithmically dependent on the cutoff but was under control numerically; here the rescaling occurs at the tree-level and is completely arbitrary. In order to insure a tachyon-free spectrum, i.e., stability, it is clear that we must have $`k_\mathrm{\Phi }>1`$.
Next we turn to the gauge fields; when the corresponding gauge group is not spontaneously broken, e.g., for the case of gluons in $`SU(3)_c`$, the action is
$$d^4x𝑑y\left[\frac{1}{4}F_{AB}F^{AB}\frac{\lambda _c}{4}\left(F_{\mu 5}F^{\mu 5}+F_{5\mu }F^{5\mu }\right)\right],$$
(20)
where color indices have been suppressed. The KK decomposition is most conveniently performed in the physical $`g_5=0`$ gauge:
$$g_\mu =\underset{n}{}g_\mu ^{(n)}(x)f_n(y),$$
(21)
and produces the standard eigenfunctions $`\mathrm{cos}ny/R`$ for $`Z_2`$-even fields normalized as usual. In a manner similar to the scalar case above, the KK masses are, however, now given by
$$m_{g_n}^2=\frac{n^2}{R^2}(1+\lambda _c),$$
(22)
where $`\lambda _c`$ can be O(1). This spectrum shift is again similar to that induced by brane term radiative corrections in the UED model but here it can rescale the spectrum arbitrarily by a large amount. Since $`\lambda `$ and $`k_\mathrm{\Phi }`$ are completely unrelated, this rescaling of the KK spectra can be performed independently for these fields.
In the electroweak sector the KK decomposition is a bit more complex due to presence of symmetry breaking, the mixing among the neutral fields and the fact that the LV coefficients for the $`SU(2)_L`$ and $`U(1)_Y`$ gauge groups, $`\lambda _{W,B}`$, respectively, can be numerically different. The case of the $`W`$ KK tower is rather straightforward since the effect of symmetry breaking here is to generate a simple bulk mass term with the usual eigenfunctions; one obtains in standard notation
$$m_{W_n}^2=\frac{1}{4}g^2v^2+\frac{n^2}{R^2}(1+\lambda _W),$$
(23)
as we might have expected. Note that we can adjust the $`W`$ and gluon towers relative to each other in an arbitrary manner; in UED the ratio of these two, loop-induced shifts is completely fixed. For the neutral fields one obtains a level-by-level mixing between the KK excitations of the $`B`$ and $`W_3`$ fields as in the case of UED. The elements of the symmetric KK mixing matrix at the $`n^{th}`$ level are given by
$`M_{W_3W_3}^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}g^2v^2+{\displaystyle \frac{n^2}{R^2}}(1+\lambda _W)`$
$`M_{W_3B}^2=M_{BW_3}^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}gg^{}v^2`$
$`M_{BB}^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}g^2v^2+{\displaystyle \frac{n^2}{R^2}}(1+\lambda _B),`$ (24)
which is somewhat similar to the case of UED. The corresponding level-dependent ‘Weinberg-angle’ is then given by
$$\mathrm{tan}2\theta _n=\frac{2gg^{}v^2}{4(\lambda _B\lambda _W)\frac{n^2}{R^2}+(g^2g^2)v^2}.$$
(25)
Since $`\lambda _{B,W}`$ are arbitrary, in principle O(1) parameters, this KK mass spectrum for these two neutral fields can be substantially different than obtained in UED.
The case for fermions proceeds in the standard fashion from the above action. Let us ignore $`SU(2)_L\times U(1)_Y`$ symmetry breaking and zero-mode mass generation for the moment; we will return to this issue below. The arbitrary 5d fermion field $`\mathrm{\Psi }`$ is decomposed into left- and right-handed pieces in the standard manner: $`\mathrm{\Psi }=P_L\mathrm{\Psi }_L+P_R\mathrm{\Psi }_R`$ using the usual projection operators and then the KK decomposition is performed, i.e., $`\mathrm{\Psi }_{L,R}=_n\psi _{L,R}^{(n)}f_{L,R}^{(n)}(y)`$ and we arrive at the the coupled equations
$`(1+k_\mathrm{\Psi })_yf_L^{(n)}`$ $`=`$ $`m_nf_R^{(n)}`$
$`(1+k_\mathrm{\Psi })_yf_R^{(n)}`$ $`=`$ $`m_nf_L^{(n)},`$ (26)
so that the fermion masses are given by
$$m_{\mathrm{\Psi }_n}^2=(1+k_\mathrm{\Psi })^2\frac{n^2}{R^2}.$$
(27)
For doublet(singlet) fields we choose the left(right)-handed fermions to be $`Z_2`$-even to obtain the conventional SM structure. Note that the mass-squared of the fermion fields are quadratic in the LV correction whereas bosons experience a linear correction. In either case we again see that we can rescale the mass spectrum as we’d like since we can choose the LV coefficients arbitrarily. In particular, there is no reason, e.g., for the left- and right-handed SM fermions to have KK towers that are in any way degenerate which can lead to new physics signatures as will be discussed below. Although KK-parity is still maintained at this point, clearly if one can rescale all of the masses of the KK excitations of the SM fields by arbitrary amounts it is no longer clear which state will be the lightest one in the spectrum. The identity of the LKKP dark matter candidate now depends on the values of the LV coefficients. We note that since we can rescale the fermion and boson spectrum as we’d like the possibility of confusion between the UED and SUSY scenarios at the LHC is now significantly increased.
The fact that the KK excitations of left-handed doublet and right-handed singlet fermions now have different tree-level masses directly leads to new phenomena. As a demonstration of this, consider for simplicity the toy model of 5d QED accompanied by LV. The KK towers of the of the left-handed and right-handed electron now having different masses will produce a signal for parity violation within a conventionally parity conserving scenario as we will now demonstrate. If one considers the coupling between the (zero-mode) photon and left- and right-handed electrons one finds that at loop level a parity violating coupling will be generated, i.e., the anapole moment, which corresponds to a tensor structure
$$<f|J_\mu ^{em}|f>_{anapole}=ieQ_f\overline{f}[q^2\gamma _\mu \gamma qq_\mu ]\gamma _5F_3(q^2)f,$$
(28)
with $`F_3(q^2)`$ being the anapole moment form factor. $`F_3(0)=a`$ is then just the standard anapole moment which has dimensions $`R^2`$. The existence of this coupling is directly related to the fact that the masses of the KK states inside the loop are different for the left- and right-handed towers; in obvious notation and summing over KK levels we find that
$$a\frac{\alpha }{\pi }\frac{\pi ^2R^2}{48}_0^1𝑑x_0^{1x}𝑑y\left[\frac{4(1x)(1y)+5xy}{1+\lambda _\gamma +[(1+k_R)^21\lambda _\gamma ](x+y)}(RL)\right].$$
(29)
Assuming that $`R^1=`$500 GeV and $`\lambda _\gamma =0`$ for purposes of demonstration we obtain the numerical results for $`a`$ shown in Fig. 1; here we see that $`|a|\alpha `$ TeV<sup>-2</sup>, which to set the scale, is comparable in magnitude to the conventional SM contribution induced by the parity-violating weak interaction. Clearly LV in 5d can lead to parity violating signatures in 4d in the absence of weak interactions.
Within the 5d UED scenario the analysis above also leads directly to a QCD color-anapole moment which also violates parity in 4d.
As we have just seen the introduction of 5d LV violating operators with O(1) coefficients allows us to modify the overall scales of the various gauge, scalar and left- and right-handed fermion KK spectra in UED in an independent fashion. Thus when such operators are present it is no longer clear that, e.g., a neutral field will be the lightest state which is odd under KK-parity and we may lose our natural dark matter candidate. The situation is actually more severe than this as we will now see.
## 4 Influence of LV Terms: KK-Parity Violation
The remaining term in the action that we have yet to examine is the Yukawa coupling between the fermion doublet and singlet fields and the Higgs boson that generates non-zero masses for the (would-be) zero-mode SM fermions. We can write this generically, dropping all generation labels, as:
$$d^4x𝑑y\lambda _5\overline{D}S\mathrm{\Phi }+h.c..$$
(30)
After rescaling by the field redefinitions employed above to rid ourselves of the unwanted ‘CPT violating’ LV terms this action becomes
$$d^4x𝑑y\lambda _5e^{i(\mathrm{\Sigma }_\mathrm{\Phi }+\mathrm{\Sigma }_S\mathrm{\Sigma }_D)y}\overline{D}S\mathrm{\Phi }+h.c.,$$
(31)
so that a position-dependent phase has crept into the Yukawa part of the action. It is important to recall that the quantities $`\mathrm{\Sigma }_i`$ are $`Z_2`$-odd, i.e., they flip their sign at the origin. To probe the influence of this term let us first extract out the all zero-mode piece and perform the $`y`$-integration. Recall that zero-mode wave functions are flat $`=1/\sqrt{2\pi R}`$; we obtain the 4d integrand
$$\frac{\lambda _5}{\sqrt{2\pi R}}\frac{v+H}{\sqrt{2}}e^{i\sigma \pi R/2}\frac{\mathrm{sin}(\sigma \pi R/2)}{\sigma \pi R/2},$$
(32)
where $`\sigma =\mathrm{\Sigma }_\mathrm{\Phi }+\mathrm{\Sigma }_S\mathrm{\Sigma }_D`$ and $`H`$ is the usual SM Higgs field. The SM 4d Yukawa coupling can then be identified as
$$\lambda _4=\frac{\lambda _5}{\sqrt{2\pi R}}e^{i\sigma \pi R/2}\frac{\mathrm{sin}(\sigma \pi R/2)}{\sigma \pi R/2}.$$
(33)
Apart from the overall phase factor the last term can substantially rescale the size of the Yukawa coupling depending on the value of $`\sigma R`$ and may lead to some interesting phenomenology.
Something even more interesting results when we do not project into the all zero-mode state. Due to the additional $`y`$-dependent phase these Yukawa terms can violate KK-parity causing, e.g., a destabilization of the LKKP. Recall that in the usual UED model KK-parity is preserved to all orders in perturbation theory. To see this effect it is useful to examine the mixing between the would-be zero mode fermion and the $`Z_2`$-even members of the KK tower; this corresponds to the off-diagonal sub-matrix linking, e.g., the zero-mode doublet field with a KK singlet state. This calculation is straightforward and, in terms of the 4d Yukawa coupling $`\lambda _4`$ is given by
$$\frac{\lambda _4v}{2}\left[e^{in\pi /2}\frac{\sigma R}{\sigma R+n}\frac{\mathrm{sin}((\sigma R+n)\pi /2)}{\mathrm{sin}(\sigma \pi R/2)}\right]+(nn),$$
(34)
which corresponds to the $`0n`$ element of the KK mass sub-matrix, $`M_{0n}`$, and is seen to be proportional to the SM zero mode mass, $`M_{0n}=\delta _{0n}m_f`$. (It is important to note that here the symbol $`\delta _{0n}`$ does not denote the Kronecker delta.) Clearly, such terms can only be significant if $`\sigma R`$ is O(1) but this might be expected. Furthermore, one finds that all of the sub-matrix elements of this type, $`M_{nm}=m_f\delta _{nm}`$, are in general found to be non-zero with a mass scale set by the conventional SM fermion mass, i.e., with $`\delta _{nm}`$ values generally of order unity and controlled by the values of $`n,m`$ and $`\sigma R`$. This is unlike the case of UED where the mixing between the $`D`$ and $`S`$ fermion fields takes place level by level; here there is also a potentially significant mixing between the various KK levels. However, light fermions, such as the electron, experience little direct KK-parity violation through such mixing whereas for heavy fields, like the top quark, this violation can be quite significant for $`R^11`$ TeV or less. The removal of the ‘CPT-violating’ LV terms in the original action via field redefinitions has thus resulted in the breakdown of KK-parity conservation. We note that KK-parity violation at some level might also occur if UED is extended higher dimensions to include gravitational effects.
This violation of KK-parity is quantifiable at the tree level by estimating the lifetime of the LKKP. To get an order of magnitude estimate we perform the calculation in the mass insertion approximation and assume that as usual the first KK photon excitation with mass $`M`$ is the LKKP. The process proceeds via $`\gamma _1\overline{f}_{L0}f_{L1}+h.c.\overline{f}_{L0}f_{L0}+(LR)`$, where the second step arises from mixing. We obtain
$$\mathrm{\Gamma }=N_cQ_f^2(2Re\delta _{01})^2\frac{\alpha M}{6}(14m_f^2/M^2)^{1/2}\left[(g_L^2+g_R^2)\left(1\frac{m_f^2}{M^2}\right)+6g_Lg_R\frac{m_f^2}{M^2}\right],$$
(35)
where $`m_f`$ is the would-be zero mode mass, $`M=M_{\gamma _1}`$,
$$g_{L,R}=\frac{m_f^2}{(m_{L,R}^2m_f^2)},$$
(36)
with $`m_{L(R)}`$ being the mass of the first KK excitation of the doublet (singlet) field $`f_{L(R)}`$ and $`\delta _{01}`$ is the dimensionless mixing parameter defined above.
Here we see again an example of induced parity violation in that the two couplings are equal, $`g_L=g_R`$, only when the fermion KK excitation masses are the same. Note that as expected this decay is very highly suppressed for light fermions, i.e., decays to heavy fermions such as top quarks, will be by far dominant. To get an idea of the size of this suppression, we take $`m_L=m_R=M`$ and $`2Re(\delta _{01})=1`$ so that
$$\mathrm{\Gamma }=\frac{N_c}{3}Q_f^2\alpha MF(x),$$
(37)
with $`x=m_f/M`$; the function $`F(x)`$ is shown in Fig. 2. As we expected, except for the closure of phase space $`F(x)`$ is larger the closer $`x`$ is to 1/2; decays to first generation fields is thus seen to be highly suppressed. Although the expression above might correspond to a very narrow width by usual collider standards, for any reasonable range of parameters the lifetime of the LKKP is quite short in comparison to the age of the universe. Clearly, if some other particle is actually the LKKP, the analogous calculation can be performed obtaining qualitatively similar results.
Another way to observe the violation of KK-parity is through loop-induced mixing among different gauge boson KK levels. This mixing is induced by the insertion of off-diagonal fermion mass matrix elements into vacuum polarization graphs connecting gauge fields with different KK number. In the 5d QED example this corresponds to a process $`\gamma _n\overline{f}_nf_0+h.c.\overline{f}_mf_0+h.c.\gamma _m`$ where the intermediate step occurs due to Yukawa induced fermion mixing. Using the notation above, mass mixing arising from this process in the photon tower mass matrix induced by a single fermion flavor is given by
$$\delta M_{mn}^2N_cQ_f^2\frac{\alpha }{\pi }2Re(\delta _{mn})m_f^2G(\frac{m_f^2}{m_{S_n}^2},\frac{m_f^2}{m_{D_m}^2})+(nm),$$
(38)
with $`m_{D_m,S_n}`$ being the masses of the KK fermions in the loop and $`G`$ is an order one loop function. Here we again see that the dominant contribution arises from the most massive SM fermion sector as we might have expected. Clearly with $`\delta _{nm}`$’s of order unity a summation over all possible fermions in the loop can lead to small yet significant mixing in the gauge boson mass matrix. This result easily generalizes to the cases of the $`W,Z`$ and gluon KK towers where gauge KK mass eigenstates will now no longer have a definite KK-parity. A similar mixing will occur among Higgs and Goldstone KK levels.
One of the other effects of KK-parity conservation in the UED model is the inability to singly produce states which are KK-parity odd at colliders, e.g., the lightest KK gauge boson excitations. The violation of KK-parity induced by Yukawa interactions leads to modifications of this conventional result though the corresponding cross sections are not necessarily large. This can be seen by the fact that the widths of the KK-odd gauge bosons into zero modes of the first two generations is quite small.
In this section we have seen that the elimination of the 5d analogs of the 4d ‘CPT violating’ operators by field redefinitions induces potentially large violations of KK-parity. We observed that the size of this violation an any given SM fermion sector is correlated with the known size of the would-be zero mode masses. As a result UED loses its dark matter candidate.
## 5 Summary and Conclusions
In this paper we have initiated a study of the influence of explicit Lorentz violation within the context of the 5d SM where all fields are in the bulk, i.e., the Universal Extra Dimensions scenario. To perform this analysis we extended the ‘conventional’ 4d model of Colladay and Kostelecky to 5d and searched for a subset of operators that can leave 4d Lorentz invariance untouched while breaking it in 5d. Two of these operators, the 5d analogs of those that violate CPT in 4d, can be (almost) removed from the action through a set of field redefinitions for fermions and scalars. One obvious result of this field redefinition is to induce a negative mass square term in the Higgs potential which may be the source of electroweak symmetry breaking. In addition, the natural scale of the induced vev would be $`R^1`$ thus linking the scale of electroweak symmetry breaking with the size of the extra dimension. The remaining LV operators lead to alterations of the various gauge, Higgs and fermionic kinetic terms and independently rescale their associated KK spectra which can increase the possible confusion of UED and SUSY at the LHC. Since, e.g., the masses of KK excitations of the left- and right-handed SM fermions need no longer be equal this induces, at loop order, parity-violating effects in previously parity-conserving parts of the SM, i.e., QED and QCD. Furthermore, we have shown that the the field redefinitions used to eliminate the 5d analogs of the 4d CPT violating terms make an important change in the nature of the Yukawa couplings responsible for generating the would-be zero-mode fermion masses. Due to an additional fifth co-ordinate-dependent phase, fermion mass terms are generated that produce mixing among all of the various KK levels thus violating KK-parity. This leads to a destabilization of the lightest KK-odd particle which is the usual dark matter candidate in UED. In addition these terms were shown to induce mixing between the various gauge KK levels at one-loop.
As we have seen the presence of LV terms in the 5d UED scenario can lead to substantial modifications from the conventional expectations. Hopefully signals for extra dimensions will be found at future colliders.
Acknowledgments
The author would like to thank J.Hewett and B. Lillie for discussions related to this work.
## Appendix: LV in 6d
It is interesting to consider what happens when LV is extended to 6d and compactified on an orthogonal torus with radii $`R_{5,6}`$. Although a detailed study lies outside the scope of the current work we would like to give some flavor here by considering for simplicity the LV 6d scalar action. From the analysis above, this is given by
$$d^4x𝑑x_5𝑑x_6\left[(D_A\mathrm{\Phi })^{}(D^A\mathrm{\Phi })V(\mathrm{\Phi }^{}\mathrm{\Phi })k_{ij}(D_i\mathrm{\Phi })^{}(D_j\mathrm{\Phi })+ih_i(\mathrm{\Phi }^{}D_i\mathrm{\Phi }\mathrm{\Phi }D_i\mathrm{\Phi }^{})\right],$$
(39)
where we take $`k_{i,j}`$, $`i,j=5,6`$, to be real and symmetric; summation over these indices when repeated is implied. Since we will be concentrating for simplicity on the pure scalar sector in our discussion below, we have ignored the possibility of new LV interaction terms that may be present in 6d which are absent in 5d. As in 5d, the ‘CPT-violating’ terms can be eliminated by a field redefinition:
$$\mathrm{\Phi }e^{i\mathrm{\Sigma }_ix_i}\mathrm{\Phi },$$
(40)
where
$$\mathrm{\Sigma }_5=\frac{h_5(1+k_{66})+h_6k_{56}}{(1+k_{55})(1+k_{66})k_{56}^2},$$
(41)
and $`\mathrm{\Sigma }_6`$ can be obtained by interchanging 5 and 6 in the expression above. As in 5d this field redefinition adds a new, likely negative term to the scalar potential:
$$\frac{h_5^2(1+k_{66})+h_6^2(1+k_{55})+2h_5h_6k_{56}}{(1+k_{55})(1+k_{66})k_{56}^2}\mathrm{\Phi }^{}\mathrm{\Phi }.$$
(42)
So far this is a rather straightforward extension of 5d; something new happens when we perform the usual KK decomposition
$$\mathrm{\Phi }(x_\mu ,x_i)=\underset{n,m}{}\varphi _{n,m}(x_\mu )\chi _{n,m}(x_i).$$
(43)
Through the usual manipulations we are led to the equation of motion for $`\chi `$ which we can write for free scalars as
$$_i\left[h^{ij}_j\chi \right]m_{n,m}^2\chi =0,$$
(44)
where the symmetric object $`h_{ij}`$ acts as a flat, constant ‘metric’ in the $`x_5x_6`$ space with elements $`h_{55}=1+k_{55}`$, $`h_{66}=1+k_{66}`$ and $`h_{56}=k_{56}`$. These satisfy $`h_{il}h^{lj}=\delta _i^j`$ and thus $`h^{ij}`$ are the elements of the inverse matrix $`h^1`$. In the $`x_5x_6`$ co-ordinate basis this equation is not generally separable; however, the metric can be diagonalized through a suitable $`x_5x_6`$ rotation to the basis $`x_\pm `$:
$`x_+`$ $`=`$ $`x_5\mathrm{cos}\theta +x_6\mathrm{sin}\theta `$
$`x_{}`$ $`=`$ $`x_6\mathrm{cos}\theta x_5\mathrm{sin}\theta ,`$ (45)
with angle $`\theta `$ given by
$$\mathrm{tan}2\theta =\frac{2k_{56}}{k_{55}k_{66}},$$
(46)
so that the now separable equation of motion for $`\chi `$ becomes
$$\lambda _+^1_+^2\chi +\lambda _{}^1_{}^2\chi +m_{n,m}^2\chi =0,$$
(47)
with $`\lambda _\pm `$ given by
$$\lambda _\pm =1+\frac{k_{55}+k_{66}}{2}\pm \frac{1}{2}\left[(k_{55}k_{66})^2+4k_{56}^2\right]^{1/2}.$$
(48)
Note that although our metric is constant, rotations no longer commute with it. The fact that there is a ‘preferred’ frame where the ‘metric’ is diagonal is the result of LV here manifest as the loss of $`x_5x_6`$ rotational invariance. We can now express $`\chi `$ as $`\chi _{n,m}=f_n(x_+)g_m(x_{})`$ in this preferred basis.
Although we have switched to the co-ordinates $`x_\pm `$, the boundary conditions will most likely be expressed in the $`x_{5,6}`$ basis. Here, for example, we consider the most simple case where we have invariance under the typical periodic conditions: $`x_{5,6}x_{5,6}+2\pi R_{5,6}`$, so that one can write $`\chi =\mathrm{exp}in_ix_i/R_i=\mathrm{exp}i[a_+x_++a_{}x_{}]`$ and thus
$$m_{n_5,n_6}^2=\frac{a_+^2}{\lambda _+}+\frac{a_{}^2}{\lambda _{}},$$
(49)
where
$`a_+`$ $`=`$ $`{\displaystyle \frac{n_6}{R_6}}\mathrm{cos}\theta {\displaystyle \frac{n_5}{R_5}}\mathrm{sin}\theta `$
$`a_{}`$ $`=`$ $`{\displaystyle \frac{n_5}{R_5}}\mathrm{cos}\theta +{\displaystyle \frac{n_6}{R_6}}\mathrm{sin}\theta .`$ (50)
Note that in this simple case, the KK mass eigenvalue equation could have been obtained without making the co-ordinate transformation above since the eigenfunctions are simple exponentials. Straightforward algebra yields
$$m_{n_5,n_6}^2=\left[(1+k_{55})(1+k_{66})+k_{56}^2\right]^1\left((1+k_{66})\frac{n_5^2}{R_5^2}+(1+k_{55})\frac{n_6^2}{R_6^2}2k_{56}\frac{n_5n_6}{R_5R_6}\right).$$
(51)
This eigenvalue equation for the KK masses is remarkably similar to that obtained by Dienes who consider tori with shift angles and shape moduli in 6d. Instead of the simple KK spectrum rescaling that we observed for LV in 5d, in 6d the KK spectrum is significantly skewed and distorted compared to conventional expectations. The shift angle of Dienes in our case arises from LV and the existence of the preferred frame.
A more detailed discussion of LV in 6d will be given elsewhere.
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# Mirror symmetry for Del Pezzo surfaces: Vanishing cycles and coherent sheaves
## 1. Introduction
The phenomenon of mirror symmetry has been studied extensively in the case of Calabi-Yau manifolds (where it corresponds to a duality between $`N=2`$ superconformal sigma models), but also manifests itself in more general situations. For example, a sigma model whose target space is a Fano variety is expected to admit a mirror, not necessarily among sigma models, but in the more general context of Landau-Ginzburg models.
For us, a Landau-Ginzburg model is simply a pair $`(M,W)`$, where $`M`$ is a non-compact manifold (carrying a symplectic structure and/or a complex structure), and $`W`$ is a complex-valued function on $`M`$ called superpotential. The general philosophy is that, when a Landau-Ginzburg model $`(M,W)`$ is mirror to a Fano variety $`X`$, the complex (resp. symplectic) geometry of $`X`$ corresponds to the symplectic (resp. complex) geometry of the critical points of $`W`$.
We place ourselves in the context of homological mirror symmetry, where mirror symmetry is interpreted as an equivalence between certain triangulated categories naturally associated to a mirror pair . In our case, B-branes on a Fano variety are described by its derived category of coherent sheaves, and under mirror symmetry they correspond to the A-branes of a mirror Landau-Ginzburg model. These A-branes are described by a suitable analogue of the Fukaya category for a symplectic fibration, namely the derived category of Lagrangian vanishing cycles. A rigorous definition of this category has been proposed by Seidel in the case where the critical points of the superpotential are isolated and non-degenerate, following ideas of Kontsevich and Hori, Iqbal, Vafa .
Therefore, for a Fano variety $`X`$ and a mirror Landau-Ginzburg model $`W:M`$, the homological mirror symmetry conjecture can be formulated as follows:
###### Conjecture 1.1.
The derived category of Lagrangian vanishing cycles $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W))`$ is equivalent to the derived category of coherent sheaves $`𝐃^b(\mathrm{coh}(X))`$.
###### Remark 1.2.
Homological mirror symmetry also predicts another equivalence of derived categories. Namely, viewing now $`X`$ as a symplectic manifold and $`M`$ as a complex manifold, the derived category of B-branes of the Landau-Ginzburg model $`W:M`$, which was defined algebraically in following ideas of Kontsevich, should be equivalent to the derived Fukaya category of $`X`$. This aspect of mirror symmetry will be addressed in a further paper; for now, we focus exclusively on Conjecture 1.1.
One of the first examples for which Conjecture 1.1 has been verified is that of $`^2`$ and its mirror Landau-Ginzburg model which is the elliptic fibration with three singular fibers determined by the superpotential $`W_0=x+y+1/xy`$ on $`(^{})^2`$ (or rather a fiberwise compactification of this fibration), see . Other examples of surfaces for which the derived category of coherent sheaves has been shown to be equivalent to the derived category of Lagrangian vanishing cycles of a mirror Landau-Ginzburg model include weighted projective planes, Hirzebruch surfaces , and toric blow-ups of $`^2`$ . For all these examples, the toric structure plays a crucial role in determining the geometry of the mirror Landau-Ginzburg model.
Our goal in this paper is to consider the case of a Del Pezzo surface $`X_K`$ obtained by blowing up $`^2`$ at a set $`K`$ of $`k8`$ points (this is never toric as soon as $`k4`$). Our proposal is that a mirror of $`X_K`$ can be constructed in the following manner. Observe that the elliptic fibration with three singular fibers determined by the superpotential $`W_0=x+y+1/xy`$ on $`(^{})^2`$ (i.e., the mirror of $`^2`$) admits a natural compactification to an elliptic fibration $`\overline{W_0}:\overline{M}^1`$ in which the fiber above infinity consists of nine rational components (see §3.1 for details). Consider a deformation of $`\overline{W_0}`$ to another elliptic fibration $`\overline{W_k}:\overline{M}^1`$, such that $`k`$ of the $`9`$ critical points in the fiber $`\overline{W_0}{}_{}{}^{1}(\mathrm{})`$ are displaced towards finite values of the superpotential. Let
$$M_k=\overline{M}\overline{W_k}{}_{}{}^{1}(\mathrm{}),$$
and denote by $`W_k:M_k`$ the restriction of $`\overline{W_k}`$ to $`M_k`$. In the generic case, $`W_k`$ is an elliptic fibration with $`k+3`$ nodal fibers, while $`\overline{W_k}{}_{}{}^{1}(\mathrm{})`$ is a singular fiber with $`9k`$ rational components. Although we will focus on the Del Pezzo case, this construction also provides a mirror in some borderline situations. For example, it can be applied without modification to the case where $`^2`$ is blown up at $`k=9`$ points which lie at the intersection of two elliptic curves (the fiber $`\overline{W_k}{}_{}{}^{1}(\mathrm{})`$ is then a smooth elliptic curve).
There are two aspects to the geometry of $`M_k`$. Viewing $`M_k`$ as a complex manifold (a Zariski open subset of a rational elliptic surface), its complex structure is closely related to the set of critical values of $`W_k`$, which has to be chosen in accordance with a given symplectic structure on $`X_K`$. A generic choice of the symplectic structure on $`X_K`$ (for which there are no homologically nontrivial Lagrangian submanifolds) determines a complex structure on $`M_k`$ for which the $`k+3`$ critical values of $`W_k`$ are all distinct (leading to a very simple category of B-branes). In the opposite situation, which we will not consider here, if we equip $`X_K`$ with a symplectic form for which there are homologically nontrivial Lagrangian submanifolds, then some of the critical values of $`W_k`$ become equal, and the topology of the singular fibers may become more complicated.
The symplectic geometry of $`M_k`$ is more important to us. Since $`H^2(M_k,)^{k+2}`$, the symplectic form $`\omega `$ on $`M_k`$, or rather its complexified variant $`B+i\omega `$, depends on $`k+2`$ moduli parameters. As we will see in §4, these parameters completely determine the derived category of Lagrangian vanishing cycles of $`W_k`$; the actual positions of the critical values are of no importance, as long as the critical points of $`W_k`$ remain isolated and non-degenerate (see Lemma 3.2). This means that we shall not concern ourselves with the complex structure on $`M_k`$; in fact, a compatible almost-complex structure is sufficient for our purposes, which makes the problem of deforming the elliptic fibration $`\overline{W_0}`$ in the prescribed manner a non-issue.
To summarize, we have:
###### Construction 1.3.
Given a Del Pezzo surface $`X_K`$ obtained by blowing up $`^2`$ at $`k`$ points, the mirror Landau-Ginzburg model is an elliptic fibration $`W_k:M_k`$ with $`k+3`$ nodal singular fibers, which has the following properties:
$`(i)`$ the fibration $`W_k`$ compactifies to an elliptic fibration $`\overline{W_k}`$ over $`^1`$ in which the fiber above infinity consists of $`9k`$ rational components;
$`(ii)`$ the compactified fibration $`\overline{W_k}`$ can be obtained as a deformation of the elliptic fibration $`\overline{W_0}:\overline{M}^1`$ which compactifies the mirror to $`^2`$.
Moreover, the manifold $`M_k`$ is equipped with a symplectic form $`\omega `$ and a B-field $`B`$, whose cohomology classes are determined by the set of points $`K`$ in an explicit manner as discussed in §5.
Our main result is the following:
###### Theorem 1.4.
Given any Del Pezzo surface $`X_K`$ obtained by blowing up $`^2`$ at $`k`$ points, there exists a complexified symplectic form $`B+i\omega `$ on $`M_k`$ for which $`𝐃^b(\mathrm{coh}(X_K))𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$.
The mirror map, i.e. the relation between the cohomology class $`[B+i\omega ]H^2(M_k,)`$ and the positions of the blown up points in $`^2`$, can be described explicitly (see Proposition 5.1).
On the other hand, not every choice of $`[B+i\omega ]H^2(M_k,)`$ yields a category equivalent to the derived category of coherent sheaves on a Del Pezzo surface. There are two reasons for this. First, certain specific choices of $`[B+i\omega ]`$ correspond to deformations of the complex structure of $`X_K`$ for which the surface contains a $`2`$-curve, which causes the anticanonical class to no longer be ample. There are many ways in which this can occur, but perhaps the simplest one corresponds to the case where a same point is blown up twice, i.e. we first blow up $`^2`$ at $`k1`$ generic points and then blow up a point on one of the exceptional curves. We then say that $`X_K`$ is obtained from $`^2`$ by blowing up $`k`$ points, two of which are infinitely close, and call this a “simple degeneration” of a Del Pezzo surface. In this case again we have:
###### Theorem 1.5.
If $`X_K`$ is a blowup of $`^2`$ at $`k`$ points, two of which are infinitely close, and a simple degeneration of a Del Pezzo surface, then there exists a complexified symplectic form $`B+i\omega `$ on $`M_k`$ for which $`𝐃^b(\mathrm{coh}(X_K))𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$.
More importantly, deformations of the symplectic structure on $`M_k`$ need not always correspond to deformations of the complex structure on $`X_K`$ (observe that $`H^2(M_k,)`$ is larger than $`H^1(X_K,TX_K)`$). The additional deformation parameters on the mirror side can however be interpreted in terms of noncommutative deformations of the Del Pezzo surface $`X_K`$ (i.e., deformations of the derived category $`𝐃^b(\mathrm{coh}(X_K))`$). In this context we have the following theorem, which generalizes the result obtained in for the case of $`^2`$:
###### Theorem 1.6.
Given any noncommutative deformation of the Del Pezzo surface $`X_K`$, there exists a complexified symplectic form $`B+i\omega `$ on $`M_k`$ for which the deformed derived category $`𝐃^b(\mathrm{coh}(X_{K,\mu }))`$ is equivalent to $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$. Conversely, for a generic choice of $`[B+i\omega ]H^2(M_k,)`$, the derived category of Lagrangian vanishing cycles $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$ is equivalent to the derived category of coherent sheaves of a noncommutative deformation of a Del Pezzo surface.
The mirror map is again explicit, i.e. the parameters which determine the noncommutative Del Pezzo surface can be read off in a simple manner from the cohomology class $`[B+i\omega ]`$.
###### Remark 1.7.
The key point in the determination of the mirror map is that the parameters which determine the composition tensors in $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$ can be expressed explicitly in terms of the cohomology class $`[B+i\omega ]`$ (see §4.3). A remarkable feature of these formulas is that they can be interpreted in terms of theta functions on a certain elliptic curve (see §4.5). As a consequence, our description of the mirror map also involves theta functions (see §5).
The rest of the paper is organized as follows. In §2 we describe the bounded derived categories of coherent sheaves on Del Pezzo surfaces, their simple degenerations, and their noncommutative deformations. In §3 we describe the topology of the elliptic fibration $`M_k`$ and its vanishing cycles. In §4.1 we recall Seidel’s definition of the derived category of Lagrangian vanishing cycles of a symplectic fibration, and in the rest of §4 we determine $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$. Finally in §5 we compare the two viewpoints, describe the mirror map, and prove the main theorems.
Acknowledgements: We are thankful to A. Kapustin, T. Pantev, P. Seidel for many helpful discussions.
## 2. Derived categories of coherent sheaves on blowups of $`^2`$
The purpose of this section is to give a description of the bounded derived categories of coherent sheaves on Del Pezzo surfaces, their simple degenerations, and their noncommutative deformations. We always work over the field of complex numbers $`.`$
### 2.1. Del Pezzo surfaces and blowups of the projective plane at distinct points
###### Definition 2.1.
A smooth projective surface $`S`$ is called a Del Pezzo surface if the anticanonical sheaf $`𝒪_S(K_S)`$ is ample (i.e., a Del Pezzo surface is a Fano variety of dimension 2).
The Kodaira vanishing theorem and Serre duality give us immediately that for any Del Pezzo surface
$`H^1(S,𝒪(mK_S))=0\text{for all }m,`$
$`H^2(S,𝒪(mK_S))=0\text{for all }m0,`$
$`H^2(S,𝒪(mK_S))=H^0(S,𝒪((m+1)K_S))\text{for all }m.`$
In particular, we obtain that $`H^1(S,𝒪_S)H^2(S,𝒪_S)=0,`$ and $`H^0(S,𝒪(mK_S))=0`$ for all $`m>0.`$ By the Castelnuovo-Enriques criterion any Del Pezzo surface is rational.
Let $`S`$ be a Del Pezzo surface. The integer $`K_S^2`$ is called the degree of $`S`$ and will be denoted by $`d`$. The Noether formula gives a relation between the degree and the rank of the Picard group of a Del Pezzo surface: $`d=K_S^2=10\mathrm{rk}\mathrm{Pic}S9.`$
We can also introduce another integer number which is called the index of $`S`$. This is the maximal $`r>0`$ such that $`𝒪(K_S)=𝒪(rH)`$ for some divisor $`H`$. The inequality $`d9`$ implies that $`r3.`$
Now recall the classification of Del Pezzo surfaces.
If $`r=3,`$ then $`S^2`$ is the projective plane and $`d=9.`$ If $`r=2,`$ then $`S^1\times ^1`$ is the quadric and $`d=8.`$ The other Del Pezzo surfaces are not minimal and can be obtained by blowing up the projective plane $`^2.`$ More precisely, if $`S`$ is a Del Pezzo surface of index $`r=1,`$ then it has degree $`1d8`$ and $`S`$ is a blowup of the projective plane $`^2`$ at $`k=9d`$ distinct points. The ampleness of the anticanonical class requires that in this set no three points lie on a line, and no six points lie on a conic; moreover, if $`k=8`$ the eight points are not allowed to lie on an irreducible cubic which has a double point at one of these points. Conversely, any surface which is a blowup of the projective plane at a set of $`k8`$ different points satisfying these constraints is a Del Pezzo surface of degree $`d=9k.`$ All these facts are well-known and can be found in any textbook on surfaces (see e.g. ).
Denote by $`𝐃^b(\mathrm{coh}(S))`$ the bounded derived category of coherent sheaves on $`S`$. It is known that the bounded derived category of coherent sheaves on any Del Pezzo surface has a full exceptional collection, which makes it possible to establish an equivalence between the category $`𝐃^b(\mathrm{coh}(S))`$ and the bounded derived category of finitely generated modules over the algebra of the exceptional collection (, see also ). This is a particular case of a more general statement about derived categories of blowups.
First, recall the notion of exceptional collection.
###### Definition 2.2.
An object $`E`$ of a $``$-linear triangulated category $`𝒟`$ is said to be exceptional if $`\mathrm{Hom}(E,E[k])=0`$ for all $`k0,`$ and $`\mathrm{Hom}(E,E)=`$. An ordered set of exceptional objects $`\sigma =(E_0,\mathrm{}E_n)`$ is called an exceptional collection if $`\mathrm{Hom}(E_j,E_i[k])=0`$ for $`j>i`$ and all $`k.`$ The exceptional collection $`\sigma `$ is said to be strong if it satisfies the additional condition $`\mathrm{Hom}(E_j,E_i[k])=0`$ for all $`i,j`$ and for $`k0.`$
###### Definition 2.3.
An exceptional collection $`(E_0,\mathrm{},E_n)`$ in a category $`𝒟`$ is called full if it generates the category $`𝒟,`$ i.e. the minimal triangulated subcategory of $`𝒟`$ containing all objects $`E_i`$ coincides with $`𝒟.`$ In this case we say that $`𝒟`$ has a semiorthogonal decomposition of the form
$$𝒟=E_0,\mathrm{},E_n.$$
The most studied example of an exceptional collection is the sequence of invertible sheaves $`𝒪_^n,\mathrm{},𝒪_^n(n)`$ on the projective space $`^n`$ (). In particular, this exceptional collection on the projective plane $`^2`$ has length $`3`$.
###### Definition 2.4.
The algebra of a strong exceptional collection $`\sigma =(E_0,\mathrm{},E_n)`$ is the algebra of endomorphisms $`B(\sigma )=\mathrm{End}()`$ of the object $`=\underset{i=0}{\overset{n}{}}E_i.`$
Assume that the triangulated category $`𝒟`$ has a full strong exceptional collection $`(E_0,\mathrm{},E_n)`$ and $`B`$ is the corresponding algebra. Denote by $`\text{mod–}B`$ the category of finitely generated right modules over $`B.`$ There is a theorem according to which if $`𝒟`$ is an enhanced triangulated category in the sense of Bondal and Kapranov , then it is equivalent to the bounded derived category $`𝐃^b(\text{mod–}B)`$. This equivalence is given by the functor $`𝐑\mathrm{Hom}(,)`$ (see ).
For example, if $`𝒟𝐃^b(\mathrm{coh}(X))`$ is the bounded derived category of coherent sheaves on a projective variety $`X,`$ then it is enhanced. Actually, the category of quasi-coherent sheaves $`\mathrm{Qcoh}`$ has enough injectives, and $`𝐃^b(\mathrm{coh}(X))`$ is equivalent to the full subcategory $`𝐃_{\mathrm{coh}}^b(\mathrm{Qcoh}(X))𝐃^b(\mathrm{Qcoh}(X))`$ whose objects are complexes with cohomologies in $`\mathrm{coh}(X).`$
Assume that $`X`$ is smooth and $`(E_0,\mathrm{},E_n)`$ is a strong exceptional collection on $`X.`$ The object $`=_{i=0}^nE_i`$ defines the derived functor
$$𝐑\mathrm{Hom}(,):𝐃^+(\mathrm{Qcoh}(X))𝐃^+(\text{Mod–}B),$$
where $`\text{Mod–}B`$ is the category of all right modules over $`B.`$ Moreover, the functor $`𝐑\mathrm{Hom}(,)`$ sends objects of $`𝐃_{\mathrm{coh}}^b(\mathrm{Qcoh}(X))`$ to objects of the subcategory $`𝐃_{\mathrm{mod}}^b(\text{Mod–}B),`$ which is also equivalent to $`𝐃^b(\text{mod–}B).`$ This gives us a functor
$$𝐑\mathrm{Hom}(,):𝐃^b(\mathrm{coh}(X))𝐃^b(\text{mod–}B).$$
The objects $`E_i`$ for $`i=0,\mathrm{},n`$ are mapped to the projective modules $`P_i=\mathrm{Hom}(,E_i).`$ Moreover, $`B=_{i=0}^nP_i.`$ The algebra $`B`$ has $`n+1`$ primitive idempotents $`e_i,i=0,\mathrm{},n`$ such that $`1_B=e_0+\mathrm{}+e_n`$ and $`e_ie_j=0`$ if $`ij.`$ The right projective modules $`P_i`$ coincide with $`e_iB.`$ The morphisms between them can be easily described since
$$\mathrm{Hom}(P_i,P_j)=\mathrm{Hom}(e_iB,e_jB)e_jBe_i\mathrm{Hom}(E_i,E_j).$$
This yields an equivalence between the triangulated subcategory of $`𝐃^b(\mathrm{coh}(X))`$ generated by the collection $`E_0,\mathrm{},E_n`$ and the derived category $`𝐃^b(\text{mod–}B).`$ Here we use the fact that the algebra $`B`$ has a finite global dimension and any right (and left) module $`M`$ has a finite projective resolution consisting of the projective modules $`P_i`$ with $`i=0,\mathrm{},n.`$ Finally, if the collection $`(E_0,\mathrm{},E_n)`$ is full, then we obtain an equivalence between $`𝐃^b(\mathrm{coh}(X))`$ and $`𝐃^b(\text{mod–}B).`$
Sometimes it is useful to represent the algebra $`B`$ as a category $`𝔅`$ which has $`n+1`$ objects, say $`v_0,\mathrm{},v_n,`$ and morphisms defined by the rule $`\mathrm{Hom}(v_i,v_j)\mathrm{Hom}(E_i,E_j)`$ with the natural composition law. Thus $`B=\underset{0i,jn}{}\mathrm{Hom}(v_i,v_j).`$
###### Theorem 2.5.
Let $`\pi :X_K^2`$ be a blowup of the projective plane $`^2`$ at a set $`K=\{p_1,\mathrm{},p_k\}`$ of any $`k`$ distinct points, and let $`l_1,\mathrm{},l_k`$ be the exceptional curves of the blowup. Let $`(F_0,F_1,F_2)`$ be a full strong exceptional collection of vector bundles on $`^2.`$ Then the sequence
(2.1)
$$(\pi ^{}F_0,\pi ^{}F_1,\pi ^{}F_2,𝒪_{l_1},\mathrm{},𝒪_{l_k}),$$
where the $`𝒪_{l_i}`$ are the structure sheaves of the exceptional $`1`$-curves $`l_i`$, is a full strong exceptional collection on $`X_K`$. Moreover, the sheaves $`𝒪_{l_i}`$ and $`𝒪_{l_j}`$ are mutually orthogonal for all $`ij.`$
In particular, there is an equivalence
(2.2)
$$𝐃^b(\mathrm{coh}(X_K))𝐃^b(\text{mod–}B_K),$$
where $`B_K`$ is the algebra of homomorphisms of the exceptional collection (2.1).
There are no restrictions on the set of points $`K=\{p_1,\mathrm{},p_k\}`$ in this theorem and, in particular, we do not need to assume that $`X_K`$ is a Del Pezzo surface.
We can easily describe the space of morphisms from $`\pi ^{}F_i`$ to the sheaf $`𝒪_{l_j}`$, since it is naturally identified with the space that is dual to the fiber of the vector bundle $`F_i`$ at the point $`p_j^2,`$ i.e.
$$\mathrm{Hom}_{X_K}(\pi ^{}F_i,𝒪_{l_j})\mathrm{Hom}_^2(F_i,𝒪_{p_j}).$$
There are various standard exceptional collections on the projective plane. One of them is the collection of line bundles $`(𝒪,𝒪(1),𝒪(2)),`$ another is the collection $`(𝒪,𝒯_^2(1),𝒪(1)),`$ where $`𝒯_^2`$ is the tangent bundle on $`^2.`$ The latter choice is the most convenient for us. It is easy to see that
$$\mathrm{Hom}(𝒪,𝒯_^2(1))\mathrm{Hom}(𝒯_^2(1),𝒪(1))V\text{and}\mathrm{Hom}(𝒪,𝒪(1))\mathrm{\Lambda }^2VV^{},$$
where $`V`$ is the 3-dimensional vector space whose projectivization $`(V)`$ is the given projective plane $`^2.`$
Let us consider the blowup $`X_K`$ of the projective plane $`(V)`$ at a set $`K=\{p_1,\mathrm{},p_k\}`$ of $`k`$ distinct points, and the exceptional collection
(2.3)
$$\sigma =(𝒪_{X_K},\pi ^{}𝒯_^2(1),\pi ^{}𝒪_^2(1),𝒪_{l_1},\mathrm{},𝒪_{l_k}).$$
Let $`𝔅_K(\sigma )`$ be the category of homomorphisms of this exceptional collection (see Figure 1). Then the surface $`X_K`$ can be recovered from the category $`𝔅_K(\sigma )`$ by means of the following procedure.
Denote by $`S_j`$ the 2-dimensional space of homomorphisms from $`\pi ^{}𝒯_^2(1)`$ to $`𝒪_{l_j}`$ and denote by $`U_j`$ the 1-dimensional space of homomorphisms from $`𝒪(1)`$ to $`𝒪_{l_j}.`$ The composition law in the category $`𝔅_K(\sigma )`$ gives a map from $`U_jV`$ to $`S_j.`$ The kernel of this map is a 1-dimensional subspace $`V_jV`$, which defines a point $`p_j(V)`$. In this way, we can determine all the points $`p_1,\mathrm{},p_k(V)`$ and completely recover the surface $`X_K`$ starting from the category $`𝔅_K(\sigma ).`$
###### Remark 2.6.
Exceptional objects and exceptional collections on Del Pezzo surfaces are well-studied objects. First, any exceptional object of the derived category is isomorphic to a sheaf up to translation. Second, any exceptional sheaf can be included in a full exceptional collection. Third, any full exceptional collection can be obtained from a given one by a sequence of natural operations on exceptional collections called mutations. All these facts can be found in the paper .
### 2.2. Simple degenerations of Del Pezzo surfaces
We now look at some simple degenerations of the situation considered above, namely when two points, for example $`p_1`$ and $`p_2,`$ converge to each other and finally coincide. More precisely, this means that we first blow up a point $`p`$ and after that we blow up some point $`p^{}`$ on the $`1`$-curve which is the pre-image of $`p`$ under the first blowup. This operation is sometimes called a blowup at two “infinitely close” points; more precisely, it corresponds to blowing up a subscheme of length $`2`$ supported at $`p`$. In this case, the pre-image $`\pi ^1(p)`$ consists of two rational curves meeting at one point. One of them is a $`1`$-curve which we denote by $`l^{}`$, and the other is a $`2`$-curve which we denote by $`l`$. The curve $`l`$ is the proper transform of the exceptional curve of the first blow up performed at the point $`p^2.`$
In this paragraph, we consider the situation where the surface $`X_K`$ is the blowup of the projective plane $`^2`$ at a subscheme $`K`$ which is supported at a set of $`k1`$ points $`\{p,p_3,\mathrm{},p_k\}`$ and has length 2 at the point $`p`$. In this case the surface $`X_K`$ is not a Del Pezzo surface, because it possesses a $`2`$-curve $`l`$. However, it follows from general results about blowups that $`𝐃^b(\mathrm{coh}(X_K))`$ still possesses a full exceptional collection .
###### Proposition 2.7.
Let $`X_K`$ be the blowup of $`^2`$ at a subscheme $`K`$ supported at a set of $`k1`$ points $`\{p,p_3,\mathrm{},p_k\}`$ and with length $`2`$ at the point $`p`$. Then the sequence
(2.4)
$$\tau =(𝒪_{X_K},\pi ^{}𝒯_^2(1),\pi ^{}𝒪_^2(1),𝒪_{\pi ^1(p)},𝒪_l^{},𝒪_{l_3},\mathrm{},𝒪_{l_k})$$
is a full exceptional collection on $`X_K.`$
As before we can see that the sheaves $`𝒪_{l_i}`$ and $`𝒪_{l_j}`$ are mutually orthogonal for all $`ij,`$ and each $`𝒪_{l_i}`$ is orthogonal to both $`𝒪_l^{}`$ and $`𝒪_{\pi ^1(p)}.`$ However, the collection $`\tau `$ is not strong, because there are non-trivial morphisms from $`𝒪_{\pi ^1(p)}`$ to $`𝒪_l^{}`$ in degrees 0 and 1. More precisely,
$$\mathrm{Hom}(𝒪_{\pi ^1(p)},𝒪_l^{})\text{and}\mathrm{Ext}^1(𝒪_{\pi ^1(p)},𝒪_l^{}).$$
Denote by $`a`$ and $`b`$ two morphisms from $`𝒪_{\pi ^1(p)}`$ to $`𝒪_l^{}`$ of degrees 0 and 1 respectively. It is easy to see that composition with the morphism $`a`$ gives isomorphisms between the spaces $`\mathrm{Hom}(F,𝒪_{\pi ^1(p)})`$ and $`\mathrm{Hom}(F,𝒪_l^{})`$ for any element $`F`$ of the exceptional collection $`\tau `$ (see Figure 2).
Two approaches can be used to obtain an analogue of equivalence (2.2) for this situation. The first possibility is to associate to the non-strong exceptional collection $`\tau `$ a differential graded algebra of homomorphisms, and obtain an equivalence between the derived category of coherent sheaves and the derived category of finitely generated (right) DG-modules over the DG-algebra of homomorphisms of the exceptional collection. (One could also try to work in the framework of $`A_{\mathrm{}}`$-algebras, which might be more appropriate here considering that the mirror situation involves an $`A_{\mathrm{}}`$-category with non-zero $`m_3`$, see §4.4).
Another approach is to change the exceptional collection $`\tau `$ to another one which is strong. There are natural operations on exceptional collections which are called mutations and which allow us to obtain new exceptional collections starting from a given one.
We omit the definition of mutations, which is classical and can be found in many places. However, we note that the left mutation of the exceptional collection (2.4) in the pair $`(\pi ^{}𝒪_^2(1),𝒪_{\pi ^1(p)})`$ gives us a new exceptional collection
(2.5)
$$\tau ^{}=(𝒪_{X_K},\pi ^{}𝒯_^2(1),,\pi ^{}𝒪_^2(1),𝒪_l^{},𝒪_{l_3},\mathrm{},𝒪_{l_k})$$
where $``$ is the line bundle on $`X_K`$ which is the kernel of the surjection $`\pi ^{}𝒪_^2(1)𝒪_{\pi ^1(p)}.`$ This new exceptional collection $`\tau ^{}`$ is strong.
In fact, we can also consider the same left mutation when the blown up points are all distinct, and obtain in that case as well a strong exceptional collection
$$\sigma ^{}=(𝒪_{X_K},\pi ^{}𝒯_^2(1),,\pi ^{}𝒪_^2(1),𝒪_{l_2},𝒪_{l_3},\mathrm{},𝒪_{l_k}),$$
which behaves very much like $`\tau ^{}`$. The distinguishing feature of the case where we blow up the point $`p`$ twice is that in the exceptional collection $`\tau ^{}`$ the composition map
(2.6)
$$\mathrm{Hom}(,\pi ^{}𝒪_^2(1))\mathrm{Hom}(\pi ^{}𝒪_^2(1),𝒪_l^{})\mathrm{Hom}(,𝒪_l^{})$$
is identically zero, whereas for $`\sigma ^{}`$ (i.e., when the points of $`K`$ are distinct) the corresponding composition is non-trivial. In this sense, the mutation allows us to give a simple description of the behaviour of the category under the degeneration process where two points of $`K`$ converge to each other. Namely, the algebra $`B_K(\tau ^{})`$ of homomorphisms of the exceptional collection $`\tau ^{}`$ is obtained as a degeneration of the algebra of homomorphisms of the exceptional collection $`\sigma ^{}`$ in which the composition (2.6) becomes zero.
###### Proposition 2.8.
Let $`X_K`$ be the blowup of $`^2`$ at a subscheme $`K`$ supported at a set of $`k1`$ points $`\{p,p_3,\mathrm{},p_k\}`$ and with length $`2`$ at the point $`p`$. Then there is an equivalence $`𝐃^b(\mathrm{coh}(X_K))𝐃^b(\text{mod–}B_K(\tau ^{})),`$ where $`B_K(\tau ^{})`$ is the algebra of homomorphisms of the exceptional collection $`\tau ^{}.`$
In this context, the surface $`X_K`$ can again be recovered from the category $`𝔅_K(\tau ^{})`$. Namely, the points $`p,p_3,\mathrm{},p_k`$ can be determined by the same method as above. To recover $`X_K`$, we also have to determine the position of the point $`p^{}`$ on the exceptional curve of the blowup of the point $`p.`$ This is equivalent to finding a tangent direction at the point $`p.`$ Consider the kernel of the composition map
$$\mathrm{Hom}(𝒪,)\mathrm{Hom}(,𝒪_l^{})\mathrm{Hom}(𝒪,𝒪_l^{}).$$
It is a one-dimensional subspace of $`\mathrm{Hom}(𝒪,).`$ The image of this subspace in the space $`V^{}=\mathrm{Hom}(𝒪,\pi ^{}𝒪(1))`$ determines a line in the projective space $`(V)`$ which passes through the point $`p`$ and hence a tangent direction at $`p`$.
### 2.3. Noncommutative deformations of Del Pezzo surfaces
As before, let $`X_K`$ be the blowup of the projective plane at a set $`K`$ of $`k`$ distinct points. Consider the strong exceptional collection
$$\sigma =(𝒪,\pi ^{}𝒯_^2(1),\pi ^{}𝒪_^2(1),𝒪_{l_1},\mathrm{},𝒪_{l_k}).$$
By the discussion in §2.1, the derived category of coherent sheaves $`𝐃^b(\mathrm{coh}(X_K))`$ is equivalent to the category of finitely generated (right) modules over the algebra $`B_K`$ of homomorphisms of $`\sigma `$. The algebra $`B_K`$ can also be represented by the category $`𝔅_K`$ associated to the exceptional collection $`\sigma `$ (see Figure 1).
The category $`𝔅_K`$ has strictly more deformations than the surface $`X_K.`$ We saw above that the surface $`X_K`$ can be reconstructed from the category $`𝔅_K`$, and that the deformation of the surface $`X_K`$ is controlled by the variation of the set $`K^2.`$
A general deformation of the category $`𝔅_K`$ can be viewed as the category of an exceptional collection on a noncommutative deformation of the surface $`X_K.`$ In other words, if $`𝔅_{K,\mu }`$ is a deformation of the category $`𝔅_K`$ then the bounded derived category $`𝐃^b(\text{mod–}B_{K,\mu })`$ of finitely generated (right) modules over the algebra $`B_{K,\mu }`$ will be viewed as the derived category of coherent sheaves on a noncommutative surface $`X_{K,\mu }.`$ Any such noncommutative surface can be represented as the blowup of a noncommutative plane $`_\mu ^2`$ at some set $`K`$ consisting of $`k`$ of its “points”. This procedure is discussed in detail in .
In the rest of this section, we describe the deformations of the category $`𝔅_K.`$ Recall that a deformation of a category is, by definition, a deformation of its composition law. We proceed in two steps. The first step is to describe the deformations of the subcategory $`𝔅(\sigma _0)`$ associated to the subcollection $`\sigma _0=(𝒪,\pi ^{}𝒯(1),\pi ^{}𝒪(1)).`$ This subcategory $`𝔅(\sigma _0)`$ is the category of homomorphisms of the full strong exceptional collection $`(𝒪,𝒯(1),𝒪(1))`$ on $`^2.`$ Therefore, considering a deformation of the subcategory $`𝔅(\sigma _0)`$ we obtain a noncommutative deformation $`_\mu ^2`$ of the projective plane. The second step is to describe the deformations of all other compositions in the category $`𝔅_K.`$ These deformations correspond to variations of the set of “points” $`K`$ on the noncommutative projective plane $`_\mu ^2.`$
Noncommutative deformations of the projective plane have been described in . Any deformation of the category $`𝔅(\sigma _0)`$ is a category with three ordered objects $`F_0,F_1,F_2`$ and with three-dimensional spaces of homomorphisms from $`F_i`$ to $`F_j`$ when $`i<j`$ (see Figure 4). Any such category $`𝔅_\mu `$ is determined by the composition tensor $`\mu :VUW.`$ We will consider only the nondegenerate (geometric) case, where the restrictions $`\mu _u=\mu (,u):VW`$ and $`\mu _v=\mu (v,):UW`$ have rank at least two for all nonzero elements $`uU`$ and $`vV`$, and the composition of $`\mu `$ with any nonzero linear form on $`W`$ is a bilinear form of rank at least two on $`VU`$. The equations $`det\mu _u=0`$ and $`det\mu _v=0`$ define closed subschemes $`\mathrm{\Gamma }_U(U)`$ and $`\mathrm{\Gamma }_V(V).`$ Namely, up to projectivization the set of points of $`\mathrm{\Gamma }_U`$ (resp. $`\mathrm{\Gamma }_V`$) consists of all $`uU`$ (resp. $`vV`$) for which the rank of $`\mu _u`$ (resp. $`\mu _v`$) is equal to $`2.`$ It is easy to see that the correspondence which associates to a vector $`vV`$ the kernel of the map $`\mu _v:UW`$ defines an isomorphism between $`\mathrm{\Gamma }_V`$ and $`\mathrm{\Gamma }_U`$. Moreover, under these circumstances $`\mathrm{\Gamma }_V`$ is either the entire projective plane $`(V)`$ or a cubic in $`(V).`$ If $`\mathrm{\Gamma }_V=(V)`$ then $`\mu `$ is isomorphic to the tensor $`VV\mathrm{\Lambda }^2V,`$ i.e. we get the usual projective plane $`^2.`$
Thus, the non-trivial case is the situation where $`\mathrm{\Gamma }_V`$ is a cubic, i.e. an elliptic curve which we now denote by $`E`$. This elliptic curve comes equipped with two embeddings into the projective planes $`(U)`$ and $`(V)`$ respectively; by restriction of $`𝒪(1)`$ these embeddings determine two line bundles $`_1`$ and $`_2`$ of degree $`3`$ over $`E`$, and it can be checked that $`_1_2.`$ This construction has a converse:
###### Construction 2.9.
The tensor $`\mu `$ can be reconstructed from the triple $`(E,_1,_2).`$ Namely, the spaces $`U,V`$ are isomorphic to $`H^0(E,_1)^{}`$ and $`H^0(E,_2)^{}`$ respectively, and the space $`W`$ is dual to the kernel of the canonical map
$$H^0(E,_1)H^0(E,_2)H^0(E,_1_2),$$
which induces the tensor $`\mu :VUW`$.
The details of these constructions and statements can be found in .
###### Remark 2.10.
Note that we can also consider a triple $`(E,_1,_2)`$ such that $`_1_2.`$ Then the procedure described above produces a tensor $`\mu `$ with $`\mathrm{\Gamma }_V(V)`$, which defines the usual commutative projective plane. Thus, in this particular case the tensor $`\mu `$ does not depend on the curve $`E.`$
Now we describe the deformations of the other compositions in the category $`𝔅_K`$. Given a category $`𝔅_\mu `$ of the form described above, corresponding to a noncommutative projective plane $`_\mu ^2`$, and given a set $`K=\{p_1,\mathrm{},p_k\}`$ of $`k`$ points on the elliptic curve $`E`$, we can construct a category $`𝔅_{K,\mu }`$ in the following manner. A point $`p_jE(U)`$ determines a one-dimensional subspace of $`U`$, generated by a vector $`u_jU`$. The map $`\mu _{u_j}:VW`$ has rank $`2`$; denote by $`v_j`$ a non-zero vector in its kernel. The category $`𝔅_{K,\mu }`$ is then constructed from the category $`𝔅_\mu `$ by adding $`k`$ mutually orthogonal objects $`𝒪_{l_j}`$ for $`j=1,\mathrm{},k`$, and defining the spaces of morphisms by the rule
$$\mathrm{Hom}(F_2,𝒪_{l_j})=,\mathrm{Hom}(F_1,𝒪_{l_j})=V/\mathrm{Ker}\mu _{u_j}=V/v_j,\mathrm{Hom}(F_0,𝒪_{l_j})=W/\mathrm{Im}\mu _{v_j}.$$
The two composition tensors involving $`\mathrm{Hom}(F_2,𝒪_{l_j})`$ are defined in the obvious manner as suggested by the notation. The only non-obvious composition is the map $`V/v_jUW/\mathrm{Im}\mu _{v_j}`$, which is by definition induced by the tensor $`\mu :VUW.`$
Conversely, if we consider a category $`𝔅_{K,\mu }`$ which is a deformation of $`𝔅_K`$ and an extension of the category $`𝔅_\mu `$, then the kernel of the composition map
$$\mathrm{Hom}(F_2,𝒪_{l_j})V\mathrm{Hom}(F_1,𝒪_{l_j})$$
defines a one-dimensional subspace $`v_jV`$. The map $`\mu _{v_j}`$ must have rank $`2`$, since otherwise $`\mu _{v_j}`$ would be an isomorphism and the composition map $`\mathrm{Hom}(F_2,𝒪_{l_j})W\mathrm{Hom}(F_0,𝒪_{l_j})`$ would vanish identically, which by assumption is not the case. Therefore, the objects $`𝒪_{l_j}`$ correspond to points on the curve $`E`$.
Thus, any category $`𝔅_{K,\mu }`$ is defined by the data $`(E,_1,_2,p_1,\mathrm{},p_k),`$ where $`E`$ is a cubic, $`_1,_2`$ are line bundles of degree $`3`$ on $`E`$, and $`p_1,\mathrm{},p_k`$ is a set of distinct points on $`E`$. If $`_1_2,`$ then we obtain the category $`𝔅_K`$ related to a blowup of the usual commutative projective plane. In the general case, the bounded derived category $`𝐃^b(\text{mod–}B_{K,\mu })`$ of finite rank modules over the algebra $`B_{K,\mu }`$ is viewed as the derived category of coherent sheaves on the non-commutative surface $`X_{K,\mu }`$, which is a blowup of $`k`$ points on the non-commutative projective plane $`_\mu ^2`$.
A standard approach to noncommutative geometry is to determine a noncommutative variety either by an abelian category of (quasi)coherent sheaves on it or by a noncommutative (graded) algebra which is considered as its (homogeneous) coordinate ring. The question of how to define the abelian category of coherent sheaves on Del Pezzo surfaces and on other blowups of surfaces is discussed in the paper . We briefly describe one of the possible approaches. It is very important to note that the category $`𝐃^b(\text{mod–}B_{K,\mu })`$ possess a Serre functor $`S`$, i.e. an additive autoequivalence for which there are bi-functorial isomorphisms
$$\mathrm{Hom}(X,SY)\stackrel{}{}\mathrm{Hom}(Y,X)^{}$$
for any $`X,Y𝐃^b(\text{mod–}B_{K,\mu })`$. In the case of the bounded derived category of finite rank modules over a finite dimensional algebra of finite homological dimension, the Serre functor is the functor which takes a complex of modules $`M^{}`$ to the complex $`𝐑\mathrm{Hom}_{B_{K,\mu }}(M,B_{K,\mu })^{}`$. The Serre functor is an exact autoequivalence.
Now we can take the projective module $`P_0`$ (corresponding to $`𝒪`$, see the discussion after Definition 2.4) and consider the sequence of objects $`R_m=S^m[2m]P_0`$ for all $`m.`$ Let us consider the subcategory $`𝒜𝐃^b(\text{mod–}B_{K,\mu })`$ consisting of all objects $`F`$ such that
$$\mathrm{Hom}(R_m,F[i])=0\text{for all}i0\text{and sufficiently large}m0.$$
If the category $`𝒜`$ is abelian and its bounded derived category $`𝐃^b(𝒜)`$ is equivalent to $`𝐃^b(\text{mod–}B_{K,\mu })`$ then $`𝒜`$ can be considered as the category of coherent sheaves on the noncommutative surface $`X_{K,\mu }`$, and $`X_{K,\mu }`$ can be called a noncommutative Del Pezzo surface.
The reason of such a definition of the abelian category of coherent sheaves on a noncommutative Del Pezzo surface is inspired by the commutative case. In the commutative case the Serre functor is isomorphic to the functor $`𝒪(K)[2],`$ where $`𝒪(K)`$ is the canonical line bundle. Hence, for usual commutative surfaces the objects $`R_m`$ are isomorphic to the invertible sheaves $`𝒪(mK).`$ Since for a Del Pezzo surface $`X`$ the anticanonical sheaf $`𝒪(K)`$ is ample, we have $`H^i(X,F(mK))=0`$ for all $`i0`$ and any coherent sheaf $`F`$ when $`m`$ is sufficiently large. This property makes it possible to separate pure coherent sheaves from other complexes of coherent sheaves.
We can also consider the graded space $`A=_{p=0}^{\mathrm{}}\mathrm{Hom}(R_0,R_p)`$ and can endow it with the structure of a graded algebra using the isomorphisms $`\mathrm{Hom}(R_0,R_p)\mathrm{Hom}(R_i,R_{ip})`$ given by the functors $`S^i[2i]`$ for all $`i.`$ This algebra can be considered as the homogeneous coordinate ring of a noncommutative Del Pezzo surface. It seems that such rings are noncommutative deformations of homogeneous commutative coordinate rings of usual Del Pezzo surfaces.
In any case, these remarks about abelian categories of coherent sheaves on noncommutative Del Pezzo surfaces will not be needed in the rest of the argument. We will only use the description of the bounded derived category of coherent sheaves on the noncommutative blowup $`X_{K,\mu }`$ in terms of finite rank modules over the algebra $`B_{K,\mu },`$ i.e. we state an equivalence of triangulated categories
(2.7)
$$𝐃^b(\mathrm{coh}(X_{K,\mu }))𝐃^b(\text{mod–}B_{K,\mu }).$$
## 3. The mirror Landau-Ginzburg models
### 3.1. Compactification of the mirror of $`^2`$
As mentioned in the introduction, the mirror of $`^2`$ is an elliptic fibration with 3 singular fibers, determined by (a fiberwise compactification of) the superpotential $`W_0=x+y+1/xy`$ on $`(^{})^2`$. This Landau-Ginzburg model compactifies naturally to an elliptic fibration $`\overline{W_0}:\overline{M}^1`$, which we now describe.
Compactifying $`(^{})^2`$ to $`^2`$, we can view $`W_0`$ as the quotient of the two homogeneous degree 3 polynomials $`P_0=X^2Y+XY^2+Z^3`$ and $`P_{\mathrm{}}=XYZ`$, which define a pencil of cubics with three base points of multiplicities respectively $`4`$, $`4`$, and $`1`$. Namely, the cubic $`C_0`$ defined by $`P_0`$ intersects the line $`X=0`$ at $`(0:1:0)`$ (with multiplicity 3), the line $`Y=0`$ at $`(1:0:0)`$ (with multiplicity 3), and the line $`Z=0`$ at $`(0:1:0)`$, $`(1:0:0)`$ and $`(1:1:0)`$. Blow up $`^2`$ three times successively at the point where the cubic $`C_0`$ and the line $`X=0`$ (or their proper transforms) intersect each other, i.e. first at the point $`(0:1:0)`$, and then twice at suitable points of the exceptional divisors (see Figure 5). Similarly, blow up three times the intersection of the cubic $`C_0`$ with the line $`Y=0`$.
Let $`\stackrel{~}{C}_0`$ be the proper transform of $`C_0`$ under these blowups, and let $`\stackrel{~}{C}_{\mathrm{}}`$ be the configuration of 9 rational curves formed by the proper transforms of the three coordinate lines and the exceptional divisors of the six blowups (so, in Figure 5, all components other than $`\stackrel{~}{C}_0`$ are eventually part of $`\stackrel{~}{C}_{\mathrm{}}`$). Then $`\stackrel{~}{C}_0`$ and $`\stackrel{~}{C}_{\mathrm{}}`$ intersect transversely at three smooth points, and define a pencil of elliptic curves representing the anticanonical class in $`^2`$ blown up six times. The complement of $`\stackrel{~}{C}_{\mathrm{}}`$ identifies with $`(^{})^2`$, and the restriction of the $`^1`$-valued map defined by the pencil to this open subset coincides with $`W_0`$. Blowing up the three points where $`\stackrel{~}{C}_0`$ and $`\stackrel{~}{C}_{\mathrm{}}`$ intersect, we obtain a rational elliptic surface $`\overline{M}`$, and the pencil becomes an elliptic fibration $`\overline{W_0}:\overline{M}^1`$, which provides a natural compactification of $`W_0:(^{})^2`$.
The meromorphic function $`\overline{W_0}`$ has 12 isolated non-degenerate critical points. Three of them are the pre-images of the points $`(1:1:1)`$, $`(j:j:1)`$, and $`(j^2:j^2:1)`$ ($`j=e^{2i\pi /3}`$), and correspond to the three critical points of $`W_0`$ in $`(^{})^2`$ (with associated critical values $`3`$, $`3j`$, and $`3j^2`$). The nine other critical points all lie in the fiber above infinity: they are the nodes of the reducible configuration $`\stackrel{~}{C}_{\mathrm{}}`$ (see Figure 6).
This compactification process can also be described in a more symmetric manner by viewing $`(^{})^2`$ as the surface $`\{xyz=1\}(^{})^3`$, and $`W_0=x+y+z`$. Compactifying $`(^{})^3`$ to $`^3`$ leads one to consider the cubic surface $`\{XYZ=T^3\}^3`$, which presents $`A_2`$ singularities at the three points $`(1:0:0:0)`$, $`(0:1:0:0)`$, and $`(0:0:1:0)`$. After blowing up $`^3`$ at these three points, we obtain a smooth cubic surface, in which the hyperplane sections $`\stackrel{~}{C}_0=\{X+Y+Z=0\}`$ and $`\stackrel{~}{C}_{\mathrm{}}=\{T=0\}`$ define a pencil of elliptic curves with three base points. As before, $`\stackrel{~}{C}_0`$ is a smooth elliptic curve, and $`\stackrel{~}{C}_{\mathrm{}}`$ is a configuration of 9 rational curves (the proper transforms of the three coordinate lines where the singular cubic surface intersects the plane $`T=0`$, and the six $`2`$-curves arising from the resolution of the singularities). Blowing up the three points of $`\stackrel{~}{C}_0\stackrel{~}{C}_{\mathrm{}}`$, we again obtain a rational elliptic surface, and an elliptic fibration with 12 isolated critical points, 9 of which lie in the fiber above infinity (as in Figure 6).
### 3.2. The vanishing cycles of $`\overline{W_0}`$
Each singular fiber of $`\overline{W_0}`$ is obtained from the regular fiber by collapsing a certain number of vanishing cycles, and the monodromy of the fibration around a singular fiber is given by a product of Dehn twists along these vanishing cycles. In this section, we determine the homology classes of the vanishing cycles associated to the critical points of $`\overline{W_0}`$.
More precisely, consider the fiber $`\mathrm{\Sigma }_0=\overline{W_0}{}_{}{}^{1}(0)`$, which is a smooth elliptic curve (in fact, the proper transform of the curve called $`\stackrel{~}{C}_0`$ in §3.1), and consider the following ordered collection of arcs $`(\gamma _i)_{0i3}`$ joining the origin to the various critical values of $`\overline{W_0}`$: $`\gamma _0,\gamma _1,\gamma _2`$ are straight line segments joining the origin to $`\lambda _0=3`$, $`\lambda _1=3j^2`$, and $`\lambda _2=3j`$ respectively, and $`\gamma _3`$ is the straight line $`e^{i\pi /3}_+`$ joining the origin to $`\lambda _3=\mathrm{}`$.
Using parallel transport (with respect to an arbitrary horizontal distribution) along the arc $`\gamma _i`$, we can associate a vanishing cycle to each critical point $`p\overline{W_0}{}_{}{}^{1}(\lambda _i)`$; this vanishing cycle is well-defined up to isotopy, and in particular we can consider its homology class in $`H_1(\mathrm{\Sigma }_0,)^2`$ (well-defined up to a choice of orientation). If we fix a symplectic structure on $`\overline{M}`$ for which the fibers of $`\overline{W_0}`$ are symplectic submanifolds, then we have a canonical horizontal distribution (given by the symplectic orthogonal to the fiber), which allows us to consider the vanishing cycles as Lagrangian submanifolds of $`\mathrm{\Sigma }_0`$, well-defined up to Hamiltonian isotopy; in §4 this will be of utmost importance, but for now we ignore the symplectic structure and only view $`\overline{W_0}`$ as a topological fibration.
###### Lemma 3.1.
In terms of a suitable basis $`\{a,b\}`$ of $`H_1(\mathrm{\Sigma }_0,)`$, the vanishing cycles $`L_0,L_1,L_2`$ associated to the critical values $`\lambda _0,\lambda _1,\lambda _2`$ $`(`$and the arcs $`\gamma _0,\gamma _1,\gamma _2)`$ represent the classes $`[L_0]=2ab`$, $`[L_1]=a+b`$, and $`[L_2]=a+2b`$, respectively; and the vanishing cycles $`L_3,\mathrm{},L_{11}`$ associated to the nine critical points in the fiber at infinity represent the class $`[L_3]=\mathrm{}=[L_{11}]=a+b`$.
###### Proof.
The vanishing cycles $`L_0,L_1,L_2`$ are exactly those of the mirror of $`^2`$, and are well-known (cf. e.g. or ). In particular it is known that, choosing a suitable homology basis $`\{a,b\}`$ for $`H_1(\mathrm{\Sigma }_0,)`$, and fixing appropriate orientations of $`L_0,L_1,L_2`$, we have $`[L_0]=2ab`$, $`[L_1]=a+b`$, and $`[L_2]=a+2b`$ (cf. e.g. Figure 14 in ).
We now consider the 9 critical points in the fiber at infinity. It is clear that $`L_3,\mathrm{},L_{11}`$ admit disjoint representatives, and hence are all homologous. Their homology class can be determined by considering the monodromy of the elliptic fibration $`\overline{W_0}`$, which is given by the product of the positive Dehn twists along the vanishing cycles. Considering the action on $`H_1(\mathrm{\Sigma }_0,)`$, and still using the basis $`\{a,b\}`$ considered above, the monodromies around the critical values $`\lambda _0,\lambda _1,\lambda _2`$ are given by
$$\tau _0=\left(\begin{array}{cc}\hfill 1& \hfill 4\\ \hfill 1& \hfill 3\end{array}\right),\tau _1=\left(\begin{array}{cc}\hfill 2& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right),\mathrm{and}\tau _2=\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 4& \hfill 3\end{array}\right),$$
while the monodromy around the fiber at infinity is given by $`\tau ^9`$, where $`\tau `$ is the positive Dehn twist along $`[L_3]=\mathrm{}=[L_{11}]`$. On the other hand, because the arcs $`\gamma _0,\mathrm{},\gamma _3`$ are ordered clockwise around the origin, we have $`\tau _0\tau _1\tau _2\tau ^9=1`$. Therefore,
$$\tau ^9=\left(\begin{array}{cc}\hfill 8& \hfill 9\\ \hfill 9& \hfill 10\end{array}\right),$$
and considering $`\mathrm{Ker}(\tau ^91)`$ we obtain $`[L_3]=\mathrm{}=[L_{11}]=a+b`$. ∎
###### Alternative proof.
(compare with §4.2 of ). Recall that $`\overline{M}`$ is obtained from $`^2`$ by successive blowups of the base points of the pencil of cubics defined by $`P_0=X^2Y+XY^2+Z^3`$ and $`P_{\mathrm{}}=XYZ`$. Consider the ruled surface $`F`$ obtained by blowing up $`^2`$ just once at the point $`(0:1:0)`$: the projection $`(X:Y:Z)(X:Z)`$ naturally extends into a fibration $`\pi _x:F^1`$, of which the exceptional divisor is a section. For $`\lambda ^1`$, denote by $`\widehat{C}_\lambda `$ the proper transform of the plane cubic $`C_\lambda `$ defined by $`P_0\lambda P_{\mathrm{}}`$, which is also the image of $`\overline{W_0}{}_{}{}^{1}(\lambda )`$ under the natural projection $`p:\overline{M}F`$.
The restriction $`\pi _{x,\lambda }`$ of $`\pi _x`$ to $`\widehat{C}_\lambda `$ has degree two, and for $`\lambda \mathrm{crit}(\overline{W_0})`$ its four branch points are associated to distinct critical values in $`^1`$, namely zero and the three roots of the equation $`x(\lambda x)^2=4`$. Indeed, since $`C_\lambda `$ always has an order 3 tangency with the line $`X=0`$ at $`(0:1:0)`$, $`\widehat{C}_\lambda `$ is always tangent to the fiber $`\pi _x^1(0)`$. The three other branch points are the critical points of the projection to the first coordinate on $`(^{})^2C_\lambda =\{(x,y)(^{})^2,xy(\lambda xy)=1\}`$; viewing $`xy(\lambda xy)=1`$ as a quadratic equation in the variable $`y`$, the discriminant is $`x(\lambda x)^24`$.
As $`\lambda `$ tends to $`\lambda _i`$ ($`i\{0,1,2\}`$), two of the critical values of $`\pi _{x,\lambda }`$ converge to each other; keeping track of the manner in which these critical values coalesce when $`\lambda `$ varies from $`0`$ to $`\lambda _i`$ along the arc $`\gamma _i`$, we obtain an arc $`\delta _i^1`$, with end points in $`\mathrm{crit}(\pi _{x,0})`$ (see Figure 7). The lift of $`\delta _i`$ under the double cover $`\pi _{x,0}`$ is (up to homotopy) the vanishing cycle $`L_i`$ (note that the projection $`p:\overline{M}F`$ allows us to implicitly identify $`\widehat{C}_\lambda `$ with $`\overline{W_0}{}_{}{}^{1}(\lambda )`$ for $`\lambda \mathrm{}`$).
Similarly, the behavior of the critical values of $`\pi _{x,\lambda }`$ as $`\lambda `$ tends to infinity describes the degeneration of $`\widehat{C}_\lambda `$ to the singular configuration $`\widehat{C}_{\mathrm{}}`$, which consists of two sections and two fibers of $`\pi _x:F^1`$ (the fibers above $`0`$ and $`\mathrm{}`$, the exceptional section, and the pre-image of the line $`Y=0`$). Namely, as $`\lambda `$ tends to infinity along the arc $`\gamma _3`$, the critical value with argument $`2\pi /3`$ approaches zero, while the two other roots of $`x(\lambda x)^24`$ tend to infinity. The manner in which pairs of critical values coalesce is encoded by the arcs $`\delta ^{}`$ and $`\delta ^{\prime \prime }`$ in Figure 7, and the four vanishing cycles associated to the degeneration are essentially the lifts under $`\pi _{x,0}`$ of closed loops which bound regular neighborhoods of the arcs $`\delta ^{}`$ and $`\delta ^{\prime \prime }`$; they all represent the same homotopy class inside $`\widehat{C}_0`$.
Recall that $`\overline{W_0}{}_{}{}^{1}(\mathrm{})\stackrel{~}{C}_{\mathrm{}}`$ is obtained from $`\widehat{C}_{\mathrm{}}`$ by repeatedly blowing up two of the nodes. Taking pre-images under these blowup operations, the vanishing cycles associated to the two other nodes of $`\widehat{C}_{\mathrm{}}`$ are naturally identified with two of the nine vanishing cycles $`L_3,\mathrm{},L_{11}`$ associated to the fiber at infinity of $`\overline{W_0}`$. In particular, these vanishing cycles represent the same homology class in $`H_1(\mathrm{\Sigma }_0,)H_1(\widehat{C}_0,)`$ as the lifts of $`\delta ^{}`$ and $`\delta ^{\prime \prime }`$.
It is then easy to check that, for suitable choices of orientations, we have $`[L_0][L_1]=[L_0][L_2]=[L_1][L_2]=3`$, $`[L_0][L_{3+i}]=[L_2][L_{3+i}]=1`$, and $`[L_1][L_{3+i}]=2`$, which completes the proof of Lemma 3.1. ∎
### 3.3. The vanishing cycles of $`(M_k,W_k)`$
Recall from the introduction that our proposal for the mirror of a Del Pezzo surface $`X_K`$ obtained from $`^2`$ by blowing up $`k8`$ generic points is an elliptic fibration $`W_k:M_k`$, obtained by deforming the fibration $`\overline{W_0}`$ to another elliptic fibration $`\overline{W_k}:\overline{M}^1`$, and considering the restriction to $`M_k=\overline{M}\overline{W_k}{}_{}{}^{1}(\mathrm{})`$. More precisely, remember that $`\overline{W_k}`$ has $`3+k`$ irreducible nodal fibers corresponding to critical values $`\lambda _0,\mathrm{},\lambda _{k+2}`$, of which the first three correspond naturally to the irreducible nodal fibers of $`\overline{W_0}`$, while the $`k`$ other finite critical values correspond to the deformation of critical points in $`\overline{W_0}{}_{}{}^{1}(\mathrm{})`$ towards finite values of the superpotential. Meanwhile, $`\overline{W_k}{}_{}{}^{1}(\mathrm{})`$ is a singular fiber with $`9k`$ components.
While the precise locations of the critical values $`\lambda _i`$ are closely related to the complex structure on $`M_k`$, they are essentially irrelevant from the point of view of symplectic topology and categories of vanishing cycles. Indeed, if we consider a family $`(M_{k,t},W_{k,t})`$ of fibrations indexed by a real parameter $`t`$, with the property that for all $`t`$ the critical points of $`W_{k,t}`$ are isolated and non-degenerate, then the vanishing cycles remain the same for all values of $`t`$, up to smooth isotopies inside the reference fiber. For this reason, we do not need to make a specific choice of $`\lambda _i`$. To fix ideas, let us say that $`\lambda _0`$ is close to $`3`$, $`\lambda _1`$ is close to $`3j^2`$, $`\lambda _2`$ is close to $`3j`$, and $`\lambda _i`$ is close to infinity for $`i3`$; we again choose the origin as base point, and note that the smooth elliptic curve $`W_k^1(0)`$ is diffeomorphic to $`\overline{W_0}{}_{}{}^{1}(0)`$, so we implicitly identify them and again call our reference fiber $`\mathrm{\Sigma }_0`$. We also choose an ordered collection of arcs $`\gamma _i`$ joining the origin to $`\lambda _i`$ which lie close to those considered in §3.2, thus ensuring that the homology classes $`[L_0],\mathrm{},[L_{k+2}]H_1(\mathrm{\Sigma }_0,)`$ of the corresponding vanishing cycles remain those given by Lemma 3.1.
Fixing a symplectic form $`\omega _k`$ on $`M_k`$ (compatible with $`W_k`$, i.e. restricting positively to the fibers), the vanishing cycles $`L_0,\mathrm{},L_{k+2}`$ associated to the arcs $`\gamma _0,\mathrm{},\gamma _{k+2}`$ naturally become Lagrangian submanifolds of the reference fiber $`(\mathrm{\Sigma }_0,\omega _{k|\mathrm{\Sigma }_0})`$ (cf. e.g. ). Indeed, the symplectic form defines a natural horizontal distribution outside of the critical points of $`W_k`$, given by the symplectic orthogonal to the fiber. Using this horizontal distribution, parallel transport induces symplectomorphisms between the smooth fibers, and the vanishing cycle $`L_i`$ is by definition the set of points in the reference fiber $`\mathrm{\Sigma }_0`$ for which parallel transport along $`\gamma _i`$ converges to the critical point in the fiber $`W_k^1(\lambda _i)`$. It is also useful to consider the Lefschetz thimble $`D_i`$, which is the set of points swept out by parallel transport of $`L_i`$ above $`\gamma _i`$; by construction, $`D_i`$ is a Lagrangian disk in $`(M_k,\omega _k)`$, fibered above the arc $`\gamma _i`$, and $`D_i=L_i`$.
We recall the following classical result (we provide a proof for completeness):
###### Lemma 3.2.
A deformation of the system of arcs $`\{\gamma _i\}`$ by an isotopy in $`\mathrm{Diff}(,\mathrm{crit}(W_k))`$ affects the vanishing cycles $`L_i`$ by Hamiltonian isotopies; moreover, the same property holds if the symplectic fibration $`(M_k,\omega _k,W_k)`$ is deformed in a manner such that the cohomology class $`[\omega _k]`$ remains constant and the critical points of $`W_k`$ remain isolated and non-degenerate.
###### Proof.
We first consider a deformation of the system of arcs $`\{\gamma _i\}`$, based at a regular value $`\lambda _{}\mathrm{crit}(W_k)`$ (in our case the origin), to an isotopic system of arcs $`\{\gamma _i^{}\}`$ based at a regular value $`\lambda _{}^{}`$. This means that we are given an arc $`\delta :[0,1]\mathrm{crit}(W_k)`$ joining $`\lambda _{}`$ to $`\lambda _{}^{}`$, and continuous families of arcs $`\{\gamma _{i,t}\}`$, $`0t1`$, with $`\gamma _{i,0}=\gamma _i`$ and $`\gamma _{i,1}=\gamma _i^{}`$, such that $`\gamma _{i,t}`$ joins the regular value $`\delta (t)`$ to the critical value $`\lambda _i`$, and $`\{\gamma _{i,t}\}_{0ik+2}`$ is an ordered collection of arcs for all $`t[0,1]`$. The vanishing cycles $`L_i^{}`$ associated to the arcs $`\gamma _i^{}`$ live inside $`\mathrm{\Sigma }_{}^{}=W_k^1(\lambda _{}^{})`$, while the original vanishing cycles $`L_i`$ associated to $`\gamma _i`$ are submanifolds of $`\mathrm{\Sigma }_{}=W_k^1(\lambda _{})`$. However, we claim that the isotopy induces a symplectomorphism $`\varphi :\mathrm{\Sigma }_{}\mathrm{\Sigma }_{}^{}`$ with the property that $`\varphi (L_i)`$ and $`L_i^{}`$ are mutually Hamiltonian isotopic for all $`i`$; this is the meaning of the statement of the lemma.
Namely, parallel transport along the arc $`\delta `$ (using the horizontal distribution described above) induces a symplectomorphism $`\varphi `$ from $`\mathrm{\Sigma }_{}=W_k^1(\lambda _{})`$ to $`\mathrm{\Sigma }_{}^{}=W_k^1(\lambda _{}^{})`$. For all $`t[0,1]`$ we can consider the vanishing cycle $`L_{i,t}W_k^1(\delta (t))`$ associated to the arc $`\gamma _{i,t}`$, and its image $`L_{i,t}^{}\mathrm{\Sigma }_{}^{}`$ under the symplectomorphism induced by parallel transport along $`\delta ([t,1])`$. The family $`L_{i,t}^{}`$, $`t[0,1]`$ defines a smooth isotopy from $`L_{i,0}^{}=\varphi (L_i)`$ to $`L_{i,1}^{}=L_i^{}`$ through Lagrangian submanifolds of $`\mathrm{\Sigma }_{}^{}`$. Moreover, each vanishing cycle $`L_{i,t}W_k^1(\delta (t))`$ bounds a Lagrangian thimble $`D_{i,t}`$, and the cylinder $`C_{i,t}`$ swept by $`L_{i,t}`$ under parallel transport along $`\delta ([t,1])`$ is also Lagrangian. By continuity, the relative cycles $`D_{i,t}C_{i,t}`$ (with boundary $`L_{i,t}^{}`$) all represent the same relative homotopy class in $`\pi _2(M_k,\mathrm{\Sigma }_{}^{})`$, and since $`D_{i,t}`$ and $`C_{i,t}`$ are Lagrangian they all have zero symplectic area. This implies that the Lagrangian submanifolds $`L_{i,t}^{}`$, $`t[0,1]`$ are all Hamiltonian isotopic inside $`\mathrm{\Sigma }_{}^{}`$; in particular, $`\varphi (L_i)`$ and $`L_i^{}`$ are Hamiltonian isotopic.
We now consider a symplectic fibration $`W_k^{}:(M_k,\omega _k^{})`$ which is isotopic to $`W_k`$ through an isotopy $`W_{k,t}:(M_k,\omega _{k,t})`$ that preserves the cohomology class of the symplectic form (i.e., $`[\omega _{k,t}]=[\omega _k]`$ for all $`t[0,1]`$). We assume that each $`W_{k,t}`$ has isolated non-degenerate critical points. This allows us to deform the system of arcs $`\{\gamma _i\}`$ through a family $`\{\gamma _{i,t}\}`$ with end points at the critical values of $`W_{k,t}`$; for $`t=1`$ we obtain a system of arcs $`\{\gamma _i^{}\}`$ based at a regular value $`\lambda _{}^{}`$ of $`W_k^{}`$. By Moser’s theorem, there exists a continuous family of symplectomorphisms $`\varphi _t`$ from $`(M_k,\omega _{k,t})`$ to $`(M_k,\omega _k^{})`$, or rather, since these are non-compact manifolds, from open subsets of $`(M_k,\omega _{k,t})`$ to an open subset of $`(M_k,\omega _k^{})`$; however, after “enlarging” $`(M_k,\omega _k^{})`$ by adding to $`\omega _k^{}`$ the pullback of a suitable area form on $``$, which affects neither the symplectic structure on the fibers nor the parallel transport symplectomorphisms, we can ensure that $`\varphi _t`$ is defined over an arbitrarily large open subset of $`M_k`$, which is good enough for our purposes. Moreover, by a relative version of Moser’s argument, we can also ensure that $`\varphi _t`$ maps the reference fiber of $`W_{k,t}`$ to the reference fiber of $`W_k^{}`$, and in particular that $`\varphi =\varphi _0`$ maps $`\mathrm{\Sigma }_{}=W_k^1(\lambda _{})`$ to $`\mathrm{\Sigma }_{}^{}=W_{k}^{}{}_{}{}^{1}(\lambda _{}^{})`$.
We now claim that $`\varphi (L_i)\mathrm{\Sigma }_{}^{}`$ is Hamiltonian isotopic to the vanishing cycle $`L_i^{}`$ of $`W_k^{}`$ associated to the arc $`\gamma _i^{}`$. Indeed, by considering the images under $`\varphi _t`$ of the vanishing cycles $`L_{i,t}`$ associated to the arcs $`\gamma _{i,t}`$, we obtain a smooth isotopy from $`\varphi (L_i)`$ to $`L_i^{}`$ through Lagrangian submanifolds of $`\mathrm{\Sigma }_{}^{}`$. Moreover, the thimbles $`D_i^{}`$ and $`\varphi (D_i)`$ represent the same relative homotopy class (as can be seen by considering the images by $`\varphi _t`$ of the thimbles $`D_{i,t}`$ associated to $`\gamma _{i,t}`$), and both are Lagrangian with respect to $`\omega _k^{}`$, which again implies that $`\varphi (L_i)`$ and $`L_i^{}`$ are Hamiltonian isotopic. ∎
### 3.4. A basis of $`H_2(M_k)`$
The manifold $`M_k`$ is simply connected, and its second Betti number is equal to $`k+2`$. A $``$-basis of $`H_2(M_k)`$ is given by considering the homology class of the fiber of $`W_k`$, $`[\mathrm{\Sigma }_0]`$, and $`k+1`$ classes $`[\overline{C}],[\overline{C}_0],\mathrm{},[\overline{C}_{k1}]`$ constructed from the vanishing cycles $`L_i`$ and Lefschetz thimbles $`D_i`$ in the following manner.
By Lemma 3.1 we have $`[L_1]=[L_0]+[L_2]`$ in $`H_1(\mathrm{\Sigma }_0,)`$, so there exists a 2-chain $`C`$ in $`\mathrm{\Sigma }_0`$ such that $`C=L_0+L_1L_2`$. Then
$$\overline{C}=C+D_0D_1+D_2$$
is a $`2`$-cycle in $`M_k`$. Note that $`[\overline{C}]`$ is in fact the image of the generator of $`H_2((^{})^2,)`$ under the inclusion map (see the proof of Lemma 4.9 in ).
Similarly, for $`0i<k`$ we have $`3[L_{3+i}]=[L_2][L_0]`$ in $`H_1(\mathrm{\Sigma }_0,)`$, so there exists a 2-chain $`C_i`$ in $`\mathrm{\Sigma }_0`$ such that $`C_i=3L_{3+i}+L_0L_2`$, and we can consider the $`2`$-cycle
$$\overline{C}_i=C_i3D_{3+i}D_0+D_2$$
in $`M_k`$. We also introduce 2-chains $`\mathrm{\Delta }_{i,j}`$ ($`i,j\{0,\mathrm{},k1\}`$) in $`\mathrm{\Sigma }_0`$ such that $`\mathrm{\Delta }_{i,j}=L_{3+j}L_{3+i}`$, and the corresponding 2-cycles
$$\overline{\mathrm{\Delta }}_{i,j}=\mathrm{\Delta }_{i,j}+D_{3+i}D_{3+j}.$$
We can choose $`C_i`$ and $`\mathrm{\Delta }_{i,j}`$ in such a way that $`C_jC_i=3\mathrm{\Delta }_{i,j}`$ (and hence $`[\overline{C}_j][\overline{C}_i]=3[\overline{\mathrm{\Delta }}_{i,j}]`$ in $`H_2(M_k)`$).
To summarize the discussion, the vanishing cycles $`L_i`$ and the 2-chains $`C`$, $`C_i`$, $`\mathrm{\Delta }_{i,j}`$ are represented on Figures 89 (compare with Figure 2 in and with ).
## 4. Categories of vanishing cycles
### 4.1. Definition
As proposed by Kontsevich and Hori-Iqbal-Vafa , the category of A-branes associated to a Landau-Ginzburg model $`W:(M,\omega )`$ is a Fukaya-type category which contains not only compact Lagrangian submanifolds of $`M`$ but also certain non-compact Lagrangians whose ends fiber in a specific way above half-lines in $``$. In the case where the critical points of $`W`$ are isolated and non-degenerate, this category admits an exceptional collection whose objects are Lagrangian thimbles associated to the critical points. Following the formalism introduced by Seidel , we view it as the derived category of a finite directed $`A_{\mathrm{}}`$-category $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ associated to an ordered collection of arcs $`\{\gamma _i\}`$. We briefly recall the definition; the reader is referred to and to §3.1 of for details.
Consider a symplectic fibration $`W:(M,\omega )`$ with isolated non-degenerate critical points, and assume for simplicity that the critical values $`\lambda _0,\mathrm{},\lambda _r`$ of $`W`$ are distinct. Pick a regular value $`\lambda _{}`$ of $`W`$, and choose a collection of arcs $`\gamma _0,\mathrm{},\gamma _r`$ joining $`\lambda _{}`$ to the various critical values of $`W`$, intersecting each other only at $`\lambda _{}`$, and ordered in the clockwise direction around $`\lambda _{}`$. Consider the horizontal distribution defined by the symplectic form: by parallel transport along the arc $`\gamma _i`$, we obtain a Lagrangian thimble $`D_i`$ and a vanishing cycle $`L_i=D_i\mathrm{\Sigma }_{}`$ (where $`\mathrm{\Sigma }_{}=W^1(\lambda _{})`$). After a small perturbation we can always assume that the vanishing cycles $`L_i`$ intersect each other transversely inside $`\mathrm{\Sigma }_{}`$. The following definition is motivated by the observation that the intersection theory of the Lagrangian thimbles $`D_iM`$ is closely related to that of the vanishing cycles $`L_i`$ inside $`\mathrm{\Sigma }_{}`$ :
###### Definition 4.1 (Seidel).
The directed category of vanishing cycles $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ is an $`A_{\mathrm{}}`$-category (over a coefficient ring $`R`$) with objects $`L_0,\mathrm{},L_r`$ corresponding to the vanishing cycles (or more accurately to the thimbles); the morphisms between the objects are given by
$$\mathrm{Hom}(L_i,L_j)=\{\begin{array}{cc}CF^{}(L_i,L_j;R)=R^{|L_iL_j|}\hfill & \mathrm{if}i<j\hfill \\ Rid\hfill & \mathrm{if}i=j\hfill \\ 0\hfill & \mathrm{if}i>j;\hfill \end{array}$$
and the differential $`m_1`$, composition $`m_2`$ and higher order products $`m_k`$ are defined in terms of Lagrangian Floer homology inside $`\mathrm{\Sigma }_{}`$. More precisely,
$$m_k:\mathrm{Hom}(L_{i_0},L_{i_1})\mathrm{}\mathrm{Hom}(L_{i_{k1}},L_{i_k})\mathrm{Hom}(L_{i_0},L_{i_k})[2k]$$
is trivial when the inequality $`i_0<i_1<\mathrm{}<i_k`$ fails to hold (i.e. it is always zero in this case, except for $`m_2`$ where composition with an identity morphism is given by the obvious formula). When $`i_0<\mathrm{}<i_k`$, $`m_k`$ is defined by fixing a generic $`\omega `$-compatible almost-complex structure on $`\mathrm{\Sigma }_{}`$ and counting pseudo-holomorphic maps from a disk with $`k+1`$ cyclically ordered marked points on its boundary to $`\mathrm{\Sigma }_{}`$, mapping the marked points to the given intersection points between vanishing cycles, and the portions of boundary between them to $`L_{i_0},\mathrm{},L_{i_k}`$ respectively.
This definition calls for several clarifications. First of all, in our case $`\mathrm{\Sigma }_{}`$ is a smooth elliptic curve and the vanishing cycles are homotopically non-trivial closed loops, we have $`\pi _2(\mathrm{\Sigma }_{})=0`$ and $`\pi _2(\mathrm{\Sigma }_{},L_i)=0`$; hence, we need not be concerned by bubbling issues in the definition of the Floer differential and products. In fact, the pseudo-holomorphic disks in $`\mathrm{\Sigma }_{}`$ that we have to consider are nothing but immersed polygonal regions bounded by the vanishing cycles, satisfying a local convexity condition at each corner point.
Also, the Maslov class vanishes identically, so we have a well-defined $``$-grading by Maslov index on the Floer complexes $`CF^{}(L_i,L_j;R)`$ once we choose graded Lagrangian lifts of the vanishing cycles. Since in our case $`c_1(\mathrm{\Sigma }_{})=0`$, we can do this by considering a nowhere vanishing 1-form $`\mathrm{\Omega }\mathrm{\Omega }^1(\mathrm{\Sigma }_{},)`$ and choosing a real lift of the phase function $`\varphi _i=\mathrm{arg}(\mathrm{\Omega }_{|L_i}):L_iS^1`$ for each vanishing cycle. The degree of a given intersection point $`pL_iL_j`$ is then determined by the difference between the phases of $`L_i`$ and $`L_j`$ at $`p.`$
Our next remark is that the pseudo-holomorphic disks appearing in Definition 4.1 are counted with appropriate weights, and with signs determined by choices of orientations of the relevant moduli spaces. The orientation is determined by the choice of a spin structure for each vanishing cycle $`L_i`$; in our case this spin structure must extend to the thimble, so it is necessarily the non-trivial one. In the one-dimensional case there is a convenient recipe for determining the correct sign factors, due to Seidel . As will be clear from the discussion in §4.2 below, we will only be interested in the specific case where all morphisms have even degree and all spin structures are non-trivial. The sign rule can then be summarized as follows: pick a marked point on each $`L_i`$, distinct from the intersections with the other vanishing cycles; then the sign associated to a pseudo-holomorphic map $`u:(D^2,D^2)(\mathrm{\Sigma }_{},L_i)`$ is $`(1)^{\nu (u)}`$, where $`\nu (u)`$ is the number of marked points that the boundary of $`u`$ passes through (, see also §4.6 of ).
Finally, the weight attributed to each pseudo-holomorphic map $`u`$ keeps track of its relative homology class, which makes it possible to avoid convergence problems. The usual approach favored by mathematicians is to work over a Novikov ring, which keeps track of the relative homology class by introducing suitable formal variables. To remain closer to the physics, we use $``$ as our coefficient ring, and assign weights according to the symplectic areas; this is in fact equivalent to working over the Novikov ring and specializing at the cohomology class of the symplectic form.
The weight formula is simplest when there is no B-field; in that case, we consider untwisted Floer theory, since any flat unitary bundle over the thimble $`D_i`$ is trivial and hence restricts to $`L_i`$ as the trivial bundle. We then count each map $`u:(D^2,D^2)(\mathrm{\Sigma }_{},L_i)`$ with a coefficient $`(1)^{\nu (u)}\mathrm{exp}(2\pi _{D^2}u^{}\omega )`$. (The normalization factor $`2\pi `$ is purely a matter of conventions, and is sometimes omitted in the literature; here we include it for convenience). Hence, given two intersection points $`pL_iL_j`$, $`qL_jL_k`$ ($`i<j<k`$), we have by definition
$$m_2(p,q)=\underset{\begin{array}{c}rL_iL_k\\ \mathrm{deg}r=\mathrm{deg}p+\mathrm{deg}q\end{array}}{}\left(\underset{[u](p,q,r)}{}(1)^{\nu (u)}\mathrm{exp}(2\pi _{D^2}u^{}\omega )\right)r$$
where $`(p,q,r)`$ is the moduli space of pseudo-holomorphic maps $`u`$ from the unit disk to $`\mathrm{\Sigma }_{}`$ (equipped with a generic $`\omega `$-compatible almost-complex structure) such that $`u(1)=p`$, $`u(\mathrm{j})=q`$, $`u(\mathrm{j}^2)=r`$ (where $`\mathrm{j}=\mathrm{exp}(\frac{2i\pi }{3})`$), and mapping the portions of unit circle $`[1,\mathrm{j}]`$, $`[\mathrm{j},\mathrm{j}^2]`$, $`[\mathrm{j}^2,1]`$ to $`L_i`$, $`L_j`$ and $`L_k`$ respectively. The other products are defined similarly. (Observe that Seidel’s definition does not involve any weights; this is because he only considers exact Lagrangian submanifolds in exact symplectic manifolds, in which case the symplectic areas are entirely determined by the primitives of the Liouville form and can be eliminated by considering suitably rescaled bases of the Floer complexes.)
In presence of a B-field, the weights are modified by the fact that we now consider twisted Floer homology. Indeed, each thimble $`D_i`$ now comes equipped with a trivial complex line bundle $`E_i=\underset{¯}{}`$ and a connection $`_i`$ with curvature $`2\pi iB`$, so its boundary $`L_i`$ is equipped with the restricted bundle and the restricted connection, whose holonomy is $`\mathrm{hol}__i(L_i)=\mathrm{exp}(2\pi i_{D_i}B)`$ by Stokes’ theorem. Since this property characterizes the connection $`_i`$ uniquely up to gauge, we can drop the line bundle and the connection from the notation when considering the objects $`(L_i,E_i,_i)`$ of $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$. Nonetheless, the holonomy of $`_i`$ modifies the weights attributed to the pseudo-holomorphic disks in the definition of the twisted Floer differentials and compositions. Namely, the weight attributed to a given pseudo-holomorphic map $`u:(D^2,D^2)(\mathrm{\Sigma }_{},L_i)`$ is modified by a factor corresponding to the holonomy along its boundary, and becomes
$$(1)^{\nu (u)}\mathrm{hol}(u(D^2))\mathrm{exp}(2\pi i_{D^2}u^{}(B+i\omega )).$$
More precisely, we fix trivializations of the line bundles $`E_i`$, so that for each intersection point $`pL_iL_j`$ we have a preferred isomorphism between the fibers $`(E_i)_{|p}`$ and $`(E_j)_{|p}`$; then it becomes possible to define the holonomy along the closed loop $`u(D^2)`$ using the parallel transport induced by $`_i`$ from one “corner” of $`u`$ to the next one, and the chosen isomorphism at each corner.
Although the category $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ depends on the chosen ordered collection of arcs $`\{\gamma _i\}`$, Seidel has obtained the following result (for the exact case, but the proof extends to our situation):
###### Theorem 4.2 (Seidel).
If the ordered collection $`\{\gamma _i\}`$ is replaced by another one $`\{\gamma _i^{}\}`$, then the categories $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ and $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i^{}\})`$ differ by a sequence of mutations.
Hence, the category naturally associated to the fibration $`W`$ is not the finite $`A_{\mathrm{}}`$-category defined above, but rather a (bounded) derived category, obtained from $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ by considering twisted complexes of formal direct sums of Lagrangian vanishing cycles, and adding idempotent splittings and formal inverses of quasi-isomorphisms (see and §5 of ). If two categories differ by mutations, then their derived categories are equivalent; hence the derived category $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W))`$ depends only on the symplectic topology of $`W`$ and not on the choice of an ordered system of arcs .
For the examples we consider, the $`A_{\mathrm{}}`$-category $`\mathrm{Lag}_{\mathrm{vc}}(W,\{\gamma _i\})`$ will in fact be an honest category (see below); the bounded derived category $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W))`$ is then by definition the bounded derived category of finite rank modules over the algebra associated to this category.
### 4.2. Objects and morphisms
We now determine the categories $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$ associated to the Landau-Ginzburg models $`(M_k,W_k)`$ mirror to Del Pezzo surfaces and the systems of arcs $`\{\gamma _i\}`$ introduced in §3.3. We start with the objects and morphisms.
Recall that $`W_k`$ has $`k+3`$ isolated critical points, giving rise to $`k+3`$ vanishing cycles $`L_0,\mathrm{},L_{k+2}`$ in the reference fiber $`\mathrm{\Sigma }_0W_k^1(0)`$. The homology classes of these vanishing cycles have been determined in §3 and are given by Lemma 3.1; these determine the vanishing cycles up to Lagrangian isotopy.
The derived category of vanishing cycles is not affected if we modify some of the vanishing cycles by Hamiltonian isotopies (more precisely, a Hamiltonian isotopy induces a chain map on the Floer complexes, which yields a quasi-isomorphism between the finite $`A_{\mathrm{}}`$-categories of vanishing cycles). Hence, equipping the elliptic curve $`\mathrm{\Sigma }_0`$ with a compatible flat metric, we can identify $`\mathrm{\Sigma }_0`$ with the quotient of $``$ by a lattice, and represent the vanishing cycles $`L_i`$ by closed geodesics parallel to those represented in Figure 8.
Assume that the cohomology class of the symplectic form $`\omega _k`$ on $`M_k`$ is generic (or more precisely, with the notations of §3.4, that $`[\omega _k][\overline{\mathrm{\Delta }}_{i,j}]`$ is never an integer multiple of $`[\omega _k][\mathrm{\Sigma }_0]`$). Then the geodesics $`L_i`$ are all distinct, and their intersections are as pictured in Figure 8, so we have:
###### Lemma 4.3.
The geometric intersection numbers between the vanishing cycles are:
$``$ $`|L_0L_1|=|L_0L_2|=|L_1L_2|=3`$;
$``$ for $`0i<k`$, $`|L_0L_{3+i}|=|L_2L_{3+i}|=1`$ and $`|L_1L_{3+i}|=2`$;
$``$ for $`0i<j<k`$, $`|L_{3+i}L_{3+j}|=0`$ as soon as $`[\omega _k][\mathrm{\Sigma }_0][\omega _k][\overline{\mathrm{\Delta }}_{i,j}]`$.
In the rest of this section, unless otherwise specified we always assume that the vanishing cycles $`L_i`$ are represented by distinct closed geodesics.
As in , we denote by $`x_0,y_0,z_0`$ (resp. $`x_1,y_1,z_1`$ and $`\overline{x},\overline{y},\overline{z}`$) the generators of $`\mathrm{Hom}(L_0,L_1)`$ (resp. $`\mathrm{Hom}(L_1,L_2)`$ and $`\mathrm{Hom}(L_0,L_2)`$) corresponding to the intersection points represented in Figure 8. Moreover, we denote by $`a_i`$ (resp. $`b_i,b_i^{}`$ and $`c_i`$) the generators of $`\mathrm{Hom}(L_0,L_{3+i})`$ (resp. $`\mathrm{Hom}(L_1,L_{3+i})`$ and $`\mathrm{Hom}(L_2,L_{3+i})`$) corresponding to the intersection points between these vanishing cycles (see Figure 9).
###### Lemma 4.4.
For suitable choices of graded lifts of the vanishing cycles, all the morphisms in $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$ have degree $`0`$.
###### Proof.
Equip $`\mathrm{\Sigma }_0`$ with a compatible flat metric and with a constant holomorphic 1-form $`\mathrm{\Omega }`$. Taking geodesic representatives of the vanishing cycles, the phase functions $`\varphi _i=\mathrm{arg}(\mathrm{\Omega }_{|L_i}):L_i/2\pi `$ are constant, and we can normalize $`\mathrm{\Omega }`$ so that it takes real negative values on the oriented tangent space to $`L_0`$, i.e. $`\varphi _0=\pi `$. Then it is possible to choose real lifts $`\stackrel{~}{\varphi }_i`$ of the phases in such a way that $`\pi =\stackrel{~}{\varphi }_0>\stackrel{~}{\varphi }_1>\stackrel{~}{\varphi }_2>\stackrel{~}{\varphi }_3=\mathrm{}=\stackrel{~}{\varphi }_{k+2}>0`$ (see Figure 8 and recall the orientations chosen in Lemma 3.1). In the 1-dimensional case, the relationship between Maslov index and phase is very simple: given a transverse intersection point $`p`$ between two graded Lagrangians $`L,L^{}\mathrm{\Sigma }_0`$, the Maslov index of $`pCF^{}(L,L^{})`$ is equal to the smallest integer greater than $`\frac{1}{\pi }(\stackrel{~}{\varphi }_L^{}(p)\stackrel{~}{\varphi }_L(p))`$. Since we only consider the Floer complexes $`CF^{}(L_i,L_j)`$ for $`i<j`$, which implies that $`\stackrel{~}{\varphi }_j\stackrel{~}{\varphi }_i(\pi ,0)`$ at every intersection point, for these choices of graded Lagrangian lifts of the vanishing cycles all morphisms in $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$ have degree 0. ∎
Since each product $`m_j`$ shifts degree by $`2j`$, it follows immediately that the $`A_{\mathrm{}}`$-category $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$ is actually an honest category:
###### Corollary 4.5.
$`m_j=0`$ for all $`j2`$.
Hence, the final step of the argument is a careful study of the various immersed triangular regions bounded by the vanishing cycles in $`\mathrm{\Sigma }_0`$ and their contributions to $`m_2`$.
### 4.3. Compositions
As before we assume that the Lagrangian vanishing cycles are realized by distinct closed geodesics in the flat torus $`\mathrm{\Sigma }_0`$, and we determine the contributions to $`m_2`$ of the various immersed triangular regions in $`(\mathrm{\Sigma }_0,L_i)`$. We use the notations introduced in §4.2 for the intersection points, and those introduced in §3.4 for various 2-chains in $`\mathrm{\Sigma }_0`$ and the corresponding 2-cycles in $`M_k`$. We also introduce the following notations:
###### Definition 4.6.
Let $`q_C=\mathrm{exp}(2\pi i[B+i\omega ][\overline{C}])`$ and $`q_F=\mathrm{exp}(2\pi i[B+i\omega ][\mathrm{\Sigma }_0])`$, and define
$$\zeta _+=\underset{n}{}(1)^nq_C^nq_F^{n(3n+1)/2},\zeta _{}=\underset{n}{}(1)^nq_C^nq_F^{n(3n1)/2},\zeta _0=\underset{n}{}(1)^nq_C^nq_F^{3n(n1)/2}.$$
Since $`\omega `$ is a symplectic form on $`\mathrm{\Sigma }_0`$, we have $`|q_F|=\mathrm{exp}(2\pi [\omega ][\mathrm{\Sigma }_0])<1`$, which ensures the convergence of the series $`\zeta _+`$, $`\zeta _{}`$ and $`\zeta _0`$.
###### Proposition 4.7.
There exist constants $`\alpha _{xy},\alpha _{yx},\alpha _{yz},\alpha _{zy},\alpha _{zx},\alpha _{xz}`$ such that
$`m_2(x_0,y_1)=\alpha _{xy}\overline{z},`$ $`m_2(y_0,x_1)=\alpha _{yx}\overline{z},`$
$`m_2(y_0,z_1)=\alpha _{yz}\overline{x},`$ $`m_2(z_0,y_1)=\alpha _{zy}\overline{x},`$
$`m_2(z_0,x_1)=\alpha _{zx}\overline{y},`$ $`m_2(x_0,z_1)=\alpha _{xz}\overline{y},`$
and these constants satisfy the relation
(4.1)
$$\frac{\alpha _{xy}\alpha _{yz}\alpha _{zx}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=q_C\left(\frac{_n(1)^nq_C^nq_F^{n(3n+1)/2}}{_n(1)^nq_C^nq_F^{n(3n1)/2}}\right)^3=q_C\left(\frac{\zeta _+}{\zeta _{}}\right)^3.$$
###### Remark 4.8.
The quantity appearing in the right-hand side of (4.1) can be understood in terms of certain theta functions; see §4.5 for details.
Before giving the proof, we make an observation which will be useful throughout this section. The geodesics $`L_i`$ are not necessarily those pictured in Figure 8, but they are parallel to them. So we can deform (in a non-Hamiltonian manner) the configuration of vanishing cycles to that of Figure 8, and all intersection points and relative 2-cycles in $`(\mathrm{\Sigma }_0,L_i)`$ can be followed through the deformation. Hence, immersed triangular regions in $`(\mathrm{\Sigma }_0,L_i)`$ are in one to one correspondence with those in the configuration of Figure 8 (but of course the deformation does not preserve areas). Moreover, we can choose the deformation in such a way that the relative 2-cycles $`C`$ and $`C_i`$ in $`\mathrm{\Sigma }_0`$ deform to those represented on Figures 89 (rather than to 2-cycles which differ by a multiple of the fundamental cycle of $`\mathrm{\Sigma }_0`$).
###### Proof of Proposition 4.7.
The composition $`m_2(x_0,y_1)`$ is the sum of an infinite series of contributions, corresponding to all immersed triangular regions in $`\mathrm{\Sigma }_0`$ with corners at the intersection points $`x_0`$, $`y_1`$, and one of the points in $`L_0L_2`$. By the above remark we can enumerate these regions by looking at Figure 8. Considering the side which lies on $`L_1`$, it is then easy to see that for every homotopy class of arc joining $`x_0`$ to $`y_1`$ inside $`L_1`$ there is a unique such immersed triangular region, and the third vertex is always $`\overline{z}`$.
These various regions can be labelled by integers $`n`$ in such a way that, denoting by $`T_{xy,n}`$ the corresponding 2-chains in $`\mathrm{\Sigma }_0`$, we have $`T_{xy,n}T_{xy,n^{}}=(nn^{})(L_0+L_1L_2)`$ for all $`n,n^{}`$. We can choose the integer labels in such a way that, after deforming to the configuration in Figure 8, $`T_{xy,0}`$ becomes the smallest triangle with vertices $`x_0,y_1,\overline{z}`$. (So, in Figure 8, $`T_{xy,1}`$ is the immersed region bounded by the portions of $`L_0L_1L_2`$ which do not belong to $`T_{xy,0}`$; and all the other $`T_{xy,n}`$ have edges which wrap more than once around the vanishing cycles).
By comparing $`T_{xy,n}`$ and $`T_{xy,0}`$, it is clear that the 2-chain represented by $`T_{xy,n}`$ can be expressed in the form $`T_{xy,n}=T_{xy,0}+nC+\varphi (n)\mathrm{\Sigma }_0`$ for some $`\varphi (n)`$. Moreover, by looking at Figure 8 one easily checks that $`\varphi (n)=\frac{1}{2}n(3n+1)`$. (So e.g. $`T_{xy,1}=T_{xy,0}C+\mathrm{\Sigma }_0`$, and $`T_{xy,1}=T_{xy,0}+C+2\mathrm{\Sigma }_0`$). Let $`\psi _{xy}`$ be the coefficient of the contribution of $`T_{xy,0}`$ to $`m_2(x_0,y_1)`$. Then, by comparing the symplectic areas and boundary holonomies for $`T_{xy,n}`$ and $`T_{xy,0}`$, one easily checks that the contribution of $`T_{xy,n}`$ is equal to
$$(1)^n\mathrm{exp}\left(2\pi i[B+i\omega ]\left(n[\overline{C}]+\frac{n(3n+1)}{2}[\mathrm{\Sigma }_0]\right)\right)\psi _{xy}=(1)^nq_C^nq_F^{n(3n+1)/2}\psi _{xy}.$$
In this expression the sign factor $`(1)^n`$ is due to the non-triviality of the spin structures (observe that $`C=L_0+L_1L_2`$ passes once through each of the three marked points on $`L_0,L_1,L_2`$); the total holonomy of the flat connections $`_i`$ along $`T_{xy,n}T_{xy,0}=nC`$ is $`\mathrm{exp}(2\pi in_{D_0D_1+D_2}B)`$ by Stokes’ theorem; and the integral of $`B+i\omega `$ over $`T_{xy,n}`$ differs from that over $`T_{xy,0}`$ by the amount $`n_C(B+i\omega )+\frac{1}{2}n(3n+1)[B+i\omega ][\mathrm{\Sigma }_0]`$.
Summing over $`n`$, and using the notation introduced in Definition 4.6, we obtain
$$\alpha _{xy}=\zeta _+\psi _{xy}.$$
The calculations of $`m_2(y_0,z_1)`$ and $`m_2(z_0,x_1)`$ are exactly identical, and lead to similar expressions. Namely, denote by $`\psi _{yz}`$ (resp. $`\psi _{zx}`$) the contribution of the triangular region $`T_{yz,0}`$ (resp. $`T_{zx,0}`$) which, after deforming to the configuration in Figure 8, corresponds to the smallest triangle with vertices $`y_0,z_1,\overline{x}`$ (resp. $`z_0,x_1,\overline{y}`$). Then one easily checks by the same argument as above that $`\alpha _{yz}=\zeta _+\psi _{yz}`$ and $`\alpha _{zx}=\zeta _+\psi _{zx}`$.
Next we consider the composition $`m_2(y_0,x_1)`$, which is again the sum of an infinite series of contributions from triangular regions $`T_{yx,n}`$, $`n`$, which all have vertices $`y_0,x_1,\overline{z}`$. We can choose the labels in such a way that, after deforming to the configuration in Figure 8, $`T_{yx,0}`$ becomes the smallest such triangle, and $`T_{yx,n}=T_{yx,0}+nC+\frac{1}{2}n(3n1)\mathrm{\Sigma }_0`$. Denoting by $`\psi _{yx}`$ the coefficient associated to $`T_{yx,0}`$, it is easy to check by the same argument as above that the contribution of $`T_{yx,n}`$ is equal to $`(1)^nq_C^nq_F^{n(3n1)/2}\psi _{yx}`$, so that
$$\alpha _{yx}=\zeta _{}\psi _{yx}.$$
Similarly, with the obvious notations we have $`\alpha _{zy}=\zeta _{}\psi _{zy}`$ and $`\alpha _{xz}=\zeta _{}\psi _{xz}`$. Finally, observe that
$$\frac{\psi _{xy}\psi _{yz}\psi _{zx}}{\psi _{yx}\psi _{zy}\psi _{xz}}=q_C.$$
Indeed, $`T_{xy,0}+T_{yz,0}+T_{zx,0}T_{yx,0}T_{zy,0}T_{xz,0}=C`$ (cf. Figure 8). Therefore, comparing the weights associated to these various triangles, the weighting by area gives a factor of $`\mathrm{exp}(2\pi i_CB+i\omega )`$, while the holonomy along the boundary $`C=L_0+L_1L_2`$ is equal to $`\mathrm{exp}(2\pi i_{D_0D_1+D_2}B)`$, and finally the minus sign is due to the orientation conventions, since $`C`$ passes once through each of the three marked points on the vanishing cycles. Hence
$$\frac{\alpha _{xy}\alpha _{yz}\alpha _{zx}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=\frac{\psi _{xy}\psi _{yz}\psi _{zx}\zeta _+^3}{\psi _{yx}\psi _{zy}\psi _{xz}\zeta _{}^3}=q_C\left(\frac{\zeta _+}{\zeta _{}}\right)^3.$$
###### Remark 4.9.
If $`[\omega +iB][\overline{C}]=0`$, then $`q_C=1`$ and the ratio between $`\alpha _{xy}\alpha _{yz}\alpha _{zx}`$ and $`\alpha _{yx}\alpha _{zy}\alpha _{xz}`$ becomes equal to $`1`$ irrespective of the value of $`q_F`$; this corresponds to a classical (commutative) Del Pezzo surface.
Moreover, in the limit where $`[\omega ][\mathrm{\Sigma }_0]\mathrm{}`$, we have $`q_F=0`$ and the ratio becomes $`q_C`$, which corresponds to the toric case studied in .
###### Proposition 4.10.
There exist constants $`\alpha _{xx},\alpha _{yy},\alpha _{zz}`$ such that
$$m_2(x_0,x_1)=\alpha _{xx}\overline{x},m_2(y_0,y_1)=\alpha _{yy}\overline{y},m_2(z_0,z_1)=\alpha _{zz}\overline{z},$$
and these constants satisfy the relation
$$\frac{\alpha _{xx}\alpha _{yy}\alpha _{zz}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=\frac{q_F}{q_C}\left(\frac{_n(1)^nq_C^nq_F^{3n(n1)/2}}{_n(1)^nq_C^nq_F^{n(3n1)/2}}\right)^3=\frac{q_F}{q_C}\left(\frac{\zeta _0}{\zeta _{}}\right)^3.$$
###### Proof.
The argument is similar to the proof of Proposition 4.7. The immersed triangular regions which contribute to $`m_2(x_0,x_1)`$ all have vertices $`\overline{x}`$ as their third vertex, and can be indexed by integers $`n`$ in a manner such that $`T_{xx,n}T_{xx,n^{}}=(nn^{})C`$ for all $`n,n^{}`$. We can choose the integer labels in such a way that, after deforming to the standard configuration, $`T_{xx,0}`$ and $`T_{xx,1}=T_{xx,0}+C`$ are the two embedded triangles with vertices $`x_0,x_1,\overline{x}`$ visible on Figure 8. It is then easy to check that $`T_{xx,n}=T_{xy,0}+nC+\frac{3}{2}n(n1)\mathrm{\Sigma }_0`$. Hence, denoting by $`\psi _{xx}`$ the coefficient associated to $`T_{xx,0}`$, we have
$$\alpha _{xx}=\zeta _0\psi _{xx},$$
by the same argument as in previous calculations. Similarly, with the obvious notations, we have $`\alpha _{yy}=\zeta _0\psi _{yy}`$ and $`\alpha _{zz}=\zeta _0\psi _{zz}`$. Moreover, $`T_{xx,0}+T_{yy,0}+T_{zz,0}T_{yx,0}T_{zy,0}T_{xz,0}=\mathrm{\Sigma }_0C`$, which implies (by the same argument as above) that
$$\frac{\psi _{xx}\psi _{yy}\psi _{zz}}{\psi _{yx}\psi _{zy}\psi _{xz}}=\frac{q_F}{q_C}.$$
Therefore
$$\frac{\alpha _{xx}\alpha _{yy}\alpha _{zz}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=\frac{\psi _{xx}\psi _{yy}\psi _{zz}\zeta _0^3}{\psi _{yx}\psi _{zy}\psi _{xz}\zeta _{}^3}=\frac{q_F}{q_C}\left(\frac{\zeta _0}{\zeta _{}}\right)^3.$$
When $`q_F=0`$ (in particular in the toric case) we have $`\alpha _{xx}\alpha _{yy}\alpha _{zz}=0`$, as in . The same conclusion also holds when $`q_C=1`$ (the commutative case). In fact, when $`q_C=1`$ each of the constants $`\alpha _{xx},\alpha _{yy},\alpha _{zz}`$ is zero, since in that case we have $`\zeta _0=0`$ (because the terms corresponding to $`n`$ and $`1n`$ in the series defining $`\zeta _0`$ exactly cancel each other).
###### Definition 4.11.
Let $`q_i=\mathrm{exp}(2\pi i[B+i\omega ][\overline{C}_i])`$, and define
$$\zeta _{i,+}=\underset{n}{}(1)^nq_i^nq_F^{n(3n+1)/2},\zeta _{i,}=\underset{n}{}(1)^nq_i^nq_F^{n(3n1)/2},\zeta _{i,0}=\underset{n}{}(1)^nq_i^nq_F^{3n(n1)/2}.$$
###### Proposition 4.12.
There exist constants $`\beta _{\overline{x},i},\beta _{\overline{y},i},\beta _{\overline{z},i}`$ such that
$$m_2(\overline{x},c_i)=\beta _{\overline{x},i}a_i,m_2(\overline{y},c_i)=\beta _{\overline{y},i}a_i,m_2(\overline{z},c_i)=\beta _{\overline{z},i}a_i,$$
and these constants satisfy the relations
$$\frac{\beta _{\overline{z},i}^2\alpha _{xy}\alpha _{zz}}{\beta _{\overline{x},i}\beta _{\overline{y},i}\alpha _{zy}\alpha _{xz}}=\left(\frac{\zeta _{i,}}{\zeta _{}}\right)^2\frac{\zeta _+\zeta _0}{\zeta _{i,+}\zeta _{i,0}},$$
$$\frac{\beta _{\overline{x},i}^2\alpha _{yz}\alpha _{xx}}{\beta _{\overline{y},i}\beta _{\overline{z},i}\alpha _{xz}\alpha _{yx}}=q_i\left(\frac{\zeta _{i,+}}{\zeta _{}}\right)^2\frac{\zeta _+\zeta _0}{\zeta _{i,0}\zeta _{i,}},\mathrm{𝑎𝑛𝑑}\frac{\beta _{\overline{y},i}^2\alpha _{zx}\alpha _{yy}}{\beta _{\overline{z},i}\beta _{\overline{x},i}\alpha _{yx}\alpha _{zy}}=\frac{q_F}{q_i}\left(\frac{\zeta _{i,0}}{\zeta _{}}\right)^2\frac{\zeta _+\zeta _0}{\zeta _{i,}\zeta _{i,+}},$$
where $`\zeta _+,\zeta _{},\zeta _0,\zeta _{i,+},\zeta _{i,},\zeta _{i,0}`$, $`q_i`$ and $`q_F`$ are as in Definitions 4.6 and 4.11.
###### Proof.
As before, the constants $`\beta _{\overline{x},i},\beta _{\overline{y},i},\beta _{\overline{z},i}`$ are the sums of infinite series corresponding to all immersed triangular regions with vertices at $`a_i`$, $`c_i`$, and one of $`\overline{x},\overline{y},\overline{z}`$. For example the coefficient $`\beta _{\overline{z},i}`$ associated to composition $`m_2(\overline{x},c_i)`$ is the sum of an infinite series of contributions associated to triangular regions $`T_{\overline{z},i,n}`$, $`n`$. The integer labels can be chosen so that $`T_{\overline{z},i,n}T_{\overline{z},i,n^{}}=(nn^{})C_i`$ and, after deforming to the configuration in Figure 9, $`T_{\overline{z},i,0}`$ becomes the smallest triangle with vertices $`\overline{z},a_i,c_i`$ (i.e., the triangle which appears with coefficient $`2`$ in the 2-chain $`C_i`$). Then one easily checks that $`T_{\overline{z},i,n}=T_{\overline{z},i,0}+nC_i+\frac{1}{2}n(3n1)\mathrm{\Sigma }_0`$. Therefore, denoting by $`\psi _{\overline{z},i}`$ the coefficient associated to $`T_{\overline{z},i,0}`$, the same argument as in the previous calculations yields the formula
$$\beta _{\overline{z},i}=\zeta _{i,}\psi _{\overline{z},i}.$$
Similarly, denote by $`T_{\overline{x},i,n}`$, $`n`$, the immersed triangles contributing to $`m_2(\overline{x},c_i)`$, in such a way that $`T_{\overline{x},i,n}T_{\overline{x},i,n^{}}=(nn^{})C_i`$, and $`T_{\overline{x},i,0}`$ corresponds to the smallest triangle with vertices $`\overline{x},a_i,c_i`$ in Figure 9 (i.e. the triangle which appears with coefficient $`+2`$ in the 2-chain $`C_i`$). Then $`T_{\overline{x},i,n}=T_{\overline{x},i,0}+nC_i+\frac{1}{2}n(3n+1)\mathrm{\Sigma }_0`$. Therefore, denoting by $`\psi _{\overline{x},i}`$ the contribution of $`T_{\overline{x},i,0}`$, we have $`\beta _{\overline{x},i}=\zeta _{i,+}\psi _{\overline{x},i}.`$
Finally, labelling the triangles with vertices $`\overline{y},a_i,c_i`$ by integers in such a way that $`T_{\overline{y},i,0}`$ and $`T_{\overline{y},i,1}=T_{\overline{y},i,0}+C_i`$ correspond to the negative and positive parts of $`C_i`$ respectively, it is easy to check that $`T_{\overline{y},i,n}=T_{\overline{y},i,0}+nC_i+\frac{3}{2}n(n1)\mathrm{\Sigma }_0`$, so denoting by $`\psi _{\overline{y},i}`$ the contribution of $`T_{\overline{y},i,0}`$ we have $`\beta _{\overline{y},i}=\zeta _{i,0}\psi _{\overline{y},i}.`$ It follows that
$$\frac{\beta _{\overline{z},i}^2\alpha _{xy}\alpha _{zz}}{\beta _{\overline{x},i}\beta _{\overline{y},i}\alpha _{zy}\alpha _{xz}}=\frac{\psi _{\overline{z},i}^2\psi _{xy}\psi _{zz}}{\psi _{\overline{x},i}\psi _{\overline{y},i}\psi _{zy}\psi _{xz}}\frac{\zeta _{i,}^2\zeta _+\zeta _0}{\zeta _{i,+}\zeta _{i,0}\zeta _{}^2}.$$
Moreover, the 2-chains $`2T_{\overline{z},i,0}+T_{xy,0}+T_{zz,0}`$ and $`T_{\overline{x},i,0}+T_{\overline{y},i,0}+T_{zy,0}+T_{xz,0}`$ are equal, which implies that $`\psi _{\overline{z},i}^2\psi _{xy}\psi _{zz}=\psi _{\overline{x},i}\psi _{\overline{y},i}\psi _{zy}\psi _{xz}`$ and completes the proof of the first identity.
The arguments are the same for
$$\frac{\beta _{\overline{x},i}^2\alpha _{yz}\alpha _{xx}}{\beta _{\overline{y},i}\beta _{\overline{z},i}\alpha _{xz}\alpha _{yx}}=\frac{\psi _{\overline{x},i}^2\psi _{yz}\psi _{xx}}{\psi _{\overline{y},i}\psi _{\overline{z},i}\psi _{xz}\psi _{yx}}\frac{\zeta _{i,+}^2\zeta _+\zeta _0}{\zeta _{i,0}\zeta _{i,}\zeta _{}^2},$$
observing that $`2T_{\overline{x},i,0}+T_{yz,0}+T_{xx,0}T_{\overline{y},i,0}T_{\overline{z},i,0}T_{xz,0}T_{yx,0}=C_i`$ (for which the corresponding weight is $`q_i`$), and for
$$\frac{\beta _{\overline{y},i}^2\alpha _{zx}\alpha _{yy}}{\beta _{\overline{z},i}\beta _{\overline{x},i}\alpha _{yx}\alpha _{zy}}=\frac{\psi _{\overline{y},i}^2\psi _{zx}\psi _{yy}}{\psi _{\overline{z},i}\psi _{\overline{x},i}\psi _{yx}\psi _{zy}}\frac{\zeta _{i,0}^2\zeta _+\zeta _0}{\zeta _{i,}\zeta _{i,+}\zeta _{}^2},$$
observing that $`2T_{\overline{y},i,0}+T_{zx,0}+T_{yy,0}T_{\overline{z},i,0}T_{\overline{x},i,0}T_{yx,0}T_{zy,0}=\mathrm{\Sigma }_0C_i`$ (for which the corresponding weight is $`q_F/q_i`$). ∎
###### Corollary 4.13.
The constants $`\beta _{\overline{x},i},\beta _{\overline{y},i},\beta _{\overline{z},i}`$ satisfy the relations: $`{\displaystyle \frac{\beta _{\overline{z},i}^3}{\beta _{\overline{x},i}^3}}{\displaystyle \frac{\alpha _{xy}\alpha _{yx}\alpha _{zz}}{\alpha _{yz}\alpha _{zy}\alpha _{xx}}}={\displaystyle \frac{1}{q_i}}\left({\displaystyle \frac{\zeta _{i,}}{\zeta _{i,+}}}\right)^3`$, $`{\displaystyle \frac{\beta _{\overline{x},i}^3}{\beta _{\overline{y},i}^3}}{\displaystyle \frac{\alpha _{yz}\alpha _{zy}\alpha _{xx}}{\alpha _{zx}\alpha _{xz}\alpha _{yy}}}={\displaystyle \frac{q_i^2}{q_F}}\left({\displaystyle \frac{\zeta _{i,+}}{\zeta _{i,0}}}\right)^3`$, and $`{\displaystyle \frac{\beta _{\overline{y},i}^3}{\beta _{\overline{z},i}^3}}{\displaystyle \frac{\alpha _{zx}\alpha _{xz}\alpha _{yy}}{\alpha _{xy}\alpha _{yx}\alpha _{zz}}}={\displaystyle \frac{q_F}{q_i}}\left({\displaystyle \frac{\zeta _{i,0}}{\zeta _{i,}}}\right)^3.`$
###### Proposition 4.14.
For all $`0i,j<k`$ we have the identities
$$\frac{\beta _{\overline{y},i}\beta _{\overline{z},j}}{\beta _{\overline{y},j}\beta _{\overline{z},i}}=\stackrel{~}{q}_{i,j}\frac{\zeta _{i,0}\zeta _{j,}}{\zeta _{j,0}\zeta _{i,}},\frac{\beta _{\overline{z},i}\beta _{\overline{x},j}}{\beta _{\overline{z},j}\beta _{\overline{x},i}}=\stackrel{~}{q}_{i,j}\frac{\zeta _{i,}\zeta _{j,+}}{\zeta _{j,}\zeta _{i,+}},\mathrm{𝑎𝑛𝑑}\frac{\beta _{\overline{x},i}\beta _{\overline{y},j}}{\beta _{\overline{x},j}\beta _{\overline{y},i}}=\stackrel{~}{q}_{i,j}^2\frac{\zeta _{i,+}\zeta _{j,0}}{\zeta _{j,+}\zeta _{i,0}},$$
where $`\stackrel{~}{q}_{i,j}=\mathrm{exp}(2\pi i[B+i\omega ][\overline{\mathrm{\Delta }}_{i,j}])`$, and $`\zeta _{i,+},\zeta _{i,},\zeta _{i,0}`$ are as in Definition 4.11.
###### Proof.
We claim that $`T_{\overline{y},i,0}+T_{\overline{z},j,0}T_{\overline{y},j,0}T_{\overline{z},i,0}=\mathrm{\Delta }_{i,j}`$. Indeed, consider first a situation in which $`L_{3+i}`$ lies in the position represented in Figure 9, and $`L_{3+j}`$ lies close to it, but is slightly shifted towards the lower-right direction. Then the intersection points $`a_j`$ and $`c_j`$ lie close to $`a_i`$ and $`c_i`$, and following the triangular regions through the small deformation which takes $`L_{3+i}`$ to $`L_{3+j}`$, we easily see that $`T_{\overline{z},j,0}`$ is obtained by slightly truncating $`T_{\overline{z},i,0}`$ on its $`L_{3+i}`$ side. Similarly, $`T_{\overline{y},j,0}`$ is obtained by slightly truncating $`T_{\overline{y},i,0}`$, and since $`\mathrm{\Delta }_{i,j}`$ is simply the thin strip in between $`L_{3+i}`$ and $`L_{3+j}`$ the claim follows.
The same property remains true if $`L_{3+i}`$ and $`L_{3+j}`$ are further apart from each other. This can be checked explicitly for example in the configuration of Figure 8, where $`\mathrm{\Delta }_{i,j}`$ is as pictured on Figure 9 (right). (In this configuration the deformation from $`L_{3+i}`$ to $`L_{3+j}`$ passes through $`\overline{y}`$ and $`\overline{z}`$, so the triangles $`T_{\overline{z},i,0}`$ and $`T_{\overline{z},j,0}`$ lie on opposite sides of $`\overline{z}`$, and similarly for $`T_{\overline{y},i,0}`$ and $`T_{\overline{y},j,0}`$; this latter triangle is now the small region to the lower-right of $`\overline{y}`$ on Figure 8).
As a consequence, we have the identity
$$\frac{\psi _{\overline{y},i}\psi _{\overline{z},j}}{\psi _{\overline{y},j}\psi _{\overline{z},i}}=\stackrel{~}{q}_{i,j},$$
which implies the first formula in the proposition. The two other formulas are proved similarly, using the equalities $`T_{\overline{z},i,0}+T_{\overline{x},j,0}T_{\overline{z},j,0}T_{\overline{x},i,0}=\mathrm{\Delta }_{i,j}`$ and $`T_{\overline{x},i,0}+T_{\overline{y},j,0}T_{\overline{x},j,0}T_{\overline{y},i,0}=2\mathrm{\Delta }_{i,j}`$. ∎
###### Remark 4.15.
The various ratios computed in Propositions 4.74.14 are intrinsic quantities attached to the symplectic geometry of $`W_k`$, i.e. they are invariant under Hamiltonian deformations, irrespective of whether the vanishing cycles are represented by geodesics or not. Equivalently, they are invariant under rescalings of the chosen generators of the morphism spaces in $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$. On the other hand, if we allow ourselves to use the fact that the vanishing cycles are geodesics in a flat torus, we can also compute some interesting non-intrinsic quantities (i.e., quantities which depend on a particular choice of scaling of the generators).
For example, the invariance of $`L_0,L_1,L_2`$ under the translation of the torus which maps $`x_0`$ to $`y_0`$ (and $`y_0`$ to $`z_0`$, $`z_0`$ to $`x_0`$) implies that, for suitable choices of the marked points associated to the spin structures and of the isomorphisms between lines used to calculate boundary holonomies, $`\alpha _{xy}=\alpha _{yz}=\alpha _{zx}`$, $`\alpha _{yx}=\alpha _{xy}=\alpha _{xz}`$, and $`\alpha _{xx}=\alpha _{yy}=\alpha _{zz}`$. In fact, going over the calculations in the proofs of Propositions 4.7 and 4.10, and observing that, in terms of areas and boundary holonomies, the contributions of $`T_{xy,0}T_{yx,0}`$ and $`T_{xx,0}T_{yx,0}`$ are equivalent to those of $`\frac{1}{3}C`$ and $`\frac{1}{3}(\mathrm{\Sigma }_0C)`$ respectively, one easily checks that there exists a constant $`s0`$ such that
(4.2)
$$\begin{array}{c}\alpha _{xy}=\alpha _{yz}=\alpha _{zx}=sq_C^{1/3}\zeta _+,\hfill \\ \alpha _{xx}=\alpha _{yy}=\alpha _{zz}=sq_F^{1/3}q_C^{1/3}\zeta _0,\hfill \\ \alpha _{yx}=\alpha _{zy}=\alpha _{xz}=s\zeta _{},\hfill \end{array}$$
where by definition $`q_C^{1/3}=\mathrm{exp}(\frac{2\pi i}{3}[B+i\omega ][\overline{C}])`$ and $`q_F^{1/3}=\mathrm{exp}(\frac{2\pi i}{3}[B+i\omega ][\mathrm{\Sigma }_0])`$. Similarly, for suitable choices we have
(4.3)
$$\beta _{\overline{x},i}=s_iq_i^{1/3}\zeta _{i,+},\beta _{\overline{y},i}=s_iq_F^{1/3}q_i^{1/3}\zeta _{i,0},\mathrm{and}\beta _{\overline{z},i}=s_i\zeta _{i,},$$
where $`s_i`$ is a non-zero constant and $`q_i^{1/3}=\mathrm{exp}(\frac{2\pi i}{3}[B+i\omega ][\overline{C}_i])`$.
The formulas (4.2) and (4.3) are only valid in the flat case, when the complexified symplectic form on $`\mathrm{\Sigma }_0`$ is translation-invariant and the vanishing cycles are geodesics; however, in the general case we can always modify our choices of generators of the various morphism spaces by suitable scaling factors (or equivalently, modify the vanishing cycles by certain Hamiltonian isotopies) in order to make these formulas hold. It is therefore these simpler formulas that we will use in order to determine the mirror map in §5 below.
### 4.4. Simple degenerations
In this section we consider the situation where the symplectic area of one of the 2-cycles $`\overline{\mathrm{\Delta }}_{i,j}`$ becomes a multiple of that of the fiber $`\mathrm{\Sigma }_0`$. The vanishing cycles $`L_{3+i}`$ and $`L_{3+j}`$ are then Hamiltonian isotopic to each other in $`\mathrm{\Sigma }_0`$, and hence cannot be represented by disjoint geodesics anymore. However we can still represent $`L_{3+i}`$ by a closed geodesic, and $`L_{3+j}`$ by a small generic Hamiltonian perturbation of $`L_{3+i}`$, intersecting it transversely in two points. These two intersection points have Maslov indices $`0`$ and $`1`$ respectively (if we choose the same graded lifts as previously), and for this configuration we have:
###### Lemma 4.16.
If there exist integers $`n`$ and $`i<j`$ such that $`[\omega ][\overline{\mathrm{\Delta }}_{i,j}]=n[\omega ][\mathrm{\Sigma }_0]`$, then $`\mathrm{Hom}(L_{3+i},L_{3+j})`$ is graded isomorphic to $`H^{}(S^1)`$. Moreover, the differential
$$m_1:\mathrm{Hom}^0(L_{3+i},L_{3+j})\mathrm{Hom}^1(L_{3+i},L_{3+j})$$
is zero if $`[B][\overline{\mathrm{\Delta }}_{i,j}]+n[B][\mathrm{\Sigma }_0]`$, and an isomorphism otherwise.
###### Proof.
The only contributions to $`m_1`$ come from the two disks $`D^{}`$ and $`D^{\prime \prime }`$ bounded by $`L_{3+i}`$ and $`L_{3+j}`$. The 2-chain $`D^{}D^{\prime \prime }`$ in $`\mathrm{\Sigma }_0`$ has symplectic area zero, and is in fact given by $`D^{}D^{\prime \prime }=\mathrm{\Delta }_{i,j}n\mathrm{\Sigma }_0`$. Hence we can compare the coefficients $`\psi ^{}`$ and $`\psi ^{\prime \prime }`$ associated to these two disks by the same argument as in §4.3. Namely, $`\psi ^{}`$ and $`\psi ^{\prime \prime }`$ differ by a sign factor, a holonomy factor, and an area factor.
In this case the sign factor is $`1`$ (the sign rule for odd degree morphisms is slightly more subtle than that for even degree morphisms ; here we can see directly that the signs for $`D^{}`$ and $`D^{\prime \prime }`$ have to be different since the untwisted Floer homology of $`L_{3+i}`$ and $`L_{3+j}`$ is non-trivial); the holonomy factor is the total holonomy along $`(D^{}D^{\prime \prime })=L_{3+j}L_{3+i}`$, i.e. $`\mathrm{exp}(2\pi i_{D_{3+i}D_{3+j}}B)`$; and the area factor is $`\mathrm{exp}(2\pi i_{D^{}D^{\prime \prime }}B+i\omega )`$. It follows that
$$\psi ^{}=\mathrm{exp}\left(2\pi i[B+i\omega ]([\overline{\mathrm{\Delta }}_{i,j}]n[\mathrm{\Sigma }_0])\right)\psi ^{\prime \prime },$$
since $`D^{}D^{\prime \prime }+D_{3+i}D_{3+j}=\overline{\mathrm{\Delta }}_{i,j}n\mathrm{\Sigma }_0`$. Since $`m_1`$ is determined by the sum $`\psi ^{}+\psi ^{\prime \prime }`$, we conclude that $`m_1=0`$ if and only if $`[B+i\omega ]([\overline{\mathrm{\Delta }}_{i,j}]n[\mathrm{\Sigma }_0])`$ is an integer. ∎
In other words, if $`[B+i\omega ][\overline{\mathrm{\Delta }}_{i,j}]([B+i\omega ][\mathrm{\Sigma }_0])`$, then $`(L_{3+i},_{3+i})`$ and $`(L_{3+j},_{3+j})`$ are essentially identical, and we have a non-cancelling pair of extra morphisms of degrees $`0`$ and $`1`$ from $`L_{3+i}`$ to $`L_{3+j}`$; this mirrors the situation in which $`^2`$ is blown up twice at infinitely close points, in which case there is a rational $`2`$-curve and the derived category of coherent sheaves is richer than in the generic case. In all other situations the intersection points between $`L_{3+i}`$ and $`L_{3+j}`$, if any, are killed by the twisted Floer differential (even when $`L_{3+i}`$ and $`L_{3+j}`$ are Hamiltonian isotopic).
###### Remark 4.17.
It is important to note that, due to the presence of immersed convex polygonal regions with two edges on $`L_0L_1L_2`$ and two edges on $`L_{3+i}L_{3+j}`$ (with a corner at the intersection point of Maslov index $`1`$), we have to consider not only the Floer differential $`m_1`$, but also the higher-order composition $`m_3`$. For example, when $`L_{3+i}`$ and $`L_{3+j}`$ are Hamiltonian isotopic the composition
$$m_3:\mathrm{Hom}(L_0,L_2)\mathrm{Hom}(L_2,L_{3+i})\mathrm{Hom}^1(L_{3+i},L_{3+j})\mathrm{Hom}(L_0,L_{3+j})$$
is in general non-zero (and similarly with $`L_1`$ instead of $`L_0`$ or $`L_2`$).
As in §2.2, it is possible to describe things in a simpler and more unified manner by considering a suitable mutation of the exceptional collection $`(L_0,\mathrm{},L_{k+2})`$. Assume for simplicity that the two vanishing cycles which may coincide are $`L_3`$ and $`L_4`$, while the others are represented by distinct geodesics. Then we can modify the system of arcs $`\{\gamma _i\}`$ considered so far to a new ordered system of arcs $`\{\gamma _i^{}\}`$ such that $`\gamma _i^{}=\gamma _i`$ for $`i\{2,3\}`$, $`\gamma _3^{}=\gamma _2`$, and $`\gamma _2^{}`$ connects the origin to $`\lambda _3\mathrm{}`$ along the negative real axis. This gives rise to a new category $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i^{}\})`$, in which all objects but one can be identified with the objects $`L_i,i3`$ of $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$; thus, we denote by $`L_0,L_1,L^{},L_2,L_4,\mathrm{},L_{k+2}`$ the objects of $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i^{}\})`$. The morphisms and compositions not involving $`L^{}`$ are as in $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$.
The new vanishing cycle $`L^{}`$ is Hamiltonian isotopic to the image of $`L_3`$ under the positive Dehn twist along $`L_2`$. In particular, with the notations of Lemma 3.1, and for a suitable choice of orientation, its homology class is $`[L^{}]=[L_2][L_3]=b`$. Choosing a geodesic representative, we have $`|L_0L^{}|=2`$, $`|L_1L^{}|=1`$, $`|L^{}L_2|=1`$, and $`|L^{}L_{3+i}|=1`$ for $`i1`$, and all morphisms in $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i^{}\})`$ have degree 0.
Because $`L^{}`$ is Hamiltonian isotopic to the image of $`L_3`$ under the Dehn twist along $`L_2`$, the fiber $`\mathrm{\Sigma }_0`$ contains a 2-chain $`\mathrm{\Delta }^{}`$ with $`\mathrm{\Delta }^{}=L^{}+L_4L_2`$ and such that $`_\mathrm{\Delta }^{}\omega =_{\mathrm{\Delta }_{3,4}}\omega `$. Capping off $`\mathrm{\Delta }^{}`$ with the appropriate Lefschetz thimbles, we obtain a 2-cycle $`\overline{\mathrm{\Delta }}^{}`$ in $`M_k`$, with $`[\overline{\mathrm{\Delta }}^{}]=[\overline{\mathrm{\Delta }}_{3,4}]`$ in $`H_2(M_k,)`$. The composition
$$\mathrm{Hom}(L^{},L_2)\mathrm{Hom}(L_2,L_4)\mathrm{Hom}(L^{},L_4)$$
corresponds to an infinite series of triangular immersed regions in $`\mathrm{\Sigma }_0`$, of which in general two are embedded. The case where the symplectic area of $`\mathrm{\Delta }^{}`$ is a multiple of that of the fiber corresponds precisely to the situation where the two embedded triangular regions have equal symplectic areas. In general, the immersed triangles contributing to the composition can be labelled $`T_n^{}`$, $`n`$, in such a way that $`T_n^{}=T_0^{}+n\mathrm{\Delta }^{}+\frac{1}{2}n(n1)\mathrm{\Sigma }_0`$. Arguing as before, one easily shows that the composition is given by the contribution of $`T_0^{}`$ multiplied by the factor
$$\underset{n}{}(1)^nq_{}^{}{}_{}{}^{n}q_F^{n(n1)/2},\mathrm{where}q^{}=\mathrm{exp}(2\pi i[B+i\omega ][\overline{\mathrm{\Delta }}^{}])=\stackrel{~}{q}_{3,4}.$$
This multiplicative factor vanishes if and only if $`q^{}=q_F^k`$ for some $`k`$ (an easy way to see this is to view this factor as a theta function, see below), i.e. iff $`[B+i\omega ][\overline{\mathrm{\Delta }}^{}]([B+i\omega ][\mathrm{\Sigma }_0])`$. Hence, as in §2.2 the mutation makes it possible to avoid dealing with a non-trivial differential, and provides an alternative description in which the simple degeneration corresponds to one of the composition maps becoming identically zero.
### 4.5. Modular invariance and theta functions
In this section we study the modularity properties of the category $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\{\gamma _i\})`$ with respect to some of the parameters governing deformations of the complexified symplectic structure, and the relation with theta functions.
###### Proposition 4.18.
Consider two complexified symplectic forms $`\kappa =B+i\omega `$ and $`\kappa ^{}=B^{}+i\omega ^{}`$ on $`M_k`$, such that $`[\kappa ^{}][\mathrm{\Sigma }_0]=[\kappa ][\mathrm{\Sigma }_0]`$ and $`[\kappa ^{}][\kappa ]H^2(M_k,)(\kappa [\mathrm{\Sigma }_0])H^2(M_k,)`$. Then the categories $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\kappa ,\{\gamma _i\})`$ and $`\mathrm{Lag}_{\mathrm{vc}}(W_k,\kappa ^{},\{\gamma _i\})`$ are equivalent.
###### Proof.
First consider the situation where $`\omega ^{}=\omega `$, and $`B^{}=B+d\chi `$ for some 1-form $`\chi `$. Then the vanishing cycles $`L_i`$ remain the same, but the associated flat connections differ, and we can e.g. take $`_i^{}=_i2\pi i\chi `$. Then the contribution of a pseudo-holomorphic map $`u:(D^2,D^2)(\mathrm{\Sigma }_0,L_i)`$ is actually the same in both cases, since the holonomy term changes by $`\mathrm{exp}(2\pi i_{u(D^2)}\chi )`$, while the weight factor changes by $`\mathrm{exp}(2\pi i_{D^2}u^{}𝑑\chi )=\mathrm{exp}(2\pi i_{u(D^2)}\chi )`$. So, in the more general situation where $`[B^{}B]H^2(M_k,)`$ and $`[B^{}B][\mathrm{\Sigma }_0]=0`$ (still assuming $`\omega ^{}=\omega `$), after modifying $`B`$ by an exact term we can assume that $`B`$ and $`B^{}`$ coincide over $`\mathrm{\Sigma }_0`$, and that the integral of $`B^{}B`$ over each thimble $`D_i`$ is a multiple of $`2\pi `$. In this situation the vanishing cycles $`L_i`$ are the same, and the associated flat connections are gauge equivalent (since their holonomies differ by multiples of $`2\pi `$), so the corresponding twisted Floer theories are identical.
Next, consider the situation where $`[B+i\omega ]`$ changes by an integer multiple of $`[B+i\omega ][\mathrm{\Sigma }_0]`$. After adding an exact term to $`\kappa =B+i\omega `$ (which does not affect the category of vanishing cycles by Lemma 3.2 and by the above remark), we can assume that $`\kappa `$ and $`\kappa ^{}`$ coincide over $`\mathrm{\Sigma }_0`$, and that the relative cohomology class of $`\kappa ^{}\kappa `$ is an element of $`(\kappa [\mathrm{\Sigma }_0])H^2(M_k,\mathrm{\Sigma }_0;)`$.
Let $`D_i`$ and $`D_i^{}`$ be the thimbles associated to the arc $`\gamma _i`$ and to the symplectic forms $`\omega `$ and $`\omega ^{}`$ respectively. The integrality assumption on $`\kappa ^{}\kappa `$ implies that there exists an integer $`n_i`$ such that $`_{D_i}\kappa ^{}=n_i[\kappa ][\mathrm{\Sigma }_0]+_{D_i}\kappa `$. Since $`D_i`$ and $`D_i^{}`$ can be deformed continuously into each other (by deforming the horizontal distribution), there exists a 2-chain $`K_i`$ in $`\mathrm{\Sigma }_0`$ such that $`[D_i+K_iD_i^{}]=0`$ in $`H_2(M_k)`$. Then $`_{K_i}\omega =_{K_i}\omega ^{}=_{D_i}\omega ^{}=n_i[\omega ][\mathrm{\Sigma }_0]`$. Since the symplectic area of the 2-chain $`K_i\mathrm{\Sigma }_0`$ is an integer multiple of that of the fiber, the two vanishing cycles $`L_i^{}=D_i^{}`$ and $`L_i=D_i`$ are mutually Hamiltonian isotopic in $`\mathrm{\Sigma }_0`$, and hence we can assume that $`L_i^{}=L_i`$. Moreover, in $`H_2(M_k,L_i)`$ we have $`[D_i^{}]=[D_i]n_i[\mathrm{\Sigma }_0]`$. Therefore, $`_{D_i^{}}B^{}=_{D_i}B^{}n_i_{\mathrm{\Sigma }_0}B^{}=(_{D_i}B+n_i[B][\mathrm{\Sigma }_0])n_i[B][\mathrm{\Sigma }_0]=_{D_i}B`$. So the flat connections $`_i`$ and $`_i^{}`$ have the same holonomy, which implies that $`(L_i,_i)`$ and $`(L_i^{},_i^{})`$ behave identically for twisted Floer theory. ∎
This property explains the invariance of the structure coefficients ($`\alpha _{xy}`$, etc.) under certain changes of variables. More precisely, one easily checks that $`\zeta _+(q_Cq_F^3,q_F)=q_C^1q_F^2\zeta _+(q_C,q_F)`$, $`\zeta _{}(q_Cq_F^3,q_F)=q_C^1q_F^1\zeta _{}(q_C,q_F)`$, and $`\zeta _0(q_Cq_F^3,q_F)=q_C^1\zeta _0(q_C,q_F)`$. This implies that the quantities considered in Propositions 4.7 and 4.10 are invariant under the change of variables $`(q_C,q_F)(q_Cq_F^3,q_F)`$; a closer examination shows that the individual constants $`\alpha _{xy}`$, etc. are also invariant under this change of variables.
On the other hand, one easily checks that $`\zeta _+(q_Cq_F,q_F)=q_C^1\zeta _0(q_C,q_F)`$, $`\zeta _0(q_Cq_F,q_F)=\zeta _{}(q_C,q_F)`$, and $`\zeta _{}(q_Cq_F,q_F)=\zeta _+(q_C,q_F)`$, which may seem surprising at first. The reason is that this change of variables corresponds to a non-Hamiltonian deformation of e.g. $`L_1`$ which sweeps exactly once through the entire fiber $`\mathrm{\Sigma }_0`$. This deformation preserves the intersection points, but induces a non-trivial permutation of their labels: namely, $`x_0,y_0,z_0`$ become $`y_0,z_0,x_0`$ respectively, and $`x_1,y_1,z_1`$ become $`z_1,x_1,y_1`$ respectively. Thus, for example, $`\alpha _{xy}(q_C,q_F)=\alpha _{yx}(q_Cq_F,q_F)=\alpha _{zz}(q_Cq_F^2,q_F)`$ (and similarly for the other coefficients).
Another way to understand these invariance properties is to relate the functions $`\zeta _+`$, $`\zeta _{}`$, and $`\zeta _0`$ to theta functions. Recall that the ordinary theta function is an analytic function defined by
$$\theta (z,\tau )=\underset{n}{}\mathrm{exp}(\pi in^2\tau +2\pi inz),$$
where $`z`$ and $`\tau `$ (here $``$ is the upper half-plane $`\{\mathrm{Im}\tau >0\}`$). This function is quasiperiodic with respect to the lattice $`\mathrm{\Lambda }_\tau `$ generated by $`1`$ and $`\tau `$, and its behavior under translation by an element of the lattice is given by the formula
$$\theta (z+u\tau +v,\tau )=\mathrm{exp}(\pi iu^2\tau 2\pi iuz)\theta (z,\tau ).$$
The zeros of the theta function are the infinite set $`\left\{z=(n+\frac{1}{2})+(m+\frac{1}{2})\tau \right|n,m\}.`$
Here we consider theta functions with rational characteristics $`a,b`$, defined by
$$\theta _{a,b}(z,\tau )=\underset{n}{}\mathrm{exp}(\pi i(n+a)^2\tau +2\pi i(n+a)(z+b)).$$
Let us introduce new variables $`q=\mathrm{exp}(\pi i\tau )`$ and $`w=\mathrm{exp}(\pi iz).`$ Now the following three $`\theta `$-functions play a very important role in our considerations:
$`\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau )`$ $`=`$ $`\mathrm{exp}(\frac{i\pi }{2})q^{3/4}{\displaystyle \underset{n}{}}(1)^nw^{6n+3}q^{3n^2+3n},`$
$`\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau )`$ $`=`$ $`\mathrm{exp}(\frac{i\pi }{6})q^{1/12}{\displaystyle \underset{n}{}}(1)^nw^{6n+1}q^{3n^2+n},`$
$`\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )`$ $`=`$ $`\mathrm{exp}(\frac{i\pi }{6})q^{1/12}{\displaystyle \underset{n}{}}(1)^nw^{6n1}q^{3n^2n}.`$
The zero set of the function $`\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau )`$ is $`\left\{\frac{n}{3}+m\tau \right|n,m\}`$, while the zero sets of the functions $`\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau )`$ and $`\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )`$ are
$$\left\{\frac{n}{3}+(m+\frac{1}{3})\tau \right|n,m\}\text{and}\left\{\frac{n}{3}+(m\frac{1}{3})\tau \right|n,m\}$$
respectively. These three theta functions can be viewed as holomorphic sections of a line bundle of degree 3 on the elliptic curve $`E=/\mathrm{\Lambda }_\tau `$; considering the zero sets, we see that this line bundle is $`𝕃=𝒪_E(3(0))`$. These three sections of $`𝕃`$ determine an embedding of the elliptic curve $`E=/\mathrm{\Lambda }_\tau `$ into the projective plane, given by
$$z(\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau ):\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau ):\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )).$$
Observe that the two functions
$$\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau )\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau )\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )\mathrm{and}\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau )^3+\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau )^3+\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )^3$$
coincide up to a constant multiplicative factor, since they both correspond to holomorphic sections of the line bundle $`𝕃^3`$ over $`E`$, and an easy calculation shows that they have the same zero set $`\{\frac{n}{3}+\frac{m}{3}\tau |n,m\}`$. Therefore, the image of the above embedding of $`E`$ into $`^2`$ is the cubic given by the equation
$$(A^3+B^3+C^3)XYZABC(X^3+Y^3+Z^3)=0,$$
where $`(A,B,C)`$ are the values of the three theta functions at any given point of $`/\mathrm{\Lambda }_\tau `$ (not in $`\frac{1}{3}\mathrm{\Lambda }_\tau `$).
Consider the function
$$\left(\frac{\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau )}{\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )}\right)^3=\left(\frac{_n(1)^nw^{6n+1}q^{n(3n+1)}}{_n(1)^nw^{6n1}q^{n(3n1)}}\right)^3.$$
Substituting $`q^2=q_F`$ and $`w^6=q_C`$, one easily checks that this coincides with the expression which appears in Proposition 4.7,
$$\frac{\alpha _{xy}\alpha _{yz}\alpha _{zx}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=q_C\left(\frac{_n(1)^nq_C^nq_F^{n(3n+1)/2}}{_n(1)^nq_C^nq_F^{n(3n1)/2}}\right)^3.$$
Similarly,
$$\left(\frac{\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau )}{\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau )}\right)^3=q^2\left(\frac{_n(1)^nw^{6n+3}q^{3n^2+3n}}{_n(1)^nw^{6n1}q^{3n^2n}}\right)^3=q^2\left(\frac{_n(1)^nw^{6n3}q^{3n^23n}}{_n(1)^nw^{6n1}q^{3n^2n}}\right)^3.$$
After the same substitution $`q^2=q_F`$ and $`w^6=q_C`$, this coincides with the expression given in Proposition 4.10,
$$\frac{\alpha _{xx}\alpha _{yy}\alpha _{zz}}{\alpha _{yx}\alpha _{zy}\alpha _{xz}}=\frac{q_F}{q_C}\left(\frac{_n(1)^nq_C^nq_F^{3n(n1)/2}}{_n(1)^nq_C^nq_F^{n(3n1)/2}}\right)^3.$$
Similarly, in the case where (4.2) holds, one easily checks that
(4.4)
$$\begin{array}{c}\alpha _{xy}=\alpha _{yz}=\alpha _{zx}=\stackrel{~}{s}e^{2i\pi /3}\theta _{\frac{1}{6},\frac{1}{2}}(3z_0,3\tau ),\hfill \\ \alpha _{xx}=\alpha _{yy}=\alpha _{zz}=\stackrel{~}{s}\theta _{\frac{1}{2},\frac{1}{2}}(3z_0,3\tau ),\hfill \\ \alpha _{yx}=\alpha _{zy}=\alpha _{xz}=\stackrel{~}{s}e^{2i\pi /3}\theta _{\frac{5}{6},\frac{1}{2}}(3z_0,3\tau ),\hfill \end{array}$$
where $`\tau =[B+i\omega ][\mathrm{\Sigma }_0]`$, $`z_0=\frac{1}{3}[B+i\omega ][\overline{C}]`$, and $`\stackrel{~}{s}=e^{i\pi /2}q_F^{1/24}q_C^{1/6}s0`$. Similar interpretations can be made for the quantities considered in Propositions 4.124.14 and in (4.3).
## 5. Proof of the main theorems
The derived categories considered in §2 depend on an elliptic curve $`E`$, two degree 3 line bundles $`_1,_2`$ over $`E`$, and $`k`$ points $`p_1,\mathrm{},p_k`$ on $`E`$. Meanwhile, the categories considered in §4 depend on a cohomology class $`[B+i\omega ]H^2(M_k,)`$. We now show how to relate these two sets of parameters.
Fix the cohomology class $`[B+i\omega ]H^2(M_k,)`$, and consider the category $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$ studied in §4. With the notations of §3.4, assume that $`[\omega ][\overline{\mathrm{\Delta }}_{i,j}]`$ is not an integer multiple of $`[\omega ][\mathrm{\Sigma }_0]`$ for any $`i,j\{0,\mathrm{},k1\}`$. Then $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))`$ admits a full strong exceptional collection $`(L_0,\mathrm{},L_{k+2})`$, whose properties have been studied in §4. In particular, the objects and morphisms in this exceptional collection are the same as for the exceptional collection $`\sigma =(𝒪_{X_K},\pi ^{}𝒯_^2(1),\pi ^{}𝒪_^2(1),𝒪_{l_1},\mathrm{},𝒪_{l_k})`$ considered in §2 for the derived category of coherent sheaves on a (possibly noncommutative) Del Pezzo surface. Hence, our goal is now to compare the composition laws and show that, for a suitable choice of the parameters $`(E,_1,_2,K)`$, the algebra of homomorphisms of the exceptional collection $`(L_0,\mathrm{},L_{k+2})`$ is isomorphic to the algebra $`B_{K,\mu }`$ considered in §2. More precisely, we claim:
###### Proposition 5.1.
Let $`E`$ be the elliptic curve $`/\mathrm{\Lambda }_\tau `$, where $`\tau =[B+i\omega ][\mathrm{\Sigma }_0]`$, realized as a plane cubic via the embedding $`j:E^2`$ given by $`z(\vartheta _+(z):\vartheta _0(z):\vartheta _{}(z))`$, where
$$\vartheta _+(z)=e^{2i\pi /3}\theta _{\frac{1}{6},\frac{1}{2}}(3z,3\tau ),\vartheta _0(z)=\theta _{\frac{1}{2},\frac{1}{2}}(3z,3\tau ),\text{and}\vartheta _{}(z)=e^{2i\pi /3}\theta _{\frac{5}{6},\frac{1}{2}}(3z,3\tau ).$$
Let $`z_0=\frac{1}{3}[B+i\omega ][\overline{C}]`$, and for $`i\{0,\mathrm{},k1\}`$ let $`p_i=\frac{1}{3}[B+i\omega ][\overline{C}_i]`$. Finally, let $`_1=𝒪_E(3(z_0))`$ and $`_2=𝒪_E(3(0))`$. Then the algebra of homomorphisms of the exceptional collection $`(L_0,\mathrm{},L_{k+2})`$ is isomorphic to $`B_{K,\mu }`$, where $`\mu `$ is determined by $`(E,_1,_2)`$ via Construction 2.9 and $`K=\{j(z_0+p_0),\mathrm{},j(z_0+p_{k1})\}`$.
###### Proof.
After a suitable rescaling of the chosen bases of the morphism spaces (or just by deforming to the situation where the fiber is flat and the vanishing cycles are geodesics), we can assume that the compositions of morphisms between the objects $`L_0,\mathrm{},L_{k+2}`$ are given by the formulas (4.2) and (4.3). We identify the vector spaces $`U=\mathrm{Hom}(L_0,L_1)`$, $`V=\mathrm{Hom}(L_1,L_2)`$, and $`W=\mathrm{Hom}(L_0,L_2)`$ with $`^3`$ by considering the bases $`(x_0,y_0,z_0)`$, $`(x_1,y_1,z_1)`$, and $`(\overline{x},\overline{y},\overline{z})`$. The composition tensor $`\mu :VUW`$ is determined by the three constants $`a=\alpha _{xy}=\alpha _{yz}=\alpha _{zx}`$, $`b=\alpha _{xx}=\alpha _{yy}=\alpha _{zz}`$, and $`c=\alpha _{yx}=\alpha _{zy}=\alpha _{xz}`$. In particular, given an element $`v=(X,Y,Z)V`$, the composition map $`\mu _v=\mu (v,):UW`$ is given by the matrix
(5.1)
$$\left(\begin{array}{ccc}\alpha _{xx}X& \alpha _{yz}Z& \alpha _{zy}Y\\ \alpha _{xz}Z& \alpha _{yy}Y& \alpha _{zx}X\\ \alpha _{xy}Y& \alpha _{yx}X& \alpha _{zz}Z\end{array}\right)=\left(\begin{array}{ccc}bX& aZ& cY\\ cZ& bY& aX\\ aY& cX& bZ\end{array}\right)$$
which has rank 2 precisely when
(5.2)
$$det(\mu _v)=(a^3+b^3+c^3)XYZabc(X^3+Y^3+Z^3)=0.$$
By (4.4), the constants $`a,b,c`$ are (up to a non-zero constant factor) the values of the theta functions $`\vartheta _+,\vartheta _0,\vartheta _{}`$ at the point $`z_0`$. Therefore, by the discussion in §4.5, there are two possibilities:
1. if $`z_0\frac{1}{3}\mathrm{\Lambda }_\tau `$, then $`abc=0`$ and $`\mu _v`$ always has rank 2; as explained in §2.3 this corresponds to a commutative situation;
2. if $`z_0\frac{1}{3}\mathrm{\Lambda }_\tau `$, then (5.2) defines a cubic $`\mathrm{\Gamma }_V(V)=^2`$, and this cubic is precisely the image of the embedding $`j`$.
The same situation holds for $`\mu _u`$; interestingly, under the chosen identifications of $`(U)`$ and $`(V)`$ with $`^2`$, the two subschemes $`\mathrm{\Gamma }_U(U)`$ and $`\mathrm{\Gamma }_V(V)`$ determined by the equations $`det(\mu _u)=0`$ and $`det(\mu _v)=0`$ coincide exactly. However, with this description, the isomorphism $`\sigma :\mathrm{\Gamma }_V\mathrm{\Gamma }_U`$ which takes $`v`$ to the point of $`\mathrm{\Gamma }_U`$ corresponding to $`\mathrm{Ker}\mu _v`$ is not the identity map. Here the reader is referred to the discussion on pp. 37–38 of , which we follow loosely.
Given a point $`v=(X:Y:Z)\mathrm{\Gamma }_V`$, the kernel of $`\mu _v`$ can be obtained as the cross-product of any two of the rows of the matrix (5.1). Taking e.g. the first two rows, we obtain that the corresponding point of $`\mathrm{\Gamma }_U`$ is
(5.3)
$$\sigma (X:Y:Z)=(a^2XZbcY^2:c^2YZabX^2:b^2XYacZ^2).$$
Observe that $`j`$ maps the origin to $`(1:0:1)\mathrm{\Gamma }_V`$, and that the corresponding point in $`\mathrm{\Gamma }_U`$ is $`\sigma (1:0:1)=(a:b:c)=j(z_0)`$. Hence, considering only the situation where $`\mathrm{\Gamma }_U\mathrm{\Gamma }_VE`$, and identifying $`E`$ with $`\mathrm{\Gamma }_V`$ by means of the embedding $`j`$, the identification of $`E`$ with $`\mathrm{\Gamma }_U`$ is given by the embedding $`\sigma j`$, which is the composition of $`j`$ with the translation by $`z_0`$. Therefore, the line bundles on $`E`$ induced by the two inclusions of $`E`$ into $`(U)`$ and $`(V)`$ are respectively $`(\sigma j)^{}𝒪_^2(1)=𝒪_E(3(z_0))=_1`$ and $`j^{}𝒪_^2(1)=𝒪_E(3(0))=_2`$. It then follows from the discussion in §2.3 that the composition tensor $`\mu `$ corresponds to the data $`(E,_1,_2)`$. This remains true even when $`z_0\frac{1}{3}\mathrm{\Lambda }_\tau `$, since in that case we have $`_1_2`$ and the composition tensor associated to the triple $`(E,_1,_2)`$ is that of the usual projective plane (see Remark 2.10).
Next we consider the composition $`\mathrm{Hom}(L_2,L_{3+i})W\mathrm{Hom}(L_0,L_{3+i})`$. Choosing generators of the lines $`\mathrm{Hom}(L_2,L_{3+i})`$ and $`\mathrm{Hom}(L_0,L_{3+i})`$ we can view this map as a linear form on $`W`$. In the given basis of $`W`$, this linear form is given by $`(\beta _{\overline{x},i},\beta _{\overline{y},i},\beta _{\overline{z},i})`$, which by (4.3) coincides up to a non-zero constant factor with
$$(\vartheta _+(p_i),\vartheta _0(p_i),\vartheta _{}(p_i)).$$
On the other hand we know from §2.3 that the kernel of this linear form should be exactly $`\mathrm{Im}\mu _{v_i}`$, where $`v_i\mathrm{\Gamma }_V`$ is the point being blown up.
For any $`v=(X:Y:Z)\mathrm{\Gamma }_V`$, the projection $`WW/\mathrm{Im}\mu _v`$ is a linear form given up to a scaling factor by the dot product of any two columns of the matrix (5.1). Taking e.g. the first two columns, we obtain that the expression of this linear form relatively to our chosen basis of $`W`$ is
$$(c^2XZabY^2,a^2YZbcX^2,b^2XYacZ^2).$$
Interestingly, if we assume that $`(X:Y:Z)=\sigma (\stackrel{~}{X}:\stackrel{~}{Y}:\stackrel{~}{Z})`$, where $`\sigma `$ is the transformation given by (5.3), then this expression simplifies to a scalar multiple of $`(\stackrel{~}{X},\stackrel{~}{Y},\stackrel{~}{Z})`$. Hence, we conclude that $`v_i=\sigma (j(p_i))=j(z_0+p_i)`$. ∎
###### Remark 5.2.
At this point the reader may legitimately be concerned that, since the homology classes $`[\overline{C}]`$ and $`[\overline{C}_i]`$ are canonically defined only up to a multiple of $`[\mathrm{\Sigma }_0]`$, and since $`[B]`$ is only defined up to an element of $`H^2(M_k,)`$, the points $`z_0`$ and $`p_i`$ of $`E`$ are canonically determined only up to translations by elements of $`\frac{1}{3}\mathrm{\Lambda }_\tau `$. However, the line bundle $`_1=𝒪_E(3(z_0))`$ is not affected by this ambiguity in the determination of $`z_0`$, and neither are the relative positions of the points $`p_i`$, since the quantity $`p_jp_i=[B+i\omega ][\overline{\mathrm{\Delta }}_{i,j}]`$ is well-defined up to an element of $`\mathrm{\Lambda }_\tau `$. Moreover, a simultaneous translation of all the blown up points by an element of $`\frac{1}{3}\mathrm{\Lambda }_\tau `$ amounts to an automorphism of the triple $`(E,_1,_2)`$, which does not actually affect the category. (From the point of view of the embedding $`j`$, this automorphism simply permutes the homogeneous coordinates $`X,Y,Z`$ and multiplies them by cubic roots of unity; this is consistent with the observation made after the proof of Proposition 4.18).
Theorems 1.4 and 1.6 now follow directly from the discussion. Namely, in the case of a blowup of $`^2`$ at a set $`K=\{p_0,\mathrm{},p_{k1}\}`$ of $`k`$ distinct points (Theorem 1.4), we consider a cubic curve $`E^2`$ which contains all the points of $`K`$, and view it as an elliptic curve $`/\mathrm{\Lambda }_\tau `$ for some $`\tau `$ with $`\mathrm{Im}\tau >0`$. This allows us to view the points $`p_i`$ as elements of $`/\mathrm{\Lambda }_\tau `$ (well-defined up to a simultaneous translation of all $`p_i`$ by an element of $`\frac{1}{3}\mathrm{\Lambda }_\tau `$, since the origin can be chosen at any of the flexes of $`E`$; however by Remark 5.2 this does not matter for our construction). Then we equip $`M_k`$ with a complexified symplectic structure such that $`[B+i\omega ][\mathrm{\Sigma }_0]=\tau `$, $`[B+i\omega ][\overline{C}]=0`$, and $`[B+i\omega ][\overline{C}_i]=3p_i`$. The existence of such a $`B+i\omega `$ follows from a standard result about symplectic structures on Lefschetz fibrations:
###### Proposition 5.3 (Gompf).
Given any cohomology class $`[\zeta ]H^2(M_k,)`$ such that $`[\zeta ][\mathrm{\Sigma }_0]>0`$, the manifold $`M_k`$ admits a symplectic structure in the cohomology class $`[\zeta ]`$, for which the fibers of $`W_k`$ are symplectic submanifolds.
###### Proof.
The map $`W_k:M_k`$ is a Lefschetz fibration, and the argument given in the proof of \[7, Theorem 10.2.18\] can be adapted in a straightforward manner to this situation, even though the base of the fibration is not compact. (Alternatively, one can also work with the compactified fibration $`\overline{W_k}:\overline{M}^1`$). The symplectic form $`\omega `$ constructed by this argument lies in the cohomology class $`t[\zeta ]+W_k^{}([\mathrm{vol}_{}])`$ for some constant $`t>0`$; since the area form on $``$ is exact, we have $`[\omega ]=t[\zeta ]`$, and scaling $`\omega `$ by a constant factor we obtain the desired result. ∎
By Proposition 5.1 the algebra of homomorphisms of the exceptional collection $`(L_0,\mathrm{},L_{k+2})`$ is then isomorphic to $`B_K`$, which implies that $`𝐃^b(\mathrm{Lag}_{\mathrm{vc}}(W_k))𝐃^b(\text{mod–}B_K)𝐃^b(\mathrm{coh}(X_K))`$.
In the case of a noncommutative blowup of $`^2`$ (Theorem 1.6), consider the triple $`(E,_1,_2)`$ associated to the underlying noncommutative $`^2`$, and view again $`E`$ as a quotient $`/\mathrm{\Lambda }_\tau `$. Choose $`z_0`$ (well-defined up to an element of $`\frac{1}{3}\mathrm{\Lambda }_\tau `$) such that $`_2_1^1𝒪_E(3(z_0)3(0))\mathrm{Pic}^0(E)`$. As explained in §2.3, the blown up points must all lie in $`\mathrm{\Gamma }_V(V)`$, and under the identification $`\mathrm{\Gamma }_VE`$ they can be viewed as elements $`p_i/\mathrm{\Lambda }_\tau `$. Equip $`M_k`$ with a complexified symplectic structure such that $`[B+i\omega ][\mathrm{\Sigma }_0]=\tau `$, $`[B+i\omega ][\overline{C}]=3z_0`$, and $`[B+i\omega ][\overline{C}_i]=3(p_iz_0)`$. By Proposition 5.1 the algebra of homomorphisms of the exceptional collection $`(L_0,\mathrm{},L_{k+2})`$ is then isomorphic to $`B_{K,\mu }`$, which yields the desired equivalence of categories.
Theorem 1.5 is proved similarly, working with the mutated exceptional collections $`\tau ^{}`$ (introduced in §2.2) and $`(L_0,L_1,L^{},L_2,L_4,\mathrm{},L_{k+2})`$ (introduced in §4.4). The details are left to the reader.
###### Remark 5.4.
The construction carried out for Theorem 1.4 also applies to some limit situations in which $`X_K`$ is actually not a Del Pezzo surface. For example, the argument applies equally well to the situation where $`^2`$ is blown up at nine points which lie at the intersection of two elliptic curves. In this case the mirror is an elliptic fibration over $``$ for which the compactification has a smooth fiber at infinity. Compared to that of $`^2`$ ($`k=0`$), this extreme case where $`k=9`$ lies at the opposite end of the spectrum that we consider.
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# Microrheological Characterisation of Anisotropic Materials
## Abstract
We describe the measurement of anisotropic viscoelastic moduli in complex soft materials, such as biopolymer gels, via video particle tracking microrheology of colloid tracer particles. The use of a correlation tensor to find the axes of maximum anisotropy without prior knowledge, and hence the mechanical director, is described. The moduli of an aligned DNA gel are reported, as an application of the technique; this may have implications for high DNA concentrations in vivo. We also discuss the errors in microrheological measurement, and describe the use of frequency space filtering to improve displacement resolution, and hence probe these typically high modulus materials.
Many important biological molecules are capable of forming liquid crystalline phases due to their inherent stiffness and aspect ratio. In this paper we demonstrate how a development of existing microrheological techniques offers a means to characterize local anisotropic properties in such media.
We believe that this technique is unique in its ability to determine the mechanical properties of typical anisotropic materials, where the director may vary over a mesoscopic lengthscale. Our approach also utilises several improvements to existing microrheology measurements, which are likely to be particularly important for such high modulus, anisotropic materials for which there are particular challenges. Although we have chosen to work with DNA, many materials, both naturally occuring and synthetic, form anisotropic phases. Examples of such materials of biological importance include aqueous biopolymer gels, including DNA, but also cellulose, xanthan and F-actin, amongst others. Anisotropy may also occur in intracellular materials, to which standard 2-D microrheology has previously been applied Lau et al. (2003); Yamada et al. (2000); Tseng et al. (2002).
Particle tracking microrheology is a technique that has attracted much interest in experimental soft condensed matter as a method for measuring the viscoelastic properties of soft materials on the micron scale. The mechanical response of such a material may be probed by monitoring the time evolution of the displacement of embedded microscopic tracer particles subject to thermal motion.
The idea of using embedded particles to probe material properties was first described in a seminal paper by Freundlich and Seifriz Freundlich and Seifriz (1922), studying gelatin gels using magnetic beads in an external field. However these experiments suffered from the use of aspherical particles, limited force control and inaccuracy in the detected particle position. The technique was rediscovered by Mason and Weitz Mason and Weitz (1995) who proposed an extension to the familiar Stokes-Einstein diffusion equation to complex, frequency dependent moduli,
$$\stackrel{~}{G}(s)=\frac{dk_BT}{3\pi asr(s)^2}$$
(1)
where the real viscoelastic shear modulus $`\stackrel{~}{G}`$ of a material at a temperature $`T`$ may be calculated as a function of Laplace frequency, $`s`$, for beads of radius $`a`$, with a displacement $`r`$, in $`d`$ dimensions, given the mean squared displacement (MSD), $`r^2(\tau )=\left|r(t+\tau )r(t)\right|^2`$. In addition to improved experimental techniques, this presentation of a generalised Stokes-Einstein relation (GSER) has encouraged renewed interest, and resulted in many advances, greatly expanding the power of microrheology as a technique (for a review see e.g. Breedveld and Pine (2003)).
From the MSD, under the same assumptions as the generalised Stokes-Einstein relation (GSER) (Eqn.1) we can calculate Xu et al. (1998) the creep compliance by multiplication with a constant,
$$\mathrm{\Gamma }(\tau )=\frac{\pi a}{k_BT}r^2(\tau )$$
(2)
or extract the frequency dependent complex viscoelastic shear modulus via a Laplace or Fourier transform Mason (2000); Mason et al. (1997); Gittes et al. (1997). If the material is a gel, i.e. a material composed of a polymer network in a fluid medium, then the network spacing (mesh size) may also be determined via particle tracking Valentine et al. (2001).
Conventional microrheology has so far used the 2-dimensional MSD in Eq. (1), but we choose to decompose the displacement into 2 orthogonal 1-dimensional displacements. Although an individual random walk is spatially anisotropic Rudnick et al. (1987), the probability density function averaged over a number of walks will be circularly symmetric in an isotropic material. In a mechanically anisotropic material, the ensemble average of the random walks will also be elliptically symmetric (Fig. 1) with major and minor axes along the mechanical director. The GSER thus allows us to calculate the dissipative and elastic properties of the gel along specific axes, subject to the same conditions as the Stokes derivation: no-slip boundary conditions on a sphere in an incompressible continuum fluid, with no inertia. In addition the implicit assumption is made that the Stokes drag for viscous fluids may be generalised to viscoelastic materials at all frequencies, and that the probe particles have a negligible effect on the material Mason (2000); Levine and Lubensky (2000). Although in general viscosity is an 81 element 4th rank tensor, many of these elements are either zero or repeated in typical cases. For example a nematic liquid crystal has 5 independent intrinsic viscosities describing diffusion. Our case is is analogous to the drag on a sphere moving though a nematic liquid crystal Ruhwandl and Terentjev (1996); Stark and Ventzki (2001); Loudet et al. (2004) and so, since we have a rotational symmetry axis parallel to the director, the Brownian motion is governed by two independent diffusion coefficients where $`D_/=k_bT/(6\pi \eta _/)a`$, as the derivation of the Stokes-Einstein relation is separable in spatial dimension. Generalising to complex viscosity, we thus expect the GSER to hold in anisotropic materials, with two independent complex viscosities $`\eta _{}^{}`$ and $`\eta _{}^{}`$. The ratio, $`\eta _{}^{}/\eta _{}^{}`$ is therefore a measure of the anisotropy of the system. Similarly to the isotropic case, it should be noted that there is no rigorous theoretical backing for the extension to complex moduli, only phenomeological justification.
The probe particles themselves may be tracked via a number of methods. We choose real-space multiparticle ‘single-point’ video tracking due to both the large amount of information available for analysis, and the relative ease of setup. Video is acquired with a Zeiss-Axioplan optical microscope with a 100x oil immersion lens (N.A. = 1.30) and conventional commercial CCD video camera (Sony SCC-DC138P PAL), directly into a computer using a Hauppauge Impact VCB capture card. The analysis is performed using custom-written MATLAB scripts. To detect the small Brownian displacements of colloidal particles in high modulus materials we need as high a spatial resolution as possible. This may be measured by determining the apparent displacement of fixed spheres, following Savin and Doyle (2005). We use a 2-dimensional bandpass filter to select appropriate spatial frequencies in Fourier space, and thus filter out noise and background variation. Regions above an intensity threshold are selected, then checks for eccentricity and shape solidity are performed to determine the quality of the probe images. Finally the centroid of the region is calculated, weighted by the original pixel intensity. By this method we have achieved a 1D displacement resolution of up to 3nm, with 4nm in typical operation.
Since the video data in conventional cameras is recorded in an interlaced manner, in which the two fields of alternate lines are recorded either 1/50 or 1/60 of a second apart in each frame, one typically de-interlaces the fields and analyses each field separately. However, by tracking fixed colloidal spheres at high resolution, we show in Fig. 2 that the two fields are displaced from each other in a time varying fashion. The source of this displacement is unclear but it means that de-interlacing may be inadvisable when tracking small displacements, and is not carried out in the current work.
High modulus gels often have some in-built shear stress when they are made, which relaxes over a very long time period. In addition one may wish to investigate the microrheological properties of sheared, flowing, or even living samples. However, probe particles in such a system will be subject to a net flow field, in addition to Brownian motion. Provided that the flow field is uniform, $`\mathrm{\Delta }r(\tau )=0`$ for a random walk, so the net drift velocity can be calculated from the ensemble average and then subtracted. If the flow field varies across the field of view this approach is clearly invalid, and if the velocity field is a function of depth the situation is complicated Zumofen et al. (1990), but in practice we find this simple approach to be adequate in this work, where we measure the drift due to relaxation as $``$ $`3\mu `$m/min.
We write the displacement in tensor form to remove coordinate system dependence from the arbitrary x (horizontal) y (vertical) Cartesian axes with the origin in the top-left in the video. The eigenvector relation (Eqn. 3) can then be written, and the resulting eigenvectors give the basis in which displacements along the two orthogonal axes are least correlated, and hence maximally anisotropic. We use the eigenvectors calculated from the shortest lag time, $`\tau `$, for which there are the most events and hence the smallest error, to find the new basis and hence recalculate displacements in these coordinates. For an anisotropic material one would expect the new axes to lie parallel and perpendicular to the average optical director. By solving the equation for the specific example shown in Fig. 3, and determining the eigenvectors (superimposed on Fig. 3), we can indeed see the correlation of the optical alignment with the directions of anisotropy. The directionality and strength of alignment in the material can thus be determined from particle tracking alone, with no prior knowledge required. We will discuss the nature of the sample below.
$$\left[\begin{array}{cc}(\mathrm{\Delta }x(\tau ))^2& \mathrm{\Delta }x(\tau )\mathrm{\Delta }y(\tau )\\ \mathrm{\Delta }x(\tau )\mathrm{\Delta }y(\tau )& (\mathrm{\Delta }y(\tau ))^2\end{array}\right]𝐀=\lambda 𝐀$$
(3)
The MSD must then be computed. This is a concise representation of the rheological data, related to the creep compliance by a constant (Eqn. 2). We need to determine the errors in microrheological data, with a particular view to judging whether anisotropy is significant, so the following expression is derived, making the assumption that a ‘typical’ value for $`x`$ is $`\sqrt{\mathrm{\Delta }x^2(\tau )}`$. Assuming Gaussian statistics, we have that $`\sigma _{x^2}=2\sigma _xx`$ and so,
$`\sigma _{\mathrm{\Delta }x^2(\tau )}`$ $`=`$ $`{\displaystyle \frac{2\sigma _xx}{\sqrt{N}}}{\displaystyle \frac{2\sigma _x\sqrt{\mathrm{\Delta }x^2(\tau )}}{\sqrt{N}}}`$ (4)
$``$ $`{\displaystyle \frac{2\mathrm{\Delta }x^2(\tau )}{\sqrt{N}}}`$
where $`N`$ is the number of events contributing to $`\mathrm{\Delta }x^2(\tau )`$. The expression is of the appropriate form as $`\sigma _x=\sqrt{\mathrm{\Delta }x^2(\tau )}`$ in this case. In addition, the noise at the resolution limit results in a ‘static error’ Savin and Doyle (2005) which is exhibited as a constant offset added to $`\mathrm{\Delta }r^2(\tau )`$, and thus can be subtracted.
A number of methods have been proposed to arrive at a complex shear modulus Gittes et al. (1997); Dasgupta et al. (2002); Mason (2000); Mason and Weitz (1995). We use a simple discrete Laplace transform,
$$\stackrel{~}{x}(s)=\frac{\tau _N}{N}\underset{n=0}{\overset{N}{}}x(\tau _n)e^{s\tau _n}$$
(5)
where $`N`$ is the number of samples spread over the interval $`\tau =0`$ to $`\tau =\tau _N`$, followed by 1st order polynomial spline fitting and analytic continuation into the complex plane, substituting $`i\omega `$ for $`s`$ in the locally fitted form Mason and Weitz (1995); Bird et al. (1977). No assumptions are made about the form of the response which is well approximated by a linear spline fit over a 5 point window, and truncation errors are found to be small for a discrete Laplace transform given a typical range of data in time (1/25s to 60s).
To demonstrate the utility of microrheology as a tool for probing anisotropic materials, we apply the method described above to shear-aligned deoxyribonucleic acid (DNA). DNA in water, with some ammonium acetate as a counter ion, forms well-documented liquid crystalline phases Strzelecka et al. (1988) at high concentrations. We look at genomic DNA (Fluka, Deoxyribonucleic acid sodium salt from herring testes. Product: 31162) at a moderate concentration (11mg/ml) at 298K in 0.25M ammonium acetate, contained in a cell made of a metal washer (0.80mm deep) sandwiched by a glass microscope slide and coverslip, holding 12$`\mu `$l of sample. The sample is left to relax for 1 hour, then shear stress is applied by moving the coverslip, to induce alignment, which may be monitored by optical birefringence (Fig. 3). We use 0.3$`\mu `$m diameter probes (polystyrene latex, Agar Scientific) which are much larger than the expected mesh size of 15nm (calculated assuming uniform volume filling), with our field of view set well away from the cell surfaces. The video is kept interlaced, as discussed previously, and 120 seconds is analyzed.
Although the sample exhibited only weak optical alignment, we find significant mechanical differences between the axes parallel and perpendicular to the optical director. Figure 4 plots the MSD, showing greater motion perpendicular to the optical director than parallel. This seems reasonable in a network of stretched chains, with the optical director lying along the chain backbone. The viscoelastic moduli are also calculated (Fig. 5). The material is clearly stiffer to shear deformation in the direction along which the chains are, on average, already stretched. This is as we would expect for a polymer network which has been stretched uniaxially, where a proportion of chains are close to fully extended between entanglements/ crosslinks in this direction. It is interesting to note that there appears to be a cross-over in the elastic modulus, where the stretched direction has a lower elastic modulus than the perpendicular direction at frequencies below 3Hz, the explanation for which is unclear. The plateau in $`G^{}`$ is typical of a gel in the elastic region of its frequency response. Some studies Goodman et al. (2002); Chen et al. (2003) have been done on DNA at lower concentrations, which indicate precursors to this elastic regime.
We have presented a method to look at the viscoelastic properties of anisotropic materials on typical (mesoscopic) lengthscales. The technique both permits the identification of the principal axes of the system, and the associated mechanical anisotropy to be characterized as a function of frequency. Some rheological data for mechanically aligned DNA is presented. These observations may have implications for high DNA concentrations in vivo. Several improvements to microrheology measurements have been described, with a particular focus on high modulus anisotropic materials. Many intracellular materials appear anisotropic, and the approach presented allows the elucidation of the mechanical properties of the materials, potentially in vivo Tseng et al. (2004). Higher modulus materials should be able to be explored using fluorescent probes with diameters in the 10-100nm range. It should be possible to extend the technique to ‘2-point’ Crocker et al. (2000) cross-correlation techniques to ensure that probe effects on the network McGrath et al. (2000); Valentine et al. (2004) can be neglected Gardel et al. (2003); Levine and Lubensky (2000). We anticipate that the technique described will also be of use in determining the anisotropic pore sizes of materials which form an anisotropic network.
###### Acknowledgements.
We would like to thank Pietro Cicuta, Tim Hosey, Heather Houghton, Mark Krebs and Salman Rogers for valuable discussions. This work was supported by BBSRC.
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# Neutral and ionic dopants in helium clusters: interaction forces for the Li₂(a³Σᵤ⁺)-He and Li₂⁺(X²Σ_g⁺)-He complexes.
## I Introduction
The accurate evaluation of the interaction potential between a dopant molecule or ion and a single $`\mathrm{He}`$ atom is an important step in the study of small-size doped helium clusters because it can be used to set up the total potential acting within them, an essential first step in determining their structure and dynamics. Nano-sized helium droplets in which a molecular or an atomic impurity has been added by means of picking up techniques provide, indeed, a unique environment for high precision spectroscopic studies of the molecular solute Toennies and Vilesov (2004) especially for weakly bound molecular species Stienkenmeier et al. (1996a); Higgins et al. (1998); Stienkemeier and Vilesov (2001) that are otherwise difficult to analyze. Among the many dopants that have been experimentally and theoretically analyzed, the alkali metal atoms and dimers have many interesting properties: from experimental Stienkenmeier et al. (1996a, b); Stienkemeier and Vilesov (2001) and theoretical Dalfovo (1994); Bodo et al. (2004a, 2005) evidence it turns out that they normally reside on the surface of an $`{}_{}{}^{4}\mathrm{He}`$ droplet, usually forming a slight “dimple” on that surface. When two or more alkali atoms are attached to the droplet they eventually meet on its surface forming a singlet or a triplet molecule. Spectroscopic identification of the molecules shows Higgins et al. (1998) that dimers in the triplet state outnumber the singlet molecules by a factor up to 10<sup>5</sup>. This is probably due to the differences in stability between the triplet dimers and the singlet molecules Higgins et al. (1998) whereby the energy release due to the formation of a singlet molecule is more likely to lead to the detachment of the molecule and to the almost complete disappearance of this kind of dopant. A similar behaviour, i.e. one in which high spin states are preferred, has been noticed also in the formation of small clusters of alkali atoms grown on helium droplet surfaces Schulz et al. (2004): the very small binding energies to the surface allows for a spontaneous “selection” of the cluster in higher spin states (it is worth pointing out that, in these experiments, lithium atoms were shown to form only dimers while heavier alkalis produced also larger structures). Very recently experiments have been also performed on alkali atoms in $`{}_{}{}^{3}\mathrm{He}`$ droplets Mayol et al. (2004) where both theoretical predictions and experiments Stienkenmeier et al. (2004) have shown a very similar behavior to that of the bosonic moieties. The possible formation of cold heteronuclear molecules has attracted the attention of researchers in the last few years (see for example the special volume edition in which Ref. Mudrich et al. (2004) can be found) because of the possibility of studying the effects of polar interactions in cold or even quantum degenerate molecular samples: helium droplets may therefore allow the spectroscopic study of weakly bound (triplet) heteronuclear dimers.
In all the experiments mentioned above, the helium droplet is used as a convenient, but inert matrix in which weakly bound molecular species are accumulated and studied. Helium droplets, however, can also be used as microscopic cryo-reactors which actively initiate series of chemical reactions. Very recently, in fact, it has been shown possible to follow chemical reactions inside a helium droplet Farnik and Toennies (2005) using, among others, simple diatomic dopants. The chain of reactions is usually initiated by the ionization of one of the helium atoms Farnik and Toennies (2005) with an ensuing charge transfer Xie et al. (2003); Scifoni et al. (2004a, b) process that moves the charge onto the impurity because its ionization potential is usually lower than He. The impurity gives rise, in turn, to a variety of further, secondary chemical reactions: the excess energy due to any exothermic reaction is rapidly thermalized within the droplet itself by helium atom evaporation. The rapid quenching of the internal degrees of freedom of the partner species afforded by the superfluid environment of the droplets has a profound effect on the various branching reactions and, in some cases, can ultimately stabilize different intermediate complexes which would instead only remain as transient species in the gas phase. The latter condition therefore becomes an interesting feature of such nanoscopic cryostats in the sense that any further step initiated by ionization of the cluster can follow specific pathways which are likely to be different from those under the corresponding gas phase conditions.
Concerning the experiments involving neutral dopants, we believe that the microscopic structure of the helium surface-dimer complex is a key piece of information for the full understanding of the experimental data. In this work we continue the analysis of the lithium-He system by presenting an accurate Potential Energy Surface (PES) for the dimer in the triplet state $`\mathrm{Li}_2(^3\mathrm{\Sigma }_u^+)`$ and one Helium atom.
When moving to the problem of ionic dopants, and therefore to the fascinating possibility of conducting ionic reactions at low temperatures in the droplet, the study of the microscopic structure of the solvation site around the ionic molecule may also help the understanding of the possible reactive pathways which are eventually followed. This is the reason why, together with the previously mentioned PES, we also want to present here similar calculations on the $`\mathrm{Li}_2^+(^2\mathrm{\Sigma }_\mathrm{g}^+)\mathrm{He}`$ system. These two PES’s may be considered as the building blocks of the complex interactions at play in the He droplets or in smaller helium clusters: they describe two rather extreme situations of either a weak, dispersion-type interaction or of a much stronger, markedly orientational ionic PES. Although the forces at play in the large systems may be extremely difficult to compute, a simplified approach where the total potential is obtained as a sum of two-body potentials represents a realistic route because the three-body forces and the higher order terms, especially for the neutral system, provide fairly negligible contributions to the total energies. For their possible relevance in ionic systems see e.g. Refs. Sebastianelli et al. (2004) and Lenzer et al. (2001) for further details.
## II The ab-initio calculations
Both surfaces have been calculated using the MP4(SDTQ) method (without freezing the lithium core) with the cc-pV5Z basis set employed within Gaussian03 M. J. Frisch et al. (1998). All energies have been corrected for BSSE using the Counterpoise procedure Boys and Bernardi (1970). Since, in either case, the distortion due to the helium atom over the electronic structure of the molecule is rather weak, we have decided to keep the molecule at its equilibrium geometry. The optimal geometry of the molecular species was determined by a separate MP4 optimization on both the triplet Li<sub>2</sub> and ionic $`\mathrm{Li}_2^+`$. The internuclear distances that we have used are $`r=4.175`$ Å and $`r=3.11`$ Å for Li<sub>2</sub> and Li$`{}_{}{}^{+}{}_{2}{}^{}`$ respectively. For both surfaces we have used a grid of Jacobi coordinates with different radial geometries and 19 angles from 0 to 90 . For the Li$`{}_{2}{}^{+}`$He case we have calculated for each angle a number of points between 50 and 64 ($`2.5R12.0`$ Å ) for a total number of 1,135. For the Li$`{}_{2}{}^{}`$He system the number of radial points for each angle varied from 30 to 80 ($`2.0R14.0`$ Å ) for a total of 847 points.
## III The Fitting Procedure
In order to make it easier to employ the two potential energy surfaces in further studies,we have decided to obtain a suitable analytical fitting of the raw points using a non-linear fitting procedure based on the minimization of the square deviation, as obtained by means of the efficient Levenberg-Marquadt method Press et al. (1986).
The full interaction can be written as
$$V_{tot}=V(R_a,\theta _a)+V(R_b,\theta _b)+V_{LR}(R,\theta )$$
(1)
where the coordinates are those reported for clarity in Figure 1.
The first two contributions represent the anisotropic interactions at short range and are written in terms of Legendre polynomials
$$V(R_a,\theta _a)=\underset{n=0}{\overset{nmax}{}}\underset{l=0}{\overset{lmax}{}}R_a^n\mathrm{exp}(\beta R_a)P_l(cos\vartheta _a)C_{nl}^a$$
(2)
where, given the symmetry of the homonuclear molecule $`C_{nl}^a=C_{nl}^b`$. The long-range contribution is instead expressed in Jacobi coordinates and given by
$$V_{LR}(R,\theta )=\underset{N}{}\underset{L=0}{\overset{N4}{}}\frac{f_N(\beta R)}{R^N}P_L(cos\vartheta )C_{NL}^{LR}$$
(3)
where the $`f_N`$ damping functions are those defined within the well-known Tang-Toennies empirical potential modeling U. Kleinekath fer et al. (1996).
For the neutral system (Li<sub>2</sub>-He) we have employed $`nmax=4`$ and $`lmax=6`$ and we have limited the long range expression to the value $`N=6`$, i.e. the latter was simply given by two terms
$$V_{LR}=f_6(\beta R)R^6[P_0(cos\theta )C_{60}^{LR}+P_2(cos\theta )C_{62}^{LR}].$$
This very simple analytical representation of the entire set of unweighted 847 ab-initio points yielded a standard deviation of 0.0046 cm<sup>-1</sup>. For the ionic system (Li$`{}_{}{}^{+}{}_{2}{}^{}`$-He) we have employed a more flexible representation given by $`nmax=8`$ and $`lmax=4`$ and two long range anisotropic terms: $`C_4`$ and $`C_6`$. The latter terms were employed in order to represent correctly the long range multipolar expansion of a neutral, polarizable atom interacting with a point-like charge ion. The long range expansion therefore becomes:
$$V_{LR}=f_4(\beta R)R^4[P_0(cos\theta )C_{40}^{LR}]+f_6(\beta R)R^6[P_0(cos\theta )C_{60}^{LR}+P_2(cos\theta )C_{62}^{LR}].$$
The total number of 1,135 points was reduced to 1,065 by excluding the repulsive energies with values above 5,000 cm<sup>-1</sup>. A simple weight function $`W(E)`$ was used in order to obtain a better representation of the interaction region and to reduce the importance of the repulsive geometries. Hence we drop the weight of the points in the less important highly repulsive regions of the interaction
$$W(E)=\{\begin{array}{cc}1\hfill & \text{for }E1,000\hfill \\ |E|^1\hfill & \text{for }1,000E5,000\hfill \end{array}$$
The final standard deviation was 0.073 cm<sup>-1</sup>. As an example of the accuracy of the two fittings we report in Figure 2 (neutral system) and Figure 3 (ionic system) two selected cuts through the PES’s at $`\theta `$=0 and $`\theta `$=90 and compare them with the ab-initio points. As is evident from the figures, we are able to describe with good accuracy both the attractive region and the lowest section of the repulsive wall. In either case, however, we were not able to provide reliable fitting values for $`R1.0`$ Å . All the coefficients are reported in Table I and Table II. Fortran77 subroutines are available on request from the authors.
## IV Results and discussion
### IV.1 The Potential energy surfaces
A 3D representation of the two PES’s is presented in Figure 4 as isoenergetic contour plots.
As can be seen from the figure, the two PES’s are completely different in shape and in strength, as one could easily have expected. The triplet interaction is very weak (very similar to the one we have already calculated for the singlet molecule in Ref. Bodo et al. (2004a, 2005)) and has its absolute minimum at $`\theta `$=90 with an interaction energy of $`2.4`$ cm<sup>-1</sup>. This is a very different anisotropy with respect to that exhibited by the singlet dimer Bodo et al. (2004a, 2005), where the collinear configuration provided the global PES minimum energy and it also exhibits deeper well values than those of the former case. The extreme weakness of the interaction of Li<sub>2</sub> with He is indirectly supported by the experimental findings Stienkemeier and Vilesov (2001) which indicate that triplet molecules also reside on the surface of the droplet and are not efficiently solvated by the helium atoms.
Given the weakness of the calculated interaction it would be also be interesting to provide a modification of it that makes use of semi-empirical guesses for the two-body contribution due to the Li-He interactions. We can, in fact, modify the calculated interaction energies by employing 2-body semiempirical potentials along the lines which we have already followed for the singlet molecule in Ref. Bodo et al. (2005). It is, in fact, possible to produce rather simply a semi-empirical potential $`V^{}(R,\theta )`$ by using the formula
$$V^{}(R,\theta )=V(R,\theta )V^{ab}(R_a)V^{ab}(R_b)+V^{sm}(R_a)+V^{sm}(R_b)$$
where $`V^{sm}(R_{a,b})`$ is the semi-empirical potential proposed by Toennies and coworkers U. Kleinekath fer et al. (1996) for the dimer Li-He, $`V^{ab}(R_{a,b})`$ is a simple fitting of the ab-initio points calculated by us at the same level of accuracy (MP4/cc-pV5Z) for the Li-He diatomic curve and $`R_a`$, $`R_b`$ are the two Li-He distances within the trimer. Our ab-initio data, the fitting curve and the model potential of Ref. U. Kleinekath fer et al. (1996) for the Li-He pair are reported in Figure 5. The resulting full potential energy surface $`V^{}(R,\theta )`$ calculated with the model potential is stronger than the one from purely ab-initio data, although the increase in well depth is still not sufficient to allow for possible solvation of the molecular impurity inside the helium clusters (see also our similar conclusions for the singlet Li<sub>2</sub> in Ref. Bodo et al. (2005)).
The ionic interaction reported in Figures 3 and 4, on the other hand, is much stronger and presents its minimum at $`\theta `$=0 . Here the situation is completely different and, although no calculations have been carried out as yet for the cluster structures, it is not difficult to see that the Li$`{}_{}{}^{+}{}_{2}{}^{}`$ impurity would be strongly solvated in liquid helium droplets and therefore it is likely to get localized at the droplet center with a solvation shell strongly bound to it (for analogous situations one should look at the existing results on the ionic Li or Na impurities described in Refs. Galli et al. (2001); Sebastianelli et al. (2004)).
### IV.2 Estimating the 3-body effects in the ionic system
As we have pointed out in the introduction, in order to treat small and medium sized clusters with 2-50 helium atoms, one should be able to rapidly evaluate the total potential energy acting within them. One of the most used and successful approaches consists in approximating the total interaction as a sum of 2-body terms, initially neglecting any non-separable 3-body contribution. While this has been proved to be a very accurate procedure for doped helium clusters with a neutral impurity, it may represent a source of error in ionic clusters where non-separable interactions among the induced multipoles in the first solvation shell may be important. However, it has been recently shown in various works on ionic dopants in rare-gas clusters including helium that these effects are small and should not, when included, alter substantially the geometries or the energies of the clusters (for a general discussion with anionic dopants see refs. Lenzer et al. (2001) and Sebastianelli et al. (2003a, b), for positive ions see Refs. Sebastianelli et al. (2004); Bodo et al. (2004c)).
In order to verify the applicability of the sum-of-potentials approximation we have carried out at a consistent level (cc-pV5Z/MP4) a series of ab-initio calculation on the Li$`{}_{}{}^{+}{}_{2}{}^{}`$He<sub>2</sub> system for which such 3-body effects may arise. The analysis carried out here is analogous to what we have already reported in Ref. Bodo et al. (2004c) for the LiH<sup>+</sup> dopant. The 3-body forces should mainly originate from the induced multipoles (attractive and repulsive contributions) and by the weaker Axilrod-Teller effects Lenzer et al. (2001). For our preliminary study, we have therefore chosen four different geometries of the complex where the two helium atoms are close enough to the ion molecule in order to contribute significantly to 3-body forces. The first geometry (geometry A) is weakly repulsive ($`70`$ cm<sup>-1</sup>) with the two helium atoms at 3.11 Å and is the upper one reported in the inset of Table 3. The second (B) and the third (C) geometries are similar in shape to the second one sketched in the same table, but they differ for the values of the distances between the atoms. The last one (geometry opt) is also similar in shape to the latter, but it comes instead from a full minimization at the MP2/cc-pv5Z level of the entire complex: all the relevant distances are reported in Table 3.
The analysis of 3-body forces is done in the following way: for each geometry we calculate the interaction energy of the entire Li$`{}_{}{}^{+}{}_{2}{}^{}`$He<sub>2</sub> complex, of one of its 2-body fragments (there are two identical ones in symmetrical geometries) Li$`{}_{}{}^{+}{}_{2}{}^{}`$He and of the remaining He<sub>2</sub> dimer by using the following formulae where the geometry is fixed, it is the same in each fragment and the counterpoise correction is used for each interaction energy:
$`V[\mathrm{Li}_2^+\mathrm{He}_2]`$ $`=`$ $`E[\mathrm{Li}_2^+\mathrm{He}_2]E[\mathrm{Li}_2^+]2E[\mathrm{He}]`$
$`V[\mathrm{Li}_2^+\mathrm{He}]`$ $`=`$ $`E[\mathrm{Li}_2^+\mathrm{He}]E[\mathrm{Li}_2^+]E[\mathrm{He}]`$
$`V[\mathrm{He}_2]`$ $`=`$ $`E[\mathrm{He}_2]2E[\mathrm{He}]`$ (4)
the residual 3-body interaction is simply calculated by using the expression
$$V[\mathrm{Li}_2^+\mathrm{He}_2]2V[\mathrm{Li}_2^+\mathrm{He}]V[\mathrm{He}_2].$$
(5)
As can be seen from Table 3, the 3-body interaction is always a small percentage of the total interaction in the various geometries except for geometry A where it represents more than 10% of the total interaction. It is however important to note that the geometry A corresponds to a repulsive geometry for the Li$`{}_{}{}^{+}{}_{2}{}^{}`$-He fragment and therefore it is not relevant for optimization purposes.
## V Conclusions
We have computed two accurate Potential Energy Surfaces for two different systems that are of interest for the experimental and theoretical study of helium droplets doped with alkali metal molecules. The two molecules considered here are the triplet state of Li<sub>2</sub> and the ground state ionic Li$`{}_{}{}^{+}{}_{2}{}^{}`$. As should be expected, we found the two PES’s to be markedly different:
* the interaction of the neutral moiety is similar to the one we have already studied for ground state singlet Li<sub>2</sub> Bodo et al. (2004a, b): a very weak interaction that confirms once more the tendency of high spin compounds of alkali metals to reside on the surface of helium droplets. A modified version of the same interaction that gives rise to slightly deeper potentials can be obtained following our earlier proposal in Ref. Bodo et al. (2005) and also found to yield weaker interaction potentials than those between He partners.
* the triplet dimer interaction with helium is showing here a markedly different anisotropy from that found earlier on for the singlet state of Li<sub>2</sub>: the minimum energy configuration is, in fact, given by T-shaped structures as opposed to the linear structures obtained for the singlet interaction.
* the ionic interaction, instead, is much stronger and more orientation-dependent: it should therefore lead to full solvation of the molecular moiety inside the droplets.
* a preliminary analysis of the three-body effects on the interactions with more He atoms indicates that such effects are relatively small and should therefore allow the use of an approximate description of the full potential energy landscapes in $`\mathrm{Li}_2^+\mathrm{He}_\mathrm{n}`$ clusters in terms of two-body potentials.
Both surfaces have been fitted using rather simple analytical expression reported in the present paper and which therefore provide working quality interactions potentials for the title systems. We believe that such potentials are an important step in the modelling of the larger clusters behavior because they may be used to set up the total interaction for much larger systems whenever using the sum-of-potential approach is found to be a realistic alternative: our present study and preliminary analysis do seem to suggest that this may be the case for both the present system. We shall verify such a possibility via our ongoing calculations for the doped $`{}_{}{}^{4}\mathrm{He}`$ clusters.
###### Acknowledgements.
We acknowledge financial support from Rome “La Sapienza” Scientific Committee, the CASPUR supercomputing Center, the MUIR National Projects FIRB and PRIN. We also acknowledge support from the INTAS grant 03-51-6170.
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# The spin-dependent semiconductor Bloch equations: a microscopic theory of Bir-Aronov-Pikus spin-relaxation
## Abstract
Semiconductor Bloch equations, in their extension including the spin degree of freedom of the carriers, are capable to describe spin dynamics on a microscopic level. In the presence of free holes, electron spins can flip simultaneously with hole spins due to electron-hole exchange interaction. This mechanism named after Bir, Aronov and Pikus, is described here by using the extended semiconductor Bloch equations Rössler (2002) and considering carrier-carrier interaction beyond the Hartree-Fock truncation. As a result we derive microscopic expressions for spin-relaxation and spin-dephasing rates.
Semiconductor Bloch equations (SBE) are a well established concept to describe the dynamics of carriers in a semiconductor or quantum structure by a scalar light field.Lindberg and Koch (1988); Haug and Koch (1993); Khitrova et al. (1999) It has been used successfully in modelling the time evolution due to carrier-carrier interaction on different time scales including the *coherent* and the *relaxation* regime.Binder and Koch (1995) Thus, SBE have become the dominating tool in the theory of semiconductor lasers and in designing even the complex structures of *vertical cavity surface emitting lasers* (VCSEL). One phenomenon, however, connected with VCSELs points to a deficiency of the SBE: these laser structures are known for their polarization instability, i.e., the uncontrolled switching of the laser modes between the two possible transverse polarizations.Choquette et al. (1993); San Miguel et al. (1995); Ando et al. (1998) In addition, the investigation of semiconductor quantum structures as model systems for coupled Rabi oscillations with electrons, heavy- and light-hole states (each spin-degenerate) required to extend the two-level SBE to six-level SBE and to take into account the polarization degree of freedom of the exciting light.Binder and Lindberg (2000)
More recently, the carrier spin and its dynamics have gained much interest in the field of spintronics.Žutić et al. (2004) Spin dynamics in semiconductors Meier and Zakharchenya (1984) and quantum structures,Sham (1993); Awschalom et al. (2002) formulated so far in a more phenomenological way, is ruled by different mechanisms: one of which is related to the electron-hole exchange interaction.Denisov and Makarov (1973) It becomes relevant if the semiconductor system contains besides electrons also holes (e.g., due to doping or high excitation). This *Bir-Aronov-Pikus* (BAP) mechanism, originally considered for bulk semiconductors,Bir et al. (1975) has been described also for semiconductor quantum structures,Sham (1993); Maialle and Degani (1997) but never by a rigorous microscopic treatment of the spin dynamics. In this perspective, the SBE have been formulated for the six-level system,Rössler (2002) considering spin-splitting of the electronic energies due to spin-orbit coupling caused by bulk inversion (BIA) Dresselhaus (1955) or structure inversion asymmetry (SIA).Rashba (1960) These extended SBE were designed only within the *Hartree-Fock truncation* leading to the coherent regime, thus neglecting scattering processes, responsible for relaxation. Recently, we have used the extended optical Bloch equations (SBE without carrier-carrier interaction Schäfer and Wegener (2002)) to provide a microscopic approach to the longitudinal ($`T_1`$) and transverse ($`T_2`$) relaxation times due to electron-phonon interaction.Lechner and Rössler (2004) In this approach we have considered scattering between electrons and phonons in second order Born approximation to provide a microscopic formulation for the D’yakonov Perel’ (DP) mechanism of spin-relaxation.D’yakonov and Perel’ (1971) The analogous concept is applied here to the electron-hole exchange interaction and yields the microscopic formulation of spin-relaxation due to the BAP mechanism.
In the following, we concentrate on spin dynamics in a semiconductor quantum well (QW) under excitation with circularly polarized light leading to a nonequilibrium spin distribution due to optical selection rules. Let the system be described by the Hamiltonian
$$=_0+_{\text{light}}+_{\text{coul}},$$
(1)
where $`_0`$ is the kinetic part including BIA and SIA spin-orbit coupling, $`_{\text{light}}`$ the interaction with the exciting light field, and $`_{\text{coul}}`$ the Coulomb interaction between the carriers. We adopt the notation of our previous works Rössler (2002); Lechner and Rössler (2004) and use the basis in which the kinetic part of the Hamiltonian for the six-level system is *diagonal*
$`_0`$ $`=`$ $`{\displaystyle \underset{𝐤^{}m_c^{}}{}}ϵ_{m_c^{}}(𝐤^{})c_{m_c^{}}^{}(𝐤^{})c_{m_c^{}}(𝐤^{})`$ (2)
$`+{\displaystyle \underset{𝐤^{}m_v^{}}{}}ϵ_{m_v^{}}(𝐤^{})v_{m_v^{}}(𝐤^{})v_{m_v^{}}^{}(𝐤^{}).`$
Here, $`c_{m_c}(𝐤)`$ $`\left[v_{m_v}(𝐤)\right]`$ are fermion operators for electrons (light- and heavy-holes) with spin quantum numbers $`m_c=\pm 1/2`$ ($`m_v=\pm 1/2,\pm 3/2`$) defined with respect to the in-plane wave vector $`𝐤`$. The time dependence of the operators is understood. The single particle energies $`ϵ_{m_c^{}}(𝐤^{})`$ $`\left[ϵ_{m_v^{}}(𝐤^{})\right]`$ describe subbands, which are spin-split due to spin-orbit interaction. In dipole approximation, the interaction with the light field reads
$`_{\text{light}}`$ $`=`$ $`{\displaystyle \underset{\begin{array}{c}m_c^{}m_v^{}\\ 𝐤^{}\end{array}}{}}[𝐄(t)𝐝_{m_c^{}m_v^{}}(𝐤^{})c_{m_c^{}}^{}(𝐤^{})v_{m_v^{}}^{}(𝐤^{})`$ (3)
$`+𝐄^{}(t)𝐝_{m_c^{}m_v^{}}^{}(𝐤^{})v_{m_v^{}}(𝐤^{})c_{m_c^{}}(𝐤^{})],`$
where $`𝐄(t)`$ is the electric field vector and $`𝐝_{m_c^{}m_v^{}}(𝐤^{})`$ is the dipole matrix element connecting valence and conduction band states (for details see Ref. Rössler, 2002).
The carrier-carrier interaction can be split up into four parts
$`_{\text{coul}}=_{\text{ee}}+_{\text{hh}}+_{\text{eh}}^\text{C}+_{\text{eh}}^\text{X}.`$ (4)
Here, $`_{\text{ee}}`$ ($`_{\text{hh}}`$) describes the Coulomb interaction between electrons (holes) in the conduction (valence) band. The remaining terms account for electron-hole interaction, the *direct* Coulomb term $`_{\text{eh}}^\text{C}`$ and the *exchange* term $`_{\text{eh}}^\text{X}`$.Denisov and Makarov (1973) In the frame of SBE,Haug and Koch (1993); Khitrova et al. (1999); Binder and Koch (1995) especially for the coherent regime, carrier-carrier interaction has been considered so far only with respect to renormalization of the single-particle energies and of the interaction with the light field, while the electron-hole exchange has been ignored. However, it is just $`_{\text{eh}}^\text{X}`$ which can cause spin-flips, thus contributing to the spin-dynamics due to the BAP mechanism and, hence, is of interest here. As derived in Ref. Rössler, 2002, the exchange term reads
$`_{\text{eh}}^\text{X}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}m_cm_c^{}\\ m_vm_v^{}\end{array}}{}}{\displaystyle \underset{\begin{array}{c}𝐤𝐤^{}\\ 𝐪\end{array}}{}}𝒱_{m_cm_vm_c^{}m_v^{}}^\text{X}(𝐤,𝐤^{},𝐪)c_{m_c}^{}(𝐤+𝐪)c_{m_c^{}}(𝐤^{}+𝐪)v_{m_v^{}}^{}(𝐤)v_{m_v}(𝐤^{}).`$ (5)
The detailed form of the interaction matrix element $`𝒱_{m_cm_vm_c^{}m_v^{}}^X(𝐤,𝐤^{},𝐪)`$ will not become important in the following. But we emphasize, that the structure of this matrix element makes simultaneous flips of electron and hole spins possible, which, in the electron system, finally contribute to spin-relaxation.Rössler (2002)
While the dynamics of the whole system is contained in the equations of motion (EOM) of the $`6\times 6`$ density matrix, we concentrate here on the dynamics of the electron spins by looking at the EOM of the $`2\times 2`$ density matrix for the electron subsystem
$$\mathit{\varrho }^{(m_c\overline{m}_c)}(𝐤)=\left(\begin{array}{cc}\varrho _{m_cm_c}(𝐤)& \varrho _{m_cm_c}(𝐤)\\ \varrho _{m_cm_c}(𝐤)& \varrho _{m_cm_c}(𝐤)\end{array}\right).$$
(6)
The single entries are expectation values of products of a creation and an annihilation operator $`\varrho _{m_c\overline{m}_c}(𝐤)=c_{m_c}^{}(𝐤)c_{\overline{m}_c}(𝐤)`$. Their EOM read
$`i\mathrm{}_t\varrho _{m_c\overline{m}_c}(𝐤)`$ $`=`$ $`\left[ϵ_{m_c}(𝐤)ϵ_{\overline{m}_c}(𝐤)\right]\varrho _{m_c\overline{m}_c}(𝐤)+{\displaystyle \underset{m_v}{}}\left[𝐄(t)𝐝_{\overline{m}_cm_v}^{cv}P_{m_cm_v}(𝐤)𝐄^{}(t)𝐝_{m_cm_v}^{cv}P_{\overline{m}_cm_v}^{}(𝐤)\right]`$ (7)
$`{\displaystyle \underset{\overline{𝐤}𝐪}{}}{\displaystyle \underset{\begin{array}{c}m_c^{}\\ \stackrel{~}{m}_v\stackrel{~}{m}_v^{}\end{array}}{}}[𝒱_{\overline{m}_c\stackrel{~}{m}_vm_c^{}\stackrel{~}{m}_v^{}}^\text{X}(𝐤+𝐪,\overline{𝐤},𝐪)c_{m_c}^{}(𝐤)c_{m_c^{}}(\overline{𝐤}+𝐪)v_{\stackrel{~}{m}_v}^{}(𝐤+𝐪)v_{\stackrel{~}{m}_v^{}}(\overline{𝐤})`$
$`𝒱_{m_c^{}\stackrel{~}{m}_vm_c\stackrel{~}{m}_v^{}}^\text{X}(\overline{𝐤},𝐤+𝐪,𝐪)c_{m_c^{}}^{}(\overline{𝐤}+𝐪)c_{\overline{m}_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(\overline{𝐤})v_{\stackrel{~}{m}_v}(𝐤+𝐪)],`$
where we have introduced the interband polarization $`P_{m_cm_v}(𝐤)=c_{m_c}^{}(𝐤)v_{m_v}^{}(𝐤)`$.Haug and Koch (1993) Due to the many-body contributions, the dynamics of $`\varrho _{m_c\overline{m}_c}(𝐤)`$ are ruled by four-point density matrices and, consequently, we run into a *hierarchy problem*, which can be solved by an appropriate truncation. The *Hartree-Fock* (HF) truncation scheme Haug and Koch (1993), as used in Ref. Rössler, 2002, factorizes the expectation values of the four-operator terms into a product of two-operator terms under the condition that they are macroscopic, namely, either electron (hole) densities or polarizations. While closing the hierarchy and renormalizing the eigenenergies and the dipole interaction the HF truncation limits the EOM to the *coherent* regime, because no scattering processes are taken into account. In order to include these processes, which are essential for spin-relaxation and -dephasing, we go beyond the HF truncation by considering the *reduced four-operator terms* Haug and Koch (1993), defined as the difference between the expectation value of the untruncated four-operator term and its HF truncated product. For $`c_{m_c^{}}^{}(\overline{𝐤}+𝐪)c_{\overline{m}_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(\overline{𝐤})v_{\stackrel{~}{m}_v}(𝐤+𝐪)`$ \[see Eq. (7)\] it reads
$`\delta c_{m_c^{}}^{}(\overline{𝐤}+𝐪)c_{\overline{m}_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(\overline{𝐤})v_{\stackrel{~}{m}_v}(𝐤+𝐪)`$ $`=`$ $`c_{m_c^{}}^{}(\overline{𝐤}+𝐪)c_{\overline{m}_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(\overline{𝐤})v_{\stackrel{~}{m}_v}(𝐤+𝐪)`$ (8)
$`c_{m_c^{}}^{}(𝐤)c_{\overline{m}_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(𝐤+𝐪)v_{\stackrel{~}{m}_v}(𝐤+𝐪)\delta _{𝐤,\overline{𝐤}+𝐪}.`$
The scattering contributions are found by solving the EOM of the reduced four-operator terms which contain the *complete* information about the scattering in expectation values of four- and six-operator terms. In analogy to the case of electron-phonon scattering Lechner and Rössler (2004) we truncate these terms by factorizing them into their macroscopic parts and taking into account only those, which contribute in *second order Born approximation*. After integrating the arising equations and applying the *adiabatic* and the *Markov* approximation Kuhn and Rossi (1992), we achieve a closed set of equations for the reduced four-operator terms, which can be solved and used in Eq. (7) (for technical details see Ref. Lechner and Rössler, 2004). Thus, the EOM for the diagonal entries of the $`2\times 2`$ density matrix due to the electron-hole exchange scattering can be cast into the form
$`_t\varrho _{m_cm_c}(𝐤)\mathbf{|}_\text{X}`$ $`=`$ $`\mathrm{\Gamma }_{m_cm_c}^{\text{out}\text{X}}(𝐤)\varrho _{m_cm_c}(𝐤)`$ (9)
$`+\mathrm{\Gamma }_{m_cm_c}^{\text{in}\text{X}}(𝐤)[1\varrho _{m_cm_c}(𝐤)],`$
with $`\mathrm{\Gamma }_{m_cm_c}^{\text{out}\text{X}}(𝐤)`$ $`\left[\mathrm{\Gamma }_{m_cm_c}^{\text{in}\text{X}}(𝐤)\right]`$ accounting for the exchange scattering out of (into) the state with spin $`m_c`$ at wave vector $`𝐤`$. The derivation of the scattering contributions for the different four-operator terms in Eq. (7) follows the same scheme. Thus, we present here only the results for the reduced four-operator term $`c_{m_c^{}}^{}(\overline{𝐤}+𝐪)c_{m_c}(𝐤)v_{\stackrel{~}{m}_v^{}}^{}(\overline{𝐤})v_{\stackrel{~}{m}_v}(𝐤+𝐪)`$. The corresponding out-scattering rate $`\mathrm{\Gamma }_{m_cm_c}^{\text{out}\text{X}}(𝐤)`$ reads
$`\mathrm{\Gamma }_{m_cm_c}^{\text{out}\text{X}}(𝐤)`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \underset{\overline{𝐤}𝐪}{}}{\displaystyle \underset{\begin{array}{c}m_c^{}\\ \stackrel{~}{m}_v\stackrel{~}{m}_v^{}\end{array}}{}}|𝒱_{m_c\stackrel{~}{m}_v^{}m_c^{}\stackrel{~}{m}_v}^\text{X}(𝐤+𝐪,\overline{𝐤},𝐪)|^2\delta [ϵ_{m_c}(𝐤)ϵ_{\stackrel{~}{m}_v}(𝐤𝐪)ϵ_{m_c^{}}(\overline{𝐤}+𝐪)+ϵ_{\stackrel{~}{m}_v^{}}(\overline{𝐤})]`$ (10)
$`\times [1\varrho _{\stackrel{~}{m}_v^{}\stackrel{~}{m}_v^{}}(\overline{𝐤})][1\varrho _{m_c^{}m_c^{}}(\overline{𝐤}+𝐪)]\varrho _{\stackrel{~}{m}_v\stackrel{~}{m}_v}(𝐤𝐪)\varrho _{m_cm_c}(𝐤),`$
with a similar expression for the in-scattering rate $`\mathrm{\Gamma }_{m_cm_c}^{\text{in}\text{X}}(𝐤)`$. These expressions represent all contributions to electron-hole scattering by Coulomb exchange interaction in second order Born approximation. It is important to note that without a macroscopic occupation of hole states (by doping or optical excitation) this scattering rate vanishes: holes are required for the mutual spin flips of the BAP mechanism.
The EOM for the off-diagonal entry of the density matrix the expressions can be written in the form
$`_t\varrho _{m_cm_c}(𝐤)\mathbf{|}_\text{X}`$ $`=`$ $`{\displaystyle \frac{1}{i\mathrm{}}}\left[\mathrm{\Sigma }_{m_cm_c}^\text{X}(𝐤)\varrho _{m_cm_c}(𝐤){\displaystyle \underset{\overline{𝐤}𝐪}{}}{\displaystyle \underset{\begin{array}{c}m_c^{}\\ \stackrel{~}{m}_v\stackrel{~}{m}_v^{}\end{array}}{}}\overline{\mathrm{\Sigma }}_{m_cm_c}^\text{X}(\overline{𝐤}+𝐪)\varrho _{m_cm_c}(\overline{𝐤}+𝐪)\right].`$ (11)
As for the electron-phonon scattering Lechner and Rössler (2004), the first self-energy term in Eq. (11) is proportional to the absolute squared value of the interaction matrix element $`𝒱_{m_c\stackrel{~}{m}_v^{}m_c^{}\stackrel{~}{m}_v}^\text{X}(𝐤+𝐪,\overline{𝐤},𝐪)`$ and can be split up into real and imaginary part connected by *Kramers-Kronig* transformation, where the imaginary part
$`\mathrm{}\{\mathrm{\Sigma }_{m_cm_c}^\text{X}(𝐤)\}`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{}}}{\displaystyle \underset{\overline{𝐤}𝐪}{}}{\displaystyle \underset{\begin{array}{c}m_c^{}\\ \stackrel{~}{m}_v\stackrel{~}{m}_v^{}\end{array}}{}}|𝒱_{m_c\stackrel{~}{m}_v^{}m_c^{}\stackrel{~}{m}_v}^\text{X}(𝐤+𝐪,\overline{𝐤},𝐪)|^2\delta [ϵ_{m_c^{}}(\overline{𝐤}+𝐪)ϵ_{m_c}(𝐤)+ϵ_{\stackrel{~}{m}_v}(𝐤+𝐪)ϵ_{\stackrel{~}{m}_v^{}}(\overline{𝐤})]`$ (12)
$`\times \{[1\varrho _{\stackrel{~}{m}_v^{}\stackrel{~}{m}_v^{}}(\overline{𝐤})][1\varrho _{m_c^{}m_c^{}}(\overline{𝐤}+𝐪)]\varrho _{\stackrel{~}{m}_v\stackrel{~}{m}_v}(𝐤𝐪)`$
$`+[1\varrho _{\stackrel{~}{m}_v\stackrel{~}{m}_v}(𝐤𝐪)]\varrho _{\stackrel{~}{m}_v^{}\stackrel{~}{m}_v^{}}(\overline{𝐤})\varrho _{m_c^{}m_c^{}}(\overline{𝐤}+𝐪)\}`$
accounts for dephasing due to scattering, while the real part contributes to the renormalization of the eigenenergies. However, the real and imaginary part of the second term $`\overline{\mathrm{\Sigma }}_{m_cm_c}^X(\overline{𝐤}+𝐪)`$ in Eq. (11) are not connected by Kramers-Kronig theorem, because they are proportional to a product of two complex valued exchange interaction matrix elements. Nevertheless it is possible to sort out two parts, one proportional to principal values and one proportional to the energy conserving $`\delta `$-functions. In order to derive the expressions corresponding to the ones given in Ref. Lechner and Rössler, 2004, we will present here the part proportional to $`\delta `$-functions denoted as $`\overline{\mathrm{\Gamma }}_{m_cm_c}^\text{X}(\overline{𝐤}+𝐪)`$
$`\overline{\mathrm{\Gamma }}_{m_cm_c}^\text{X}(\overline{𝐤}+𝐪)`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{}}}𝒱_{m_c\stackrel{~}{m}_v^{}m_c^{}\stackrel{~}{m}_v}^\text{X}(𝐤+𝐪,\overline{𝐤},𝐪)𝒱_{m_c\stackrel{~}{m}_v^{}m_c^{}\stackrel{~}{m}_v}(𝐤+𝐪,\overline{𝐤},𝐪)`$ (13)
$`\times \delta [ϵ_{m_c^{}}(\overline{𝐤}+𝐪)ϵ_{m_c}(𝐤)ϵ_{\stackrel{~}{m}_v^{}}(\overline{𝐤})+ϵ_{\stackrel{~}{m}_v}(𝐤+𝐪)]`$
$`\times \{[1\varrho _{\stackrel{~}{m}_v\stackrel{~}{m}_v}(𝐤+𝐪)][1\varrho _{m_cm_c}(𝐤)]\varrho _{\stackrel{~}{m}_v^{}\stackrel{~}{m}_v^{}}(\overline{𝐤})`$
$`+[1\varrho _{\stackrel{~}{m}_v^{}\stackrel{~}{m}_v^{}}(\overline{𝐤})]\varrho _{\stackrel{~}{m}_v\stackrel{~}{m}_v}(𝐤+𝐪)\varrho _{m_cm_c}(𝐤)\}.`$
With the expressions given in Eqs. (9), (11),(12) and (13) we can define the $`T_1`$ and $`T_2`$ times of the Bloch equations as in Ref. Lechner and Rössler, 2004, with the $`T_1`$ time being ruled by the in- and out-scattering terms of Eq. (9) while the $`T_2`$ time is given by the imaginary part of the self-energies of Eq. (11). The interpretation of these times with respect to spin-relaxation and -dephasing remains the same as in Ref. Lechner and Rössler, 2004 apart from the fact that the spin-relaxation mechanism is different.
In conclusion, we have derived microscopic expressions for the scattering rates due to electron-hole exchange interaction in a semiconductor QW in the frame of extended SBE. As it turns out, the expressions for these rates show the same qualitative structure as found for carrier-phonon scattering Lechner and Rössler (2004). The particular scattering mechanism considered here is the one responsible for the BAP mechanism of spin relaxation. Thus, the presented results are a microscopic formulation of the BAP spin relaxation in the frame of the extended semiconductor Bloch equations.
We thankfully acknowledge financial support from the DFG via Forschergruppe 370 *Ferromagnet-Halbleiter-Nanostrukturen*.
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# Some half-BPS solutions of M-theory
## 1 Introduction
It was recently shown by Lin, Lunin and Maldacena (LLM) that half BPS-solutions of M-theory can be expressed in terms of a single function $`D`$ satisfying the 3d continuum Toda equation. Solutions to this equation are not easily found. There is however a simple class of solutions which turns out to lead to physically interesting geometries, namely separable solutions which can locally be expressed in terms of a quadratic form in a single variable and solutions of the Liouville equation.
Solutions of 11-dimensional supergravity are of great interest and have been discussed extensively in the literature (see for example - and references therein). The approach described in provides a way to study them in a unified framework. The moduli space of half-BPS solutions is of great interest from the point of view of the AdS/CFT correspondence. Already some very interesting physical questions (such as topology change ) have been addressed using this approach (in particular the relation to the phase space of free fermions).
The LLM description applies to both regular and singular solutions. For a regular solution the function $`D`$ has to satisfy specific boundary conditions. Generic solutions for $`D`$ will not satisfy these boundary conditions and will give rise to singular geometries, which may be interpreted as a sign that important degrees of freedom have been neglected.
It turns out that among the geometries corresponding to separable solutions of the Toda equation there is only one non-singular example – the Maldacena-Nunez solution, which describes the near-horizon limit of M5-branes wrapping holomorphic curves in Calabi-Yau threefolds. All the other possibilities which stem from separable solutions of the Toda equation turn out to have singularities if standard signature is imposed.
Two special cases are discussed in this note. One is the Alishahiha-Oz solution, which describes a system of intersecting M5-branes. This solution is a warped product of $`AdS_5`$ and a 6-dimensional space and as such is dual to a superconformal field theory in four dimensions. The other example discussed below is a warped product with an $`AdS_2`$ factor, which is expected to be dual to a model of superconformal quantum mechanics. This geometry describes the near-horizon limit of a system of intersecting M2-branes. Warped products with $`AdS`$ factors are of great interest from various points of view, including the AdS/CFT correspondence as well as flux compactifications to four dimensions.
This note begins by citing the relevant formulas from , after which the separable solutions of the Toda equations are described. Rather than analyzing all the cases systematically, only some general features are given and a few special cases are discussed, possibly leaving a more complete presentation for the future.
## 2 LLM solutions
The general solution of the supergravity equations of motion described by LLM has the form
$`ds_{11}^2`$ $`=`$ $`4e^{2\lambda }(1+y^2e^{6\lambda })(dt+V_idx^i)^2+4e^{2\lambda }d\mathrm{\Omega }_5^2+y^2e^{4\lambda }d\stackrel{~}{\mathrm{\Omega }}_2^2+`$ (1)
$`+`$ $`{\displaystyle \frac{e^{4\lambda }}{1+y^2e^{6\lambda }}}[dy^2+e^D(dx_1^2+dx_2^2)]`$
$`e^{6\lambda }`$ $`=`$ $`{\displaystyle \frac{_yD}{y(1y_yD)}}`$ (2)
$`V_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}ϵ_{ij}_jD,`$ (3)
where $`i,j=1,2`$. This, supplemented by the fluxes given in , provides an M-theory background preserving 16 of the original 32 supersymmetries.
The function $`D`$ which determines the solution obeys the equation
$$(_{x_1}^2+_{x_2}^2)D+_y^2e^D=0.$$
(4)
This is the 3-dimensional continuous version of the Toda equation. To obtain regular supergravity solutions one has to impose specific boundary conditions on $`D`$ so that potential singularities inherent in (1) do not occur. These conditions can be found in .
Note that the form of the ansatz is preserved under $`y`$-independent conformal transformations of the $`x_1x_2`$ plane provided $`D`$ is shifted appropriately:
$$x_1+ix_2g(x_1+ix_2),DD\mathrm{log}|g|^2.$$
(5)
The solutions examined above can be Wick-rotated (as discussed by LLM) to yield solutions of 11-dimensional supergravity which contain an $`AdS_5`$ factor and a compact 6-dimensional manifold. These solutions can be interpreted as dual to conformal field theories in four dimensions. Specifically, one can get solutions of the form of an $`AdS_5`$ warped product by performing the analytic continuation
$$\psi \tau \alpha i\rho .$$
This maps
$$\mathrm{cos}^2\alpha d\psi ^2+d\alpha ^2+\mathrm{sin}^2\alpha d\mathrm{\Omega }_3^2(\mathrm{cosh}^2\rho d\tau ^2+d\rho ^2+\mathrm{sinh}^2\rho d\mathrm{\Omega }_3^2)$$
(6)
i.e.
$$d\mathrm{\Omega }_5^2ds_{AdS_5}^2.$$
(7)
In addition one has to take
$$\lambda =\stackrel{~}{\lambda }+i\frac{\pi }{2},$$
(8)
with the remaining coordinates unchanged. For real $`\stackrel{~}{\lambda }`$ one finds a metric with the correct signature<sup>1</sup><sup>1</sup>1Since $`t`$ is now a spacelike coordinate, it is denoted by $`\chi `$ below to avoid confusion.. This way one arrives at
$$ds_{11}^2=e^{2\stackrel{~}{\lambda }}\left(4ds_{AdS_5}^2+y^2e^{6\stackrel{~}{\lambda }}d\stackrel{~}{\mathrm{\Omega }}_2^2+ds_4^2\right),$$
(9)
where
$$ds_4^2=4(1y^2e^{6\stackrel{~}{\lambda }})(d\chi +V_idx^i)^2+\frac{e^{6\stackrel{~}{\lambda }}}{1y^2e^{6\stackrel{~}{\lambda }}}[dy^2+e^D(dx_1^2+dx_2^2)]$$
(10)
and
$$e^{6\stackrel{~}{\lambda }}=\frac{_yD}{y(1y_yD)},$$
(11)
where $`D`$ satisfies the Toda equation (4). Note that due to the analytic continuation regular solutions in this case have to satisfy different boundary conditions than those for (1). The form (9) characterizes all M-theory compactifications to $`AdS_5`$ which preserve $`N=2`$ supersymmetry in four dimensions.
Instead of the Wick rotation (2), (8) one can instead do
$$yiy,x_kix_k,$$
(12)
which leads to metrics of the form of a warped product of $`AdS_2`$ with a 9-dimensional manifold:
$$ds_{11}^2=e^{2\stackrel{~}{\lambda }}\left(4ds_{S^5}^2+y^2e^{6\stackrel{~}{\lambda }}ds_{AdS_2}^2ds_4^2\right),$$
(13)
where $`ds_4^2`$ and $`e^{6\stackrel{~}{\lambda }}`$ are as above. The form (13) characterizes all M-theory compactifications to $`AdS_2`$ which preserve $`N=2`$ supersymmetry in four dimensions.
## 3 Separable solutions of the Toda equation
Separable solutions of (4) are of the form<sup>2</sup><sup>2</sup>2In a different context such solutions were discussed in .:
$$D(x_1,x_2,y)=F(x_1,x_2)+G(y).$$
(14)
Using this in the Toda equation (4) one finds that
$$G(y)=\mathrm{log}(\alpha y^2+\beta y+\gamma ),$$
(15)
while $`F`$ has to satisfy the Liouville equation:
$$(_{x_1}^2+_{x_2}^2)F+2\alpha e^F=0.$$
(16)
Here $`\alpha `$, $`\beta `$ and $`\gamma `$ are constants.
In terms of $`\xi x_1+ix_2`$ one has the well known general solution of the Liouvile equation:
$$e^F=\frac{4|f^{}(\xi )|^2}{(1+\alpha |f(\xi )|^2)^2},$$
(17)
where $`f`$ is a holomorphic function.
Thus the function $`D`$ is given by
$$D(x_1,x_2,y)=\mathrm{log}\frac{4Q(y)|f^{}(\xi )|^2}{(1+\alpha |f(\xi )|^2)^2}$$
(18)
with
$$Q(y)=\alpha y^2+\beta y+\gamma .$$
(19)
The parameters $`\alpha ,\beta ,\gamma `$ parameterize the possible solutions. It is convenient to discuss separately the following cases:
1. $`\alpha =0`$, i.e. $`Q=\beta y+\gamma `$;
2. $`\alpha <0`$ and $`Q=|\alpha |(yy_1)(yy_2)`$ for real $`y_1,y_2`$;
3. $`\alpha >0`$ and $`Q=\alpha (yy_1)(yy_2)`$ for real $`y_1,y_2`$;
4. $`\alpha >0`$ and $`Q=\alpha |yu|^2`$ for complex $`u`$.
and it will be assumed that $`y_1<y_2`$. In the first case $`\beta `$ can be scaled away by redefining the function $`f`$ in (18). In the remaining cases $`\alpha `$ can be scaled away in the same way once its sign is fixed.
Since the Toda equation is invariant under
$$\xi g(\xi )$$
(20)
for any holomorphic function $`g`$, one can locally choose a convenient canonical form for the function $`G`$. In view of this, one can take the following solutions for $`D`$:
$`D_0(x_1,x_2,y)`$ $`=`$ $`\mathrm{log}(y+\gamma ),`$ (21)
$`D_1(x_1,x_2,y)`$ $`=`$ $`\mathrm{log}{\displaystyle \frac{(yy_1)(y_2y)}{x_2^2}},`$ (22)
$`D_2(x_1,x_2,y)`$ $`=`$ $`\mathrm{log}{\displaystyle \frac{4(yy_1)(yy_2)}{(1+(x_1^2+x_2^2)^2)^2}},`$ (23)
$`D_3(x_1,x_2,y)`$ $`=`$ $`\mathrm{log}{\displaystyle \frac{4|yu|^2}{(1+(x_1^2+x_2^2)^2)^2}}.`$ (24)
## 4 M-theory solutions
Each of the solutions (21)–(24) for $`D`$ leads to an $`N=2`$ supersymmetric solution of the supergravity equations of motion. However to interpret the resulting metrics in physical terms one has to properly define the ranges of coordinates so that the signature is correct. All but one of the metrics arising from (21)–(24) are singular at some points. Such singularities are assumed to indicate that some degrees of freedom which generically decouple become light and have to be accounted for if a non-singular description is to ensue.
A case which is regular was already pointed out in , where the following solution of the Toda equation is discussed:
$$e^D=\frac{1}{x_2^2}(14y^2).$$
(25)
This leads to the regular geometry found earlier by Maldacena and Nunez. The form (25) is clearly separable and is in fact of the form $`D_1`$ given in (22). Substituting the solution (22) in (9) yields (note that $`y_1<y<y_2`$):
$`ds^2`$ $`=`$ $`(y_1y_2y^2)^{1/3}(4({\displaystyle \frac{y}{y_1+y_22y}})^{1/3}ds_{AdS_5}^2+{\displaystyle \frac{y^{4/3}(y_1+y_22y)^{2/3}}{y_1y_2y^2}}d\mathrm{\Omega }_2^2+`$ (26)
$`+`$ $`4({\displaystyle \frac{y}{y_1+y_22y}})^{1/3}{\displaystyle \frac{(yy_1)(yy_2)}{y_1y_2y^2}}(d\chi {\displaystyle \frac{dx_2}{x_2}})^2+`$
$`+`$ $`({\displaystyle \frac{y_1+y_22y}{y}})^{2/3}({\displaystyle \frac{dy^2}{(yy_1)(yy_2)}}{\displaystyle \frac{dx_1^2+dx_2^2}{x_2^2}})).`$
This metric is singular if the roots $`y_1,y_2`$ are arbitrary. However if one sets $`y_1=s,y_2=s`$ (for some real $`s`$), then the metric is regular (once $`\chi `$ is identified with period $`2\pi `$):
$`ds^2`$ $`=`$ $`2^{2/3}(s^2+y^2)^{1/3}(2ds_{AdS_5}^2+{\displaystyle \frac{y^2}{s^2+y^2}}d\mathrm{\Omega }_2^2`$ (27)
$`+`$ $`2{\displaystyle \frac{s^2y^2}{s^2+y^2}}(d\chi {\displaystyle \frac{dx_1}{x_2}})^2+{\displaystyle \frac{dy^2}{s^2y^2}}+{\displaystyle \frac{dx_1^2+dx_2^2}{x_2^2}}).`$
By rescaling $`y`$ one can reduce the $`s`$-dependence to an overall factor. Substituting $`y=s\mathrm{cos}\theta `$ one finds the metric
$`ds^2`$ $`=`$ $`(2s)^{2/3}\mathrm{\Delta }^{1/3}(2ds_{AdS_5}^2+\mathrm{\Delta }^1\mathrm{cos}^2\theta d\mathrm{\Omega }_2^2+`$ (28)
$`+`$ $`2\mathrm{\Delta }^1\mathrm{sin}^2\theta (d\chi {\displaystyle \frac{dx_1}{x_2}})^2+d\theta ^2+{\displaystyle \frac{dx_1^2+dx_2^2}{x_2^2}}),`$
where $`\mathrm{\Delta }=1+\mathrm{cos}^2\theta `$, which is up to an overall factor the same as the metric appearing in . The metric (28) has the form of an $`AdS_5`$ fibration, and so is expected to be the supergravity dual of a superconformal field theory in four dimensions defined on the boundary of $`AdS_5`$. It may be interesting to study the singular deformations of (28) described by (27).
Another previously known case is a singular solution first discussed in as the description of the near-horizon limit of a system of intersecting M5-branes. In the present context it arises by using solution $`D_0`$, eq. (22), in the metric (9). To obtain the correct signature one needs $`\gamma >0`$ and $`\gamma <y<0`$. The metric can be written as
$$ds^2=\gamma ^{1/3}y^{1/3}\left(4ds_{AdS_5}^2\frac{y}{\gamma }d\mathrm{\Omega }_2^2+\frac{y+\gamma }{\gamma }d\chi ^2\frac{dy^2}{y(y+\gamma )}\frac{dx_1^2+dx_2^2}{y}\right).$$
(29)
The change of variable $`y=\gamma \mathrm{sin}^2\alpha `$ results in
$$ds^2=4\gamma ^{2/3}\mathrm{sin}^{2/3}\alpha \left(ds_{AdS_5}^2+\frac{1}{4}\mathrm{sin}^2\alpha d\mathrm{\Omega }_2^2+d\alpha ^2+\mathrm{cos}^2\alpha d\chi ^2+\frac{dx_1^2+dx_2^2}{\gamma \mathrm{sin}^2\alpha }\right),$$
(30)
which is of the form given in . The singularity at $`\alpha =0`$ was interpreted there as being due to M2-branes ending on the M5-branes.
A somewhat similar example arises from using solution $`D_0`$ (21) in the metric (13). To have the correct signature one needs $`\gamma <0`$ and $`y>|\gamma |`$. The metric reads
$`ds^2`$ $`=`$ $`{\displaystyle \frac{y^{4/3}}{\gamma ^{2/3}}}(ds_{AdS_2}^2+4|\gamma |d\mathrm{\Omega }_5^2+{\displaystyle \frac{|\gamma |}{y^2}}(dx_1^2+dx_2^2)+`$ (31)
$`+`$ $`4{\displaystyle \frac{y|\gamma |}{y}}d\chi ^2+{\displaystyle \frac{|\gamma |}{y^2(y|\gamma |)}}dy^2).`$
Setting $`\gamma =1`$ and $`y=1/\mathrm{sin}^2\alpha `$ leads to
$$ds^2=\mathrm{sin}^{8/3}\alpha \left(ds_{AdS_2}^2+d\alpha ^2+\mathrm{cos}^2\alpha d\chi ^2+\mathrm{sin}^2\alpha d\mathrm{\Omega }_5^2+\mathrm{sin}^4\alpha (dx_1^2+dx_2^2)\right),$$
(32)
which is the form found in for the near-horizon geometry of a system of semilocalised intersecting M2-branes. As this is an example of an $`AdS_2`$ fibration it can be expected to be dual to a superconformal quantum mechanics.
## 5 Conclusions
Although very simple, separable solutions of the Toda equation lead to a large family of half-BPS solutions of M-theory which include at least some physically interesting cases. All but one of the solutions arising this way are singular. It would be interesting to establish whether these geometries have interpretations in terms of branes. It could also be of interest to explore these solutions more closely, in particular, to understand the origin and interpretation of their singularities. Another natural question is whether there is a geometric interpretation of parameters appearing in the quadratic form $`Q`$.
It is perhaps disappointing that this class of supergravity solutions does not include any regular cases beyond the well-known Maldacena-Nunez solution. To find new regular cases one has to understand more general (in particular non-separable) solutions of the Toda equation. One example of a non-separable solution is in fact determined by the functions $`D_3`$ in eq. (24) if one allows the parameter $`u`$ appearing there to depend holomorphically on $`\xi x_1+ix_2`$. Such solutions were introduced and studied by Calderbank and Tod. While non-separable, these solutions also do not lead to regular metrics. It appears that for this purpose the construction described by Ward (which is discussed by LLM), may be more promising.
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# Glueball spectrum based on a rigorous three-dimensional relativistic equation for two-gluon bound states II: calculation of the glueball spectrum
## I Introduction
As mentioned in the preceding paper, searching for glueballs, nowadays, is an challenging task in particle physics. Since there are numerous technical difficulties of giving unambiguous identifications of the glueballs in experiment \[1-15\], it is expected that the existence of glueballs and their properties could be precisely predicted by theoretical investigations so as to guide the experimental searches. Various methods were proposed in the past to serve such investigations \[16-31\]. However, the predictions given by different methods are not consistent sometimes and even contradictory with each other \[32-33\]. Of these methods, the lattice simulation \[26-31\] is considered to be faithful. Nevertheless, even for this method, there still are controversies on the results given by different calculations. Apart from the lattice computation, the nonrelativistic potential model \[16-18\], the relativistic Dirac equation and the Bethe-Salpeter (B-S) equation \[23-25\] have recently been applied to evaluate the glueball spectrum. In these methods, the interaction between gluons is constructed by two parts: the short-range part which is described mainly by the one-gluon exchange interaction and the long-range part which is represented by a phenomenological confining potential. In the nonrelativistic potential model, the short-range interaction is simulated by a potential which is derived in the approximation of order $`v^2/c^2`$ where $`v`$ is the gluon velocity and $`c`$ the light velocity with the assumption that the gluons in a glueball move not too fast. In a recent work by using this model , with the choice of the lightest glueball masses given in the lattice computation as input, the authors obtained a series of two-gluon glueball states with masses below $`3GeV`$. Except for some glueball masses which are in pretty good agreement with the lattice predictions, the other calculated masses are apparently different from the lattice results. As was emphasized in Ref., to gain a physical solution to the lightest scalar glueball, it is necessary to additionally introduce a phenomenological smearing function to replace the $`\delta `$-function in the attractive contact terms of the potential. Otherwise, the Hamiltonian would be unbounded from below. This probably is an unnatural feature of the nonrelativistic potential model. In Ref., the calculation of the glueball spectrum was performed by employing the relativistic Dirac equation and showed only three theoretical results for the lowest glueball states $`0^{++},2^{++}`$ and $`3^{++}`$ some of which are not in so good agreement with those given by lattice investigations. In the calculation, a Fermi-Breit potential (the t-channel one-gluon exchange potential) was inserted into the Dirac equation with the assumption that the nature and the force between two gluons are the same as between two quarks. It seems that this assumption ignores the difference between the potential for quarks and the one for gluons. In addition, it would be mentioned that the Fermi-Breit potential is derived in the nonrelativistic approximation of order $`v^2/c^2`$. Therefore, the calculation is not fully relativistic. The relativistic calculation of the glueball spectrum was carried out in the framework of B-S equation . Owing to the difficulty of solving a relativistic equation, only a few states were predicted in these calculations. It is noted that in all the previous applications of the B-S equation, the four-dimensional equation was recast in a three-dimensional form in the instantaneous approximation in which the retardation effect is completely neglected. Another point we would like to note here is that in the aforementioned works, the gluons are all viewed as massive. Each of such gluons in general has three degrees of freedom of polarization. Correspondingly, a gluon field should includes three independent spatial components: two transverse fields and one longitudinal field in the three-dimensional space. In this sense, we can say that the Coulomb gauge as taken in Ref. is inappropriate for the massive gluons because in this gauge the longitudinal mode of the field is completely eliminated. Similarly, in Ref., only the transverse gluons are taken into account even though the temporal gauge adopted in the work allows existence of the longitudinal gluons.
In this paper, we intend to investigate the glueball spectrum based on the three-dimensional relativistic equation for two-gluon bound states which was derived in our former paper in the angular momentum representation. This equation is actually a coupled set of equations satisfied by the four B-S amplitudes for a glueball state: one is related to the positive energy states of two gluons, the other three are related to the two gluon states in which there is at least one gluon in the negative energy state. In the next section, we will derive from this coupled equations an equivalent equation obeyed by only the B-S amplitude of the glueball state for which the two gluons are in the positive energy states and give the effective interaction Hamiltonian in the equation a complete form. Since we are unable to compute all the terms in the Hamiltonian at present, we are limited ourself to work in a semi-phenomenological model by which the interaction Hamiltonian in the equation is given by the one-gluon exchange kernel plus the phenomenological linear confining kernel as was usually done in the previous literature \[16-18,23-25\]. The new aspects of this paper which distinguish from the previous works are: (1) The calculation is fully relativistic and hence includes the contribution arising from all the relativistic effects to the glueball masses; (2) The retardation effect of the one-gluon exchange interaction is completely taken into account; (3) Apart from the transverse modes of the gluon fields, the contribution from the longitudinal mode of the field to the glueball spectrum is appropriately considered; (4) The renormalization effect is considered by the effective QCD coupling constant which was derived in the one-loop approximation and in a mass-dependent subtraction in our previous work . This coupling constant is not only suitable in the high energy domain, but also in the low energy regime; (5) We work in the angular momentum representation. In the this representation, the glueball states are easily constructed. In particular, with completing the radial integrals containing three and four spherical Bessel functions, the gluon vertices are given explicit and analytical expressions which greatly facilitate the numerical task of solving the equation. The theoretical results obtained in this calculation are in quite good agreement with those given in the lattice study \[29-31\], In addition, some new predictions are presented.
The remainder of this paper is organized as follows. In Section II, we will derive the three-dimensional equation satisfied by the gluon positive energy state B-S amplitude from the coupled equations derived in the preceding paper. In Section III, the interaction Hamiltonian obtained in the tree diagram approximation will be discussed and its explicit expression will be given. Section IV serves to derive the expression of the linear-wise potential which is used to simulate the gluon confinement and incorporated in the glueball equation for numerical calculations. In the last section, the calculated results will be presented and some discussions will be made. In Appendices A and B, the analytical expressions of three-line and four-line gluon vertices are derived respectively.
## II The three-dimensional equation for the gluon positive energy state amplitude
The three-dimensional equations derived in the preceding paper (which will be called paper I later on) are shown in the following
$$(E_n\omega _\alpha \omega _\beta )\chi _{\alpha \beta }(n)=\underset{\rho \sigma }{}K(\alpha \beta ;\rho \sigma ;E_n)\chi _{\rho \sigma }(n),$$
(1)
where $`E`$ is the total energy of a glueball state, $`\omega _\alpha `$ and $`\omega _\beta `$ represent the energies of free gluons 1 and 2 respectively, $`\chi _{\alpha \beta }(n)`$ stands for the B-S amplitude describing the glueball state which is defined by
$$\chi _{\alpha \beta }(n)=0^+|𝐚_\alpha 𝐚_\beta |n,$$
(2)
and $`K(\alpha \beta ;\rho \sigma ;E)`$ designates the interaction kernel whose closed expression was derived in paper I. In the matrix notation, it is of the form
$`K`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{3}{}}}\mathrm{\Lambda }_j^{(i)}`$ (3)
$`=`$ $`\{{\displaystyle \underset{i=1}{\overset{3}{}}}g_iS_i{\displaystyle \underset{i,j=1}{\overset{3}{}}}g_iG_{ij}g_j+{\displaystyle \underset{i,j=1}{\overset{3}{}}}g_iG_iG^1G_jg_j\}S^1,`$ ()
where the matrices in the above expression were clearly defined in paper I. Noticing the definitions of $`\omega _\alpha `$ and $`𝐚_\alpha `$ (see Eqs.(4.9) and (5.16) in paper I),
$$\omega _\alpha =\{\begin{array}{c}\omega (k)\text{ if }\xi _\alpha =1\\ \omega (k)\text{ if }\xi _\alpha =1\end{array}$$
(4)
and
$$𝐚_\alpha (k)=\{\begin{array}{c}𝐚_\alpha (k)\text{ if }\xi _\alpha =1\\ \text{ }𝐚_\alpha ^+(k)\text{ if }\xi _\alpha =1\text{ }\end{array},$$
(5)
where the subscript $`\alpha `$ on the right hand side (RHS) of Eq.(2.5) is defined without including $`\xi _\alpha `$and hence $`𝐚_\alpha (k)`$ and $`𝐚_\alpha ^+(k)`$ represent the annihilation and creation operators respectively, the equation in Eq.(2.1) may be separately written as
| $`[E_n\omega (k_1)\omega (k_2)]\chi _{\alpha ^+\beta ^+}(n)=\underset{\rho \sigma }{}K(\alpha ^+\beta ^+;\rho \sigma ;E_n)\chi _{\rho \sigma }(n),`$ |
| --- |
| $`[E_n\omega (k_1)+\omega (k_2)]\chi _{\alpha ^+\beta ^{}}(n)=\underset{\rho \sigma }{}K(\alpha ^+\beta ^{};\rho \sigma ;E_n)\chi _{\rho \sigma }(n),`$ |
| $`[E_n+\omega (k_1)\omega (k_2)]\chi _{\alpha ^{}\beta ^+}(n)=\underset{\rho \sigma }{}K(\alpha ^{}\beta ^+;\rho \sigma ;E_n)\chi _{\rho \sigma }(n),`$ |
| $`[E_n+\omega (k_1)+\omega (k_2)]\chi _{\alpha ^{}\beta ^{}}(n)=\underset{\rho \sigma }{}K(\alpha ^{}\beta ^{};\rho \sigma ;E_n)\chi _{\rho \sigma }(n),`$ |
(6)
where the superscripts $`\pm `$ in $`\alpha ^\pm `$ and $`\beta ^\pm `$ denote $`\xi _\alpha ,\xi _\beta =\pm 1,`$the indices $`\rho `$ and $`\sigma `$ still include the indices $`\xi _\rho `$ and $`\xi _\sigma `$ respectively and
| $`\chi _{\alpha ^+\beta ^+}(n)=0^+|𝐚_\alpha 𝐚_\beta |n,\chi _{\alpha ^+\beta ^{}}(n)=0^+|𝐚_\alpha 𝐚_\beta ^+|n,`$ |
| --- |
| $`\chi _{\alpha ^{}\beta ^+}(n)=0^+|𝐚_\alpha ^+𝐚_\beta |n,\chi _{\alpha ^{}\beta ^{}}(n)=0^+|𝐚_\alpha ^+𝐚_\beta ^+|n.`$ |
(7)
Following the procedure described in Ref. for fermion systems, the coupled equations in Eq.(2.6) can be reduced to an equivalent equation satisfied by the B-S amplitude $`\chi _{\alpha ^+\beta ^+}(n)`$ for the glueball state in which each of gluons is in its positive energy state. For later convenience of derivation, we define
$$\mathrm{\Delta }_{ab}(E)=Ea\omega (k_1)b\omega (k_2),$$
(8)
where the subscript $`n`$ in $`E_n`$ has been suppressed, $`a,b=\pm 1`$,
| $`\varphi _{++}(\alpha \beta ;E)=\chi _{\alpha ^+\beta ^+}(n),\varphi _+(\alpha \beta ;E)=\chi _{\alpha ^+\beta ^{}}(n),`$ |
| --- |
| $`\varphi _+(\alpha \beta ;E)=\chi _{\alpha ^{}\beta ^+}(n),\varphi _{}(\alpha \beta ;E)=\chi _{\alpha ^{}\beta ^{}}(n)`$ |
(9)
and
$$K_{abcd}(\alpha \beta ;\rho \sigma ;E)=K(\alpha ^a\beta ^b;\rho ^c\sigma ^d;E),$$
(10)
in which the $`\alpha ,\beta ,\rho ,\sigma `$ are defined without including the index $`\xi `$. With the definitions given in Eqs.(2.8)-(2.10), the equations in Eq.(2.6) can compactly be written as
$$\mathrm{\Delta }_{ab}(E)\varphi _{ab}(\alpha \beta ;E)=\underset{cd}{}\underset{\rho \sigma }{}K_{abcd}(\alpha \beta ;\rho \sigma ;E)\varphi _{cd}(\rho \sigma ;E),$$
(11)
where $`a,b,c,d=\pm 1.`$ In the product space of momentum $`k_{1,}`$ $`k_2`$ and angular momentum marked by $`\alpha `$, the above equations may be written in the matrix form
$$\mathrm{\Delta }_{++}(E)\varphi _{++}(E)=K_{++++}(E)\varphi _{++}(E)+\underset{cd++}{}K_{++cd}(E)\varphi _{cd}(E),$$
(12)
$$\mathrm{\Delta }_{ab}(E)\varphi _{ab}(E)=K_{ab++}(E)\varphi _{++}(E)+\underset{cd++}{}K_{abcd}(E)\varphi _{cd}(E),$$
(13)
where $`ab++`$ and the terms related to $`\varphi _{++}(E)`$ have been separated out from the others. Furthermore, In the space spanned by $`\varphi _{ab}(E)`$ with $`ab++`$, we use the matrix representation defined as follows
| $`\psi (E)=\varphi _{++}(E),\varphi (E)=\{\varphi _{ab}(E)\},`$ |
| --- |
| $`\mathrm{\Delta }_+(E)=\mathrm{\Delta }_{++}(E),\mathrm{\Delta }(E)=\{\mathrm{\Delta }_{ab}(E)\},`$ |
| $`K_+(E)=K_{++++}(E),\overline{K}^t(E)=\{K_{++cd}(E)\},`$ |
| $`\overline{G}(E)=\{K_{ab++}(E)/\mathrm{\Delta }_{ab}(E)\},\},`$ |
| $`G(E)=\{K_{abcd}(E)/\mathrm{\Delta }_{ab}(E)\}.`$ |
(14)
According to these definitions, Eqs.(2.12) and (2.13) may be written in the full matrix form
$$\mathrm{\Delta }_+(E)\psi (E)=K_+(E)\psi (E)+\overline{K}^t(E)\varphi (E),$$
(15)
$$\varphi (E)=\overline{G}(E)\psi (E)+G(E)\varphi (E).$$
(16)
Solving the equation (2.16), we obtain
$$\varphi (E)=\frac{1}{1G(E)}\overline{G}(E)\psi (E).$$
(17)
Substituting the above expression into Eq.(2.15), we finally arrive at
$$\mathrm{\Delta }_+(E)\psi (E)=V(E)\psi (E),$$
(18)
where
$$V(E)=K_+(E)+\overline{K}^t(E)\frac{1}{1G(E)}\overline{G}(E),$$
(19)
which is identified itself with the interaction Hamiltonian. Noticing the definition
$$\frac{1}{1G(E)}=\underset{n=0}{}G^{(n)}(E).$$
(20)
Eq.(2.19) can be written as
$$V(E)=\underset{n=0}{}V^{(n)}(E),$$
(21)
where
| $`V^{(0)}(E)=K_+(E),`$ |
| --- |
| $`V^{(1)}(E)=\overline{K}^t(E)\overline{G}(E),`$ |
| $`V^{(2)}(E)=\overline{K}^t(E)G(E)\overline{G}(E)`$ |
| $`\mathrm{}\mathrm{}.`$ |
(22)
According to the definitions in Eq.(2.14), Eqs.(2.21) and (2.22) may be explicitly written as
$$V(\alpha \beta ;\gamma \delta ;E)=\underset{n=0}{}V^{(n)}(\alpha \beta ;\gamma \delta ;E),$$
(23)
where
$$V^{(0)}(\alpha \beta ;\gamma \delta ;E)=K(\alpha ^+\beta ^+;\gamma ^+\delta ^+;E),$$
(24)
$$V^{(1)}(\alpha \beta ;\gamma \delta ;E)=\underset{ab++}{}\underset{\rho \sigma }{}\frac{K(\alpha ^+\beta ^+;\rho ^a\sigma ^b;E)K(\rho ^a\sigma ^b;\gamma ^+\delta ^+E)}{Ea\omega (k_1)b\omega (k_2)},$$
(25)
| $`V^{(2)}(\alpha \beta ;\gamma \delta ;E)`$ |
| --- |
| $`=\underset{ab++}{}\underset{cd++}{}\underset{\rho \sigma }{}\underset{\mu \nu }{}\frac{K(\alpha ^+\beta ^+;\rho ^a\sigma ^b;E)K(\rho ^a\sigma ^b;\mu ^c\nu ^d;E)K(\mu ^c\nu ^d;\gamma ^+\delta ^+;E)}{(Ea\omega (k_1)b\omega (k_2))(Ec\omega (k_1)d\omega (k_2))}`$ |
| $`\mathrm{}\mathrm{},`$ |
(26)
where we have used the notation shown in Eqs.(2.8) and (2.10). From Eqs.(2.19)-(2.26), it is clear to see that the negative energy states of two gluons act as intermediate states to appear in the effective interaction Hamiltonian and give contributions to the higher order terms in the Hamiltonian. The equation (2.18) written in an explicit form is such that
$$[E\omega (k_1)\omega (k_2)]\psi (\alpha \beta ;E)=\underset{\gamma \delta }{}V(\alpha \beta ;\gamma \delta ;E)\psi (\gamma \delta ;E).$$
(27)
This is just the equation satisfied by the gluon positive energy state amplitudes for the glueball states in which
$$\psi (\alpha \beta ;E)=\chi _{\alpha ^+\beta ^+}(n).$$
(28)
## III The Hamiltonian given in the lowest order approximation
In this section, we plan to discuss the interaction Hamiltonian in the tree diagram approximation of the order of $`g^2`$ here $`g`$ is the QCD coupling constant. This Hamiltonian can only be given by the term shown in Eq.(2.24) because the other terms in the effective Hamiltonian give the contributions which are all higher than $`g^2`$. In general, the $`K(\alpha ^+\beta ^+;\gamma ^+\delta ^+;E)`$ should be calculated according to the expression denoted in Eq.(2.3) which includes three parts. The last part in Eq.(2.3) plays the role of cancelling the B-S reducible diagrams contained in the first two parts and gives no contribution of the order $`g^2`$. This is because (1) the coefficients $`g_1(\alpha \beta ;\rho \sigma \lambda )`$ and $`g_3(\alpha \beta ;\rho \sigma \lambda )`$ are proportional to $`g`$ and $`g_2(\alpha \beta ;\rho \sigma \tau \lambda )`$ is proportional to $`g^2`$; (2) the Green’s functions $`G_1(\rho \sigma \lambda ;\gamma \delta ;t_1t_2)`$ and $`G_3(\rho \sigma \lambda ;\gamma \delta ;t_1t_2)`$ vanish in the lowest order approximation. As for the first part in Eq.(2.3), in the lowest order approximation, it is easy to verify that
| $`S_1(\rho \sigma \lambda ,\gamma \delta )=0|[:𝐚_\rho 𝐚_\sigma :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0=0,`$ |
| --- |
| $`S_2(\rho \sigma \tau \lambda ,\gamma \delta )=0|[:𝐚_\rho 𝐚_\sigma 𝐚_\tau :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0=0,`$ |
| $`S_3(\rho \sigma \lambda ,\gamma \delta )=0|[:𝐜_\rho ^+𝐜_\sigma :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0=0,`$ |
(29)
where $`|0`$ denotes the bare vacuum state. Therefore, we only need to consider the second part in Eq.(2.3). In this part, the term related to $`g_2(\alpha \beta ;\rho \sigma \tau \lambda )`$ gives the contribution of order $`g^4`$ in the lowest approximation and hence is beyond our consideration. For the terms associated with $`g_3(\alpha \beta ;\rho \sigma \lambda ),`$the relevant Green’s functions vanish in the lowest order approximation. Thus, we are only left with terms in the second part of Eq.(2.3) such that
$$K^0(\alpha ^+\beta ^+;\gamma ^+\delta ^+;E)=\underset{\rho \sigma }{}\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho \sigma ;E)S^1(\rho \sigma ;\gamma ^+\delta ^+),$$
(30)
where
| $`\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho \sigma ;E)`$ |
| --- |
| $`=\underset{\xi \eta \lambda }{}\underset{\mu \nu \tau }{}g_1(\alpha ^+\beta ^+;\xi \eta \lambda )G_{11}(\xi \eta \lambda ;\mu \nu \tau ;E)g_1(\mu \nu \tau ;\rho \sigma )`$ |
(31)
and the indices $`\rho ,\sigma ,\mathrm{}`$ should be understood as $`\rho ^\pm ,\sigma ^\pm ,\mathrm{}.`$
Let us first compute the inverse $`S^1(\rho \sigma ;\gamma ^+\delta ^+)`$. For this purpose, we operate on the both sides of Eq.(3.2) with $`S`$ from the right and get
$$\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho \sigma ;E)=\underset{\gamma \delta }{}K^0(\alpha ^+\beta ^+;\gamma ^+\delta ^+;E)S(\gamma ^+\delta ^+;\rho \sigma ).$$
(32)
It is easy to verify that except for $`S(\gamma ^+\delta ^+;\rho ^{}\sigma ^{}),`$ the $`S(\gamma ^+\delta ^+;\rho ^{}\sigma ^+),`$ $`S(\gamma ^+\delta ^+;\rho ^+\sigma ^{})`$ and $`S(\gamma ^+\delta ^+;\rho ^+\sigma ^+)`$ are all vanishing in the lowest order approximation. As for the $`S(\gamma ^+\delta ^+;\rho ^{}\sigma ^{})`$, we have
$$S(\gamma ^+\delta ^+;\rho ^{}\sigma ^{})=\delta _{\gamma \rho }\delta _{\delta \sigma }+\delta _{\gamma \sigma }\delta _{\delta \rho }.$$
(33)
Substituting Eq.(3.5) in Eq.(3.4), we find
$$\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^{}\sigma ^{};E)=K^0(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)+K^0(\alpha ^+\beta ^+;\sigma ^+\rho ^+;E).$$
(34)
Since we may interchange the indices $`\rho `$ and $`\sigma `$ in Eq.(2.1) or in Eq.(2.27), noticing $`\chi _{\rho \sigma }(n)=\chi _{\sigma \rho }(n)`$ or $`\psi (\rho \sigma ;E)=\psi (\sigma \rho ;E)`$, we can write from Eq.(3.6) the following relation
$$K^0(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)=\frac{1}{2}\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^{}\sigma ^{};E),$$
(35)
which means that we may set
$$S^1(\rho \sigma ;\gamma ^+\delta ^+)=S^1(\rho ^{}\sigma ^{};\gamma ^+\delta ^+)=\frac{1}{2}\delta _{\gamma \rho }\delta _{\delta \sigma }.$$
(36)
Combining Eqs.(3.3) and (3.7), we have
| $`K^0(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)`$ |
| --- |
| $`=\frac{1}{2}\underset{\gamma \delta \lambda }{}\underset{\mu \nu \tau }{}g_1(\alpha ^+\beta ^+;\gamma \delta \lambda )G_{11}(\gamma \delta \lambda ;\mu \nu \tau ;E)g_1(\mu \nu \tau ;\rho ^{}\sigma ^{}).`$ |
(37)
In accordance with the definition of $`g_1(\alpha \beta ;\gamma \delta \lambda )`$ and $`g_1(\mu \nu \tau ;\rho \sigma )`$ (see Eqs.(5.7), (5.9), (5.10) and (6.4) in paper I), we can write
| $`g_1(\alpha ^+\beta ^+;\gamma \delta \lambda )=\underset{\tau }{}f_1(\gamma \delta \tau )\mathrm{\Delta }_{\alpha ^+\beta ^+;\tau \lambda },`$ |
| --- |
| $`g_1(\mu \nu \tau ;\rho ^{}\sigma ^{})=\underset{\lambda }{}f_1(\mu \nu \lambda )\mathrm{\Delta }_{\rho ^{}\sigma ^{};\lambda \tau },`$ |
(38)
where
| $`\mathrm{\Delta }_{\alpha ^+\beta ^+;\tau \lambda }=\mathrm{\Delta }_{\alpha ^+\tau }\delta _{\beta ^+\lambda }+\mathrm{\Delta }_{\beta ^+\tau }\delta _{\alpha ^+\lambda },`$ |
| --- |
| $`\mathrm{\Delta }_{\rho ^{}\sigma ^{};\lambda \tau }=\mathrm{\Delta }_{\rho ^{}\lambda }\delta _{\sigma ^{}\tau }+\mathrm{\Delta }_{\sigma ^{}\lambda }\delta _{\rho ^{}\tau }`$ |
(39)
and
$$f_1(\alpha \beta \gamma )=A(\alpha \beta \gamma )+A(\alpha \gamma \beta )+A(\gamma \alpha \beta ),$$
(40)
here $`A(\alpha \beta \gamma )`$ is the three-line gluon vertex given in the angular momentum representation. Considering the expressions in Eq.(3.11) and the fact that only $`\mathrm{\Delta }_{\alpha ^+\beta ^{}}=\mathrm{\Delta }_{\alpha ^{}\beta ^+}=\delta _{\alpha \beta }`$ are nonvanishing for $`\mathrm{\Delta }_{\alpha \beta }`$, Eq.(3.10) can be represented as
$$\begin{array}{c}g_1(\alpha ^+\beta ^+;\gamma \delta \lambda )=f_1(\gamma \delta \alpha ^{})\delta _{\beta ^+\lambda }+f_1(\gamma \delta \beta ^{})\delta _{\alpha ^+\lambda },\hfill \\ g_1(\mu \nu \tau ;\rho ^{}\sigma ^{})=f_1(\mu \nu \rho ^+)\delta _{\sigma ^{}\tau }f_1(\mu \nu \sigma ^+)\delta _{\rho ^{}\tau }.\hfill \end{array}$$
(41)
On inserting Eq.(3.13) into Eq.(3.9), one gets
| $`K^0(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)=\frac{1}{2}[\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)+(\alpha \beta )`$ |
| --- |
| $`+(\rho \sigma )+(\alpha \beta ,\rho \sigma )],`$ |
(42)
where
$$\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)=\underset{\gamma \delta ,\mu \nu }{}f_1(\gamma \delta \alpha ^{})G_{11}(\gamma \delta \beta ^+;\mu \nu \sigma ^{};E)f_1(\mu \nu \rho ^+)$$
(43)
and the other terms in Eq.(3.14) can be obtained from the first term by exchanging the indices as shown in Eq.(3.14). Now, let us calculate the Green’s functions $`G_{11}(\gamma \delta \beta ^+;\mu \nu \sigma ^{};E)`$ in the lowest order approximation which are the Fourier transform of the Green functions $`G_{11}(\gamma \delta \beta ^+;\mu \nu \sigma ^{};t_1t_2)`$. With the aid of Wick theorem, it can be found that only the following Green’s function is nonvanishing
| $`G_{11}(\gamma ^+\delta ^+\beta ^+;\mu ^{}\nu ^{}\sigma ^{};t_1t_2)`$ |
| --- |
| $`=0|T\{:𝐚_\gamma (t_1)𝐚_\delta (t_1):𝐚_\beta (t_1):𝐚_\mu ^+(t_2)𝐚_\nu ^+(t_2):𝐚_\sigma ^+(t_2)\}|0,`$ |
(44)
where $`𝐚_\gamma (t_1)`$ and $`𝐚_\sigma ^+(t_2)`$ are the annihilation and creation operators in the interaction picture. Noticing
$$𝐚_\gamma (t_1)=𝐚_\gamma e^{i\omega _\gamma t_1},𝐚_\sigma ^+(t_2)=𝐚_\sigma ^+e^{i\omega _\sigma t_2}$$
(45)
and applying the Wick theorem, we find
$$\begin{array}{c}G_{11}(\gamma ^+\delta ^+\beta ^+;\mu ^{}\nu ^{}\sigma ^{};t_1t_2)\hfill \\ =\theta (t_1t_2)e^{i(\omega _\gamma +\omega _\delta +\omega _\beta )(t_1t_2)}[\delta _{\gamma \mu }\delta _{\delta \nu }\delta _{\beta \sigma }+\delta _{\gamma \nu }\delta _{\delta \mu }\delta _{\beta \sigma }\hfill \\ +\delta _{\gamma \sigma }\delta _{\delta \mu }\delta _{\beta \nu }+\delta _{\gamma \sigma }\delta _{\delta \nu }\delta _{\beta \mu }+\delta _{\gamma \mu }\delta _{\delta \sigma }\delta _{\beta \nu }+\delta _{\gamma \nu }\delta _{\delta \sigma }\delta _{\beta \mu }].\hfill \end{array}$$
(46)
By the Fourier transformation and using the familiar integral representation of the step function, we obtain
$$\begin{array}{c}G_{11}(\gamma ^+\delta ^+\beta ^+;\mu ^{}\nu ^{}\sigma ^{};E)\hfill \\ =\frac{1}{i}_{\mathrm{}}^+\mathrm{}𝑑te^{iEt}G_{11}(\gamma ^+\delta ^+\beta ^+;\mu ^{}\nu ^{}\sigma ^{};t)\hfill \\ =\frac{1}{E\omega _\gamma \omega _\delta \omega _\beta +iϵ}[\delta _{\gamma \mu }\delta _{\delta \nu }\delta _{\beta \sigma }+\delta _{\gamma \nu }\delta _{\delta \mu }\delta _{\beta \sigma }+\delta _{\gamma \sigma }\delta _{\delta \mu }\delta _{\beta \nu }\hfill \\ +\delta _{\gamma \sigma }\delta _{\delta \nu }\delta _{\beta \mu }+\delta _{\gamma \mu }\delta _{\delta \sigma }\delta _{\beta \nu }+\delta _{\gamma \nu }\delta _{\delta \sigma }\delta _{\beta \mu }].\hfill \end{array}$$
(47)
Substituting the above expression in Eq.(3.15), we are led to
$$\begin{array}{c}\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)\hfill \\ =\underset{\gamma \delta }{}\frac{1}{E_n\omega _\gamma \omega _\delta \omega _\beta +iϵ}f_1(\gamma ^+\delta ^+\alpha ^{})\{[f_1(\gamma ^{}\delta ^{}\rho ^+)\hfill \\ +f_1(\delta ^{}\gamma ^{}\rho ^+)]\delta _{\beta \sigma }+[f_1(\delta ^{}\beta ^{}\rho ^+)+f_1(\beta ^{}\delta ^{}\rho ^+)]\delta _{\gamma \sigma }\hfill \\ +[f_1(\gamma ^{}\beta ^{}\rho ^+)+f_1(\beta ^{}\gamma ^{}\rho ^+)]\delta _{\delta \sigma }\}.\hfill \end{array}$$
(48)
Observing the above expression of $`\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)`$, we see, the first term containing $`\delta _{\beta \sigma }`$ is an unconnected term. It gives the one-loop correction to the gluon propagator whose effect will be included in the QCD effective coupling constant by the renormalization procedure. The second term proportional to $`\delta _{\gamma \sigma }`$ describes the one-gluon exchange interaction between two gluons. The third term is the exchanged term for the one-gluon exchange interaction. By dropping the first term, we have
| $`\mathrm{\Lambda }(\alpha ^+\beta ^+;\rho ^+\sigma ^+;E)`$ |
| --- |
| $`=\underset{\tau }{}\frac{[f_1(\tau ^+\sigma ^+\alpha ^{})+f_1(\sigma ^+\tau ^+\alpha ^{})][f_1(\tau ^{}\beta ^{}\rho ^+)+f_1(\beta ^{}\tau ^{}\rho ^+)]}{E\omega _\tau \omega _\sigma \omega _\beta +iϵ},`$ |
(49)
where the summation over $`\tau `$ is performed with respect to the gluon intermediate states and the function $`f_1`$ was represented in Eq.(3.12) in terms of the function $`A(\alpha \beta \gamma )`$whose explicit expression is derived in Appendix A and shown in the following.
$$\begin{array}{c}A(\alpha _1\alpha _2\alpha _3)=\frac{g}{2}(\frac{2}{\pi })^{\frac{3}{2}}f^{abc}k_3\underset{i=1}{\overset{3}{}}\frac{1}{\sqrt{2\omega (k_i)}}k_iB^{\xi _i}(l_i)_{\lambda _i\tau _i}T_{\lambda _3\lambda _j}\hfill \\ \times J_{l_1^{^{}}l_2^{^{}}l_3^{^{}}}(k_1,k_2,k_3)\mathrm{\Gamma }(l_i,l_i^{},m_i,\xi _i),\hfill \end{array}$$
(50)
where
$$\begin{array}{c}J_{l_1^{}l_2^{}l_3^{^{}}}(k_1,k_2,k_3)=\frac{\pi ^{\frac{5}{2}}}{4}(1)^{\frac{1}{2}(l_1^{}+l_2^{}+l_3^{^{}})}\hfill \\ \times \underset{\mu _1,\mu _2,\mu _3=0}{\overset{\mathrm{}}{}}\delta _{2(\mu _1+\mu _2+\mu _3),l_1^{}+l_2^{}+l_3^{}}\underset{i=1}{\overset{3}{}}\frac{k_i^{2\mu _il_i^{^{}}1}}{\mathrm{\Gamma }(\mu _i+1)\mathrm{\Gamma }(\mu _il_i^{}+\frac{1}{2})},\hfill \end{array}$$
(51)
| $`\mathrm{\Gamma }(l_i,l_i^{},m_i,\eta _i)`$ |
| --- |
| $`=\frac{i}{\sqrt{2\pi }}\underset{i=1}{\overset{3}{}}(1)^{(l_i+l_i^{}+m_i+1)\mathrm{sin}[\frac{(1\eta _i)\pi }{4}]}[(2l_i+1)(2l_i^{}+1)]^{\frac{1}{2}}`$ |
| $`\times \left(\begin{array}{ccc}l_1^{}& l_2^{}& l_3^{}\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)\left\{\begin{array}{ccc}1& 1& 1\\ l_1^{}& l_2^{}& l_3^{}\\ l_1& l_2& l_3\end{array}\right\},`$ |
(52)
$`T_{\lambda _3\lambda _j}`$ and $`B^{\xi _i}(l_i)_{\lambda _i\tau _i}`$ are defined respectively in (A.5) and (A.15). It is noted that for a given set of angular momenta, due to the restriction of $`\delta _{2(\mu _1+\mu _2+\mu _3),l_1^{}+l_2^{}+l_3^{}}`$, only a few terms in the series of Eq.(3.23) survive.
## IV The glueball equation with inclusion of the confinement
As shown in Sec.II, the interaction Hamiltonian in the exact relativistic equation for the glueball states is much complicated. In practical investigations, usually, one only considers the lowest order term in the Hamiltonian which was explicitly derived in the angular momentum representation in the preceding section. Obviously, in order to get reasonable theoretical results, it is necessary to introduce a certain confining potential to phenomenologically simulate all the higher order terms in the Hamiltonian \[16-18, 23-25\]. How to determine the form of confining interaction in the relativistic equation ? In this paper, as was similarly done for the quark-antiquark system , we introduce the confining Hamiltonian operator in such a manner
$`H_c`$ (53)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x_1d^3x_2f^{abe}f^{cde}\stackrel{}{A^a}(\stackrel{}{x}_1)\stackrel{}{A^b}(\stackrel{}{x}_1)V(|\stackrel{}{x}_1\stackrel{}{x}_2|)\stackrel{}{A^c}(\stackrel{}{x}_2)\stackrel{}{A^d}(\stackrel{}{x}_2)},`$ ()
where $`\stackrel{}{A^a}(\stackrel{}{x})`$ stands for the gluon field operator and $`V(|\stackrel{}{x}_1\stackrel{}{x}_2|)`$ the confining potential which will be taken to be the linear one
$$V(|\stackrel{}{x}_1\stackrel{}{x}_2|)=\gamma |\stackrel{}{x}_1\stackrel{}{x}_2|,$$
(54)
here the parameter $`\gamma `$ designates the strength of the confining potential. When the gluon field operators in Eq.(4.1) are replaced by their expansions in terms of the multipole fields (see the expansion given in Eq.(4.14) in paper I), Eq.(4.1) will be represented as
$$H_c=\underset{\alpha _1\alpha _2\alpha _3\alpha _4}{}V_c(\alpha _1\alpha _3;\alpha _2\alpha _4):𝐚_{\alpha _1}𝐚_{\alpha _2}𝐚_{\alpha _3}𝐚_{\alpha _4}:,$$
(55)
where
| $`V_c(\alpha _1\alpha _3;\alpha _2\alpha _4)`$ |
| --- |
| $`=\frac{1}{2}f^{abe}f^{cde}d^3x_1d^3x_2\stackrel{}{A}_{\beta _1}^{\lambda _1}(\stackrel{}{x}_1)\stackrel{}{A}_{\beta _2}^{\lambda _2}(\stackrel{}{x}_1)V(|\stackrel{}{x}_1\stackrel{}{x}_2|)\stackrel{}{A}_{\beta _3}^{\lambda _3}(\stackrel{}{x}_2)\stackrel{}{A}_{\beta _4}^{\lambda _4}(\stackrel{}{x}_2).`$ |
(56)
Here the symbols $`\lambda _i`$ and $`\beta _i`$ were defined in Appendix A. This is just the wanted confining potential written in the angular momentum representation which will be inserted into the relativistic equation. By using the Fourier transformation
$$|\stackrel{}{x}_1\stackrel{}{x}_2|=\frac{d^3q}{(2\pi )^3}\frac{8\pi }{\stackrel{}{q}^4}e^{i\stackrel{}{q}(\stackrel{}{x}_1\stackrel{}{x}_2)},$$
(57)
the expansion for the plane wave function
$$e^{\stackrel{}{p}\stackrel{}{x}}=4\pi \underset{lm}{}i^lj_l(pr)Y_{lm}^{}(\widehat{p})Y_{lm}(\widehat{x})$$
(58)
and the expression for the scalar multipole field
$$A_{lm}^S(k\stackrel{}{x})=\sqrt{\frac{2}{\pi }}kj_l(kr)Y_{lm}(\widehat{x}),$$
(59)
we can write
$$V(|\stackrel{}{x}_1\stackrel{}{x}_2|)=8\pi \gamma \underset{lm}{}_0^{\mathrm{}}𝑑q\frac{1}{q^4}A_{lm}^S(q\stackrel{}{x}_1)A_{lm}^S^{}(q\stackrel{}{x}_2).$$
(60)
Substitution of this expression into Eq.(4.4) leads to
$$V_c(\alpha _1\alpha _3;\alpha _2\alpha _4)=4\pi \gamma f^{abe}f^{cde}\underset{lm}{}_0^{\mathrm{}}𝑑q\frac{1}{q^4}D_{\alpha _1\alpha _2\tau }D_{\alpha _3\alpha _4\tau }^{},$$
(61)
where
$$\begin{array}{c}D_{\alpha _1\alpha _2\tau }=d^3x_1\stackrel{}{A}_{\beta _1}^{\lambda _1}(\stackrel{}{x}_1)\stackrel{}{A}_{\beta _2}^{\lambda _2}(\stackrel{}{x}_1)A_\tau ^S(\stackrel{}{x}_1),\\ D_{\alpha _3\alpha _4\tau }=d^3x_2\stackrel{}{A}_{\beta _3}^{\lambda _3}(\stackrel{}{x}_2)\stackrel{}{A}_{\beta _4}^{\lambda _4}(\stackrel{}{x}_2)A_\tau ^S^{}(\stackrel{}{x}_2),\end{array}$$
(62)
here $`\tau =(q,l,m)`$.
Completely analogous to the calculations described in Appendix A, it is easy to derive the following expression
$$\begin{array}{c}V_c(\alpha _1\alpha _3;\alpha _2\alpha _4)\hfill \\ =\frac{32}{\pi ^2}\lambda f^{abe}f^{cde}\underset{lm}{}_0^{\mathrm{}}𝑑q\frac{1}{q^2}\underset{i=1}{\overset{4}{}}k_iB^{\xi _i}(l_i)_{\lambda _i\tau _i}\hfill \\ \times J_{l_1^{}l_2^{}l}(k_1,k_2,q)J_{l_3^{}l_4^{}l}(k_3,k_4,q)\hfill \\ \times \widehat{\mathrm{\Gamma }}(l_i,l_i^{},m_i,\eta _i,l,m,1)\widehat{\mathrm{\Gamma }}(l_j,l_j^{},m_j,\eta _j,l,m,1),\hfill \end{array}$$
(63)
where $`i=1,2`$ , $`j=3,4`$ ,
$$\begin{array}{c}\widehat{\mathrm{\Gamma }}(l_i,l_i^{},m_i,\eta _i,l,m,\eta )\hfill \\ 𝑑\mathrm{\Omega }(\widehat{x})\stackrel{}{Y}_{l_1l_1^{}m_1}^{\eta _1}(\widehat{x})\stackrel{}{Y}_{l_2l_2^{}m_2}^{\eta _2}(\widehat{x})Y_{lm}^\eta (\widehat{x})\hfill \\ =\underset{i=1}{\overset{2}{}}(1)^{(l_i+l_i^{}+m_i+1)\mathrm{sin}[\frac{(1\eta _i)\pi }{4}]}(1)^{m\mathrm{sin}[\frac{(1\eta )\pi }{4}]}\hfill \\ \times \stackrel{~}{\mathrm{\Gamma }}(l_i,l_i^{},\eta _im_i,l,\eta m)\hfill \end{array}$$
(64)
with
$$\begin{array}{c}\stackrel{~}{\mathrm{\Gamma }}(l_i,l_i^{},m_i,l,m)𝑑\mathrm{\Omega }(\widehat{x})\stackrel{}{Y}_{l_1l_1^{}m_1}(\widehat{x})\stackrel{}{Y}_{l_2l_2^{}m_2}(\widehat{x})Y_{lm}(\widehat{x})\hfill \\ =(1)^{l_1+l_2^{}+l}[\frac{(2l_1+1)(2l_2+1)(2l_1^{}+1)(2l_2^{}+1)(2l+1)}{4\pi }]^{\frac{1}{2}}\hfill \\ \times \left(\begin{array}{ccc}l_1^{}& l_2^{}& l\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l\\ m_1& m_2& m\end{array}\right)\left\{\begin{array}{ccc}l_1& l_2& l\\ l_2^{}& l_1^{}& 1\end{array}\right\},\hfill \end{array}$$
(65)
and $`J_{l_il_j^{}l}(k_i,k_j,q)`$ is defined as the same as given in Eq.(3.23). With the introduction of the above confining potential, the total interaction Hamiltonian in Eq.(2.27) is now taken to be
$$V(\alpha \beta ;\gamma \delta )=V_g(\alpha \beta ;\gamma \delta )+V_c(\alpha \beta ;\gamma \delta ),$$
(66)
where
$$V_g(\alpha \beta ;\gamma \delta )=K^0(\alpha \beta ;\gamma \delta ),$$
(67)
which was formulated in Eqs.(3.9)-(3.24) and the $`V_c(\alpha \beta ;\gamma \delta )`$ was given in Eqs.(4.11)-(4.13).
Now we turn to discuss the wave function $`\psi (\alpha \beta )`$ in Eq.(2.27) which was defined in Eq.(2.28). In the lowest order approximation, the two-gluon bound states can be written in the form
$$|n=\underset{\alpha \beta }{}f_{\alpha \beta }^n𝐚_\alpha ^+𝐚_\beta ^+|0,$$
(68)
where
$$f_{\alpha \beta }^n=\frac{1}{\sqrt{8}}\delta _{c_1c_2}C_{l_1m_1l_2m_2}^{JM}f_{\lambda _1l_1,\lambda _2l_2}^{J\pi }(k_1,k_2),$$
(69)
in which $`\delta _{c_1c_2}`$ represents the color singlet, $`C_{l_1m_1l_2m_2}^{JM}`$ is the C-G coupling coefficient, $`\alpha =(c_1,\lambda _1,l_1,m_1,k_1,\xi _\alpha =+1),`$ $`\beta =(c_2,\lambda _2,l_2,m_2,k_2,\xi _\beta =+1)`$, $`J,M`$ are the total angular momentum and its third component of a glueball and $`\pi `$ denotes the spatial parity and charge conjugation parity. With the introduction of the cluster coordinates
$$\stackrel{}{K}=\stackrel{}{k}_1+\stackrel{}{k}_2,\stackrel{}{k}=\frac{1}{2}(\stackrel{}{k}_1\stackrel{}{k}_2),$$
(70)
where $`\stackrel{}{K}`$ and $`\stackrel{}{k}`$ are the total momentum and relative momentum respectively, we see, in the center of mass system ($`\stackrel{}{K}=0`$), Eq.(4.17) reads
$$f_{\alpha \beta }^n=\frac{1}{\sqrt{8}}\delta _{c_1c_2}C_{l_1m_1l_2m_2}^{JM}f_{\lambda _1l_1,\lambda _2l_2}^{J\pi }(k)\delta (k_1k_2),$$
(71)
where $`k=|\stackrel{}{k}|=k_1=k_2`$. Substituting Eq.(4.16) into Eq.(2.2), we find
$$\psi (\alpha \beta )=\chi _{\alpha ^+\beta ^+}(n)=\frac{1}{\sqrt{8}}\delta _{c_1c_2}C_{l_1m_1l_2m_2}^{JM}g_{\lambda _1l_1,\lambda _2l_2}^{J\pi h}(k)\delta (k_1k_2),$$
(72)
where
$$g_{\lambda _1l_1,\lambda _2l_2}^{J\pi h}(k)=f_{\lambda _1l_1,\lambda _2l_2}^{J\pi }+(1)^hf_{\lambda _2l_2,\lambda _1l_1}^{J\pi },$$
(73)
in which $`h=l_1+l_2J`$. Evidently, if $`\lambda _1=\lambda _2,l_1=l_2`$ and $`h=odd`$, we have $`g_{\lambda _1l_1,\lambda _2l_2}^{J\pi h}(k)=0`$. This gives a new selection rule for the glueball states. Upon substituting Eq.(4.20) into Eq.(2.27) and noticing
$$\begin{array}{c}\alpha =(c,\lambda _1,l_1,m_1,k,+1),\beta =(c,\lambda _2,l_2,m_2,k,+1),\hfill \\ \rho =(c^{},\lambda _3,l_3,m_3,q,+1),\sigma =(c^{},\lambda _4,l_4,m_4,q,+1)\hfill \end{array}$$
(74)
and
$$\underset{\alpha }{}\underset{c\lambda lm}{}_0^{\mathrm{}}𝑑k,$$
(75)
we finally arrive at
$$\begin{array}{c}(E2\omega )g_{\lambda _1l_1,\lambda _2l_2}^{J\pi h}(k)\hfill \\ =\frac{1}{8(2J+1)}\underset{\lambda _3l_3}{}_0^{\mathrm{}}𝑑qV(\lambda _1l_1k;\lambda _2l_2k;\lambda _3l_3q;\lambda _4l_4q;E)g_{\lambda _3l_3,\lambda _4l_4}^{J\pi h}(q),\hfill \end{array}$$
(76)
where $`\omega =\sqrt{k^2+\mu ^2}`$, $`E`$ is the mass of a glueball state given in the center of mass frame and
$`V(\lambda _1l_1k;\lambda _2l_2k;\lambda _3l_3q;\lambda _4l_4q;E)`$ (77)
$`=`$ $`{\displaystyle \underset{_{m_1m_3m_2m_4cc^{^{}}M}}{}}V(\alpha \beta ;\rho \sigma ;E)C_{l_1m_1l_2m_2}^{JM}C_{l_3m_3l_4m_4}^{JM},`$ ()
here $`V(\alpha \beta ;\rho \sigma ;E)`$ was given in Eq.(4.14). When the explicit expression of $`V(\alpha \beta ;\rho \sigma ;E)`$ is substituted in Eq.(4.25), one can see that the summation over $`m_{1,}m_{2,}m_3,m_4`$ and $`M`$ is easily carried out by utilizing the well-known formula for the angular momentum coupling and the summation over the color indices $`c`$ and $`c^{}`$ can be completed by noticing $`f^{abc}f^{abc}=24`$. We think, it is unnecessary to show here the result given by these summations.
The equation in Eq.(4.24) is the eigenvalue equation used to calculate the glueball spectrum. In the calculation, the QCD coupling constant $`g`$ contained in the part of Hamiltonian $`V_g(\alpha \beta ;\gamma \delta )`$ is replaced by the running one which was derived in Ref. recently in the one-loop approximation and in a mass-dependent momentum space subtraction. The coupling constant used in this calculation is of the form
$$\alpha _s(\lambda )=\frac{\alpha _s^0}{1+\frac{\alpha _s^0}{2\pi }G(\lambda )},$$
(78)
where $`\alpha _s(\lambda )=g^2(\lambda )/4\pi `$,
$$G(\lambda )=11\mathrm{ln}\lambda \frac{2}{3}N_f[2+\sqrt{3}\pi \frac{2}{\lambda ^2}+(1+\frac{2}{\lambda ^2})\frac{1}{\lambda }\eta (\lambda )],$$
(79)
in which $`N_f`$ is the number of quark flavors,
$$\eta (\lambda )=\sqrt{\lambda ^24}\mathrm{ln}\frac{\lambda +\sqrt{\lambda ^24}}{2}$$
(80)
and $`\lambda =\sqrt{\frac{p^2+\mu ^2}{\mathrm{\Lambda }_{QCD}^2}}`$ here $`p`$ is chosen to be the transfer momentum of the exchanged gluon which may simply be taken as $`p=kq`$ for simplicity and $`\mathrm{\Lambda }_{QCD}`$ is the QCD scale parameter. The running coupling constant shown above is applicable not only in the high energy domain, but also in the low energy regime. Particularly, in the large momentum limit, it immediately goes over to the result obtained previously in the minimal subtraction scheme.
## V Numerical results and discussions
In this section, we first show the theoretical glueball masses calculated from the equation given in Eq.(4.24) and then make some discussions. Our calculation is performed by using the standard program of Mathematica which allows us to compute the effective Hamiltonian in Eq.(4.24) analytically. In this paper, we confine ourself to investigate the low-lying glueball states including the $`0^{++},0^+,1^{++},1^+,2^{++}`$ and $`2^+`$ ground and lower excited states whose masses are less than $`4.0GeV`$. Some of these states have been investigated before in various models. We also examine the effects of the longitudinal mode of the multipole fields and the different sets of free parameters on the glueball masses. In our calculation, the theoretical parameters are adjusted so as to be able to compare our results to those presented recently by the lattice simulation . The parameters taken are: the gluon mass $`\mu =0.42GeV`$ which is comparable with $`\mu =(0.5\pm 0.2)GeV`$ taken previously in the nonperturbative continuum studies, the scale parameter $`\mathrm{\Lambda }_{QCD}=0.45GeV`$ and the strength of the confining potential $`\gamma =0.18GeV^2`$ , which satisfies the relation $`\mu \mathrm{\Lambda }_{QCD}`$ and $`\gamma \mathrm{\Lambda }_{QCD}^2`$ which make the parameters essentially depend on a single dimensional quantity. Moreover, the value of $`\gamma =0.18GeV^2`$ is consistent with that of the string tension in lattice simulations. The coupling constant $`\alpha _s^0=0.3`$ and quark flavor $`N_f=3`$. The calculated masses of glueball states are displayed in table I. In the table, the case I and the case II respectively denote the results obtained with and without considering the contribution arising from the longitudinal mode of the multipole fields which appears in the intermediate states of the matrix elements of the interaction Hamiltonian. In the last column of the table, we quote the results shown in Ref. which were calculated by the lattice simulation.
Table I. The mass spectrum of two-gluon glueballs.
| | Mass(GeV) | | Mass(MeV) |
| --- | --- | --- | --- |
| Glueball states ($`J^{PC}`$) | Case I | Case II | Lattice results |
| 0<sup>++</sup> | 1.73 | 2.18 | 1730(50)(80) |
| | 2.66 | 3.59 | 2670(180)(130) |
| | 3.59 | | |
| 0<sup>-+</sup> | 2.60 | 2.30 | 2590(40)(130) |
| | 3.65 | 3.78 | 3640(60)(180) |
| 1<sup>++</sup> | 2.73 | 2.42 | |
| | 3.45 | 3.51 | |
| 1<sup>-+</sup> | 2.67 | 2.59 | |
| | 2.87 | 3.01 | |
| 2<sup>++</sup> | 2.43 | 2.43 | 2400(25)(120) |
| 2<sup>-+</sup> | 3.32 | 2.26 | 3100(30)(150) |
As seen from Eq.(4.24), each glueball state is not only assigned by its spin $`J`$ and parity $`\pi ,`$ but also related to the mode marked by $`(\lambda _1,l_1)(\lambda _2,l_2)`$. In this paper, we take low-lying modes to perform the calculation. For the scalar glueballs of quantum numbers $`J^{PC}=0^{++}`$ and the tensor ones $`2^{++}`$, according to the angular momentum and the parities of the multipole fields, we take the mixture of the modes $`(TE1TE1)`$ and $`(TM1TM1)`$. For the glueball states $`0^+`$ and $`2^+`$, the modes are taken to be $`(TE1TM1)`$ for every glueball, as was similarly done in the investigation within the bag model . This means that these glueballs are mainly constructed by the gluons with transverse polarization. But, this does not imply no contribution of the longitudinally polarized gluons to these glueballs. The longitudinal gluons may, as virtual particles, appear in the intermediate states in the effective interaction Hamiltonian. It is emphasized here that for the transverse mode of gluons, as mentioned in Appendix A, the mode $`l_i=0`$ is not permitted. This mode can only exist for the longitudinal gluons. Different from the case of massless gluons, the longitudinal mode of massive gluons is possible to take part in formation of some glueballs. For example, the glueball states $`1^+`$ can be formed not only by a combination of modes $`(M1E1)`$ and $`(E1L0)`$ which gives the states with masses as listed in the table, the modes $`(M1E1)`$ and $`(L1L0)`$ can also form the glueball states with masses $`3.23GeV`$ and $`3.82GeV`$ respectively. For the states $`1^{++},`$we only take the mode $`(L0M1)`$ in our calculation because according to the B-S amplitude constructed in Sec.IV (see Eq.(4.21)), the modes ($`E1E1),(M1M1)`$ are forbidden.
In order to determine the parameter dependence and errors in our calculation, we take three sets of parameters in which we set $`\mu =\mathrm{\Lambda }_{QCD}`$ and $`\gamma `$ (in unit GeV<sup>2</sup>)$`=\mathrm{\Lambda }_{QCD}^2`$. The results of case I in table I with different parameters (the parameters $`\alpha _s^0`$ and $`N_f`$ remain unchanged) are presented in Table II. We find that the masses of glueballs increase gradually when the set of parameters increase. This indicates that our numerical calculation is stable and the results are reliable.
Table II. The mass spectrum with different parameters.
| | Mass(GeV) | | |
| --- | --- | --- | --- |
| Glueball states ($`J^{PC}`$) | $`\mu =0.35GeV`$ | $`\mu =0.45GeV`$ | $`\mu =0.55GeV`$ |
| 0<sup>++</sup> | 1.66 | 1.74 | 1.81 |
| | 2.59 | 2.68 | 2.75 |
| | 3.52 | 3.60 | 3.66 |
| 0<sup>-+</sup> | 2.51 | 2.60 | 2.68 |
| | 3.57 | 3.66 | 3.73 |
| 1<sup>++</sup> | 2.64 | 2.73 | 2.82 |
| | 3.37 | 3.46 | 3.54 |
| 1<sup>-+</sup> | 2.59 | 2.69 | 2.78 |
| | 2.78 | 2.88 | 2.96 |
| 2<sup>++</sup> | 2.35 | 2.45 | 2.53 |
| 2<sup>-+</sup> | 3.25 | 3.34 | 3.41 |
The calculated results show that the gluon mass give an appreciable effect on the glueball masses. In our calculation, the mass of gluon should be around 0.45 GeV. This fact indicates the reasonability of the QCD with massive gluons which is chosen to be the starting point in our calculation. For this kind of QCD, it is necessary to take the longitudinal mode of gluons into account. As shown in Table I, when the longitudinal mode is considered in the intermediate states, the theoretical masses for the states mentioned above would be greatly improved. Otherwise, there would occur a considerable discrepancy between the results given by this paper and the lattice calculation. In addition, as illustrated before, the longitudinal mode allows us to investigate more glueball states which possibly exist in the world.
In comparison with the previous theoretical glueball masses obtained from the B-S equation and the Dirac equation, our results give more support to the lattice predictions \[29-31\] which are believed to be more reliable because the lattice calculation is based on the QCD first principle and essentially nonperturbative. Within the statistical errors existing in the lattice calculations, our results shown in the first column of Table I can be considered to be consistent with lattice predictions for the low-lying glueballs with masses less than $`3.5GeV`$, especially, for the lowest scalar glueball state $`0^{++}`$ with mass about $`1.7GeV`$ and the tensor glueball state $`2^{++}`$ with mass about $`2.4GeV`$. The achievement of the better consistence is obviously attributed to the fact that our calculation is fully relativistic and is able to include the contributions arising from the retardation effect and longitudinal mode of the gluon field which could not be considered in the previous investigations \[23-25\]. Now let us analyze our results in some more detail. It is mentioned that the lowest scalar glueball state $`0^{++}`$ and the tensor one $`2^{++\text{ }}`$ have been investigated in many models and the theoretical masses are almost the same even though there are a little difference between different calculations \[29-31\]. For example, the mass of the lowest state $`0^{++}`$ was given by $`1754\pm 65\pm 86MeV`$ in a recent lattice calculation which is different from that given in Ref.. It is expected that these states may be identified with the pure glueball states and searched out first in future experiments. Aside from the two states mentioned above, the first radial excited state $`0^{++}`$ with mass $`2.66GeV`$ should correspond to the state $`2670(180)(130)MeV`$ given in the lattice simulation even though whether the latter is a pure glueball or not is still in question . The next radial excited state $`0^{++}`$ with mass $`3.59GeV`$ is a new prediction given in this paper which was not predicted in the lattice simulation and the other calculations. As for the pseudoscalar states $`0^+`$ and pseudotensor state $`2^+`$ are all comparable with the corresponding states presented in the lattice calculation. But, the mass of the state $`2^+`$ is little higher than the lattice one. In addition, we note that the states $`1^+`$ and $`1^{++}`$ were not predicted in the lattice simulation, but the states $`1^+`$ with masses $`2.67GeV`$ and the state $`1^{++\text{ }}`$with mass $`2.73GeV`$ are compatible with the recent calculation by the nonrelativistic potential model .
In conclusion, it is emphasized that different from some previous investigations, the calculation in this paper is based on the rigorous three-dimensional relativistic equation satisfied by the two-gluon glueball states which is derived from the QCD with massive gluons and represented in the angular momentum representation. Especially, the interaction Hamiltonian in the equation is given a complete expression which provides a firm basis for further study. In this paper, even though we work in the relativistic potential model with introducing phenomenologically a confining potential, the new consideration of the retardation effect and the longitudinal mode of the gluon fields allows us to get the improved theoretical results which are well consistent with the lattice predictions. The only uncertainty in our calculation arises from the introduction of confining potential. Certainly, if a sophisticated confining potential could be found from the exact interaction Hamiltonian derived in this and former papers, it would be anticipated that a relativistic calculation may give more accurate theoretical predictions.
## VI Acknowledgment
This work was supported in part by National Natural Science Foundation of China.
## VII Appendix A: The expression of gluon three-line vertex
This appendix is used to derive the explicit expression of the gluon three-line vertex in the angular momentum representation. As shown in Eqs.(4.7), (4.29) and (4.30) in paper I, the gluon three-line vertex in the interaction Hamiltonian is
$$H_g^3=\frac{g}{2}f^{abc}d^3x(\stackrel{}{A^a}\times \stackrel{}{A^b})(\times \stackrel{}{A^c}).$$
(81)
It is represented in the angular momentum representation as follows
$$H_g^3=\underset{\alpha _1\alpha _2\alpha _3}{}A(\alpha _1\alpha _2\alpha _3):𝐚_{\alpha _1}𝐚_{\alpha _2}𝐚_{\alpha _3}:,$$
(82)
where
$$A(\alpha _1\alpha _2\alpha _3)=\frac{g}{2}f^{abc}d^3x(\stackrel{}{A}_{\beta _1}^{\lambda _1}\times \stackrel{}{A}_{\beta _2}^{\lambda _2})(\times \stackrel{}{A}_{\beta _3}^{\lambda _3}).$$
(83)
Here we have set $`\alpha _i=(\lambda _i,\beta _i)`$ in which $`\lambda _i=TE,TM,L`$ mark the transverse electric, transverse magnetic and longitudinal modes of the multipole fields and $`\beta _i=(l_i,m_i,k_i,\xi _i).`$ From now on, we use the symbols $`l_im_i`$ to represent the total angular momentum. For later convenience, the relations between the multipole fields which were mentioned in paper I are represented in the matrix form
$$\times \left(\begin{array}{c}\stackrel{}{A}_{lm}^{TE}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^{TM}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^L(k\stackrel{}{x})\end{array}\right)=k\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)\left(\begin{array}{c}\stackrel{}{A}_{lm}^{TE}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^{TM}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^L(k\stackrel{}{x})\end{array}\right).$$
(84)
When we define
$$T=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right),$$
(85)
(A.4) can be rewritten as
$$\times \stackrel{}{A}_{lm}^{\lambda _i}(k\stackrel{}{x})=kT_{\lambda _i\lambda _j}\stackrel{}{A}_{lm}^{\lambda j}(k\stackrel{}{x}).$$
(86)
Thus, (A.3) reads
$$A(\alpha _1\alpha _2\alpha _3)=\frac{g}{2}f^{abc}k_3T_{\lambda _3\lambda _j}d^3x(\stackrel{}{A}_{\beta _1}^{\lambda _1}\times \stackrel{}{A}_{\beta _2}^{\lambda _2})\stackrel{}{A}_{\beta _3}^{\lambda _j}.$$
(87)
With the definition
$$\stackrel{}{y}_{lm}^{(\tau )}(k\stackrel{}{x})=\sqrt{\frac{2}{\pi }}kj_{l+\tau }(kr)\stackrel{}{Y}_{l,l+\tau ,m}(\widehat{x})\stackrel{}{y}_{ll^{}m}(k\stackrel{}{x}),$$
(88)
where $`l^{}=l+\tau `$ and $`\tau =0,\pm 1`$, the relations shown in Eqs.(3.2)-(3.4) in paper I for the multipole fields can also be written in the matrix form
$$\left(\begin{array}{c}\stackrel{}{A}_{lm}^{TE}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^{TM}(k\stackrel{}{x})\\ \stackrel{}{A}_{lm}^L(k\stackrel{}{x})\end{array}\right)=\left(\begin{array}{ccc}i\sqrt{\frac{l+1}{2l+1}}& 0& i\sqrt{\frac{l}{2l+1}}\\ 0& 1& 0\\ i\sqrt{\frac{l}{2l+1}}& 0& i\sqrt{\frac{l+1}{2l+1}}\end{array}\right)\left(\begin{array}{c}\stackrel{}{y}_{lm}^{(1)}(k\stackrel{}{x})\\ \stackrel{}{y}_{lm}^{(0)}(k\stackrel{}{x})\\ \stackrel{}{y}_{lm}^{(1)}(k\stackrel{}{x})\end{array}\right).$$
(89)
If we define $`(TE,TM,L)=(1,0,1)`$ and
$$B(l)=\left(\begin{array}{ccc}i\sqrt{\frac{l+1}{2l+1}}& 0& i\sqrt{\frac{l}{2l+1}}\\ 0& 1& 0\\ i\sqrt{\frac{l}{2l+1}}& 0& i\sqrt{\frac{l+1}{2l+1}}\end{array}\right),$$
(90)
where $`l0`$ and
$$B(0)=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& i\end{array}\right)$$
(91)
which means only the longitudinal mode survives when $`l=0`$, then (A.9) can concisely be written as
$$\stackrel{}{A}_{lm}^\lambda (k\stackrel{}{x})=B(l)_{\lambda \tau }\stackrel{}{y}_{lm}^{(\tau )}(k\stackrel{}{x}).$$
(92)
After inserting (A.12) into (A.7) and noticing the definition given in Eq.(4.16) in paper I, we have
$$\begin{array}{c}A(\alpha _1\alpha _2\alpha _3)\hfill \\ =\frac{g}{2}f^{abc}k_3\frac{1}{\underset{i=1}{\overset{3}{}}\sqrt{2\omega (k_i)}}T_{\lambda _3\lambda _j}B^{\xi _1}(l_1)_{\lambda _1\tau _1}B^{\xi _2}(l_2)_{\lambda _2\tau _2}B^{\xi _3}(l_3)_{\lambda _j\tau _3}\hfill \\ \times d^3x[\stackrel{}{y}_{l_1l_1^{}m_1}^{\xi _1}(k_1\stackrel{}{x})\times \stackrel{}{y}_{l_2l_2^{}m_2}^{\xi _2}(k_2\stackrel{}{x})]\stackrel{}{y}_{l_3l_3^{}m_3}^{\xi _3}(k_3\stackrel{}{x}),\hfill \end{array}$$
(93)
where we have defined
$$\stackrel{}{y}_{ll^{}m}^\xi (k\stackrel{}{x})=\{\begin{array}{c}\stackrel{}{y}_{ll^{}m}(k\stackrel{}{x})\text{ if }\xi =1\\ \text{ }\stackrel{}{y}_{ll^{}m}^{}(k\stackrel{}{x})\text{ if }\xi =1\end{array}$$
(94)
and
$$B^\xi (l_i)_{\lambda \tau }=\{\begin{array}{c}B(l_i)_{\lambda \tau }\text{ if }\xi =1\\ \text{ }B^{}(l_i)_{\lambda \tau }\text{ if }\xi =1\end{array}.$$
(95)
In light of the the expression in (A.8), the integral over $`\stackrel{}{x}`$ in (A.13) can be represented as
$$\begin{array}{c}d^3x[\stackrel{}{y}_{l_1l_1^{}m_1}^{\xi _1}(k_1\stackrel{}{x})\times \stackrel{}{y}_{l_2l_2^{}m_2}^{\xi _2}(k_2\stackrel{}{x})]\stackrel{}{y}_{l_3l_3^{}m_3}^{\xi _3}(k_3\stackrel{}{x})\hfill \\ =(\frac{2}{\pi })^{\frac{3}{2}}k_1k_2k_3J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)\mathrm{\Gamma }(l_i,l_i^{},m_i,\xi _i),\hfill \end{array}$$
(96)
where
$$\mathrm{\Gamma }(l_i,l_i^{},m_i,\xi _i)=𝑑\mathrm{\Omega }(\widehat{x})[\stackrel{}{Y}_{l_1l_1^{}m_1}^{\xi _1}(\widehat{x})\times \stackrel{}{Y}_{l_2l_2^{}m_2}^{\xi _2}(\widehat{x})]\stackrel{}{Y}_{l_3l_3^{}m_3}^{\xi _3}(\widehat{x})$$
(97)
and
$$J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)=𝑑rr^2j_{l_1^{}}(k_1r)j_{l_2^{}}(k_2r)j_{l_3^{}}(k_3r).$$
(98)
On substituting (A.16) in (A.13), we just give the formula denoted in Eq.(3.22).
First, let us calculate the $`\mathrm{\Gamma }(l_i,l_i^{},m_i,\xi _i)`$ in the case of $`\xi _i=1`$. In this case, we set
$$\mathrm{\Gamma }(l_i,l_i^{},m_i)=𝑑\mathrm{\Omega }(\widehat{x})[\stackrel{}{Y}_{l_1l_1^{}m_1}(\widehat{x})\times \stackrel{}{Y}_{l_2l_2^{}m_2}(\widehat{x})]\stackrel{}{Y}_{l_3l_3^{}m_3}(\widehat{x}).$$
(99)
By using the following formulas :
$$\begin{array}{c}\stackrel{}{Y}_{ll^{}m}(\widehat{x})=\underset{m^{}q}{}C_{l^{}m^{}1q}^{lm}Y_{l^{}m^{}}(\widehat{x})\stackrel{}{e}_q,\hfill \\ \stackrel{}{e}_{q_1}\times \stackrel{}{e}_{q_2}=i\sqrt{2}\underset{s}{}C_{1q_11q_2}^{1s}\stackrel{}{e}_s,\hfill \\ \stackrel{}{e}_q=(1)^q\stackrel{}{e}_q^{},\text{ }\stackrel{}{e}_q^{}\stackrel{}{e}_q^{}=\delta _{qq^{}},\hfill \end{array}$$
(100)
we find
$$\begin{array}{c}\mathrm{\Gamma }(l_i,l_i^{},m_i)=i\sqrt{2}\underset{n_1n_2n_3q_1q_2q_3}{}(1)^{q_3}C_{l_1^{}n_11q_1}^{l_1m_1}C_{l_2^{}n_21q_2}^{l_2m_2}C_{l_3^{}n_31q_3}^{l_3m_3}C_{1q_11q_2}^{1,q_3}\hfill \\ \times d\mathrm{\Omega }(\widehat{x})Y_{l_1^{}n_1}(\widehat{x})Y_{l_2^{}n_2}(\widehat{x})Y_{l_3^{}n_3}(\widehat{x}).\hfill \end{array}$$
(101)
Employing the familiar formula for the above integral and the definition and property of 3-j and 9-j symbols for the angular momentum couplings , one can get
$$\begin{array}{c}\mathrm{\Gamma }(l_i,l_i^{},m_i)=i[\frac{(2l_1+1)(2l_2+1)(2l_3+1)(2l_1^{}+1)(2l_2^{}+1)(2l_3^{}+1)}{2\pi }]^{\frac{1}{2}}\\ \times \left(\begin{array}{ccc}l_1^{}& l_2^{}& l_3^{}\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)\left\{\begin{array}{ccc}1& 1& 1\\ l_1^{}& l_2^{}& l_3^{}\\ l_1& l_2& l_3\end{array}\right\}.\end{array}$$
(102)
In the case of $`\xi =1`$, it is easy to find
$$\stackrel{}{Y}_{ll^{}m}^{}(\widehat{x})=(1)^{l+l^{}+m+1}\stackrel{}{Y}_{ll^{}m}(\widehat{x}).$$
(103)
If we define
$$\stackrel{}{Y}_{ll^{}m}^\eta (\widehat{x})=\{\begin{array}{c}\stackrel{}{Y}_{ll^{}m}(\widehat{x})\text{ if }\xi =1\\ \text{ }\stackrel{}{Y}_{ll^{}m}^{}(\widehat{x}))\text{ if }\xi =1\end{array},$$
(104)
the vector spherical harmonics $`\stackrel{}{Y}_{ll^{}m}(\widehat{x})`$ and $`\stackrel{}{Y}_{ll^{}m}^{}(\widehat{x})`$ can be represented in an unified form
$$\stackrel{}{Y}_{ll^{}m}^\eta (\widehat{x})=(1)^{(l+l^{}+m+1)\mathrm{sin}[\frac{(1\eta )\pi }{4}]}\stackrel{}{Y}_{ll^{},\eta m}(\widehat{x}).$$
(105)
Thus, according to (A.22)- (A.25), we have
$$\begin{array}{c}\mathrm{\Gamma }(l_i,l_i^{},m_i,\eta _i)𝑑\mathrm{\Omega }(\widehat{x})[\stackrel{}{Y}_{l_1l_1^{}m_1}^{\eta _1}(\widehat{x})\times \stackrel{}{Y}_{l_2l_2^{}m_2}^{\eta _2}(\widehat{x})]\stackrel{}{Y}_{l_3l_3^{}m_3}^{\eta _3}(\widehat{x})\hfill \\ =\underset{i=1}{\overset{3}{}}(1)^{(l_i+l_i^{}+m_i+1)\mathrm{sin}[\frac{(1\eta _i)\pi }{4}]}\mathrm{\Gamma }(l_i,l_i^{},\eta _im_i),\hfill \end{array}$$
(106)
where $`\mathrm{\Gamma }(l_i,l_i^{},\eta _im_i)`$, as expressed in Eq.(3.24), is directly written out from (A.22) with replacing $`m_i`$ by $`\eta _im_i.`$
Let us turn to compute the integral in (A.18) following the method proposed by one of the authors in this paper and his coworker in their early publications . As we know, there is a momentum conservation in the gluon three-line vertex: $`\stackrel{}{k_1}+\stackrel{}{\text{ }k_2}+\stackrel{}{\text{ }k_3}=0`$ which gives a certain restriction on the magnitudes of the three momenta. In fact, from $`k_1^2=(\stackrel{}{\text{ }k_2}+\stackrel{}{\text{ }k_3})^2=k_2^2+k_3^2+2k_2k_3\mathrm{cos}\theta _{12},`$ it is seen that when $`\mathrm{cos}\theta _{12}=\pm 1,`$ we have $`k_1=k_2+k_3`$ and $`k_1=\left|k_2k_3\right|`$. This implies that only when the conditions $`k_1+k_2k_3,k_2+k_3k_1`$ and $`k_1+k_3k_2`$ or
$$k_1+k_2k_30,k_2+k_3k_10,k_1+k_3k_20$$
(107)
are simultaneously satisfied, the momentum conservation holds; whereas, when $`k_1>k_2+k_3`$, the momentum conservation is violated. In addition, adding any two inequalities in (A.27) together, we find $`k_i0`$, therefore, each $`k_i`$ varies from zero to infinity. In later derivations, the following relations are useful
$$\begin{array}{c}h_l^{(1)}(x)=j_l(x)+in_l(x),\text{ }h_l^{(2)}(x)=j_l(x)in_l(x),\hfill \\ j_l(x)=(1)^lj_l(x),\text{ }h_l^{(1)}(x)=(1)^lh_l^{(2)}(x),\hfill \end{array}$$
(108)
where $`j_l(x)`$ is the spherical Bessel function, $`n_l(x)`$ the spherical Neumann function, $`h_l^{(1)}(x)`$ and $`h_l^{(2)}(x)`$ are the first class spherical Hankel function and the second class one respectively. The asymptotic behaviors of these functions are as follows. When $`x0`$,
$$\begin{array}{c}j_l(x)\frac{2^ll!x^l}{(2l+1)!},\text{ }n_l(x)\frac{(2l1)!!}{x^{l+1}},\hfill \\ h_l^{(1)}(x)\frac{i(2l)!}{2^ll!x^{l+1}},\text{ }h_l^{(2)}(x)\frac{i(2l1)!!}{x^{l+1}},\hfill \end{array}$$
(109)
and when $`x\mathrm{},`$
$$\begin{array}{c}j_l(x)\frac{i}{2x}[e^{i(x\frac{l\pi }{2})}e^{i(x\frac{l\pi }{2})}],\hfill \\ h_l^{(1)}(x)\frac{i}{x}e^{i(x\frac{l\pi }{2})},h_l^{(2)}(x)\frac{i}{x}e^{i(x\frac{l\pi }{2})}.\hfill \end{array}$$
(110)
First, we prove that the integral $`J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)`$ vanishes in the case of $`k_1>k_2+k_3`$. In this case, considering the analytical property of the functions $`j_l(x)`$ and $`h_l^{(1)}(x)`$ as shown in (A.29) and (A.30), the following integral along the contour $`C`$ on the upper half complex plane of $`r`$ as depicted in Fig.1 is zero
$$_C𝑑rr^2h_{l_1}^{(1)}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)=0.$$
(111)
The contour $`C`$ can be divided into four parts, $`C=C_0+(\mathrm{},0^{})+C_1+(0^+,+\mathrm{})`$. Clearly, the integral along the large half circle $`C_1`$ vanishes when $`\left|r\right|`$ tends to infinity. Thus, noticing the relations in (A.28), $`l_i0`$ and $`l_1+l_2+l_3=even`$ which is implied by the first 3-j symbol in (A.22), one can get from (A.31)
$$J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)=\frac{1}{2}_{C_0}𝑑rr^2h_{l_1}^{(1)}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r).$$
(112)
Substituting the series expansions
| $`j_l(kr)=\frac{\sqrt{\pi }}{2}\underset{\mu =0}{\overset{\mathrm{}}{}}\frac{(1)^\mu (\frac{kr}{2})^{2\mu +l}}{\mathrm{\Gamma }(\mu +1)\mathrm{\Gamma }(\mu +l+\frac{3}{2})},`$ |
| --- |
| $`n_l(kr)=\frac{\sqrt{\pi }}{2}\underset{\mu =0}{\overset{\mathrm{}}{}}\frac{(1)^{l+\mu +1}(\frac{kr}{2})^{2\mu l1}}{\mathrm{\Gamma }(\mu +1)\mathrm{\Gamma }(\mu l+\frac{1}{2})}`$ |
(113)
into the right hand side of (A.32), it is easy to find that the integral along the circle $`C_0`$ around the origin also vanishes when $`\left|r\right|`$ goes to zero. Thus, we reach the following result
$$J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)=0.$$
(114)
Next, we compute the integral under the conditions shown in (A.27). In view of these conditions and the asymptotic behaviors of $`h_l^{(1)}(x)`$ and $`h_l^{(2)}(x)`$ shown in (A.30), the function $`f(r)`$ defined by
$$\begin{array}{c}f(r)=h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(2)}(k_3r)+h_{l_1}^{(1)}(k_1r)h_{l_2}^{(2)}(k_2r)h_{l_3}^{(1)}(k_3r)\hfill \\ +h_{l_1}^{(2)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(1)}(k_3r)+h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(1)}(k_3r)\hfill \end{array}$$
(115)
is analytical on the upper half complex plane of $`r`$ except for at the origin. Therefore, we have
$$_C𝑑rr^2f(r)=0,$$
(116)
where the contour $`C`$ is still represented in Fig.1. Due to the conditions in (A.27), the integral along$`C_1`$ still vanishes. Thus, from the above integral, we get
| $`J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)=\frac{1}{8}\{_0^{\mathrm{}}𝑑rr^2f(r)+_{\mathrm{}}^0𝑑rr^2f(r)\}`$ |
| --- |
| $`=\frac{1}{8}_{C_0}𝑑rr^2f(r).`$ |
(117)
In accordance with (A.28), the function $`f(r)`$ can be written as
| $`f(r)`$ |
| --- |
| $`=4j_{l_1}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)+2in_{l_1}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)`$ |
| $`+2ij_{l_1}(k_1r)n_{l_2}(k_2r)j_{l_3}(k_3r)+2ij_{l_1}(k_1r)j_{l_2}(k_2r)n_{l_3}(k_3r)`$ |
| $`+2in_{l_1}(k_1r)n_{l_2}(k_2r)n_{l_3}(k_3r).`$ |
(118)
Inserting this expression into (A.37) and using the series representation in (A.33), it can be found that except for the last term, the other terms in (A.38) all give no contribution to the integral. Therefore, we have
$$\begin{array}{c}J_{l_1^{}l_2^{}l_3^{}}(k_1,k_2,k_3)=\frac{i}{4}_{C_0}𝑑rr^2n_{l_1}(k_1r)n_{l_2}(k_2r)n_{l_3}(k_3r)\hfill \\ =\frac{i\pi ^{\frac{3}{2}}}{32}\underset{\mu _1,\mu _2,\mu _3}{}\underset{i=1}{\overset{3}{}}\frac{(1)^{\mu _i+l_i}(\frac{k_i}{2})^{2\mu _il_i1}}{\mathrm{\Gamma }(\mu _i+1)\mathrm{\Gamma }(\mu _il_i+\frac{1}{2})}_{C_0}𝑑rr^{2(\mu _1+\mu _2+\mu _3)l_1l_2l_31}.\hfill \end{array}$$
(119)
Setting $`r=\rho e^{i\theta }`$ and noticing $`2(\mu _1+\mu _2+\mu _3)l_1l_2l_3=even`$ and $`_\pi ^0𝑑\theta e^{i2m\theta }=\pi \delta _{m,0}`$ here $`m`$ is an integer, we finally obtain the expression as shown in Eq.(3.23).
## VIII Appendix B: The expression of gluon four-line vertex
In this appendix we would like to derive the explicit expression of the gluon four-line vertex in the angular momentum representation for completeness although the vertex gives no contribution to the equation (2.27) in the lowest order approximation due to $`S_2(\rho \sigma \tau \lambda ,\gamma \delta )=0`$ as shown in Eq.(3.1). The four-line vertex in the interaction Hamiltonian at $`t=0`$ which was described in Eqs.(4.7), (4.29) and (4.31) in paper I may be rewritten as
| $`H_g^4=\frac{g^2}{4}f^{abe}f^{cde}d^3x(\stackrel{}{A^a}\stackrel{}{A^c})(\stackrel{}{A^b}\stackrel{}{A^d})`$ |
| --- |
| $`=\underset{\alpha \beta \gamma \delta }{}B(\alpha _1\alpha _2\alpha _3\alpha _4):𝐚_{\alpha _1}𝐚_{\alpha _2}𝐚_{\alpha _3}𝐚_{\alpha _4}:,`$ |
(120)
where the second equality is obtained by substituting the expansion of gluon fields in terms of the multipole fields into the first equality and the coefficient function is of the form
$$B(\alpha _1\alpha _2\alpha _3\alpha _4)=\frac{g^2}{4}f^{abe}f^{cde}d^3x(\stackrel{}{A}_{\alpha _1}^{\lambda _1}\stackrel{}{A}_{\alpha _2}^{\lambda _2})(\stackrel{}{A}_{\alpha _3}^{\lambda _3}\stackrel{}{A}_{\alpha _4}^{\lambda _4}).$$
(121)
By making use of the representation in (A.8)-(A.12), (A.14) and (A.15) and noticing the definition given in Eq.(4.16) in paper I, the $`B(\alpha _1\alpha _2\alpha _3\alpha _4)`$ can be represented as
| $`B(\alpha _1\alpha _2\alpha _3\alpha _4)=\frac{g^2}{\pi ^2}f^{abe}f^{cde}\underset{i=1}{\overset{4}{}}\frac{k_i}{\sqrt{2\omega (k_i)}}B^{\xi _i}(l_i)_{\lambda _i\tau _i}`$ |
| --- |
| $`\times J_{l_1^{^{}}l_2^{^{}}l_3^{^{}}l_4^{}}(k_1,k_2,k_3,k_4)\stackrel{~}{\mathrm{\Gamma }}(l_i,l_i^{},m_i,\eta _i),`$ |
(122)
where
$$J_{l_1^{}l_2^{}l_3^{}l_4}(k_1,k_2,k_3,k_4)=𝑑rr^2j_{l_1^{}}(k_1r)j_{l_2^{}}(k_2r)j_{l_3^{}}(k_3r)j_{l_4^{}}(k_4r)$$
(123)
and
$$\stackrel{~}{\mathrm{\Gamma }}(l_i,l_i^{},m_i,\eta _i)=𝑑\mathrm{\Omega }(\widehat{x})\stackrel{}{Y}_{l_1l_1^{}m_1}^{\eta _1}(\widehat{x})\stackrel{}{Y}_{l_2l_2^{}m_2}^{\eta _2}(\widehat{x})\stackrel{}{Y}_{l_3l_3^{}m_3}^{\eta _3}(\widehat{x})\stackrel{}{Y}_{l_4l_4^{}m_4}^{\eta _4}(\widehat{x}),$$
(124)
here the notation in (A.24) has been used. Inserting (A.20) and (A.25) into (B.5) and employing the familiar formulas for the integrals of spherical harmonics and for the angular momentum coupling , it is not difficult to get
| $`\stackrel{~}{\mathrm{\Gamma }}(l_i,l_i^{},m_i,\eta _i)`$ |
| --- |
| $`=\frac{1}{4\pi }_{lm}(1)^{l_1+l_2+l_3^{}+l_4^{}+m}(2l+1)_{i=1}^4(1)^{(l_i+l_i^{}+m_i+1)\mathrm{sin}[\frac{(1\eta _i)\pi }{4}]}`$ |
| $`\times [(2l+1)(2l^{}+1)]^{\frac{1}{2}}\left(\begin{array}{ccc}l_1^{}& l_3^{}& l\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_2^{}& l_4^{}& l\\ 0& 0& 0\end{array}\right)\left\{\begin{array}{ccc}l_1& l_3& l\\ l_3^{}& l_1^{}& 1\end{array}\right\}`$ |
| $`\times \left\{\begin{array}{ccc}l_2& l_4& l\\ l_4^{}& l_2^{}& 1\end{array}\right\}\left(\begin{array}{ccc}l_1& l_3& l\\ \eta _1m_1& \eta _3m_3& m\end{array}\right)\left(\begin{array}{ccc}l_2& l_4& l\\ \eta _2m_2& \eta _4m_4& m\end{array}\right).`$ |
(125)
To compute the integral in (B.4), we first examine the conditions satisfied by the magnitudes of the momenta. From the momentum conservation $`\stackrel{}{k_1}+\stackrel{}{k_2}+\stackrel{}{k_3}+\stackrel{}{k_4}=0`$, we have
| $`k_1^2=(\stackrel{}{k_2}+\stackrel{}{k_3}+\stackrel{}{k_4})^2=k_2^2+k_3^2+k_4^2`$ |
| --- |
| $`+2k_2k_3\mathrm{cos}\theta _{23}+2k_2k_4\mathrm{cos}\theta _{24}+2k_3k_4\mathrm{cos}\theta _{34}.`$ |
(126)
From the above equality, we may find the maximum and minimum of $`k_1`$ by setting $`\theta _{ij}=0`$ or $`\pi `$. There are only four ways which permit us to take the values of $`\theta _{ij}=0`$ or $`\pi `$. These are (1) $`\theta _{23}=\theta _{24}=\theta _{34}=0`$; (2) $`\theta _{23}=\theta _{24}=\pi ,\theta _{34}=0`$; (3) $`\theta _{23}=\theta _{34}=\pi ,\theta _{24}=0`$; (4) $`\theta _{24}=\theta _{34}=\pi ,\theta _{23}=0`$. Correspondingly, we get from (B.7) the equalities $`k_1=k_2+k_3+k_4,k_1=\left|k_2k_3k_4\right|,k_1=\left|k_3k_2k_4\right|`$ and $`k_1=\left|k_4k_2k_3\right|.`$ From these equalities we find the following restriction conditions which are consistent with the momentum conservation
| $`k_2+k_3+k_4k_10,k_1+k_3+k_4k_20,`$ |
| --- |
| $`k_1+k_2+k_4k_30,k_1+k_2+k_3k_40,`$ |
| $`k_1+k_2k_3k_40,k_1+k_3k_2k_40,`$ |
| $`k_1+k_4k_2k_30,k_1+k_2+k_3+k_40.`$ |
(127)
By adding some two of the above inequalities, one may see $`k_i0.`$ And similar to the proof described in the preceding appendix, it can be proved that the integral in (B.4) is merely nonvanishing provided that the conditions in (B.8) are respected. According to the above conditions, it is obvious that the function defined below is analytical on upper half complex plane of $`r`$
$$\begin{array}{c}F(r)\hfill \\ =h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(1)}(k_3r)h_{l_4}^{(2)}(k_4r)+h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(2)}(k_3r)h_{l_4}^{(1)}(k_4r)\hfill \\ +h_{l_1}^{(1)}(k_1r)h_{l_2}^{(2)}(k_2r)h_{l_3}^{(1)}(k_3r)h_{l_4}^{(1)}(k_4r)+h_{l_1}^{(2)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(1)}(k_3r)h_{l_4}^{(1)}(k_4r)\hfill \\ +h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(2)}(k_3r)h_{l_4}^{(2)}(k_4r)+h_{l_1}^{(1)}(k_1r)h_{l_2}^{(2)}(k_2r)h_{l_3}^{(1)}(k_3r)h_{l4}^{(2)}(k_4r)\hfill \\ +h_{l_1}^{(1)}(k_1r)h_{l_2}^{(2)}(k_2r)h_{l_3}^{(2)}(k_3r)h_{l_4}^{(1)}(k_4r+h_{l_1}^{(1)}(k_1r)h_{l_2}^{(1)}(k_2r)h_{l_3}^{(1)}(k_3r)h_{l_4}^{(1)}(k_4r)\hfill \end{array}$$
(128)
where the $`l_i^{^{}}`$ have been replaced by $`l_i`$ for convenience. Therefore, based on the Cauchy theorem, the integral along the contour $`C`$ as depicted in Fig.1 is zero,
$$_C𝑑rr^2F(r)=0$$
(129)
From this equation, noticing $`l_1+l_2+l_3+l_4=even`$ as implied by the first two 3-j symbols in (B.6), we obtain
| $`J_{l_1^{}l_2^{}l_3^{}l_4}(k_1,k_2,k_3,k_4)=\frac{1}{16}\{_0^{\mathrm{}}𝑑rr^2F(r)+_{\mathrm{}}^0𝑑rr^2F(r)\}`$ |
| --- |
| $`=\frac{1}{16}_{C_0}𝑑rr^2F(r).`$ |
(130)
In view of the relations in (A.28), the function $`F(r)`$ can be represented as
| $`F(r)`$ |
| --- |
| $`=8j_{l_1}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)j_{l_4}(k_4r)+2j_{l_1}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)n_{l_4}(k_4r)`$ |
| $`+2j_{l_1}(k_1r)j_{l_2}(k_2r)n_{l_3}(k_3r)j_{l_4}(k_4r)+2ij_{l_1}(k_1r)n_{l_2}(k_2r)j_{l_3}(k_3r)j_{l_4}(k_4r)`$ |
| $`+6in_{l_1}(k_1r)j_{l_2}(k_2r)j_{l_3}(k_3r)j_{l_4}(k_4r)+2in_{l_1}(k_1r)j_{l_2}(k_2r)n_{l_3}(k_3r)n_{l_4}(k_4r)`$ |
| $`+2in_{l_1}(k_1r)n_{l_2}(k_2r)j_{l_3}(k_3r)n_{l_4}(k_4r)+2in_{l_1}(k_1r)n_{l_2}(k_2r)n_{l_3}(k_3r)j_{l_4}(k_4r)`$ |
| $`2ij_{l_1}(k_1r)n_{l_2}(k_2r)n_{l_3}(k_3r)n_{l_4}(k_4).`$ |
(131)
Upon inserting the above expression into (B.11), using the series representation in (A.33) and considering the relations among the angular momenta which are implied by the 3-j and 9-j symbols in (B.6), it is easily verified that the first five terms in (B.12) give vanishing contributions to the integral. The nonvanishing contributions given by the last four terms can be calculated by the same procedure as described in the last part of Appendix A. The result is
| $`J_{l_1^{}l_2^{}l_3^{}l_4}(k_1,k_2,k_3,k_4)`$ |
| --- |
| $`=\frac{\pi ^3}{16}()^{\frac{1}{2}(l_1+l_2+l_3+l_4)}\stackrel{\mathrm{}}{\underset{\mu _1,\mu _2,\mu _3,\mu _4=0}{}}\delta _{2(\mu _1+\mu _2+\mu _3+\mu _4),l_2^{}+l_3^{}+l_4^{}l_1^{}}`$ |
| $`\times \underset{i=1}{\overset{4}{}}\frac{(12\delta _{i1})}{\mathrm{\Gamma }(\mu _i+1)\mathrm{\Gamma }(\mu _i+l_i+\frac{3}{2})}\underset{i=1(ji)}{\overset{4}{}}\frac{(k_j/k_1)^{2\mu _jl_j}}{k_j\mathrm{\Gamma }(\mu _j+1)\mathrm{\Gamma }(\mu _jl_j+\frac{1}{2})}.`$ |
(132)
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## X <br>Figure caption
Fig.1. The contour for the integrals containing three and four spherical Bessel functions.
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# Hall coefficient in heavy fermion metals
## Abstract
Experimental studies of the antiferromagnetic (AF) heavy fermion metal $`\mathrm{YbRh}_2\mathrm{Si}_2`$ in a magnetic field $`B`$ indicate the presence of a jump in the Hall coefficient at a magnetic-field tuned quantum state in the zero temperature limit. This quantum state occurs at $`BB_{c0}`$ and induces the jump even though the change of the magnetic field at $`B=B_{c0}`$ is infinitesimal. We investigate this by using the model of heavy electron liquid with the fermion condensate. Within this model the jump takes place when the magnetic field reaches the critical value $`B_{c0}`$ at which the ordering temperature $`T_N(B=B_{c0})`$ of the AF transition vanishes. We show that at $`BB_{c0}`$, this second order AF phase transition becomes the first order one, making the corresponding quantum and thermal critical fluctuations vanish at the jump. At $`T0`$ and $`B=B_{c0}`$, the Grüneisen ratio as a function of temperature $`T`$ diverges. We demonstrate that both the divergence and the jump are determined by the specific low temperature behavior of the entropy $`S(T)S_0+a\sqrt{T}+bT`$ with $`S_0`$, $`a`$ and $`b`$ are temperature independent constants.
The most outstanding puzzle of heavy fermion (HF) metals is what determines their universal behavior which drastically differs from the behavior of ordinary metals. It is wide accepted that the fundamental physics observed in the HF metals is controlled by quantum phase transitions. A quantum phase transition is driven by control parameters such as composition, pressure, number density $`x`$ of electrons (holes), magnetic field $`B`$, etc, and takes place at a quantum critical point (QCP) when the temperature $`T=0`$. In the case of conventional quantum phase transitions (CQPT) the physics is dominated by thermal and quantum fluctuations near CQP. This critical state is characterized by the absence of quasiparticles. It is believed that the absence of quasiparticle-like excitations is the main cause of the non-Fermi liquid (NFL) behavior, see e.g. voj . However, theories based on CQPT fail to explain the experimental observations of the universal behavior related to the divergence of the effective mass $`M^{}`$ at the magnetic field tuned QCP, the specific behavior of the spin susceptibility, its scaling properties, etc.
It is possible to explain the observed universal behavior of the HF metals on the basis of the fermion condensation quantum phase transition (FCQPT) which takes place at $`x=x_{FC}`$ and allows the existence of the Landau quasiparticles down to the lowest temperatures shag2 . It is the quasiparticles which define the universal behavior of the HF metals at low temperatures shag2 ; ckhz . In contrast to the conventional Landau quasiparticles, these are characterized by the effective mass which strongly depends on temperature $`T`$, applied magnetic field $`B`$ and the number density $`x`$ of the heavy electron liquid of HF metal. Thus, we come back again to the key role of the of the effective mass.
On the other hand, it is plausible to probe the other properties of the heavy electron liquid which are not directly determined by the effective mass. Behind the point of FCPT when $`x<x_{FC}`$, the heavy electron liquid possesses unique features directly determined by its quasiparticle distribution function $`n_0(𝐩)`$ formed by the presence of the fermion condensate (FC) ks . Therefore, the function $`n_0(𝐩)`$ drastically differs from the quasiparticle distribution function of a typical Landau Fermi liquid (LFL) lanl1 . For example, it was predicted that at low temperatures the tunneling differential conductivity between HF metal with FC and a simple metallic point can be noticeably dissymmetrical with respect to the change of voltage bias tun . As we shall see below, the magnetic field dependence of the Hall coefficient $`R_H(B)`$ can also provide information about electronic systems with FC.
Recent experiments have shown that the Hall coefficient in the antiferromagnetic (AF) HF metal $`\mathrm{YbRh}_2\mathrm{Si}_2`$ in a magnetic field $`B`$ undergoes a jump in the zero temperature limit upon magnetic-field tuning the metal from AF to a paramagnetic state pash . At some critical value $`B_{c0}`$, the magnetic field $`B`$ induces the jump even though the change of the magnetic field at at the critical value $`B_{c0}`$ is infinitesimal.
In this letter, we show that the abrupt change in the Hall coefficient is determined by the presence of FC and investigate this jump by using the model of the heavy electron liquid with FC which is represented by an uniform electron liquid near FCQPT. Within this model the jump takes place when magnetic field reaches the critical value $`B_{c0}`$ at which the Néel temperature $`T_N(B=B_{c0})`$ of the AF transition vanishes. At some temperature $`T_{crit}`$ when $`BB_{c0}`$, this second order AF phase transition becomes the first order one, making the corresponding quantum and thermal critical fluctuations vanish at the point where $`T_N(B=B_{c0})0`$. At $`T0`$ and $`B=B_{c0}`$, the Grüneisen ratio $`\mathrm{\Gamma }(\mathrm{T})=\alpha (\mathrm{T})/\mathrm{C}(\mathrm{T})`$ as a function of temperature $`T`$ diverges. Here, $`\alpha (T)`$ is the thermal expansion coefficient and $`C(T)`$ is the specific heat. We show that both the divergence and the jump are determined by the specific low temperature behavior of the entropy $`S(T)S_0+a\sqrt{T/T_f}+bT/T_f`$ with $`S_0`$, $`a`$ and $`b`$ being temperature independent constants, and $`T_f`$ is the temperature at which the influence of FC vanishes.
To study the universal behavior of the HF metals at low temperatures, we use the heavy electron liquid model in order to get rid of the specific peculiarities of a HF metal. It is possible since we consider processes related to the power-low divergences of the corresponding physical quantities. These divergences are determined by small momenta transferred as compared to momenta of the order of the reciprocal lattice, therefore, the contribution coming from the lattice can be ignored. On the other hand, we can simply use the common concept of the applicability of the LFL theory when describing electronic properties of metals lanl1 . Thus, we may safely ignore the complications due to the anisotropy of the lattice regarding the medium as the homogeneous heavy electron isotropic liquid.
At first, we briefly describe the heavy electron liquid with FC. Dealing with FCQPT, we have to put $`T=0`$. In that case, the ground state energy $`E_{gs}`$ of a system in the superconducting state is given by the BSC theory formula
$$E_{gs}[\kappa (𝐩)]=E[n(𝐩)]+E_{sc}[\kappa (𝐩)],$$
(1)
where the occupation numbers $`n(𝐩)`$ are connected to the factors $`v(𝐩)`$, $`u(𝐩)`$ and the order parameter $`\kappa (𝐩)`$
$$n(𝐩)=v^2(𝐩);v^2(𝐩)+u^2(𝐩)=1;$$
$$\kappa (𝐩)=v(𝐩)u(𝐩)=\sqrt{n(𝐩)(1n(𝐩))}.$$
(2)
The second term $`E_{sc}[\kappa _p]`$ on the right hand side of Eq. (1) is defined by the superconducting contribution which in the simplest case of the weak coupling regime is of the form
$$E_{sc}[\kappa _p]=\lambda V_{pp}(𝐩_1,𝐩_2)\kappa (𝐩_1)\kappa ^{}(𝐩_2)\frac{d𝐩_1d𝐩_2}{(2\pi )^4},$$
(3)
where $`\lambda V_{pp}(𝐩,𝐩_1)`$ is the pairing interaction. Varying $`E_{gs}`$ given by Eq. (1) with respect to $`v(𝐩)`$ one finds
$$\epsilon (𝐩)\mu =\mathrm{\Delta }(𝐩)\frac{12v^2(𝐩)}{2\kappa (𝐩)}.$$
(4)
Here $`\epsilon (𝐩)`$ is defined by the Landau equation $`\delta E[n(𝐩)]/\delta n(𝐩)=\epsilon (𝐩),`$ $`\mu `$ is chemical potential, and the gap
$$\mathrm{\Delta }(𝐩)=\lambda V_{pp}(𝐩,𝐩_1)\sqrt{n(𝐩_1)(1n(𝐩_1))}\frac{d𝐩_1}{4\pi ^2}.$$
(5)
If $`\lambda 0`$, then $`\mathrm{\Delta }(𝐩)0`$, and Eq. (4) reduces to the equation
$$\frac{\delta E[n(𝐩)]}{\delta n(𝐩)}\mu =\epsilon (𝐩)\mu =0,\mathrm{if}\kappa (𝐩)0.$$
(6)
As a result, at $`x<x_{FC}`$, the function $`n(𝐩)`$ is determined by the standard equation to search the minimum of functional $`E[n(𝐩)]`$ dkss ; vsl . Equation (6) determines the quasiparticle distribution function $`n_0(𝐩)`$ which delivers the minimum value to the ground state energy $`E`$. The function $`n_0(𝐩)`$ being the signature of the new state of quantum liquids vol does not coincide with the step function in the region $`(p_fp_i)`$ where $`\kappa (𝐩)0`$, so that $`0<n_0(𝐩)<1`$ and $`p_i<p_F<p_f`$, with $`p_F=(3\pi ^2x)^{1/3}`$ is the Fermi momentum. We note the remarkable peculiarity of FCQPT at $`T=0`$: this transition is related to spontaneous breaking of gauge symmetry, when the superconducting order parameter $`\kappa (𝐩)=\sqrt{n_0(𝐩)(1n_0(𝐩))}`$ has a nonzero value over the region occupied by the fermion condensate, with the entropy $`S=0`$ vsl ; shag2 , while the gap $`\mathrm{\Delta }(𝐩)`$ vanishes provided that $`\lambda 0`$ dkss ; vsl . Thus the state with FC cannot exist at any finite temperatures and driven by the parameter $`x`$: at $`x>x_{FC}`$ the system is on the disordered side of FCQPT; at $`x=x_{FC}`$, Eq. (6) possesses the non-trivial solutions $`n_0(𝐩)`$ with $`p_i=p_F=p_f`$; at $`x<x_{FC}`$, the system is on the ordered side shag2 .
At finite temperatures $`0<TT_f`$, the function $`n_0(𝐩)`$ determines the entropy $`S_{NFL}(T)`$ of the heavy electron liquid in its NFL state
$$S_{NFL}[n(p)]=2[n(𝐩,T)\mathrm{ln}n(𝐩,T)+(1n(𝐩,T))$$
$$\times \mathrm{ln}(1n(𝐩,T))]\frac{d𝐩}{(2\pi )^3},$$
(7)
with $`T_f`$ being the temperature at which the influence of FC vanishes dkss ; vsl . Inserting into Eq. (7) the function $`n_0(𝐩)`$, one can check that behind the point of FCQPT there is a temperature independent contribution $`S_0(r)(p_fp_F)|r|`$, where $`r=x_{FC}x`$. Another specific contribution is related to the spectrum $`\epsilon (𝐩)`$ which insures the connection between the dispersionless region $`(p_fp_i)`$ occupied by FC and the normal quasiparticles located at $`p<p_i`$ and at $`p>p_f`$, and therefore it is of the form $`\epsilon (𝐩)(pp_f)^2(p_ip)^2`$. Such a form of the spectrum can be verified in exactly solvable models for systems with FC and leads to the contribution of this spectrum to the specific heat $`C\sqrt{T/T_f}`$ ks . Thus at $`0<TT_f`$, the entropy can be approximated as
$$S_{NFL}(T)S_0(r)+a\sqrt{\frac{T}{T_f}}+b\frac{T}{T_f},$$
(8)
with $`a`$ and $`b`$ are constants. The third term on the right hand side of Eq. (8) comes from the contribution of the temperature independent part of the spectrum $`\epsilon (𝐩)`$ and gives a relatively small contribution to the entropy.
The temperature independent term $`S_0(r)`$ determines the specific NFL behavior of the system. For example, the thermal expansion coefficient $`\alpha (T)x(S/x)/x`$ determined mainly by the contribution coming from $`S_0(r)`$ becomes constant at $`T0`$ zver , while the specific heat $`C=TS(T)/T(a/2)\sqrt{T/T_f})`$. As a result, the Grüneisen ratio $`\mathrm{\Gamma }(T)`$ diverges as $`\mathrm{\Gamma }(T)=\alpha (T)/C(T)\sqrt{T_f/T}`$.
We see that at $`0<TT_f`$, the heavy electron liquid behaves as if it were placed at QCP, in fact it is placed at the quantum critical line $`x<x_{FC}`$, that is the critical behavior is observed at $`T0`$ for all $`xx_{FC}`$. At $`T0`$, the heavy electron liquid undergoes a first-order quantum phase transition because the entropy is not a continuous function: at finite temperatures the entropy is given by Eq. (8), while $`S(T=0)=0`$. Therefore, the entropy undergoes a sudden jump $`\delta S=S_0(r)`$ in the zero temperature limit. We make up a conclusion that due to the first order phase transition, the critical fluctuations are suppressed at the quantum critical line and the corresponding divergences, for example the divergence of $`\mathrm{\Gamma }(T)`$, are determined by the quasiparticles rather than by the critical fluctuations as one could expect in the case of CQPT, see e.g. voj . Note that according to the well known inequality, $`\delta QT\delta S`$, the heat $`\delta Q`$ of the transition from the ordered phase to the disordered one is equal to zero, because $`\delta QS_0(r)T0`$ at $`T0`$.
To study the nature of the abrupt change in the Hall coefficient, we consider the case when the LFL behavior arises by the suppression of the AF phase upon applying a magnetic field $`B`$, for example, as it takes place in the HF metals $`\mathrm{YbRh}_2\mathrm{Si}_2`$ and YbRh<sub>2</sub>(Si<sub>0.95</sub>Ge<sub>0.05</sub>)<sub>2</sub> geg ; geg1 . The AF phase is represented by the heavy electron LFL, with the entropy vanishing as $`T0`$. For magnetic fields exceeding the critical value $`B_{c0}`$ at which the Néel temperature $`T_N(BB_{c0})0`$ the weakly ordered AF phase transforms into weakly polarized heavy electron LFL. At $`T=0`$, the application of the magnetic field $`B`$ splits the FC state occupying the region $`(p_fp_i)`$ into the Landau levels and suppresses the superconducting order parameter $`\kappa (𝐩)`$ destroying the FC state. Such a state is given by the multiconnected Fermi sphere, where the smooth quasiparticle distribution function $`n_0(𝐩)`$ in the $`(p_Fp_i)`$ range is replaced by a multiconnected distribution. Therefore the LFL behavior is restored being represented by the weakly polarized heavy electron LFL and characterized by quasiparticles with the effective mass $`M^{}(B)`$ shag2 ; shag
$$M^{}(B)\frac{1}{\sqrt{BB_{c0}}}.$$
(9)
At elevated temperatures $`T>T^{}(BB_{c0})\sqrt{BB_{c0}}`$, the NFL state is restored and the entropy of the heavy electron liquid is given by Eq. (8). This behavior is displayed in the $`TB`$ phase diagram shown in Fig. 1.
In accordance with experimental facts we assume that at relatively high temperatures $`T/T_{NO}1`$ the AF phase transition is of the second order geg . Where $`T_{NO}`$ is the Néel temperature in the absence of the magnetic field. In that case, the entropy and the other thermodynamic functions are continuous functions at the transition temperature $`T_N(B)`$. This means that the entropy of the AF phase $`S_{AF}(T)`$ coincides with the entropy of the NFL state given by Eq. (8),
$$S_{AF}(TT_N(B))=S_{NFL}(TT_N(B)).$$
(10)
Since the AF phase demonstrates the LFL behavior, that is $`S_{AF}(T0)0`$, Eq. (10) cannot be satisfied at sufficiently low temperatures $`TT_{crit}`$ due to the temperature-independent term $`S_0(r)`$, see Eq. (8). Thus, the second order AF phase transition becomes the first order one at $`T=T_{crit}`$ as it is shown in Fig. 1. At $`T=0`$, the critical field $`B_{c0}`$ at which the AF phase becomes the heavy LFL is determined by the condition that the ground state energy of the AF phase coincides with the ground state energy $`E[n_0(𝐩)]`$ of the heavy LFL, that is the ground state of the AF phase becomes degenerated at $`B=B_{c0}`$. Therefore, the Néel temperature $`T_N(BB_{c0})0`$, and the behavior of the effective mass $`M^{}(BB_{c0})`$ is given by Eq. (9), that is $`M^{}(B)`$ diverges when $`BB_{c0}`$. We note that the corresponding quantum and thermal critical fluctuations vanish at $`T<T_{crit}`$ because we are dealing with the first order AF phase transition. We can also safely conclude that the critical behavior observed at $`T0`$ and $`BB_{c0}`$ is determined by the corresponding quasiparticles rather than by the critical fluctuations accompanying second order phase transitions. When $`r0`$ the heavy electron liquid approaches FCQPT from the ordered phase. Obviously, $`T_{crit}0`$ at the point $`r=0`$, and we are led to the conclusion that the Néel temperature vanishes at the point when the AF second order phase transition becomes the first order one. As a result, one can expect that the contributions coming from the corresponding critical fluctuations can only lead to the logarithmic corrections to the Landau theory of the phase transitions lanl2 , and the power low critical behavior is again defined by the corresponding quasiparticles.
Now we are in position to consider the recently observed jump in the Hall coefficient at $`BB_{c0}`$ in the zero temperature limit pash . At $`T=0`$, the application of the critical magnetic field $`B_{c0}`$ suppressing the AF phase (with the Fermi momentum $`p_{AF}p_F`$) restores the LFL with the Fermi momentum $`p_f>p_F`$. At $`B<B_{c0}`$, the ground state energy of the AF phase is lower then that of the heavy LFL, while at $`B>B_{c0}`$, we are dealing with the opposite case, and the heavy LFL wins the competition. At $`B=B_{c0}`$, both AF and LFL have the same ground state energy being degenerated . Thus, at $`T=0`$ and $`B=B_{c0}`$, the infinitesimal change in the magnetic field $`B`$ leads to the finite jump in the Fermi momentum. In response the Hall coefficient $`R_H(B)1/x`$ undergoes the corresponding sudden jump. Here we have assumed that the low temperature $`R_H(B)`$ can be considered as a measure of the Fermi volume and, therefore, as a measure of the Fermi momentum pash . As a result, we obtain
$$\frac{R_H(B=B_{c0}\delta )}{R_H(B=B_{c0}+\delta )}1+3\frac{p_fp_F}{p_F}1+d\frac{S_0(r)}{x_{FC}}.$$
(11)
Here $`\delta `$ is infinitesimal magnetic field, $`S_0(r)/x_{FC}`$ is the entropy per one heavy electron, and $`d`$ is a constant, $`d5`$. It follows from Eq. (11) that the abrupt change in the Hall coefficient tends to zero when $`r0`$ and vanishes when the system in question is on the disordered side of FCQPT.
As an application of the above consideration we study the $`TB`$ phase diagram for the HF metal $`\mathrm{YbRh}_2\mathrm{Si}_2`$ geg shown in Fig. 2. The LFL behavior is characterized by the effective mass $`M^{}(B)`$ which diverges as $`1/\sqrt{BB_{c0}}`$ geg . We can conclude that Eq. (9) gives good description of this experimental fact, and $`M^{}(B)`$ diverges at the point $`BB_{c0}`$ with $`T_N(B=B_{c0})=0`$. It is seen from Fig. 2, that the line separating the LFL state and NFL can be approximated by the function $`c\sqrt{BB_{c0}}`$ with $`c`$ being a parameter. Taking into account that the behavior of YbRh<sub>2</sub>Si<sub>2</sub> strongly resembles the behavior of YbRh<sub>2</sub>(Si<sub>0.95</sub>Ge<sub>0.05</sub>)<sub>2</sub> geg1 ; cust ; geg2 , we can conclude that in the NFL state the thermal expansion coefficient $`\alpha (T)`$ does not depend on $`T`$ and the Grüneisen ratio as a function of temperature $`T`$ diverges geg1 . We are led to the conclusion that the entropy of the NFL state is given by Eq. (8). Taking into account that at relatively high temperatures the AF phase transition is of the second order geg , we predict that at lower temperatures this becomes the first order phase transition. Then, the described behavior of the Hall coefficient $`R_H(B)`$ is in good agreement with experimental facts pash .
In summary, we have shown that the $`TB`$ phase diagram of the heavy electron liquid with FC is in good agreement with the experimental $`TB`$ phase diagram obtained in measurements on the HF metals $`\mathrm{YbRh}_2\mathrm{Si}_2`$ and YbRh<sub>2</sub>(Si<sub>0.95</sub>Ge<sub>0.05</sub>)<sub>2</sub>. We have also demonstrated that the abrupt jump in the Hall coefficient $`R_H(B)`$ is determined by the presence of FC. We observed that at decreasing temperatures $`TT_{crit}`$, the second order AF phase transition becomes the first order one, making the corresponding quantum and thermal critical fluctuations vanish at the jump. Therefore, the abrupt jump and the divergence of the effective mass taking place at $`T_N0`$ are defined by the behavior of quasiparticles rather than by the corresponding thermal and quantum critical fluctuations.
This work was supported by Russian Foundation for Basic Research, Grant No 05-02-16085.
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# Popescu-Rohrlich correlations as a unit of nonlocality
(June 19, 2005)
## Abstract
A set of nonlocal correlations that have come to be known as a PR box suggest themselves as a natural unit of nonlocality, much as a singlet is a natural unit of entanglement. We present two results relevant to this idea. One is that a wide class of multipartite correlations can be simulated using local operations on PR boxes only. We show this with an explicit scheme, which has the interesting feature that the number of PR boxes required is related to the computational resources necessary to represent a function defining the multipartite box. The second result is that there are quantum multipartite correlations, arising from measurements on a cluster state, that cannot be simulated with $`n`$ PR boxes, for any $`n`$.
By performing measurements on an entangled quantum system, two separate observers can obtain correlations that are nonlocal, in the sense that the joint probabilities $`P(a_1a_2|x_1x_2)`$ for the observers to get the outcomes $`a_1`$ and $`a_2`$ given the measurements $`x_1`$ and $`x_2`$ cannot be written in the product form $`P(a_1a_2|x_1x_2)=_jp_jP_j(a_1|x_1)P_j(a_2|x_2)`$ with $`p_j0`$ and $`_jp_j=1`$ Bell (1987). The nonlocal character of the correlations implies that two parties who wish to simulate the experiment with classical resources only cannot do so without communication. Nonlocal correlations, although they cannot be used to signal from one observer to the other, can be exploited in various information processing tasks, such as in communication complexity Brassard (2003), or for the distribution of a secret key between two parties Barrett et al. (2004).
Nonlocality can thus be viewed as an information-theoretic resource and investigated as such Barrettet al. (2005). Forty years after Bell’s seminal paper, however, we still lack a proper theoretical framework — analogous, e.g., to the framework that has been developed for the study of entanglement — that would allow us to answer unambiguously questions such as, can two given sets of nonlocal correlations be considered equivalent resources, or, what is a good measure of nonlocality. In particular, we have not yet identified what would constitute a unit of nonlocality, in the same way that the singlet state constitutes the unit of entanglement.
To progress on these issues, drawing analogies with entanglement is a natural way to proceed. But note a conceptual difference between entanglement and nonlocality: while entanglement is intimately related to the tensorial structure of quantum mechanics, nonlocality, on the contrary, can be defined without reference to quantum theory. In particular, it is possible to write down sets of nonsignaling correlations that are more nonlocal than allowed by quantum mechanics Khalfi and Tsirelson (1985). Why is quantum mechanics not more nonlocal Popescu and Rohrlich (1994)? What are the implications of the quantum restrictions, and what are the principles at the origin of these restrictions van Dam (2005)? A quantitative approach to nonlocality, as required in an information-theoretic perspective, may help us answer these questions.
In this spirit, it is useful to consider nonlocal correlations in the abstract, and not necessarily as arising from a set of measurements on a quantum state. Suppose that two observers have access to a black box. When an observer $`i`$ introduces an input $`x_i`$, the box produces an output $`a_i`$. The box is characterized by the joint probability $`P(a_1a_2|x_1x_2)`$ of obtaining the output pair $`(a_1,a_2)`$ given the input pair $`(x_1,x_2)`$. Compatibility with special relativity requires that these joint probabilities satisfy the nosignaling conditions
$$\underset{a_2}{}P(a_1a_2|x_1x_2)=\underset{a_2}{}P(a_1a_2|x_1x_2^{})P(a_1|x_1)$$
(1)
for all $`a_1,x_1,x_2,x_2^{}`$, as well as a similar set of conditions obtained by summing over the first observer’s outputs. This ensures that one observer cannot signal to the other via his choice of input in the box. Apart from these constraints, the joint distribution can be arbitrary, and in particular nonlocal. The definition of nonlocal boxes generalizes to more parties in a straightforward way.
Some comparisons with entangled quantum states are as follows. Nonlocal boxes and entangled states both represent undirected resources that can be shared between two or more parties. In both cases, the set of allowed states is convex. Extremal elements of this convex set can be thought of as pure states, whereas others are mixed states. There is a notion of monogamy of nonlocality analogous to the monogamy of entanglement Barrettet al. (2005).
Entanglement theory is based on the notion that the entanglement contained in different quantum states can be compared by interconverting them through local operations and classical communication (LOCC), where a basic premise is that LOCC cannot on average increase the entanglement. In the case of bipartite pure states, reversible interconversion is possible, at least asymptotically, and this leads to a unique measure of entanglement. In the case of multipartite states, or bipartite mixed states, reversibility is not in general possible, and this makes the situation more complicated. But still, we do have that any entangled state (bipartite or multipartite) can be obtained from sufficient copies of the singlet state. Similarly, it is possible to study interconversions between nonlocal boxes, i.e., the simulation of nonlocal boxes using other boxes as a resource. In this context, the parties are allowed unlimited access to shared randomness and have the ability to perform local operations, such as relabelling inputs and outputs, or using the output for one box as the input for another box. Communication, however, is not allowed, since it enables trivially the simulation of any nonlocal box. We can now ask the following question. Is there an elementary nonlocal box that allows the simulation of all other boxes, and which could thus be viewed as a unit of nonlocality?
A plausible candidate is the following box, which was described by Kalfhi and Tsirelson Khalfi and Tsirelson (1985), and was independently introduced in a more physical context by Popescu and Rohrlich Popescu and Rohrlich (1994). It takes two inputs $`x_1,x_2\{0,1\}`$ and produces two outputs $`a_1,a_2\{0,1\}`$ according to the joint distribution
$$P(a_1a_2|x_1x_2)=\{\begin{array}{cc}1/2:\hfill & a_1+a_2=x_1x_2mod2\hfill \\ 0:\hfill & \mathrm{otherwise}.\hfill \end{array}$$
(2)
The marginals thus satisfy $`P(0|x_1)=P(1|x_1)=1/2`$ and similarly for $`x_2`$. By Tsirelson’s theorem tsi80 , these correlations cannot be obtained from measurements on any quantum state. Following the denomination used in other works, we will refer to this box as a Popescu-Rohrlich (PR) box. The PR box is maximally nonlocal for the class of two-input two-output boxes, the simplest class of nonlocal boxes Barrettet al. (2005), and is a natural primitive for communication complexity as it allows the solution of any problem with one bit of communication van Dam (2005). In analogy, the singlet state is the maximally entangled state of two qubits, the simplest family of entangled states, and is a natural primitive for entanglement consuming information processing tasks, such as teleportation or dense coding.
To answer our previous question, it is important to determine if all nonlocal correlations can be obtained from PR boxes, or if they are any which cannot. It was shown in Barrettet al. (2005), that all two-input bipartite boxes can be simulated with PR boxes (at least in an approximate sense). It is also known that one PR box is sufficient to reproduce correlations arising from arbitrary von Neumann measurements on the singlet state Cerf et al. (2004). Other examples of correlations that can be simulated with PR boxes are given in Broadbent and Méthot (2005).
In this article, we investigate further the potential of the PR box. We first present a simple protocol that allows the simulation of a large class of two-output boxes. These boxes are natural generalizations of the PR box to more inputs and more parties. It follows from our construction that any $`n`$-partite communication complexity problem can be solved with $`n1`$ bits of communication and a number of PR boxes related to the computational resources necessary to represent its objective function. A second consequence of our result is that any bipartite box with binary outputs can be simulated with PR boxes. We then consider a box that does not belong to the previous class. It arises from measurements of Pauli operators on a cluster state of five qubits. We demonstrate that the corresponding correlations cannot be simulated with PR boxes, or even with arbitrary bipartite boxes.
We now show how PR boxes can simulate a large class of multipartite correlations. These boxes are $`n`$-partite boxes with an arbitrary number of inputs for each party and with binary outputs. We denote the $`n`$-tuple of inputs as $`\stackrel{}{x}=(x_1\mathrm{},x_n)`$, where, without loss of generality, each input $`x_i\{0,\mathrm{},2^m1\}`$ and can thus be represented by a $`m`$-bit string. The $`n`$-tuple of outputs is $`\stackrel{}{a}=(a_1\mathrm{},a_n)`$, where $`a_i\{0,1\}`$. We consider boxes characterized by the following joint distribution
$$P(\stackrel{}{a}|\stackrel{}{x})=\{\begin{array}{cc}1/2^{n1}:\hfill & _ia_i=f(\stackrel{}{x})mod2\hfill \\ 0:\hfill & \mathrm{otherwise},\hfill \end{array}$$
(3)
where $`f(\stackrel{}{x})`$ is a Boolean function of the inputs. Note that the outputs for any subset of $`n1`$ parties are completely random. The only nontrivial correlations thus involve the full set of $`n`$ parties. Boxes of the form (3) are the most general two-output boxes with this property. We will refer to them as full-correlation boxes.
###### Theorem 1
Any full-correlation box can be simulated with PR boxes.
*Proof.* Our proof uses the fact that a Boolean function can be represented as a Boolean circuit, and that the NAND gate, whose action on two input bits $`q`$ and $`r`$ is $`\text{NAND}(q,r)=qr+1`$, constitutes a universal gate for Boolean circuits Papadimitriou (1994). We begin by supposing that the $`n`$ parties have already succeeded in simulating a full correlation $`n`$-partite box with outputs $`\beta _i`$ such that $`_i\beta _i=g_1(\stackrel{}{x})`$, and a full correlation $`n`$-partite box with outputs $`\gamma _i`$ such that $`_i\gamma _i=g_2(\stackrel{}{x})`$. We show how, by using these two boxes along with $`n(n1)`$ PR boxes, they can simulate a single $`n`$-partite box with outputs $`a_i`$ such that $`_ia_i=\mathrm{NAND}(g_1(\stackrel{}{x}),g_2(\stackrel{}{x}))`$. The simulation of a general full-correlation box, using PR boxes only, consists of iterations of this basic building block - one for each NAND gate in a circuit that evaluates $`f(\stackrel{}{x})`$.
Suppose then that each party possesses two bits $`\beta _i`$ and $`\gamma _i`$, such that the $`n`$-bit strings $`\stackrel{}{\beta }`$ and $`\stackrel{}{\gamma }`$ satisfy respectively $`_i\beta _i=g_1(\stackrel{}{x})`$ and $`_i\gamma _i=g_2(\stackrel{}{x})`$. Each pair of parties shares two PR boxes between them (making a total of $`n(n1)`$). Let $`B_{ij}`$, $`ij`$, denote a PR box shared between party $`i`$ and $`j`$. (Our notation is such that the two PR boxes shared between parties $`i`$ and $`j`$ are $`B_{ij}`$ and $`B_{ji}`$.) In box $`B_{ij}`$ party $`i`$ inputs $`\beta _i`$ and gets an output $`b_{ij}`$, while party $`j`$ inputs $`\gamma _j`$ and gets an output $`c_{ij}`$. It thus follows that $`b_{ij}+c_{ij}=\beta _i\gamma _j`$. The final output of party $`i`$ is given by $`a_i=_{ji}\left(b_{ij}+c_{ji}\right)+\beta _i\gamma _i+r_i`$, where $`r_i=1`$ if $`i=1`$, $`r_i=0`$ otherwise. If we sum (modulo 2) the $`n`$ outputs we thus get
$`{\displaystyle \underset{i}{}}a_i`$ $`={\displaystyle \underset{i}{}}{\displaystyle \underset{ji}{}}\left(b_{ij}+c_{ji}\right)+{\displaystyle \underset{i}{}}\beta _i\gamma _i+{\displaystyle \underset{i}{}}r_i`$
$`={\displaystyle \underset{i}{}}{\displaystyle \underset{ji}{}}\left(b_{ij}+c_{ij}\right)+{\displaystyle \underset{i}{}}\beta _i\gamma _i+{\displaystyle \underset{i}{}}r_i`$
$`={\displaystyle \underset{i}{}}{\displaystyle \underset{j}{}}\beta _i\gamma _j+1={\displaystyle \underset{i}{}}\beta _i{\displaystyle \underset{j}{}}\gamma _j+1`$
$`=g_1(\stackrel{}{x})g_2(\stackrel{}{x})+1=\text{NAND}(g_1(\stackrel{}{x}),g_2(\stackrel{}{x})).`$ (4)
Moreover, because the outputs of a PR box are locally completely random, the outputs for any subset of $`n1`$ parties take each of the possible values in $`\{0,1\}^{n1}`$ with equal probability. It follows that each possible value of $`\stackrel{}{a}`$ consistent with $`_ia_i=\text{NAND}(g_1(\stackrel{}{x}),g_2(\stackrel{}{x}))`$ occurs with probability $`1/2^{n1}`$ as required in (3).
Consider again a circuit comprised of NAND gates that evaluates $`f(\stackrel{}{x})`$, such that the inputs to the final NAND gate in the circuit are $`g_1(\stackrel{}{x})`$ and $`g_2(\stackrel{}{x})`$. If we assume that we can already simulate the boxes characterized by $`g_1`$ and $`g_2`$, then we can simulate the $`f(\stackrel{}{x})`$ box. But $`g_1(\stackrel{}{x})`$ is itself the output of a NAND gate with inputs $`h_1(\stackrel{}{x})`$ and $`h_2(\stackrel{}{x})`$, so we can simulate the $`g_1`$ box with the same construction. We keep working backwards until we reach the point where inputs to NAND gates are simply the input bits themselves. But the corresponding boxes are local and can be simulated without PR boxes. $`\mathrm{}`$
###### Corollary 1
Any $`n`$-partite communication complexity problem can be solved with $`n1`$ bits of communication and at most $`kn(n1)`$ PR boxes, where $`k`$ is the size of the smallest circuit comprised of NAND gates that computes $`f(\stackrel{}{x})`$.
In communication complexity, $`n`$ parties are each given an input $`x_i`$ and must compute a function $`f(\stackrel{}{x})`$ of their joint inputs while communicating as little as possible. This problem can easily be solved if the parties share a box of the form (3). It suffices that each party introduces his input into the box and communicates his output to the first party, who can recover the value of $`f(\stackrel{}{x})`$ by summing all the outputs. Corollary 1 then simply follows from Theorem 1.
###### Corollary 2
Any two-output bipartite box can be simulated with PR boxes.
We sketch the proof.
*Proof:* For a fixed number of inputs and outputs, the set of nonlocal bipartite boxes is a convex polytope Barrettet al. (2005). To simulate any two-output bipartite box, it is thus sufficient to simulate every box which is a vertex of the corresponding polytope, since the others can be obtained as mixtures of vertices. Further, it is sufficient to focus on genuine two-output boxes, that is on boxes such that for every input $`x_1`$, $`P(0|x_1)>0`$ and $`P(1|x_1)>0`$, and similarly for every input $`x_2`$. Indeed, if a box satisfies $`P(0|x_1)=0`$ or $`P(1|x_1)=0`$ for some $`x_1`$, it is straightforward that we can simulate this box if we can simulate the box obtained from it by removing input $`x_1`$. Finally, it is then easy to show that every genuine two-output extremal box is of the form (3). (This can be done, for example, by adapting a proof of Barrettet al. (2005), where the polytope of two-input $`d`$-output boxes is characterized.) $`\mathrm{}`$
We have noted that the outputs for any subset of fewer than $`n`$ parties are completely random in boxes of the form (3) and that these are the most general boxes with this property. We now give an example of correlations that are not of that form and that cannot be simulated with PR boxes. The fact that there exist nontrivial correlations between subsets of the parties is crucial to prove this fact. The correlations arise from spin measurements on a one dimensional cluster state of five qubits in a ring br01 . Cluster states are remarkable in that they can act as a universal substratum for measurement based computation rb01 , as well as playing a role in quantum error correction sw01 . Their nonsimulability by PR boxes is thus another interesting property. The correlations can be described as follows. Let inputs $`x=0`$ and $`x=1`$ correspond to spin measurements in the $`\sigma _z`$ and $`\sigma _x`$ bases, and let $`a_i`$ be the output of party $`i`$ for input $`x_i=0`$, and $`a_i^{}`$ for input $`x_i=1`$. Then we have
$$a_1+a_2^{}+a_3=0mod2,$$
(5)
and cyclic permutations of the parties, together with
$$a_1^{}+a_2^{}+a_3^{}+a_4^{}+a_5^{}=1mod2.$$
(6)
These correlations were described in DiVincenzo and Peres (1997) and constitute a GHZ-type paradox. (If values are assigned locally to $`a_i`$ and $`a_i^{}`$ for each party, then on summing the left hand sides together one obtains $`0`$; summing the right hand sides, on the other hand, gives $`1`$.)
###### Theorem 2
The correlations of Eqs. (5) and (6) cannot be simulated exactly by parties who share random data and $`n`$ PR boxes, for any $`n`$.
In order to prove this result, it is useful to describe how a general protocol that aims to simulate these correlations with boxes will work. Let the value of the shared random data on a given round be $`\lambda `$, with probability $`P(\lambda )`$. Consider a particular party, Alice say, and denote her measurement $`x`$. Suppose that Alice shares $`m`$ boxes with other parties. Label these boxes $`B_1,\mathrm{},B_m`$. Alice proceeds as follows.
1. She puts an input $`y_1`$ into box $`B_{i_1}`$, where $`i_1=i_1(\lambda ,x)`$, and $`y_1=y_1(\lambda ,x)`$. She obtains an output $`\alpha _1`$.
2. She puts an input $`y_2`$ into box $`B_{i_2}`$, where $`i_2=i_2(\lambda ,x,\alpha _1)`$, and $`y_2=y_2(\lambda ,x,\alpha _1)`$. She obtains an output $`\alpha _2`$.
3. She continues in this fashion until all $`m`$ boxes have been used. Her final output is a function $`a=a(\lambda ,x,\alpha _1,\mathrm{},\alpha _m)`$.
Finally, a different strategy along these lines may of course be defined for each party.
Note that the nosignaling conditions ensure that the correlations $`P(a_2\mathrm{}a_n|x_2\mathrm{}x_n)`$ restricted to the other parties do not depend on the specifics of Alice’s protocol. In particular they do not depend on wether Alice uses her boxes or not.
*Proof of Theorem 2:* We begin by supposing that there is a protocol for simulating the above correlations exactly using $`n`$ PR boxes. Then we show that this implies that there is a protocol for simulating the correlations exactly using shared random data alone. This we know to be impossible.
A protocol for simulating the correlations must of course produce outputs correlated according to the six equations described by (5) and (6). Consider the case in which Alice’s measurement is $`x_1=0`$. She follows some strategy as described above. Suppose that there is some set of values of $`\lambda ,\alpha _1,\mathrm{},\alpha _{m1}`$, occurring with non-zero probability, such that $`a_1(\lambda ,x_1=0,\alpha _1,\mathrm{},\alpha _{m1},\alpha _m=0)a_1(\lambda ,x_1=0,\alpha _1,\mathrm{},\alpha _{m1},\alpha _m=1)`$. Suppose that the $`m^{\mathrm{th}}`$ PR box in this sequence is shared with party 5. Eq. (5) implies that the outputs of parties 1, 2 and 3 should be correlated according to $`a_1+a_2^{}+a_3=0`$. Since this equation does not involve party 5, we may as well assume, by nosignaling, that party 5 does not use her half of the box. Then $`\alpha _m`$ is random and uncorrelated with the rest of the protocol. This means that, conditioned on the specified values of $`\lambda ,\alpha _1,\mathrm{},\alpha _{m1}`$ occurring, Alice has a $`1/2`$ chance of outputting a value $`a_1`$ that is not correctly correlated with the outputs of parties 2 and 3. Furthermore, it does not matter whom the $`m^{\mathrm{th}}`$ box is shared with. For any of the other parties, there is one equation (5) such that Alice’s input is $`x_1=0`$ and the other party is not involved.
We can conclude from the above that $`a_1(\lambda ,x_1=0,\alpha _1,\mathrm{},\alpha _{m1},\alpha _m=0)=a_1(\lambda ,x_1=0,\alpha _1,\mathrm{},\alpha _{m1},\alpha _m=1)`$ for all values of $`\lambda ,\alpha _1,\mathrm{},\alpha _{m1}`$. This means that in fact, Alice never needs to know the value of $`\alpha _m`$, and her strategy may as well terminate before putting an input into the $`m^{\mathrm{th}}`$ box. But now we can run an identical argument for the $`(m1)^{\mathrm{st}}`$ box, concluding that Alice’s strategy terminates after the $`(m2)^{\text{nd}}`$ box, and so on. We conclude that if errors are not tolerated, Alice’s output, in the event her measurement is $`x_1=0`$, must be fixed by $`\lambda `$ alone, and she does not use her boxes. We can then run this argument for each of the five parties.
Finally, consider Alice’s strategy in the event her measurement is $`x_1=1`$. Then her output $`a_1^{}`$ should satisfy the constraint $`a_1^{}+a_2+a_5=0`$. We have already established that in this case, parties 2 and 5 do not use PR boxes, and output values that depend only on $`\lambda `$. Thus Alice’s output must also be deterministic, and fixed by $`\lambda `$. The PR boxes in her possession simply output values that are random and uncorrelated with the rest of the protocol, so she cannot use them. This argument may now be run for each of the five parties.
We have established that none of the five parties in fact use the PR boxes in their possession, for any of the measurements $`x_i=0,1`$. Thus if they are producing exactly the required correlations, they are doing so using only shared random data. This we know to be impossible, which concludes the proof. $`\mathrm{}`$
It is straightforward to modify the proof we have just given to show that the correlations (5) and (6) cannot be reproduced with any bipartite box (including boxes with more inputs or outputs than a PR box).
*Conclusion.* We have shown that PR boxes can be used to simulate a large class of correlations. These include all bipartite boxes with binary outputs, and therefore, for example, von Neumann measurements on non-maximally entangled states of two qubits (Cerf et al. have shown that a single PR box can simulate a maximally entangled state Cerf et al. (2004)). This encourages the idea that PR boxes should be considered as a proper unit of bipartite nonlocality. In a multipartite scenario, we have seen that PR boxes cannot simulate correlations arising from measurements on a five-qubit cluster state. One could consider the approximate simulation of these correlations, where one demands that the error can be made arbitrary small. It seems reasonable to think that even in this case, PR boxes cannot simulate the cluster state correlations we described, in which case they cannot be considered as a unit of nonlocality in a multipartite scenario.
It thus appears that the structure of nonlocal correlations is rather different from the structure of entanglement. There are many open questions, for example, are there $`n`$-partite boxes such that any set of $`(nk)`$-partite boxes are not sufficient for their simulation? Can we define a finite set of boxes that would be sufficient for the simulation of all $`n`$-partite boxes? One obvious reason for the difference between entanglement and nonlocal boxes may be the fact that we have not allowed classical communication in our protocols for manipulating boxes. An interesting extension of these ideas would be to introduce secrecy, and to see how boxes can be transformed in the presence of public, but not private communication.
*Note added.* After completion of this work, N. Jones and Ll. Masanes informed us that they independently derived the result of Corollary 2 Masanes et al (2005).
###### Acknowledgements.
S. P. acknowledges support by the David and Alice Van Buuren fellowship of the Belgian American Educational Foundation and by the National Science Foundation under Grant No. EIA-0086038.
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# Vibration-assisted tunneling through competing molecular states
## I Model
Molecule. We consider the minimal model $`H=H_M+_rH_r+H_T`$ incorporating the molecule ($`M`$), the electrodes $`r=L,R`$ and the tunneling ($`T`$) in units $`\mathrm{}=k_B=1`$:
$`H_M`$ $`=`$ $`{\displaystyle \underset{i}{}}\left(ϵ_in_i+u_in_{}n_{}\right)+vn_1n_2`$ (1)
$`+{\displaystyle \frac{\omega }{2}}\left[P^2+\left(Q{\displaystyle \underset{i}{}}\sqrt{2}\lambda _in_i\right)^2\right],`$
$`H_r`$ $`=`$ $`{\displaystyle \underset{ki\sigma }{}}ϵ_{kr}a_{k\sigma r}^{}a_{k\sigma r},`$ (2)
$`H_T`$ $`=`$ $`{\displaystyle \underset{ki\sigma r}{}}t_{ir}a_{k\sigma r}^{}c_{i\sigma }+h.c..`$ (3)
The Hamiltonian $`H_M`$ describes a molecule with two electron-accepting orbitals $`i=1,2`$ (operators $`c_{i\sigma }`$, energies $`ϵ_i`$) with an energy splitting $`\mathrm{\Delta }=ϵ_2ϵ_1`$. Here $`n_i=_\sigma n_{i\sigma }`$ , $`n_{i\sigma }=c_{i\sigma }^{}c_{i\sigma }`$. We assume throughout that double occupation of each orbital is completely suppressed due to strong local Coulomb interactions $`u_i`$ i.e. $`n_i1`$. Simultaneous occupation of the two orbitals is similarly suppressed by the interaction $`v`$ which introduces an important correlation: $`_in_i1`$. Since $`u_i`$ and $`v`$ are the largest energy scales they do not enter explicitly in any further way. The nuclear configuration of the molecule is assumed to be altered with respect to some coordinate $`Q`$ when either orbital is occupied. The effective potentials for the nuclear motion in these charged states (here in the harmonic approximation, frequency $`\omega `$) are shifted relative to the neutral state by $`\sqrt{2}\lambda _i,i=1,2`$. The coordinate $`Q=(b+b^{})/\sqrt{2}`$ is normalized to the nuclear zero-point motion by $`(M\omega )^{1/2}`$ ($`M=`$ nuclear mass involved) and $`\lambda _i`$ is dimensionless. Here $`b^{}`$ excites the vibrational mode by one quantum $`\omega `$ and $`P=(bb^{})/\sqrt{2}i`$ is the conjugate nuclear momentum. The energy scale characterizing the electron-vibration coupling associated with orbital $`i=1,2`$ is the change in the elastic energy at fixed nuclear configuration $`\omega \lambda _i^2`$. Vibration-assisted processes are thus expected to lead to a progression of conductance resonances spread over a bias voltage range of at least $`\mathrm{max}\{\lambda _i^2\omega \}`$. Typical energies $`\omega `$ of internal vibrations observed experimentally range up to a few tens of meV Park et al. (2002); Pasupathy et al. (2005); Yu et al. (2004); Yu and Natelson (2004a, b). In general electron-vibration coupling is expected to be particularly strong for many-particle states of molecules which are characterized to a good approximation by occupation of a particular orbital Koeppel et al. (1984), which is typically the case for charged states of otherwise neutral molecules. Large relative displacements $`|\lambda _i|>1`$ of the nuclear potentials may be expected, for instance, when $`Q`$ is an angle coordinate and the nuclear configuration of the charged molecule is internally twisted relative to the neutral one Cizek et al. .
The electronic energy parameters in Eq. 1 are the relevant effective parameters for finite $`\lambda _i`$. These are related to their values for $`\lambda _1=\lambda _2=0`$ (indicated by a superscript (0)) by $`ϵ_i=ϵ_i^{(0)}\lambda _i^2\omega ,u_i=u_i^{(0)}2\lambda _i^2\omega ,v=v^{(0)}2\lambda _1\lambda _2\omega `$. This is seen by diagonalizing the molecular Hamiltonian through a translation of the nuclear coordinates Lang and Firsov (1963), $`U=_{i=1,2}e^{\lambda _in_i(b^{}b)}`$. The resulting Hamiltonian has the form (1) where $`\lambda _i,i=1,2`$ is eliminated. The electron becomes “dressed” with vibrational excitations (polaron) resulting in the renormalization of the energy parameters which we anticipated in writing Eq. 1. The renormalization of the charging energies is irrelevant here since we assume them to be largest energy scales (i.e. $`u_i^{(0)}2|\lambda _i|^2\omega ,i=1,2`$ and $`v^{(0)}2|\lambda _1\lambda _2|\omega `$). The correlations $`n_1+n_21,n_i1`$ are thus not weakened. Cases where a strong renormalization of the interaction becomes relevant were discussed in Cornaglia et al. (2004, ). In contrast to the one-orbital model the renormalization of the orbital energies is important when the coupling to the vibration is asymmetric, $`\lambda _1\lambda _2`$. Then the electronic splitting is renormalized to an effective value $`\mathrm{\Delta }=ϵ_2ϵ_1=\mathrm{\Delta }^{(0)}+\omega (\lambda _1^2\lambda _2^2)`$ which can even have a different sign as $`\mathrm{\Delta }^{(0)}`$. Since only the excitation energy $`\mathrm{\Delta }`$ is observable in the transport characteristics, we use it as an independent positive parameter i.e. the state $`1`$ by definition has the lowest renormalized energy $`ϵ_1`$. We are interested in the case where resonances related to orbital and vibrational excitations occur on the same voltage scale, i.e. $`\omega \mathrm{\Delta }\mathrm{max}\{\lambda _i^2\omega \}`$. The transport mechanism which we wish to illustrate operates in the limit of asymmetric coupling. This requires that either the lowest orbital couples strongly to the vibration or the excited orbital, see Fig. 1. We point out that in Hamiltonian (1) we have not written intramolecular terms which couple the two nuclear potentials 1 and 2. Such terms become important, for instance, when the electronic energy splitting is an integer multiple $`p`$ of the vibrational energy quantum, $`\mathrm{\Delta }=p\omega `$. This has been discussed in Ref. Kaat and Flensberg for the case $`\mathrm{\Delta }=0`$ and $`\lambda _1=\lambda _2`$. Here we avoid such degeneracies, i.e. we assume $`tmin_p\{\mathrm{\Delta }p\omega \}`$, where $`t`$ is a tunneling amplitude between the electronic states. Furthermore, we can safely disregard electronic transitions induced by the nuclear motion near the crossing of the potentials 1 and 2 since for large asymmetric coupling the barrier separating the minima of potentials $`i=1,2`$ is $`(\mathrm{\Delta }/\omega (\lambda _1\lambda _2)^1\pm (\lambda _1\lambda _2))^2\omega /4\omega `$. Below we will also present results for cases of moderate asymmetry of the vibrational couplings where such effects may start to play a role. These will serve as a simple starting point for the discussion of the strong asymmetry case and also illustrates the enhancement of NDC effects when multiple orbitals (instead of just one) are competing in the transport.
The electrodes $`r=L,R`$ are modeled by $`H_r`$, Eq. (2), as non-interacting quasi-particle reservoirs at electro-chemical potential $`\mu _r`$. The electrode-molecule tunneling $`H_T`$, Eq. (3), picks up the shift of the nuclear coordinate from the unitary transformation of the molecular operators: $`H_T=_{ki\sigma r}t_{ir}a_{k\sigma r}^{}e^{\lambda _i(b^{}b)}\overline{c}_{i\sigma }+h.c.`$ Here $`\overline{c}_{i\sigma }^{}`$ creates a polaron state associated with the effective potential of electronic orbital $`i=1,2`$. Since we consider here an intramolecular vibration we do not include a dependence of the bare tunneling matrix elements $`t_{ir}`$ on the coordinate $`Q`$ (shuttle-effect, cf. Ref. McCarthy et al. (2003); Koch et al. (c)).
Master equations. We are interested in the weak tunneling regime, $`\mathrm{\Gamma }T`$, where in addition the vibrational excitations can be resolved, $`T\omega `$. We can describe the transport using diagonal density matrix elements $`P_q^i`$ (occupation probabilities, $`_{i=0}^2_qP_q^i=1`$) , where $`q=0,1,2,\mathrm{}`$ is the vibrational number and $`i=1,2`$ denotes the charged state with only orbital $`i`$ occupied and the neutral state is labeled by $`i=0`$. The transitions including the vibrational excitations $`q,q^{}`$ are denoted by $`0_q1_q^{}`$ and $`0_q2_q^{}`$. The occupation probabilities are coupled by the stationary master equations
$`\dot{P}_q^0=0`$ $`={\displaystyle \underset{i}{}}{\displaystyle \underset{rq^{}}{}}\left(W_{0qiq^{}}^rP_q^{}^i2W_{iq^{}0q}^rP_q^0\right),`$ (4)
$`\dot{P}_q^i=0`$ $`={\displaystyle \underset{rq^{}}{}}\left(2W_{iq0q^{}}^rP_q^{}^0W_{0q^{}iq}^rP_q^i\right),`$
where $`i=1,2`$, with transition rates ($`f_r(ϵ)(e^{(ϵ\mu _r)/T}+1)^1`$)
$`W_{iq^{}0q}^r`$ $`=`$ $`\mathrm{\Gamma }_{qq^{}}^{ri}f_r(\mu _{q^{}q}^i),`$
$`W_{0qiq^{}}^r`$ $`=`$ $`\mathrm{\Gamma }_{qq^{}}^{ri}(1f_r(\mu _{q^{}q}^i)).`$ (5)
The addition energies for the transition $`i_q^{}0_q`$ are
$$\mu _{q^{}q}^i=ϵ_i+(q^{}q)\omega \alpha V_g.$$
(6)
The gate voltage $`V_g`$ effectively varies $`\mu `$ relative to the ground-state transition energy $`\mu _0^1`$ ($`\alpha =`$ capacitance ratio) and the bias voltage $`V>0`$ is applied symmetrically, $`\mu _{L,R}=\mu \pm V/2`$. The stationary current flowing out of reservoir $`r=L,R`$ is given by ($`I_L+I_R=0`$)
$$I_r=\underset{qq^{}}{}\underset{i}{}\left(2W_{iq0q^{}}^rP_q^{}^0W_{0qiq^{}}^rP_q^{}^i\right).$$
(7)
The equations for the one-orbital case with coupling $`\lambda _1`$ are obtained by simply discarding all $`P_q^{}^2`$ in Eqs. (4) and (7) and are equivalent to those in Refs. Boese and Schoeller (2001); McCarthy et al. (2003); Mitra et al. (2004); Koch and von Oppen (2005); Koch et al. (a). The current will change whenever a line in the $`(\mu ,V)`$ plane is crossed corresponding to a right-electrode resonance $`\mu _R=\mu _{q^{}q}^i`$ (positive slope in $`(\mu ,V)`$ plane) or a left-electrode resonance $`\mu _L=\mu _{q^{}q}^i`$ (negative slope). Importantly, due to the harmonic excitation spectrum only the change in vibrational number $`q^{}q`$ enters in the resonance condition: transitions between all states $`i_q^{}`$ and $`0_q`$ with fixed difference $`q^{}q`$ become allowed at a single resonance. A cascade of single-electron transitions can then lead to a significant population of high vibrational excitations, e.g. $`i_00_0i_10_1i_20_2\mathrm{}`$ is a possible cascade for $`\mu _L>\mu _1^i>\mu _0^i>\mu _R`$.
Franck-Condon factors and Condon parabola. The tunneling rates consist of two factors: $`\mathrm{\Gamma }_{qq^{}}^{ri}=\mathrm{\Gamma }^{ri}F_{qq^{}}^i`$. The tunneling rates $`\mathrm{\Gamma }^{ri}=2\pi |t_{ir}|^2\rho _r`$ between electrode $`r=L,R`$ (density of states $`\rho _r`$) and orbital $`i=1,2`$ determine the overall current scale. The FC factors $`F_{q^{}q}^i=F_{qq^{}}^i`$ take into account that the stable nuclear geometry is changed when occupying orbital $`i`$:
$`F_{qq^{}}^i=|q|X_i|q^{}|^2=e^{\lambda _i^2}{\displaystyle \frac{q!}{q^{}!}}\lambda _i^{2|qq^{}|}\left(L_q^{|qq^{}|}(\lambda _i^2)\right)^2,`$ (8)
where $`L`$ is the associated Laguerre-polynomial and $`q<q^{}`$. Note that the sign of $`\lambda _i`$ does not play a role. The general sum rule $`_qF_{qq^{}}^i=_q^{}F_{qq^{}}^i=1`$, guarantees that the current will saturate at large bias voltage to the value it would have without the vibrations (“electronic limit”). This holds only when the $`\lambda _i`$ do not depend strongly on the bias voltage (cf. McCarthy et al. (2003)).
Without vibrations, asymmetry of the tunneling rates with respect to the orbital- and electrode- index gives rise to NDC Hettler et al. (2002a) and super-poissonian current noise Thielmann et al. (2004), see also Cottet et al. (2004a, b); Belzig (2005). Below we show that qualitatively different dependence of the FC factors $`F_{qq^{}}^i`$ on the vibrational numbers for state $`i=1,2`$ and the effective energy splitting $`\mathrm{\Delta }`$ give rise to NDC effects which can dominate the transport. We therefore set $`\mathrm{\Gamma }^{ir}=\mathrm{\Gamma },i=1,2,r=L,R`$ and restrict our discussion to $`V>0`$ since in this case $`I(V)=I(V)`$. We symmetrize the stationary current $`I=(I_LI_R)/2`$ and decompose it into a sum of positive partial currents of the states weighted with their occupation:
$`I`$ $`=`$ $`{\displaystyle \underset{q}{}}I_q^0P_q^0+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \underset{q^{}}{}}I_q^{}^iP_q^{}^i`$ (9)
$`I_q^0`$ $`=`$ $`\mathrm{\Gamma }{\displaystyle \underset{q^{}}{}}{\displaystyle \underset{i}{}}F_{qq^{}}^i\left(f_L(\mu _{q^{}q}^i)f_R(\mu _{q^{}q}^i)\right)`$ (10)
$`I_q^{}^i`$ $`=`$ $`\mathrm{\Gamma }{\displaystyle \underset{q}{}}F_{qq^{}}^i\left(f_L(\mu _{q^{}q}^i)f_R(\mu _{q^{}q}^i)\right)`$ (11)
For low $`T\omega `$ the partial currents are the FC factors $`F_{qq^{}}^i`$ summed over the transitions $`0_qi_q^{}`$ inside the bias window $`\mu _L>\mu _{q^{}q}^i>\mu _R`$ in Fig. 2. Note that the partial current of the neutral state $`I_q^0`$ has contributions from both charged states $`i=1,2`$ into which it can decay. One can understand the numerical results in almost all detail using the following simple graphical scheme. This approach works for multiple orbitals and also for a more general shape of the nuclear potential. Without going into the details, we comment on the basic points in the procedure. In Fig. 2 the FC factor associated with the transition between a pair of states $`0_qi_q^{}`$ is depicted as function of $`q`$ and $`q^{}`$. The change in the partial currents in (9) with increasing bias can be understood by drawing the bias voltage window in this figure. Only qualitative features of the FC factors are of importance which follow from simple quasi-classical arguments. The FC-factor $`F_{qq^{}}^i`$ is basically non-zero only in the classically allowed region delimited by the tilted Condon parabola Herzberg (1950):
$$q+q^{}\frac{|qq^{}|^2}{\lambda _i^2}+\frac{\lambda _i^2}{2}$$
(12)
In this region the classical orbits of the nuclear motion in the shifted potentials intersect in phase-space. The FC-factor oscillates with $`q,q^{}`$ with a quasi-classical envelope which varies algebraically on the scale of $`\lambda _i^2`$ Perelomov (1986). The global maxima $`1/\lambda _i^2`$ are attained where the parabola touches the axes ( $`q\lambda _i^2,q^{}=0`$ and $`q=0,q^{}\lambda _i^2`$). In the classically forbidden region where the opposite of condition (12) holds the FC factor $`F_{qq^{}}^i`$ is exponentially small. In the forbidden regions where $`q^{}q`$ or $`q^{}q`$ the nuclear momenta of the two motions are incompatible. These regions exist for any value of $`\lambda _i`$. Starting from the allowed region the FC-factors eventually decrease exponentially with increasing $`q^{}`$ or $`q`$. On the other hand, the forbidden region where $`q,q^{}\lambda _i^2`$ is only well defined for $`\lambda _i^21`$. In this case $`F_{qq^{}}^i`$ initially increases exponentially with increasing $`q^{}`$ or $`q\lambda _i^2`$. This “inverted” regime exists for nuclear potentials of more general shape with large relative displacements of the potential minima.
As the bias window in Fig. 2 widens, the partial currents increase. For weak coupling $`\lambda _i^21`$ the Condon-parabola is very narrow i.e. transitions which conserve the vibrational number have the largest amplitude. When the bias window reaches the vertex of this narrow parabola nearly all partial currents reach their maximal value at once. For strong coupling, $`\lambda _i^21`$, the parabola is very broad and the partial currents show a slow exponential increase as the bias window widens. The complex transport characteristics of the two orbital model considered here follow from two intramolecular asymmetries between the orbitals: (1) the bias window covers different parts of the Condon-parabola due to the electronic splitting $`\mathrm{\Delta }`$ and (2) the Condon-parabolas are qualitatively different due to asymmetric coupling to the vibration.
Strong relaxation. We now prove an important restriction on the occurrence of NDC in the limit where the vibrational excitations completely relax before each tunneling event due to some dissipative environment. This limit implies the factorization ansatz $`P_q^i=P^iP_q`$ with the vibrational equilibrium distribution $`P_q=e^{q\omega /T}(1e^{\omega /T})`$. We can then reduce the equations (4) to an effective electronic three-level problem with effective bias-voltage dependent rates (cf. Braig and Flensberg (2004); Mitra et al. (2004)) obtained by averaging over the equilibrium distribution:
$$W_{0i}^r=\underset{qq^{}}{}W_{0qiq^{}}^rP_q^{},W_{i0}^r=\underset{qq^{}}{}W_{iq0q^{}}^rP_q^{}$$
(13)
(these vary monotonically with $`V`$) and $`W_{0i}=_rW_{0i}^r`$ ,$`W_{i0}=_rW_{i0}^r`$. The stationary current in electrode $`r=L,R`$ reads ( $`I_L+I_R=0`$):
$`I_r=2{\displaystyle \frac{_i\left(W_{i0}^rW_{0i}^r\frac{W_{i0}}{W_{0i}}\right)}{1+2_i\frac{W_{i0}}{W_{0i}}}}.`$ (14)
In Appendix A we show that in this case for two orbitals the current can be reduced by increasing the bias voltage at resonances related to the left electrode, $`\mu _L=\mu _k^i,k=1,\mathrm{}`$ (for $`V>0`$), i.e. lines with positive slope in the $`(\mu ,V)`$ plane. Here the transition $`i_{q+k}0_q`$ becomes allowed. Any NDC along a resonance line with positive slope is thus a proof of a non-equilibrium vibrational distribution on the molecule. For one orbital the current can never decrease with $`V`$ in this limit. Finally, we consider the limit where in addition to the strong relaxation, the transitions $`0_qi_q^{}`$, $`i=1,2`$ are not correlated i.e. the renormalized Coulomb interaction Cornaglia et al. (2004, ) is zero: $`v^{}=0`$. It is readily shown (Appendix B) that in this case the current increases monotonically with bias voltage for any number of orbitals. The strong Coulomb correlations are thus essential for NDC effects.
Intermediate relaxation. Basically all the physics is captured by considering the opposite limits of negligible and strong relaxation. We have confirmed this by considering intermediate regimes where we add a relaxation term $`_q^{}W_{qq^{}}P_q^{}^i`$ to the right-hand side of the equation for $`\dot{P_q^i}`$, Eq. (4). We considered an environment Boese and Schoeller (2001) with either ohmic ($`s=1`$) or sub-ohmic ($`s=0`$) spectral function $`J(E)=\gamma |E/\omega |^s`$ for which the relaxation rates are $`W_{qq^{}}=J(\omega (qq^{}))[\pm N(\omega (qq^{}))]`$ for $`qq^{}`$ where $`N(E)=(e^{\beta E}1)^1`$. We briefly discuss the results for intermediate relaxation when interesting deviations from a simple interpolation between non-equilibrium and strong relaxation limit occur.
## II Results
We now present results for the stationary current $`I`$ (Eq. (7)) and differential conductance $`dI/dV`$ for symmetric tunneling rates $`\mathrm{\Gamma }^{ir}=\mathrm{\Gamma },r=L,R,i=1,2`$. Throughout the paper we set the temperature to $`T=0.025\omega `$. Gray-scale plots of $`dI/dV(\mu ,V)`$ have been given different linear scale factors for $`dI/dV0`$ to clarify the voltage conditions under which NDC occurs. The NDC magnitude can be inferred from the presented $`I(V)`$ curves. We note that the results for a molecule with two neutral states and one charged state are simply obtained by inverting the sign of $`\mu \mu _0^1`$ and modifying the discussion accordingly.
### II.1 Strongly coupled single orbital \- Feedback mechanism and weak NDC
In Ref. Koch and von Oppen (2005) the one-orbital model in the limit of strong coupling to the vibration ($`\lambda _1^21`$) was shown to exhibit a current suppression which was related to “avalanches” of vibration assisted tunneling processes which also leads to super-poissonian current noise effects Koch and von Oppen (2005); Koch et al. (b). Additionally, for asymmetric gate voltages a weak NDC effect was noted in the absence of relaxation. Here we focus on this weak NDC effect and additional small current peaks (not discussed in Ref. Koch and von Oppen (2005)) which are visible in Fig. 3 as black lines and white-black double lines, respectively. The mechanism responsible for this is based on an energy asymmetry (induced by the gate potential) and the qualitative features of the FC-factors. This mechanism will also play a role in the transport involving two (or more) orbitals where it results in much stronger effects. We therefore consider this simple case in some detail.
Current suppression. The differential conductance plotted in Fig. 3 is symmetric about $`\mu =\mu _0^1`$ up to unimportant differences in amplitude due to spin degeneracy of the orbital (without which we would have exact symmetry) so we only discuss the case $`\mu \mu _0^1`$; the opposite case follows from interchanging the roles of the electronic states $`01`$. For low voltages $`V2\omega \lambda _1^2`$ states with low vibrational number are predominantly occupied and the current is exponentially suppressed both in the limit of weak and strong relaxation Koch and von Oppen (2005); Koch et al. (b). This is related to the classically forbidden region $`q,q^{}\lambda _1^2`$ in Fig. 2 where the FC factors depend exponentially on the vibrational numbers. In the limit of strong relaxation the current suppression is simply due to the exponentially small partial currents of the few vibrational states which are thermally occupied at low $`T\omega `$. In the absence of relaxation the FC-factors also prevent the excited states from actually becoming occupied for $`V2\omega \lambda _1^2`$. The resulting non-equilibrium vibrational distribution induced by the tunneling is “equilibrium-like” as was noted before Mitra et al. (2004); Koch and von Oppen (2005). This is due a type of feedback mechanism in the tunneling transitions, as we will now explain. In principle at finite bias voltage arbitrarily high vibrational excitations can be accessed via cascades of single-electron tunneling processes. However, at low bias voltage transitions which lie outside the bias window in Fig. 2 (i.e. $`0_q1_q^{}`$ resp. $`0_q1_q^{}`$ ) correspond to large changes of the vibrational energy and have exponentially larger rates. Once the initial states for these transitions start to be occupied, the total rate for populating the lowest states (transitions outside the bias window) becomes much larger than the total rates of its decay (transitions inside bias window). As a result only the low-lying states are occupied due to the large asymmetry between the rates. Compared with the strong relaxation limit vibrational excitations are slightly more favored and therefore the current suppression is less severe Koch and von Oppen (2005) in this limit (not shown). The central observation is that although the occupations decrease strongly with vibrational number this is compensated by the exponential increase of the partial currents. The main contribution to the current comes from the excited states. The “inverted” dependence of the FC-factors on energy (vibrational numbers) thus stabilizes the lowest vibrational state and enhances the sensitivity of the current to the small occupations of the vibrational excitations. This is at the basis of the weak NDC and small current peaks which we will discuss now.
Weak NDC and current peak. Apart from the quantitative effect on the current suppression, there appear interesting qualitative new features in the absence of relaxation. The many visible resonance lines in Fig. 3 form a Franck-Condon progression which extends beyond voltages $`2\omega \lambda _1^2`$. Interestingly, for asymmetric gate energy $`\mu \mu _0^1>\omega /2`$ the current can be reduced at the resonances $`\mu _L=\mu _k^1,k=1,2,\mathrm{}`$ (black features on right in Fig. 3). Here the transitions $`0_q1_{q+k},q=0,1,\mathrm{}`$ become allowed and one could naively expect the current to increase since this favors the population of the excited states $`1_{q+k},q>0`$ which are responsible for the main current contributions. Actually, the opposite happens. Due to the asymmetric gate-energy $`\mu \mu _0^1>0`$ a feedback mechanism is active which involves a cascade of single-electron transitions. This is illustrated in Fig. 5. In the regime of voltages where NDC occurs the charged vibrational ground state $`1_0`$ is stabilized, not only relative to the neutral states (due to the Coulomb blockade) but also relative to the vibrationally excited charged states $`1_q,q>0`$. Due to the asymmetric gate energy state $`1_0`$ has less transitions which depopulate it than which populate it, see Fig.5. Due to the strong increase of the FC-factors with vibrational number this asymmetry in the rates causes a nearly complete occupation of $`1_0`$ for asymmetric gate energies. Upon increasing the bias, at resonances where the excitations $`1_{q+k}`$ become accessible other excitations $`0_q^{}`$ with $`q^{}q+k`$ are also favored via subsequent tunneling processes involving small changes in the vibrational number. These subsequent transitions are allowed already at low $`V`$. The excitations $`0_q^{}`$ decay with large rates back to the the state $`1_0`$ and its first few excitations, since they change the vibrational number by a large amount. Importantly, the reverse of the latter transition, $`0_q^{}1_0`$ is not allowed at low $`T`$, see Fig. 5. Therefore the occupation of $`1_0`$ is effectively increased at the expense of the excited states which contribute most to the current and NDC occurs.
Small current peaks of width $`T`$ occur in the intermediate region $`0<\mu \mu _0^1\omega /2`$ (white-black double lines in Fig. 3 and inset of Fig. 4). These signal a redistribution of the vibrational energy when the bias is tuned through the resonance. In this case the above feedback mechanism can only become effective when the transition energy lies sufficiently close to $`\mu _L`$. Thus initially the current rises but once the excited states become sufficiently populated they start to relax via the feedback and the current drops again.
Finally, Fig. 4 shows that even though the absolute current-step amplitude exponentially increases with increasing bias voltage Koch and von Oppen (2005), the NDC becomes relatively less pronounced. Careful inspection reveals that at sufficiently large gate energy $`|\mu \mu _0^1|`$ the resonances initially correspond to current drops but with increasing voltage turn into peaks and finally become current steps. The increasing bias eventually compensates for the gate-asymmetry and the feedback mechanism becomes ineffective.
In the limit of strong relaxation at low $`T\omega `$ only transitions inside the bias window along the $`q=0`$ and $`q^{}=0`$ axis in Fig. 2 play a role (cf. Eq. 13). The charged vibrational ground state is stabilized due to a gate voltage $`\mu \mu _0^1>0`$ and the strong relaxation. However, no NDC or current peaks can occur at the resonances discussed above since the vibrational distribution is not affected by the tunneling in this limit. The feedback mechanism is cut off: after each single-electron tunneling process the excitation relaxes on a much shorter time scale and the next tunneling process starts from a vibrational ground-state again. Indeed one can show explicitly that in the limit of complete relaxation the NDC and the current peaks disappear in the one-orbital model for arbitrary spin and orbital degeneracy (see Appendix A). We note that at resonance lines $`\mu _R=\mu _k^1`$ for $`k=1,2,\mathrm{}`$ the current always increases, independent of the relaxation. Here the transitions $`1_q0_{q+k}`$ become allowed whereby $`1_0`$ can decay and repopulate the excited states which carry the current.
### II.2 Strongly coupled ground state
We now demonstrate that in a two-orbital model with asymmetric coupling to the vibration already at moderate coupling a new but weak NDC effect occurs which is robust against strong relaxation, in contrast to the weak NDC in the one-orbital model. Additionally, the non-equilibrium feedback mechanism responsible for the weak NDC in the one-orbital model leads to strong NDC effects for two competing orbitals. These general statements carry over to the case of multiple non-degenerate orbitals with asymmetric coupling occupied by at most one electron due to Coulomb blockade. We start our discussion with an intermediate case where only the weak NDC occurs.
#### II.2.1 $`\lambda _1^21\lambda _2^2`$ Weak NDC - Current oscillations
The differential conductance and typical I-V curve in Fig. 6 and 7, respectively, display a number of features which can be understood by considering the two orbitals invidually.
Vibrational excitations of the stronger coupled state 1 form a FC progression of resonances far beyond $`V2\lambda _1^2\omega `$. In contrast, the excitations of the weakly coupled orbital 2 only show up as a single $`dI/dV`$ peak at $`\mu _L=\mu _0^2`$. The resonances $`\mu _L=\mu _k^1`$ of the stronger coupled state 1, correspond to current steps for $`k[\mathrm{\Delta }/\omega ]`$. Interestingly, once the weakly coupled state 2 has started to contribute to the current, $`\mu _L>\mu _0^2`$, these turn into anti-resonances, $`k>[\mathrm{\Delta }/\omega ]`$, (dark lines in Fig. 6 with negative slope). The resonances, $`\mu _R=\mu _k^1,k0`$, always correspond to current steps (white lines in Fig. 6 with positive slope). A distinctive feature is that the current steps and drops due to the excitations of the stronger coupled state have opposite gate voltage dependence. Since they are of the same order of magnitude they give rise to current oscillations on the slowly saturating background. We note that for identical parameters the one-orbital model (i.e. with either $`\lambda _1=0.1`$ or $`1.4`$) produces negligible NDC effects. The origin of the enhancement of NDC is the following. For moderate gate energy $`\mu \mu _0^1`$ the partial currents of states $`1_q`$ are small compared with those of states $`2_q`$ and $`0_q`$, cf. Fig. 2. At resonance lines $`\mu _L=\mu _k^1`$ the transitions $`1_{q+k}0_q`$ become allowed, and the occupations of the states with the larger partial currents are reduced due to the Coulomb correlation between the two orbitals. The asymmetry between the partial currents is only present when the weakly coupled orbital 2 is accessible, therefore the current drops at these resonances only once both orbitals are accessible, i.e. $`\mu _L\mu _0^2\mu _R`$. In contrast, the resonances $`\mu _R=\mu _k^1,k>0`$ correspond to current steps since here states $`1_q`$ are depopulated in favor of the states with larger partial currents.
Although the current oscillations in Fig. 6 seem very similar to those for the one-orbital model in Fig. 3, there is an important difference: here the NDC is not completely suppressed in the strong relaxation limit, although larger values of the dominant coupling $`\lambda _1`$ are required for visibility comparable with the limit of no relaxation. In Appendix A we prove that in the strong relaxation limit a drop of the current can only occur along resonance lines $`\mu _L=\mu _k^i`$ (negative slope) if it occurs. A condition for the visibility of NDC is that the total current below the resonance is larger than compared with the partial current of the orbital causing it. This requires a weakly coupled orbital, $`\lambda _2^21`$, with large partial currents in combination with a stronger coupled orbital, $`\lambda _1^21`$. These are roughly the same conditions as for the visibility of the oscillating current in the non-equilibrium limit. However, we point out that under identical bias and gate voltage conditions the NDC need not to be visible in both the non-equilibrium and equilibrium limit, see for instance Fig. 7 and Fig. 10 below.
In summary, the current oscillation occurs due to the competition between transport channels with significantly different partial currents. It is a Coulomb repulsion effect: the current calculated without the effective correlations ($`v=0`$), always increases with $`V`$ (see Appendix B).
#### II.2.2 $`\lambda _1^21\lambda _2^2`$ Strong NDC
When the charged ground state couples strongly to the vibration , $`\lambda _1^21`$, and the electronic excitation lies low, $`\mathrm{\Delta }<\omega `$, we find the typical structure of Fig. 9. The current oscillations are clearly visible again and this set of resonances needs no further discussion. An obvious difference with Fig. 6 is the finite gap $`\mathrm{\Delta }`$ for any gate voltage. In fact, in addition to the Coulomb blockade regime where $`\mu _L<\mu _0^1`$ or $`\mu _0^1<\mu _R`$ resp., the current is suppressed in the entire strip $`\mu _0^2\mu _L\mu _0^1,\mu _R>\mu _0^1`$ (i.e. where the excited orbital 2 is not yet accessible). This is due to the exponentially small FC-factors $`F_{qq^{}}^1,q,q^{}\lambda _1^2`$ of the lowest orbital which now couples strongly to the vibration (cf. Sec. II.1). When increasing the bias voltage above the gap the resulting current depends strongly on the order in which additional transitions become allowed i.e. on the gate voltage. Basically four different situations can occur which are labeled (a)-(d) in Fig. 9 and the relevant transitions are schematically indicated in Fig. 9.
(a) Stabilization of the charged ground state. In the upper-right region in Fig. 9, the current may be expected to flow: the transitions $`2_q^{}0_q`$ are energetically allowed for at least $`q=q^{}=0`$, for which the FC-factor is large, $`F_{00}^21`$. Instead the current is strongly suppressed and increases in small steps of increasing height with increasing bias. This is very similar to the situation in Fig. 4 where the transport is dominated by a single orbital with $`\lambda _1^21`$. Between states $`1_q^{}`$ and $`0_q`$ the feedback mechanism discussed in Sec. II.1 is operative which keeps state $`1_0`$ almost fully occupied with increasing bias, see Fig. 9(a) and compare with Fig. 5. The presence of orbital 2 further enhances the feedback since the states $`0_q`$ (which feed back into $`1_0`$) can now also be populated via cascades of transitions involving orbital 2.
(b,c) Isolated region. The small diamond-shaped region in Fig. 9 at finite voltage $`\mathrm{\Delta }<V<2\omega `$ is remarkable. Inside it the current is non-zero and beyond any of its four defining boundaries in the plane of gate- and bias voltage the current is completely suppressed. A typical $`I(V)`$ curve through this region is shown in Fig. 10. The fact that the region is bounded by a NDC line with positive slope proves that it must be caused by non-equilibrium vibrational effects (Sec. I), see also the strong relaxation result in Fig. 10. This region can be reproduced by truncating the spectrum to 5 states: $`0_q,2_q,q=0,1`$ and $`1_0`$, see also Fig. 9. At the low bias side, Coulomb blockade and the small FC factor $`F_{00}^1`$ discussed above are responsible for the current suppression. Inside this region only the transitions $`2_00_01_0`$ are allowed. The stationary occupations follow from Eq. (4): $`P_0^0=1/5,P_0^1=P_0^2=2/5`$ and the current is:
$`I{\displaystyle \frac{1}{2}}I_0^2P_0^2+I_0^0P_0^0={\displaystyle \frac{2}{5}}\mathrm{\Gamma }`$ (15)
This is less than $`2\mathrm{\Gamma }/3`$ and $`4\mathrm{\Gamma }/5`$, the maximal current through one and two orbitals (without the vibration), respectively, due to the partial occupation of the strongly coupled state $`1_0`$ with suppressed partial currents. Note that state $`1_0`$ is not yet blocking the transport. At the high bias side, the current becomes suppressed when the first neutral excited state $`0_1`$ can be reached either via the cascade $`0_02_10_1`$ (NDC line with negative slope, case (b) in Figs. 9 and 9) or $`2_00_1`$ (NDC line with positive slope, case (c)). State $`0_1`$ can decay to $`1_0`$ when an electron enters the molecule through either junction i.e. the reverse transition $`0_11_0`$ is suppressed at low temperature $`T\omega \mathrm{\Delta }`$. Now state $`1_0`$ is almost fully occupied since it is populated much faster than it can decay: it is blocking the transport since the current it limited by the very small sum of its decay rates. The feedback loop thus involves both the weakly and strongly coupled state: the competition between the two orbitals which couple asymmetrically to the vibration causes the NDC to be much stronger compared with the one-orbital case (Sec.II.1). From the above it follows that the diamond-shaped region disappears for excitation energy $`\mathrm{\Delta }>\omega `$. However, the feature (a) and (d), which we discuss now, will still be present.
(d) Current peak. At the resonance line $`\mu _L=\mu _0^2`$ where the weakly coupled orbital 2 starts to participate a single large current step could be expected. Remarkably, NDC occurs in the middle of this resonance, producing a current peak (white-black double line with negative slope in Fig. 9) whose width is proportional to the electron temperature $`T`$. This sharp features stands out between the thermally broadened plateaus in Fig. 11.
The origin of the peak is a strong competition between the charged states $`1_0`$ and $`2_0`$ in the narrow energy window $`|\mu _L\mu _0^2|T`$, involving the feedback via neutral excited states $`0_q`$. This is most simply illustrated by considering gate energies $`\mathrm{\Delta }/2\omega <\mu \mu _0^1<\mathrm{\Delta }/2\omega /2`$ where a minimal set of 6 states is sufficient to understand the peak, see Fig. 9(d). At the rising side of the peak the rate of the transition $`2_00_0`$ (through the $`V`$ dependence in the Fermi-function, cf. Eq. 5) has increased sufficiently to enhance the current relative to the very small value supported only by the transitions to / from the blocking state, $`1_00_0`$. Due to the gate energy the simplest feedback loop involves a cascade of 6 transitions: $`0_02_00_12_10_21_0`$. This feedback initially remains ineffective since the excited state $`0_2`$ is not yet sufficiently populated. The vibrational distribution is equilibrium-like and the current follows the result for equilibrated vibrations, see Fig. 11. As one increases $`V`$ the occupations of the excited states in the feedback loop increases and a maximum is reached. Here the feedback dynamically starts to trap the molecule in the state $`1_0`$. The occupations of the excited states and the current now start to decrease and reach a lower value (although higher than before the peak). The current peak thus signals this redistribution of the vibrational energy in a small bias window. We note that the current is not completely suppressed: this only happens beyond the second strong NDC line $`\mu _L=\mu _1^2`$ (cf. Fig. 11) as may be understood by considering Fig. 9(d). The peak can thus be considered a precursor of the full onset of the feedback mechanism. For lower $`\mu \mu _0^1<\mathrm{\Delta }/2\omega `$ a similar argument involving more than 6 states explains why the peak becomes a step and why simultaneously the strong NDC along $`\mu _L=\mu _1^2`$ is further enhanced.
Intermediate relaxation. Upon increasing the vibrational relaxation rate $`\gamma `$ (cf. Sec. I) starting from zero, the strong relaxation result is approached as expected. The relaxation cuts off the cascade of transitions leading to the blocking state and reduces the importance of the feedback for the transport. Now the NDC becomes more pronounced at resonances $`\mu _L=\mu _k^1,k=1,2,\mathrm{}`$ where the transitions $`0_q1_{q+k}`$ become allowed. These enhance the occupation of $`1_0`$ due to relaxation and suppress the current. However, this approach is rather slow at low bias voltages: the NDC lines marking the isolated region remain clearly visible.
In summary, the strong NDC lines $`\mu _L=\mu _1^2`$ and $`\mu _R=\mu _1^2`$ are associated with excitations of the state with weakly coupled to the vibration (state 2). However, the strongly coupled state 1 is actually blocking the transport. The weakly coupled state allows an excess vibrational energy to accumulate on the molecule (through a cascade of tunneling processes) which is subsequently spent to trap the molecule in the strongly coupled state (in a single tunneling process). The blocking state can thus be reached under very general energetic conditions. Therefore NDC effects become strong when 2 (or more) orbitals which couple asymmetrically to the vibration compete in the transport. Finally we note that Fig. 9 is reminiscent of the signatures of spin-blockade of tunneling. There the resonance line marking the transition between ground states can be terminated at finite bias and the current is only recovered when excited states become accessible Weinmann et al. (1995). The resonance line thus shows a kink. Here the kink in the resonance line is more drastic since also transitions to many vibrationally excited states are suppressed. Such details are of importance to distinguish NDC due to spin excitations in molecules Romeike et al. from effects due to vibrational excitations.
### II.3 Stronger coupled excited state
We now consider the case opposite to Section II.2 where the excited orbital 2 takes up the role of the blocking orbital due to a strong(er) coupling to the vibration and the lower orbital 1 is weakly coupled. Vibration-assisted tunneling processes now stabilize the electronically excited state of the charged molecule i.e. an excess charge and energy may be stored on the molecule by the feedback mechanism. We focus on the resulting qualitative differences with respect to Sec. II.2 resulting from this population-inversion controlled by the bias voltage. As a simple reference point, we start again from an intermediate case.
#### II.3.1 $`\lambda _1^21\lambda _2^2`$ Weak NDC - current oscillations
Fig. 12 shows the typical differential conductance for $`\mathrm{\Delta }>\omega `$.
Similar to Fig. 6 the weakly coupled state (state 1) basically shows up as one large current step and the many resonance lines correspond to the excitations of the strongly coupled state (here state 2). However, at the resonances $`\mu _L=\mu _k^2,k[\mathrm{\Delta }/\omega ]`$ we have current drops for small $`|\mu \mu _0^2|`$ (i.e. away from the charge degeneracy point) and current steps for large $`|\mu \mu _0^2|`$. Therefore, around the charge degeneracy point $`\mu =\mu _0^1`$ the first few excitations beyond the electronic excited line $`\mu _L=\mu _0^2`$ show up as current steps and only at higher bias voltage turn into current drops, in contrast to the case where the lower orbital couples stronger to the vibration (Fig. 6). Up to now we have only found features in the differential conductance (either positive or negative) when incoming electrons excite the vibration with their excess energy. This is expected at low temperatures $`T\omega `$. Interestingly, in Fig. 12 a small current step at the resonance $`\mu _L=\mu _1^2`$ with negative slope is visible which corresponds to absorption of the vibrational energy by an incoming electron despite the low temperature $`T\omega `$ and moderate bias. The slightly enhanced current may be understood as an effect of significant heating of the molecule by the vibration assisted-tunneling.
#### II.3.2 $`\lambda _1^21\lambda _2^2`$ Strong NDC
When the coupling to the charged excited state becomes strong, $`\lambda _2^21`$ the case of most interest is that of a higher excited orbital, $`\omega <\mathrm{\Delta }\lambda _2^2\omega `$. The typical structure is shown in Fig. 13.
In contrast to Fig. 9 we do not have a gap here since the change in the nuclear configuration of the two ground states is now small: the current starts to flow at the edges of the Coulomb blockade region $`\mu _L=\mu _0^1`$ and $`\mu _0^1=\mu _R`$. Two NDC lines stand out in Fig. 13 (dark) where the current is significantly suppressed: $`\mu _L=\mu _1^1`$ and $`\mu _R=\mu _1^1`$. Again, the appearance of the latter NDC line with positive slope is a proof that non-equilibrium vibrations play a role (Sect. I). At these lines the transitions $`1_{q+1}0_q`$ and $`1_q0_{q+1}`$ between the neutral and the weakly coupled charged state become allowed which enhance the occupation of the excited states $`0_q`$. The latter feed back to the strongly coupled state $`2_0`$ with rates which increase exponentially with $`q`$ and the electronic excited orbital is predominantly occupied. We thus have a bias-controlled population inversion between ground- and excited state of the charged molecule due to their asymmetric coupling to the vibration. We discuss the four different situations labeled (a)-(d) in Fig. 13 for which the relevant transitions are schematically indicated in Fig. 14.
(a) Current peak. A current peak appears along the line $`\mu _R=\mu _0^1`$ for $`\mu \mu _0^1>\omega /2`$ (white-black double line on the right in Fig. 13). Compared with Fig. 9 this peak occurs on the opposite side and remains visible up to much higher values of $`\mu \mu _0^1`$. The mechanism causing the current peak is analogous to that discussed in Section II.2.2 with the roles of 1 and 2 interchanged. However, more vibrational excitations are involved in the feedback mechanism depicted in Fig. 14 since $`\mathrm{\Delta }>\omega `$. For moderate gate energy $`\mu \mu _0^1`$ the current peak is a precursor to the population inversion: the feedback is only completely activated beyond the resonance above it, $`\mu _R=\mu _1^1`$. There the transitions $`1_q0_{q+1}`$ become allowed and the strongly coupled state $`2_0`$ is predominantly occupied suppressing the current (upper-right region in Fig. 13). For large asymmetric gate energy $`\mu \mu _0^1\omega /2`$ the peak actually marks the population inversion: the NDC at the peak gains in amplitude at the expense of the NDC above it at $`\mu _1^1=\mu _R`$ (opposite to Fig. 9). In this case state $`2_0`$ can already be reached by many feedback cascades $`0_01_k2_0,1k[(\mu \mu _0^1)/\omega ]`$, once the escape from orbital 1 ($`0_01_0`$) becomes possible. Therefore the population inversion and current suppression are complete at the peak. For a low-lying excited orbital, $`\mathrm{\Delta }<\omega `$, this is in fact the general situation since the cascade of transitions involved in the feedback is shorter.
(b,c) Isolated region. For $`|\mu \mu _0^1|<\omega /2`$ and $`V<2\omega `$ we have an isolated region in the sense that the current reaches a local maximum value (diamond shaped region at bottom of Fig. 13). This region does not occur for $`\mathrm{\Delta }<\omega `$ (opposite to Sec. II.2.2 where $`\mathrm{\Delta }>\omega `$ suppresses the isolated region). Within this region only the transitions $`1_00_0`$ are allowed, $`P_0^0=1/3,P_0^1=2/3`$, and the current
$`II_0^0P_0^0+{\displaystyle \frac{1}{2}}I_0^1P_0^1={\displaystyle \frac{2}{3}}\mathrm{\Gamma }`$ (16)
equals the maximum current which a single orbital (without the vibration) can carry. When going along (b) and (c) in Fig. 13 we next cross the resonance lines $`\mu _L=\mu _1^1`$ and $`\mu _R=\mu _1^1`$ respectively and the current decreases as discussed above. The current suppression is complete once both transitions $`0_q1_{q+1}`$ and $`1_q0_{q+1}`$ are allowed.
(d) Absorption by incoming electrons. In Fig. 13 resonances $`\mu _L=\mu _k^2,k=1,2`$ (negative slope) are visible which correspond to absorption of the vibrational energy by an incoming electron ($`2_q0_{q+k}`$). What is remarkable compared to Fig. 12 is that the current decreases here. This is another signature of the population inversion due to the feedback mechanism. Vibrational energy can accumulate on the molecule through previous sequences of tunneling events involving only the weakly coupled orbital 1. The molecule is then “brought to a standstill” when, in a single tunneling process, an electron with an energy deficit matching the total accumulated vibrational energy enters: $`2_00_k`$. For $`\mu _k^2<\mu _L<\mu _0^2`$ there are $`[\mathrm{\Delta }/\omega ]`$ such resonance lines where such a new trapping process becomes possible. Note that this can not be understood as an effect of heating since at these resonances the occupation of vibrational excited states is suppressed due to the feedback and the current is reduced.
Intermediate relaxation. Compared with Section II.2.2 the non-equilibrium effects are more sensitive to relaxation since here roughly $`2\mathrm{\Delta }/\omega `$ tunneling events comprise the feedback mechanism instead of 2 (cf. Figs. 9 and 14). This restricts the minimal vibrational quality factor $`𝒬`$ (writing the vibrational relaxation rate as $`\gamma =\omega /𝒬`$, cf. Sec. I) for the observation of effects due the feedback mechanism to $`𝒬>2\mathrm{\Delta }/\omega `$ for $`\lambda _1^21\lambda _2^2`$. In contrast, for $`\lambda _1^21\lambda _2^2`$ the requirement is $`𝒬>2`$. For the cases discussed here only $`𝒬5`$ is required. This is confirmed by our calculations for intermediate values of the relaxation rate $`\gamma `$. Interestingly, when increasing the vibrational relaxation rate $`\gamma `$ starting from zero, the dependence of the amplitude NDC line with positive slope $`\mu _R=\mu _1^1`$ (proof of a non-equilibrium distribution) is non-monotonic: it is initially weakened and then regains amplitude and remains clearly visible with increasing $`\gamma `$. Also the current peak shifts to larger values of $`\mu \mu _0^1`$ and $`V`$ but remains visible. The low-bias effects of the excited state are however suppressed since cascades responsible for the population inversion effect are cut off by the relaxation.
In summary, it is remarkable that in all discussed cases (a)-(d) the excited state $`2_0`$ dominates the current at low bias where it can not be reached from the neutral ground state $`0_0`$ by a single-electron tunneling process. The strong deviation from equilibrium is induced by cascades of single electron tunneling processes.
### II.4 Non-degenerate strongly coupled states
Having analyzed the above cases in detail, we can now restrict ourselves to a brief classification of the results for non-degenerate states $`\mathrm{\Delta }>0`$ where both $`\lambda _1^2\lambda _2^2>1`$. Here the feedback mechanism discussed above produces more complex results. The basic change is that when the weakest coupling is increased the feedback mechanism becomes less efficient at populating vibrational excitations of the neutral molecule (compare with the feedback mechanism for a single orbital Sec. II.1). More vibrational excitations of the weakly coupled orbital and the neutral state must be accessible by a single tunneling process in order to fully activate the feedback and trap the molecule in strongest coupled orbital. The patterns of NDC lines will therefore extend over a broader range of applied voltages.
Indeed, a glance at Figs. 16-19 already shows that more NDC lines are visible. Also, there are more NDC lines with positive slope, which proves that the deviations from an equilibrium vibrational distribution are stronger. For $`\lambda _1^2\lambda _2^2>1`$ the NDC effects are strongest for the case $`\mathrm{\Delta }<\omega `$ presented in Fig. 16. Compared with Fig. 9 the isolated region defined by strong NDC lines is repeated a number of times to the left and it extends further to the right. Also, the current peak at the resonance line $`\mu _L=\mu _0^2`$ has shifted further to the left. The extended $`I(V)`$ plateaus have a width fixed by $`\omega \mathrm{\Delta }`$, independent of the gate voltage (compare with the weak NDC in Sec. II.2.1 where the current steps and drops have opposite gate-voltage dependence). For the case $`\mathrm{\Delta }>\omega `$ presented in Fig. 17 the low bias structure disappears and the current peak along $`\mu _L=\mu _0^2`$ becomes the dominant feature. For $`\lambda _2^2\lambda _1^2>1`$ the NDC effects are strongest for the case $`\mathrm{\Delta }>\omega `$ presented in Fig. 19. The two strong NDC lines in Fig. 13 have developed into a “checkerboard” pattern of such lines. These correspond to excitations of the weaker coupled state. In addition more resonances due to the strongly coupled state appear. For $`\mathrm{\Delta }<\omega `$ the current is more suppressed at positive gate energies, Fig. 18.
## III Summary and discussion
We have calculated the non-linear current through a molecule with two non-degenerate electron-accepting orbitals coupled asymmetrically to an internal vibration in the limit of weak tunneling to the electrodes. We found that due to the interplay of Coulomb blockade and non-equilibrium vibration-assisted tunneling NDC effects become amplified and pervasive in comparison with a one-orbital model. The only resonances where we consistently find current steps correspond to an electron tunneling off the molecule starting from the charged state coupled strongly to the vibration. At all other resonance lines the $`dI/dV`$ may become negative depending on the electron-vibration couplings and applied voltages. A weak and strong NDC effect may be distinguished, which require the larger of the two electron-vibration couplings to be moderate and strong respectively.
The weak NDC effect is found at resonance lines where an electron can tunnel onto the molecule resulting in the charged state coupled stronger to the vibration. This effect only occurs when two (or more) orbitals are competing in the transport. The current steps and drops occur at bias positions with an opposite gate-voltage dependence and give rise to current oscillations. We proved that this type of NDC is robust against strong relaxation of the vibrational distribution on the molecule due to a dissipative environment. Any NDC at other resonances conditions is a proof of a non-equilibrium vibrational distribution on the molecule. The current oscillations may even become amplified by strong relaxation depending on the applied voltages.
Strong NDC effects appear at the first few resonance lines associated with the state weakly coupled to the vibration. This is a non-equilibrium effect which is typically weakened by relaxation processes. Cascades of single-electron tunneling processes involving the vibrational excitations of the weakly coupled state provide a feedback which rapidly populates the strongly coupled state. The latter thus acts as a blocking state which is almost fully occupied. The few, strong NDC lines correspond to the activation of the feedback mechanism and can have the same gate voltage dependence as the current steps, in contrast to the weak NDC above. In a one-orbital model the feedback mechanism is also active but produces only a weak effect due to the absence of a competing orbital. An anomalous current peak of width $`T`$ appears when the feedback mechanism becomes effective only sufficiently deep inside a resonance. The peak signals the crossover of the vibrational distribution from equilibrium to non-equilibrium.
Interestingly, the blocking state can be the vibrational ground state of either charged state whichever couples stronger to the vibration. When the electronic excited orbital couples most strongly the NDC signals a voltage-controlled population inversion between the charged states induced by the vibration-assisted tunneling. Also, new resonances appear associated with an electron entering the molecule and absorbing vibrational energy stored on it (despite the low temperature) where the current is suppressed.
The NDC effects are due to asymmetry of the orbital energies and couplings to the internal vibration which are intrinsic properties of the molecule. One can thus tailor the electronic response of the device by molecular engineering. In contrast to other NDC effects, we do not require detailed assumptions of orbital- and/or electrode- specific electronic wave-function overlap with the electrodes Hettler et al. (2002b, a); Hettler et al. (2003); Thielmann et al. (2003) nor bias-voltage dependent coupling to the vibration McCarthy et al. (2003).
For the interpretation of transport spectroscopy experiments on molecular devices an important result of our work is that multiple orbitals may be relevant for effects at voltages where only a single orbital would seem to matter. For instance, an excited orbital may already dominate the transport by cascades of single tunneling events at a low voltage where this state is not directly accessible from the neutral ground state, cf. Fig. 13. Similarly the charged ground state may completely dominate the transport even at a high bias voltage where a far “better conducting” excited orbital is already directly accessible cf. Fig. 9.
We have used a basic parameterization of the nuclear potential surface of the electronic states and some comments are appropriate. For one, the nuclear potential shape may be anharmonic in the coordinate $`Q`$ considered here. The resulting qualitative changes may be determined from the FC-factors for these potentials, when plotted in similar fashion as Fig. 2. The main results are not sensitive to the fine details of these factors but only to their large-scale dependence on the vibrational numbers due to the shift of the potentials, which can be established by quasi-classical considerations. Secondly, anharmonic terms in the nuclear potential may also couple the mode $`Q`$ to other internal modes which we have not considered here. When this coupling is strong for a large number of such other modes or a Fermi- (nonlinear-)resonance is involved, intramolecular vibrational energy redistribution may relax the vibration. Generally, this will become more important for large molecules. The effect of relaxation has been discussed. If however, only a few other modes couple strongly to $`Q`$, say one, an interesting two mode problem occurs. A treatment of the effects of such multi-mode dynamics Koeppel et al. (1984) on the tunneling transport, lies outside the scope of the present paper.
###### Acknowledgements.
We acknowledge stimulating discussions with H. Schoeller, K. Flensberg, W. Belzig, I. Sandalov, and M. Hettler. M. R. W. acknowledges the financial support provided through the European Community’s Research Training Networks Program under contract HPRN-CT-2002-00302, Spintronics.
## Appendix A Strong relaxation
We consider the limit $`\mathrm{\Gamma }\gamma T`$ where the vibrational excitations have completely relaxed before each tunneling event due to a dissipative environment. The tunneling rates $`\mathrm{\Gamma }^{ir}`$ are not assumed to be symmetric. With the factorization ansatz $`P_q^i=P^iP_q`$, $`P_q=e^{q\omega /T}(1e^{\omega /T})`$ we can reduce equations (4) to an effective electronic three-level problem with voltage dependent rates (13). The stationary probabilities and the current can be explicitly given ($`r=R,L`$):
$`{\displaystyle \frac{P^i}{P^0}}=2{\displaystyle \frac{W_{i0}}{W_{0i}}},P^0=\left[1+2{\displaystyle \underset{i}{}}{\displaystyle \frac{W_{i0}}{W_{0i}}}\right]^1,`$ (17)
$`I_r={\displaystyle \underset{i}{}}\left(2W_{i0}^rP^0W_{0i}^rP^i\right).`$ (18)
We can now find a simple explicit condition for the occurrence of current steps or drops with increasing bias voltage. Consider an increase of the positive bias $`VV^{}`$ such that one additional transition involving orbital $`i`$ comes into the bias window through a resonance with electrode $`r=L,R`$, $`\mu _r=\mu _k^i`$ for some $`k=0,\pm 1,\pm 2,\mathrm{}`$ (for $`V<0`$ interchange $`LR`$ below). For simplicity we consider the values of the current on the two subsequent plateaus: only two transition rates are then changed, $`W_{i0}W_{i0}^{}`$ and $`W_{0i}W_{0i}^{}`$ the changes being related by $`\delta W_{0i}^r=e^{k\beta \omega }\delta W_{i0}^r`$ (cf. Eqs. 5, 13). The change in the stationary current $`\delta I_r=I_r^{}I_r`$ may be calculated at either electrode $`r=L,R`$ from (14) (since $`I_L+I_R=I_L^{}+I_R^{}=0`$):
$`{\displaystyle \frac{\delta I_r}{I_r}}=2\left({\displaystyle \frac{W_{i0}^{}}{W_{0i}^{}}}{\displaystyle \frac{W_{i0}}{W_{0i}}}\right)P^0\left({\displaystyle \frac{W_{0i}^{\overline{r}}}{I^r}}1\right),`$ (19)
Here $`\overline{r}=R,L`$ denotes the electrode opposite to $`r=L,R`$. In order to have NDC at a resonance where $`\mu _L`$ becomes larger than $`\mu _k^i`$ we require $`\delta I_L/I_L<0`$, which, using $`r=L`$ in Eq. (19), gives $`W_{0i}^RI_L<0`$. The rate of escape through junction $`R`$ at voltage $`V^{}`$ (including the new transition in increased bias window $`V^{}`$) must thus be smaller than the current $`I_L`$ at initial voltage $`V`$. It is readily seen that this condition cannot be fulfilled in the case of only one orbital: $`W_{01}^RI_L=W_{01}^R+I_R=2W_{10}^RP^0+W_{01}^R(1P^1)>0`$. However, for two (or more) orbitals it is possible to satisfy this requirement. To see if NDC may occur at a resonance where $`\mu _R`$ drops below $`\mu _k^i`$ we use $`r=R`$ in Eq (19), leading to the requirement $`W_{0i}^LI_R=W_{0i}^L+I_L<0`$. which can not be satisfied for any applied voltages since $`I_L>0`$ for $`V>0`$. The current must increase at such resonances. Thus for fully equilibrated vibrations NDC can only occur at resonances related to the left electrode, $`\mu _L=\mu \mu _k^i,k=0,\mathrm{}`$, where electrons can enter the molecule by an additional tunneling process $`i_{q+k}0_q`$. These resonances correspond to lines with positive slope in the $`(\mu ,V)`$ plane. Any NDC occurring at an other resonance is a proof of a non-equilibrium vibrational distribution. This proof can be trivially extended to $`N`$ orbitals correlated by Coulomb charging (maximally one extra electron). It also does not depend on the FC-factors involved, although the amplitude of the possible NDC may be small for a particular choice.
## Appendix B Uncorrelated vibration-assisted tunneling
We consider the special case where the renormalization of the interaction (due to the polaron effect) compensates the Coulomb repulsion effects, i.e. $`v=v^{(0)}2\omega \lambda _1\lambda _2=0`$. Now we have to include the doubly charged state of the molecule $`n_1=1,n_2=1`$ (di-anion). We denote the diagonal density matrix elements by $`P_q^{n_1n_2}`$, where $`q`$ is the vibrational number and $`n_1n_2`$ denotes an electronic state with occupations $`n_i=0,1`$ of orbital $`i=1,2`$. The occupations are coupled by the stationary master equations (cf. Wegewijs and Nazarov (1999))
$`{\displaystyle \underset{i}{}}\left[n_i{\displaystyle \underset{q^{}}{}}W_{0q^{}iq}+(1n_i){\displaystyle \underset{q^{}}{}}2W_{iq^{}0q}\right]P_q^{n_1n_2}`$
$`+{\displaystyle \underset{q^{}}{}}[2n_1W_{1q0q^{}}P_q^{}^{0n_2}+2n_2W_{2q0q^{}}P_q^{}^{n_10}`$
$`+(1n_1)W_{0q1q^{}}P_q^{}^{1n_2}+(1n_2)W_{0q2q^{}}P_q^{}^{n_11}]=0`$
together with $`_{n_1n_2}P^{n_1n_2}=1`$. In the case where the vibration is assumed to be completely equilibrated, $`\overline{P}_q^{n_1n_2}=\overline{P}^{n_1n_2}P_q`$ the kinetic equations can be decoupled into equations for the occupations of two uncorrelated “channels”, $`\overline{P}^i_{n_1n_2}\delta _{1n_i}P^{n_1n_2}`$ with the rates (13): $`\overline{P}^i=0=2W_{i0}(1\overline{P}^i)W_{0i}\overline{P}^i`$. The current is then a sum of independent contributions of the individual orbitals: $`I_r=_iI_{ri}`$ where $`I_{ri}=(W_{0i}^{r}{}_{}{}^{1}+2W_{i0}^{r}{}_{}{}^{1})^1`$. For a single orbital the current monotonically increases in the limit of strong relaxation (Appendix A) and the same thus holds for two (or more) uncorrelated orbitals.
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# Sparse equidistribution problems, period bounds, and subconvexity
## 1. Introduction.
### 1.1. General introduction.
Let $`\mathrm{\Gamma }G`$ be a lattice in an $`S`$-arithmetic group. Let $`Y\mathrm{\Gamma }\backslash G`$ be a subset endowed with a probability measure $`\nu `$, and $`f`$ a function on $`\mathrm{\Gamma }\backslash G`$. Fixing a basis $`\{\psi _j^{(Y)}\}`$ for $`L^2(Y,\nu )`$, we shall refer to the numbers $`f\psi _j^{(Y)}𝑑\nu `$, as the periods of $`f`$ along $`Y`$. Evidently, the periods depend heavily on the choice of basis for $`L^2(Y,\nu )`$. They play a major role in the theory of automorphic forms, in significant part because they often express information about $`L`$-functions.
The present paper is centered around a geometric method yielding upper bounds for these periods. It is applicable, roughly speaking, when considering the periods of a fixed function $`f`$ along a sequence of subsets $`(Y_i,\nu _i)`$, with the property that the $`Y_i`$ are becoming equidistributed; that is to say, the $`\nu _i`$ approach weakly the $`G`$-invariant measure on $`\mathrm{\Gamma }\backslash G`$. The key inputs of this method are, firstly, the equidistribution of the $`\nu _i`$, and secondly, the mixing properties of certain auxiliary flows. More precisely, we shall need these properties in a quantitative form; in the cases we consider, this will follow eventually from an appropriate spectral gap.
This situation might seem rather restrictive. However, it arises often in many natural equidistribution questions (“sparse equidistribution problems,” as we discuss below) as well as in the analytic theory of automorphic forms (especially, subconvexity results for $`L`$-functions). There are applications besides those discussed in the present paper; our aim has not been to give an exhaustive discussion, but rather just to present a representative sample of interesting cases. We shall explain the method abstractly in Sec. 1.3 and will carry out, in the body of the paper, one example of each of the following cases: $`Y_i`$ is the orbit of a unipotent, a semisimple, and a toral subgroup of $`G`$.
In the present paper, we have focused mostly on the case of $`\mathrm{PGL}_2`$ and $`\mathrm{GL}_2`$ over number fields. All our results pertain to this setting, except for Thm. 3.2, which applies to a general semisimple group. The geometric methods of this paper are general and we hope to analyze further higher rank examples in a future paper.
Throughout the present methods we have tried to use “soft” techniques as a substitute for explicit spectral expansions. However, there still seem to be instances where the explicit spectral expansions are important. In a future paper , joint with P. Michel, we shall combine ideas drawn from this paper with ideas from Michel’s paper ; in that paper, we shall make much more explicit use of spectral decomposition.
We shall use the term “sparse equidistribution problems” to describe questions of the following flavor: Suppose $`Z_iY_i`$ is a subset endowed with a measure $`\nu _i^Z`$, and we would like to prove that the $`\nu _i^Z`$ are becoming equidistributed. In other words, we wish to deduce the equidistribution of the “sparse” subset $`Z_i`$ from the known equidistribution of $`Y_i`$. Examples of this type of question are Shah’s conjecture (where the $`Z_i`$ are discrete subsets of $`Y_i`$, a full horocycle orbit) as well as Michel’s results on subsets of Heegner points (where the $`Z_i`$ are subsets of the $`Y_i`$, the set of all Heegner points.) The connection to period integrals is as follows: one can spectrally expand the measure $`\nu _i^Z`$ in terms of the basis for $`L^2(Y_i,\nu _i)`$. Using our results for periods along $`Y_i`$, it will sometimes be possible to deduce the equidistribution of $`\nu _i^Z`$.
We now briefly summarize our results.
1. Sec. 3 considers where the $`Y_i`$s are orbits, or pieces of orbits, of unipotent groups. The mixing flow is the horocycle flow along $`Y_i`$.
In Thm. 3.1 (p. 3.1) we show that certain sparse subsets of horocycles on compact quotients of $`\mathrm{SL}_2()`$ become equidistributed. This is progress towards a conjecture of N. Shah. In Thm. 3.2 (p. 3.2) we give a fairly general bound (in the context of an arbitrary semisimple group) on the Fourier coefficients of automorphic forms. In the case of $`G=\mathrm{SL}_2()`$ it recovers results of Good and Sarnak , which resolved a problem of Selberg. The present proof is more direct, avoiding in particular the triple product bounds for eigenfunctions.
2. Sec. 4 considers the case when $`G=\mathrm{PGL}_2(F)`$, where $`F`$ is a number field, and $`\mathrm{\Gamma }`$ is a congruence subgroup thereof. The $`Y_i`$ are a sequence of closed diagonal $`G`$-orbits on $`\mathrm{\Gamma }\backslash G\times \mathrm{\Gamma }\backslash G`$. The mixing flow (after lifting to the adeles) is the diagonal action of $`\mathrm{PGL}_2(𝔸_{F,f})`$, where $`𝔸_{F,f}`$ is the ring of finite adeles of $`F`$.
Prop. 4.1 and Prop. 4.2 give period bounds in this context. We refer to Prop. 4.1 as a subconvex bound for the triple product period, for the reason that it should be in fact be equivalent to subconvexity for the triple product $`L`$-function, in the level aspect as one factor varies, but the necessary computation of $`p`$-adic integrals (Hypothesis 11.1) has apparently not yet been done in sufficient generality. In Thm. 5.1 (p. 5.1) it is shown that these results yield subconvex bounds, in the level aspect, for standard and Rankin-Selberg $`L`$-functions attached to $`\mathrm{PGL}_2`$.
The results on standard and Rankin-Selberg $`L`$-functions generalize results of Duke-Friedlander-Iwaniec and Kowalski-Michel-Vanderkam from the case $`F=`$.<sup>1</sup><sup>1</sup>1We have not attempted to address the issue of varying the central character. This, in a sense, is the most subtle point, as is shown by Michel’s recent work on Rankin-Selberg convolutions. Our aim in the present paper has been to show that one can derive a coherent theory for $`\mathrm{PGL}_2`$ from the triple product bound of Prop. 4.1. The case of varying central character will be discussed in a future paper with Michel. The third result, concerning subconvexity of the triple product period in the level aspect, was not known even over $``$; however, Bernstein and Reznikov have shown subconvexity for the triple product period in the eigenvalue aspect.
3. Sec. 6 considers the case when $`Y_i`$ is a certain family of noncompact torus orbits on $`\mathrm{\Gamma }\backslash G`$, where $`(\mathrm{\Gamma },G)`$ is as in Sec. 4. (In fact, the $`Y_i`$ are obtained by taking a fixed noncompact torus orbit, and translating by a $`p`$-adic unipotent, where $`p`$ varies.) The mixing flow is the action of the adelic points of the torus.
We establish in Thm. 6.1 (p. 6.1) subconvexity for character twists of $`\mathrm{GL}(2)`$ in the level aspect. This was established for $`F=`$ by Duke-Friedlander-Iwaniec, and the special case where $`F`$ is totally real and the form holomorphic at all infinite places was treated by Cogdell, Piatetski-Shapiro and Sarnak. In particular, (6.2) gives a subconvex bound for Grössencharacter $`L`$-functions over $`F`$, in the level aspect; this was known over $``$ by work of Burgess, and some special cases were known in the general case.
4. In Sec. 7 we consider the case where $`Y_i`$ is a (union of) compact torus orbits on $`\mathrm{\Gamma }\backslash G`$, where $`(\mathrm{\Gamma },G)`$ are as in Sec. 4. The equidistribution of such $`Y_i`$ will amount to the equidistribution of Heegner points, and we deduce it from Thm. 6.1 in Thm. 7.1 (p. 7.1). This result generalizes work of Duke over $``$ and was proven, conditionally on GRH, by Zhang, Cohen, and Ullmo-Clozel (independently). The present work makes this result unconditional.
Applying mixing properties of the adelic torus flow, we obtain in Thm. 7.2 (p. 7.2) we obtain, under a condition of splitting of enough small primes, the equidistribution of certain sparse subsets of Heegner points. In the case $`F=`$, an unconditional result of this nature is due to Michel.<sup>2</sup><sup>2</sup>2 Our method is different to Michel’s: we do not deduce our result from results on Rankin-Selberg convolutions, and indeed it is possible to deduce a subconvexity result from ours. However, there seem to be some curious parallels between the methods. In fact, the method of Thm. 7.2 is even more closely related – as Michel has pointed out to me – to the work of Duke, Friedlander and Iwaniec. In that paper they amplify class group $`L`$-functions but obtain only a conditional result for precisely the same reason that Thm. 7.2 fails to be unconditional, namely, one cannot guarantee unconditionally the existence of enough small split primes.
In that context of $`L`$-functions, one pleasing feature of the present method is that it is geometric: it proceeds not via Fourier coefficients but via the integral representation. In practice, this means that there is no difference between Maass or holomorphic forms, nor between $``$ and an arbitrary base field. Moreover, we do not make use of either the trace formula or the Kuznetsov formula; indeed, we make no explicit use of families.
The recent work of Bernstein-Reznikov is of a similar flavor. They establish a “subconvex” bound for the triple product when the eigenvalue of one factor varies, whereas we have treated the case where the level of one factor varies. Their method is also geometric in nature, and moreover their result applies to a nonarithmetic group. By contrast, the level aspect question is not well-posed if one leaves the arithmetic setting.
Throughout the paper we have not attempted to optimize the results. The input to our method is an equidistribution result. As far as possible we have tried to establish these results by relatively “geometric” methods, deriving in the end from the mixing properties of a certain flow. Of course, it is in many contexts better to use spectral methods, but this would involve departing from the geometric method that is intended to be the central theme of this paper. As remarked, we will pursue such “spectral” approaches in a forthcoming paper with P. Michel ; some of the results of this have been discussed in .
Finally, implicit in various parts of the paper is “adelic analysis”, i.e. the analytic theory of functions on adelic quotients, in the quantitative sense needed for analytic number theory. There seems to be considerable scope to develop this theory fully.
### 1.2. Other applications.
The method of this paper has other applications not elaborated here. We discuss some of them here.
There are other subconvexity results that are naturally approached by the same method: for instance, a subconvex estimate for $`L(\pi ,\frac{1}{2}+it)`$ where $`t`$ varies and $`\pi `$ is a fixed cuspidal representation of $`\mathrm{GL}(2)`$ over a number field $`F`$. In such a context it is natural to use the fact that the horocycle flow is (quantifiably) weakly $`k`$-mixing, for certain $`k>1`$; the use of this higher order mixing is closely related to Weyl’s “successive squaring” approach to $`\zeta (1/2+it)`$. Of course, this particular instance of subconvexity is approachable by standard methods also; an intriguing question in the subconvexity context is how to combine the present methods with those such as Bernstein-Reznikov.
There are certain applications to effective equidistribution theorems: for instance, it is also possible to establish some new effective cases of Ratner’s theorem by the same ideas, see Rem. 3.1. The question of giving such “nontrivial” cases was raised by Margulis in his talk at the American Institute of Mathematics, June 2004. Unfortunately, these new cases are rather artificial.
One can give certain analytic applications: let $`\mathrm{\Gamma }`$ be a cocompact subgroup of $`\mathrm{SL}(2,)`$, and let $`\pi L^2(\mathrm{\Gamma }\backslash \mathrm{SL}(2,))`$ be an irreducible $`\mathrm{SL}(2,)`$-subrepresentation. For $`m`$, let $`e_m`$ be the $`m`$th weight vector in $`\pi `$, if defined; i.e. a vector which transforms under the character $`\left(\begin{array}{cc}\mathrm{cos}(\theta )& \mathrm{sin}(\theta )\\ \mathrm{sin}(\theta )& \mathrm{cos}(\theta )\end{array}\right)e^{2\pi im\theta }`$. We normalize it (up to a complex scalar of absolute value $`1`$) by requiring that $`e_m_{L^2}=1`$. Bernstein and Reznikov proved the bound $`e_m_L^{\mathrm{}}(1+|m|)^{1/2}`$, and asked \[2, Remark 2.5(4)\] if any improvement of the exponent $`1/2`$ is possible. It is quite easy to deduce from Lem. 3.1 such a bound; indeed, the analytic properties of the $`e_m`$, as $`|m|\mathrm{}`$, is connected with the long time behavior of the horocycle flow in the same fashion that the analytic behavior of Laplacian eigenfunctions are connected to the long time behavior of the geodesic flow. In the time during which this paper was being revised for submission, Reznikov has proven independently a result of this type . Since the result he obtains is most likely sharper than that obtained by the technique indicated above, we will not pursue this further, noting only that an advantage of the method we have indicated above is that it is likely to generalize to higher rank.
Moving slightly away from the main subject of the present paper, the idea of using equidistribution theorems to produce mean value results for $`L`$-functions seems capable of application in a variety of settings. In particular, equidistribution results are readily available on $`\mathrm{GL}(n)`$, owing to Ratner’s work, whereas trace formulae are extremely unwieldy for $`n>2`$. It would be interesting to see what mean-value statements can be deduced from Ratner-type equidistribution results.
Historically, one application of such results has been to nonvanishing results; here the most spectacular results (e.g. ) have been achieved through the so-called mollifier technique. It would be quite interesting to understand if there is a geometric interpretation of the mollifier technique.
### 1.3. Discussion of method: equidistribution, mixing, and periods.
We now turn to a discussion of the specifics of the method used in this paper. This method itself is quite easy to describe. It consists in essence of two simple steps (see (1.2) and (1.3) below).
We also remark that the discussion that follows is a relatively faithful rendition of the method of the paper. The body of this paper does not really utilize any new ideas beyond the ones indicated below. Most of the bulk consists of the technical details necessary to connect periods with other objects of interest (e.g. equidistribution questions or $`L`$-functions), as well as setting up the machinery to quantify some standard equidistribution results. As much as possible, we have tried to give a self-contained treatment of all these technical details in Sections 811.
We hope the ensuing discussion serves as a unifying thread for the rest of the paper. We explain the method first in an abstract setting (Sec. 1.3.1). We then explain (Sec. 1.3.2 and 1.3.3) these ideas in a a more down-to-earth fashion, emphasizing the parallel with the analytic techniques for studying $`L`$-functions. Finally, Sec. 1.3.4 illustrates these ideas in a simple example – that of Fourier coefficients of modular forms.
#### 1.3.1. Abstract setting.
Let $`G_2G_1`$ be locally compact groups, $`\mathrm{\Gamma }G_1`$ a lattice, $`X=\mathrm{\Gamma }\backslash G_1`$. Let $`x_iX`$ and put $`Y_i=x_iG_2`$. We shall suppose that there exists a $`G_2`$-invariant probability measure $`\nu _i`$ on $`Y_i`$. (This does not precisely cover all the contexts we consider – at some points we will consider $`Y_i`$ which are “long pieces” of a $`G_2`$-orbit rather than a single $`G_2`$-orbit, but the ideas in that case will be identical to those discussed here).
Let $`f`$ be a function on $`X`$ and $`\psi _i`$ a function on $`Y_i`$ such that $`_{Y_i}|\psi _i|^2𝑑\nu _i=1`$. We will give a bound for the period $`_{Y_i}f\psi _i𝑑\nu _i`$.
In words, the idea will be to find certain correlations between the values of $`\psi `$ at different points; and then show that the values of $`f`$ at these same points are “uncorrelated,” in some quantifiable sense. Putting these together will show that the period must be small. The “hard” ingredient here is some version of the spectral gap, i.e. quantitative mixing, which is what will show that the “uncorrelated-ness” property of $`f`$.
We will suppose that there exists $`\sigma `$, a measure on $`G_2`$, such that
(1.1)
$$\psi _i\sigma =\lambda _i\psi _i,$$
for some $`\lambda _i`$. Here $`\sigma `$ denotes the action of $`\sigma `$ by right convolution. Let $`\stackrel{ˇ}{\sigma }`$ be the image of $`\sigma `$ by the involution $`gg^1`$ of $`G_2`$.
Then
(1.2) $`\left|{\displaystyle f\psi _i𝑑\nu _i}\right|^2=\left|\lambda _i^1{\displaystyle _{Y_i}}f(\psi _i\sigma )𝑑\nu _i\right|^2=\left|\lambda _i^1{\displaystyle _{Y_i}}(f\stackrel{ˇ}{\sigma })\psi _i𝑑\nu _i\right|^2`$
$`|\lambda _i|^2{\displaystyle _{Y_i}}|f\stackrel{ˇ}{\sigma }|^2𝑑\nu _i,`$
where we have applied Cauchy-Schwarz at the final step. Now, we are assuming that the $`Y_i`$ are becoming equidistributed, and so $`\nu _i\nu `$, the $`G_1`$ invariant measure on $`\mathrm{\Gamma }\backslash G_1`$. Thus
(1.3) $`{\displaystyle _{Y_i}}|f\stackrel{ˇ}{\sigma }|^2𝑑\nu _i{\displaystyle _X}|f\stackrel{ˇ}{\sigma }|^2𝑑\nu `$
$`={\displaystyle _{g,g^{}G_2}}gg^1f,f_{L^2(X)}𝑑\sigma (g)𝑑\sigma (g^{}),`$
where $`gg^1f`$ denotes the right translate of $`f`$ by $`gg^1`$.
If the $`G_2`$-action on $`X`$ is mixing in a quantifiable way – i.e., one has strong bounds on the decay of matrix coefficients – one obtains good upper bounds on the right-hand side of (1.3); in combination with (1.2) this gives an upper bound for the period $`|_{Y_i}f\psi _i𝑑\nu _i|`$.
The strength of the information required about the mixing varies. In the cases we study where $`G_2`$ is amenable, any nontrivial information will suffice. In the one case where $`G_2`$ is semisimple, a strong bound towards Ramanujan is needed. For instance, in the case of triple products, we need any improvement of the bound that the $`p`$th Hecke eigenvalue of a cusp form on $`\mathrm{GL}(2)`$ is bounded in absolute value by $`p^{1/4}+p^{1/4}`$. (In this normalization, the trivial bound is $`p^{1/2}+p^{1/2}`$).
In the rest of this paper, we shall merely apply this argument many times, with various different choices for $`\mathrm{\Gamma },G_1,G_2`$. The part of the argument which will vary is quantifying the equidistribution of the $`\nu _i`$, i.e. keeping track of the error in the first approximation of (1.3). Thus we make heavy use of Sobolev norms (Sec. 8), which are an efficient method of bounding this error.
In each instance, the proof of the equidistribution result $`\nu _i\nu `$ will always be rather straightforward, except for the result of Sec. 7. The equidistribution result needed for the proof of Thm. 7.2 is essentially equivalent to the subconvexity result proved in Sec. 6. A rather striking point is that a similar logical dependence (although manifested very differently) is present in the work of Michel. The meaning of this is unclear to the author.
In certain specific cases, the above technique is quite familiar. When $`G_2`$ is a one-parameter real group, the above argument is quite closely related to standard techniques of analytic number theory. <sup>3</sup><sup>3</sup>3 For example, in certain contexts when $`G_2`$ is abelian, one can push this method further by squaring multiple times, that is to say, considering $`|f\psi _i𝑑\nu _i|^4,|f\psi _i𝑑\nu _i|^8`$ and so forth. In this context, one replaces the mixing property of the $`G_2`$ action with results about higher order mixing of the $`G_2`$-flow. Although we will not carry this out in the present paper, this seems rather closely connected to Weyl’s proof of subconvexity for $`\zeta (1/2+it)`$. On the other hand, when $`G_2`$ is an adelic group, and $`\sigma `$ a measure on $`G_2`$ that corresponds to the action of Hecke operators (this is carried out in Sec. 4, for instance), the above argument will be essentially “amplification” in the sense of Friedlander-Iwaniec .
In the following two sections, we shall attempt to explain more colloquially the main idea that is at work here, and also discuss how the method described above fits into the framework of analytic number theory. Modern proofs of subconvexity, following the path-breaking work of Friedlander-Iwaniec , have roughly speaking consisted of a mean-value theorem and an amplification step. We shall discuss how the proof indicated above may be viewed as geometrizing this strategy, where the mean-value step is replaced by an equidistribution theorem, and the amplification step is controlled using mixing.
Note, in particular, that in the work of Friedlander-Iwaniec, families of $`L`$-functions play a central role, whereas the method above has in a certain sense eliminated the family. Although in the discussion below we rephrase matters so as to make clear the connection with the work of Friedlander-Iwaniec, it seems that from the perspective of the present paper the phrasing in terms of families is rather artificial.
#### 1.3.2. Connection with analytic number theory: Equidistribution, and mean-value theorem for periods.
Follow the notations of the previous section. We choose an orthonormal basis $`\{\psi _{i,j}\}_{j=1}^{\mathrm{}}`$ for $`L^2(Y_i,\nu _i)`$ so that $`\psi _{i,1}:=\psi _i`$.
By Plancherel’s formula, $`_{j=1}^{\mathrm{}}\left|f\psi _{i,j}𝑑\nu _i\right|^2=|f|^2𝑑\nu _i`$. Since $`\nu _i\nu `$ weakly, and we are holding $`f`$ fixed, it follows that:
(1.4)
$$\underset{j=1}{\overset{\mathrm{}}{}}\left|f\psi _{i,j}𝑑\nu _i\right|^2_{\mathrm{\Gamma }\backslash G}|f|^2𝑑\nu ,$$
as $`i\mathrm{}`$. Thus the equidistribution property of $`\nu _i`$ underlies a mean-value theorem for the $`Y_i`$-periods.
In many cases involving automorphic forms, the periods will essentially be special values of $`L`$-functions and (1.4) amounts to a mean-value theorem for $`L`$-functions. This is fairly well-known; for example, the mean-value theorem $`_T^T|\zeta (1/2+it)|^4dtT\mathrm{log}(T)^4`$ is rather closely connected with the equidistribution properties of the cycle $`\{(1+i/T)x,x\}`$, when projected to $`\mathrm{SL}_2()\backslash `$. A more striking example is Vatsal’s use of equidistribution to prove nonvanishing results . In general, it seems that there are many interesting mean value theorems for $`L`$-functions that are connected to equidistribution results.
In any case, (1.4) is not unrelated to the standard methods of obtaining such results; however, its primary advantage is that it is often technically much simpler, for example when working over a number field.
#### 1.3.3. Connection with analytic number theory (II): Mixing, and bounds for a single period.
We now wish to pass from (1.4) to nontrivial upper bounds for a single period. It is clear that (1.4) implies at once – by omitting all terms but one – that $`|f\psi _{i,j}𝑑\nu _i|\stackrel{<}{}f_{L^2(X)}`$; we shall refer to an improvement of this bound as nontrivial. It is evident that one must have some further information about $`\{\psi _{i,j}\}`$ in order to do this; otherwise one could simply take $`\psi _{i,1}`$ to be a multiple of $`f|_{Y_i}`$.
In the context of analytic number theory, this is often carried out by “shortening the family,” that is to say: proving a sharp mean-value theorem of the form of (1.4), but over some subfamily of $`\{\psi _{i,j}\}_{j=1}^{\mathrm{}}`$; then omitting all terms but $`\psi _{i,1}=\psi _i`$ will often give a nontrivial upper bound. In the work of Friedlander-Iwaniec, a weighted mean-value theorem is derived, which has the same effect as shortening the family.
Such a weighted mean-value theorem is also implicit in our context. Following the notation of Sec. 1.3.1, suppose that there is a fixed measure $`\sigma `$ on $`G_2`$ such that for all $`i,j`$, we have $`\psi _{i,j}\sigma =\lambda _{i,j}\psi _{i,j}`$ (some $`\lambda _{i,j}`$). Then, by Plancherel’s formula, and using the fact $`\nu _i\nu `$, we conclude:
(1.5)
$$\underset{j=1}{\overset{\mathrm{}}{}}|\lambda _{i,j}|^2\left|f\psi _{i,j}𝑑\nu _i\right|^2f\stackrel{ˇ}{\sigma },f\stackrel{ˇ}{\sigma }_{L^2(X)}.$$
This gives a weighted mean value theorem, which for appropriate choices of $`\sigma `$ amounts to shortening the effective range of summation in (1.4). Moreover, the mixing of the $`G_2`$-flow bounds the right hand side of (1.5). In this phrasing, it becomes clear that the measure $`\sigma `$ has played the role of an “amplifier” and the orthonormal basis for $`L^2(Y_i,\nu _i)`$ has played the role of the family.
Having now explained the method in an abstract context and indicated its equivalence with other methods, we now indicate more informally the source of cancellation in periods that is at the center of our results.
In many natural situations, one obtains a basis for $`L^2(Y_i,\nu _i)`$ by diagonalizing a geometrically defined algebra of operators on $`Y_i`$. The result of this process is that the functions $`\{\psi _j\}`$ exhibit correlations between their values at different points of $`Y`$. For instance (for example when $`G_2`$ is semisimple), it often will occur that there is a correspondence $`𝒞:YY`$ such the value of each $`\psi _j`$ at $`PY`$ and at the collection of points $`𝒞(P)`$ are correlated in some way. On the other hand (and we shall now speak quite imprecisely) if the correspondence $`𝒞`$ “extends” to a correspondence $`\stackrel{~}{𝒞}:XX`$ one can often show, using mixing properties of $`\stackrel{~}{𝒞}`$, that the values of $`f`$ at $`P`$ and $`\stackrel{~}{𝒞}(P)`$ will be uncorrelated, at least if $`P`$ is chosen at random w.r.t the the uniform measure on $`X`$.
However, since the $`Y_i`$ are becoming equidistributed, it amounts to almost the same thing to choose $`P`$ at random w.r.t. $`\nu _i`$ and w.r.t. the uniform measure on $`X`$. Thus, for $`\nu _i`$-typical $`PY_i`$, the values of $`f`$ at $`P`$ and $`𝒞(P)`$ are uncorrelated, whereas the values of $`\psi _j`$ at $`P`$ and $`𝒞(P)`$ are correlated. One can then play these phenomena against each other to obtain cancellation in the period integral $`f\psi _j𝑑\nu _i`$.
#### 1.3.4. A concrete example.
We shall now discuss how to bound Fourier coefficients of a modular form by the methods just described. Although the material below is essentially redone - with $`\mathrm{SL}(2,)`$ replaced by a general group – in Sec. 3, the example below was very important in motivating the author’s intuition, and it seems worthwhile to include it in the introduction.
Let $`\mathrm{\Gamma }\mathrm{SL}(2,)`$ be a lattice containing the element $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$. Let $`f(z)`$ be a holomorphic form of weight $`2`$ w.r.t. $`\mathrm{\Gamma }`$, which we write in a Fourier expansion $`f(z)=_{n=1}^{\mathrm{}}a_ne^{2\pi inz}`$. Hecke proved the bound $`|a_n|Cn`$, a bound which was only improved (for a general – possibly nonarithmetic – $`\mathrm{\Gamma }`$) much later, to $`|a_n|Cn^{5/6}`$, by A. Good. We shall sketch a simple proof of a nontrivial bound $`|a_n|Cn^{1\delta }`$ along the lines just indicated; for further details, we refer the reader to Sec. 3.2, where the procedure outlined is implemented for a general semisimple group.
We note that the ideas that will enter here are exactly those that will enter into the proof of equidistribution of sparse subsets of horocycles (see Sec. 3.1), or for the nontrivial bound for $`L^{\mathrm{}}`$ norms in the weight aspect that is discussed in Sec. 1.2. The proof below also works for Maass forms (in that case the result is due to Sarnak).
The Fourier expansion implies that
(1.6)
$$a_n=e^{2\pi }_{x/}f(x+\frac{i}{n})e^{2\pi inx}.$$
In words, the idea is as follows: the function $`e^{2\pi inx}`$ takes the same values at $`x,x+\frac{1}{n},x+\frac{2}{n},\mathrm{}`$. On the other hand, the values of the function $`f`$ at these points are (in a quantifiable sense) uncorrelated, as we shall deduce from the mixing properties of the horocycle flow. Playing these two properties against each other will yield an improvement of the Hecke bound for $`|a_n|`$. <sup>4</sup><sup>4</sup>4Underlying this is the usual “van der Corput” trick: to bound $`_{k=1}^Kc_k`$ it suffices to bound correlations $`_{k=1}^Kc_kc_{k+h}`$; in effect we apply this with $`c_k=f(\frac{k+i}{n}),K=n`$.)
Let $`\stackrel{~}{f}`$ be the lift of $`f`$ to $`\mathrm{\Gamma }\backslash \mathrm{SL}(2,)`$; that is to say,
$$\stackrel{~}{f}:\mathrm{\Gamma }\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)f(\frac{ai+b}{ci+d})(ci+d)^2.$$
Let $`x_n=\mathrm{\Gamma }\left(\begin{array}{cc}n^{1/2}& 0\\ 0& n^{1/2}\end{array}\right)`$, and put $`n(t)=\left(\begin{array}{cc}1& t\\ 0& 1\end{array}\right)`$. Then the definitions show that $`\stackrel{~}{f}(x_nn(t))=n^1f(\frac{i+t}{n})`$; consequently, we see that
(1.7)
$$a_n=e^{2\pi }_{t=0}^n\stackrel{~}{f}(x_nn(t))e^{2\pi it}𝑑t$$
(1.7) expresses the $`n`$th Fourier coefficient of $`f`$ as the integral of $`\stackrel{~}{f}`$ over a closed horocycle of length $`n`$. Moreover, (1.7) falls into the pattern of Sec. 1.3.1, with $`G_1=\mathrm{SL}_2()`$, $`G_2=\{n(t):t\}`$, $`Y_n=\{x_nn(t):t\}`$, and $`\psi _n:Y_n`$ the function given by $`x_nn(t)e^{2\pi it}`$. The fact that the $`Y_n`$ are becoming equidistributed amounts to the “equidistribution of low horocycles,” cf. . In the language of Sec. 1.3.1, we will take $`\sigma `$ to be the measure on $`G_2`$ that is a sum of point masses $`\delta _i`$, for integers $`i=1,\mathrm{},K`$. We now carry out the procedure of Sec. 1.3.1 in an explicit fashion in the paragraphs that follow.
Let $`T`$ be the operation of right translation by $`n(1)`$ on $`C^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$: that is to say, for $`FC^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$, we put $`TF(g)=F(gn(1))`$.
The value of the right-hand side of (1.7) remains unchanged if we replace $`\stackrel{~}{f}`$ by $`T\stackrel{~}{f}`$; consequently, for any integer $`K1`$, we have
$$a_n=\frac{e^{2\pi }}{K}_{t=0}^n(\underset{i=0}{\overset{K1}{}}T^i\stackrel{~}{f}(x_nn(t))e^{2\pi it}dt.$$
Applying the Cauchy-Schwarz inequality we deduce that
(1.8)
$$|a_n|^2\frac{n}{e^{4\pi }K^2}_{t=0}^n\left|\underset{i=0}{\overset{K1}{}}T^i\stackrel{~}{f}(x_nn(t))\right|^2𝑑t.$$
We now use come to the equidistribution part of the argument. The equidistribution of long closed horocycles asserts that the closed horocycle $`\{x+iy:0x1\}`$ becomes equidistributed in $`\mathrm{\Gamma }\backslash `$ as $`y\mathrm{}`$. Quantitatively, for any $`FC^{\mathrm{}}(\mathrm{\Gamma }\backslash )`$, we have
(1.9)
$$\left|_0^1F(x+iy)𝑑x_{\mathrm{\Gamma }\backslash }F\right|C_Fy^\delta ,$$
for some $`C_F`$ depending on $`F`$, and some $`\delta `$ depending only on $`\mathrm{\Gamma }`$. This assertion, originally proved by Sarnak by spectral methods, can be deduced quite easily from the mixing properties of the geodesic flow; this is done, in a somewhat more general context, in Lem. 9.6.
We note that – a special case of the discussion in Sec. 1.3.2 – the equidistribution statement (1.9) above reflects a mean-value theorem for periods. Indeed, if one applies it to $`F=y^2|f|^2`$, one deduces the asymptotic for $`_{n<X}|a_n|^2`$.
In any case, what will be more useful is the version of (1.9) that is lifted to $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$. This asserts that for any $`FC^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$, we have
(1.10)
$$\left|\frac{1}{n}_{t=0}^nF(x_nn(t))𝑑t_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}F(g)𝑑g\right|C_Fn^\delta ,$$
where $`\delta `$ is an explicit constant depending only on $`\mathrm{\Gamma }`$, and $`C_F`$ is a constant depending on $`F`$.
From (1.8) and (1.10) we conclude that
(1.11)
$$|a_n|^2\frac{e^{4\pi }n^2}{K^2}\left(\underset{i=0}{\overset{K1}{}}T^i\stackrel{~}{f}_{L^2(\mathrm{\Gamma }\backslash \mathrm{SL}_2())}^2+C_{f,K}n^\delta \right).$$
On the other hand, the explicit derivation of (1.10) shows that $`C_F`$ may be bounded by a Sobolev norm of $`F`$, and consequently the constant $`C_{f,K}`$ that appears in (1.11) is bounded by $`O_f(K^A)`$ for some $`A>0`$. Thus
$$|a_n|^2_fn^2K^2\left(\underset{i=0}{\overset{K1}{}}T^i\stackrel{~}{f}_{L^2(\mathrm{\Gamma }\backslash \mathrm{SL}_2())}^2+K^An^\delta \right).$$
We now use the fact that the horocycle flow is mixing, in a quantifiable way. This amounts to the assertion that there is an explicit $`\delta ^{}>0`$ and constant $`C_f^{}`$ such that, for $`i`$, $`|T^i\stackrel{~}{f},\stackrel{~}{f}|C_f^{}(1+|i|)^\delta ^{}`$. It follows easily that
$$\underset{i=0}{\overset{K1}{}}T^i\stackrel{~}{f}_{L^2}^2_fK^{2\delta ^{}}.$$
We conclude that $`|a_n|n(K^{\delta ^{}/2}+K^{A/21}n^{\delta /2})`$. Taking $`K`$ to be a sufficiently small power of $`n`$, we conclude that $`a_n`$ is bounded by $`n^{1\delta ^{\prime \prime }}`$ for some $`\delta ^{\prime \prime }>0`$ depending only on $`\mathrm{\Gamma }`$.
Clearly $`\delta ^{\prime \prime }`$ depends only on the spectral gap of $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$. Of course, this dependence does not arise in the “spectral” methods. It can be removed in the above method, but this seems to require some extra input, e.g. the finite-dimensionality of the space of functionals on an irreducible $`\mathrm{SL}_2()`$-representation that are invariant under the subgroup $`\{n(t):t\}`$.
#### 1.3.5. Two other viewpoints on the method of 1.3.4
There are two other viewpoints which might be helpful in thinking about Section 1.3.4. Both of these viewpoints do not literally generalize to the other situations we consider (e.g. triple products) but may be helpful for intuition.
1. The first is based on the following simple principle: suppose that $`T`$ is a measure-preserving transformation of the probability space $`(Y,\nu )`$, and that $`T`$ is ergodic. If $`\mu _1,\mu _2`$ are two $`T`$-invariant probability measures with average $`\frac{\mu _1+\mu _2}{2}=\nu `$, then $`\mu _1=\mu _2=\nu `$; this follows because $`\nu `$ is an extreme point of the convex set of $`T`$-invariant probability measures. More generally, given any family of probability measures averaging to $`\nu `$, they must almost all equal $`\nu `$.
We will apply this to $`Y=\mathrm{SL}_2()\backslash \mathrm{SL}_2()`$ and $`T`$ the operation of translation by $`n(1)`$.
Let $`n`$ be large; for $`t/`$, let $`\mu _t`$ be the probability measure that corresponds to normalized counting measure on $`\{x_nn(t+k):k,0k<n\}`$. Here notation is as prior to (1.7).
Then $`_0^1\mu _t`$ is the measure on the closed horocycle $`\{x_nn(t):0tn\}`$. Thus the family of measures $`\mu _t`$ averages to the measure on a long closed horocycle which, as we remarked earlier (see (1.10)) approximates the $`\mathrm{SL}_2()`$-invariant measure $`dg`$ on $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$. But this latter measure $`dg`$ is ergodic w.r.t. $`T`$. So applying a more quantitative form of the principle discussed above shows that, for almost all $`t[0,1]`$, $`\mu _t`$ must be close to $`dg`$. It is simple to see that one can use this to deduce bounds for the Fourier coefficients, via (1.7).
2. We will phrase the second rather imprecisely. Consider (1.7). Our strategy of proof can be rephrased as:
The function $`t\stackrel{~}{f}(x_nn(t))`$ is weak-mixing, whereas the function $`te^{2\pi it}`$ is periodic, and a weak-mixing function cannot correlate with a periodic function.
To explain this statement, we need to explain what it means for a function on the real line to be weak-mixing. Consider instead the case of a function $`h:`$. Furstenberg’s correspondence principle asserts that one can (loosely speaking) associate to this a dynamical system $`(Y,\nu ,T)`$ in such a way that (again loosely speaking) $`h`$ arises by sampling a function on $`Y`$ along a generic trajectory $`y_0,T(y_0),T^2(y_0),\mathrm{}`$. We then say that $`h`$ is weak-mixing if the system $`(Y,\nu ,T)`$ is so. The fact that our function $`t\stackrel{~}{f}(x_nn(t))`$ (say, when restricted to integer times) is weak-mixing follows from the equidistribution of long closed horocycles together with the fact that the horocycle flow is, itself, mixing.
For more on this point of view, see e.g. \[36, Section 4\] and \[36, Lemma 5.2\] for a version of the statement that weak-mixing functions cannot correlate with periodic ones.
### 1.4. Connection to existing methods.
The following comments pertain to the results of the present paper that concern subconvex bounds for $`L`$-functions. As we have emphasized above, the methods presented here are, upon examination, seen to be closely related to existing methods: in particular, “Sarnak’s spectral method,” which gave the only hitherto known instance of subconvexity over a base other than $``$.
Indeed, as we have already indicated, the equidistribution step of our method can be seen as the geometric version of a mean value theorem, and the rest of the method can be seen as an amplification step (or “shortening the family.”) Nevertheless, the key features of the present method are that it is essentially geometric (in that it avoids Fourier coefficients) and adelic (which allows us to import much from the modern theory of automorphic forms); it also does not use families in any explicit way. Once the notation is established – admittedly a nontrivial overhead – the method allows for very considerable technical simplification.
It is perhaps also noteworthy that the method given here does not require any exponential decay information for triple products. Although such exponential decay information is central only to subconvexity in the eigenvalue aspect, it has thus far entered as a technical device even in treatments of the level aspect.
In a sense, the present method bears the same relation to existing methods as adelic methods do to classical methods in the theory of automorphic forms. The classical situation has the advantage of concreteness, and whatever can be carried out in the adelic setting can be (in principle) carried out in the classical setting. However there is often a considerable technical and conceptual advantage in working adelically.
As we have discussed, the connection between equidistribution results and mean-value theorems for periods – implicitly exploited throughout this paper – appears in the work of V. Vatsal.
D. Hejhal considered ideas similar to that of Sec. 3.2 in the context of proving bounds towards Fourier coefficients; see . In the language of this paper, his method used a measure $`\sigma `$ (notation of Sec. 1.3) with much larger support, and consequently he was unable to get unconditional results.
Finally, as was remarked in Sec. 1.1, the main result of Sec. 4.1 is the analogue in the level aspect of a recent result of Bernstein-Reznikov : they establish a “subconvex bound” on triple products as the eigenvalue of one factor varies. Their methods also are geometric, avoiding the use of Fourier coefficients.
### 1.5. Acknowledgements.
This paper grew out of my proof of Thm. 3.1. The original proof was significantly more complicated, and I am indebted to Elon Lindenstrauss for his insistence that Thm. 3.1 should amount to nothing more than equidistribution and mixing. It was thus his intuition that led to a simplification of the proof and an important step in my understanding. The idea that the methods for Thm. 3.1 might be applicable in a more general setting arose during conversations with Andreas Strömbergsson, who also made many valuable suggestions about an early version of this paper. I thank them both for their significant contributions.
I am very grateful to Gergely Harcos and Philippe Michel for their encouragement of this project. Philippe read carefully an early draft of this paper and pointed out many points where the argument and results could be significantly improved. I am also grateful to Peter Sarnak, from whom I learned much of what I know about this subject.
I have also benefited from several conversations with Joseph Bernstein and André Reznikov. I thank for their generosity in sharing and discussing their elegant ideas.
I have learned many of the methods that appear here from the work of others. I mention in particular Peter Sarnak’s paper , which uses the idea of changing the test vector; the Friedlander-Iwaniec idea of amplification and the geometric version of it that appears in Bourgain-Lindenstrauss ; and the recent work of Bernstein-Reznikov , in particular their elegant use of Sobolev norms.
This paper suffered a considerable delay before submission. I would like to thank Philippe Michel for his encouragement and insistence that it be revised and submitted, without which the delay would have likely been considerably longer. I also thank Nicolas Bergeron and Marina Ratner for comments that improved the exposition and correctness.
The ideas of this paper were worked out during the workshop “Emerging applications of measure rigidity,” AIM, San Francisco and at the Isaac Newton Institute.
I was supported by the Clay Mathematics Institute during much of the writing of the paper, and I thank them for their generous support. I also thank the Institute for Advanced Study for providing excellent working conditions during the academic year 2005-2006. I was also partially supported by NSF grants DMS-0111298 and DMS-0245606.
### 1.6. Structure of paper.
The logical structure of this paper is as follows: Sec. 2 introduces all necessary notation. The heart of the paper are Sec. 3 (unipotent periods), Sec. 4 (the triple product period), Sections 6 and 7 (torus periods). The remaining Sections 811 are of a technical nature, proving various technical results required in the main text; at a first reading (or even later) they should perhaps be referred to only as necessary.
The two examples that best convey the flavor of the paper are Theorem 3.1 and Prop. 4.1. The proofs of these results are relatively self-contained, and we advise that the reader start with them.
## 2. Notation.
### 2.1. General notation.
We use the symbol $``$ as is standard in analytic number theory: namely, $`AB`$ means that there exists a constant $`c`$ such that $`AcB`$. The notation $`A_{f,g,h}B`$ means that the constant $`c`$ may depend on the quantities $`f,g,h`$; the notation $`A_ϵB`$ or $`A_\epsilon B`$ will mean, unless otherwise indicated, that the stated bound holds for all $`ϵ`$ or $`\epsilon >0`$. In general, we will never explicate the dependence of implicit constants on the number field over which we work; and, by an abuse of terminology, we will sometimes use the phrase “absolute constant” to mean a constant that depends only on this number field.
If $`Z`$ is a space we denote by $`\delta _z`$ the point measure at $`zZ`$, i.e. $`\delta _z(f)=f(z)`$ for $`f`$ a continuous function on $`Z`$.
Now let $`Z`$ be a right $`G`$-space. For $`f`$ a function on $`Z`$ and $`gG`$, we write $`gf`$ for the right translate of $`f`$ by $`g`$, i.e. $`gf(z)=f(zg)`$. If $`\mu `$ is a measure on $`Z`$, we define the translate $`g\mu `$ by the rule $`g\mu (gf)=\mu (f)`$. In particular, if $`\mu =\delta _z`$ is the point mass at $`zZ`$, then $`g\mu =\delta _{zg^1}`$ is the point mass at $`zg^1`$.
If $`\sigma `$ is a compactly supported measure on $`G`$, we set $`f\sigma \stackrel{\mathrm{def}}{=}_g(gf)𝑑\sigma (g)`$, i.e. $`f\sigma (z)=_{gG}f(zg)𝑑\sigma (g)`$. In particular, if $`\delta _{g_0}`$ is the point-mass at $`g_0`$, then $`f\delta _{g_0}=g_0f`$ is the right translate of $`f`$ by $`g_0`$.
If $`\sigma _1,\sigma _2`$ are two compactly supported measures on $`G`$, we define the convolution $`\sigma _1\sigma _2`$ to be the pushforward to $`G`$ of $`\sigma _1\times \sigma _2`$ on $`G\times G`$, under the multiplication map $`(g_1,g_2)G\times Gg_1g_2`$. Notations as above, one has the (somewhat unfortunate) compatibility relation $`(f\sigma _2)\sigma _1=f(\sigma _1\sigma _2)`$.
For $`\sigma `$ a measure on a group $`G`$, we denote by $`\stackrel{ˇ}{\sigma }`$ the image of $`\sigma `$ by the involution $`gg^1`$, and by $`\sigma `$ the total variation of $`\sigma `$.
If $`G`$ is a Lie group, we denote by $`\mathrm{Ad}(g)`$ the endomorphism “$`XgXg^1`$” of its Lie algebra.
If $`BA`$ is a finite index subgroup of the group $`A`$, then we denote by $`[A:B]`$ the index of $`B`$ in $`A`$.
If $`h`$ is an entire function, the notation $`_{\mathrm{}(s)=\sigma }h(s)𝑑s`$ denotes the line integral along the line $`\mathrm{}(s)=\sigma `$ from $`\sigma i\mathrm{}`$ to $`\sigma +i\mathrm{}`$. The notation $`_{\mathrm{}(s)1}h(s)𝑑s`$ denotes $`_{\mathrm{}(s)=\sigma }h(s)𝑑s`$ for sufficiently large $`\sigma `$; in the contexts where we use this notation, the answer will be constant when $`\sigma `$ is sufficiently large.
### 2.2. Classical modular forms.
As usual $``$ denotes the upper half plane, i.e. $`\{z:\mathrm{Im}(z)>0\}`$. It admits the usual action of $`\mathrm{SL}(2,)`$ by fractional linear transformations.
### 2.3. Number fields and associated notations.
Let $`F`$ be a number field. Throughout the paper we shall regard $`F`$ as fixed: that is to say, we allow implicit constants in $`,`$ may depend on $`F`$ without explicit statement.
We set $`F_{\mathrm{}}=F`$, $`𝔸_F`$ the ring of adeles of $`F`$, $`𝔸_{F,f}`$ the ring of finite adeles. Thus $`𝔸_F=F_{\mathrm{}}\times 𝔸_{F,f}`$. We will fix once and for all an additive character $`e_F:𝔸_F/F`$, and denote by $`e_{F_v}`$ the induced additive character of $`F_v`$.
For each place $`v`$ we have a canonical “absolute value” $`x|x|_v`$ on $`F_v^\times `$, namely, $`|x|_v=\mathrm{meas}(xS)/\mathrm{meas}(S)`$ for any Haar measure, $`\mathrm{meas}`$, on $`F_v^\times `$, and any subset $`S`$ of positive measure.
The same definition defines a character $`𝔸_F^\times /F^\times _{>0}`$, which we denote by $`a|a|_𝔸`$, or simply by $`a|a|`$ if it is clear from context. We denote by $`𝔸_F^1`$ the subgroup of $`𝔸_F^\times `$ consisting of adeles of norm $`1`$; then the quotient $`𝔸_F^1/F^\times `$ is compact. For a finite place $`v`$ of $`F`$, we denote by $`𝔬_{F_v}`$ the maximal compact subring of the completion $`F_v`$, by $`𝔮_v`$ the maximal ideal of $`𝔬_{F_v}`$, and by $`q_v`$ the cardinality of the residue field.
We shall generally denote ideals of $`𝔬_F`$ by gothic letters $`𝔩,𝔮,𝔫`$, etc. If $`𝔣`$ is an integral ideal of $`𝔬_F`$, we set $`\mathrm{N}(𝔣):=|𝔬_F/𝔣|`$ to be its norm. Moreover, we shall denote $`𝔬_𝔣:=_{𝔮|𝔣}𝔬_𝔮`$. Here $`𝔬_𝔮`$ denotes the completed ring, not the localized ring, i.e. $`𝔬_𝔣`$ is the inverse limit of the rings $`𝔬_F/𝔣^N`$.
We denote by $`𝔡`$ the different of the character $`e_F`$, i.e. $`𝔡`$ is a fractional ideal so that $`𝔡_v^1`$ is, for every finite place $`v`$, the largest $`𝔬_{F_v}`$-submodule of $`F_v`$ upon which $`e_F`$ is trivial.
### 2.4. Adele groups and their function spaces
Let $`𝐆`$ be a connected reductive algebraic group over a number field $`F`$, and let $`𝐙`$ be its center. Denote by $`𝔸_{F,f}`$ the ring of finite adeles, and fix for each finite place $`v`$ a maximal open compact subgroup $`K_{v,𝐆}𝐆(F_v)`$ with the property that $`K_{\mathrm{max},𝐆}:=_{v\mathrm{finite}}K_{v,𝐆}`$ is a maximal open compact subgroup of $`𝐆(𝔸_{F,f})`$. Put $`𝐗_𝐆=𝐆(F)\backslash 𝐆(𝔸_F)`$, $`𝐗_{𝐆,\mathrm{ad}}=𝐙(𝔸_F)𝐆(F)\backslash 𝐆(𝔸_F)`$. Then $`𝐗_{𝐆,\mathrm{ad}}`$ has finite volume with respect to any $`𝐆(𝔸_F)`$-invariant measure.
Let $`\omega :𝐙(𝔸_F)^\times `$ be a unitary character. We define the space $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ to be the space of functions on $`𝐗_𝐆`$ whose stabilizer in $`K_{\mathrm{max},𝐆}`$ has finite index, which transform under $`𝐙(𝔸_F)`$ by $`\omega `$, and so that the function $`gf(xg)`$ is a $`C^{\mathrm{}}`$ function of $`g𝐆(F_{\mathrm{}})`$, for each $`x𝐗_𝐆`$. Similarly one defines an $`L^2`$-space $`L_\omega ^2(𝐗_𝐆)`$, or simply $`L^2`$ if the central character $`\omega `$ is clear from context, by completing the space of compactly supported functions in $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ with respect to the Hilbert norm $`f_2:=\left(_{𝐗_{𝐆,\mathrm{ad}}}|f(g)|^2𝑑g\right)^{1/2}`$.
For $`\psi C_\omega ^{\mathrm{}}(𝐗_𝐆)`$, we denote by $`K_{v,\psi }`$ the stabilizer of $`\psi `$ in $`K_{v,𝐆}`$, and put
(2.1)
$$K_\psi =\underset{v\mathrm{finite}}{}K_{v,\psi }.$$
We note that $`K_\psi `$ is, in general, a proper subgroup of the stabilizer of $`\psi `$ in $`K_{\mathrm{max},𝐆}`$.
For $`\psi C_\omega ^{\mathrm{}}(𝐗_𝐆)`$, we define the finite set of places $`\mathrm{Supp}(\psi )`$ to be those finite $`v`$ for which $`K_{v,𝐆}`$ does not fix $`\psi `$, i.e.
(2.2)
$$\mathrm{Supp}(\psi )\stackrel{\mathrm{def}}{=}\{v:K_{v,𝐆}K_{v,\psi }\}.$$
It is convenient to introduce some notions of “size” on $`𝐆(𝔸_F)`$. Let $`𝔤`$ be the Lie algebra of $`𝐆(F_{\mathrm{}})`$. It is a finite dimensional real vector space; fix an arbitrary norm on it. For $`g_{\mathrm{}}𝐆(F_{\mathrm{}})`$, we denote by $`g_{\mathrm{}}`$ the operator norm of the adjoint endomorphism $`\mathrm{Ad}(g_{\mathrm{}}^1):𝔤𝔤`$. If $`v`$ is a finite place of $`F`$ and $`g_v𝐆(F_v)`$, we set $`g_v=[K_{v,𝐆}g_vK_{v,𝐆}:K_{v,𝐆}]`$, i.e. the number of right- $`K_{v,𝐆}`$ cosets in $`K_{v,𝐆}g_vK_{v,𝐆}`$. For $`g_f=(g_v)_{v\mathrm{finite}}𝐆(𝔸_{F,f})`$ we put $`g_f=_vg_v`$. Finally for $`g_𝔸=(g_{\mathrm{}},g_f)𝐆(F_{\mathrm{}})\times 𝐆(𝔸_{F,f})`$, set $`g_𝔸=g_{\mathrm{}}g_f`$.
We remark that $`g_{\mathrm{}},g_f,g_𝔸`$ are all invariant by the center of $`𝐆`$.
### 2.5. The groups $`𝐆=\mathrm{GL}(2)`$ and $`𝐆=\mathrm{PGL}(2)`$ and some of their subgroups.
We will deal most often with the cases of $`𝐆=\mathrm{GL}(2)`$ (resp. $`𝐆=\mathrm{PGL}(2)`$). In that setting we shall write $`𝐗_{\mathrm{GL}(2)}`$ (resp. $`𝐗`$) for $`𝐗_𝐆`$.
We will make use of the following algebraic subgroups of $`\mathrm{GL}_2`$, which we will often also regard as algebraic subgroups of $`\mathrm{PGL}_2`$ in the obvious way:
$$N=\left(\begin{array}{cc}1& \\ 0& 1\end{array}\right),B=\left(\begin{array}{cc}& \\ 0& \end{array}\right),A=\left(\begin{array}{cc}& 0\\ 0& \end{array}\right),Z=\left(\begin{array}{cc}x& 0\\ 0& x\end{array}\right).$$
If $`R`$ is any ring and $`xR,yR^\times `$, we denote<sup>5</sup><sup>5</sup>5 (In Section 3.1 alone, we will use slightly different notation for $`a(y)`$ to accomodate the fact that we deal with $`\mathrm{SL}_2`$ rather than $`\mathrm{GL}_2`$. We make the relevant notation clear in that section.
(2.3)
$$\begin{array}{c}n(x)=\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right),\overline{n}(x)=\left(\begin{array}{cc}1& 0\\ x& 1\end{array}\right),a(y)=\left(\begin{array}{cc}y& 0\\ 0& 1\end{array}\right),\hfill \\ \hfill a^{}(y)=\left(\begin{array}{cc}1& 0\\ 0& y\end{array}\right),z(y)=\left(\begin{array}{cc}y& 0\\ 0& y\end{array}\right),w=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).\end{array}$$
all elements of $`\mathrm{GL}_2(R)`$.
If $`v`$ is a place of $`F`$ and $`xF_v,yF_v^\times `$, we denote by $`n_v(x)`$ (resp. $`a_v(x)`$) the element $`n(x)`$ (resp. $`a(y)`$) considered as an element of $`\mathrm{GL}_2(𝔸_F)`$ via the natural inclusion $`\mathrm{GL}_2(F_v)\mathrm{GL}_2(𝔸_F)`$.
For each place $`v`$, we let $`K_v`$ be the standard maximal compact subgroup of $`\mathrm{GL}_2(F_v)`$, i.e. $`K_v`$ is the stabilizer of the norm on $`F_v^2`$ given by $`\sqrt{|x|_v^{2/\mathrm{deg}(v)}+|y|_v^{2/\mathrm{deg}(v)}}`$ if $`v`$ is archimedean, where $`\mathrm{deg}(v)`$ is the degree<sup>6</sup><sup>6</sup>6Recall that $`|x|_v`$, for a complex place $`v`$ and $`xF_v`$, is the square of the usual absolute value on $``$! of $`F_v`$ over $``$; and $`\mathrm{max}(|x_v|,|y_v|)`$ if $`v`$ is nonarchimedean. Thus, in particular, $`K_v=\mathrm{GL}_2(𝔬_{F,v})`$ if $`v`$ is nonarchimedean. We put $`K_{\mathrm{max}}=_{v\mathrm{finite}}K_v`$. $`K_v`$ (respectively $`K_{\mathrm{max}}`$) is a maximal compact subgroup of $`\mathrm{GL}_2(F_v)`$ (respectively $`\mathrm{GL}_2(𝔸_{F,f})`$), and (by projection) can also be regarded as a maximal compact subgroup of $`\mathrm{PGL}_2(F_v)`$ (respectively $`\mathrm{PGL}_2(𝔸_{F,f})`$). Similarly $`K_{\mathrm{max}}\times K_{\mathrm{}}`$ is a maximal compact subgroup of $`\mathrm{GL}_2(𝔸_F)`$, and may also be regarded as a maximal compact subgroup of $`\mathrm{PGL}_2(𝔸_F)`$.
For $`𝔮`$ a finite prime of $`F`$, we denote by $`\varpi _𝔮F_𝔮`$ a uniformizer, and by $`[\varpi _𝔮]`$ the element of $`𝔸_F^\times `$ that is the image of $`\varpi _𝔮`$ under the natural inclusion $`F_𝔮^\times 𝔸_F^\times `$.
Let $`𝔮`$ be a finite prime of $`F`$. It will be convenient to define certain open compact subgroups of $`K_𝔮`$. For each $`e_𝔮>0`$, we define $`K[𝔮^{e_𝔮}]K_𝔮`$ (resp. $`K_0[𝔮^{e_𝔮}]K_𝔮`$) to be the be the kernel of $`\mathrm{GL}_2(𝔬_𝔮)\mathrm{GL}_2(𝔬_𝔮/\varpi _𝔮^m)`$ (resp. the preimage, under this map, of the upper triangular matrices). Thus
$$K_0[𝔮^{e_𝔮}]=\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):a,b,d𝔬_𝔮,c𝔮^{e_𝔮},adbc𝔬_𝔮^\times \}.$$
Now let $`𝔣`$ be a fractional ideal, not necessarily prime, of $`F`$. Factorize $`𝔣=_𝔮𝔮^{e_𝔮}`$ into prime ideals. We define elements $`[𝔣]𝔸_F^\times ,a([𝔣]),n([𝔣])\mathrm{GL}_2(𝔸_F)`$ via:
(2.4)
$$[𝔣]=\underset{𝔮|𝔣}{}[\varpi _𝔮]^{e_𝔮},n([𝔣]):=\underset{𝔮|𝔣}{}n_𝔮(\varpi _𝔮^{e_𝔮}),a([𝔣])=\underset{𝔮|𝔣}{}a_𝔮(\varpi _𝔮^{e_𝔮}).$$
Suppose $`\chi :𝔸_F^\times /F^\times ^\times `$ is a character. We define
$$\chi (𝔣)=\{\begin{array}{cc}0,\chi \text{ ramified at any place dividing }𝔣,\hfill & \\ _{𝔮|𝔣}\chi (\varpi _𝔮)^{e_𝔮},\text{ else.}\hfill & \end{array}$$
### 2.6. Measures.
The choice of measure is not especially important, as we are only interested in upper bounds; thus, so long as we are consistent, the precise selection does not matter. We choose a “standard” set of measures here; at times in the text, especially when carrying out equidistribution arguments, it will be more convenient to use probability measures, and we will indicate when this is the case.
We denote by $`\mu _𝐗`$ the $`\mathrm{PGL}_2(𝔸_F)`$-invariant probability measure on $`𝐗`$. We shall sometimes simply denote it by $`dx`$.
Let $`v`$ be a finite place of $`F`$. Unless explicitly stated otherwise, the measures on $`\mathrm{GL}_2(F_v)`$, $`\mathrm{PGL}_2(F_v)`$, $`F_v`$ and $`F_v^\times `$ are the Haar measure which assigns $`\mathrm{GL}_2(𝔬_{F_v})`$ (resp. $`\mathrm{PGL}_2(𝔬_{F_v})`$, $`𝔬_{F_v}`$, $`𝔬_{F_v}^\times `$) the total mass $`1`$.
For $`v`$ archimedean, endow $`F_v`$ with a multiple of Lebesgue measure $`c_vdx`$, where the constants $`c_v`$ are fixed arbitrarily in such a way that the induced product measure on $`F_{\mathrm{}}`$ satisfies $`\mathrm{vol}(F_{\mathrm{}}/𝔬_F)=1`$; equivalently, the product measure on $`𝔸_F`$ satisfies $`\mathrm{vol}(𝔸_F/F)=1`$. In particular, this product measure on $`𝔸_F`$ is self-dual with respect to $`e_F`$. We endow $`F_v^\times `$ with the measure $`d^\times x=\frac{dx}{|x|_v}`$, where $`dx`$ is Lebesgue measure.
These choices induce a Haar measure on $`N(F_v)`$, by means of the identification $`xn(x)`$; similarly, the identifications $`(y,y^{})a(y)a^{}(y^{})`$ and $`yz(y)`$ induce Haar measures on $`A(F_v)`$ and $`Z(F_v)`$. Equip $`K_v`$ with the measure of mass $`1`$, and give $`\mathrm{GL}_2(F_v)`$ the measure arising from the Iwasawa decomposition $`N(F_v)\times A(F_v)\times K_v`$. Equip $`\mathrm{PGL}_2(F_v)=\mathrm{GL}_2(F_v)/Z(F_v)`$ with the “quotient” measure.
We then take the measures on $`\mathrm{GL}_2(𝔸_F),\mathrm{PGL}_2(𝔸_F),𝔸_F,𝔸_F^\times `$ to be the corresponding product measures.
The measure on any discrete group (e.g. $`\mathrm{PGL}_2(F)`$, considered as a subgroup of $`\mathrm{PGL}_2(𝔸_F)`$) will be counting measure.
Usually (indeed, unless otherwise specified) we shall use the $`\mathrm{PGL}_2(𝔸_F)`$-invariant probability measure on $`𝐗=\mathrm{PGL}_2(F)\backslash \mathrm{PGL}_2(𝔸_F)`$. This does not coincide with the quotient measure induced from $`\mathrm{PGL}_2(𝔸_F)`$, but they differ by some constant depending only on $`F`$. On the few occasions we shall have occasion to use the latter measure, we will indicate this.
### 2.7. Projection onto locally constant functions.
For equidistribution questions it is usually convenient to deal with the constant function and its orthogonal complement separately. Some minor complications arise in our case since the ambient spaces are not connected. In fact: The space $`C^{\mathrm{}}(𝐗_𝐆)`$ is a direct limit of function spaces $`C^{\mathrm{}}(𝐗_𝐆/K)`$ where $`KK_{\mathrm{max},𝐆}`$ has finite index. Unless $`𝐆`$ is simply connected, the manifolds $`𝐗_𝐆/K`$ need not be connected.
Of course, to deal with this, one can (if $`𝐆`$ is semisimple) simply replaces the notion of constant function by locally constant function. However, in the general case of $`𝐆`$ reductive, matters are slightly complicated by the necessity of dealing with central characters.
Since we will only use this definition when $`𝐆`$ is a product of $`\mathrm{GL}(2)`$s, we restrict ourselves to that setting. First suppose that $`𝐆=\mathrm{GL}(2)`$. We define a projection $`𝒫:C_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})C_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$ via
(2.5)
$$𝒫f(x)=_{hSL_2(F)\backslash \mathrm{SL}_2(𝔸_F)}f(hx)𝑑h=\underset{\chi ^2=\omega }{}\chi (x)_𝐗f(y)\overline{\chi (y)}𝑑y,$$
where $`dh`$ is the $`\mathrm{SL}_2(𝔸_F)`$-invariant probability measure, $`dy`$ the $`\mathrm{GL}_2(𝔸_F)`$-invariant probability measure on $`𝐗`$, $`\chi `$ ranges over characters of $`𝔸_F^\times /F^\times `$ with square $`\omega `$, $`\chi (y)`$ the function on $`𝐗_{\mathrm{GL}(2)}`$ defined by $`g\chi (det(g))`$, and the second equality is easily verified. We note, in particular, that the $`\chi `$ sum is finite (any $`\chi `$ for which the corresponding term is nonvanishing must be unramified outside $`\mathrm{Supp}(f)`$).
Then $`𝒫f_L^{\mathrm{}}f_L^{\mathrm{}}`$, as is clear from the first equality of (2.5), and $`𝒫`$ is a self-adjoint projection w.r.t. $`L^2`$, as is clear from the second equality. We say a function $`f`$ is totally nondegenerate if $`𝒫f=0`$.
If $`𝐆=\mathrm{GL}(2)\times \mathrm{GL}(2)`$, and $`\omega =(\omega _1,\omega _2)`$ is a character of the center $`𝐙(𝔸_F)=𝔸_F^\times \times 𝔸_F^\times `$, we denote by $`𝒫_1`$ the operator on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ given by
$$𝒫_1f(x_1,x_2)=_{hSL_2(F)\backslash \mathrm{SL}_2(𝔸_F)}f(hx_1,x_2)𝑑h=\underset{\chi ^2=\omega _1}{}\chi (x_1)_𝐗f(y,x_2)\overline{\chi (y)}𝑑y.$$
We define $`𝒫_2`$ similarly, interchanging the role of the first and second coordinate. The operators $`𝒫_j`$ for $`j=1,2`$ commute, satisfy $`𝒫_jf_L^{\mathrm{}}f_L^{\mathrm{}}`$ and are commuting self-adjoint projections on $`L^2`$. We say that a function $`f`$ is totally nondegenerate if $`𝒫_1f=𝒫_2f=0`$.
###### Lemma 2.1.
Let $`v`$ be a place of $`F`$. The projection $`𝒫`$ acts by the identity on the subspace $`WL_\omega ^2(𝐗_𝐆)`$ spanned by one-dimensional representations of $`\mathrm{GL}_2(F_v)`$ occurring in $`L_\omega ^2(𝐗_𝐆)`$.
Similarly, $`𝒫_1`$ (resp. $`𝒫_2`$) acts by the identity on the space $`W_1`$ (resp. $`W_2`$) spanned by one-dimensional representations of $`\mathrm{GL}_2(F_v)`$ occurring in $`L^2(𝐗\times 𝐗)`$ for the action on the first (resp. second) factor.
###### Proof.
This follows from the spectral decomposition for $`\mathrm{GL}(2)`$. For instance, it is known that the space $`W`$ is precisely the span of functions of the form $`g\chi (det(g))`$, where $`\chi `$ ranges over characters of $`𝔸_F^\times /F^\times `$ satisfying $`\chi ^2=\omega `$. ∎
### 2.8. Hecke operators and bounds towards the Ramanujan conjecture.
Let $`𝔩`$ be a prime ideal of $`𝔬_F`$ and $`r`$ an integer $`1`$. Let $`F_𝔩`$ be the completion of $`F`$ at the prime $`𝔩`$. Take the Haar measure on $`\mathrm{GL}_2(F_𝔩)`$ so that it assigns mass $`1`$ to $`\mathrm{GL}_2(𝔬_{F_𝔩})`$. Define the measure $`\mu _{𝔩^r}^{}`$ on $`\mathrm{GL}_2(F_𝔩)`$ to the restriction of Haar measure to the set $`\mathrm{GL}_2(𝔬_{F_𝔩})\left(\begin{array}{cc}\varpi _𝔩^r& 0\\ 0& 1\end{array}\right)\mathrm{GL}_2(𝔬_{F_𝔩})`$, so that the total mass of $`\mu _{𝔩^r}^{}`$ is $`\mathrm{N}(𝔩)^{r1}(\mathrm{N}(𝔩)+1)`$.
Moreover, set
(2.6)
$$\mu _{𝔩^r}=\frac{1}{\mathrm{N}(𝔩)^{r/2}}\underset{k\frac{r}{2}}{}\mu _{r2k}^{},\overline{\mu }_{𝔩^r}:=\frac{\mu _{𝔩^r}}{\mu _{𝔩^r}},$$
where $``$ denotes total variation. Thus $`\overline{\mu }_{𝔩^r}`$ is a probability measure. Via the natural inclusion of $`\mathrm{GL}_2(F_𝔩)`$ in $`\mathrm{GL}_2(𝔸_{F,f})`$, we may regard $`\mu _{𝔩^r}`$ as a compactly supported measure on $`\mathrm{GL}_2(𝔸_{F,f})`$; by abuse of notation, we will not introduce a different symbol for this measure. If $`𝔫`$ is an integral ideal of $`𝔬_F`$, factorize $`𝔫=_i𝔩_i^{r_i}`$ and put $`\mu _𝔫=\mu _{𝔩_i^{r_i}},\overline{\mu }_𝔫=\overline{\mu }_{𝔩_i^{r_i}}`$. Here $``$ is taken to mean convolution of measures on $`\mathrm{GL}_2(𝔸_{F,f})`$.
Convolution by $`\mu _𝔫`$ on $`L^2(𝐗)`$ corresponds to the $`𝔫`$th Hecke operator; in this normalization the Ramanujan conjecture corresponds to it having eigenvalues $`2`$ in absolute value.
The adelic measures $`\mu _𝔫`$ satisfy the usual multiplication laws, appropriately interpreted: if $`𝔫`$ and $`𝔪`$ are ideals, then
(2.7)
$$_{\mathrm{PGL}_2(𝔸_F)}h(x)d(\mu _𝔫\mu _𝔪)(x)=\underset{𝔡|(𝔪,𝔫)}{}_{\mathrm{PGL}_2(𝔸_F)}h(x)𝑑\mu _{𝔫𝔪𝔡^2}(x),$$
whenever $`h`$ is a function on $`\mathrm{PGL}_2(𝔸_F)`$ that is invariant under $`\mathrm{PGL}_2(𝔬_v)`$ for all $`v|𝔫𝔪`$.
###### Definition 2.1.
Set $`\alpha `$ be a bound towards Ramanujan for $`\mathrm{GL}_2`$ over $`F`$, i.e. $`\alpha `$ is so that $`\mu _𝔩`$ acts on any $`\mathrm{PGL}_2(𝔬_{F_𝔩})`$-invariant cuspidal eigenfunction by an eigenvalue $`\mathrm{N}(𝔩)^\alpha +\mathrm{N}(𝔩)^\alpha `$ in absolute value.
Thus $`\alpha =0`$ corresponds to the Ramanujan conjecture, $`\alpha =1/2`$ the trivial bound. By work of Kim and Kim-Shahidi, we can take $`\alpha =3/26`$. For our applications, any value of $`\alpha `$ less than $`1/4`$ would suffice.
Note that we shall slightly vary this notation (but in a reasonably compatiable way) in Section 3.1 and Section 9.3.1. In those parts, we shall deal with a (not necessarily arithmetic) quotient $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$, and $`\alpha `$ will denote a number so that $`L^2(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$ does not contain any complementary series with parameter $`\alpha `$. (Here the complementary series is understood to be parameterized by $`(0,1/2)`$). This is compatible with the above notation, however; e.g. if $`\mathrm{\Gamma }`$ were a congruence subgroup, $`\alpha =3/26`$ would again be admissible.
### 2.9. Sobolev-type norms on real and adelic quotients.
#### 2.9.1. General comments.
Let $`M`$ be a real manifold. Recall that the Sobolev norm on $`C^{\mathrm{}}(M)`$ controls, roughly speaking, the $`L^p`$ norm of a function together with the $`L^p`$ norms of certain derivatives. These norms will be tremendously useful throughout the paper to control equidistribution rates. We shall use both the relatively simple definition when $`M=\mathrm{\Gamma }\backslash G`$ and an adelic variant.
First let us remark on the use of $`L^p`$-Sobolev norms for $`p>2`$. This is solely to do with noncompactness. If we were to deal only with compact quotients, then the $`L^2`$-Sobolev theory would always suffice. However, in the noncompact case, the $`L^2`$-Sobolev norms do not (e.g.) give good bounds on the size of a function high in a cusp. There are, of course, various ways to rectify this; for example we could include weights that measure the height into the cusp. We have chosen instead to use $`L^p`$-norms with $`p>2`$, which is technically very simple, but has some disadvantages (e.g. it does not induce a Hilbert space structure).
Note that we will allow our seminorms and norms to take the value $`\mathrm{}`$. Thus a seminorm on a complex vector space $`V`$ will be a function from $`V`$ to $`_0\{\mathrm{}\}`$ satisfying
1. $`\lambda v=|\lambda |v`$, for any $`vV`$ such that $`v<\mathrm{}`$;
2. $`v_1+v_2v_1+v_2`$ if both $`v_1`$ and $`v_2`$ are not infinite.
It is a norm if additionally $`v=0`$ implies $`v=0`$. Note that giving such a seminorm on $`V`$ is equivalent to giving a subspace $`V_fV`$ together with a finite-valued seminorm on $`V_f`$. Indeed take $`V_f=\{vV:v<\mathrm{}\}`$, equipped with the restriction of $``$.
We remark that we do not require that our norms be complete.
#### 2.9.2. Non-adelic setting.
Suppose $`\mathrm{\Gamma }G`$ is a lattice in a connected semisimple Lie group. Fix for all time a basis $``$ for the Lie algebra $`𝔤`$ of $`G`$ and a norm $``$ on $`𝔤`$. For $`gG`$, we denote by $`g`$ the operator norm of $`\mathrm{Ad}(g^1):𝔤𝔤`$, i.e. the map $`Xg^1Xg`$.
For $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$, and $`1p\mathrm{}`$, we put
(2.8)
$$S_{p,d}=\underset{\mathrm{ord}(𝒟)d}{}𝒟f(g)_{L^p(\mathrm{\Gamma }\backslash G)}.$$
Here $`𝒟`$ ranges over all monomials in $``$ of order $`d`$, and $`𝒟`$ acts on $`f`$ by right differentiation. (For example, $`X𝔤`$ acts on $`f`$ via $`Xf(g)=\frac{d}{dt}f(ge^{tX})`$.)
Changing $``$ only distorts $`S_{p,d}`$ by a bounded factor. (That is to say, if $`S_{p,d}^{}`$ is the norm obtained by replacing $``$ by another basis, then there are positive reals $`c_1,c_2`$, possibly depending on $`d`$, such that $`c_1S_{p,d}S_{p,d}^{}c_2S_{p,d}`$. )
We will often use the following simple remark: Fix a Riemannian metric $`d(,)`$ on $`G`$ and suppose $`gG`$ belongs to some fixed compact set. Then, for $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash G),x\mathrm{\Gamma }\backslash G`$, we have $`|f(xg)f(x)|S_{\mathrm{},1}(f)d(g,1)`$. Indeed, we may assume that $`g`$ is close to the identity and write $`g=\mathrm{exp}(X)`$, with $`X𝔤`$; now apply the mean value theorem to $`tf(xe^{tX})`$.
Moreover, the following elementary properties are easily verified (we only need them in the case $`p=\mathrm{}`$).
###### Lemma 2.2.
Let $`f_1,f_2C^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$ and $`gG`$. Then
(2.9) $`S_{\mathrm{},d}(f_1f_2)_dS_{\mathrm{},d}(f_1)S_{\mathrm{},d}(f_2)`$
$`S_{\mathrm{},d}(gf_1)_dg^dS_{\mathrm{},d}(f_1)`$
We remark that, in the case $`p=2`$, the rule (2.8) also defines a system of Sobolev norms on any unitary $`G`$-representation; the case discussed above corresponds to the unitary representation $`L^2(\mathrm{\Gamma }\backslash G)`$.
#### 2.9.3. Adelic Sobolev norms
Let’s first describe what the point is intended to be (evidently there are many ways of implementing it, cf. Rem. 2.1). We would like to put a norm on the adelic function space, suitable for controlling e.g. period integrals. Consider $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ in the case of $`𝐆=\mathrm{SL}_2`$, $`F=`$, $`\omega =1`$ as a direct limit of spaces $`C^{\mathrm{}}(\mathrm{\Gamma }_i\backslash \mathrm{SL}_2())`$, where $`\mathrm{\Gamma }_i`$ ranges over some class of congruence subgroups of $`\mathrm{\Gamma }_0:=\mathrm{SL}_2()`$. We equip each quotient $`\mathrm{\Gamma }_i\backslash \mathrm{SL}_2()`$ with the $`\mathrm{SL}_2()`$ invariant probability measure. Then, on each space $`C^{\mathrm{}}(\mathrm{\Gamma }_i\backslash \mathrm{SL}_2())`$ we have the norm $`S_{p,d}`$ defined in the previous section. On the other hand, typical bounds on automorphic forms have an implicit dependence on the “level”, i.e. the index $`[\mathrm{\Gamma }:\mathrm{\Gamma }_i]`$, so one would like to have a norm that increases with the level. The most naive candidate is, fixing a real number $`\beta >0`$, to define the “norm” of $`fC^{\mathrm{}}(\mathrm{\Gamma }_i\backslash \mathrm{SL}_2())`$ to be $`[\mathrm{\Gamma }:\mathrm{\Gamma }_i]^\beta S_{p,d}(f)`$. This unfortunately does not quite make sense when we pass to the direct limit: however, we can “force it to make sense” by considering the maximal norm on the direct limit whose restriction to each $`C^{\mathrm{}}(\mathrm{\Gamma }_i\backslash \mathrm{SL}_2())`$ is bounded above by $`[\mathrm{\Gamma }:\mathrm{\Gamma }_i]^\beta S_{p,d}(f)`$. This will suffice for our purposes.
Let us formalize these ideas. In what follows we return to the setting of $`𝐆`$ a reductive group over $`F`$. The adelic Sobolev norms will be a family of norms on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ indexed by a triple $`(p,d,\beta )`$. The $`d`$ and $`\beta `$ indicate, approximately speaking, how stringently one should “penalize” rapid variation at the infinite and finite places respectively.
Let $`p1,k,\beta 0`$. Fix a basis $`=\{X_i\}`$ for the real Lie group $`\mathrm{Lie}(𝐆(F_{\mathrm{}}))`$. Recalling the definition of $`K_\psi `$ from (2.1), we define the pre-Sobolev functions $`PS_{p,d,\beta }`$ on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ via:
(2.10)
$$\begin{array}{c}PS_{p,d,\beta }(\psi )=[K_{\mathrm{max},𝐆}:K_\psi ]^\beta \underset{\mathrm{ord}(𝒟)d}{}𝒟\psi _{L^p(𝐗_{𝐆,\mathrm{ad}})}\hfill \\ \hfill =\underset{v\mathrm{finite}}{}[K_{v,𝐆}:K_{v,\psi }]^\beta \underset{\mathrm{ord}(𝒟)d}{}𝒟\psi _{L^p(𝐗_{𝐆,\mathrm{ad}})},\end{array}$$
where the sum ranges over $`𝒟`$ that are monomials in $``$ of order $`d`$.
The function $`PS_{p,d,\beta }`$ does not satisfy the triangle inequality. We define the $`(p,d,\beta )`$-Sobolev norm $`S_{p,d,\beta }`$ to be the maximal seminorm on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ satisfying $`S_{p,d,\beta }(\psi )PS_{p,d,\beta }(\psi )`$. Explicitly,
(2.11)
$$S_{p,d,\beta }(\psi )=inf\{\underset{i=1}{\overset{n}{}}PS_{p,d,\beta }(\psi _j):\underset{i=1}{\overset{n}{}}\psi _j=\psi ,\psi _jC_\omega ^{\mathrm{}}(𝐗_𝐆).\}$$
In fact, it is clear that the right-hand side of (2.11) defines a seminorm that is dominated by $`PS_{p,d,\beta }`$ (take the collection $`\{\psi _i\}`$ to consist of $`\{\psi \}`$ alone); moreover, it is evidently maximal in the class of such seminorms. Finally, as $`PS_{p,d,\beta }(\psi )\psi _{L^p}`$, the $`S_{p,d,\beta }`$ are in fact norms on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$.
It will often be useful to omit the argument $`\beta `$ and set it to a “default” value of $`1/p`$. We therefore define $`S_{p,d}:=S_{p,d,1/p}`$, for $`p0`$, and $`S_{\mathrm{},d}:=S_{\mathrm{},d,0}`$.
Notational convention: We will very often have cause to bound linear functionals $`L`$ on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ by Sobolev norms. In writing statements of the form $`|L(f)|S_{p,d,\beta }(f)`$, we will always allow the implicit constant to depend on $`p,d`$ and $`\beta `$ without explicitly saying so.
### 2.10. Adelic Sobolev norms – a slight generalization.
The notations of this section will only be required in Sec. 7. We recommend it be omitted at a first reading.
In the discussion at the start of Sec. 2.9.3, we did not address what class of subgroups $`\mathrm{\Gamma }_i`$ to consider (should we take all finite index subgroups of a fixed $`\mathrm{\Gamma }_0`$ or some subclass?) Implicitly, such a choice was made in defining the Sobolev norms of the previous section. The Sobolev norms introduced in the previous section are good for most of our purposes. However, roughly speaking, they have the following defect: they only measure the index of a stabilizer of a function $`fC_\omega ^{\mathrm{}}(𝐗_𝐆)`$.
That this definition might lead to some peculiar results can be already seen in the case $`𝐆=\mathrm{PGL}(2)`$, $`F=`$. Let $`\chi _p`$ be the character of $`𝔸_{}^\times /^\times `$ that corresponds to the quadratic Dirichlet character of $``$ with conductor $`p`$, a prime number. Then the function $`g\chi _p(det(g))`$ descends to a function $`f`$ on $`𝐗`$, and it is easy to check that $`[K_{\mathrm{max}}:K_f]=2`$, for any $`p`$. Thus the index of this stabilizer does not reflect the conductor of the underlying representation (which, by any reasonable definition of conductor, should grow as $`p`$ increases). In this section we shall introduce a slight modification of the definitions which avoids this problem. (This problem would not occur for $`\mathrm{SL}(2)`$).
This is a purely technical matter, and it seems there is much scope for giving better and more natural definitions. We restrict ourselves to the case $`𝐆=\mathrm{GL}(n)`$. For a finite place $`𝔮`$ and $`m0`$, we put $`K[𝔮^m]\stackrel{\mathrm{def}}{=}\mathrm{ker}(\mathrm{GL}(n,𝔬_𝔮)\mathrm{GL}(n,𝔬_𝔮/\varpi _𝔮^m𝔬_{F_𝔮})`$, where $`\varpi _𝔮`$ is a uniformizer in $`F_𝔮`$. Now, for $`\psi C_\omega ^{\mathrm{}}(𝐗_𝐆)`$, put $`K_{𝔮,\psi }^{}`$ to be the largest subgroup $`K[𝔮^m]`$ which stabilizes $`\psi `$, and put $`K_\psi ^{}=_𝔮K_{𝔮,\psi }^{}`$,.
We define the $``$-pre-Sobolev norm $`PS_{p,d,\beta }^{}`$ by the rule
$$PS_{p,d,\beta }^{}(\psi )=[K_{\mathrm{max},𝐆}:K_\psi ^{}]^\beta \underset{\mathrm{ord}(𝒟)d}{}𝒟\psi _{L^p(𝐗_{𝐆,\mathrm{ad}})}$$
and we define the $``$-Sobolev norm $`S_{p,d,\beta }^{}`$ to be the maximal seminorm dominated by $`PS_{p,d,\beta }^{}`$. Clearly $`S_{p,d,\beta }^{}S_{p,d,\beta }`$.
The eventual purpose of this is that $`S_{p,d,\beta }^{}`$ (unlike $`S_{p,d,\beta }`$) will never be “too small” on an automorphic representation whose conductor is large. This can be quantified, although we do not do so in the present document.
###### Remark 2.1.
Evidently the definitions of this section and the previous are not the only “sensible” way of defining a notion of adelic Sobolev norms. The results of this paper do not require any more sophisticated definition, although this would certainly be of help in optimizing the results.
However, it would be interesting to impose a system of Sobolev-type norms in a less ad hoc fashion. Moreover, it would be pleasant if the system of norms had nice interpolation properties (this often is very helpful for getting sharp results). For example it would be nice if as one varied $`\beta `$ one got a family of interpolation spaces.
We remark on a simple way of defining Hilbertian norms which seems (more) appropriate to the adelic context. Let $`K[𝔮^m]`$ be as above, and let $`E_{𝔮^m}`$ be the averaging projection onto the $`K[𝔮^m]`$-fixed vectors, i.e. $`E_{𝔮^m}(v)=_{kK[𝔮^m]}kv𝑑k`$, where the measure is the Haar probability measure. Then $`e_{𝔮^m}:=E_{𝔮^m}E_{𝔮^{m1}}`$ is a projection. If $`𝔣=_i𝔮_i^{m_i}`$ is an arbitrary integral ideal, put $`e_𝔣:=_ie_{𝔮_i^{m_i}}`$. Now put $`P(s)=_𝔣e_𝔣\mathrm{N}(𝔣)^s`$. Then $`f_{\mathrm{ord}(𝒟)d}𝒟P(s)f_{L^2}`$ defines a Hilbert norm which seems to have reasonably pleasant formal properties. In fact it is majorized (up to constants) by a norm of the type described above,.
J. Bernstein has a more canonical notion of norms on representation spaces of $`p`$-adic groups, and he has informed me that these norms have adelic analogues. I do not know the relation. The norms arising from his constructions are Hilbertian.
### 2.11. Some properties and uses of the Sobolev norms.
We briefly summarize certain results that will be used in the text. Detailed proofs are given in Sec. 8.
For general $`𝐆,\omega `$ we have:
(2.12) $`S_{p,d,\beta }(F_1F_2)_dS_{2p,d,\beta }(F_1)S_{2p,d,\beta }(F_2).`$
(2.13) $`S_{p,d,\beta }(gF)g_{\mathrm{}}^dg_f^\beta S_{p,d,\beta }(F).`$
(2.12), proved in Lem. 8.1, and (2.13), proved in Lem. 8.2, give some basic stability properties of Sobolev norms.
Now we specialize to some results for $`\mathrm{GL}(2)`$ and $`\mathrm{PGL}(2)`$. Let $`FC^{\mathrm{}}(𝐗\times 𝐗)`$, let $`𝔮`$ be a prime ideal of $`𝔬_F`$, and suppose $`F`$ is invariant by $`\mathrm{PGL}_2(𝔬_{F_𝔮})\times \mathrm{PGL}_2(𝔬_{F_𝔮})`$. Then:
(2.14)
$$\begin{array}{c}\left|_𝐗F(x,xa([𝔮]))𝑑x\underset{\chi ^2=1}{}\chi ([𝔮])_𝐗F(x,y)\chi (x)\chi (y)𝑑\mu _𝐗(x)𝑑\mu _𝐗(y)\right|\hfill \\ \hfill _ϵ\mathrm{N}(𝔮)^{\frac{2\alpha 1}{p}+ϵ}S_{p,d}(F).\end{array}$$
(2.14), proved in Lem. 9.8, quantifies Hecke equidistribution. To understand the relation, take $`F`$ to be a pure tensor: $`F(x,y)=f_1(x)f_2(y)`$. Then (2.14) in effect bounds the inner product $`T_𝔮f_1,f_2`$, where $`T_𝔮`$ is the Hecke operator corresponding to $`𝔮`$.
### 2.12. Cusp forms, $`L`$-functions and the analytic conductor.
As a general remark on notation – and a mild abuse of notation– by cuspidal representation we shall always mean unitary cuspidal representation. This is automatic for $`\mathrm{PGL}(2)`$ but not for $`\mathrm{GL}(2)`$.
#### 2.12.1. $`L`$-functions.
Let $`\pi =_v\pi _v`$ be an automorphic cuspidal representation of $`\mathrm{GL}(n)`$ over $`F`$. We denote by $`L_v(s,\pi _v)`$ the local $`L`$-factor of the representation $`\pi _v`$; when it causes no confusion, we will sometimes abbreviate this to $`L(s,\pi _v)`$.
We write $`L(s,\pi ):=_{v\mathrm{finite}}L_v(s,\pi _v)`$ for the (finite part of) the global $`L`$-function attached to $`\pi `$, and $`\mathrm{\Lambda }(s,\pi ):=_vL_v(s,\pi _v)`$ for the (completed) $`L`$-function attached to $`\pi `$.
#### 2.12.2. The analytic conductor of Iwaniec-Sarnak.
We recall the definition in the context where it will arise. Let $`\pi =\pi _v`$ be a cuspidal representation of $`\mathrm{GL}(n)`$ over $`F`$.
For each finite place $`v`$ we denote by $`\mathrm{Cond}_v(\pi )`$ the conductor, in the sense of Jacquet, Piatetski-Shapiro, and Shalika, of $`\pi _v`$; thus $`\mathrm{Cond}_v(\pi )=q_v^{m_v}`$, where $`m_v`$ is the smallest non-negative integer such that $`\pi _v`$ possesses a fixed vector under the subgroup of $`\mathrm{GL}_n(𝔬_{F_v})`$ consisting of matrices whose bottom row is congruent to $`(0,0,\mathrm{},0,1)`$ modulo $`\varpi _v^m`$.
For each infinite place $`v`$, let $`\mathrm{\Gamma }_v(s)=\pi ^{s/2}\mathrm{\Gamma }(s/2)`$ or $`(2\pi )^s\mathrm{\Gamma }(s)`$ according to whether $`v`$ is real or complex respectively, and put $`\mathrm{deg}(v)=[F_v:]`$. Let $`\mu _{j,v}`$ satisfy $`L(s,\pi _v)=\mathrm{\Gamma }_v(s+\mu _{j,v})`$, and put $`\mathrm{Cond}_v(\pi )=_v(1+|\mu _{j,v}|)^{\mathrm{deg}(v)}`$. We then put $`\mathrm{Cond}(\pi )=_v\mathrm{Cond}_v(\pi )`$ (this is within a constant factor of the Iwaniec-Sarnak definition). Moreover, we put $`\mathrm{Cond}_{\mathrm{}}(\pi )=_{v\mathrm{infinite}}\mathrm{Cond}_v(\pi )`$ and $`\mathrm{Cond}_f(\pi )=_{v\mathrm{finite}}\mathrm{Cond}_v(\pi )`$ (the “infinite” and “finite” parts of the conductor).
We will occasionally refer to the “finite conductor” of $`\pi `$ as the ideal $`_v𝔮_v^{m_v}`$, where $`𝔮_v`$ is the prime ideal corresponding to the finite place $`v`$; then $`\mathrm{Cond}_f(\pi )`$ is the norm of this ideal. Hopefully the distinction between the two usages will be clear from context.
###### Remark 2.2.
(Explication for $`\mathrm{GL}(1)`$ in the archimedean case) Let us be slightly more explicit in the case of a unitary character $`\omega `$ of $`𝔸_F^\times /F^\times `$. If $`v`$ is real, then there is $`t`$ such that $`\omega (x)=|x|^{it}`$ for $`x>0`$; then $`\mathrm{Cond}_v(\omega )(1+|t|)`$. If $`v`$ is complex, then there is $`t,N`$ such that $`\omega (re^{i\theta })=|r|^{it}e^{iN\theta }`$; then $`\mathrm{Cond}_v(\omega )(1+|t|+N)^2`$.
We can heuristically summarize this: in the real case, $`\omega `$ is approximately constant in a neighbourhood of the identity of size $`\mathrm{Cond}_v(\omega )^1`$; in the complex, case $`\omega `$ is approximately constant in a disc around the identity of area $`\mathrm{Cond}_v(\omega )^1`$.
#### 2.12.3. Cusp forms
If $`\pi `$ is a cuspidal representation of $`\mathrm{GL}_2(𝔸_F)`$ or $`\mathrm{PGL}_2(𝔸_F)`$, it will be convenient to denote by $`\pi _{\mathrm{}}`$ the archimedean representation (of $`\mathrm{GL}_2(F_{\mathrm{}})`$ or $`\mathrm{PGL}_2(F_{\mathrm{}})`$) that corresponds to $`\pi `$.
By the dual $`\widehat{\mathrm{GL}_2(F_{\mathrm{}})}`$ or $`\widehat{\mathrm{PGL}_2}(F_{\mathrm{}})`$, we shall mean the space of irreducible, admissible representations. We say a subset of this dual is bounded if the corresponding set of Langlands parameters is bounded. We may define, in an evident way, the conductor $`\mathrm{Cond}(\pi _{\mathrm{}})`$ for $`\pi _{\mathrm{}}\widehat{\mathrm{GL}_2(F_{\mathrm{}})}`$; with this definition, a subset is bounded exactly when $`\mathrm{Cond}`$ takes bounded values on it.
In a similar fashion, we define the notion of a bounded subset of $`\widehat{GL_2(F_v)}`$ or $`\widehat{PGL_2(F_v)}`$ for any place $`v`$, where, again $`\widehat{\mathrm{GL}_2(F_v)}`$ denotes the set of irreducible, admissible representations.
## 3. Unipotent periods.
In this section, we will make systematic use of the (nonadelic) Sobolev norms $`S_{\mathrm{},d}`$ on homogeneous spaces $`\mathrm{\Gamma }\backslash G`$. In rough terms, $`S_{\mathrm{},d}`$ controls the $`L^{\mathrm{}}`$ norm of the first $`d`$ derivatives. See Sec. 2.9.2. Note in particular that $`S_{\mathrm{},0}`$ is just the $`L^{\mathrm{}}`$ norm.
### 3.1. Equidistribution of sparse subsets of horocycles.
Let $`\mathrm{\Gamma }\mathrm{SL}_2()`$ be a cocompact lattice. For this section alone, we will use mildly different notation to that of Sec. 2.5, to accommodate the fact we deal with $`\mathrm{SL}_2`$ and not with $`\mathrm{PGL}_2`$. For $`x`$, put
(3.1)
$$n(x)=\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right),a(x)=\left(\begin{array}{cc}x^{1/2}& 0\\ 0& x^{1/2}\end{array}\right),\overline{n}(x)=\left(\begin{array}{cc}1& 0\\ x& 1\end{array}\right).$$
We denote by $`C(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$ (resp. $`C^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$) the space of continuous (resp. smooth) functions on the compact real manifold $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$. We denote by $`dg`$ the measure on $`\mathrm{SL}_2()`$ that descends to a probability measure on the quotient $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$. Finally, we denote by $``$ the usual upper half plane $`\{z:\mathrm{}(z)>0\}`$ with the standard action of $`\mathrm{SL}_2()`$.
###### Theorem 3.1.
There exists $`\gamma _{\mathrm{max}}>0`$, depending on $`\mathrm{\Gamma }`$, such that $`\{x_0n(j^{1+\gamma }):j\}`$ is equidistributed, for any $`x_0\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$ and any $`0\gamma <\gamma _{\mathrm{max}}`$. In other words, for any $`fC(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$,
$$\underset{N\mathrm{}}{lim}\frac{_{j=1}^Nf(x_0n(j^{1+\gamma }))}{N}=_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g.$$
If $`\lambda _1`$ is the smallest nonzero eigenvalue of the Laplacian on $`\mathrm{\Gamma }\backslash `$, put $`\alpha =\{\begin{array}{cc}0,\hfill & \lambda _11/4\hfill \\ \sqrt{1/4\lambda _1},\hfill & \text{ else.}\hfill \end{array}.`$ Then we can take $`\gamma _{\mathrm{max}}=\frac{(12\alpha )^2}{16(32\alpha )}`$.
This result represents (extremely modest) progress towards a conjecture of N.Shah, which asserts that the statement should remain valid for any $`\gamma >0`$. The method is not restricted to sequences of the specific type in Thm. 3.1, and we have also not optimized the maximal value for $`\gamma _{\mathrm{max}}`$. Nevertheless the method is fundamentally limited. As it presently stands, it does not seem capable of achieving even $`\gamma =1`$. See also Remark 3.1 (page 3.1).
The dependence of $`\gamma _{\mathrm{max}}`$ on $`\mathrm{\Gamma }`$ can likely be removed, but this seems to require using further input (cf. last paragraph of Section 1.3.4).
The proof follows the line of Sec. 1.3.1, with $`G_1=\mathrm{SL}_2()`$, $`G_2=\{n(x):x\}`$. The $`Y_i`$ are not quite closed $`G_2`$-orbits, but rather long pieces of general $`G_2`$-orbits. The basis $`\{\psi _{i,j}\}`$ for $`Y_i`$ will correspond to additive characters of $`G_2`$.
Let $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$, and let $`\alpha `$ be as in the statement of Thm. 3.1. Let $`T1`$. Let $`\psi `$ be a fixed nontrivial character of the additive group of $``$. Let $`g`$ be a fixed smooth function of compact support on $``$ satisfying $`_{\mathrm{}}^{\mathrm{}}g(x)𝑑x=1`$. We denote by $`,_{L^2(\mathrm{\Gamma }\backslash G)}`$ the inner product in the Hilbert space $`L^2(\mathrm{\Gamma }\backslash G)`$.
We set:
(3.2)
$$\nu _T(f)=\frac{1}{T}_0^Tf(x_0n(t))𝑑t,\mu _{T,\psi }(f)=\frac{1}{T}_0^T\psi (t)f(x_0n(t))𝑑t.$$
Remark first that the measures $`\nu _T`$ are equidistributed as $`T\mathrm{}`$, in the following quantitative sense: for $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$,
(3.3)
$$\left|\nu _T(f)_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g\right|T^{\kappa _1}S_{\mathrm{},1}(f),$$
for any $`\kappa _1<\frac{1/2\alpha }{2}`$. This is proven in Lem. 9.4, without taking any pains to optimize the exponent. (We prove it to keep the paper self-contained. However, we emphasize that neither result nor proof is new; see and . A precise analysis of the equidistribution of long horocycles is carried out in .)
###### Lemma 3.1.
Suppose $`_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g=0`$. Then:
(3.4)
$$|\mu _{T,\psi }(f)|T^bS_{\mathrm{},1}(f),$$
whenever $`b<\frac{(12\alpha )^2}{8(32\alpha )}`$ and the implicit constant is independent of $`\psi `$.
Remark that if $`\psi `$ is wildly oscillatory, cancellation in $`\mu _{T,\psi }`$ can be proved directly by integration by parts; on the other hand, if $`\psi `$ is almost constant, the cancellation in $`\mu _{T,\psi }`$ arises from the equidistribution of the horocycle $`x_0G_2`$. It is therefore the intermediate case in which (3.4) is of interest.
###### Proof.
Let $`H1`$, and let $`\sigma _H`$ be the measure on $`N()`$ defined by $`\sigma _H(g)=\frac{1}{H}_0^H\psi (x)g(n(x))𝑑x`$, for $`g`$ a function on $`N()`$.
For $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$, we denote by $`f\sigma _H`$ the right convolution of $`f`$ by $`\sigma _H`$. Then it is easy to verify that $`\left|\mu _{T,\psi }(f)\mu _{T,\psi }(f\sigma _H)\right|\frac{H}{T}S_{\mathrm{},0}(f).`$ Here $`\sigma _H`$ denotes right convolution by $`\sigma _H`$. On the other hand, by Cauchy-Schwarz, $`|\mu _{T,\psi }(f\sigma _H)|^2\nu _T(|f\sigma _H|^2)`$. Thus, expanding $`|f\sigma _H|^2`$ and applying (3.3), we conclude
(3.5) $`|\mu _{T,\psi }(f)|{\displaystyle \frac{H}{T}}S_{\mathrm{},0}(f)+\left({\displaystyle \frac{1}{H^2}}{\displaystyle _{(h_1,h_2)[0,H]^2}}\left|\nu _T(n(h_1)f\overline{n(h_2)f})\right|𝑑h_1𝑑h_2\right)^{1/2}`$
$`{\displaystyle \frac{H}{T}}S_{\mathrm{},0}(f)+\left({\displaystyle \frac{1}{H^2}}{\displaystyle _{(h_1,h_2)[0,H]^2}}\left|n(h_1h_2)f,f_{L^2(\mathrm{\Gamma }\backslash G)}\right|𝑑h_1𝑑h_2\right)^{1/2}`$
$`+\left(T^{\kappa _1}\underset{(h_1,h_2)[0,H]^2}{sup}S_{\mathrm{},1}(n(h_1)f\overline{n(h_2)f})\right)^{1/2}`$
Utilising bounds towards matrix coefficients – see Sec. 9.1.2, esp. (9.6) – and basic properties of Sobolev norms (see <sup>7</sup><sup>7</sup>7Lem. 2.2 would actually give the exponent $`(1+|h_1|+|h_2|)^4`$ in the latter inequality, but it is easy to see directly the stronger result. Lem. 2.2), we note:
(3.6)
$$\begin{array}{c}n(h)f,f_ϵ(1+|h|)^{2\alpha 1+ϵ}S_{\mathrm{},1}(f)^2,\hfill \\ \hfill |S_{\mathrm{},1}(n(h_1)f\overline{n(h_2)f})|(1+|h_1|+|h_2|)^2S_{\mathrm{},1}(f)^2.\end{array}$$
Thus, $`|\mu _{T,\psi }(f)|_ϵ\left(\frac{H}{T}+H^{\alpha 1/2+ϵ}+T^{\kappa _1/2}H\right)S_{\mathrm{},1}(f).`$ Choose $`H`$ so that $`H^{\alpha 1/2}=HT^{\kappa _1/2}`$ to obtain the claimed result. ∎
###### Proof.
(of Thm. 3.1). Given Lem. 3.1, the Theorem follows quite readily by Fourier-expanding the measure on $``$ that is a sum of point masses at $`j^\gamma `$, for $`j`$. The argument that follows formalizes a minor variant of this argument (we first consider instead a sum of point masses along arithmetic progressions which approximate $`\{j^\gamma :j\}`$).
Let $`x_0\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$, $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$. We first claim that, if $`b`$ is as in the previous Lemma, $`f`$ so that $`_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g=0`$, and $`K1`$, then:
(3.7)
$$\frac{_{0j<K^{1/b1}}f(x_0n(Kj))}{K^{1/b1}}0,$$
as $`K\mathrm{}`$. In other words, $`K^{1/b1}`$ points, distributed along a horocycle with spacing $`K`$, become equidistributed.
This follows from Lem. 3.1: put $`g_\delta (x)=\mathrm{max}(\delta ^2(\delta |x|),0)`$, a function on $``$. For $`\lambda `$, write $`a_\lambda =K^1_{}\mathrm{exp}(2\pi iK^1\lambda t)g_\delta (t)𝑑t`$. Then $`_jg_\delta (t+Kj)=_k\mathrm{exp}(2\pi iK^1kt)a_k`$. Moreover, a simple computation shows that $`_k|a_k|\delta ^1`$.
Choose $`\epsilon >0`$ so that $`b+\epsilon `$ still satisfies the inequality of Lem. 3.1. By Lem. 3.1,
(3.8)
$$\left|_{t=0}^T𝑑t\left(\underset{j}{}g_\delta (t+Kj)\right)f(x_0n(t))\right|_fT^{1b\epsilon }\underset{k}{}|a_k|T^{1b\epsilon }\delta ^1.$$
$`g_\delta `$ has integral $`1`$ and is supported in a $`\delta `$-neighbourhood of $`0`$; in particular, the left-hand side of (3.8) differs from $`_{j,0KjT}f(x_0n(Kj))`$ by an error that is $`_f(1+TK^1\delta )`$. Thus $`\left|_{j,0KjT}f(x_0n(Kj))\right|_f(1+TK^1\delta +T^{1b\epsilon }\delta ^1)`$ from which (3.7) readily follows.
We now deduce the Theorem from (3.7). Let $`T_0`$ be large. Then, for $`t`$ small, we have $`(T_0+t)^{1+\gamma }=T_0^{1+\gamma }+(1+\gamma )T_0^\gamma t+O(t^2T_0^{\gamma 1})`$. In particular, $`(T_0+t)^{1+\gamma }`$ is well-approximated by the linear function $`T_0^{1+\gamma }+(1+\gamma )T_0^\gamma t`$ in the range where $`|t|T_0^{\frac{1\gamma }{2}}`$.
The claim of Thm. 3.1 follows from (3.7) as long as $`\frac{1\gamma }{2\gamma }>1/b1`$; in particular, any $`\gamma <b/2`$ will do. ∎
###### Remark 3.1.
The applicability of this method is not, of course, restricted to sequences of the form $`t^{1.01},t`$. In certain (very artificial) settings, one can prove some effective instances of Ratner’s theorem in non-horospherical cases by the same technique. The question of proving such results, even in very special cases, was raised by Margulis in his talk in the workshop “Emerging applications of measure rigidity.”
For instance, let $`r>1`$ be an integer, and let $`\mathrm{\Gamma }`$ be a cocompact lattice in $`\mathrm{SL}_2()^r`$. For concreteness, let us suppose that $`\mathrm{\Gamma }`$ is a congruence lattice. Let $`𝐧(x_1,\mathrm{},x_r)=(n(x_1),n(x_2),\mathrm{},n(x_r))\mathrm{SL}_2()^r`$. Then it is well-known that one can give an effective version of the equidistribution of $`𝐧(^r)`$, since $`𝐧(^r)`$ is horospherical. Using the method described above, one can extend this slightly: let $`V^r`$ be a linear subspace. One may show that $`𝐧(V)`$-orbits are effectively equidistributed if $`\mathrm{dim}(V)/r`$ is sufficiently close to $`1`$. However, the small codimension condition is crucial for this method and it does not seem that it would extend beyond this case.
### 3.2. Fourier coefficients of automorphic forms.
In the notations of Sec. 1, if we take for $`Y_i`$ the closed orbits of a unipotent group, the resulting periods are so-called “Fourier coefficients of automorphic forms.” We shall give a general nontrivial bound in that context. (The word “nontrivial” must be interpreted with care; see the discussion at the end of this section.) Our methods are restricted to the case of horospherical unipotent subgroups.
We have made no effort to optimize the exponents of the results, nor even to state a result of maximal generality. In fact, one can considerably increase the scope of Thm. 3.2, since we deal in the present section only with closed orbits of horospherical subgroups, one can profitably apply spectral theory. We do not carry this out, instead using to give equidistribution statements in a fairly soft fashion.
Let $`G`$ be a connected semisimple real Lie group, $`\mathrm{\Gamma }G`$ a lattice, $`KG`$ the maximal compact subgroup, $`𝔤`$ the Lie algebra of $`G`$, and $`H𝔤`$ a semisimple element. Fix a norm $``$ on the real vector space $`𝔤`$. Let $`\mathrm{exp}:𝔤G`$ be the exponential map. Fix a Haar measure on $`G`$ so that $`\mathrm{\Gamma }\backslash G`$ has volume $`1`$. Let $`𝔲`$ be the sum of all negative root spaces for $`H`$ and let $`U=\mathrm{exp}(𝔲)G`$. Let $`x_0\mathrm{\Gamma }\backslash G`$ be so that $`x_0U`$ is compact. Let $`x_t=x_0\mathrm{exp}(tH)`$, and let $`\mathrm{\Delta }_t`$ be the stabilizer of $`x_t`$ in $`U`$. We denote by $`,_{L^2(\mathrm{\Gamma }\backslash G)}`$ the inner product in the Hilbert space $`L^2(\mathrm{\Gamma }\backslash G)`$.
We shall analyze periods of a fixed function along $`x_tU`$ as $`t`$ varies. The proofs follow Sec. 1.3.1 with $`G_1=G,G_2=U`$, $`Y_i=x_tU`$, and $`\psi _{i,j}`$ corresponding to characters of $`U`$.
Let $`T>0`$ and let $`\psi `$ be any character of $`U`$ trivial on $`\mathrm{\Delta }_T`$ (:= $`\mathrm{exp}(TH)\mathrm{\Delta }_0\mathrm{exp}(TH)`$). We define $`\nu _T,\mu _{T,\psi }`$ in a closely analogous fashion to (3.2):
(3.9)
$$\nu _T(f)=\frac{_{\mathrm{\Delta }_T\backslash U}f(x_Tu)𝑑u}{\mathrm{vol}(\mathrm{\Delta }_T\backslash U)},\mu _{T,\psi }(f)=\frac{_{\mathrm{\Delta }_T\backslash U}f(x_Tu)\psi (u)𝑑u}{\mathrm{vol}(\mathrm{\Delta }_T\backslash U)}.$$
Let $`f,gC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$. Let $`E𝔲`$ have unit length w.r.t. the fixed norm $``$ on $`𝔤`$. It is proven by Kleinbock-Margulis in – see also Lem. 9.5 and Lem. 9.6 – that there are $`\kappa _1,\kappa _2>0`$ such that:
(3.10) $`\left|\mathrm{exp}(sE)f,g_{L^2(\mathrm{\Gamma }\backslash G)}{\displaystyle _{\mathrm{\Gamma }\backslash G}}f{\displaystyle _{\mathrm{\Gamma }\backslash G}}g\right|(1+|s|)^{\kappa _1}S_{\mathrm{},dim(K)}(f)S_{\mathrm{},dim(K)}(g)`$
(3.11)
$$|\nu _T(f)_{\mathrm{\Gamma }\backslash G}f|e^{\kappa _2T}S_{\mathrm{},dim(K)}(f).$$
(3.10) and (3.11) assert, respectively, quantitative mixing of the flow generated by $`U`$ on $`\mathrm{\Gamma }\backslash G`$, and the equidistribution of the orbit $`x_TU`$ as $`T\mathrm{}`$. We may assume that $`\kappa _1<1`$ (since making $`\kappa _1`$ smaller does not change the truth of (3.10)). This will ease the notation in the proof of the Theorem.
###### Theorem 3.2.
There exists $`\kappa _3>0`$ such that, for any $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$ satisfying $`_{\mathrm{\Gamma }\backslash G}f=0`$, we have:
(3.12)
$$|\mu _{T,\psi }(f)|\mathrm{exp}(T\kappa _3)S_{\mathrm{},dim(K)}(f)$$
for all $`T0`$, and for all characters $`\psi `$ of $`U`$ trivial on $`U_T`$.
Indeed, if $`o`$ is the order of the polynomial map $`\mathrm{End}(𝔤)`$ defined by $`s\mathrm{Ad}(\mathrm{exp}(sH))`$, then any $`\kappa _3<\frac{\kappa _1\kappa _2}{2(2odim(K)+\kappa _1)}`$ is admissible, $`\kappa _{1,2}`$ being as in (3.10) and (3.11).
The relevance to this to “Fourier coefficients” in the classical sense may not be immediately clear; after the proof, we give the example of $`\mathrm{SL}_2()`$ to illustrate.
Also, observe that the estimate (3.12) is uniform in $`\psi `$. In fact, just as in Lem. 3.1, the case when $`\psi `$ is constant amounts to (3.11), whereas the case where $`\psi `$ is highly oscillatory could be handled by integration by parts. It is, again, the intermediate case where (3.12) has content.
###### Proof.
We first remark that, other than being in a slightly more general setting, the proof is almost exactly the same as the proof of Lem. 3.1.
The signed measure $`\mu _{T,\psi }`$ satisfies $`\mu _{T,\psi }(uf)=\overline{\psi (u)}\mu _{T,\psi }(f)`$, for $`uU`$.
Take $`E𝔲`$ of unit length w.r.t. the norm $``$ on $`𝔤`$. Fix $`H1`$. Let $`\sigma `$ be the measure on $`U`$ defined via the rule
$$\sigma (h)=\frac{1}{H}_{s=0}^H\psi (\mathrm{exp}(sE))h(\mathrm{exp}(sE))𝑑s,$$
for $`h`$ any continuous compactly supported function on $`U`$. Then $`\mu _{T,\psi }(f)=\mu _{T,\psi }(f\sigma )`$; thus:
(3.13)
$$\begin{array}{c}|\mu _{T,\psi }(f)|^2=|\mu _{T,\psi }(f\sigma )|^2\nu _T(|f\sigma |^2)\hfill \\ \hfill _{\mathrm{\Gamma }\backslash G}|f\sigma |^2+e^{\kappa _2T}S_{\mathrm{},dim(K)}(|f\sigma |^2)\\ \hfill _{\mathrm{\Gamma }\backslash G}|f\sigma |^2+H^{2odim(K)}\mathrm{exp}(\kappa _2T)S_{\mathrm{},dim(K)}(f)^2.\end{array}$$
where we have applied Cauchy-Schwarz followed by (3.11), noting that by Lem. 2.2, we have that $`S_{\mathrm{},dim(K)}(|f\sigma |^2)S_{\mathrm{},dim(K)}(f\sigma )^2H^{2odim(K)}S_{\mathrm{},dim(K)}(f)^2`$. Here $`o`$ is chosen as in the statement of the Theorem.
By (3.10), we see that
(3.14)
$$\begin{array}{c}_{\mathrm{\Gamma }\backslash G}|f\sigma |^2\left(\frac{1}{H}_0^H(1+|t|)^{\kappa _1}𝑑t\right)S_{\mathrm{},dim(K)}(f)^2\hfill \\ \hfill H^{\kappa _1}S_{\mathrm{},dim(K)}(f)^2.\end{array}$$
Indeed, this follows simply by expanding the leftmost expression. Thus
$$|\mu _{T,\psi }(f)|^2(H^{\kappa _1}+H^{2odim(K)}\mathrm{exp}(\kappa _2T))S_{\mathrm{},dim(K)}(f)^2.$$
We choose $`H`$ so that $`H^{2odim(K)+\kappa _1}=\mathrm{exp}(\kappa _2T)`$ to conclude. ∎
###### Remark 3.2.
We now explain, when we specialize $`G=\mathrm{SL}_2()`$, why this recovers Sarnak’s result , which was the first improvement of the Hecke bound for nonarithmetic groups. Take $`H=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)𝔰𝔩_2`$, so that $`U=\left(\begin{array}{cc}1& \\ 0& 1\end{array}\right)`$. Let $`\mathrm{\Gamma }G`$ be a nonuniform lattice so that $`\mathrm{\Gamma }U=\{\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right),n\}`$, and take $`x_0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. Let $`f(x+iy)=_{n0}a_n\sqrt{y}K_{i\nu }(2\pi ny)e^{2\pi inx}`$ be a Maass cusp form of eigenvalue $`1/4+\nu ^2`$ on $`\mathrm{\Gamma }\backslash `$, where $``$ denotes the upper half-plane; it lifts to a function on $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$, viz. $`gf(g.i)`$.
Then the Theorem implies (in concrete language) that there exists $`\delta >0`$ such that, for any $`y1`$ and any $`n`$,
(3.15)
$$_{0x1}f(x+iy)e(nx)𝑑xy^\delta .$$
Taking $`yn^1`$ in (3.15), one easily deduces that the Fourier coefficients $`a_n`$ satisfy the “nontrivial” bound $`|a_n|n^{1/2\delta }`$.
###### Remark 3.3.
The bound Thm. 3.2 is nontrivial in that it improves, as $`t\mathrm{}`$ on the trivial bound:
$$\left|\frac{_{\mathrm{\Delta }_t\backslash U}f(xu)\psi (u)𝑑u}{\mathrm{vol}(\mathrm{\Delta }_t\backslash U)}\right|_f1,$$
which follows from Cauchy-Schwarz.
However, if there is another interpretation for the Fourier coefficients, it is not always the case that Thm. 3.2 improves on “trivial” bounds arising from that interpretation.
For instance, Fourier coefficients of cusp forms on $`\mathrm{GL}_n`$ admit a spectral interpretation, that is to say, they are connected to the eigenvalues of Hecke operators. In that case, Thm. 3.2 does not give anything even approaching the bounds of Jacquet-Piatetski-Shalika. Another example is when $`G=\stackrel{~}{\mathrm{SL}}_2()`$, and one takes for $`f`$ the Shimura lift of a cusp form of integral weight. In that case the (absolute values of the squares of) square-free Fourier coefficients of $`f`$ are given by special values of a twisted $`L`$-function; but the estimate above does not even recover the convexity bound (in fact, the method as indicated cannot recover this bound, even under optimal assumptions.)
It seems as though, in these cases, there is extra cancellation in the unipotent integrals for subtle arithmetic reasons. The crude methods indicated above do not detect this.
###### Remark 3.4.
We remark that, in the proof just given, the constant $`\kappa _3`$ depends on the spectral gap for $`\mathrm{\Gamma }\backslash G`$. This dependence can very likely be removed in many cases, including the case of $`G=\mathrm{SL}_2()`$, but we do not carry this out; again, cf. the last paragraph of Section 1.3.4. In the higher rank case, if $`G`$ has property (T), one has in any case a uniform spectral gap and this point becomes irrelevant.
It seems worthwhile to remark that, whereas the proof above is clearly not unrelated to that of Sarnak, it does not require any information on the decay of triple products (in particular, the deep “exponential decay” results proved by Sarnak and Bernstein-Reznikov).
We also remark that the proof indicated above, although it can be optimized in various ways, probably does not lead to as good an exponent as the work of Good, and the later refinement of Sarnak’s result due to Bernstein-Reznikov. Its advantage lies, rather, in its robustness and general applicability.
## 4. Semisimple periods: triple products in the level aspect.
In this section, we will give bounds for the triple product period on $`\mathrm{PGL}_2`$ over a number field $`F`$. We will use the notation of Sec. 2; in particular, $`𝔸_F`$ is the adele ring of $`F`$ and and $`𝐗=\mathrm{PGL}_2(F)\backslash \mathrm{PGL}_2(𝔸_F)`$.
In Sec. 4.1, we will give a special case of the triple product bound (Prop. 4.1) which does not require Sobolev norms to state. For some applications we will require a slight generalization, which will require the Sobolev norms of Sec. 2.9.3. This will be carried out in Sec. 4.2 (see Prop. 4.2).
### 4.1. Period bound for triple products.
We now give a period bound for triple products on $`\mathrm{PGL}_2`$. Although it is unfortunately somewhat disguised in the adelic language, the situation and method corresponds to that of Sec. 1.3.1 with $`G_1=\mathrm{PGL}_2(F_S)\times \mathrm{PGL}_2(F_S)`$, $`G_2=\mathrm{PGL}_2(F_S)`$ embedded diagonally. Here $`S`$ is a set of places of $`F`$ containing all infinite places, and $`F_S=_{vS}F_v`$.
For $`1p\mathrm{}`$ we will write $`L^p`$ for $`L^p(𝐗)`$. Thus, e.g., $`f_1_{L^4}`$ denotes $`\left(_𝐗|f_1(x)|^4𝑑x\right)^{1/4}`$.
###### Proposition 4.1.
(“Subconvexity for the triple product period.”) Let $`\pi `$ be an automorphic cuspidal representation of $`\mathrm{PGL}_2(𝔸_F)`$ with prime finite conductor $`𝔭`$. Let $`f_1,f_2C^{\mathrm{}}(𝐗)`$ be totally nondegenerate<sup>8</sup><sup>8</sup>8See Sec. 2.7; equivalent to “orthogonal to locally constant functions” in this case. and such that $`f_1,f_2`$ are $`\mathrm{PGL}_2(𝔬_{F_𝔭})`$-invariant. Let $`\phi \pi `$ and suppose<sup>9</sup><sup>9</sup>9Recall that a finite place $`v`$ belongs to the support of $`fC^{\mathrm{}}(𝐗)`$ exactly when $`\mathrm{PGL}_2(𝔬_v)`$ does not fix $`f`$. that there exists $`b`$ such that<sup>10</sup><sup>10</sup>10This assumption (4.1) is purely technical and the reader may safely assume that $`\phi `$ is spherical at all places away from $`𝔭`$ and $`f_1,f_2`$ are everywhere spherical without losing the gist of the argument. It is not used in the present document, but will probably be of use in establishing polynomial dependence of subconvex bounds. (4.1) ensures, among other things, that there are many places when all of $`f_1,f_2,\phi `$ are unramified, so that we can use the Hecke operators at those places.
(4.1)
$$\underset{𝔮\mathrm{Supp}(\phi )\mathrm{Supp}(f_1)\mathrm{Supp}(f_2)}{}\mathrm{N}(𝔮)\mathrm{N}(𝔭)^b$$
Put $`I(\phi )=_𝐗f_1(g)f_2(ga([𝔭]))\phi (g)𝑑g,`$ where $`a([𝔭])`$ is as in (2.4) and $`dg`$ is the $`\mathrm{PGL}_2(𝔸_F)`$-invariant probability measure. Then
(4.2)
$$|I(\phi )|_{b,ϵ,F}f_1_{L^4}f_2_{L^4}\phi _{L^2}\mathrm{N}(𝔭)^{ϵ\frac{(14\alpha )(12\alpha )}{4(34\alpha )}}$$
We refer to Prop. 4.1 as subconvexity for the triple product period, cf. first assertion of Theorem 5.1. We note that, with $`\alpha =3/26`$ (Kim’s bound) we have $`\frac{(14\alpha )(12\alpha )}{4(34\alpha )}>1/26`$. As we will see, Prop. 4.1 is a very strong result that implies many subconvexity results on $`\mathrm{PGL}(2)`$.
First let us explain the content of Prop. 4.1 in a classical setting, and how it can be regarded as the type of period bound discussed in the introduction. Suppose $`F=`$; let $`p1`$ and let
$$\mathrm{\Gamma }_0(p)=\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{GL}_2():p|c\},$$
and put $`Y(p)=\mathrm{\Gamma }_0(p)\backslash \mathrm{PGL}_2()`$. Then there is an embedding $`Y(p)Y(1)\times Y(1)`$ which corresponds to the graph of the $`p`$th Hecke correspondence on $`Y(1)`$; the image is a certain closed orbit of the diagonal $`\mathrm{PGL}_2()`$. Let $`f_1,f_2`$ be fixed functions on $`Y(1)`$ and $`\phi `$ a Maass form on $`Y_0(p)`$. Then the function $`f_1\times f_2:(x_1,x_2)f_1(x_1)f_2(x_2)`$ is a function on $`Y(1)\times Y(1)`$, and we can construct its restriction $`f_1\times f_2|_{Y(p)}`$ by means of the embedding indicated above. Then, translating from adelic to classical, one finds that (4.2) furnishes precisely an estimate for $`_{Y(p)}(f_1\times f_2)|_{Y(p)}\phi `$. When we vary $`p,\phi `$ and hold $`(f_1,f_2)`$ fixed, such an estimate falls precisely into the pattern described in the introduction: we are computing the periods of the fixed function $`f_1\times f_2`$ along the varying sequence of sets $`Y(p)`$. The fact that the $`Y(p)`$ become equidistributed in $`Y(1)\times Y(1)`$ is precisely equivalent to the equidistribution of $`p`$-Hecke orbits on $`Y(1)`$. Moreover, the key property of $`\phi `$ that is used is the fact that $`\phi `$ is an eigenfunction of many Hecke operators; this is used to construct the measure $`\sigma `$, in the notation of Sec. 1.3.
Prior to beginning the proof, we make some comments about applications and generalizations; for details, we refer to Sec. 5. The implicit constant of (4.2) is independent of $`f_1`$, $`f_2`$. Taking $`f_1,f_2`$ to be a pair of cusp forms Prop. 4.1 implies (conditional on some computations of $`p`$-adic integrals that we state as Hypothesis Prop. 11.1) subconvexity for certain triple product $`L`$-functions. Similarly, taking $`f_1,f_2`$ to be a cusp form and an Eisenstein series, resp. a pair of Eisenstein series, Prop. 4.1 implies subconvexity for Rankin-Selberg convolutions and standard $`L`$-functions.
The latter applications are rather delicate because Eisenstein series are not in $`L^4`$. To get around this we will eventually replace the Eisenstein series by an appropriate wave-packet (cf. proof of Thm. 5.1).
###### Proof.
Clearly we may assume that $`\phi _{L^2}=f_1_{L^4}=f_2_{L^4}=1`$. It follows that $`f_1_{L^2}1`$ and $`f_2_{L^2}1`$.
We shall moreover assume, for simplicity, that $`\phi `$ is spherical at all finite places $`v𝔭`$. The reader may verify that the proof carries through to the more general situation of Prop. 4.1 without modification.
Put $`q=\mathrm{N}(𝔭)`$. Let $`\sigma `$ be a (signed real) measure on $`\mathrm{PGL}_2(𝔸_F)`$ such that $`\phi \stackrel{ˇ}{\sigma }=\lambda \phi `$, for some $`\lambda `$. We shall assume that $`\mathrm{supp}(\sigma )`$ commutes with $`\mathrm{PGL}_2(_𝔭)`$; we will choose $`\sigma `$ later. Set further $`\mathrm{\Psi }(x)=f_1(x)f_2(xa([𝔭]))C^{\mathrm{}}(𝐗)`$. Then
(4.3) $`\lambda I(\phi )={\displaystyle _𝐗}\mathrm{\Psi }(x)(\phi \stackrel{ˇ}{\sigma })(x)𝑑x={\displaystyle _𝐗}(\mathrm{\Psi }\sigma )(x)\phi (x)𝑑x\left({\displaystyle _𝐗}|\mathrm{\Psi }\sigma |^2𝑑x\right)^{1/2}`$
$`=\left({\displaystyle _𝐗}{\displaystyle _{(g,g^{})\mathrm{PGL}_2(𝔸_F)^2}}\left(g\mathrm{\Psi }\right)\overline{\left(g^{}\mathrm{\Psi }\right)}𝑑\sigma (g)𝑑\sigma (g^{})𝑑x\right)^{1/2}`$
$`=\left({\displaystyle _𝐗}{\displaystyle _{(g,g^{})\mathrm{PGL}_2(𝔸_F)^2}}f_1(xg)f_2(xa([𝔭])g)\overline{f_1(xg^{})f_2(xa([𝔭])g^{})}𝑑\sigma (g)𝑑\sigma (g^{})𝑑x\right)^{1/2}`$
$`=\left({\displaystyle _𝐗}{\displaystyle _{(g,g^{})\mathrm{PGL}_2(𝔸_F)^2}}f_1(xg)f_2(xga([𝔭]))\overline{f_1(xg^{})f_2(xg^{}a([𝔭]))}𝑑\sigma (g)𝑑\sigma (g^{})𝑑x\right)^{1/2}`$
In the last step, we have used the fact that $`\mathrm{PGL}_2(_𝔭)`$, and thus $`a([𝔭])`$, commutes with $`\mathrm{supp}(\sigma )`$.
For any two functions $`h_1,h_2`$ on $`𝐗`$, both right invariant by $`\mathrm{PGL}_2(𝔬_{F_𝔭})`$, the assumed bound on Ramanujan (Def. 2.1) implies:
(4.4) $`\left|{\displaystyle _𝐗}h_1(x)h_2(xa([𝔭]))𝑑x{\displaystyle \underset{\chi ^2=1}{}}\chi (𝔭){\displaystyle _𝐗}h_1(x)\chi (x)𝑑x{\displaystyle _𝐗}h_2(x)\chi (x)𝑑x\right|`$
$`2q^{\alpha 1/2}h_1_{L^2}h_2_{L^2}`$
Here $`\chi `$ ranges over all characters of $`𝔸_F^\times /F^\times `$ such that $`\chi ^2=1`$, and $`\chi (x)`$ denotes the function $`g\chi (det(g))`$ on $`\mathrm{PGL}_2(F)\backslash \mathrm{PGL}_2(𝔸_F)`$. Indeed, to see (4.4), we note that the quantity inside the absolute value on the left hand side of (4.4) equals $`h_1𝒫h_1,a([𝔭])(h_2𝒫h_2)_{L^2}`$, where $`𝒫`$ is as in Sec. 2.7 (in this case, the orthogonal projection onto the locally constant functions). But the $`L^2`$ orthogonal projection $`\mathrm{Id}𝒫`$ kills all one-dimensional $`\mathrm{PGL}_2(F_𝔭)`$ representations, from which the result follows easily.
The functions $`h_j(x)=f_j(xg)\overline{f_j(xg^{})}(j=1,2)`$ are $`\mathrm{PGL}_2(𝔬_{F_𝔭})`$-invariant for $`g,g^{}\mathrm{supp}(\sigma )`$ since $`\mathrm{supp}(\sigma )`$ commutes with $`\mathrm{PGL}_2(F_𝔭)`$ and $`𝔭\mathrm{Supp}(f_1)\mathrm{Supp}(f_2)`$. Moreover, $`h_j_{L^2}f_j_{L^4}^2`$. Apply (4.4) to these $`h_j`$, and substitute in (4.3). It results:
(4.5)
$$\begin{array}{c}\left|\lambda I(\phi )\right|^2q^{\alpha 1/2}\sigma ^2\hfill \\ \hfill +\underset{\chi ^2=1}{}_{(g,g^{})}|g^1g^{}f_1,f_1\chi |g^1g^{}f_2,f_2\chi d|\sigma |(g)d|\sigma |(g^{})\end{array}$$
where $`|\sigma |`$ is the total variation measure associated to $`\sigma `$, $`\sigma =|\sigma |(𝐗)`$ is the total variation of $`\sigma `$, $`f_i\chi `$ is the function $`xf_i(x)\chi (det(x))`$, and brackets $`,`$ denote inner product in the Hilbert space $`L^2(𝐗)`$; we will suppress the reference to $`L^2(𝐗)`$ here and in the rest of the argument.
Put $`\sigma ^{(2)}=\stackrel{ˇ}{|\sigma |}|\sigma |`$. We may rewrite the previous result as:
(4.6)
$$\begin{array}{c}|\lambda |^2|I(\phi )|^2q^{\alpha 1/2}\sigma ^2\hfill \\ \hfill +\left(_g\underset{\chi ^2=1}{}|gf_1,f_1\chi ||gf_2,f_2\chi |d\sigma ^{(2)}(g)\right)\end{array}$$
We shall take $`\sigma `$ in (4.6) to be a linear combination of Hecke operators. We follow the notations introduced in Sec. 2.8. For $`𝔫\mathrm{Supp}(\phi )`$, we denote by $`\lambda (𝔫)`$ be the $`𝔫`$th Hecke eigenvalue of $`\phi `$, i.e. $`\phi \mu _𝔫=\lambda (𝔫)\phi `$. With our normalizations, the Ramanujan conjecture amounts to $`|\lambda (𝔩)|2`$ for $`𝔩`$ prime.
Let $`a_𝔫`$ be a sequence of complex numbers indexed by integral ideals of $`𝔬_F`$. Assume moreover that $`a_𝔫=0`$ whenever $`𝔫`$ is divisible by any place in $`\mathrm{Supp}(\phi )\mathrm{Supp}(f_1)\mathrm{Supp}(f_2)`$. Let $`\sigma `$ be the measure on $`\mathrm{PGL}_2(𝔸_{F,f})`$ defined by $`_𝔫a_𝔫\mu _𝔫`$. Then $`\sigma `$ is symmetric under $`gg^1`$, and $`|\sigma |=_𝔫|a_𝔫|\mu _𝔫`$. Moreover, $`\phi \sigma =\lambda \phi _3`$, where $`\lambda =_na_𝔫\lambda (𝔫)`$. From the assumed bound on Ramanujan (see Sec. 9.1, esp. equation (9.1))<sup>11</sup><sup>11</sup>11 A small caution here is that the vectors $`f_1\chi `$ and $`f_2\chi `$ need not be invariant by $`\mathrm{GL}_2(𝔬_{F_𝔭})`$, if $`𝔭|𝔫`$, because $`\chi `$ may be ramified. However, $`\mathrm{GL}_2(𝔬_{F_𝔭})`$ always fixes the line spanned by either of these vectors, and the bound of (9.1) depends only on the dimension of the $`\mathrm{GL}_2(𝔬_{F_𝔭})`$-span of the vectors in question., an elementary computation shows:
(4.7)
$$\left|_{g\mathrm{PGL}_2(𝔸_F)}\right|gf_1,f_1\chi gf_2,f_2\chi \left|d\mu _𝔫(g)\right|_ϵ\mathrm{N}(𝔫)^{2\alpha 1/2+ϵ}$$
Moreover, for fixed $`g\mathrm{Supp}(\mu _𝔫)`$, the inner product $`gf_1,f_1\chi `$ is nonvanishing only if $`\chi `$ is unramified at those places at all places not in $`\mathrm{Supp}(f_1)`$ and not dividing $`𝔫`$. The number of such quadratic characters is $`O_ϵ(\mathrm{N}(𝔫)^ϵq^ϵ)`$, where the implicit constant (as always) is allowed to depend on the base field $`F`$. Thus:
(4.8)
$$\left|\underset{\chi ^2=1}{}_g\right|gf_1,f_1\chi gf_2,f_2\chi \left|d\mu _𝔫(g)\right|_ϵq^ϵ\mathrm{N}(𝔫)^{2\alpha 1/2+ϵ}$$
The total variation of $`\sigma `$ may be computed:
(4.9)
$$\sigma _ϵ\underset{𝔫}{}\mathrm{N}(𝔫)^{1/2+ϵ}|a_𝔫|$$
Using (2.7) to compute $`\sigma ^{(2)}`$, and combining (4.6), (4.7) and (4.9), we conclude:
(4.10)
$$\left|I(\phi )\right|_ϵq^ϵ\frac{\left(\left(_𝔫\mathrm{N}(𝔫)^{1/2+ϵ}|a_𝔫|\right)^2q^{\alpha 1/2}+_{𝔫,𝔪}_{𝔡|(𝔫,𝔪)}\left(\mathrm{N}\left(\frac{𝔫𝔪}{𝔡^2}\right)\right)^{2\alpha 1/2}|a_𝔫||a_𝔪|\right)^{1/2}}{\left|_na_𝔫\lambda (𝔫)\right|}.$$
The choice of $`a_𝔫`$ follows an idea of Iwaniec; we slightly modify the standard choice so that we do not need to appeal to Ramanujan on average. <sup>12</sup><sup>12</sup>12The argument that follows was improved by a suggestion of P. Michel. Fix $`K`$ with $`q^{1/1000}Kq^{1000}`$. Let $`S`$ be the set of prime ideals $`𝔩`$ such that $`\mathrm{N}(𝔩)[K,2K]`$ and $`𝔩\mathrm{Supp}(f_1)\mathrm{Supp}(f_2)\mathrm{Supp}(\phi )`$. In view of the assumptions, $`|S|_ϵK^{1ϵ}`$.
For $`z`$ we put $`\mathrm{sign}(z)=z/|z|`$ for $`z0`$ and $`\mathrm{sign}(0)=1`$. Put
(4.11)
$$a_𝔫=\{\begin{array}{cc}\overline{\mathrm{sign}(\lambda (𝔫))},𝔫S\hfill & \\ \overline{\mathrm{sign}(\lambda (𝔫^2))},𝔫=𝔩^2,𝔩S\hfill & \\ 0,\text{ else}.\hfill & \end{array}$$
Then $`\left|_na_𝔫\lambda (𝔫)\right|_{ϵ,F}K^{1ϵ}`$, $`(_𝔫\mathrm{N}(𝔫)^{1/2+ϵ}|a_𝔫|)_ϵK^{2+ϵ}`$, and
(4.12)
$$\underset{𝔫,𝔪}{}\underset{𝔡|(𝔫,𝔪)}{}\left(\mathrm{N}\left(\frac{𝔫𝔪}{𝔡^2}\right)\right)^{2\alpha 1/2}|a_𝔫||a_𝔪|K^{4\alpha +1}.$$
We deduce from (4.10) that
(4.13)
$$|I(\phi )|_{ϵ,F}(qK)^ϵ\frac{\left(K^4q^{\alpha 1/2}+K^{1+4\alpha }\right)^{1/2}}{K}$$
Taking $`K=q^{\frac{1/2\alpha }{34\alpha }}`$, we obtain $`|I(\phi )|_{ϵ,F}q^{ϵ+\frac{(2\alpha 1/2)(1/2\alpha )}{34\alpha }}`$. ∎
### 4.2. A technical generalization.
For certain applications, we shall require a slight generalization of Prop. 4.1 in which the role of $`gf_1(g)f_2(ga([𝔭]))`$ is replaced by $`gF(g,ga([𝔭]))`$, where $`F`$ is a function on $`𝐗\times 𝐗`$ that is not necessarily of product type. Although the method of proof is identical to Prop. 4.1 the details are slightly more technical; in particular, to state the result we will have need of the adelic Sobolev norms discussed in Sec. 2.9.3. We shall also use the notion of totally nondegenerate for functions on $`𝐗\times 𝐗`$: see Sec. 2.4.
###### Proposition 4.2.
Suppose $`FC^{\mathrm{}}(𝐗\times 𝐗)`$ is totally nondegenerate. Suppose moreover that there is $`b`$ with
(4.14)
$$\underset{𝔮\mathrm{Supp}(F)\mathrm{Supp}(\phi )}{}\mathrm{N}(𝔮)\mathrm{N}(𝔭)^b$$
Let $`\pi `$ be a cuspidal representation of $`\mathrm{PGL}_2`$ over $`F`$, with conductor $`𝔭`$, and put $`I(\phi )=_𝐗F(x,xa([𝔭]))\phi (x)𝑑x`$, for $`\phi \pi `$. Then, for any $`p>4,d1`$,
$$|I(\phi )|_{b,ϵ}\mathrm{N}(𝔭)^{\beta +ϵ}\phi _{L^2}S_{p,d,2/p}(F),$$
where $`\beta =\frac{(12\alpha )(14\alpha )}{p(74\alpha )}`$.
With Kim’s bound $`\alpha =3/26`$, we obtain $`\frac{(12\alpha )(14\alpha )}{74\alpha }>1/17`$.
###### Proof.
The proof follows closely the proof of Prop. 4.1; the only difference is that we apply (2.14) (proved in Lem. 9.8) in place of (4.4). Again, we may freely assume that $`\phi _{L^2}=1`$; again we put $`q=\mathrm{N}(𝔭)`$.
Let notations be as in the proof of Prop. 4.1; in particular, $`\sigma `$ is a signed real measure on $`\mathrm{PGL}_2(𝔸_{F,f})`$ whose support commutes with $`\mathrm{PGL}_2(F_𝔭)`$, and $`\lambda `$ satisfies $`\phi \stackrel{ˇ}{\sigma }=\lambda \phi `$. Proceeding as in that proof, and in particular as in (4.3), we obtain:
(4.15)
$$\begin{array}{c}|\lambda I(\phi )|^2_𝐗_{(g,g^{})\mathrm{PGL}_2(𝔸_{F,f})^2}F((x,x)(g,g)(1,a([𝔭])))\hfill \\ \hfill \overline{F((x,x)(g^{},g^{})(1,a([𝔭])))}d\sigma (g)d\sigma (g^{})\end{array}$$
For any $`g,g^{}\mathrm{PGL}_2(𝔸_F)`$, set $`F_{g,g^{}}(x_1,x_2)=F((x_1,x_2)(g,g))\overline{F((x_1,x_2)(g^{},g^{}))}`$. Then $`F_{g,g^{}}`$ is invariant by $`\mathrm{PGL}_2(𝔬_{F_𝔭})\times \mathrm{PGL}_2(𝔬_{F_𝔭})`$ for $`g,g^{}\mathrm{supp}(\sigma )`$. By Hecke equidistribution in the form of (2.14) (proved in Lem. 9.8) we see that for $`p>2,d1`$:
(4.16)
$$\begin{array}{c}\left|_𝐗F_{g,g^{}}((x,x)(1,a([𝔭])))𝑑x\underset{\chi ^2=1}{}\chi (𝔭)_{𝐗\times 𝐗}F_{g,g^{}}(x_1,x_2)\chi (x_1)\chi (x_2)𝑑x_1𝑑x_2\right|\hfill \\ \hfill _ϵq^{(2\alpha 1)/p+ϵ}S_{p,d}(F_{g,g^{}}).\end{array}$$
Here, as before, $`\chi (x)`$ denotes the function on $`𝐗`$ defined by $`g\chi (det(g))`$. By definition $`|_{𝐗\times 𝐗}F_{g,g^{}}((x_1,x_2))\chi (x_1)\chi (x_2)|=|(g^1g,g^1g)F,F(\chi ,\chi )_{L^2(𝐗\times 𝐗)}|`$.
By the basic properties of adelic Sobolev norms ((2.12) and (2.13), proofs in Lem. 8.1 and Lem. 8.2)
(4.17) $`S_{p,d}(F_{g,g^{}}):=S_{p,d,1/p}(F_{g,g^{}})S_{2p,d,1/p}((g,g)F)S_{2p,d,1/p}((g^{},g^{})F)`$
$`g^{2/p}g^{}^{2/p}S_{2p,d,1/p}(F)^2.`$
Let us remark that the factors $`g^{2/p}`$ and $`g^{}^{2/p}`$ arises in the following way: Lem. 8.2 actually gives a factor $`(g,g)^{1/p}`$, where the norm $``$ (as in Sec. 2.4) is computed in $`\mathrm{PGL}_2(𝔸_{F,f})^2`$; this equals $`g^{1/p}`$ where the norm is computed in $`\mathrm{PGL}_2(𝔸_{F,f})`$.
Choose $`\sigma `$ as in the proof of Prop. 4.1 (see esp. paragraph before (4.7)) and choose the coefficients $`a_𝔫`$ as in that proof (see (4.11)). In particular, $`\sigma K^{2+ϵ}`$. Since $`F`$ is totally nondegenerate, the matrix coefficients $`(g^1g,g^1g)F,F`$ satisfy bounds that are of the same quality as in the proof of Prop. 4.1; in particular, as in (4.8):
$$\underset{\chi ^2=1}{}_g\left|(g,g)F,F(\chi ,\chi )\right|𝑑\mu _𝔫(g)q^ϵ\mathrm{N}(𝔫)^{2\alpha 1/2+ϵ}F^2$$
Finally $`g_ϵK^{2+ϵ}`$ for all $`g\mathrm{supp}(\sigma )`$. Proceeding just as in the previous proof,
$$|I(\phi )|_ϵ(qK)^ϵ\frac{\left(K^4q^{(2\alpha 1)/p}K^{8/p}S_{2p,d,1/p}(F)^2+K^{1+4\alpha }F_{L^2(𝐗\times 𝐗)}^2\right)^{1/2}}{K}.$$
Consequently, for any $`p>2`$,
$$|I(\phi )|_ϵ(qK)^ϵ\frac{\left(K^8q^{(2\alpha 1)/p}S_{2p,d,1/p}(F)^2+K^{1+4\alpha }F_{L^2}^2\right)^{1/2}}{K}$$
To conclude, choose $`K=q^{\frac{12\alpha }{p(74\alpha )}}`$ and replace $`p`$ by $`p/2`$ (thus, e.g., $`p>2`$ becomes $`p>4`$). ∎
## 5. Application to $`L`$-functions.
We now present the first applications to subconvexity. The rough idea is simply that certain $`L`$-functions are expressed as period integrals of the type that are bounded by Prop. 4.1 and Prop. 4.2. There is one significant issue in implementing this (rather evident) idea: namely, the integral representation that we use for Rankin-Selberg and the standard $`L`$-functions involve Eisenstein series, which are not in $`L^2`$; this causes problems in applying Prop. 4.1!
Thus we need to regularize. Two natural ways of doing this are to replace an Eisenstein series by a “wave-packet”; or to use a suitable form of truncation in the defining integrals. In the present paper we will use the wave-packet technique; in the paper we shall also use truncation.
Let us briefly describe the wave packet technique in a classical language. Roughly speaking we can express the Rankin-Selberg $`L`$-function of two classical forms $`f,g`$ via an integral of the form $`L(s,f\times g)=_zf(z)g(z)E(s,z)`$, for some Eisenstein series $`E(s)`$. We now regularize, replacing $`E(s,z)`$ by a wave packet. Let $`h(s)`$ be any holomorphic function: then
(5.1)
$$_{Re(s)=1/2}h(s)L(s,f\times g)𝑑s=_zf(z)g(z)_{\mathrm{}(s)=1/2}h(s)E(s,z).$$
We wish to eventually recover an upper bound for $`L(1/2,f\times g)`$ (say) from the left-hand side, so we take $`h(s)=\overline{L(1\overline{s},f\times g)}`$. Then $`h(s)L(s,f\times g)`$ is positive along $`\mathrm{}(s)=1/2`$. To apply Prop. 4.1 to the right-hand side of (5.1), we shall moreover need to control the behavior of the regularized Eisenstein series $`E_h=_{\mathrm{}(s)=1/2}h(s)E(s,z)`$; this type of analysis is carried out in Sec. 10.2 and Sec. 10.3, the main point being that the divergence of the Eisenstein series comes entirely from the constant term.
It is worth remarking that Iwaniec’s bounds for the $`L`$-function near $`1`$ enter rather crucially in this analysis: in effect, we bound $`E_h`$ by an easy argument involving shift of contours; this necessitates that $`h`$ be estimated on a line $`\mathrm{}(s)=\epsilon `$, which amounts to estimating $`L(s,f\times g)`$ for $`\mathrm{}(s)=1+\epsilon `$.
In what follows we have not attempted to obtain polynomial dependence in all parameters. This is not hard to do — and, at its essence, a statement that one can find analytically suitable test vectors in a Rankin-Selberg integral; but we have not done so here. On the other hand, we give full details of this procedure in the proof of Thm. 6.1 (in which the polynomial dependence is particularly useful for applications).
###### Theorem 5.1.
Let $`\pi _1,\pi _2`$ be fixed automorphic cuspidal representations of $`\mathrm{PGL}_2`$ over $`F`$; fix $`t`$. Let $`\pi `$ be an automorphic cuspidal representation with conductor $`𝔭`$, a prime ideal that is prime to the conductors of $`\pi _1`$ and $`\pi _2`$.
Then, assuming Hypothesis 11.1:
(5.2)
$$L(\frac{1}{2},\pi _1\pi _2\pi )_\pi _{\mathrm{}}\mathrm{N}(𝔭)^{1\frac{1}{13}}$$
and, unconditionally:
(5.3)
$$|L(\frac{1}{2}+it,\pi _1\pi )|^2_\pi _{\mathrm{}}\mathrm{N}(𝔭)^{1\frac{1}{100}}$$
(5.4)
$$|L(\frac{1}{2}+it,\pi )|^4_\pi _{\mathrm{}}\mathrm{N}(𝔭)^{1\frac{1}{600}}$$
In these statements, the notation $`_\pi _{\mathrm{}}`$ indicates an implicit constant that depends continuously on the local archimedean representation $`\pi _{\mathrm{}}`$ of $`\mathrm{GL}_2(F_{\mathrm{}})`$ underlying $`\pi `$.
Note we make no claim about the dependency of the implicit constant on $`t,\pi _1,\pi _2`$; as remarked above, this dependence could be made polynomial in the conductors, but this would require more careful analysis of the archimedean integrals. <sup>13</sup><sup>13</sup>13 It is important to note, however, that this is an entirely local problem; it is intended that this will be carried out in a more general context in . Both for applications and to illustrate procedure, we have carried out this type of analysis for the results on subconvexity of character twists in Section 6. Those results are proved with polynomial dependence on all parameters.
We remark that we have used H.Kim’s bound $`\alpha =3/26`$; any value of $`\alpha `$ less than $`1/4`$ would give subconvexity and under Ramanujan one obtains for (5.2) the exponent $`5/6`$. The exponents for (5.3) and (5.4) can be improved, e.g. the present proof does not take into account the fact that unitary Eisenstein series satisfy Ramanujan!
### 5.1. Results relating periods and integral representations.
For the convenience of the reader, we summarize here the results that relate periods and integral representations (proved in later sections). Roughly speaking, any integral representation for an $`L`$-function expresses it as a period integral with certain test vectors belonging to the space of an automorphic cuspidal representation.
A delicate point, which is quite relevant to issues of polynomial dependence in auxiliary parameters, is precisely which test vectors. In principle, the proofs of results about integral representation give explicit test vectors. In practice, it is tedious to extract these explicit test vectors. Our policy throughout this paper is the lazy one: to deduce results, as far as possible, by formal arguments and without choosing explicit test vectors. The price of this is that we will not obtain not quite the $`L`$-function, but rather a holomorphic function that differs from the $`L`$-function by some harmless factors.
More precisely, the content of the Proposition (Prop. 5.1) that follows is that one can write down an integral representation $`I(s)`$ for the $`L`$-functions of interest, so that:
1. $`I(1/2)`$ is not too much smaller than $`L(1/2)`$ – or with $`1/2`$ replaced by the point of interest – so that a bound for $`I(1/2)`$ gives a bound for $`L(1/2)`$.
2. $`I(s)`$ is not too much bigger than $`L(s)`$ for any $`s`$. This type of control will be useful in shifting contours.
One might prefer to get $`I(s)=L(s)`$ but we don’t need this stronger statement.
As is discussed at length in Sec. 10, to a Schwarz function $`\mathrm{\Psi }`$ on $`𝔸_F^2`$ is associated a family of Eisenstein series $`E_\mathrm{\Psi }(s,g)`$ on $`𝐗`$, which varies meromorphically in the parameter $`s`$.
###### Proposition 5.1.
Let $`s_0,t_0,t_0^{}`$. Let $`\pi _1`$ be a fixed automorphic cuspidal representation of $`\mathrm{PGL}_2(𝔸_F)`$ and $`\pi `$ an automorphic cuspidal representation of prime conductor $`𝔭`$; assume that the finite conductors of $`\pi ,\pi _1`$ are coprime.
Denoting by $`\pi _{\mathrm{}}`$ the representation of $`\mathrm{PGL}_2(F_{\mathrm{}})`$ corresponding to $`\pi `$, suppose that $`\mathrm{Cond}(\pi _{\mathrm{}})`$ is bounded above; equivalently, $`\pi _{\mathrm{}}`$ belongs to a bounded subset<sup>14</sup><sup>14</sup>14See Section 2.12.3 for definition of the dual $`\widehat{\mathrm{PGL}_2}(F_{\mathrm{}})`$ (in what follows the implicit constants may depend on these bounds).
There exists a fixed finite set $``$ of Schwarz Bruhat functions on $`𝔸_F^2`$ and a real number $`C>0`$ so that:
There exist vectors $`\phi _1\pi _1,\phi \pi `$ and $`\mathrm{\Psi }`$ so that
$$\mathrm{\Phi }(s):=\mathrm{N}(𝔭)^{1s}\frac{_𝐗\phi (g)\phi _1(ga([𝔭]))E_\mathrm{\Psi }(s,g)𝑑g}{\mathrm{\Lambda }(s,\pi _1\pi )}$$
is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(s_0)|1`$ and $`|\mathrm{\Phi }(s)|C^{|\mathrm{}(s)|}(1+|s|)^C`$;
2. At any nonarchimedean place $`v`$ such that $`\pi _1`$ is unramified, $`\mathrm{\Psi }_v`$ is invariant by $`\mathrm{PGL}_2(𝔬_{F_v})`$; at any nonarchimedean place $`v`$ such that $`\pi _1`$ and $`\pi `$ are both unramified, $`\phi `$, $`\phi _1`$ are both invariant by $`\mathrm{PGL}_2(𝔬_{F_v})`$.
3. $`\phi _1_L^{\mathrm{}}1`$ and $`\phi _{L^2(𝐗)}_ϵ\mathrm{N}(𝔭)^ϵ`$.
Moreover, there exist vectors $`\phi \pi ,\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ so that:
(5.5)
$$\mathrm{\Phi }(t,t^{})=\mathrm{N}(𝔭)^{1/2t}\frac{_𝐗\phi (g)E_{\mathrm{\Psi }_1}(g,\frac{1}{2}+t)E_{\mathrm{\Psi }_2}(ga([𝔭]),\frac{1}{2}+t^{})𝑑g}{\mathrm{\Lambda }(\frac{1}{2}+t+t^{},\pi )\mathrm{\Lambda }(\frac{1}{2}+tt^{},\pi )}$$
is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(t_0,t_0^{})|1`$ and $`|\mathrm{\Phi }(t,t^{})|C^{|\mathrm{}(t)|+C|\mathrm{}(t^{})|}(1+|t|+|t^{}|)^C`$.
2. For any nonarchimedean place $`v`$, each $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ is invariant by $`\mathrm{PGL}_2(𝔬_{F_v})`$; for each place at which $`\pi `$ is unramified, the same is true of $`\phi `$.
3. $`\phi _{L^2(𝐗)}_ϵ\mathrm{N}(𝔭)^ϵ`$.
###### Proof.
Lem. 11.4 and Lem. 11.5. ∎
In effect, we could achieve “$`\mathrm{\Phi }=1`$” in Prop. 5.1 by a more careful choice of local data; but this is irrelevant for the purpose of global estimation.
### 5.2. Proof of Thm. 5.1.
###### Proof.
(of (5.2)). It follows from Hypothesis 11.1 and Proposition 4.1.
###### Proof.
(of (5.3)). The basic idea is that the Rankin-Selberg convolution is a triple product, with one factor being an Eisenstein series. However, one cannot naively apply Prop. 4.1 since Eisenstein series do not belong to $`L^4(𝐗)`$. To avoid this, we will use a wave-packet of Eisenstein series.
First, we can assume from the start that $`\pi _{\mathrm{}}`$ belongs to a bounded subset of the dual $`\widehat{\mathrm{PGL}_2(F_{\mathrm{}})}`$. The implicit constants in the proof that follow depend on this subset. We denote by $`\mathrm{\Lambda }`$ the completed $`L`$-function. We begin by remarking that since $`\mathrm{N}(𝔭)\mathrm{}`$ we may assume that $`\pi _1`$ is not isomorphic to $`\pi `$, or to any quadratic twist of $`\pi `$. In particular, we are free to assume that $`\mathrm{\Lambda }(s,\pi _1\pi )`$ has no poles. Moreover, the finite conductor of $`\mathrm{\Lambda }(s,\pi _1\pi )`$ differs from $`\mathrm{N}(𝔭)^2`$ by an absolutely bounded constant.
Fixing $`t_0`$, let $`\mathrm{\Psi },\phi ,\phi _1,E_\mathrm{\Psi }(g,s),\mathrm{\Phi }`$ be as in Prop. 5.1 with $`s_0=1/2+it_0`$, so that $`|\mathrm{\Phi }(1/2+it_0)|1`$. For simplicity we write simply $`E(g,s)`$ for $`E_\mathrm{\Psi }(g,s)`$. Fix $`\kappa >0`$. In the rest of the proof we omit the subscript $`\kappa ,ϵ`$ from $``$, with the understanding that all implicit constants depend on $`\kappa `$ and $`ϵ`$. Put
$$I(s)=\mathrm{N}(𝔭)^{s1}\mathrm{\Lambda }(s,\pi _1\pi )\mathrm{\Phi }(s)=_𝐗\phi _1(ga([𝔭]))E(s,g)\phi (g)𝑑g$$
From Iwaniec’s upper bounds for $`L`$-functions \[15, Chapter 8\], the functional equation for $`\mathrm{\Lambda }`$, and the bounds on $`\mathrm{\Phi }`$ furnished by Prop. 5.1,
(5.6)
$$|I(1+\kappa +it)|(1+|t|)^4\mathrm{N}(𝔭)^{\kappa +ϵ},|I(\kappa +it)|(1+|t|)^4\mathrm{N}(𝔭)^{\kappa +ϵ}.$$
Put $`h(s)=s(1s)(s\frac{1}{2})^2\overline{I(1\overline{s})}`$. Then $`h(s)`$ is holomorphic in $`\kappa \mathrm{}(s)1+\kappa `$ and $`h(\frac{1}{2})=0`$.<sup>15</sup><sup>15</sup>15 The fact that we impose $`h(1/2)=0`$ has a very concrete meaning in classical terms. Fix, for example, a form $`f`$ and $`t`$. Consider $`_g|L(\frac{1}{2}+it,fg)|^2`$ where the sum is taken over a basis of holomorphic Hecke eigenforms of level $`N`$ and trivial Nebentypus. If $`t=0`$, this has the asymptotic behaviour $`N\mathrm{log}(N)^3`$. On the other hand, if $`t0`$, it behaves like $`N\mathrm{log}(N)`$. Forcing $`h(1/2)=0`$ “counteracts” this extra singularity. Moreover $`h(s)`$ has rapid decay as $`\mathrm{}(s)\mathrm{}`$, in view of the $`\mathrm{\Gamma }`$-factors of the completed $`L`$-function. Put $`E_h(g)=_{\mathrm{}(s)=1+\kappa }h(s)E(s,g)`$. It is proved in Lem. 10.6 that, for such $`h`$, $`E_h(g)_L^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\left(|h(\kappa +it)|+|h(1+\kappa +it)|\right)𝑑t.`$ Applying (5.6), we conclude that $`E_h(g)_L^{\mathrm{}}\mathrm{N}(𝔭)^{\kappa +ϵ}`$.
On the other hand, we see from the definition of $`I(s)`$ that
(5.7)
$$_{\mathrm{}(s)=1+\kappa }h(s)I(s)=_{\mathrm{}(s)=1+\kappa }h(s)𝑑s_𝐗\phi _1(ga([𝔭]))E(s,g)\phi (g)𝑑g.$$
The double integral on the right hand side of (5.7) is absolutely convergent and orders may be switched; thus
$$_{\mathrm{}(s)=1+\kappa }h(s)I(s)=_𝐗\phi _1(ga([𝔭]))E_h(g)\phi (g)𝑑g=_𝐗\phi _1(ga([𝔭]))E_h^0(g)\phi (g)𝑑g,$$
where $`E_h^0:=𝒫(E_h)`$ is totally nondegenerate (see Section 2.7) and satisfies $`E_h^0_L^{\mathrm{}}_ϵ\mathrm{N}(𝔭)^{\kappa +ϵ}`$.
We then deduce from Prop. 4.1 that:
(5.8)
$$\left|_{\mathrm{}(s)=1+\kappa }h(s)I(s)\right|\phi _1_{L^4(𝐗)}\mathrm{N}(𝔭)^{\frac{1}{26}+\kappa +ϵ}$$
$`I(s)`$ and $`h(s)`$ both decay exponentially rapidly as $`\mathrm{}(s)\mathrm{}`$. It is therefore simple to justify shifting the line of integration in (5.8) to $`\mathrm{}(s)=1/2`$. We deduce thereby that
(5.9)
$$\left|_{\mathrm{}(s)=\frac{1}{2}}t^2|I(\frac{1}{2}+it)|^2𝑑t\right|\phi _1_{L^4(𝐗)}\mathrm{N}(𝔭)^{\kappa \frac{1}{26}+ϵ}.$$
From (5.6) we deduce bounds on $`I`$ and $`I^{}`$ inside the strip $`0\mathrm{}(s)1`$ by the maximal modulus principle. In particular,
(5.10)
$$|I^{}(\frac{1}{2}+it)|_t\mathrm{N}(𝔭)^{\kappa +ϵ}.$$
Combining (5.9) and (5.10), and recalling that $`\kappa `$ is arbitrary, we obtain $`|I(1/2+it)|_t\mathrm{N}(𝔭)^{\frac{1}{526}+ϵ}`$. Thus $`|\mathrm{\Lambda }(\frac{1}{2}+it_0,\pi _1\pi )|_{ϵ,t}\mathrm{N}(𝔭)^{\frac{1}{2}\frac{1}{130}+ϵ}`$. ∎
###### Proof.
(of (5.4).) The proof is similar to that of (5.3), but a slightly more elaborate regularization is required, since we shall proceed from the expression (5.5) of $`L(s,\pi )`$ as a triple product against two Eisenstein series. Again we may assume from the start that $`\pi _{\mathrm{}}`$ is confined to a bounded subset of $`\widehat{\mathrm{GL}_2(F_{\mathrm{}})}`$; the implicit constants will, again, depend on this subset.
Let $`\mathrm{\Lambda }(s,\pi )`$ be the completed $`L`$-function attached to $`\pi `$. Fixing $`t_0,t_0^{}i`$, Prop. 5.1 gives the existence of Eisenstein series $`E_{\mathrm{\Psi }_1}(g,s)=E_1(g,s),E_{\mathrm{\Psi }_2}(g,s)=E_2(g,s)`$ on $`𝐗`$, and $`\phi \pi `$ so that
$$\mathrm{\Phi }(t,t^{}):=\mathrm{N}(𝔭)^{1/2t}\frac{_𝐗\phi (g)E_1(g,\frac{1}{2}+t)E_2(ga([𝔭]),\frac{1}{2}+t^{})𝑑g}{\mathrm{\Lambda }(\frac{1}{2}+t+t^{},\pi )\mathrm{\Lambda }(\frac{1}{2}+tt^{},\pi )}$$
satisfies $`|\mathrm{\Phi }(t_0,t_0^{})|1`$ and $`\mathrm{\Phi }(t,t^{})C^{|\mathrm{}(t)|+|\mathrm{}(t^{})|}(1+|t|+|t|^{}|)^C`$.
We put
(5.11) $`I(z_1,z_2)=\mathrm{\Phi }(z_1,z_2)\mathrm{N}(𝔭)^{z_11/2}\mathrm{\Lambda }({\displaystyle \frac{1}{2}}+z_1+z_2,\pi )\mathrm{\Lambda }({\displaystyle \frac{1}{2}}+z_1z_2)`$
$`=\mathrm{\Phi }(z_1,z_2)\mathrm{N}(𝔭)^{\frac{z_1+z_2}{2}1/4}\mathrm{\Lambda }({\displaystyle \frac{1}{2}}+z_1+z_2,\pi )\mathrm{N}(𝔭)^{\frac{z_1z_2}{2}1/4}\mathrm{\Lambda }({\displaystyle \frac{1}{2}}+z_1z_2,\pi )`$
Then $`I(z_1,z_2)`$ is a holomorphic function of $`(z_1,z_2)^2`$. $`I(z_1,z_2)`$ has rapid decay along “vertical lines”, that is, for $`\sigma ,\sigma ^{}`$ in a fixed compact set and $`(t,t^{})`$ we have $`I(\sigma +it,\sigma ^{}+it^{})_N(1+|t|+|t^{}|)^N`$.
Let $`\kappa >0`$ be fixed. From (5.11), Iwaniec’s bounds for $`L`$-functions near $`1`$, and the rapid decay of $`I`$ along “vertical lines,” we obtain by the maximal modulus principle:
$$(1+|z_1|+|z_2|)^N\mathrm{max}(|I(z_1,z_2)|,|_1I(z_1,z_2)|,|_2I(z_1,z_2)|)_N\mathrm{N}(𝔭)^\kappa ,|\mathrm{}(z_1)|+|\mathrm{}(z_2)|1/2+\kappa ,$$
where $`_1`$ (resp. $`_2`$) is the operator of differentiation w.r.t. $`z_1`$ (resp. $`z_2`$). Put
(5.12)
$$h(z_1,z_2)=z_1^2z_2^2(1/4z_1^2)(1/4z_2^2)\overline{I(\overline{z_1},\overline{z_2})}.$$
Then, in the notation of Sec. 10.3 (esp. Def. 10.1) , $`h`$ belongs to the space $`^{(2)}(\kappa )`$ and satisfies $`h_N_N\mathrm{N}(𝔭)^\kappa `$.
Put $`I=_{\mathrm{}(z_1)=\mathrm{}(z_2)=0}h(z_1,z_2)I(z_1,z_2)𝑑z_1𝑑z_2`$. Then:
(5.13) $`I={\displaystyle _{\mathrm{}(z_1)=0,\mathrm{}(z_2)=0}}h(z_1,z_2)𝑑z_1𝑑z_2{\displaystyle _𝐗}\phi (g)E_1(g,1/2+z_1)E_2(ga([𝔭]),1/2+z_2)𝑑g`$
$`={\displaystyle _𝐗}\phi (x)E_h((x,x)(1,a([𝔭])))𝑑x`$
where the function $`E_h`$ on $`𝐗\times 𝐗`$ is defined by
$$E_h(g_1,g_2)=_{\mathrm{}(z_1)=0,\mathrm{}(z_2)=0}h(z_1,z_2)E_1(g_1,1/2+z_1)E_2(g_2,1/2+z_2),$$
and the interchange of orders is justified by the (easily verified) absolute convergence of the double integral defining $`I`$. Note that $`E_h(g_1,g_2)`$ is totally nondegenerate (see Section 2.7 for definition). We now apply Prop. 4.2 to conclude that $`|I|_{p,d}S_{p,d,2/p}(E_h)\phi _{L^2}\mathrm{N}(𝔭)^{\frac{1}{17p}}`$ for any $`p>4,d1`$. We note at this point that the requirement $`p>4`$ makes it critical that the regularized Eisenstein series $`E_h`$ belong to $`L^4`$; the trivial fact that Eisenstein series belong to $`L^{2ϵ}`$ is far from sufficient.
On the other hand, by Lem. 10.9, $`S_{p,d,2/p}(E_h)h_N`$ for some $`N`$, possibly depending on $`p,d`$. By Prop. 5.1, $`\phi _{L^2}_ϵ\mathrm{N}(𝔭)^ϵ`$.
Thus $`|I|_ϵ\mathrm{N}(𝔭)^{ϵ1/68}`$. Now, by the definition of $`h`$ (5.12) we have $`I=\mathrm{N}(𝔭)^1_{(t,t^{})^2}(1/4+t_1^2)(1/4+t_2^2)t_1^2t_2^2|I(it_1,it_2)|^2𝑑t_1𝑑t_2`$. Thus we obtain:
(5.14)
$$_{(t,t^{})^2}|I(it_1,it_2)|^2t_1^2t_2^2𝑑t_1𝑑t_2_ϵ\mathrm{N}(𝔭)^{11/68+ϵ}.$$
Using (5.14), and the given properties of $`\mathrm{\Phi }`$, we deduce that $`|\mathrm{\Lambda }(\frac{1}{2}+t_0+t_0^{})\mathrm{\Lambda }(\frac{1}{2}+t_0t_0^{})|^2_{t_0,t_0^{}}\mathrm{N}(𝔭)^{11/600}`$ in a similar fashion to the conclusion of the proof of (5.3). We take $`t_0^{}=0`$ to conclude. ∎
## 6. Torus periods (I): subconvex bounds for character twists over a number field.
In this section we shall work in considerable generality; we shall derive subconvex bounds without any assumptions of prime or squarefree conductor, and obtaining polynomial dependence in all auxiliary parameters. This is useful for applications, but will involve some notational overhead. As a result, we have sacrificed good exponents for simplicity at many steps.
###### Theorem 6.1.
Let $`\pi `$ be a (unitary) cuspidal representation of $`\mathrm{GL}_2(𝔸_F)`$, and $`\chi `$ a unitary character of $`𝔸_F^\times /F^\times `$, with finite conductor $`𝔣`$. Then there is $`N>0`$ such that
(6.1) $`L({\displaystyle \frac{1}{2}},\pi \times \chi )\mathrm{Cond}(\pi )^N\mathrm{Cond}_{\mathrm{}}(\chi )^N\mathrm{N}(𝔣)^{1/2\frac{1}{24}},`$
(6.2) $`L({\displaystyle \frac{1}{2}},\chi )\mathrm{Cond}_{\mathrm{}}(\chi )^N\mathrm{N}(𝔣)^{1/4\frac{1}{200}}`$
Note that the result also implies a corresponding statement for the $`L`$-functions evaluated at $`\frac{1}{2}+it`$, since one may replace $`\chi `$ by $`\chi ||^{it}`$.
Since it is perhaps hidden in the proof where the polynomial dependence on conductor arises, we would like to explicate it now. If $`\pi `$ is an automorphic cuspidal representation with analytic conductor $`\mathrm{Cond}(\pi )`$, there exists a vector $`\psi \pi `$ with Sobolev norms $`S_{2,d,\beta }(\psi )\mathrm{Cond}(\pi )^{\mathrm{const}\mathrm{max}(\beta ,d)}`$. Moreover, one can choose such a $`\psi `$ to be a “good” test vector w.r.t. certain toral periods. Thus the analytic conductor enters precisely through the minimal Sobolev norm of a suitable vector belonging to the space of $`\pi `$. We note that the test vectors we choose are smooth but not $`K`$-finite at infinite places; this idea has been heavily exploited in the previous work of Bernstein and Reznikov.
Note that some cases of Thm. 6.1 – where $`\pi `$ has trivial central character and $`\pi `$ is quadratic – are subsumed by the previous result Thm. 5.1. Nevertheless, we have chosen to give a distinct presentation since the method is entirely different, it is simpler in the present method to deal with the case of noncuspidal $`\pi `$. Also, we shall consistently deal in the present section with $`\mathrm{GL}(2)`$, rather than $`\mathrm{PGL}_2`$. Thus $`\omega `$ will be a unitary character of $`𝔸_F^\times /F^\times `$, and $`C_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$ the space of functions on $`𝐗_{\mathrm{GL}(2)}=\mathrm{GL}_2(F)\backslash \mathrm{GL}_2(𝔸_F)`$ with central character $`\omega `$.
In the case $`F=`$, the subconvexity result (6.2) for characters is due to Burgess. Burgess’ method gives a much better exponent; of course there is considerable scope for improvement in the present technique also.
For the ease of the reader, we briefly explain in advance the points of our proof in classical language. The discussion that follows is not a completely faithful rendition of the proof, but it hopefully conveys the main ideas. While it follows the pattern of all the proofs in this paper, one minor complication is that we deal with integrals w.r.t. certain measures of infinite mass.
1. If $`f`$ is a Maass form on $`\mathrm{SL}_2()\backslash `$, the integral
(6.3)
$$\frac{1}{q}_{y=0}^{\mathrm{}}\underset{x=1}{\overset{q}{}}\chi (x)f(\frac{x}{q}+iy)\frac{dy}{y}$$
equals, up to some $`\mathrm{\Gamma }`$-factors, $`\frac{1}{\sqrt{q}}L(\frac{1}{2},f\times \chi )`$. This is an exercise in Hecke-Jacquet-Langlands theory. The version of this equality that we shall use is proved in Lem. 11.8 when $`f`$ is a cusp form and Lem. 11.10 when $`f`$ is Eisenstein.
2. It will then suffice to bound $`_{x=1}^q\chi (x)f(\frac{x}{q}+iy)`$ for each fixed value of $`y`$. As it turns out, the crucial range of $`y`$ is around $`y=q^1`$; the contribution of other $`y`$s are small for relatively trivial reasons (use the Fourier expansion). This is roughly a geometric form of the approximate functional equation: it says that the Fourier coefficients $`a_n(f)`$ with $`nq`$ are most important to determining the $`L`$-function. The general version of this fact is proven in Lem. 11.9.
3. In the range when $`yq^1`$, the set $`\{\frac{x}{q}+iy\}_{\{1xq1\}}`$ is roughly equidistributed, because it is (with the exception of two points) the orbit of $`iqy`$ by the $`q`$th Hecke operator. This is easy to quantify and actually can be regarded as a statement about equidistribution of $`p`$-adic horocycles.The general version of this is proved in Lem. 9.10.
4. We are now in a situation where we are trying to bound the period of $`f`$ along the roughly equidistributed set $`\{\frac{x}{q}+iy\}_{\{1xq1\}}`$. To do this, we apply mixing properties of the adelic torus flow, in the same fashion as the previous proofs of this paper. This shows that $`_{x=1}^q\chi (x)f(\frac{x}{q}+iy)`$ is small.
The computations that underlie steps (1), (2) and (3) are fairly routine but technically complicated. We have therefore carried them out in Sec. 11.4. In the sections that follow, we merely quote the results and carry out what amounts to step (4).
### 6.1. Relating integral representations and periods.
Let $`z`$ and let $`\mu _z,\nu _z,\mu ,\nu `$ be the measures on $`𝐗_{\mathrm{GL}(2)}`$ defined by
(6.4)
$$\begin{array}{c}\mu _z(f)=_{|y|=z}f(a(y)n([𝔣]))\chi (y)d^\times y,\mu =_{z^\times }\mu _zd^\times z,\hfill \\ \hfill \nu _z(f)=_{|y|=z}f(a(y)n([𝔣]))d^\times y,\nu =_{z^\times }\nu _zd^\times z.\end{array}$$
In both cases, the measure $`d^\times y`$ is the probability measure invariant by $`𝔸_F^1/F^\times `$ and the measure $`d^\times z`$ is a Haar measure on $`^\times `$. Thus $`\mu _z,\nu _z`$ are probability measures, whereas $`\mu ,\nu `$ have infinite mass. It is simple to see that the integrals defining $`\mu (f),\nu (f)`$ converge absolutely if $`f`$ is a function decaying rapidly enough at the cusps, e.g. satisfying $`|f(x)|0pt(x)^\epsilon `$ (notation of Sec. 8.2), for any $`\epsilon >0`$. Note also the analogy between these measures and those used in the analysis of unipotent periods (cf. (3.2).) Classically, $`\nu _z(f)`$ should be thought of the measure on $`\mathrm{SL}_2()\backslash `$ defined by $`_{0xq1}f(\frac{x}{q}+iz)`$, and $`\mu _z(f)`$ the measure on $`\mathrm{SL}_2()\backslash `$ defined by $`_{0xq1}f(\frac{x}{q}+iz)\chi (x)`$. (These statements are not to be interpreted precisely; they are for intuition only).
Here is the Proposition that formalizes (1) and (2) of the discussion above, in the cuspidal case.
###### Proposition 6.1.
Let $`\pi `$ be a cuspidal representation on $`\mathrm{GL}_2(𝔸_F)`$, $`\chi `$ a character of $`𝔸_F^\times /F^\times `$ of finite conductor $`𝔣`$. Write $`L_{unr}(s,\pi \times \chi )=_{\chi _v\mathrm{unram}.}L_v(s,\pi \times \chi )`$, where the product is taken over all finite places at which $`\chi `$ is not ramified.
Let $`d,\beta 0`$. Let $`g_+,g_{}`$ be positive smooth functions on $`_0`$ such that $`g_++g_{}=1`$, $`g_+(t)=1`$ for $`t2`$ and $`g_{}(t)=1`$ for all $`t1/2`$.
Then there exists $`\phi \pi `$ such that, with
(6.5)
$$\mathrm{\Phi }(s)=\mathrm{N}(𝔣)^{1/2}\frac{_z\mu _z(\phi )|z|^{s1/2}d^\times z}{L_{unr}(s,\pi \times \chi )}$$
then $`\mathrm{\Phi }(s)`$ is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(s)|_{\mathrm{}(s),ϵ}\mathrm{N}(𝔣)^ϵ`$ and $`|\mathrm{\Phi }(\frac{1}{2})|_ϵ\mathrm{N}(𝔣)^ϵ`$.
2. $`\phi `$ is new at every finite place (i.e., for each finite prime $`𝔮`$ it is invariant by $`K_0[𝔮^{s_𝔮}]`$, where $`s_𝔮`$ is the local conductor of $`\pi `$).
3. The Sobolev norms of $`\phi `$ satisfy the bounds
(6.6)
$$S_{2,d,\beta }(\phi )_ϵ\mathrm{Cond}_{\mathrm{}}(\pi )^{2d+ϵ}\mathrm{Cond}_f(\pi )^{\beta +ϵ}\mathrm{Cond}_{\mathrm{}}(\chi )^{1/2+2d}$$
4. The integration of (6.5) may be “truncated without significant change” to the region $`z`$ around $`\mathrm{N}(𝔣)^1`$; more formally:
$$\left|_z\mu _z(\phi )g_+(z/T)d^\times z\right|(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\pi )\mathrm{Cond}(\chi ))^ϵ$$
$$\left|_z\mu _z(\phi )g_{}(z/T)d^\times z\right|(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵ(\mathrm{Cond}_{\mathrm{}}(\chi )\mathrm{Cond}(\pi ))^{1+ϵ}$$
###### Proof.
Lem. 11.8 and Lem. 11.9. ∎
We next give the corresponding result for the “noncuspidal case.” We recall that the Eisenstein series $`E_\mathrm{\Psi }(s,g)`$ associated to a Schwarz function $`\mathrm{\Psi }`$ on $`𝔸_F^2`$ are discussed in Sec. 10. The normalization is so that the functional equation interchanges $`s`$ and $`1s`$. $`\overline{E}(s,g)`$ denotes, as explained in that section (cf. (10.9)) the truncated Eisenstein series obtained by subtracting the constant term; it is a function on $`B(F)\backslash \mathrm{GL}_2(𝔸_F)`$.
###### Proposition 6.2.
Let $`s_0,s_0^{}`$, and suppose that $`\chi `$ is ramified at at least one finite place. Let $`g_\pm `$ be as in Prop. 6.1. There is an absolute $`C>0`$ (i.e. depending only on $`F`$) and a choice of $`K_{\mathrm{max}}`$-invariant Schwarz function $`\mathrm{\Psi }`$ (depending on $`\chi `$) so that if we put
$$\mathrm{\Phi }(s,s^{}):=\mathrm{N}(𝔣)^{1/2}\frac{_{y𝔸_F^\times /F^\times }\overline{E}_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)|y|^s^{}d^\times y}{L(\chi ,s+s^{})L(\chi ,1s+s^{})}$$
where $`\overline{E}`$ is defined as in (10.9), then the integral defining $`\mathrm{\Phi }`$ is absolutely convergent when $`\mathrm{}(s),\mathrm{}(s^{})1`$. Moreover, $`\mathrm{\Phi }`$ extends from $`\mathrm{}(s),\mathrm{}(s^{})1`$ to a holomorphic function on $`^2`$, satisfying
1. $`|\mathrm{\Phi }(s_0,s_0^{})|1`$ and $`|\mathrm{\Phi }(s,s^{})|C^{1+|\mathrm{}(s)|+|\mathrm{}(s^{})|}(1+|s|+|s^{}|)^C`$.
Moreover, given $`N>0`$ we have that
(6.7)
$$|\mathrm{\Phi }(s,s^{})|(1+|s|+|s^{}|)^N_{\mathrm{}(s),\mathrm{}(s^{}),N}\mathrm{Cond}_{\mathrm{}}(\chi )^N^{}$$
where $`N^{}`$ and the implicit constant may be taken to depend continuously on $`N,\mathrm{}(s),\mathrm{}(s^{})`$.
2. $`\mathrm{\Psi }`$, and so also $`E_\mathrm{\Psi }(s,g)`$ is invariant by $`K_{\mathrm{max}}`$;
3. Let $`h(\kappa )`$ be as in (10.18), and put $`E_h:=_{\mathrm{}(s)1}h(s)E_\mathrm{\Psi }(s,g)𝑑g`$. Then, for each $`d,\beta `$, there is $`N>0`$ such that $`S_{\mathrm{},d,\beta }(E_h)_\kappa h_0\mathrm{Cond}_{\mathrm{}}(\chi )^N`$, where the norm $`h_0`$ is defined in (10.18).
4. We have $`\mu _z(E_h)_{K,\mathrm{\Psi },h}\mathrm{min}(|z|^K,|z|^K)`$ for each<sup>16</sup><sup>16</sup>16The implicit constants here are totally unimportant; this estimate will be used only to verify that certain integrals converge. $`K1`$. Moreover, there is $`N>0`$ such that
$$\left|_z\mu _z(E_h)g_+(z/T)d^\times z\right|(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵh_N$$
$$\left|_z\mu _z(E_h)g_{}(z/T)d^\times z\right|(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵ\mathrm{Cond}_{\mathrm{}}(\chi )^{1+ϵ}h_N$$
###### Proof.
Lem. 11.10 and Lem. 11.11. ∎
### 6.2. Proof of theorem – cuspidal case.
Let $`\chi `$ be a character of $`𝔸_F^\times /F^\times `$, of varying conductor $`𝔣`$. Put $`q=\mathrm{N}(𝔣)`$.
We shall need the following estimate, proved in Lem. 9.10. It amounts in essence to a statement about the equidistribution of $`p`$-adic horocycles (classically, these roughly correspond to a statement about the equidistribution of $`\{\frac{x}{q}+iz\}_{0xq1}`$, if $`zq^1`$).
For any function $`f`$ that is invariant by $`_𝔮K_0[𝔮^{s_𝔮}]`$, we have (with $`𝔪=_𝔮𝔮^{s_𝔮}`$)
(6.8)
$$|\nu _z(f)_𝐗f|_ϵq^{\alpha 1/2+ϵ}\mathrm{N}(𝔪)^{3/2+ϵ}\mathrm{max}(qz,\frac{1}{qz})^{1/2}S_{2,d}(f),fC^{\mathrm{}}(𝐗)$$
###### Proof.
(of Thm. 6.1 – cuspidal case.) Choose $`f\pi `$ to be the “$`\phi `$” of Prop. 6.1, so that $`|\mu (f)|_ϵ\mathrm{N}(𝔣)^{1/2ϵ}|L_{unr}(1/2,\pi \times \chi )|`$. For each ramified prime $`𝔮`$ of $`\pi `$, let $`𝔮^{s_𝔮}`$ be the local conductor of the local representation $`\pi _𝔮`$. Set $`𝔪:=_𝔮𝔮^{s_𝔮}`$, the finite conductor of $`\pi `$.
Let $`K1`$ be an integer satisfying $`K\mathrm{N}(𝔣)`$. Let $`𝒮`$ be the set of prime ideals of $`𝔬_F`$, with norm lying in $`[K,2K]`$, and satisfying $`(𝔫,𝔣)=1`$ and $`(𝔫,\mathrm{Supp}(f))=1`$. Fix $`𝔫_0𝒮`$. For each prime ideal $`𝔫𝒮`$, let $`\varpi _𝔫F_𝔫`$ be a uniformizer. We define a measure $`\sigma `$ on $`\mathrm{GL}_2(𝔸_F)`$ so that
(6.9)
$$\sigma =|𝒮|^1\underset{𝔫𝒮}{}\chi (\varpi _𝔫)\chi (\varpi _{𝔫_0})^1\delta _{a_𝔫(\varpi _𝔫)a_{𝔫_0}(\varpi _{𝔫_0})^1}.$$
Clearly $`\sigma `$ has total mass $`1`$. Moreover, $`\mu (f)=\mu (f\sigma )`$ and $`f\sigma `$ is invariant by $`K_0[𝔮^{s_𝔮}]`$ for each $`𝔮|𝔪`$.
Choose $`\kappa `$ “slightly smaller than $`1`$,” to be specified later. Our aim is now to cut the $`z`$ integration in $`\mu =_z\mu _zd^\times z`$ into three ranges, the crucial range of which will be $`q^{2+\kappa }zq^\kappa `$; this avoids the pain of dealing with the infinite mass measure $`\mu `$. Let $`g_+,g_{}`$ be as in Prop. 6.1. Define $`h(t)`$ by the rule $`g_{}(\frac{z}{q^{2+\kappa }})+h(t)+g_+(\frac{t}{q^\kappa })=1`$.
(6.10)
$$\begin{array}{c}|\mu (f)|^2=|\mu (f\sigma )|^2=\left|_z\mu _z(f\sigma )d^\times z\right|^2\hfill \\ \hfill \left|g_{}(\frac{z}{q^{2+\kappa }})\mu _z(f\sigma )\right|^2+\left|h(z)\mu _z(f\sigma )d^\times z\right|^2+\left|g_+(\frac{z}{q^\kappa })\mu _z(f\sigma )d^\times z\right|^2\end{array}$$
By Prop. 6.1, the first and last term (without the square, e.g. $`\left|g_+(\frac{z}{q^\kappa })\mu _z(f\sigma )d^\times z\right|`$) are $`\mathrm{Cond}(\pi )^{1+ϵ}\mathrm{Cond}_{\mathrm{}}(\chi )^{1+ϵ}q^{\frac{\kappa 1}{2}+ϵ}`$. More explicitly, we note that
(6.11)
$$\mu _z(f\sigma )=|S|^1\underset{𝔫𝒮}{}\mu _{\mathrm{N}(𝔫)^1\mathrm{N}(𝔫_0)z}(f)$$
Now $`1/2\mathrm{N}(𝔫)\mathrm{N}(𝔫_0)^12`$ for all $`𝔫`$ – this was the purpose of the factors involving $`𝔫_0`$ in (6.9) – and so one easily deduces a bound on $`g_{}(\frac{z}{q^{2+\kappa }})\mu _z(f\sigma )d^\times z`$ from the final assertion of Prop. 6.1. Similarly for the term involving $`g_+`$.
That the first and last terms of (6.10) should be less significant may be seen in the classical setting from the Fourier expansion; it should be regarded as a geometric version of the approximate functional equation).
As for the intermediate term, we note
$$_zh(z)\mu _z(f)=_{y𝔸_F^\times /F^\times }h(|y|)\chi (y)f(a(y)n([𝔣]))d^\times y.$$
Applying Cauchy-Schwarz, and the fact that $`_{𝔸_F^\times /F^\times }h(|y|)d^\times y\mathrm{log}(q)`$, we get:
(6.12)
$$\begin{array}{c}\left|_zh(z)\mu _z(f\sigma )\right|^2_ϵq^ϵ_zh(z)\nu _z(|f\sigma |^2)d^\times z\hfill \\ \hfill _ϵq^ϵ_𝐗|f\sigma |^2𝑑\mu _𝐗+q^{\alpha \kappa /2+ϵ}\mathrm{N}(𝔪)^{3/2+ϵ}S_{2,d}(|f\sigma |^2),\end{array}$$
where we have applied (6.8). By Lemmas 8.1 and 8.2,
(6.13)
$$\begin{array}{c}S_{2,d,\beta }(|f\sigma |^2)S_{4,d,\beta }(f\sigma )^2\hfill \\ \hfill (\mathrm{sup}_{g\mathrm{supp}(\sigma )}g)^{2\beta }S_{4,d,\beta }(f)^2K^{4\beta }S_{2,d^{},\beta +3/2}(f)^2.\end{array}$$
where the last line holds for $`d^{}d`$, and we have used Lem. 9.3 (which bounds the $`L^{\mathrm{}}`$ norm of a cusp form in terms of $`L^2`$ norms), together with the easily verified fact that $`\mathrm{sup}_{g\mathrm{supp}(\sigma )}gK^2`$.
By bounds towards Ramanujan,
(6.14)
$$f\sigma _{L^2}^2=_{g,g^{}}g^1g^{}f,f𝑑\sigma (g)𝑑\sigma (g^{})K^{2\alpha 1}f_{L^2}^2.$$
Thus:
(6.15)
$$\begin{array}{c}|\mu (f)|_ϵ\mathrm{Cond}(\pi )\mathrm{Cond}_{\mathrm{}}(\chi )(\mathrm{Cond}(\chi )\mathrm{Cond}(\pi ))^ϵq^{\frac{\kappa 1}{2}}\hfill \\ \hfill +\left(K^{\alpha 1/2}q^ϵ+\mathrm{N}(𝔪)^{3/4+ϵ}q^{\kappa /4+\alpha /2+ϵ}K\right)S_{2,d,2}(f)\\ \hfill q^ϵ(q^{(\kappa 1)/2}+K^{\alpha 1/2}+q^{\alpha /2\kappa /4}K)\mathrm{Cond}_{\mathrm{}}(\chi )^N\mathrm{Cond}(\pi )^N,\end{array}$$
for appropriate $`N>0`$. We have used Prop. 6.1, (3) at the last step.
Prop. 6.1 guarantees that $`|L_{unr}(1/2,\pi \times \chi )|_ϵq^{1/2+ϵ}|\mu (f)|`$. From this, optimizing $`\kappa ,K`$, and applying trivial bounds at ramified places, we obtain the conclusion, taking $`\alpha =3/26`$. ∎
### 6.3. Proof of theorem – noncuspidal case.
We turn to the proof of (6.2). This is very similar, but we implement a mild regularization procedure to deal with the Eisenstein series, just as in the case of Rankin-Selberg $`L`$-functions.
###### Proof.
(of (6.2).) We may assume that $`\chi `$ ramifies at least at one finite place.
Let $`\mathrm{\Psi }`$ be a Schwarz function on $`𝔸_F^2`$, $`E(g,s):=E_\mathrm{\Psi }(g,s)`$ the corresponding Eisenstein series, chosen as in Prop. 6.2 with $`s_0=1/2,s_0^{}=0`$.
Let $`\kappa ^{}>0`$, let $`h`$ be holomorphic in an open neighbourhood of the vertical strip $`\kappa ^{}\mathrm{}(s)1+\kappa ^{}`$ and put $`E_h(s)=_{\mathrm{}(s)=1+\kappa ^{}}h(s)E(g,s)𝑑s`$. Then if $`h(0)=h(\frac{1}{2})=h(1)=0`$, it follows from the third assertion of Prop. 6.2 that
$$S_{\mathrm{},d,\beta }(E_h)_{ϵ,d}\mathrm{Cond}_{\mathrm{}}(\chi )^Nh_0$$
for appropriate $`N=N(d,\beta )>0`$. Here, as in (10.18) with $`\kappa `$ replaced by $`\kappa ^{}`$, the norm $`h_N`$ is defined to be $`_{\mathrm{}}^{\mathrm{}}\left(|h(1+\kappa ^{}+it)|+|h(\kappa ^{}+it)|\right)(1+|t|)^N𝑑t`$.
Put, in the notation of Prop. 6.2, $`I(s)=\mathrm{\Phi }(s,0)L(\chi ,s)L(\chi ,1s)`$. Then:
(6.16)
$$_z\mu _z(E_h)d^\times z=\mathrm{N}(𝔣)^{1/2}_{\mathrm{}(s)=1/2}h(s)I(s)𝑑s.$$
This is established in (11.31); for now, we remark that that this is “almost” obvious from Proposition 6.2, the only additional point being that one can replace $`E`$ by $`\overline{E}`$, and this is exactly where the fact that $`\chi `$ is ramified at a finite place comes in – to kill the constant term of the Eisenstein series.<sup>17</sup><sup>17</sup>17The classical version of this fact – see (6.3) – it is clear that the $`\chi `$-sum will kill any constant term of $`f`$, as long as $`\chi `$ is not trivial.
Take $`h=(s1/2)^2s(1s)\overline{I(1\overline{s})}`$. The “good” analytic properties of $`h`$, e.g. rapid decay along vertical lines, follow<sup>18</sup><sup>18</sup>18This point was not clear in a previous version; thanks to N. Bergeron for pointing this out. from (6.7). In particular, $`h`$ belongs to the function spaces $`(\kappa )`$ defined in (10.18) for any $`\kappa >0`$, and the norms $`h_N`$ are all bounded by suitable powers of $`\mathrm{Cond}_{\mathrm{}}(\chi ).q`$.
Then (6.16) becomes
(6.17)
$$_{t=\mathrm{}}^{\mathrm{}}t^2(1/4+t^2)|I(\frac{1}{2}+it)|^2=q^{1/2}_z\mu _z(E_h)d^\times z$$
To bound the right-hand side, we proceed as in Sec. 6.2, but with $`f`$ replaced by $`E_h`$. We use notation as in that Section, except replacing the “$`h`$” defined before (6.10) by $`1g_{}g_+`$ to avoid clashing with its alternate usage here.
One proves as in that Section, that for $`d1`$:
(6.18)
$$\begin{array}{c}\left|_z(1g_{}g_+)\mu _z(E_h\sigma )\right|_ϵq^ϵ\left(K^{\alpha 1/2}+q^{\alpha /2\kappa /4}(\mathrm{sup}_{g\mathrm{supp}(\sigma )}g)^{1/2}\right)S_{4,d,1/2}(E_h)\hfill \\ \hfill q^ϵ\left(K^{\alpha 1/2}+q^{\alpha /2\kappa /4}K\right)h_0\mathrm{Cond}_{\mathrm{}}(\chi )^N\end{array}$$
for some appropriate $`N>0`$. At the last stage we have applied Prop. 6.2 to control the Sobolev norm.
Prop. 6.2 also guarantees that, for appropriate $`N>0`$, we have:
(6.19)
$$\begin{array}{c}\left|g_{}(\frac{z}{q^{2+\kappa }})\mu _z(E_h\sigma )\right|+\left|g_+(\frac{z}{q^\kappa })\mu _z(E_h\sigma )d^\times z\right|\hfill \\ \hfill \mathrm{Cond}_{\mathrm{}}(\chi )^{1+ϵ}q^{\frac{\kappa 1}{2}+ϵ}h_N,\end{array}$$
Combining (6.17) and (6.19), we obtain as in the previous Section the bound, for sufficiently large $`N`$:
(6.20)
$$\begin{array}{c}_{t=\mathrm{}}^{\mathrm{}}t^2(\frac{1}{4}+t^2)|L(\frac{1}{2}+it,\chi )L(\frac{1}{2}it,\chi )|^2|\mathrm{\Phi }(1/2+it,0)|^2\hfill \\ \hfill _ϵ\left(q^{\frac{\kappa 1}{2}}+K^{\alpha 1/2}+q^{\alpha /2\kappa /4}K\right)h_Nq^{1/2+ϵ}\mathrm{Cond}_{\mathrm{}}(\chi )^N\end{array}$$
One applies the convexity bound to bound $`h_N`$, obtaining
$$_{\mathrm{}}^{\mathrm{}}t^2|L(\frac{1}{2}+it,\chi )|L(\frac{1}{2}it,\chi )|^2|\mathrm{\Phi }(1/2+it,0)|^2\mathrm{Cond}_{\mathrm{}}(\chi )^Nq^{24/25},$$
where we have increased $`N`$ as necessary. From this we get $`L(\frac{1}{2},\chi )\mathrm{Cond}_{\mathrm{}}(\chi )^Nq^{1/41/200}`$. ∎
## 7. Torus periods (II): equidistribution of compact torus orbits.
It has been independently shown by Zhang , Clozel-Ullmo and P. Cohen that the subconvexity result Thm. 6.1 implies the equidistribution of Heegner points over totally real fields; in particular, they pointed out that GRH implies this equidistribution. Thm. 6.1 makes this result unconditional.
The main aim of this section is to explain how one can obtain certain conditional results about equidistribution of subsets of Heegner points, and how this fits into the general framework of “sparse equidistribution questions.” In particular, this approach does not rely on reducing questions about subsets of Heegner points to subconvexity, but rather approaches the equidistribution question directly.
The proofs of the results (and various supporting Lemmas) will only be sketched, and we will confine ourselves for simplicity to the case of narrow class number $`1`$; we will in any case present an unconditional approach, based on combining the ideas of this paper with the ideas of Michel, in the paper (joint with P. Michel). We nevertheless feel that the ideas presented here may be of use in other contexts. Indeed, this section is of a different flavor to the other Sections; it uses “adelic analysis” more genuinely.
In fact, we shall need a mild refinement of the results of , which allow better control of the dependence on the test vectors. We state this refinement without proof in Thm. 7.1; the proof is an exercise in explicating some of the proofs in .
### 7.1. Equidistribution of Heegner points.
We recall the definition of Heegner points. Let $`F`$ be a totally real number field of degree $`d`$ over $``$. For simplicity we shall confine ourselves to the case where the ring of integers of $`F`$ has narrow class number $`1`$. This assumption does not change any of the technical details, which are in any case carried out adelically; it simply allows us to be a little more explicit about the torus orbits we consider. Let $`E=F(\sqrt{𝐝})`$ be a totally imaginary quadratic extension of $`F`$, where $`𝐝𝔬_F`$ is totally positive and squarefree. Here “squarefree” means that it is of valuation $`1`$ at all finite places.
Let $`T_E`$ be the torus $`\mathrm{Res}_{E/F}(𝔾_m)/𝔾_m`$; we embed $`T_E`$ in $`\mathrm{PGL}_2`$ via (in obvious notation):
(7.1)
$$\iota _E:x+y\sqrt{𝐝}\left(\begin{array}{cc}x& y\\ y𝐝& x\end{array}\right).$$
Regard $`𝐝`$ as an element of $`F`$ via the inclusion $`FF`$. Since it is totally positive, it possesses a unique totally positive square root, $`\sqrt{𝐝}F`$. Set $`[𝐝]_{\mathrm{}}=\left(\begin{array}{cc}1& 0\\ 0& \sqrt{𝐝}\end{array}\right)\mathrm{PGL}_2(F)`$. We define a map $`:T_E(𝔸_F)/T_E(F)𝐗`$ via
(7.2)
$$:x\iota _E(x)[𝐝]_{\mathrm{}},$$
where we regard $`[𝐝]_{\mathrm{}}\mathrm{PGL}_2(F)\mathrm{PGL}_2(𝔸_F)`$ acting by right translation on $`𝐗`$. Denote by $`\mathrm{N}(𝐝)`$ the absolute norm of $`𝐝`$, i.e. $`\mathrm{N}(𝐝)=|𝔬_F/𝐝𝔬_F|`$.
The $`F`$-torus $`T_E`$ is anisotropic, and there is a unique $`T_E(𝔸_F)`$-invariant probability measure on $`T_E(𝔸_F)/T_E(F)`$. Let $`\nu _E`$ be its image by the map $``$.
###### Theorem 7.1.
Set $`E=F(\sqrt{𝐝})`$, where $`𝐝𝔬_F`$ is totally positive and squarefree. The measures $`\nu _E`$ become equidistributed as $`\mathrm{N}(𝐝)\mathrm{}`$. Indeed, there exist $`\delta >0,d,\beta `$ such that for $`fC^{\mathrm{}}(𝐗)`$ we have
$$\left|f𝑑\nu _E_𝐗f(x)𝑑x\right|\mathrm{N}(𝐝)^\delta S_{\mathrm{},d,\beta }^{}(f).$$
Recall the definition of $`S^{}`$ from Sec. 2.10. We do not give the proof; as we have remarked it can be obtained by following the computations of a little more explicitly.
One recovers from Thm. 7.1 the equidistribution of certain Heegner points associated to $`E=F(\sqrt{𝐝})`$ as $`𝐝`$ varies. Thm. 7.1 also gives an effective rate of equidistribution for Heegner points with polynomial dependence on the level and the eigenvalue of a test function. This rather innocuous polynomial dependence (in the level aspect, at least) will in fact play a crucial role in our deduction of the equidistribution of sparse subsets in the following section.
### 7.2. Equidistribution of subsets of Heegner points.
We turn to certain conditional results on equidistribution of sparse subsets. $`F`$ being as in Sec. 7.1, let $`E_i=F(\sqrt{𝐝_i})`$ be a sequence of distinct quadratic, totally imaginary, extensions of $`F`$. For each $`E_i`$, let $`S_iT_{E_i}(𝔸_F)/T_{E_i}(F)`$ be a subgroup of finite index $`m_i`$. Let $`\mu _{E_i}^{S_i}`$ be the image of the Haar probability measure on $`S_i`$ by the map $``$.
The import of the next theorem is that, if $`E_i`$ has enough small split primes, one can obtain the equidistribution of the measures $`\mu _{E_i}^{S_i}`$ as $`i\mathrm{}`$. This result is quite similar to the results of Duke-Friedlander-Iwaniec in the case $`F=`$, although the method is at least superficially rather different. One can also contrast with Michel’s striking result, for $`F=`$, that gives a comparable result but without the condition on enough small split primes. Our method is different to Michel, who deduces the result from his subconvexity bound for Rankin-Selberg $`L`$-functions. <sup>19</sup><sup>19</sup>19We note that this bound of Michel is considerably deeper than (5.3), since it deals with varying central character. For some speculative discussion on the “reason” that Michel’s method can avoid this condition, see the last paragraph of . In the present approach, we prove the equidistribution theorem directly. In a sequel to this paper, the author and P. Michel combine the methods here with some methods developed by Michel to make the results of this section unconditional.
To quantify the existence of enough small split primes, one might impose the condition (as does Linnik ) that the $`E_i`$ vary through a sequence of quadratic extensions that split at a fixed prime of $`F`$. We will prefer to take a more quantitative approach, which will yield a stronger result at the price of a stronger assumption. In that regard we introduce the following notation: For $`\delta >0`$, we put
$$\mathrm{wt}(E_i,\delta )=\mathrm{\#}\{𝔮𝔬_F\text{ prime and split in }E_i,\mathrm{N}(𝐝_i)^\delta \mathrm{N}(𝔮)2\mathrm{N}(𝐝_i)^\delta \}.$$
###### Theorem 7.2.
There exists $`\delta _1>0`$ such that, if $`\frac{m_i}{\mathrm{min}(\mathrm{N}(𝐝_i)^{\delta _1(1/2\alpha )},\mathrm{wt}(E_i,\delta _1)^{1/2})}0`$, the sequence $`\mu _{E_i}^{S_i}`$ converges, as $`i\mathrm{}`$, to the invariant measure on $`𝐗`$.
###### Proof.
This is deduced from Thm. 7.1 by using the mixing properties of the $`T_E(𝔸_F)`$-flow. Indeed, we fix an index $`i`$ and a corresponding field $`E_i`$. Let $`\delta _1>0`$ be fixed. Let $`𝒮`$ be the set of prime ideals of $`F`$ which split in $`E_i`$ and with norm in $`[\mathrm{N}(𝐝_i)^{\delta _1},2\mathrm{N}(𝐝_i)^{\delta _1}]`$. For each $`𝔮𝒮`$, the torus $`T_{E_i}(F_𝔮)`$ is isomorphic to $`F_𝔮^\times `$. Fix an isomorphism $`\mathrm{{\rm Y}}_𝔮:T_{E_i}(F_𝔮)F_𝔮^\times `$, and let $`\varpi _𝔮`$ be an element in $`T_{E_i}(F_𝔮)`$ such that $`\mathrm{{\rm Y}}_𝔮(\varpi _𝔮)`$ has valuation $`\pm 1`$ in $`F_𝔮^\times `$.
Let $`\chi `$ be a character of $`T_{E_i}(𝔸_F)/T_{E_i}(F)`$, trivial on $`S_i`$. Let $`\nu _{E_i}`$ be as defined prior to Thm. 7.1, and define
$$\mu _{E_i}(f)=_{tT_{E_i}(𝔸_F)/T_{E_i}(F)}f((t))\chi (t)𝑑t,$$
where $`dt`$ is the Haar probability measure on $`T_{E_i}(𝔸_F)/T_{E_i}(F)`$. Let $`\sigma `$ be the probability measure $`\frac{1}{|𝒮|}_{𝔮𝒮}\chi (\varpi _𝔮)\delta _{\varpi _𝔮}`$ on $`T_{E_i}(𝔸_F)`$. Then $`\mu _{E_i}(f)=\mu _{E_i}(f_{}\sigma )`$, where $`_{}\sigma `$ denotes the image of $`\sigma `$ by the map $``$.
By Cauchy-Schwarz, and Thm. 7.1,
(7.3)
$$\begin{array}{c}|\mu _{E_i}(f_{}\sigma )|^2\nu _{E_i}(|f_{}\sigma |^2)\hfill \\ \hfill f_{}\sigma _{L^2}^2+O\left(\mathrm{N}(𝐝_i)^\delta S_{\mathrm{},d,\beta }^{}(|f_{}\sigma |^2)\right),\end{array}$$
where $`\delta ,d,\beta `$ are as in Thm. 7.1. Now, appropriate variants of Lem. 8.1 and 8.2 (for $`S^{}`$ instead of $`S`$) show that
(7.4)
$$\begin{array}{c}S_{\mathrm{},d,\beta }^{}(|f_{}\sigma |^2)S_{\mathrm{},d,\beta }^{}(f_{}\sigma )^2\hfill \\ \hfill \underset{g\mathrm{supp}_{}\sigma }{sup}g^{6\beta }S_{\mathrm{},d,\beta }^{}(f)^2\mathrm{N}(𝐝_i)^{6\delta _1\beta }S_{\mathrm{},d,\beta }^{}(f)^2\end{array}$$
and bounds towards Ramanujan show that
(7.5)
$$f_{}\sigma _{L^2}^2\left(\mathrm{N}(𝐝_i)^{\delta _1(2\alpha 1)}+|𝒮|^1\right)f_{L^2}^2.$$
We note that (7.4) and (7.5) are very closely analogous to (6.13) and (6.14), with $`K`$ replaced by $`\mathrm{N}(𝐝_i)^{\delta _1}`$. In the context of (6.14), the set $`𝒮`$ has size $`K^{1ϵ}`$; thus the term $`|𝒮|^1`$ that appears in (7.5) could be neglected.
Recalling the definition of $`\mu _{E_i}`$, we conclude
(7.6)
$$\begin{array}{c}\left|_{tT_{E_i}(𝔸_F)/T_{E_i}(F)}f((t))\chi (t)𝑑t\right|\hfill \\ \hfill \left(\mathrm{N}(𝐝_i)^{3\delta _1\beta \delta /2}+\mathrm{N}(𝐝_i)^{\delta _1(\alpha 1/2)}+|𝒮|^{1/2}\right)S_{\mathrm{},d,\beta }^{}(f).\end{array}$$
Summing the left-hand side of (7.6) over all $`m_i`$ characters $`\chi `$ of $`T_{E_i}(𝔸_F)/T_{E_i}(F)`$ that are trivial on $`S_i`$, and substituting $`|𝒮|=\mathrm{wt}(E_i,\delta _1)`$, we obtain:
$$\left|\mu _{E_i}^{S_i}(f)\right|m_i\left(\mathrm{N}(𝐝_i)^{3\delta _1\beta \delta /2}+\mathrm{N}(𝐝_i)^{\delta _1(\alpha 1/2)}+\mathrm{wt}(E_i,\delta _1)^{1/2}\right)S_{\mathrm{},d,\beta }^{}(f).$$
Choosing $`\delta _1`$ sufficiently small (the exact value will depend on the value of $`\beta ,\delta `$ from Thm. 7.1) we obtain the claimed conclusion. ∎
## 8. Background on Sobolev norms and reduction theory.
The rest of the paper consists of technical Lemmas. The sections that follow are arranged to be used as a reference, rather than to be read through.
### 8.1. Formal properties of the Sobolev norms.
We begin by explicating certain formal properties of the Sobolev norms defined in Sec. 2.9.3.
###### Remark 8.1.
The following properties of this definition are formal and will be repeatedly used:
1. Translations by $`K_{\mathrm{max},𝐆}`$ preserve $`S_{p,d,\beta }`$, i.e. $`S_{p,d,\beta }(kf)=S_{p,d,\beta }(f)`$ for $`kK_{\mathrm{max},𝐆}`$.
2. If $`L:C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ is a linear functional and $`|L(\psi )|PS_{p,d,\beta }(\psi )`$, then also $`|L(\psi )|S_{p,d,\beta }(\psi )`$. Indeed $`\psi |L(\psi )|`$ is itself a seminorm on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$.
3. Suppose that $`E:C_\omega ^{\mathrm{}}(𝐗_𝐆)C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ is a linear endomorphism satisfying $`PS_{p,d,\beta }(Ef)APS_{p,d,\beta }(f)`$, for some $`A`$. Then also $`S_{p,d,\beta }(Ef)AS_{p,d,\beta }(f)`$. Indeed, $`fA^1S_{p,d,\beta }(Ef)`$ is a seminorm dominated by $`PS_{p,d,\beta }`$.
4. We shall need a slight variant of (3) in the case where we are studying only the space of $`f`$ with some invariance property.
Suppose that $`E:C_\omega ^{\mathrm{}}(𝐗_𝐆)C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ is a linear endomorphism, $`M`$ is a finite set of finite places, and for each $`vM`$ we are given an open compact $`K_{1,v}K_v`$. Suppose moreover that $`PS_{p,d,\beta }(Ef)APS_{p,d,\beta }(f)`$ for some $`A`$ and for all $`f`$ which are $`_{vM}K_{1,v}`$-fixed. Then, for all $`f`$ which are $`_{vM}K_{1,v}`$-fixed, we have in fact
$$S_{p,d,\beta }(Ef)A\underset{vM}{}[K_v:K_{1,v}]^\beta S_{p,d,\beta }(f).$$
Indeed, put $`K_{1,M}=_{vM}K_{1,v}`$ and let $`\mathrm{\Pi }`$ be the averaging operator $`_{kK_{1,M}}\pi (k)𝑑k`$, where $`K_{1,M}`$ is endowed with the Haar probability measure. Then apply (3) above to the operator $`fE(\mathrm{\Pi }f)`$.
###### Lemma 8.1.
Let $`F_1C_{\omega _1}^{\mathrm{}}(𝐗_𝐆)`$, $`F_2C_{\omega _2}^{\mathrm{}}(𝐗_𝐆)`$. Then
$$S_{p,d,\beta }(F_1F_2)_dS_{2p,d,\beta }(F_1)S_{2p,d,\beta }(F_2).$$
Note that $`F_1F_2C_{\omega _1\omega _2}^{\mathrm{}}(𝐗_𝐆)`$.
###### Proof.
Put $`F=F_1F_2`$. For any monomial $`𝒟`$ of degree $`d`$ in $``$, we can write $`𝒟(F_1F_2)=_\alpha (𝒟_{\alpha ,1}F_1)(𝒟_{\alpha ,2}F_2)`$, where $`\alpha `$ range over an index set $``$ whose size is bounded by a constant depending only on $`d`$, and the $`𝒟_{\alpha ,}`$ are certain monomials in $``$ satisfying $`\mathrm{ord}(𝒟_{\alpha ,1})+\mathrm{ord}(𝒟_{\alpha ,2})=d`$. It follows that
$$𝒟F_{L^p(𝐗_{𝐆,\mathrm{ad}})}\underset{\alpha }{}\left(_{𝐗_{𝐆,\mathrm{ad}}}|𝒟_{\alpha ,1}F_1|^p|𝒟_{\alpha ,2}F_2|^p\right)^{1/p}.$$
Applying Cauchy-Schwarz, we conclude
(8.1) $`𝒟F_{L^p(𝐗_{𝐆,\mathrm{ad}})}{\displaystyle \underset{\alpha }{}}𝒟_{\alpha ,1}F_1_{L^{2p}(𝐗_{𝐆,\mathrm{ad}})}𝒟_{\alpha ,2}F_2_{L^{2p}(𝐗_{𝐆,\mathrm{ad}})}`$
Clearly, for each finite place $`v`$, we have $`K_{v,F}K_{v,F_1}K_{v,F_2}`$; in particular $`[K_{\mathrm{max},𝐆}:K_F][K_{\mathrm{max},𝐆}:K_{F_1}][K_{\mathrm{max},𝐆}:K_{F_2}]`$. It follows that
(8.2)
$$\begin{array}{c}[K_{\mathrm{max},𝐆}:K_F]^\beta \underset{𝒟}{}𝒟F_{L^p(𝐗_{𝐆,\mathrm{ad}})}\hfill \\ \hfill ([K_{\mathrm{max},𝐆}:K_{F_1}]^\beta \underset{𝒟}{}𝒟F_1_{L^{2p}})([K_{\mathrm{max},𝐆}:K_{F_2}]^\beta \underset{𝒟}{}𝒟F_2_{L^{2p}}),\end{array}$$
where the implicit constant depends only on $`d`$, and in all three instances $`𝒟`$ varies over the set of monomials in $``$ of degree $`d`$.
That is to say, there is a constant $`C=C(d)`$ such that
$$PS_{p,d,\beta }(F_1F_2)CPS_{2p,d,\beta }(F_1)PS_{2p,d,\beta }(F_2).$$
From (2.11) we deduce
$$S_{p,d,\beta }(F_1F_2)CS_{2p,d,\beta }(F_1)S_{2p,d,\beta }(F_2),$$
as required. ∎
We recall the definition of $`g`$ for $`g𝐆(F_{\mathrm{}}),𝐆(𝔸_F)`$ etc. from Sec. 2.4.
###### Lemma 8.2.
Let $`FC^{\mathrm{}}(𝐗_𝐆)`$ and $`g=(g_{\mathrm{}},g_f)𝐆(𝔸_F)`$.
$$S_{p,d,\beta }(gF)g_{\mathrm{}}^dg_f^\beta S_{p,d,\beta }(F).$$
###### Proof.
Put $`F^{}=(g_{\mathrm{}},g_f)F`$, where $`g_f=(g_v)_{v\mathrm{finite}}`$. For each finite place $`v`$, we note that $`K_{v,F^{}}g_vK_{v,F}g_v^1K_{v,𝐆}`$. The index $`[K_{v,𝐆}:K_{v,F^{}}]`$ is therefore bounded above by the number of cosets $`xg_vK_{v,F}`$ in $`K_{v,𝐆}g_vK_{v,F}`$. Clearly this is bounded above by the number of left $`K_{v,F}`$ cosets in $`K_{v,𝐆}g_vK_{v,𝐆}`$; but the number of such cosets is precisely $`g_v[K_{v,𝐆}:K_{v,F}]`$. It now follows easily from the definitions that $`PS_{d,\beta ,f}(F^{})g_{\mathrm{}}^dg_f^\beta PS_{d,\beta ,f}(F)`$. Applying Rem. 8.1 to the endomorphism $`F(g_{\mathrm{}},g_f)F`$, we obtain the claim. ∎
The following crude Lemma is as much of interpolation as we need. It will be applied, in practice, where $`E`$ is a composite of a Hecke operator and a certain $`L^2`$-projection.
###### Lemma 8.3.
Let $`E`$ be a linear endomorphism of $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ which commutes with $`𝐆(F_{\mathrm{}})\times K_{\mathrm{max},𝐆}`$. Suppose there are real numbers $`A,B>0`$ such that for any $`fC_\omega ^{\mathrm{}}(𝐗_𝐆)`$, we have $`Ef||_{L^2}Af_{L^2},Ef_L^{\mathrm{}}Bf_L^{\mathrm{}}`$. Then for $`2p\mathrm{}`$:
$$S_{p,d,\beta }(Ev)A^{2/p}B^{1\frac{2}{p}}S_{p,d,\beta }(v).$$
(We admit also $`B=\mathrm{}`$, in which case the $`L^{\mathrm{}}`$ hypothesis should be seen as void, and the result becomes $`S_{2,d,\beta }(Ev)AS_{2,d,\beta }(v)`$.)
###### Proof.
By interpolation, the operator norm of $`E`$ w.r.t. the $`L^p`$ norm on $`C_\omega ^{\mathrm{}}(𝐗_𝐆)`$ is $`A^{2/p}B^{12/p}`$. Moreover, the assumption on $`E`$ shows that $`K_{Ef}K_f`$.
It follows that for $`fC_\omega ^{\mathrm{}}(𝐗_𝐆)`$ we have the inequality
$$PS_{p,d,\beta }(Ef)A^{2/p}B^{12/p}PS_{p,d,\beta }(f).$$
Rem. 8.1 implies the conclusion. ∎
#### 8.1.1. Computing Sobolev norms in the Kirillov model.
In the present section, let $`v`$ be an archimedean place of $`F`$.
Let $`\pi _v`$ be a generic unitary irreducible representation of $`\mathrm{GL}_2(F_v)`$. Recall that this means that $`\pi _v`$ is realized in a space of functions $`𝒦`$ (the Kirillov model, consisting of restrictions of functions in the Whittaker model to the diagonal torus) on $`F_v^\times `$. Recall also the definition of the local conductor $`\mathrm{Cond}_v(\pi _v)`$ from Sec. 2.12.2.
In this model, the diagonal torus acts by translation and upper triangular matrices act through multiplication by characters: that is to say, for $`f𝒦`$, $`y_1,y_2F_v^\times `$, $`zF_v`$ we have the rules
(8.3)
$$\pi (a(y_1))f:y_2f(y_1y_2),\pi (n(z))f:y_2f(y_2)e_{F_v}(zy_2).$$
From these facts it is easy to verify that the space of smooth vectors in $`\pi _v`$ contains all compactly supported smooth functions on $`F_v^\times `$. Moreover,
(8.4)
$$f_2^2=_{F_v^\times }|f(y)|^2d^\times y$$
defines a $`\mathrm{GL}_2(F_v)`$-invariant inner product on $`𝒦`$.
We will eventually have occasion to choose test vectors in $`\pi _v`$ in this model, and wish to evaluate the “Sobolev norms” of the resulting vectors.
###### Lemma 8.4.
Suppose $`F_v`$. Let $`f𝒦`$ be $`C^{\mathrm{}}`$ and compactly supported. Then
$$\underset{\mathrm{ord}(𝒟)k}{}𝒟f_2\mathrm{Cond}_v(\pi _v)^{2k}\left(\underset{j=0}{\overset{2k}{}}_^\times (|y|+|y|^1)^{2k}\left|\frac{d^jf}{d^jy}\right|^2d^\times y\right)^{1/2},$$
where the $`𝒟`$ sum ranges over all monomials in a fixed basis for $`\mathrm{Lie}(\mathrm{GL}_2(F_v))`$ of degree $`k`$.
Suppose $`F_v`$, and suppose $`f𝒦`$ is $`C^{\mathrm{}}`$ and compactly supported. Then
$$\underset{\mathrm{ord}(𝒟)k}{}𝒟f_2\mathrm{Cond}_v(\pi _v)^k\left(\underset{0i+j2k}{}_^\times (|z|+|z|^1)^{2k}\left|\frac{^{i+j}f}{^iz^j\overline{z}}\right|^2d^\times z\right)^{1/2}$$
###### Proof.
We prove only the case with $`F_v`$, the complex case being similar.
Let $`h,e`$, $`f`$, $`z`$ be nonzero elements of the (real) Lie algebras of $`\mathrm{GL}_2()`$, defined via
$$h=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),e=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),f=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),z=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
These satisfy the usual commutation relations $`[h,e]=2e,[h,f]=2f,[e,f]=h`$. Let $`\lambda `$ be the scalar by which the Casimir operator $`\frac{1}{2}h^2+ef+fe`$ acts, and $`\nu `$ the scalar by which $`z`$ acts; then $`1+|\lambda |+|\nu |^2\mathrm{Cond}_v(\pi _v)^2`$.
It is easy to see how $`h,e`$ act on $`𝒦`$: $`h`$ acts by a multiple of the differential operator $`c_1\nu +y\frac{d}{dy}`$ and $`e`$ acts by multiplication by $`c_3y`$, for some constants $`c_1,c_2,c_3`$. The Casimir operator $`\frac{1}{2}h^2+ef+fe=\frac{1}{2}h^2+2efh`$ acts by the scalar $`\lambda `$; so it follows that for $`v𝒦`$ we have $`efv=\frac{1}{2}(\lambda +hh^2)v`$. In particular, $`e`$ acts on any compactly supported function via the differential operator $`c_1^{}y^1+c_2^{}\frac{d}{dy}+c_3^{}y\frac{d^2}{dy^2}`$, for certain constants $`c_1^{},c_2^{},c_3^{}`$, satisfying $`|c_1^{}|,|c_2^{}|,|c_3^{}|\mathrm{Cond}_v(\pi _v)^2`$. (In fact, $`|c_i^{}|\mathrm{Cond}_v(\pi _v)^{3i}`$.)
Any monomial of degree $`k`$ in $`h,e,f,z`$ is therefore a sum of terms $`c_{\gamma \delta }y^\gamma _y^\delta `$, where $`|c_{\gamma \delta }|\mathrm{Cond}_v(\pi _v)^{2k},|\gamma |k,\delta 2k`$. The claimed result follows in the case $`F_v`$.
A similar proof holds for $`F_v`$. ∎
### 8.2. Reduction theory.
Recall that $`F_{\mathrm{}}:=F_{}`$. Let $`K_{\mathrm{}},K_v,K_{\mathrm{max}}`$ be as in Sec. 2.5. Then $`K_{\mathrm{}}\times K_{\mathrm{max}}`$ is a maximal compact subgroup of $`\mathrm{GL}_2(𝔸_F)`$. Given $`g\mathrm{GL}_2(𝔸_F)`$ we may always write $`g=\left(\begin{array}{cc}1& t\\ 0& 1\end{array}\right)\left(\begin{array}{cc}x& 0\\ 0& y\end{array}\right)k`$, with $`t𝔸_F,x,y𝔸_F^\times ,kK_{\mathrm{}}\times K_{\mathrm{max}}`$. We set $`0pt(g)=|xy^1|_𝔸`$; this is well-defined, although $`x,y`$ are not unique.
Then $`0pt`$ descends to a function $`B(F)\backslash \mathrm{GL}_2(𝔸_F)_{>0}`$. Explicitly,
(8.5)
$$0pt\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\frac{|adbc|_𝔸}{_v(c_v,d_v)^2},$$
where one defines $`(c_v,d_v)=\mathrm{max}(|c_v|_v,|d_v|_v)`$ for $`v`$ finite, and
(8.6)
$$(c_v,d_v)=\left(|c_v|_v^{2/\mathrm{deg}(v)}+|d_v|_v^{2/\mathrm{deg}(v)}\right)^{\mathrm{deg}(v)/2}$$
for $`v`$ infinite, where $`\mathrm{deg}(v)=[F_v:]`$.
Define $`𝔖(T)B(F)\backslash \mathrm{GL}_2(𝔸_F)`$ to be $`𝔖(T):=\{g:0pt(g)T\}`$. Then, for all $`T>0`$ the natural projection $`\mathrm{\Pi }:𝔖(T)\mathrm{GL}_2(F)\backslash \mathrm{GL}_2(𝔸_F)`$ has finite fibers; for sufficiently large $`T`$, it is injective, and for sufficiently small $`T`$ it is surjective. This is the content of reduction theory for $`\mathrm{GL}_2`$. As a consequence, the complement of $`\mathrm{\Pi }(𝔖(T))`$ has compact closure, modulo the center, for each $`T`$.
Fix $`T_0`$ such that $`\mathrm{\Pi }:𝔖(T_0)\mathrm{GL}_2(F)\backslash \mathrm{GL}_2(𝔸_F)`$ is injective. Then we define a function $`0pt:\mathrm{GL}_2(F)\backslash \mathrm{GL}_2(𝔸_F)`$ via the rule
$$0pt(g)=\{\begin{array}{cc}0pt(g^{}),\text{if }g=\mathrm{\Pi }(g^{})\text{ for some }g^{}𝔖(T_0),\hfill & \\ T_0,\text{ else.}\hfill & \end{array}$$
In fact, it is clear that $`0pt`$ descends to a function $`𝐗=\mathrm{PGL}_2(F)\backslash \mathrm{PGL}_2(𝔸_F)_{T_0}`$.
###### Lemma 8.5.
Let $`U\mathrm{GL}_2(F_{\mathrm{}})`$ be compact and $`x𝐗_{\mathrm{GL}(2)}`$. The fibers of the map $`U\times K_{\mathrm{max}}𝐗_{\mathrm{GL}(2)}`$ defined by $`(u,k)xuk`$ have size bounded by $`O(0pt(x))`$, where the implicit constant depends on $`U`$.
###### Proof.
Suppose $`g\mathrm{GL}_2(𝔸_F)`$ is a lift of $`x𝐗_{\mathrm{GL}(2)}`$. Consider the map $`U\times K_{\mathrm{max}}𝐗_{\mathrm{GL}(2)}`$ given by $`(u,k)guk`$, as above. Let $`(u,k)`$ belong to a fiber of maximal size. Call this size $`M`$. Then
(8.7)
$$\begin{array}{c}M=\mathrm{\#}\{\gamma \mathrm{GL}_2(F):guk=\gamma gu^{}k^{},u^{}U,k^{}K_{\mathrm{max}}\}\hfill \\ \hfill \mathrm{\#}\{\gamma :gu^{\prime \prime }k^{\prime \prime }=\gamma g,u^{\prime \prime }UU^1,k^{\prime \prime }K_{\mathrm{max}}\}.\end{array}$$
Set $`V=UU^1`$, a compact subset of $`\mathrm{GL}_2(F_{\mathrm{}})`$. The definition of $`𝔖(T)`$ shows that there exists a constant $`c<1`$, depending on $`V`$, such that $`𝔖(T)VK_{\mathrm{max}}𝔖(cT)`$. Choose $`T`$ so large that the projection $`𝔖(cT)𝐗_{\mathrm{GL}(2)}`$ is injective. It will suffice to show, for each $`g𝔖(T)`$, that
(8.8)
$$\mathrm{\#}\{\gamma \mathrm{GL}_2(F):\gamma ggVK_{\mathrm{max}}\}0pt(g).$$
Both $`g`$ and $`gVK`$ belong entirely to $`𝔖(cT)`$. By the choice of $`T`$, $`\gamma ggVK_{\mathrm{max}}`$ implies $`\gamma `$ in $`B(F)`$. Write $`\gamma =a_\gamma n_\gamma `$, with $`a_\gamma A(F)`$ and $`n_\gamma N(F)`$; also, write $`g=n_ga_gk_g`$ with $`n_gN(𝔸_F),a_gA(𝔸_F),k_gK_{\mathrm{}}\times K_{\mathrm{max}}`$. We are free to adjust $`g`$ on the left by an element of $`N(F)`$, since doing so will not affect the cardinality of the set $`\{\gamma \mathrm{GL}_2(F):\gamma ggVK_{\mathrm{max}}\}`$. We may thereby assume that $`n_g`$ lies in a fixed compact subset of $`N(𝔸_F)`$. Thus we can write $`g=a_gk_g^{}`$, where $`k_g^{}:=a_g^1n_ga_gk_g`$ lies in a certain fixed compact subset $`\mathrm{\Omega }`$ of $`\mathrm{GL}_2(𝔸_F)`$.
Now, $`\gamma ggVK`$ implies that $`a_g^1a_\gamma n_\gamma a_g\mathrm{\Omega }VK_{\mathrm{max}}\mathrm{\Omega }^1`$. Noting that $`a_g^1a_\gamma n_\gamma a_g=a_\gamma a_g^1n_\gamma a_g`$, we deduce that $`a_\gamma `$ lies in a fixed compact subset of $`\mathrm{GL}_2(𝔸_F)`$, depending only on $`U`$; thus the number of possibilities for $`a_\gamma `$ are $`_U1`$. Moreover, it now follows that $`a_g^1n_\gamma a_g`$ lies in a compact subset of $`𝔸_F`$ depending only on $`U`$.
Thus, if we write $`a_g=\left(\begin{array}{cc}x& 0\\ 0& y\end{array}\right),n_\gamma =\left(\begin{array}{cc}1& \beta \\ 0& 1\end{array}\right)`$, then $`\beta xy^1\mathrm{\Omega }^{}`$, where $`\mathrm{\Omega }^{}𝔸_F`$ is a compact subset that depends only on $`U`$. It is easy to see that the number of possibilities for $`\beta `$ is $`_U1+|xy^1|_{𝔸_F}`$. But $`|xy^1|_{𝔸_F}=0pt(g)`$, which is a function that is bounded away from zero, and we are done. ∎
###### Lemma 8.6.
Let notations be as in the previous Lemma 8.5. Consider the composite map $`UK_{\mathrm{max}}\stackrel{\mathrm{\Pi }}{}𝐗_{\mathrm{GL}(2)}𝐗`$. Each fiber of this map may be written as the union of at most $`O(0pt(x))`$ sets each of the form $`yZ(𝔸_F)UK_{\mathrm{max}}`$, where $`Z`$ is the center of $`\mathrm{GL}_2`$ and $`y\mathrm{GL}_2(𝔸_F)`$.
###### Proof.
Let $`\overline{x}`$ be the image of $`x`$ in $`𝐗`$. Let $`u,u^{}U,k,k^{}K_{\mathrm{max}}`$. Suppose that $`\overline{x}uk=\overline{x}u^{}k^{}`$ in $`𝐗`$. Then there is $`z𝔸_F^\times `$ and $`\gamma \mathrm{GL}_2(F)`$ such that
(8.9)
$$xuk=\gamma xu^{}k^{}a(z,z),\text{ equality in }\mathrm{GL}_2(𝔸_F)$$
For fixed $`u,k`$ and $`\gamma `$, the set of $`u^{}k^{}`$ satisfying (8.9) is visibly the intersection of $`UK_{\mathrm{max}}`$ with a fixed $`Z(𝔸_F)`$-coset. This coset depends only on the class of $`\gamma `$ in $`\mathrm{PGL}_2(F)`$, so it suffices to show that those $`\gamma \mathrm{GL}_2(F)`$ that occur in equalities such as (8.9) for varying $`u,k,u^{},k^{}`$ represent at most $`O(0pt(x))`$ distinct cosets $`\gamma Z(F)`$ in $`\mathrm{PGL}_2(F)`$.
Taking determinant followed by the norm $`𝔸_F^\times /F^\times `$, we conclude that $`|z|_𝔸`$ belongs to a compact subset of $`^\times `$ that depends only on $`U`$. The norm map $`𝔸_F^\times /F^\times ^\times `$ being proper, it follows that $`z`$ itself belongs to a compact subset $`𝔸_F^\times /F^\times `$ that depends only on $`U`$.
In particular, there is a compact subset $`\mathrm{\Omega }F_{\mathrm{}}^\times `$, depending only on $`U`$, and a finite subset $`P𝔸_F^\times `$, containing $`1`$ and also depending only on $`U`$, such that $`zF^\times \mathrm{\Omega }.P._{v\mathrm{finite}}𝔬_{F,v}^\times `$. Let $`\stackrel{~}{U}=U\{a(z_{\mathrm{}},z_{\mathrm{}}):z_{\mathrm{}}\mathrm{\Omega }\}`$. Given a solution to (8.9), write $`z=\delta z_{\mathrm{}}po`$, with $`\delta F^\times ,z_{\mathrm{}}\mathrm{\Omega },pP,o_{v\mathrm{finite}}𝔬_{F,v}^\times `$. Then
$$xuk=\gamma a(\delta ,\delta )xa(p,p)u^{}a(z_{\mathrm{}},z_{\mathrm{}})k^{}a(o,o),$$
in particular, taking $`\stackrel{~}{u}=u^{}a(z_{\mathrm{}},z_{\mathrm{}})\stackrel{~}{U},k^{\prime \prime }=k^{}a(o,o)K_{\mathrm{max}}`$, the image of $`xa(p,p)\stackrel{~}{u}k^{\prime \prime }`$ in $`𝐗_𝐆`$ coincides with $`xuk`$. So the number of possibilities for the $`Z(F)`$-coset of $`\gamma `$ is bounded above by the fibers of the map $`P\times \stackrel{~}{U}\times K_{\mathrm{max}}𝐗_𝐆`$ given by $`(p,\stackrel{~}{u},k)xa(p,p)\stackrel{~}{u}k`$. The result follows from Lem. 8.5.
We shall now need a quantitative version of certain statements in reduction theory. The subsequent Lemma is a fancier version of the following statement: the number of $`\gamma \mathrm{SL}(2,)`$ that map a fixed $`z`$ to the Siegel set
$$\{x+iy:0x1,yT\}$$
is $`1+T^1.`$
###### Lemma 8.7.
Let $`g\mathrm{GL}_2(𝔸_F)`$ and $`Y>0`$ a positive real number. Then
(8.10)
$$\mathrm{\#}\{\gamma B(F)\backslash \mathrm{GL}_2(F):0pt(\gamma g)Y\}_ϵ1+Y^{1ϵ}.$$
Here the implicit constant is independent of $`g`$. Moreover, suppose $`g𝔖(T)`$ with $`T1`$. Then:
(8.11)
$$\mathrm{sup}\{0pt(\gamma g):\gamma B(F)\}T^1.$$
###### Proof.
The proof of (8.10) is not difficult, generalizing in a straightforward way the proof with $`F=`$. However, it is somewhat notationally tedious; the (hypothetical) reader may wish to simply work out the proof for $`F=`$, where it is equivalent to the following fact: the number of primitive vectors in a unimodular sublattice of $`^2`$ that are contained in an $`R`$-ball is $`(1+R^2)`$, uniformly in the lattice. (The result also be deduced if one admits some basic facts from the theory of Eisenstein series over $`F`$, but we wish to rather deduce these basic facts from the present Lemma). We also remark that the entire content of (8.10) lies in the uniformity in $`g`$.
Without loss of generality, we take $`g𝔖(T_0)`$, where $`T_0`$ is sufficiently small that the map $`𝔖(T_0)𝐗_{\mathrm{GL}(2)}`$ is surjective. So $`g=\left(\begin{array}{cc}1& t\\ 0& 1\end{array}\right)\left(\begin{array}{cc}x& 0\\ 0& y\end{array}\right)k`$ with $`|xy^1|_𝔸T_0`$. Moreover, replacing $`g`$ by $`gz`$, for any $`zZ(𝔸_F)`$ does not affect the problem, so we may take $`y=1`$. Then, for $`\gamma =\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{GL}_2(F)`$, we have
(8.12)
$$0pt(\gamma g)=\frac{|x|_𝔸}{_v(x_vc,ct_v+d)_v^2},$$
The equivalence class of $`\gamma `$ in $`B(F)\backslash \mathrm{GL}_2(F)`$ depends only the pair $`(c,d)F^2`$, considered up to $`F^\times `$ equivalence (i.e., it depends only on $`c/dF\{\mathrm{}\}`$.) It suffices, then, to estimate the number
(8.13)
$$\mathrm{\#}\{[c:d]^1(F),\underset{v}{}x_vc,ct_v+d_v^2Y^1|x|_𝔸\},$$
If $`\mathrm{\Omega }`$ is any fixed compact subset of $`𝔸_F^\times `$, then for $`\omega \mathrm{\Omega }`$ we have $`_vx_v\omega c,ct_v+d_v_\mathrm{\Omega }_vx_vc,ct_v+d_v.`$ Consider $`_{>0}`$ as embedded in $`𝔸_F^\times `$ via $`_{>0}𝔸_{}^\times 𝔸_F^\times `$. Then there is a compact subset $`\mathrm{\Omega }𝔸_F^\times `$ such that $`𝔸_F^\times =F^\times \mathrm{\Omega }_{>0}`$.
The size of (8.13) is unaffected by the substitution $`(x,t)(x\tau ,t\tau )`$, for any $`\tau F^\times `$. In view of the above remarks we may assume – decreasing $`Y`$ by a constant that depends only on $`F`$ – that $`x_{>0}`$. Moreover, the size of (8.13) is also unaffected by the substitution $`tt+\tau `$, for $`\tau F`$. We may therefore assume that $`|t|_v1`$ for all finite places $`v`$.
Fix a set of representatives $`𝔍_1,\mathrm{},𝔍_h`$ for the class group of $`𝔬_F`$; we will assume each $`𝔍_i`$ is integral. For any $`[c:d]^1(F)`$, we may find a representative $`(c,d)`$ so that the ideal $`c𝔬_F+d𝔬_F`$ is one of the $`𝔍_i`$; moreover, replacing $`(c,d)𝔍_i^2`$ by $`(ϵc,ϵd)`$ for $`ϵ𝔬_F^\times `$ does not change the class $`[c:d]`$.
The restrictions on $`x,t`$ imply that $`(x_vc,ct_v+d)_v=(c,d)_v`$ for all finite $`v`$. Then $`_{v\mathrm{finite}}(c,d)_v=\mathrm{N}(𝔍_i)^1`$, the inverse of the norm of $`𝔍_i=c𝔬_F+d𝔬_F`$. So it will suffice to bound, for each $`1ih`$, the quantity
$$\mathrm{\#}\{(c,d)𝔍_i^2/𝔬_F^\times :c𝔬_F+d𝔬_F=𝔍_i:\underset{\mathrm{}|v}{}(|x_v|^2|c|_v^2+|ct_v+d|_v^2)Y^1|x|_𝔸\mathrm{N}(𝔍)_i^2\}.$$
Since $`𝔍_i`$ belongs to a finite set, the quantity $`\mathrm{N}(𝔍)`$ is bounded; thus, decreasing $`Y`$ again as necessary, it suffices to estimate for each $`1ih`$
$$\mathrm{\#}\{(c,d)𝔍_i^2/𝔬_F^\times :c𝔬_F+d𝔬_F=𝔍_i,\underset{\mathrm{}|v}{}(x_vc,ct_v+d)_v^2Y^1|x|_𝔸\}$$
There is only one term corresponding to $`c=0`$. Otherwise, $`(c)`$ is a principal ideal divisible by $`𝔍_i`$; let $`𝒫`$ be the set of integral principal ideals. Then the size of the set above is precisely
(8.14)
$$\underset{(c)𝒫}{}\mathrm{\#}\{d𝔍_i:(c)+d𝔬_F=𝔍_i,\underset{\mathrm{}|v}{}(x_vc,ct_v+d)_v^2Y^1|x|_𝔸\}$$
We note that the size of the inner set is independent of the choice of generator for the principal ideal $`(c)`$. Moreover, the inequality of (8.14) implies that the norm $`\mathrm{N}((c))`$ of the principal ideal $`(c)`$ satisfies $`\mathrm{N}((c))^2Y^1|x|_𝔸^1`$.
Let us estimate the number of $`d`$ that can correspond to a fixed principal ideal $`(c)`$ in (8.14). Recall that $`|x|_𝔸T_0`$ and that $`x_𝔸`$ is in the image of the embedding $`_{>0}𝔸_Q^\times 𝔸_F^\times `$. In particular, $`|x|_v`$ is bounded below at each infinite place. Moreover, since $`𝔬_F^\times `$ is a cocompact subgroup of the elements of $`F_{\mathrm{}}^\times `$ with norm $`1`$, we can choose a representative for the principal ideal $`(c)`$ so the same is true of $`|c|_v`$. Note that (cf. 8.6) that $`(x_vc,ct_v+d)_v(|x_vc|_v+|ct_v+d_v|_v)`$.
So in fact, again decreasing $`Y`$ as necessary, it will suffice to estimate
(8.15)
$$\underset{(c)𝒫:\mathrm{N}(c)Y^{1/2}|x|_𝔸^{1/2}}{}\mathrm{\#}\{d𝔍_i:\underset{\mathrm{}|v}{}(1+|ct+d|_v)^2Y^1|x|_𝔸\}$$
To estimate the right-hand side, first observe that if $`\{M_v\}_{\mathrm{}|v}`$ is any set of positive real numbers indexed by the infinite places of $`F`$, then $`\mathrm{\#}\{d𝔍_i:|ct+d|_vM_v\text{ for }\mathrm{}|v\}_{\mathrm{}|v}(1+M_v)`$. Indeed, by subtraction, it will suffice to estimate $`\mathrm{\#}\{d𝔍_i:|d|_v2M_v\text{ for }\mathrm{}|v\}`$; this amounts to counting points in the lattice $`𝔍_iF_{\mathrm{}}`$ in a region that is the product of a box and a disc; the result is then clear.
Next, if $`T1`$, then the subset $`\{(y_1,\mathrm{},y_d):_i(1+y_i)T\}`$ in $`_{>0}^d`$ is contained in the union of $`O_ϵ(T^ϵ)`$ boxes $`\{(y_1,\mathrm{},y_d):y_iM_i\}`$, where $`_i(1+M_i)T`$. We may assume $`Y^1|x|_𝔸1`$, else (8.15) has no solutions. We conclude that the number of $`d`$ attached to each principal ideal $`(c)`$ in (8.15) is $`_ϵ(Y^{1/2}|x|_𝔸^{1/2})^{1+ϵ}`$.
The number of possibilities for $`(c)`$ is bounded by the number of integral ideals with norm $`Y^{1/2}|x|_𝔸^{1/2}`$, which is $`_ϵ(Y^{1/2}|x|_𝔸^{1/2})^{1+ϵ}`$. Finally there is one class with $`c=0`$. We conclude that the number of pairs $`(c,d)`$ up to equivalence is $`Y^{1ϵ}+1`$. This proves (8.10).
As for (8.11), suppose $`g𝔖(T)`$, so we may write $`g=\left(\begin{array}{cc}x& z\\ 0& y\end{array}\right)k`$ with $`kK_{\mathrm{}}\times K_{\mathrm{max}}`$, and $`|xy^1|_𝔸T`$. Suppose $`\gamma =\left(\begin{array}{cc}\alpha & \beta \\ \alpha ^{}& \beta ^{}\end{array}\right)`$. If $`\gamma B(F)`$, then $`\alpha ^{}0`$. In that case, following the notation of (8.5), we have:
$$\underset{v}{}(\alpha _v^{},\beta _v^{})g_v\underset{v}{}|\alpha _v^{}x_v|_v=|x|_𝔸.$$
and therefore, by (8.5), $`0pt(\gamma g)|det(g)x^2|_𝔸T^1`$. ∎
## 9. Background on quantitative equidistribution results.
The aim of this section is to quantify various standard equidistribution results (equidistribution of long horocycles, Hecke points, etc.), using the adelic Sobolev norms. As such neither the results nor the methods are new; we just collect together those results we need and provide brief proofs.
As regards the origin of the ideas used here, we have drawn in particular from the work of Clozel-Ullmo, Linnik, Oh, Margulis, Ratner and Sarnak.
### 9.1. Decay of matrix coefficients.
#### 9.1.1. Local setting.
Our fundamental tool in establishing all these results is the spectral gap, i.e., quantitative mixing properties of real and $`p`$-adic flow. As such, we begin by recalling the basic relevant bound on matrix coefficients.
Let $`0\alpha 1/2`$. Let $`v`$ be a place of $`F`$, and suppose that $`(V,\pi )`$ is a unitary representation of $`\mathrm{GL}_2(F_v)`$ which does not contain, in its spectral decomposition, any complementary series with parameter $`\alpha `$. (More formally: $`V`$ does not weakly contain such a representation). Thus $`\alpha =0`$ corresponds to $`V`$ being tempered, and $`\alpha =1/2`$ corresponds to $`V`$ having no almost invariant vectors.
Then for $`w_1,w_2`$ any two $`K_v`$-finite elements of $`V`$, satisfying $`w_1,w_1=w_2,w_2=1`$, and any $`xF_v`$ we have the bound on matrix coefficients given by
(9.1)
$$\pi (a(x)w_1,w_2_{ϵ,F}dim(K_vw_1)^{1/2}dim(K_vw_2)^{1/2}(1+|x|_v)^{\alpha 1/2+ϵ}.$$
The implicit constant of (9.1) depends only on $`ϵ`$. Since we do not know of an available reference, we briefly sketch an argument for (9.1). In the case where $`\alpha =0`$, i.e. $`V`$ is tempered, then (9.1) is proven in . In the general case, let $`(\sigma _{1/2\alpha },W)`$ be the complementary series with parameter $`1/2\alpha `$; let $`v^0W`$ be a unitary spherical vector. Then the representation $`VW`$ is tempered. Indeed it suffices – again by – to verify that a dense set of matrix coefficients are in $`L^{2+ϵ}`$, which follows by direct computation. Now one may estimate the matrix coefficient $`a(x)w_1v^0,w_2v^0`$ by appealing again to . On the other hand $`a(x)w_1v^0,w_2v^0=a(x)w_1,w_2a(x)v^0,v^0`$, and an easy computation shows that $`a(x)v^0,v^0_ϵ(1+|x|_v)^{\alpha ϵ}`$. Thus (9.1) follows. (This argument is a variant of an argument that appears at the end of .)
Let us record a useful further variant. Suppose $`v`$ is finite. Let $`K_1,K_2K_v`$ be subgroups and let $`\sigma `$ be the $`(K_1,K_2)`$-bi-invariant probability measure supported on $`K_1a(x)K_2`$.
Then
(9.2)
$$v\sigma _2[K_v:K_1]^{1/2}[K_v:K_2]^{1/2}(1+|x|_v)^{\alpha 1/2+ϵ}v_2.$$
Indeed, for $`i=1,2`$ let $`\mathrm{\Pi }_{K_i}`$ be the projection operator $`w_{K_i}kw`$ on $`V`$, where $`K_i`$ is endowed with the Haar probability measure. Then:
(9.3)
$$\begin{array}{c}v\sigma _2=\underset{wV}{sup}\frac{v\sigma ,w}{w_2}\hfill \\ \hfill =\underset{wV}{sup}\frac{a(x)\mathrm{\Pi }_{K_1}v,\mathrm{\Pi }_{K_2}w}{w_2}[K_v:K_1]^{1/2}[K_v:K_2]^{1/2}(1+|x|_v)^{\alpha 1/2+ϵ}v_2\end{array}$$
#### 9.1.2. Variant for $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$.
Let $`0\alpha 1/2`$, suppose $`G=\mathrm{SL}_2()`$, and let $`V`$ be a unitary representation of $`G`$ such that $`V`$ does not weakly contain any complementary series with parameter $`\alpha `$. The normalization is again so that $`\alpha =0`$ corresponds to tempered and $`\alpha =1/2`$ corresponds to $`V`$ not having almost invariant vectors.
Then one has the following variant of (9.1), proved by the same method:
(9.4)
$$\pi (\left(\begin{array}{cc}y^{1/2}& 0\\ 0& y^{1/2}\end{array}\right)w_1,w_2_ϵdim(\mathrm{SO}(2)w_1)^{1/2}dim(\mathrm{SO}(2)w_2)^{1/2}(1+|y|)^{\alpha 1/2+ϵ}.$$
It is convenient to extend the validity of (9.4) beyond the $`K`$-finite space by replacing $`dim(\mathrm{SO}(2)w_i)`$ by appropriate Sobolev norms. We confine ourselves to the case of main interest, where $`V`$ is the orthogonal complements of the constants in $`L^2(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$, where $`\mathrm{\Gamma }`$ is a lattice in $`\mathrm{SL}_2()`$. The estimates we are about to describe are, again, not new; estimates for effective mixing of geodesic and horocycle flows in this setting are contained in .
For our purposes it would be optimal to use fractional Sobolev norms; since we have not defined these, we shall use a rather crude form of interpolation instead.
Thus let $`f_1,f_2C^{\mathrm{}}(\mathrm{\Gamma }\backslash \mathrm{SL}_2())`$. One expands both $`f_1`$ and $`f_2`$ into a sum of $`\mathrm{SO}(2)`$-types and applies (9.4). Indeed, write for $`i\{1,2\}`$ an expansion $`f_i=_{n=\mathrm{}}^{\mathrm{}}f_i^{(n)}`$, where $`f_i^{(n)}`$ transforms under the character $`\left(\begin{array}{cc}\mathrm{cos}(\theta )& \mathrm{sin}(\theta )\\ \mathrm{sin}(\theta )& \mathrm{cos}(\theta )\end{array}\right)e^{in\theta }`$. Expanding:
(9.5)
$$\begin{array}{c}\left(\begin{array}{cc}y^{1/2}& 0\\ 0& y^{1/2}\end{array}\right)f_1,f_2=\underset{n,m}{}\left(\begin{array}{cc}y^{1/2}& 0\\ 0& y^{1/2}\end{array}\right)f_1^{(n)},f_2^{(m)}\hfill \\ \hfill _ϵ(1+|y|)^{\alpha 1/2+ϵ}\underset{n,m}{}f_1^{(n)}_2f^{(m)}_2=(1+|y|)^{\alpha 1/2+ϵ}\left(\underset{n}{}f_1^{(n)}_2\right)\left(\underset{m}{}f_2^{(m)}_2\right)\end{array}$$
Our definitions of the Sobolev norms (Sec. 2.9.2) are so that $`S_{2,1}(f_1)^2_n(1+|n|)^2f_1^{(n)}_2^2`$, and similarly for $`f_2`$. On the other hand, it is an elementary estimate that
$$\left(\underset{n}{}f_1^{(n)}_2\right)^2_ϵ\left(\underset{n}{}f_1^{(n)}_2^2(1+|n|)^2\right)^{1/2+ϵ}\left(\underset{n}{}f_1^{(n)}_2^2\right)^{1/2ϵ}$$
It follows from this that for any $`k,k^{}\mathrm{SO}(2)`$ we have the matrix coefficient bound:
(9.6)
$$\begin{array}{c}|k\left(\begin{array}{cc}y^{1/2}& 0\\ 0& y^{1/2}\end{array}\right)k^{}f_1,f_2|\hfill \\ \hfill (1+|y|)^{\alpha 1/2+ϵ}(S_{2,1}(f_1)S_{2,1}(f_2))^{1/2+ϵ}f_1^{1/2ϵ}f_2^{1/2ϵ},\end{array}$$
at least for $`f_1,f_2`$ which are $`\mathrm{SO}(2)`$-finite. But the general case of smooth $`f_1,f_2`$ follows from density.
Note that in (9.6) that the factor $`f_1^{1/2ϵ}S_{2,1}(f_1)^{1/2+ϵ}`$ is a crude substitute for the fractional ($`1/2+ϵ`$-) Sobolev norm of $`f_1`$.
### 9.2. Pointwise bounds.
In this section, we make free use of the adelic Sobolev norms introduced in Sec. 2.9.3. We recall the definition $`S_{p,d}:=S_{p,d,1/p}`$. We also recall that in statements of the form $`|L(f)|S_{p,d}(f)`$, for certain linear functionals $`L`$, we shall allow the implicit constant of $``$ to depend on $`p`$ and $`d`$ without explicit mention.
###### Lemma 9.1.
Let $`fC_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$ and let $`x𝐗_{\mathrm{GL}(2)}`$. Then, for any $`p2`$ and $`d1`$,
(9.7)
$$|f(x)|0pt(x)^{1/p}S_{p,d}(f)$$
Moreover, if $`FC^{\mathrm{}}(𝐗\times 𝐗)`$, and $`p>2,d1`$,
(9.8)
$$_𝐗|F(x,x)|𝑑xS_{p,d}(F)$$
###### Proof.
As in (2.1), set $`K_f=_{v\mathrm{finite}}K_{v,f}`$, where $`K_{v,f}`$ is the stabilizer of $`f`$ in $`K_v`$. Fix an open neighbourhood of the identity $`U\mathrm{GL}_2(F_{\mathrm{}})`$. Consider the map $`\mathrm{\Pi }:UK_f𝐗`$ defined by $`(u,k)xuk`$. By Lem. 8.6, the fibers are unions of at most $`O(0pt(x))`$ sets, each of the form $`yZ(𝔸_F)UK_f`$. Moreover, for any $`yUK_f`$, the measure of $`\{z𝔸_F^\times :ya(z,z)UK_f\}`$ is bounded above by a constant depending only on $`U`$. Indeed, the set of such $`z`$ is contained in a fixed compact subset of $`𝔸_F^\times `$ that depends only on $`U`$.
Equip $`UK_f`$ with the restriction of Haar measure from $`\mathrm{GL}_2(𝔸_F)`$. From the preceding paragraph, one easily deduces that the push-forward of this measure to $`𝐗`$, under $`(u,k)xuk`$, is bounded above by $`C0pt(x)`$ times the measure on $`𝐗`$, where the constant $`C`$ depends only on $`U`$.
Then:
(9.9)
$$\begin{array}{c}_{uU}|f(xu)|^p=\mathrm{vol}(K_f)^1_{uU,kK_f}|f(xuk)|^p𝑑u𝑑k\hfill \\ \hfill 0pt(x)[K_{\mathrm{max}}:K_f]_𝐗|f(x)|^pd\mu _𝐗(x)\end{array}$$
(9.9) holds with $`f`$ replaced by $`𝒟f`$, for $`𝒟`$ any fixed monomial in $`\mathrm{Lie}(\mathrm{GL}_2(F_{\mathrm{}}))`$. The standard Sobolev estimate, applied to the function $`uf(xu)`$ on the real manifold $`U`$, implies that $`|f(x)|0pt(x)^{1/p}PS_{p,d,1/p}(f)`$ for sufficiently large $`d`$. (Indeed, it suffices to take any $`d>\mathrm{dim}(U)/2=2[F:]`$.) Then Remark 8.1, (2) implies the conclusion.
As for the second conclusion, we proceed in a similar fashion as above (with $`𝐗`$ replaced by $`𝐗\times 𝐗`$) to obtain the estimate $`|F(x,y)|0pt(x)^{1/p}0pt(y)^{1/p}S_{p,d}(F)`$. It is easy to see that $`_𝐗0pt(x)^{2/p}𝑑x<\mathrm{}`$ for $`p>2`$, and the conclusion follows. ∎
The next lemma quantifies the rapid decay of a cuspidal function, or more generally a truncated automorphic function, in the cusp. Recall that for $`T_0>0`$ we have defined the Siegel domain $`𝔖(T_0)`$ in Section 8.2.
###### Lemma 9.2.
Let $`fC_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$. Put $`f^N(g)=_{N(F)\backslash N(𝔸_F)}f(ng)𝑑n`$, where the measure on $`N(F)\backslash N(𝔸_F)`$ is the $`N(𝔸_F)`$-invariant probability measure. Then for $`x𝔖(T_0)`$, $`p2,k0`$ and $`d1`$,
(9.10)
$$|f(x)f^N(x)|_{T_0}0pt(x)^{1/pk}S_{p,d,1/p+k}(f)$$
###### Proof.
We may assume that $`xN(𝔸_F)A()\mathrm{\Omega }(K_{\mathrm{}}\times K_{\mathrm{max}})`$, for some fixed compact set $`\mathrm{\Omega }A(𝔸_F)`$. Here $`A()`$ is regarded as a subset of $`A(F_{\mathrm{}})`$ via the natural inclusion $`F_{\mathrm{}}`$. Write accordingly $`x=na\omega k`$, where $`\omega \mathrm{\Omega }`$.
Consider the function on $`F_{\mathrm{}}`$ defined by $`g(t)=f(n(t)x)f^N(x)`$. It is invariant by the lattice $`\mathrm{\Lambda }=\{t𝔬_F:n(t)\mathrm{GL}_2(F_{\mathrm{}})\omega K_f\omega ^1\}`$, , where $`K_f`$ is again as in $`(\text{2.1})`$. One sees that, since $`\omega `$ belongs to the fixed compact $`\mathrm{\Omega }`$, the covolume bound $`\mathrm{vol}(F_{\mathrm{}}/\mathrm{\Lambda })[K_{\mathrm{max}}:K_f]`$. Moreover, since $`\mathrm{\Lambda }`$ may be regarded as a fractional ideal of $`𝔬_F`$, the homothety class of $`\mathrm{\Lambda }`$ lies in a fixed compact set in the space of homothety classes of lattices in $`F_{\mathrm{}}`$. Also, $`g(t)`$ defines a function on $`F_{\mathrm{}}/\mathrm{\Lambda }`$, with integral $`0`$.
Suppose now that $`G`$ is a smooth function on $`^d/L`$, for some $`d>1`$ and some lattice $`L^d`$, with integral $`0`$. Let $`G_{(i)}=sup_{z^d/L}|𝒟G|`$, where $`𝒟`$ varies over all monomials in $`_1,\mathrm{},_d`$ of exact order $`i`$. Then an elementary argument shows that $`G_{(0)}\mathrm{vol}(^d/L)^{i/d}G_{(i)}`$, and the implicit constant may be taken to vary continuously with the homothety class of $`L`$.
Apply this lemma to the function $`g`$ on $`F_{\mathrm{}}/\mathrm{\Lambda }`$, with $`i=k[F:]`$ for some $`k1`$. The norm $`g_{(k[F:])}`$, in the sense of the above paragraph, is bounded, by Lem. 9.1 and an elementary computation, by $`0pt(x)^k(0pt(x)^{1/p}PS_{p,d^{}}(f))`$, for some $`d^{}1`$. It follows that
$$sup|g(t)|0pt(x)^{1/pk}PS_{p,d^{}}(f)[K_{\mathrm{max}}:K_f]^k=0pt(x)^{1/pk}PS_{p,d^{},1/p+k}(f).$$
Applying Rem. 8.1, (2), we conclude $`|f(x)f^N(x)|0pt(x)^{1/pk}S_{p,d^{},1/p+k}(f)`$. ∎
###### Lemma 9.3.
Suppose $`fC_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$ is cuspidal. Then $`S_{\mathrm{},d,\beta }(f)S_{2,d^{},\beta +3/2}(f)`$, for sufficiently large $`d^{}`$.
###### Proof.
By Lem. 9.2 for $`f`$ cuspidal, applied with $`p=2,k=1`$, we see that $`|f(x)|S_{2,d,3/2}(f)`$ for $`d1`$. Applying this inequality to $`𝒟f`$, for $`𝒟`$ in the universal enveloping algebra of $`\mathrm{GL}_2(F_{\mathrm{}})`$, we see that $`PS_{\mathrm{},d,\beta }(f)PS_{2,d^{},\beta +3/2}(f)`$, for $`d^{}`$ sufficiently large. This equality holds for cuspidal $`f`$.
Let $`\mathrm{\Pi }`$ be the $`L^2`$-orthogonal projection onto the space of cuspidal functions; then $`\mathrm{\Pi }`$ commutes with $`\mathrm{GL}_2(𝔸_F)`$, and it follows that
$$PS_{\mathrm{},d,\beta }(\mathrm{\Pi }f)PS_{2,d^{},\beta +3/2}(\mathrm{\Pi }f)PS_{2,d^{},\beta +3/2}(f)$$
for arbitrary $`fC_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$. Now Rem. 8.1, (3) (or, more precisely, a trivial modification thereof) implies the conclusion. ∎
### 9.3. Equidistribution of long horocycles and closed horospheres.
Let $`G`$ be a semisimple group, $`\mathrm{\Gamma }G`$ a lattice, $`U`$ a unipotent subgroup of $`G`$. It is well-known that one can prove, in a quantitative fashion, the equidistribution of $`U`$-orbits on $`\mathrm{\Gamma }\backslash G`$ if $`U`$ is a horospherical subgroup, i.e. the unipotent radical of a proper parabolic subgroup. We shall quantify two instances of this that will be of interest to us.
We emphasize that neither the results nor the techniques of this section are new; we have included proofs only to keep the present paper as self-contained as possible.
Effective estimates for equidistribution of long horocycles on quotients of $`\mathrm{SL}_2()`$ are already implicit in the work of Ratner and , where the effective mixing of the horocycle flow is used. We will proceed in a closely related fashion, using the mixing property of the Cartan action; again, this is definitely not new and appears already, although in a different context, in the doctoral thesis of Margulis (reprinted in ).
#### 9.3.1. Equidistribution of long horocycles in hyperbolic $`2`$-space.
Let $`\mathrm{\Gamma }\mathrm{SL}(2,)`$ be a lattice such that $`L^2(\mathrm{\Gamma }\backslash \mathrm{SL}(2,))`$ does not contain any complementary series representation with parameter $`\alpha `$, for any $`0\alpha <1/2`$. (That is: $`\alpha [0,1/2)`$ is such that all nonzero eigenvalues of the hyperbolic Laplacian $`y^2(_{xx}+_{yy})`$ on $`\mathrm{\Gamma }\backslash ^2`$ are bounded below by $`1/4\alpha ^2`$).
We define $`n,a,\overline{n}`$ as in (3.1).
The following Lemma quantifies the equidistribution of long horocycles. Results of this type are already implicit in and . This problem is analyzed in much more detail than we go into, in and .
###### Lemma 9.4.
Assume $`\mathrm{\Gamma }`$ is cocompact, and let $`x_0\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$.
(9.11)
$$\left|\frac{1}{T}_{t=0}^Tf(x_0n(t))𝑑t_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g\right|_ϵT^{\frac{\alpha 1/2}{2}+ϵ}S_{\mathrm{},1}(f).$$
###### Proof.
The idea (which is certainly not new – cf. remarks at start of Section 9.3) is that, upon flowing a small ball in $`\mathrm{\Gamma }\backslash G`$ for a long time by the geodesic flow, it turns into a narrow neighbourhood of a long horocycle. One thereby can deduce the equidistribution of the long horocycle from the mixing properties of the geodesic flow.
Let $`N,A,\overline{N}`$ be the images of $`n,a,\overline{n}`$ respectively. Let $`g_1`$ be a smooth function of compact support on the real line, with integral $`_{\mathrm{}}^{\mathrm{}}g_1(x)𝑑x=1`$. It will remain fixed for all time throughout our arguments. Fix $`1>\delta >0`$ and let $`g_\delta :`$ be the convolution of the characteristic function of $`[0,1]`$ with $`g_1(x/\delta )\delta ^1`$; that is to say
$$g_\delta (x)=\delta ^1_{t=0}^1g_1(\frac{xt}{\delta })𝑑t.$$
Then $`g_\delta `$ is a smooth function of integral $`1`$, which is supported in a small interval around $`[0,1]`$.
Define a probability measure $`\mu _\delta `$ on $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$ via the rule
$$\mu _\delta (f)=\delta ^1_{x,y,z}f(x_0n(x)a(e^y)\overline{n}(z))g_\delta (x)g_1(y/\delta )g_1(z)𝑑x𝑑y𝑑z.$$
In words, $`\mu _\delta `$ is a measure supported on a small box around $`x_0`$; this box has width $`O(1)`$ in the $`N`$ and $`\overline{N}`$ directions, and $`O(\delta )`$ in the $`A`$ direction. When we flow this by $`A`$, it will become a measure supported along a box that closely approximates an $`N`$-orbit.
We observe that
$$\mu _\delta (a(T^1)f)=\frac{1}{\delta }_{x,y,z}f(x_0a(T)^1n(x)a(e^y)\overline{n}(z))g_\delta (x/T)g_1(y/\delta )g_1(Tz)𝑑x𝑑y𝑑z.$$
On the other hand, for any fixed $`x_1\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$, we note that
(9.12)
$$\left|\delta ^1Tf(x_1a(e^y)\overline{n}(z))g_1(y/\delta )g_1(Tz)𝑑y𝑑zf(x_1)\right|\mathrm{max}(T^1,\delta )S_{\mathrm{},1}(f).$$
Indeed, (9.12) merely quantifies the fact that the right-hand side integral is against a probability measure supported in a very small ball (of size $`\mathrm{min}(\delta ,T^1)`$) around $`x_1`$. Since $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$ is assumed compact, the implicit constant of (9.12) may be taken independent of $`x_1`$.
Consequently,
(9.13)
$$\left|\frac{1}{T}_tg_\delta (t/T)f(x_0a(T)^1n(t))𝑑t\mu _\delta (a(T)^1f)\right|\mathrm{max}(T^1,\delta )S_{\mathrm{},1}(f).$$
On the other hand, the measure $`\mu _\delta `$ has a continuous distribution function $`h_\delta `$, i.e. $`\mu _\delta (f)=_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)h_\delta (g)𝑑g`$, and $`\mu _\delta (a(T)^1f)`$ may be estimated using (9.6), i.e. the decay of matrix coefficients.
A routine computation shows that $`h_\delta _{L^2}\delta ^{1/2}`$ and $`S_{2,1}(h_\delta )\delta ^{3/2}`$; on account of the cocompactness of $`\mathrm{\Gamma }\backslash \mathrm{SL}_2()`$, both these are estimates are uniform in $`x_0`$.
Using (9.6) now yields:
(9.14)
$$|\mu _\delta (a(T)^1f)_{\mathrm{\Gamma }\backslash \mathrm{SL}_2()}f(g)𝑑g|_ϵT^{\alpha 1/2}S_{2,1}(f)\delta ^{1ϵ}.$$
Finally, note that if $`\chi _{[0,T]}`$ denotes the characteristic function of $`[0,T]`$ in the real line, then $`\frac{1}{T}_t|g_\delta (t/T)\chi _{[0,T]}(t)|𝑑t\delta `$. It follows that
(9.15)
$$\frac{1}{T}\left|_0^Tf(x_0a(T)^1n(t))𝑑t_tg_\delta (t/T)f(x_0a(T)^1n(t))𝑑t\right|\delta S_{\mathrm{},0}(f).$$
Combining (9.13), (9.14) and (9.15), and replacing $`x_0`$ by $`x_0a(T)`$, we conclude that the left hand side of (9.11) is bounded by
$$O_ϵ\left(S_{\mathrm{},1}(f)(\mathrm{max}(T^1,\delta )+T^{\alpha 1/2}\delta ^{1ϵ}+\delta )\right).$$
We choose $`\delta ^2=T^{\alpha 1/2}`$ to obtain the claimed conclusion. ∎
#### 9.3.2. Equidistribution of large horospheres on higher rank groups.
We now prove quantitative equidistribution of large closed horospheres. This result is well-known and generalizes the result of Sarnak, that the closed horocycle $`\{x+iy\}_{0x1}`$ is equidistributed in $`\mathrm{SL}_2()\backslash `$, as $`y0`$.
We shall follow the notation of Sec. 3.2, which we briefly reprise. Let $`G`$ be a connected semisimple (real) Lie group, $`\mathrm{\Gamma }G`$ a lattice, $`KG`$ the maximal compact subgroup, $`𝔤`$ the Lie algebra of $`G`$, and $`H𝔤`$ a semisimple element. Fix arbitrarily a norm $``$ on $`𝔤`$. We equip $`G`$ with the Haar measure in which $`\mathrm{\Gamma }\backslash G`$ has volume $`1`$. Let $`\mathrm{exp}:𝔤G`$ be the exponential map. Let $`𝔲`$ be the sum of all negative root spaces for $`H`$, and let $`U=\mathrm{exp}(𝔲)G`$. Let $`x_0\mathrm{\Gamma }\backslash G`$ be so that $`x_0U`$ is compact; note that the existence of such $`x_0`$ implies that $`\mathrm{\Gamma }\backslash G`$ is noncompact.
Let $`x_t=x_0\mathrm{exp}(tH)`$, and let $`\mathrm{\Delta }_t`$ be the stabilizer of $`x_t`$ in $`U`$. We denote by $`,_{L^2(\mathrm{\Gamma }\backslash G)}`$ the inner product in the Hilbert space $`L^2(\mathrm{\Gamma }\backslash G)`$.
###### Lemma 9.5.
There is $`\kappa _1>0`$ such that, for any $`f,gC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$ and for any $`U𝔲`$ with unit length (w.r.t. the fixed norm $``$ on $`𝔤`$) we have:
(9.16) $`\left|\mathrm{exp}(tH)f,g{\displaystyle _{\mathrm{\Gamma }\backslash G}}f{\displaystyle _{\mathrm{\Gamma }\backslash G}}g\right|\mathrm{exp}(\kappa _1|t|)S_{\mathrm{},dim(K)}(f)S_{\mathrm{},dim(K)}(g)`$
$`\left|\mathrm{exp}(sU)f,g{\displaystyle _{\mathrm{\Gamma }\backslash G}}f{\displaystyle _{\mathrm{\Gamma }\backslash G}}g\right|(1+|s|)^{\kappa _1}S_{\mathrm{},dim(K)}(f)S_{\mathrm{},dim(K)}(g)`$
Of course the constant $`\kappa _1`$ will depend on the choice of the norm $``$.
###### Proof.
This follows from a nice result of Kleinbock and Margulis: see . (The orthogonal complement $`L_0^2`$ of the identity representation in $`L^2(\mathrm{\Gamma }\backslash G)`$ is isolated, by \[18, Thm 1.12\], from the trivial representation in the unitary dual of $`\widehat{G}`$. A sufficiently high tensor power of $`L_0^2`$ is therefore tempered, whereupon one applies the bounds of .) Note that only claims the result (in effect) with $`S_{\mathrm{},d}`$ for some $`d`$; the fact that we can take $`d=dim(K)`$ follows by explicating the argument just sketched. ∎
Recall the definition of $`\nu _T`$ from Sec. 3.2, that is to say: $`\nu _T(f)=\frac{_{\mathrm{\Delta }_T\backslash U}f(x_Tu)𝑑u}{\mathrm{vol}(\mathrm{\Delta }_T\backslash U)}.`$ Thus $`\nu _T`$ is the measure supported on a closed horosphere, and this horosphere expands as $`T\mathrm{}`$. One deduces from Lem. 9.5 that the measures $`\nu _T`$ are equidistributed as $`T\mathrm{}`$:
###### Lemma 9.6.
Set $`\kappa _2=\frac{\kappa _1}{dim(G)+dim(K)+1}`$, $`\kappa _1`$ being as in the previous Lemma. Then, for $`T0`$ and $`fC^{\mathrm{}}(\mathrm{\Gamma }\backslash G)`$,
$$|\nu _T(f)_{\mathrm{\Gamma }\backslash G}f|e^{\kappa _2T}S_{\mathrm{},d}(f).$$
###### Proof.
The idea is identical to Lem. 9.4 and we refer to the first paragraph of that proof for a description of it.
Fix a left-invariant Riemannian metric on $`G`$. This descends to a metric on $`\mathrm{\Gamma }\backslash G`$. We first choose some “smoothing kernels” on $`G`$. For each $`ϵ>0`$, choose a function $`k_ϵC^{\mathrm{}}(G)`$ such that $`k_ϵ`$ is positive, supported in an $`ϵ`$-neighbourhood of the identity, $`_Gk_ϵ=1`$, and so that for any $`X_1,X_2,\mathrm{},X_l𝔤`$ we have:
(9.17)
$$\underset{gG}{sup}|X_1\mathrm{}X_lk_ϵ|_{X_1,\mathrm{},X_l}ϵ^{ldim(G)}.$$
It is easy to see this is possible (for example: choose an appropriate sequence of functions on $`𝔤`$ and transport to $`G`$ via the exponential map.)
The measure $`\nu _0`$ is a $`U`$-invariant probability measure supported on the closed orbit $`x_0U`$. $`\nu _0k_ϵ`$ is supported in an $`ϵ`$-neighbourhood of $`x_0U`$ and is given by integration against a $`C^{\mathrm{}}`$ density function $`g_ϵ`$, that is: $`\nu _0k_ϵ(f)=_{\mathrm{\Gamma }\backslash G}fg_ϵ.`$
Moreover, it follows from (9.17) that $`g_ϵ`$ satisfies the bounds $`S_{\mathrm{},l}(g_ϵ)ϵ^{ldim(G)}`$, for any $`l0`$.
The translate of $`\nu _0k_ϵ`$ by $`\mathrm{exp}(TH)`$ is supported in an $`ϵ`$-neighbourhood of $`x_TU`$; note it is essential that $`T0`$ for this. (Recall – Sec. 2.1 – our conventions are such that the right translate of the point mass at $`x`$ by $`gG`$ is the point mass at $`xg^1`$.)
In fact, one verifies that
(9.18)
$$\left|\nu _T(f)\nu _0k_ϵ(\mathrm{exp}(TH)f)\right|ϵS_{\mathrm{},1}(f).$$
(Indeed, let $`g\mathrm{supp}(k_ϵ)`$ and let $`\delta _g`$ be the point mass at $`g`$. It suffices to check that the identical bound holds for $`\left|\nu _T(f)\nu _0\delta _g(\mathrm{exp}(TH)f)\right|`$, which equals $`\left|\nu _T(f)\nu _0(g\mathrm{exp}(TH)f)\right|`$. Let $`𝔟`$ the sum of non-negative root spaces for $`H`$ on $`𝔤`$. If $`ϵ`$ is sufficiently small, we may write $`g=um`$, with $`u\mathrm{exp}(𝔲)`$ and $`m\mathrm{exp}(𝔟)`$. Moreover, again if $`ϵ`$ is sufficiently small, $`u,m`$ lie in a $`Cϵ`$-neighbourhood of the identity, for some fixed constant $`C`$. Then $`\mathrm{exp}(TH)g\mathrm{exp}(TH)=u^{}m^{}`$, with $`u^{}\mathrm{exp}(𝔲)`$ and where $`m^{}\mathrm{exp}(𝔟)`$ is in a $`C^{}ϵ`$-neighbourhood of the identity, for some absolute $`C^{}`$. Also, $`\nu _0(g\mathrm{exp}(TH)f)=\nu _T(m^{}f)`$. Thus it suffices to bound $`\left|\nu _T(f)\nu _T(m^{}f)\right|`$. But the $`L^{\mathrm{}}`$ norm of $`fm^{}f`$ is $`ϵS_{\mathrm{},1}(f)`$.)
On the other hand, by Lem. 9.5, for $`T0`$: we have
(9.19)
$$\begin{array}{c}\left|\nu _0k_ϵ(\mathrm{exp}(TH)f)_{\mathrm{\Gamma }\backslash G}f\right|=\left|g_ϵ,\mathrm{exp}(TH)f_{L^2(\mathrm{\Gamma }\backslash G)}_{\mathrm{\Gamma }\backslash G}f\right|\hfill \\ \hfill \mathrm{exp}(\kappa _1T)S_{\mathrm{},dim(K)}(f)ϵ^{dim(G)dim(K)}.\end{array}$$
It follows from this and (9.18) that
$$|\nu _T(f)_{\mathrm{\Gamma }\backslash G}f|(ϵ+\mathrm{exp}(\kappa _1T)ϵ^{dim(G)dim(K)})S_{\mathrm{},dim(K)}(f).$$
To conclude, take $`ϵ=\mathrm{exp}(\frac{\kappa _1T}{dim(G)+dim(K)+1})`$. ∎
### 9.4. The equidistribution of Hecke orbits and $`p`$-adic horocycles.
In this section, we prove some “$`p`$-adic” equidistribution statements, pertaining to the equidistribution of Hecke points and $`p`$-adic horocycles.
In the Lemmas that follow, $`𝔣`$ will be a prime ideal of $`F`$, $`\overline{\mu }_𝔣`$, the normalized Hecke measure defined subsequent to (2.6), and $`[𝔣]`$ as defined in Sec. 2.5.
The first Lemma is an adelic version of the fact that the Hecke orbit $`T_q(z)`$ of a point $`z\mathrm{SL}(2,)\backslash `$ is equidistributed, as $`z\mathrm{}`$.
###### Lemma 9.7.
Let $`fC_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$ and $`𝔣`$ an ideal of $`F`$. Then, for $`x_0𝐗_{\mathrm{GL}(2)}`$, $`d1`$,
(9.20)
$$\left|f\overline{\mu }_𝔣(x_0)\underset{\chi ^2=\omega }{}\chi ([𝔣])\chi (x_0)_{x𝐗}f(x)\chi (x)𝑑\mu _𝐗(x)\right|\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}0pt(x)^{1/2}S_{2,d}(f).$$
Here $`\chi (x)`$ denotes the function $`g\chi (det(g))`$ on $`𝐗`$.
###### Proof.
Let $`𝒫`$ be the projection defined in Sec. 2.7. Let $`E`$ be the endomorphism $`f(f𝒫f)\overline{\mu }_𝔣`$ of $`C_\omega ^{\mathrm{}}(𝐗_{\mathrm{GL}(2)})`$. The operator $`E`$ has norm $`_ϵ\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}`$ w.r.t. the $`L^2`$ norm (this follows from Lem. 2.1 and the bounds of Sec. 9.1). By Lem. 8.3 it follows that the operator norm of $`E`$ w.r.t $`S_{2,d,\beta }`$ is also $`_ϵ\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}`$.
The left hand side of (9.20) is exactly $`Ef(x_0)`$. Now apply Lem. 9.1, with $`p=2`$, to conclude. ∎
The next Lemma is an adelic version of the following (again closely connected to equidistribution of Hecke points). Let $`Y(p)`$ be embedded in $`Y(1)\times Y(1)`$ (notation of discussion after Prop. 4.1). Then $`Y(p)`$ is equidistributed as $`p\mathrm{}`$. The quantification of this is slightly complicated by noncompactness; in particular, we must use Sobolev norms $`S_{p,d}`$ for $`p>2`$. (Cf. discussion in Sec. 2.9.1).
###### Lemma 9.8.
Let $`𝔮`$ be a prime ideal of $`𝔬_F`$. Let $`FC^{\mathrm{}}(𝐗\times 𝐗)`$ be $`\mathrm{PGL}_2(𝔬_{F_𝔮})\times \mathrm{PGL}_2(𝔬_{F_𝔮})`$ invariant.Then, for any $`d1,p>2`$,
(9.21)
$$\begin{array}{c}\left|_𝐗F(x,xa([𝔮]))𝑑x\underset{\chi ^2=1}{}\chi ([𝔮])_𝐗F(x,y)\chi (x)\chi (y)𝑑\mu _𝐗(x)𝑑\mu _𝐗(y)\right|\hfill \\ \hfill _ϵ\mathrm{N}(𝔮)^{\frac{2\alpha 1}{p}+ϵ}S_{p,d}(F).\end{array}$$
###### Proof.
Let $`\sigma `$ be the measure $`\delta _1\times \overline{\mu }_𝔮`$ on $`\mathrm{PGL}_2(F_𝔮)\times \mathrm{PGL}_2(F_𝔮)`$, where $`\delta _1`$ is the measure consisting of a point mass at the identity. Recalling (see (2.6) in the case of a prime ideal, and Section 2.5 for the definition of $`K_𝔮`$) that $`\overline{\mu }_𝔮`$ is the $`K_𝔮`$-bi-invariant probability measure supported on $`K_𝔮a([𝔮])K_𝔮`$, we note that
$$(F\sigma )(x,x)=_{k_1,k_2K_𝔮}F(x,xk_1a([𝔮])k_2)𝑑k_1𝑑k_2=_{K_𝔮}F(xk,xka([𝔮]))𝑑k,$$
where we equip $`K_𝔮`$ with the Haar measure of mass $`1`$, and we use the $`\mathrm{PGL}_2(𝔬_{F_𝔮})`$-invariance of $`F`$ at the second step. It follows that
(9.22)
$$_𝐗F(x,xa([𝔮]))𝑑x=_𝐗(F\sigma )(x,x)𝑑x.$$
Let $`𝒫_2`$ be as in Sec. 2.7. Let $`E`$ be the endomorphism of $`C^{\mathrm{}}(𝐗\times 𝐗)`$ defined by $`E(F)=(F𝒫_2F)\sigma `$. Combining (9.22) and the easily verified equality
$$_𝐗(𝒫_2F\sigma )(x,x)𝑑x=\underset{\chi ^2=1}{}\chi ([𝔮])_𝐗F(x,y)\chi (x)\chi (y)𝑑\mu _𝐗(x)𝑑\mu _𝐗(y),$$
we see that the left hand side of (9.21) is precisely $`_𝐗EF(x,x)𝑑x`$.
Since $`𝒫_2`$ does not increase $`L^{\mathrm{}}`$ norms, and $`\sigma `$ is a probability measure, it follows that the operator norm of $`E`$ w.r.t the $`L^{\mathrm{}}`$ norm is $`2`$. Moreover, the operator norm of $`E`$ w.r.t the $`L^2`$ norm is $`\mathrm{N}(𝔮)^{\alpha 1/2}`$, as follows from Lem. 2.1.
Lem. 8.3 now implies that for $`2p\mathrm{}`$ we have the majorization $`S_{p,d}(EF)\mathrm{N}(𝔮)^{\frac{2\alpha 1}{p}+ϵ}S_{p,d}(F)`$. Now Lem. 9.1 shows that
$$|_𝐗EF(x,x)𝑑x|S_{p,d}(EF)\mathrm{N}(𝔮)^{\frac{2\alpha 1}{p}+ϵ}S_{p,d}(F)$$
for $`p>2,d1`$; whence the conclusion of the Lemma. ∎
The next Lemma shows the equidistribution of certain $`p`$-adic horocycle orbits, as $`p`$ varies. The idea will be as follows: (speaking very loosely, in the case of $`\mathrm{SL}_2`$) a typical $`p`$-adic horocycle orbit, when projected to $`\mathrm{SL}(2,)\backslash `$, looks like $`\{z+\frac{i}{p}\}_{0ip1}`$. This set looks very much like the image, under the $`p`$-Hecke operator, of the point $`pz`$. Thus one can deduce distribution properties of the $`p`$-adic horocycle orbit from some standard facts about Hecke operators.
This is a rather ad hoc argument. Let us say a few words about why this problem does not quite fit into the usual setup of such questions. We are proving statements about the distribution of e.g. $`p`$-adic horocycles when $`p`$ varies. This does not fit easily into the usual context of such matters, where one considers e.g. a fixed unipotent flow on an $`S`$-arithmetic homogeneous space. It would be interesting to have a more conceptual and natural way of treating such questions, in the aspect where “$`p`$ varies.”
###### Lemma 9.9.
Let $`fC^{\mathrm{}}(𝐗)`$ and let $`𝔣`$ be an integral ideal of $`𝔬_F`$, factorizing as $`𝔣=_𝔮𝔮^{e_𝔮}`$. For each $`𝔮|𝔣`$, let $`s_𝔮0`$ be a non-negative integer, and suppose $`f`$ is invariant by $`_{𝔮|𝔣}K_0[𝔮^{s_𝔮}]`$. Put $`𝔪=_{𝔮|𝔣}𝔮^{s_𝔮}`$. Let $`\eta _𝔣`$ be the Haar probability measure on $`_{𝔮|𝔣}N(𝔮^{e_𝔮}𝔬_𝔮)`$ and $`dh`$ the Haar probability measure on $`\mathrm{SL}_2(F)\backslash \mathrm{SL}_2(𝔸_F)`$.
Then, for $`yF^\times \backslash 𝔸_F^\times `$,
(9.23)
$$\begin{array}{c}\left|f\eta _𝔣(a(y))_{hSL_2(F)\backslash \mathrm{SL}_2(𝔸_F)}f(ha(y))𝑑h\right|\hfill \\ \hfill _ϵ\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}\mathrm{max}(\mathrm{N}(𝔣)|y|,\frac{1}{\mathrm{N}(𝔣)|y|})^{1/2}\mathrm{N}(𝔪)^{3/2+ϵ}S_{2,d}(f)\end{array}$$
###### Proof.
As usual let $`K_{v,f}`$ be the stabilizer of $`f`$ in $`K_v=\mathrm{GL}_2(𝔬_v)`$, so that $`K_{𝔮,f}`$ contains $`K_0[𝔮^{s_𝔮}]`$ for each $`𝔮|𝔣`$.
We now define a measure $`\stackrel{~}{\eta }_𝔮`$ on $`\mathrm{PGL}_2(F_𝔮)`$ for each $`𝔮|𝔣`$. It will “approximate” $`\eta _𝔣`$ but will be composed of Hecke operators.
For those $`𝔮`$ such that $`s_𝔮=0`$, put
(9.24)
$$\stackrel{~}{\eta }_𝔮=\mathrm{N}(𝔮)^{e_𝔮/2}\delta _{a(\varpi ^{e_𝔮})}\mu _{𝔮^{e_𝔮}}\mathrm{N}(𝔮)^{\frac{e_𝔮1}{2}}\delta _{a(\varpi ^{e_𝔮1})}\mu _{𝔮^{e_𝔮1}}.$$
(We refer to Sec. 2.8 for definitions of $`\mu _\mathrm{?}`$ appearing above.) For $`𝔮`$ such that $`s_𝔮1`$, we set $`\sigma _𝔮`$ to be the unique bi-$`K_0[𝔮^{s_𝔮}]`$-invariant probability measure on $`K_0[𝔮^{s_𝔮}]a(\varpi ^{e_𝔮})K_0[𝔮^{s_𝔮}]`$, normalized to have mass $`1`$, and we put $`\stackrel{~}{\eta }_𝔮=\delta _{a(\varpi ^{e_𝔮})}\sigma _𝔮`$. Finally, set $`\stackrel{~}{\eta }_𝔣=_{𝔮|𝔣}\stackrel{~}{\eta }_𝔮`$.
One then verifies by a direct computation that
(9.25)
$$f\eta _𝔣=f\stackrel{~}{\eta }_𝔣$$
The intuition for this statement, in the classical setting, as as follows: let $`z\mathrm{SL}_2()\backslash `$. Then (for a prime number $`p`$) the set $`\{z+i/p\}_{0ip1}`$ is the $`p`$-Hecke orbit of $`pz`$, with the point $`p^2z`$ removed. In the case $`e_𝔮=1`$, the first term on the right hand side of (9.24) corresponds to the $`p`$-Hecke orbit of $`pz`$, and the second term corresponds to removing the point $`p^2z`$.
More formally, to verify (9.25), the unramified computation, at those places where $`s_𝔮=0`$, is easy; the ramified computation is just Hecke theory at ramified primes, see e.g. \[33, Prop 3.33\]. <sup>20</sup><sup>20</sup>20For the unramified assertion, let $`=\mathrm{PGL}_2(F_𝔮)/K_𝔮`$ and let $`x_0`$ be the identity coset. The set $``$ has the structure of the vertices of a $`q_v+1`$-valent tree. Let $`S_1`$ be the set of all vertices at distance $`e_𝔮2i`$ (some $`i0`$) from $`a(\varpi ^{e_𝔮})x_0`$. Let $`S_2`$ be the set of all vertices at even distance $`e_𝔮12i`$ (some $`i0`$) from $`a(\varpi ^{e_𝔮1})x_0`$. Then $`S_2S_1`$ and $`S_1S_2`$ is precisely the $`n(𝔮^{e_𝔮})`$-orbit of $`x_0`$. As for the ramified case: one notes that, if $`s_𝔮>0`$, the Haar measure on $`K_0[𝔮^{s_𝔮}]`$ is just the pushforward of the Haar measure on $`n(𝔬_𝔮)\times a(𝔬_𝔮^\times )\times \overline{n}(𝔮^{s_𝔮})`$ by the product map $`(n,a,\overline{n})na\overline{n}`$.
The projection $`𝒫`$ of Sec. 2.7 commutes with the action of $`\mathrm{GL}_2(𝔸_F)`$, and so (9.25) holds also with $`f`$ replaced by $`𝒫f`$ or $`f𝒫f`$. Moreover, $`𝒫f\eta _𝔣=𝒫f`$. It follows that
(9.26)
$$f\eta _𝔣(x)=_{hSL_2(F)\backslash \mathrm{SL}_2(𝔸_F)}f(hx)𝑑h+(f𝒫f)\stackrel{~}{\eta }_𝔣(x).$$
Set $`\overline{f}=f𝒫f`$. Then, expanding the term $`\overline{f}\stackrel{~}{\eta }_𝔣`$:
(9.27)
$$\begin{array}{c}\overline{f}\stackrel{~}{\eta }_𝔣=\underset{S\{𝔮|𝔣,s_𝔮=0\}}{}\overline{f}\underset{𝔮|𝔣:s_𝔮1}{}\delta _{a(\varpi ^{e_𝔮})}\sigma _𝔮\hfill \\ \hfill \underset{𝔮|𝔣:s_𝔮=0,𝔮S}{}\left(\mathrm{N}(𝔮)^{e_𝔮/2}\delta _{a(\varpi ^{e_𝔮})}\mu _{𝔮^{e_𝔮}}\right)\underset{𝔮S}{}\left(\mathrm{N}(𝔮)^{\frac{e_𝔮+1}{2}}\delta _{a(\varpi ^{e_𝔮1})}\mu _{𝔮^{e_𝔮1}}\right)\end{array}$$
We now specialize to the case under consideration where $`x=a(y)`$ for some $`y𝔸_F^\times `$. For $`S\{𝔮|𝔣,s_𝔮=0\}`$ set
$$\sigma _S=\underset{𝔮|𝔣:s_𝔮1}{}\sigma _𝔮\underset{𝔮|𝔣:s_𝔮=0,𝔮S}{}\mathrm{N}(𝔮)^{e_𝔮/2}\mu _{𝔮^{e_𝔮}}\underset{s_𝔮=0,𝔮S}{}\mathrm{N}(𝔮)^{\frac{e_𝔮+1}{2}}\mu _{𝔮^{e_𝔮1}}.$$
With this notation, we have:
(9.28)
$$\overline{f}\stackrel{~}{\eta }_𝔣(a(y))=\underset{S\{𝔮|𝔣,s_𝔮=0\}}{}\overline{f}\sigma _S\left(a(y[𝔣])\underset{s_𝔮=0,𝔮S}{}a([𝔮])\right)$$
Now apply Lem. 9.1 to see that, for any $`z𝔸_F^\times `$ and $`d1`$, we have
(9.29)
$$\overline{f}\sigma _S(a(z))\mathrm{max}(|z|,|z|^1)^{1/2}PS_{2,d}(\overline{f}\sigma _S),$$
where we have used the easily verified fact that $`0pt(a(z))\mathrm{max}(|z|,|z|^1)`$.
Now, for any $`fC^{\mathrm{}}(𝐗)`$, we have
$$[K_{\mathrm{max}}:K_{\overline{f}\sigma _S}]\underset{s_𝔮1}{}[K_𝔮:K_0[𝔮^{s_𝔮}]][K_{\mathrm{max}}:K_f]_ϵ\mathrm{N}(𝔪)^{1+ϵ}[K_{\mathrm{max}}:K_f].$$
By the bounds on matrix coefficients (9.2), and recalling that $`𝔪=_{𝔮|𝔣}𝔮^{s_𝔮}`$, we compute that
$$PS_{2,d}(\overline{f}\sigma _S)_ϵ\mathrm{N}(𝔪)^{3/2+ϵ}(\mathrm{N}(𝔣)\underset{𝔮S}{}\mathrm{N}(𝔮)^1)^{\alpha 1/2+ϵ}\underset{𝔮S}{}\mathrm{N}(𝔮)^1PS_{2,d}(f).$$
Combining this with (9.28) and (9.29), we find that for each $`S\{𝔮:s_𝔮1\}`$
(9.30)
$$\begin{array}{c}|\overline{f}\stackrel{~}{\eta }_𝔣(a(y))|_ϵ\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}\mathrm{N}(𝔪)^{3/2+ϵ}\mathrm{max}(\mathrm{N}(𝔣)|y|,\mathrm{N}(𝔣)^1|y|^1)^{1/2}PS_{2,d}(f)\hfill \end{array}$$
This bound is valid for all $`fC^{\mathrm{}}(𝐗)`$, not merely those $`f`$ that are invariant by $`_{𝔮|𝔣}K_0[𝔮^{s_𝔮}]`$. Apply Rem. 8.1, (3) to the endomorphism $`f\overline{f}\stackrel{~}{\eta }_𝔣`$; this shows that (9.30) remains valid, for any $`fC^{\mathrm{}}(𝐗)`$, if we replace $`PS_{2,d}`$ by $`S_{2,d}`$ on the right hand side. Now, specialize to the case where $`fC^{\mathrm{}}(𝐗)`$ is actually $`_{𝔮|𝔣}K_0[𝔮^{s_𝔮}]`$-invariant and apply (9.26) to obtain the conclusion of the Lemma. ∎
The Lemma that follows states an adelic version of the following fact: the measure on $`\mathrm{SL}_2()\backslash `$ defined by $`\nu _y:=q^1_{0xq1}\delta _{\frac{x}{q}+iy}`$, approximates the uniform measure if $`yq^1`$; more precisely we have an inequality that $`\left|\nu _y(f)_{\mathrm{SL}_2()\backslash }f\right|`$ is bounded by $`\mathrm{max}(qy,\frac{1}{qy})^{1/2}q^\delta S(f)`$, where $`S`$ is an appropriate Sobolev norm and $`\delta >0`$.
###### Lemma 9.10.
Let $`fC^{\mathrm{}}(𝐗)`$ and let notations be as in Sec. 6(see esp. (6.4)). In particular, $`𝔣`$ is an integral ideal of $`𝔬_F`$, $`q=\mathrm{N}(𝔣)`$, $`[𝔣]`$ is as in (2.4) and
$$\nu _z(f)=_{|y|=z,y𝔸_F^\times /F^\times }f(a(y)n([𝔣]))d^\times y.$$
Suppose $`f`$ is invariant by $`K_0[𝔮^{s_𝔮}]`$, for each $`𝔮|𝔣`$, and put $`𝔪=_{𝔮|𝔣}𝔮^{s_𝔮}`$. Then
(9.31)
$$\begin{array}{c}\left|\nu _z(f)_𝐗f(x)𝑑\mu _𝐗(x)\right|\hfill \\ \hfill _ϵ\mathrm{N}(𝔣)^{\alpha 1/2+ϵ}\mathrm{N}(𝔪)^{3/2+ϵ}\mathrm{max}(\mathrm{N}(𝔣)z,\frac{1}{\mathrm{N}(𝔣)z})^{1/2}S_{2,d}(f).\end{array}$$
###### Proof.
For each $`𝔮|𝔣`$ and integer $`0e`$, let $`\eta _{𝔮^e}`$ be the Haar probability measure on the group $`N(𝔮^e𝔬_𝔮)`$. Then, since the assumption implies that $`f`$ is right invariant by $`a_𝔮(𝔬_𝔮^\times )`$, for each $`𝔮`$ dividing $`𝔣`$, we see that for any $`x𝐗`$:
(9.32)
$$\begin{array}{c}_{y𝔬_{F_𝔮}^\times }f(xa(y)n_𝔮(\varpi _𝔮^e))d^\times y=_{y𝔬_{F_𝔮}^\times }f(xn_𝔮(y\varpi _𝔮^e))d^\times y\hfill \\ \hfill =\frac{f(\eta _{𝔮^e}\mathrm{N}(𝔮)^1\eta _{𝔮^{e1}})(x)}{1\mathrm{N}(𝔮)^1}\end{array}$$
It follows that
$$\nu _z(f)=\underset{𝔮|𝔣}{}(1\mathrm{N}(𝔮)^1)^1_{y𝔸_F^\times /F^\times ,|y|=z}f\underset{𝔮|𝔣}{}(\eta _{𝔮_𝔮^e}\mathrm{N}(𝔮)^1\eta _{𝔮^{e_𝔮1}})(a(y)).$$
We conclude by applying the previous Lemma. ∎
## 10. Background on Eisenstein series.
This section essentially develops the theory of Eisenstein series on $`\mathrm{PGL}_2`$ over a number field. This is needed for the Rankin-Selberg method that we reprise in the next section, which in turn is used in the text to relate a period integral with an $`L`$-function.
Let $`Z`$ be a topological space. In this section, we will often speak – in various contexts, often with $`Z=𝐗`$ or $`\mathrm{GL}_2(𝔸_F)`$ – of a function $`F(s,z)`$ on $`\times Z`$ being “holomorphic” or “holomorphic in $`s`$.” For the purposes of this document, this can be assumed to mean that the function is jointly continuous and holomorphic for each $`z`$ individually.
Note that $`s_ZF(s,z)𝑑z`$, if absolutely convergent and uniformly so in $`s`$, defines a holomorphic function. Indeed, it suffices to verify that its integral over a closed curve in the $`s`$-variable is zero, which follows by Fubini’s theorem.
Similarly, we will say that $`F(s,z)`$ is meromorphic if there exists a holomorphic function $`h(s)`$ so that $`h(s)F(s,z)`$ is holomorphic.
### 10.1. Construction and basic properties of the Eisenstein series.
We recall the Eisenstein series that we shall have need of and its basic properties, following Jacquet \[17, §19\]. We will need Eisenstein series only on $`\mathrm{PGL}_2`$.
#### 10.1.1. Schwarz functions.
Let $`\mathrm{\Psi }`$ be a Schwarz-Bruhat function on $`𝔸_F^2`$, i.e. $`\mathrm{\Psi }`$ is a finite linear combination of functions $`_v\mathrm{\Psi }_v`$, where each $`\mathrm{\Psi }_v`$ is locally constant of compact support, for $`v`$ finite, $`\mathrm{\Psi }_v`$ is a Schwarz function on $`F_v^2`$ for $`v`$ infinite, and $`\mathrm{\Psi }_v`$ is the characteristic function of $`𝔬_v^2`$ for almost all $`v`$.
If $`v`$ is a real place, choose $`a_vF_v`$ so that $`e_{F_v}(x)=e^{2\pi ia_vx}`$, and say a Schwarz function $`\mathrm{\Psi }_v`$ on $`F_v^2`$ is standard if it is the product of a polynomial and $`e^{\pi |a_v|_v(|x|_v^2+|y|_v^2)}`$. If $`v`$ is a complex place, choose $`a_vF_v`$ so that $`e_{F_v}(x)=e^{2\pi i\mathrm{Tr}_/(a_vx)}`$; we say that a Schwarz function $`\mathrm{\Psi }_v`$ is standard if it is the product of a polynomial and $`e^{2\pi |a|_v^{1/2}(|x|_v+|y|_v)}`$. The significance of this normalization is twofold: a standard function is automatically $`K_v`$-finite and also the class of standard functions is self-dual under the Fourier transform corresponding to the character $`e_{F_v}`$.
If $`V`$ is a real vector space, then by a Schwarz norm on the space of Schwarz functions on $`V`$, we shall mean a norm $`𝒮`$ of the form
(10.1)
$$𝒮(\mathrm{\Psi })=\underset{𝒟}{sup}\underset{x}{sup}|(1+x)^M𝒟\mathrm{\Psi }|,$$
for some finite collection of constant-coefficients differential operators $`𝒟`$ on $`V`$ and some norm $`x`$ on $`V`$.
Put, for $`g\mathrm{GL}_2(𝔸_F)`$,
$$f_\mathrm{\Psi }(s,g)=|det(g)|^s_{t𝔸_F^\times }\mathrm{\Psi }((0,t)g)|t|^{2s}d^\times t.$$
The integral converges absolutely for $`\mathrm{}(s)>1/2`$ and extends to a meromorphic function of $`s`$ with possible poles at most at $`s=0,1/2`$. Moreover, for all $`s`$,
(10.2)
$$f_\mathrm{\Psi }(\left(\begin{array}{cc}a& x\\ 0& b\end{array}\right)g)=|a/b|^sf(g).$$
Put $`E_\mathrm{\Psi }(s,g)=_{\gamma B(F)\backslash \mathrm{GL}_2(F)}f(s,\gamma g)`$. This converges when $`\mathrm{Re}(s)>1`$, extends to a meromorphic function of $`s`$ with a simple pole at $`s=0,1`$ and satisfies the functional equation
$$E_\mathrm{\Psi }(s,g)=E_{\widehat{\mathrm{\Psi }}}(1s,g),$$
where $`\widehat{\mathrm{\Psi }}`$ is the Fourier transform
(10.3)
$$\widehat{\mathrm{\Psi }}(x_1,y_1)=_{𝔸_F^2}\mathrm{\Psi }(x,y)e_F(x_1yy_1x)𝑑x𝑑y.$$
Moreover, the pole at $`s=1`$ is the constant function with value $`c_1_{𝔸_F^2}\mathrm{\Psi }(x,y)𝑑x𝑑y`$, and the pole at $`s=0`$ is the constant function with value $`c_2\mathrm{\Psi }(0)`$, where $`c_1,c_2`$ are constants (depending only on the choice of measure). Finally, for any fixed $`g`$ the function $`ss(1s)E_\mathrm{\Psi }(s,g)`$ decays rapidly in vertical strips, i.e. $`(1+|s|)^N|s(1s)E_\mathrm{\Psi }(s,g)|`$ is bounded in any strip $`A\mathrm{}(s)B`$. The proof of all these properties follows from “Poisson summation” for $`F^2𝔸_F^2`$, and we omit them.
Moreover, the association $`\mathrm{\Psi }E_\mathrm{\Psi }`$ is twisted-equivariant for the natural $`\mathrm{GL}_2(𝔸_F)`$-action on the space of Schwarz functions and on $`C^{\mathrm{}}(𝐗)`$: that is to say,
(10.4)
$$E_{h.\mathrm{\Psi }}(s,g)=|det(h)|^s(hE_\mathrm{\Psi }(s,g)),$$
where $`h`$ denotes right translation by $`h`$.
We give an example with $`F=`$ (cf. \[15, (3.29)\]).
###### Example 10.1.
Suppose $`F=`$, $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$ where, for each finite $`v`$, $`\mathrm{\Psi }_v`$ is the characteristic function of the maximal compact of $`F_v`$, and $`\mathrm{\Psi }_{\mathrm{}}(x,y)=e^{\pi (x^2+y^2)}`$. Then $`E_\mathrm{\Psi }(g)`$ is determined by its restriction to $`\mathrm{SL}_2()`$.
Moreover, $`E_\mathrm{\Psi }(s,g)`$ descends from a function of $`g\mathrm{SL}_2()`$ to a function $`E^{}(s,z)`$ on $`=\mathrm{SL}_2()/\mathrm{SO}_2`$, where the identification is $`ggi`$. In fact,
(10.5)
$$E^{}(s,z)=\pi ^s\mathrm{\Gamma }(s)\zeta (2s)\underset{[c:d]^1()}{}\frac{y^s}{|cz+d|^{2s}}.$$
If we put $`\xi (s)=\pi ^{s/2}\mathrm{\Gamma }(s/2)\zeta (s)`$, then $`E^{}(s,z)`$ has the Fourier expansion
(10.6)
$$E^{}(s,z)=\xi (2s)y^s+\xi (22s)y^{1s}+4\sqrt{y}\underset{n}{}K_{s1/2}(2\pi ny)\mathrm{cos}(2\pi ny)\underset{ab=n}{}\left(\frac{a}{b}\right)^{s1/2}$$
It satisfies the functional equation $`E^{}(s,z)=E^{}(1s,z)`$. Moreover it is a meromorphic function of $`s`$ with poles precisely at $`s=0`$ and $`s=1`$. In both cases the residue is the constant function.
Motivated by this example, the reader may find it helpful to keep in mind the “dictionary”: $`f_\mathrm{\Psi }(s,g)`$ corresponds to $`\pi ^s\mathrm{\Gamma }(s)\zeta (2s)y^s=\xi (2s)y^s`$, and $`E_\mathrm{\Psi }(s,g)`$ to $`E^{}(s,z)`$ as defined in (10.5).
###### Remark 10.1.
Suppose $`\mathrm{\Psi }`$ is invariant by $`K_{\mathrm{}}\times K_{\mathrm{max}}`$. Then $`f_\mathrm{\Psi }`$ is a multiple of $`g0pt(g)^s`$, as follows from the uniqueness of spherical functions satisfying (10.2). Thus, for $`\mathrm{}(s)>1`$, $`E_\mathrm{\Psi }(s,g)=c(s)_{\gamma B(F)\backslash \mathrm{GL}_2(F)}0pt(\gamma g)^s`$.
We now proceed to establish the “standard” properties of the Eisenstein series for $`E_\mathrm{\Psi }`$. It is convenient to first recall an explicit bound for archimedean Mellin transforms; the first part is Tate’s thesis, and the second will only be needed much later.
###### Lemma 10.1.
Let $`v`$ be archimedean and let $`\mathrm{\Psi }_v`$ be a Schwarz function on $`F_v`$. The integral $`G(s):=_{xF_v^\times }\mathrm{\Psi }_v(x)|x|^sd^\times x`$ extends to a meromorphic function and:
1. $`\frac{G(s)}{\zeta _{F,v}(s)}`$ is holomorphic, where $`\zeta _{F,v}(s)`$ is the local factor of the Dedekind $`\zeta `$-function of $`F`$ at $`v`$.
2. For any $`N0`$, the function $`G_N(s):=_{i=0}^N(s+i)G(s)`$ is holomorphic in $`\mathrm{}(s)N`$, and the absolute value of $`(1+|s|)^MG_N(s)`$ in any strip $`N\mathrm{}(s)A`$ is bounded by some Schwarz norm (depending on $`A,N,M`$; see (10.1) for the definition) of $`\mathrm{\Psi }_v`$.
###### Proof.
The first assertion is Tate’s thesis, and we leave the second to the reader (if any). ∎
###### Lemma 10.2.
The function $`ss(1/2s)f_\mathrm{\Psi }(s,g)`$ extends to a holomorphic function of $`s`$. It decays rapidly along vertical lines:
(10.7)
$$|(1+|\mathrm{}(s)|)^Ns(1/2s)f_\mathrm{\Psi }(s,g)|_\mathrm{\Psi }0pt(g)^{\mathrm{}(s)},$$
where the implicit constant is uniform for $`\mathrm{}(s)`$ in a compact set.
###### Proof.
By (10.2) and the Iwasawa decomposition, it will suffice to prove the assertions in the special case $`gK_{\mathrm{}}\times K_{\mathrm{max}}`$. So we write $`g=kK_{\mathrm{}}\times K_{\mathrm{max}}`$ and denote by $`k_v`$ the component of $`k`$ in $`\mathrm{PGL}_2(F_v)`$. Moreover, without loss of generality, we may assume $`\mathrm{\Psi }`$ is a product of Schwarz functions at each place, i.e. $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$. Then
$$f_\mathrm{\Psi }(s,g)=\underset{v\mathrm{infinite}}{}_{F_v^\times }\mathrm{\Psi }_v((0,t)k)|t|^{2s}d^\times t\underset{v\mathrm{finite}}{}_{F_v^\times }\mathrm{\Psi }_v((0,t)k_v)|t|^{2s}d^\times t$$
By Tate’s thesis, it follows that that the product over finite places is of the form $`\zeta _F(2s)h(s)`$, where $`\zeta _F()`$ is the (finite part of the) Dedekind $`\zeta `$-function of the number field $`F`$ and $`h(s)`$ is a holomorphic function with at most polynomial growth in vertical strips (indeed, a polynomial in $`q^{\pm s}`$ for various $`q`$). All the assertions of the Lemma now follow from Lem. 10.1, and standard facts about the analytic properties of $`\zeta _F`$.
In fact, if the $`\mathrm{\Psi }_v`$ for $`v`$ finite are regarded as fixed, then the implicit constant in (10.7) is bounded by an appropriate Schwarz norm, depending on $`N`$ and the compact set to which $`\mathrm{}(s)`$ is constrained, of $`_{v\mathrm{infinite}}\mathrm{\Psi }_v`$. This follows from the second assertion of Lem. 10.1. ∎
###### Lemma 10.3.
The constant term $`E_\mathrm{\Psi }^N(s,g):=_{xF\backslash 𝔸_F}E_\mathrm{\Psi }(s,n(x)g)𝑑x`$ equals $`f_\mathrm{\Psi }(s,g)+f_{\widehat{\mathrm{\Psi }}}(1s,g)`$.
###### Proof.
(Sketch). A double coset decomposition shows that, for $`s1`$, $`E_\mathrm{\Psi }^N(s,g)=f_\mathrm{\Psi }(s,g)+_{nN(𝔸_F)}f_\mathrm{\Psi }(s,wng)𝑑n`$. So it will suffice to show that $`_{nN(𝔸_F)}f_\mathrm{\Psi }(s,wng)=f_{\widehat{\mathrm{\Psi }}}(1s,g)`$. The left-hand side may be expressed as
(10.8)
$$\begin{array}{c}|det(g)|^s_{t𝔸_F^\times }_{x𝔸_F}\mathrm{\Psi }((t,tx)g)|t|^{2s}d^\times t𝑑x\hfill \\ \hfill =|det(g)|^s_{t𝔸_F^\times /F^\times }_{x𝔸_F}\underset{\delta F^\times }{}\mathrm{\Psi }(t(\delta ,x)g)|t|^{2s}d^\times tdx\end{array}$$
For any Schwarz function $`\mathrm{\Psi }`$ on $`𝔸_F^2`$, one has $`_{\alpha F}_{y𝔸_F}\mathrm{\Psi }(\alpha ,y)=_{\beta F}\widehat{\mathrm{\Psi }}(0,\beta )`$. The result follows from routine manipulation and use of Tate’s functional equation. ∎
We set
(10.9)
$$\overline{E}_\mathrm{\Psi }(s,g)=E_\mathrm{\Psi }(s,g)f_\mathrm{\Psi }(s,g)f_{\widehat{\mathrm{\Psi }}}(1s,g),$$
so $`\overline{E}_\mathrm{\Psi }`$ defines a function on $`B(F)\backslash PGL_2(𝔸_F)`$. It is a “truncated” Eisenstein series where we have removed the constant term. Moreover, $`\overline{E}_\mathrm{\Psi }(s,g)`$ is holomorphic in $`s`$ (this follows, for example, by computing residues at each of the points $`s=0,1/2,1`$ and seeing they are all zero). By definition, for $`g\mathrm{GL}_2(𝔸_F)`$ we have an equality
(10.10)
$$E_\mathrm{\Psi }(s,g)=\overline{E}_\mathrm{\Psi }(s,g)+f_\mathrm{\Psi }(s,g)+f_{\widehat{\mathrm{\Psi }}}(1s,g).$$
###### Lemma 10.4.
Let $`T,N>0`$ and let $`\mathrm{}(s)`$ lie in a fixed compact subset of $``$. Then
(10.11)
$$(1+|s|)^4\overline{E}_\mathrm{\Psi }(s,g)_{\mathrm{\Psi },N,T}0pt(g)^N,$$
for $`g𝔖(T)`$. In particular, if $`\mathrm{\Omega }𝐗`$ is compact, then $`s(1s)E_\mathrm{\Psi }(s,g)`$ is uniformly bounded in $`|\mathrm{}(s)|2,g\mathrm{\Omega }`$.
###### Proof.
We first claim that, for $`t`$, we have $`|(1+t^4)\overline{E}_\mathrm{\Psi }(N+1+it,g)|_\mathrm{\Psi }0pt(g)^{N+ϵ}`$. Indeed, by definition,
$$\overline{E}_\mathrm{\Psi }(s,g)=\underset{\gamma B(F)\backslash \mathrm{PGL}_2(F),\gamma B(F)}{}f_\mathrm{\Psi }(s,\gamma g)f_{\widehat{\mathrm{\Psi }}}(1s,g).$$
In view of Lem. 10.2, it will suffice to show that
(10.12)
$$\underset{\gamma B(F)\backslash \mathrm{PGL}_2(F),\gamma B(F)}{}0pt(\gamma g)^\sigma _ϵ0pt(g)^{1\sigma +ϵ},$$
which follows from (8.10) and (8.11).
Now (10.11) follows at once from the functional equation $`\overline{E}_\mathrm{\Psi }(s,g)=\overline{E}_{\widehat{\mathrm{\Psi }}}(1s,g)`$, the maximal modulus principle in the strip $`|\mathrm{}(s)|N+1`$, and the previous Lemma.<sup>21</sup><sup>21</sup>21To apply the maximal modulus principle in this context, one needs some a priori decay of $`\overline{E}_\mathrm{\Psi }`$, which follows easily from the corresponding properties of $`E_\mathrm{\Psi }`$ and $`f_\mathrm{\Psi }`$.
The second assertion (involving $`\mathrm{\Omega }`$) follows from (10.7) and (10.11). ∎
We now compute the Fourier coefficients of the Eisenstein series in general. Recall that $`e_F`$ is a fixed additive character of $`𝔸_F/F`$.
###### Lemma 10.5.
Set $`W_\mathrm{\Psi }(s,g)=_{xF\backslash 𝔸_F}E_\mathrm{\Psi }(s,n(x)g)e_F(x)𝑑x`$. Then, for $`\mathrm{}(s)>1`$,
(10.13)
$$W_\mathrm{\Psi }(a(y))=|y|^{1s}_{t𝔸_F^\times ,x𝔸_F}\mathrm{\Psi }(t,tx)e_F(xy)|t|^{2s}𝑑xd^\times t.$$
In particular, if $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$, then $`W_\mathrm{\Psi }=_vW_{\mathrm{\Psi }_v}`$, where for $`\mathrm{}(s)>1`$,
$$W_{\mathrm{\Psi }_v}(a(y))=|y|_v^{1s}_{tF_v^\times ,xF_v}\mathrm{\Psi }_v(t,tx)e_F(xy)|t|_v^{2s}𝑑xd^\times t,$$
for $`yF_v`$. Finally, if $`\mathrm{\Psi }_v(x,y)=\phi _1(x)\phi _2(y)`$, $`\omega _v`$ a character of $`F_v`$, and $`\mathrm{}(s^{})+|\mathrm{}(s)|1`$,
(10.14)
$$\begin{array}{c}_{yF_v^\times }W_{\mathrm{\Psi }_v}(a(y))|y|^s^{}\omega _v(y)d^\times y\hfill \\ \hfill =_{yF_v^\times }\phi _1(y)|y|^{s^{}+s}\omega _v(y)d^\times y_{yF_v^\times }\widehat{\phi _2}(y)|y|^{1+s^{}s}\omega _v(y)d^\times y,\end{array}$$
where $`\widehat{\phi _2}`$ is the Fourier transform, defined by $`\widehat{\phi _2}(y)=_{F_v}\phi _2(y)e_{F_v}(yt)𝑑t`$.
###### Proof.
By the Bruhat decomposition,
(10.15)
$$\begin{array}{c}W_\mathrm{\Psi }(s,g)=_{F\backslash 𝔸_F}e_F(x)\underset{\gamma B(F)\backslash \mathrm{PGL}_2(F)}{}f_\mathrm{\Psi }(\gamma n(x)g)\hfill \\ \hfill =_{𝔸_F}f_\mathrm{\Psi }(wn(x)g)e_F(x)𝑑x.\end{array}$$
Thus
(10.16)
$$\begin{array}{c}W_\mathrm{\Psi }(s,g)=|det(g)|^s_{x𝔸_F}_{t𝔸_F^\times }\mathrm{\Psi }((t,0)n(x)g)d^\times t|t|^{2s}e_F(x)\hfill \\ \hfill =|det(g)|^s_{t𝔸_F,x𝔸_F^\times }\mathrm{\Psi }((t,tx)g)|t|^{2s}e_F(x)𝑑xd^\times t.\end{array}$$
The claimed conclusion follows upon substituting $`g=a(y)`$, together with some routine computations. ∎
###### Remark 10.2.
Remark that $`W_\mathrm{\Psi }(s,g)`$ belongs to the Whittaker model of a certain induced representation of $`\mathrm{PGL}_2(𝔸_F)`$, namely that representation $`\pi (s)`$ induced from the character $`a(y)|y|^{s1/2}`$ of the maximal torus (unitary induction, so $`\pi (s)`$ is tempered for $`\mathrm{}(s)=1/2`$). This representation is the tensor product of local representations $`\pi _v(s)`$, analogously defined; these local representations are irreducible and generic for all $`s`$.
Thus (10.13) determines $`W_\mathrm{\Psi }`$ uniquely (the theory of the Kirillov model). Similarly the condition $`W_{\mathrm{\Psi }_v,s}(1)=1`$ uniquely determines the (spherical) vector $`W_{\mathrm{\Psi }_v,s}`$.
We finally remark that $`W_{\mathrm{\Psi }_v,s}`$, as $`\mathrm{\Psi }_v`$ ranges over all Schwarz-Bruhat functions on $`F_v^2`$ if $`v`$ is nonarchimedean, or over all standard functions if $`v`$ is archimedean, exhausts the Whittaker model of $`\pi (s)`$. Indeed, the set of such functions $`W_{\mathrm{\Psi }_v,s}`$ is a subspace of the Whittaker model of $`\pi (s)`$ that is stable under the action of the Hecke algebra of $`\mathrm{PGL}_2(F_v)`$; this action is irreducible, whence the result.
We recall that $`𝔡`$ denotes the different (Sec. 2.3) and we denote by $`\zeta _{F,v}(s)`$ or simply $`\zeta _v(s)`$ the local factor of the Dedekind $`\zeta `$-function of $`F`$ at the place $`v`$.
###### Corollary 10.1.
Suppose $`v`$ is nonarchimedean, and $`\mathrm{\Psi }_v`$ the characteristic function of $`𝔬_v^2`$. Then $`W_v(a(y))`$ satisfies
(10.17)
$$_{F_v^\times }W_v(a(y))|y|^s^{}d^\times y=q_v^{d_v(1+s^{}s)}\zeta _v(s+s^{})\zeta _v(1s+s^{}),$$
with $`d_v=v(𝔡)`$. Note that this specifies $`W_v`$, because it is $`K_v`$-invariant.
In particular, for each finite $`v`$ with $`v(𝔡)=0`$, the function $`W_v(g)`$ is the unique spherical Whittaker function on $`\mathrm{GL}_2(F_v)`$ with Hecke eigenvalue $`q_v^s+q_v^{1s}`$, and with $`W_v(1)=1`$.
As is evident from (10.6), the Eisenstein series themselves are not bounded. They belong to $`L^{2\epsilon }`$, but not $`L^2`$. To avoid some difficulties with growth, we shall use wave-packets of Eisenstein series. We now turn to their analysis.
### 10.2. Regularization of Eisenstein series on $`\mathrm{PGL}_2`$.
Our aim in this section is to show that an appropriate “wave packet” of the Eisenstein series $`E_\mathrm{\Psi }(g,s)`$ constructed in the previous section lies in $`L^{\mathrm{}}`$.
Note that, in Example 10.1 above $`E^{}(s,z)`$ differs from the usual unitary Eisenstein series by a factor $`\xi (2s)`$. This factor ensures that $`E^{}(s,z)`$ is holomorphic, but this causes an inconvenience at $`s=1/2`$, which will manifest itself in our construction of bounded wave-packets. Recall that this pole can be interpreted rather naturally: see footnote on p. 15.
Let $`\kappa >0`$, and let $`(\kappa )`$ be the family of functions holomorphic in an open neighbourhood of the strip $`\kappa \mathrm{}(s)1+\kappa `$, with rapid polynomial decay in vertical strips (i.e. $`sup_t(1+|t|)^N|h(\sigma +it)|`$ is bounded, for each $`N`$, by a continuous function of $`\sigma `$) and satisfying $`h(0)=h(\frac{1}{2})=h(1)=0`$. For each $`N`$ we have a norm $`_N`$ on $`(\kappa )`$ defined via:
(10.18)
$$h_N=_{\mathrm{}}^{\mathrm{}}\left(|h(1+\kappa +it)|+|h(\kappa +it)|\right)(1+|t|)^N𝑑t.$$
###### Lemma 10.6.
Let $`h(\kappa )`$, and set $`E_{h,\mathrm{\Psi }}(g)=_{\mathrm{}(s)=1+\kappa }h(s)E_\mathrm{\Psi }(g,s)𝑑s`$. Then:
$$E_{h,\mathrm{\Psi }}(g)_L^{\mathrm{}}_{\mathrm{\Psi },\kappa ,F}h_0$$
###### Proof.
In the notation of (10.10)
(10.19)
$$\begin{array}{c}E_{h,\mathrm{\Psi }}(g)=_{\mathrm{}(s)=1+\kappa }\overline{E}_\mathrm{\Psi }(s,g)h(s)𝑑s+_{\mathrm{}(s)=1+\kappa }h(s)f_\mathrm{\Psi }(s,g)𝑑s\hfill \\ \hfill +_{\mathrm{}(s)=1+\kappa }h(s)f_{\widehat{\mathrm{\Psi }}}(1s,g)𝑑s.\end{array}$$
Fix $`T>0`$ so that $`𝔖(T)`$ surjects onto $`𝐗`$ (see Section 8.2 for definitions). We will bound each term on the right-hand side of the above equation for $`g𝔖(T)`$.
By Lem. 10.4, the first term on the right-hand side is $`O_{\mathrm{\Psi },\kappa }(h_0)`$. By Lem. 10.2, the function $`f_{\widehat{\mathrm{\Psi }}}(1s,g)`$ is uniformly bounded above in the region $`\mathrm{}(s)=1+\kappa ,g𝔖(T)`$; thus the third term on the right-hand side is also $`O_{\mathrm{\Psi },\kappa }(h_0)`$.
As for the second term, we shift contours to the line $`\mathrm{}(s)=\kappa `$. The shift of contours is justified by the rapid decay of $`h(s)`$ along vertical lines and Lem. 10.2. Moreover, since $`h(0)=h(1/2)=0`$, the function $`sh(s)f_\mathrm{\Psi }(s,g)`$ has no poles in between the contours.
Applying Lem. 10.2 one more time to control the contour integral along $`\mathrm{}(s)=\kappa `$, we conclude. ∎
###### Remark 10.3.
Suppose $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$, and the $`\mathrm{\Psi }_v`$ are regarded as fixed for $`v`$ finite. Put $`\mathrm{\Psi }_f=_{v\mathrm{finite}}\mathrm{\Psi }_v`$, a Schwarz function on $`𝔸_{F,f}^2`$, and $`\mathrm{\Psi }_{\mathrm{}}=_{v\mathrm{infinite}}\mathrm{\Psi }_v`$. Then the above argument gives the slightly more explicit bound
(10.20)
$$E_{h,\mathrm{\Psi }}_L^{\mathrm{}}_{\kappa ,\mathrm{\Psi }_f}h_0𝒮(\mathrm{\Psi }_{\mathrm{}})$$
where $`𝒮`$ is a Schwarz norm on $`F_{\mathrm{}}^2`$. This follows by explicating the above argument, taking into account the last sentence of the proof of Lem. 10.2. Indeed, one obtains even the corresponding bound for Sobolev norms, namely
(10.21)
$$S_{\mathrm{},d,\beta }(E_{h,\mathrm{\Psi }})_{\kappa ,\mathrm{\Psi }_f}h_0𝒮(\mathrm{\Psi }_{\mathrm{}})$$
for an appropriate Schwarz norm of $`F_{\mathrm{}}^2`$. One deduces this from (10.20) upon noting that, if $`𝒟`$ belongs to the universal enveloping algebra of $`\mathrm{SL}_2(F_{\mathrm{}})`$, then, by (10.4), $`𝒟E_\mathrm{\Psi }(s,g)=E_{𝒟\mathrm{\Psi }}(s,g)`$, so also $`𝒟E_{h,\mathrm{\Psi }}=E_{h,𝒟\mathrm{\Psi }}`$. It is then easy to check that a Schwarz norm of $`𝒟\mathrm{\Psi }_{\mathrm{}}`$ is bounded by a Schwarz norm of $`\mathrm{\Psi }_{\mathrm{}}`$.
### 10.3. Regularization of Eisenstein series on $`\mathrm{PGL}_2\times \mathrm{PGL}_2`$.
In this section we carry out the analogue of Lem. 10.6 in the context of $`\mathrm{PGL}_2\times \mathrm{PGL}_2`$ (this amounts to regularizing the rank $`2`$ Eisenstein series on $`\mathrm{PGL}_2\times \mathrm{PGL}_2`$).
To ease the reader’s path, we briefly mention what the point of this section is in classical notation: Suppose $`h(s_1,s_2)`$ is holomorphic in two variables inside the square $`|\mathrm{}(s_1)|+|\mathrm{}(s_2)|1/2+\kappa `$, and, moreover, $`h(s_1,s_2)`$ has zeroes along the six planes defined by any of the linear constraints $`s_1=0,s_1=1/2,s_1=1/2,s_2=0,s_2=1/2,s_2=1/2`$.
Define the wave-packet $`E_h(z_1,z_2)`$ on $`\mathrm{SL}_2()\backslash \times \mathrm{SL}_2()\backslash `$ via
$$E_h(z_1,z_2)=_{t,t^{}}h(it_1,it_2)E^{}(1/2+it,z_1)E^{}(1/2+it^{},z_2)𝑑t𝑑t^{}.$$
Here $`E^{}`$ is as in Example 10.1. We shall show – under mild decay conditions on $`h`$ –that $`E_h(z_1,z_2)`$ is majorized, on the product of two fundamental regions, by $`A(y_1,y_2):=\frac{\sqrt{y_1y_2}}{y_1^{1/2+\kappa }+y_2^{1/2+\kappa }}`$. Since $`_{y_11,y_21}A(y_1,y_2)^4\frac{dy_1dy_2}{y_1^2y_2^2}`$ is finite, $`E_h`$ lies in $`L^4`$, and even in $`L^{4+ϵ}`$ for $`ϵ`$ small.
As the reader may verify at this point, the majorization is little more than an exercise in complex integration, using the fact that the large contribution to the Eisenstein series comes from the constant term.
We will need to repeatedly shift contours in the setting of a function of two complex variables. To clarify matters, we state the following Lemma, which we will use repeatedly without explicitly invoking it.
###### Lemma 10.7.
Suppose $`U^2`$ is an open domain and $`f(z_1,z_2)`$ a holomorphic function on the complex domain $`\{(z_1,z_2)^2:(\mathrm{}(z_1),\mathrm{}(z_2))U\}`$. Suppose moreover that there is, a continuous function $`M:U`$ such that
(10.22)
$$\underset{(t_1,t_2)^2}{sup}|f(\sigma _1+it_1,\sigma _2+it_2)|(1+|t_1|+|t_2|)^3M(\sigma _1,\sigma _2).$$
Then the function
(10.23)
$$(\sigma _1,\sigma _2)_{\mathrm{}(z_1)=\sigma _1,\mathrm{}(z_2)=\sigma _2}f(z_1,z_2)𝑑z_1𝑑z_2$$
is locally constant on $`U`$.
We omit the easy proof.
We will now introduce a family of normed spaces $`^{(2)}(\kappa )`$. In fact, the spaces themselves are independent of $`\kappa `$, but the norm depends on $`\kappa `$. These are spaces of holomorphic functions in two variables $`z_1,z_2`$; and the norm, roughly speaking, controls the behavior of $`h`$ when the real parts of $`(z_1,z_2)`$ lie in the square $`|\mathrm{}(z_1)|+|\mathrm{}(z_2)|1/2+\kappa `$.
###### Definition 10.1.
Let $`0<\kappa <1`$. Let $`^{(2)}(\kappa )`$ be the family of functions $`h(z_1,z_2)`$ in two complex variables, holomorphic in a neighbourhood of $`(0,0)`$, and satisfying:
1. Write $`h^{}=\frac{h(z_1,z_2)}{z_1z_2(1/4z_1^2)(1/4z_2^2)}`$. Then $`h^{}`$, originally a meromorphic function in a neigbourhood of $`0`$, extends to a holomorphic function in the strip $`\{z_1:|\mathrm{}(z_1)|2\}\times \{z_2:|\mathrm{}(z_2)|2\}`$.
2. Growth condition: for every $`N0`$,
$$\underset{(\sigma ,\sigma ^{})[2,2]^2}{sup}\underset{t,t^{}^2}{sup}(1+|t|+|t^{}|)^Nh(\sigma +it,\sigma ^{}+it)<\mathrm{}$$
For each $`N`$ we introduce a norm on $`H^{(2)}(\kappa )`$ via:
(10.24) $`h_N={\displaystyle _{(t,t^{})^2}}{\displaystyle \underset{ϵ_1,ϵ_2\{\pm 1\}}{}}(1+|t|+|t^{}|)^N`$
$`\left(|h^{}(ϵ_1(1/2+\kappa )+it,it^{})|+|h^{}(it,ϵ_2(1/2+\kappa )+it^{})|\right)dtdt^{}.`$
###### Lemma 10.8.
For $`hH^{(2)}(\kappa )`$, put
$$E_{h,\mathrm{\Psi },\mathrm{\Psi }^{}}(g_1,g_2)=_{\mathrm{}(t)=0}_{\mathrm{}(t^{})=0}h(t,t^{})E_\mathrm{\Psi }(g_1,1/2+t)E_\mathrm{\Psi }^{}(g_2,1/2+t^{}).$$
Then
$$E_{h,\mathrm{\Psi },\mathrm{\Psi }^{}}(x_1,x_2)_{\mathrm{\Psi },\mathrm{\Psi }^{}}h_0\frac{0pt(x_1)^{1/2}0pt(x_2)^{1/2}}{0pt(x_1)^{1/2+\kappa }+0pt(x_2)^{1/2+\kappa }}.$$
###### Proof.
We may assume that $`h_0=1`$. Let notations be as established prior to Lem. 10.4. We will proceed as in Lem. 10.6, expanding $`E_\mathrm{\Psi }`$ via (10.10).
It will suffice to give an upper bound, in absolute value, for each of:
(10.25) $`I_0(g_1,g_2)={\displaystyle _{t,t^{}}}h(t,t^{})\overline{E}_{\mathrm{\Psi }_1}(g_1,1/2+t)\overline{E}_{\mathrm{\Psi }_2}(g_2,1/2+t^{})𝑑t𝑑t^{}`$
(10.26) $`I_1(g_1,g_2)={\displaystyle _{t,t^{}}}h(t,t^{})\overline{E}_{\mathrm{\Psi }_1}(g_1,1/2+t)f_{\mathrm{\Psi }_2}(g_2,1/2\pm t^{})𝑑t𝑑t^{}`$
(10.27) $`I_2(g_1,g_2)={\displaystyle _{t,t^{}}}h(t,t^{})f_{\mathrm{\Psi }_1}(g_1,1/2\pm t)\overline{E}_{\mathrm{\Psi }_2}(g_2,1/2+t^{})𝑑t𝑑t^{}`$
(10.28) $`I_3(g_1,g_2)={\displaystyle _{t,t^{}}}h(t,t^{})f_{\mathrm{\Psi }_1}(g_1,1/2\pm t)f_{\mathrm{\Psi }_2}(g_2,1/2\pm t^{})𝑑t𝑑t^{},`$
whenever $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ are Schwarz functions on $`𝔸_F^2`$, and in each case the contour of integration is the surface $`\mathrm{}(t)=\mathrm{}(t^{})=0`$. Moreover, in view of condition (1) in Def. 10.1, each integrand extends to a holomorphic function of $`(t,t^{})`$ in the region $`|\mathrm{}(t)|1/2+\kappa ,|\mathrm{}(t^{})|1/2+\kappa `$.
The bound $`|I_0|0pt(g_1)^N0pt(g_2)^N`$ follows from Lem. 10.4, whereas the bounds $`|I_1|0pt(g_1)^N0pt(g_2)^\kappa `$ and $`|I_2|0pt(g_2)^N0pt(g_1)^\kappa `$ follow from moving the $`t^{}`$ (in the case of $`I_1`$) integral to the contour $`\mathrm{}(t^{})=\pm (1/2+\kappa )`$, applying Lem. 10.4 and Lem. 10.2.
We now turn to $`I_3`$. We shall consider the case where both signs are $`+`$, the other cases being similar with appropriate interchanges of sign. Thus set
(10.29) $`Z(t,t^{})=h(t,t^{})f_{\mathrm{\Psi }_1}(g_1,1/2+t)f_{\mathrm{\Psi }_2}(g_2,1/2+t^{})`$
$`=h^{}(t,t^{})t(1/4t^2)f_{\mathrm{\Psi }_1}(g_1,1/2+t)t^{}(1/4t^2)f_{\mathrm{\Psi }_2}(g_2,1/2+t^{})`$
In view of Lem. 10.2, the function $`Z(t,t^{})`$ satisfies the conditions for $`f`$ in Lem. 10.7. We apply Lem. 10.7 to shift the contour to $`\mathrm{}(t)=1/2\kappa ,\mathrm{}(t^{})=0`$. Now Lem. 10.2 implies that $`|_{\mathrm{}(t)=1/2\kappa ,\mathrm{}(t^{})=0}Z(t,t^{})|0pt(g_2)^{1/2}0pt(g_1)^\kappa `$. A similar bound holds with $`(g_1,g_2)`$ interchanged, so in fact we have the stronger bound $`|Z(t,t^{})|\mathrm{min}(0pt(g_2)^{1/2}0pt(g_1)^\kappa ,0pt(g_1)^{1/2}0pt(g_2)^\kappa )`$. This may also be written $`|Z(t,t^{})|\frac{0pt(g_1)^{1/2}0pt(g_2)^{1/2}}{0pt(g_1)^{1/2+\kappa }+0pt(g_2)^{1/2+\kappa }}`$.
Similar considerations apply to the terms in $`I_3`$ corresponding to other choices of sign, so we conclude that $`|I_3|\frac{0pt(g_1)^{1/2}0pt(g_2)^{1/2}}{0pt(g_1)^{1/2+\kappa }+0pt(g_2)^{1/2+\kappa }}`$. ∎
###### Lemma 10.9.
Let notations be as in the previous Lemma. For any $`p<\frac{4}{12\kappa }`$, any $`d,\beta >0`$, there exists $`N`$ such that
$$S_{p,d,\beta }(E_{h,\mathrm{\Psi }})_{\mathrm{\Psi },\kappa ,p,\beta }h_N$$
###### Proof.
Indeed, we note that
$$_{y_1,y_21}\left(\frac{\sqrt{y_1y_2}}{y_1^{1/2+\kappa }+y_2^{1/2+\kappa }}\right)^p\frac{dy_1dy_2}{y_1^2y_2^2}<\mathrm{}$$
whenever $`p<\frac{4}{12\kappa }`$. We apply the previous Lemma and reduction theory to conclude. ∎
## 11. Background on integral representations of $`L`$-functions.
The purpose of this section is as follows. The geometric method we have explained in the text yields upper bounds for certain periods; to obtain subconvexity, we need to know that $`L`$-functions can be expressed in terms of these periods. This is the whole point of the theory of integral representations of $`L`$-functions; however, we cannot quite simply quote from that theory, as we often need e.g. some analytic control on the choice of test vector for which there is no readily available reference.
On occasion we have only sketched proofs in this section, as they amount to simple explications of standard techniques such as the Rankin-Selberg method, and moreover they are in some sense irrelevant to the main point of this paper (which is to bound periods, not $`L`$-functions!)
### 11.1. Cuspidal triple product $`L`$-functions.
###### Hypothesis 11.1.
Let $`\pi _2`$ and $`\pi _3`$ be fixed automorphic cuspidal representations of $`\mathrm{PGL}_2`$ over $`F`$. Let $`\pi _1`$ be an automorphic cuspidal representation, whose finite conductor is a prime ideal, prime to the finite conductors of $`\pi _2`$ and $`\pi _3`$. Suppose that $`\pi _{1,\mathrm{}}`$ (the representation of $`\mathrm{PGL}_2(F_{\mathrm{}})`$ underlying $`\pi _1`$) is restricted to a bounded set; let $`\phi _1`$ be the new vector in $`\pi _1`$.
Then there exists finite collections of vectors $`_2\pi _2,_3\pi _3`$ so that, for any such $`\pi _1`$, there exist $`\phi _j_j(j=2,3)`$ with
(11.1)
$$\frac{L(\frac{1}{2},\pi _1\pi _2\pi _3)}{\left|_𝐗\phi _2(xa([𝔭]))\phi _3(x)\phi _1(x)𝑑x\right|^2}_{ϵ,F,\pi _{1,\mathrm{}}}N(𝔭)^{1+ϵ}$$
Note that no claim is made about the dependence of the constants in (11.1) on $`\pi _2,\pi _3`$ or the bounded set containing $`\pi _{1,\mathrm{}}`$; presumably with enough effort one could obtain polynomial dependence on the conductors.
The proof of Hypothesis 11.1 should be, we believe, an elaborate but routine computation of certain $`p`$-adic integrals; this has not carried out, but we expect it to be valid.
In the case when $`F=`$ and $`\pi _1,\pi _2,\pi _3`$ holomorphic Hypothesis 11.1 may follow (in a slightly modified form, replacing $`\mathrm{PGL}_2`$ by a division algebra) from the work of Böcherer and Schulze-Pillot. In any case, there exist good heuristic reasons to believe the Hypothesis: based on a computation of the size of the relevant family, or alternately it is true if one of the $`\pi _j`$ is Eisenstein.
### 11.2. Rankin-Selberg convolutions.
#### 11.2.1. The Rankin-Selberg integral representation.
Let $`\pi _1,\pi _2`$ be two automorphic representations, with $`\pi _2`$ cuspidal.
Let $`\mathrm{\Psi }_v`$ be a Schwarz-Bruhat function on $`F_v^2`$ such that, for almost all $`v`$, $`\mathrm{\Psi }_v`$ is the characteristic function of $`𝔬_{F_v}^2`$. Put $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$, a Schwarz function on $`𝔸_F^2`$. Let $`\phi _j`$ belong to the space of $`\pi _j`$ for $`j=1,2`$ and put
(11.2)
$$I(\phi _1,\phi _2,\mathrm{\Psi },s)=_𝐗\phi _1(g)\phi _2(g)E_\mathrm{\Psi }(s,g)𝑑g$$
Unwinding, we see that for $`\mathrm{}(s)>1`$:
(11.3)
$$\begin{array}{c}I(\phi _1,\phi _2,\mathrm{\Psi },s)=c_F_{B(F)\backslash \mathrm{PGL}_2(𝔸_F)}f_\mathrm{\Psi }(s,g)\phi _1(g)\phi _2(g)\hfill \\ \hfill =c_F_{B(F)\backslash \mathrm{PGL}_2(𝔸_F)}f_\mathrm{\Psi }(s,g)\left(_{nN(F)\backslash N(𝔸_F)}\phi _1(ng)\phi _2(ng)𝑑n\right)𝑑g\end{array}$$
Here the constant $`c_F`$ arises from change of measure: the measure on $`𝐗`$ is the $`\mathrm{PGL}_2(𝔸_F)`$-invariant probability measure, which is not the same as the quotient measure from $`\mathrm{PGL}_2(𝔸_F)`$. Note that $`c_F`$ will be unimportant in our arguments, as it depends only on $`F`$ and we are only interested in bounds.
Put $`W_1(g)=_{F\backslash 𝔸_F}\phi _1(n(x)g)e_F(x)𝑑x`$, and define $`W_2`$ similarly but with $`e_F`$ replaced by $`\overline{e_F}`$. Recall that our normalizations are so that the volume of $`𝔸_F/F`$ is $`1`$. Fourier inversion shows that $`\phi _i(g)=_{\alpha F^\times }W_i(a(\alpha )g)`$ if $`\phi _i`$ is cuspidal. Thus, as long as one of $`\phi _1,\phi _2`$ is cuspidal, we see that:
(11.4)
$$\begin{array}{c}I(\phi _1,\phi _2,\mathrm{\Psi },s)=c_F_{B(F)\backslash \mathrm{PGL}_2(𝔸_F)}f_\mathrm{\Psi }(s,g)\left(\underset{\alpha F^\times }{}W_1(a(\alpha )g)W_2(a(\alpha )g)\right)\hfill \\ \hfill =c_F_{N(𝔸_F)\backslash \mathrm{PGL}_2(𝔸_F)}W_1(g)W_2(g)f_\mathrm{\Psi }(s,g)𝑑g\end{array}$$
If $`\phi _1,\phi _2`$ are pure tensors, then there is a corresponding product decomposition $`W_1=_vW_{1,v},W_2=_vW_{2,v}`$, where $`W_{j,v}`$ belongs to the local Whittaker model of $`\pi _{j,v}`$, a representation of $`\mathrm{PGL}_2(F_v)`$. In that case,
(11.5)
$$I(\phi _1,\phi _2,\mathrm{\Psi },s)=c_F\underset{v}{}I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s),$$
where
(11.6)
$$\begin{array}{c}I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s)=\hfill \\ \hfill _{N(F_v)\backslash \mathrm{PGL}_2(F_v)}W_1(g_v)W_2(g_v)\left(|det(g_v)|_v^s_{tF_v^\times }\mathrm{\Psi }((0,t)g_v)|t|^{2s}d^\times t\right)\end{array}$$
We note that the bracketed quantity, defined a priori for $`g_v\mathrm{GL}_2(F_v)`$, descends to $`\mathrm{PGL}_2(F_v)`$.
###### Lemma 11.1.
Suppose $`W_{1,v},W_{2,v}`$ the new vectors associated to spherical representations $`\pi _{1,v},\pi _{2,v}`$, $`\mathrm{\Psi }_v`$ is the characteristic function of $`𝔬_v^2`$, and $`e_{F_v}`$ is unramified. Then $`I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v)=L_v(s,\pi _{1,v}\pi _{2,v})`$.
If $`W_{1,v},W_{2,v}`$ are nonzero and $`\mathrm{PGL}_2(𝔬_v)`$-invariant, $`\mathrm{\Psi }_v`$ as above, but $`e_{F_v}`$ is possibly ramified, then $`I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v)=aq_v^{ks}L_v(s,\pi _{1,v}\times \pi _{2,v})`$ where $`k`$ is so that $`e_{F_v}`$ is trivial on $`\varpi _v^k𝔬_v`$ but not on $`\varpi _v^{k1}𝔬_v`$. Moreover, $`a=1`$ if $`W_{1,v}(\varpi _v^k)=W_{2,v}(\varpi _v^k)=1`$.
Suppose $`W_{1,v},W_{2,v}`$ are the new vectors associated to $`\pi _{1,v}`$ spherical and $`\pi _{2,v}`$ a Steinberg representation, and that $`e_{F_v}`$ is unramified. Then, with $`\mathrm{\Psi }_v`$ the characteristic function of $`𝔬_v^2`$, we have:
$$I_v(\pi _{1,v}\left(\begin{array}{cc}1& 0\\ 0& \varpi _v\end{array}\right)W_{1,v},W_{2,v},\mathrm{\Psi }_v)=\pm \frac{q_v^s}{q_v+1}L(s,\pi _{1,v}\pi _{2,v}).$$
###### Proof.
See \[17, Thm 15.9\] for the first assertion. The second assertion is an easy consequence. See Sec. 11.3 for the final assertion. ∎
Applying the Iwasawa decomposition to (11.6) yields the equivalent
(11.7)
$$\begin{array}{c}I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s)=_{yF_v^\times ,kK_v}W_1(a(y)k)W_2(a(y)k)|y|^{s1}d^\times y\hfill \\ \hfill \left(_{tF_v^\times }\mathrm{\Psi }((0,t)k)|t|^{2s}d^\times t\right)\end{array}$$
#### 11.2.2. Topologizing the space of local representations.
The results in provide “good” test vectors for the functionals $`I_v`$ when the local representations $`\pi _{1,v},\pi _{2,v}`$ are fixed. On the other hand, we will need such results with some mild uniformity in $`\pi _{1,v}`$. One can certainly extract the stronger results from the proofs in . For now, we will proceed by deducing the results “by continuity”; to do this, we will need to define the topology on the space of possible $`\pi _{1,v}`$. The considerations that follow are not really very crucial; it would be better simply to explicate the implicit dependences in .
Let $`𝒮`$ be a finite set of irreducible (continuous) representations of $`K_v=\mathrm{PGL}_2(𝔬_v)`$. For any representation $`W`$ of $`K_v`$, we denote by $`W^𝒮`$ that subspace of $`W`$ consisting of vectors whose $`K_v`$-span contains only irreducibles that belong to $`𝒮`$. We shall say that elements of $`W^𝒮`$ are of type $`𝒮`$.
Let $`𝒢_v`$ be the set of isomorphism classes of generic irreducible representations of $`\mathrm{PGL}_2(F_v)`$. If $`\pi `$ is a generic irreducible representation which is a discrete series or supercuspidal, we shall define it to be isolated. Otherwise, $`\pi `$ is induced from two quasicharacters $`\mu ,\nu :F_v^\times `$. For $`s`$, let $`(\pi (s),V(s))`$ be the representation induced from the quasicharacters $`\mu ||_v^s,\nu ||_v^s`$. Then $`\pi (s)`$ is generic for all $`s`$ and irreducible in a neighbourhood of $`0`$. We shall topologize $`𝒢_v`$ in such a way that sets of the form $`[\pi (s)]`$, for $`|s|<\epsilon `$ form a basis.
If $`E𝒢_v`$ is a closed subset that is bounded (when considered as a subset of the set of isomorphism classes of irreducible admissible representations, and bounded in the sense of Sec. 2.12.3), then $`E`$ is compact, as one checks by direct verification.
For each $`\pi 𝒢`$, we have a Whittaker model $`𝒲(\pi )`$. Consider a function $`\pi W_\pi `$, that assigns to each $`\pi 𝒢_v`$ an element $`W_\pi `$ of its Whittaker model. We shall say that such an assignment $`\pi W_\pi `$ is continuous if there exists a neighbourhood of each $`\pi `$, which we may assume to be of the form, $`\{\pi (s):|s|<\epsilon \}`$, and a set $`𝒮`$ of irreducible representations of $`K_v`$ so that:
1. $`W_{\pi (s)}`$ is of type $`𝒮`$, for each $`|s|<\epsilon `$
2. The assignment $`sW_{\pi (s)}(g)`$ is continuous for each $`g\mathrm{PGL}_2(F_v)`$, uniformly for $`g`$ in any fixed compact.
It can be verified that if $`W_0`$ is an element of the Whittaker model of $`\pi _0`$, there exists a continuous assignment $`\pi W_\pi `$ in a neighbourhood of $`\pi _0`$ which has the value $`W_0`$ at $`\pi _0`$.
The requirement (2) is not very strong, as it does not impose any uniformity on all of $`\mathrm{PGL}_2(F_v)`$. However, in every context we shall consider, the necessary uniformity in $`g`$ is automatic. Let us sketch how one can prove such results. Assume that $`v`$ is finite; the infinite case is similar although more technically involved. One first observes that if $`\pi W_\pi `$ is a continuous assignment on some open set, then, for a fixed character $`\chi _v`$ of $`F_v^\times `$, the quotient $`\frac{_{yF_v^\times }W_\pi (a(y))\chi _v(y)|y|^{s1/2}d^\times y}{L_v(s,\pi _v\chi _v)}`$ is a polynomial of the form $`_{k=N}^Nc_kq_v^{ks}`$; moreover, the degree $`N`$ is locally bounded as $`\pi `$ varies, and all the coefficients $`c_k`$ can be taken to depend continuously on $`\pi `$. To verify the local boundedness of the degree – which requires only property (1) above – one just notes that there is (locally) a fixed $`M`$ such that $`W_\pi (a(y))`$ vanishes for $`|y|_v>M`$; this, together with the functional equation, gives the local boundedness. To see that the coefficients vary continuously, it suffices to check that, for any fixed integer $`t`$, the integral $`_{v(y)=t}W_\pi (a(y))\chi _v(y)|y|^{s1/2}d^\times y`$ varies continuously, which follows from the definition of continuity for the assignment $`\pi W_\pi `$. The archimedean case proceeds similarly, but one replaces the role of polynomials in $`q_v^{\pm s}`$ by functions of the form $`c^sP(s)`$, where $`P`$ is a polynomial and $`c`$.
In the next few pages, we will make certain claims regarding the continuity of various integrals involving $`W_\pi `$, if $`\pi W_\pi `$ is a continuous assignment. One can reduce all the claimed continuity statements (by standard “Mellin transform” arguments) to the result just discussed. We will omit the details.
#### 11.2.3. Choice of test vectors
###### Lemma 11.2.
Let notation be as above.
1. The quotient
$$\mathrm{\Xi }_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s):=\frac{I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s)}{L_v(s,\pi _{1,v}\pi _{2,v})}$$
is holomorphic in $`s`$. If $`v`$ is nonarchimedean, $`\mathrm{\Xi }_v`$ is a polynomial in $`q_v^{\pm s}`$; if $`v`$ is archimedean and $`\mathrm{\Psi }_v`$ is standard, then $`\mathrm{\Xi }_v|a_v|_v^{2s}`$ is a polynomial in $`s`$.
2. For any fixed $`s_0`$ we may choose data $`(W_{1,v},W_{2,v},\mathrm{\Psi }_v)`$ of the type described in (1) with $`\mathrm{\Xi }_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s_0)0`$.
3. If $`\pi _{2,v}`$ is regarded as fixed, and $`\pi _{1,v}`$ remains within a fixed compact subset of $`𝒢_v`$ consisting entirely of unitarizable representations, then there exists a constant $`C`$ depending on the compact set so that one may choose data as in (2) in such a way that:
1. $`\mathrm{\Psi }_v`$ and $`W_{2,v}`$ may both be chosen from a finite list of size $`C`$;
2. $`_{F_v^\times }|W_{1,v}(a(y))|^2d^\times yC`$;
3. $`|\mathrm{\Xi }_v(s_0)|1`$ and, for all $`s`$, we have $`|\mathrm{\Xi }_v(s)|C^{|\mathrm{}(s)|}(1+|s|)^C`$.
###### Proof.
The first two assertions are in .
We will only sketch the last assertion. It can be also be proved directly by exhibiting such data by explicating the arguments of . In any case, we start by noting: Given any continuous assignment $`\pi _{1,v}W_{1,v},\pi _{2,v}W_{2,v}`$, the function $`\mathrm{\Xi }_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s)`$ and $`|W_{i,v}(a(y)|^2d^\times y`$ all varies continously in $`\pi _{1,v},\pi _{2,v}`$. This assertion can be deduced by the methods explained in Section 11.2.2. Here, when we speak of $`\mathrm{\Xi }_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v,s)`$ varying continuously, we mean this in the “strong sense”, i.e the statement that $`\mathrm{\Xi }_v`$ can be expressed as a polynomial in $`q_v^s`$ (nonarchimedean case) or $`b^sP(s)`$ where $`P`$ is polynomial (archimedean case), so that all coefficients vary continuously with $`\pi _{1,v},\pi _{2,v}`$.
Now, given fix momentarily $`\pi _{1,v}`$ and $`\pi _{2,v}`$ and suppose we have chosen data $`(W_{1,v},W_{2,v},\mathrm{\Psi }_v)`$ as in (2). Extend $`W_{1,v}`$ to a continuous assignment $`\pi W_\pi `$ in a neighbourhood of $`\pi _{1,v}`$. By the remarks above, $`(W_\pi ,W_{2,v},\mathrm{\Psi }_v)`$ will satisfy (3b) and (3c), for a suitable constant $`C`$, whenever $`\pi `$ belongs to a sufficiently small neighbourhood of $`\pi _{1,v}`$. Now a compactness argument demonstrates (3). ∎
We emphasize again that (11.6) is valid so long as one of $`\pi _1,\pi _2`$ is cuspidal.
###### Lemma 11.3.
Let $`\pi `$ be an automorphic cuspidal representation in $`L^2(𝐗)`$, and let $`\phi \pi `$ be so that $`W_\phi :=_{F\backslash 𝔸_F}e_F(x)\phi (n(x)g)`$ factorizes as a product $`_vW_v(g)`$. Then, for a certain constant absolute constant $`c`$
(11.8)
$$_𝐗|\phi (g)|^2𝑑g=c\mathrm{Res}_{s=1}\mathrm{\Lambda }(s,\pi \stackrel{~}{\pi })\underset{v}{}\frac{_{F_v^\times }|W_v(a(y))|^2d^\times y}{L_v(s,\pi _v\stackrel{~}{\pi }_v)}$$
###### Proof.
This follows by taking the residue of $`I(\phi ,\overline{\phi },\mathrm{\Psi },s)`$ at $`s=1`$. Indeed, this residue equals, up to a constant depending only on the measure normalization, $`\left(_𝐗|\phi (g)|^2𝑑g\right)\left(_{𝔸_F^2}\mathrm{\Psi }(x,y)𝑑x𝑑y\right)`$ (see discussion of properties of $`E_\mathrm{\Psi }`$ after (10.3).
On the other hand, by (11.5) and (11.6) $`I(\phi ,\overline{\phi },\mathrm{\Psi },s)`$ may be written as a product $`c_F_vI_v(W_v,\overline{W_v},\mathrm{\Psi }_v,s)`$, where each $`I_v`$ is given by (11.7). The integral $`_{F_v^\times }|W_v(a(y)k)|^2d^\times y`$ is independent of $`kK_v`$, so $`I_v(W_v,\overline{W_v},\mathrm{\Psi }_v,1)`$ factors as the product of $`_{yF_v^\times }|W_v(a(y))|^2d^\times y`$ and $`_{tF_v^\times ,kK_v}\mathrm{\Psi }_v((0,t)k)|t|^2𝑑t`$. The latter integral differs from $`_{F_v^2}\mathrm{\Psi }_v(x,y)𝑑x𝑑y`$ by a factor that depends only on the normalizations of measure; moreover, this factor equals $`(1q_v^2)^1`$ for almost all $`v`$, so the product of these factors is convergent. The conclusion easily follows. ∎
We now specialize to the cases of interest. Fix $`\pi _1`$. We vary $`\pi _2:=\pi `$ through a sequence of automorphic cuspidal representations with prime conductor $`𝔭`$, prime to the conductor of $`\pi _1`$. In particular, the local constituent of $`\pi `$ at $`𝔭`$ is a special representation. We denote by $`\pi _{\mathrm{}}`$ the representation of $`\mathrm{GL}_2(F_{\mathrm{}})`$ underlying the representation $`\pi `$.
###### Lemma 11.4.
Suppose the archimedean constituent $`\pi _{\mathrm{}}`$ belongs to a bounded subset of $`\widehat{\mathrm{PGL}_2(F_{\mathrm{}})}`$ (in what follows the implicit constants may depend on this subset) and regard $`\pi _1`$ as being fixed.
Let $`s_0`$. There exists a fixed finite set $``$ of Schwarz Bruhat functions and a real number<sup>22</sup><sup>22</sup>22depending on $`\pi _1`$ and the choice of bounded subset of $`\widehat{\mathrm{PGL}_2(F_{\mathrm{}})}`$ $`C>0`$ so that, for any such $`\pi `$,
There exist vectors $`\phi _1\pi _1,\phi \pi `$ and $`\mathrm{\Psi }`$ so that
$$\mathrm{\Phi }(s):=\mathrm{N}(𝔭)^{1s}\frac{I(a([𝔭])\phi _1,\phi ,\mathrm{\Psi },s)}{\mathrm{\Lambda }(s,\pi _1\pi )}$$
is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(s_0)|1`$ and $`|\mathrm{\Phi }(s)|C^{|\mathrm{}(s)|}(1+|s|)^C`$;
2. At any nonarchimedean place $`v`$ such that $`\pi _1`$ and $`\pi `$ are both unramified, both $`\phi `$ and $`\phi _1`$ are invariant by $`\mathrm{PGL}_2(𝔬_{F_v})`$;
3. $`\phi _1_L^{\mathrm{}}`$ is $`O(1)`$.
4. $`\phi _{L^2(𝐗)}_ϵ\mathrm{N}(𝔭)^ϵ`$.
###### Proof.
We first choose local data. For each place where $`e_{F,v}`$ and $`\pi _1`$ are not ramified, we take $`W_v`$ (resp. $`W_{v,1}`$) to be the new vector in the Whittaker model of $`\pi _v`$ (resp $`\pi _{1,v}`$). We put $`\mathrm{\Psi }_v`$ to be the characteristic function of $`𝔬_v^2`$.
Let $``$ be the set of remaining $`v`$. For $`v`$, the assumptions show that $`\pi _v`$ is restricted to a bounded set. We choose $`W_v`$, $`W_{v,1},\mathrm{\Psi }_v`$ for $`v`$ according to Lem. 11.2. Finally we choose $`\phi `$ so that $`_{xF\backslash 𝔸_F}e_F(x)\phi (n(x)g)=_vW_v(g)`$, and similarly for $`\phi _1`$, and take $`\mathrm{\Psi }=_v\mathrm{\Psi }_v`$. The first two assertions of the Lemma are immediate.
To bound the $`L^2`$ norm of $`\phi `$, use Lem. 11.2 (3b), Lem. 11.3, and Iwaniec’s bounds on $`L`$-functions near $`1`$. As for $`\phi _1`$, it in fact belongs to a fixed finite set of cusp forms, so the third assertion is immediate. ∎
We continue to keep $`\pi `$ an automorphic cuspidal representation of $`\mathrm{PGL}_2(𝔸_F)`$ with prime conductor.
###### Lemma 11.5.
Suppose $`\pi _{\mathrm{}}`$ belongs to a bounded subset of $`\widehat{\mathrm{PGL}_2(F_{\mathrm{}})}`$ (in what follows the implicit constants may depend on this bounded subset).
Let $`t_0,t_0^{}`$. There exists a fixed finite set $``$ of Schwarz Bruhat functions and a real number $`C>0`$ so that:
There exist vectors $`\phi \pi ,\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ so that:
(11.9)
$$\mathrm{\Phi }(t,t^{})=\mathrm{N}(𝔭)^{1/2t}\frac{_𝐗\phi (g)E_{\mathrm{\Psi }_1}(g,\frac{1}{2}+t)E_{\mathrm{\Psi }_2}(ga([𝔭]),\frac{1}{2}+t^{})𝑑g}{\mathrm{\Lambda }(\frac{1}{2}+t+t^{},\pi )\mathrm{\Lambda }(\frac{1}{2}+tt^{},\pi )}$$
is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(t_0,t_0^{})|1`$ and $`|\mathrm{\Phi }(t,t^{})|(1+|t|+|t^{}|)^CC^{|\mathrm{}(t)|+|\mathrm{}(t^{})|}`$.
2. For any nonarchimedean place $`v`$, each $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ is invariant by $`\mathrm{PGL}_2(𝔬_{F_v})`$.
3. $`\phi _{L^2(𝐗)}_ϵ\mathrm{N}(𝔭)^ϵ`$.
###### Proof.
The proof is similar to that of the previous Lemma; recall that (11.6) was valid as long as one of $`\pi _1,\pi _2`$ were cuspidal.
Let $`\mathrm{\Psi }_2^{}`$ be the translate of the Schwarz function $`\mathrm{\Psi }_2`$ by $`a([𝔭])`$. Then by (10.4),
$$E_{\mathrm{\Psi }_2^{}}(s,g)=\mathrm{N}(𝔭)^sE_{\mathrm{\Psi }_2}^{a([𝔭])}(s,g)$$
Suppose $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ factorize as $`_v\mathrm{\Psi }_{1,v},_v\mathrm{\Psi }_{2,v}`$, and define $`W_{\mathrm{\Psi }_{1,v}}(s,g)`$ and $`W_{\mathrm{\Psi }_{2,v}}(s,g)`$ as in Lem. 10.5. Suppose moreover that $`_{xF\backslash 𝔸_F}e_F(x)\phi (n(x)g)`$ factorizes as $`_vW_v(g)`$. Then we can express the global integral of (11.9) as a product in two different ways, depending on whether we let $`E_{\mathrm{\Psi }_1}`$ or $`E_{\mathrm{\Psi }_2}`$ play the role of $`\pi _2`$. Namely, as in (11.7):
(11.10)
$$\begin{array}{c}_𝐗\phi (g)E_{\mathrm{\Psi }_1}(g,\frac{1}{2}+t)E_{\mathrm{\Psi }_2}(ga([𝔭]),\frac{1}{2}+t^{})𝑑g\hfill \\ \hfill =c_F\mathrm{N}(𝔭)^{1/2+t^{}}\underset{v}{}I_v(W_v,W_{\mathrm{\Psi }_{1,v}}(1/2+t,),\mathrm{\Psi }_{2,v}^{},1/2+t^{})\\ \hfill =c_F\underset{v}{}I_v(W_v,W_{\mathrm{\Psi }_{2,v}}(1/2+t^{},)^{a([𝔭])_v},\mathrm{\Psi }_{1,v},1/2+t)\end{array}$$
Here $`W_{\mathrm{\Psi }_{2,v}}(1/2+t^{},)^{a([𝔭])_v}`$ denotes the translate of $`W_{\mathrm{\Psi }_{2,v}}`$ by the $`v`$th component of $`a([𝔭])`$.
For $`v`$ nonarchimedean (notations being similar to that of the previous Lemma) we take $`\mathrm{\Psi }_{1,v}`$ and $`\mathrm{\Psi }_{2,v}`$ to be the characteristic function of $`𝔬_v^2`$ for every finite $`v`$, and $`W_v`$ to be the new vector.
For $`v`$ archimedean we first apply Lem. 11.2 with $`s_0=1/2+t_0`$, and $`\pi _{2,v}`$ the representation of $`\mathrm{PGL}_2(F_v)`$ spanned by $`E_{\mathrm{\Psi }_2}(1/2+t_0^{},g)`$, i.e. the representation unitarily induced from the character $`a(y)|y|_v^{it_0^{}}`$. Lem. 11.2 provides $`W_v`$ in the Whittaker model of $`\pi _v`$, $`W_{2,v}`$ in the Whittaker model of $`\pi _{2,v}`$, and a Schwarz function $`\mathrm{\Psi }_{1,v}`$ with $`|I_v(W_v,W_{2,v},\mathrm{\Psi }_{1,v},1/2+t_0)|1`$. The last comment of Rem. 10.2 shows that there is a standard $`\mathrm{\Psi }_{2,v}`$ so that $`W_{\mathrm{\Psi }_{2,v}}(1/2+t_0^{},g_v)=W_{2,v}(g_v)`$ (notation of Lem. 10.5). Moreover, Lemma 11.2 also shows that $`\mathrm{\Psi }_{1,v}`$ and $`W_{2,v}`$ (so also $`\mathrm{\Psi }_{2,v}`$) may be chosen from a fixed finite set of possibilities (depending, of course, on the original bounded set to which $`\pi _{\mathrm{}}`$ belongs, as well as $`t_0`$ and $`t_0^{}`$).
Again we put $`\mathrm{\Psi }_i=_v\mathrm{\Psi }_{i,v}`$ for $`i=1,2`$ and take $`\phi `$ with $`_{xF\backslash 𝔸_F}\phi (n(x)g)=_vW_v(g)`$. From (11.10) we deduce that, with our choices, $`|\mathrm{\Phi }(t_0,t_0^{})|1`$. The assertion about $`\phi _{L^2}`$ follows as in the proof of the previous Lemma. The second assertion of the Lemma (concerning invariance of $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$) is immediate.
It remains to prove that $`\mathrm{\Phi }`$ is actually holomorphic in $`(t,t^{})`$ and that $`|\mathrm{\Phi }(t,t^{})|(1+|t|+|t^{}|)^Ce^{C|\mathrm{}(t)|+C|\mathrm{}(t^{})|}`$. Put $`\mathrm{\Xi }_v=\frac{I_v}{L_v(\frac{1}{2}+t+t^{},\pi _v)L_v(\frac{1}{2}+tt^{},\pi _v)}`$. It is simple to explicitly compute $`\mathrm{\Xi }_v`$ for nonarchimedean $`v`$, using Cor. 10.1 and Lem. 11.1. One thereby sees that it will suffice to check, by similar arguments to those used in Lem. 11.2, the following statement for $`v`$ archimedean: $`\mathrm{\Xi }_v=c^sc^s^{}P(s,s^{})`$, where $`P`$ is a polynomial, and moreover $`c,c^{},P`$ vary continuously in $`\pi _v`$, if $`\pi _vW_v`$ is a continuous assignment. We only sketch the proof of this. From (11.7) and the fact that $`W_v,\mathrm{\Psi }_{1,v},\mathrm{\Psi }_{2,v}`$ are all $`K_v`$-finite, it suffices to prove the corresponding assertions for $`_{F_v^\times }W_v(a(y))W_{\mathrm{\Psi }_{1,v}}(s,a(y))|y|^{s^{}1}d^\times y`$. For this we use Barnes’ formula as in . ∎
### 11.3. Local Rankin-Selberg convolutions.
Let $`v`$ be a nonarchimedean place of $`F`$ with residue characteristic $`q_v`$. Let $`\pi _1,\pi _2`$ be generic irreducible admissible representations of $`\mathrm{GL}(2,F_v)`$ with trivial central character. (Since we shall work purely locally over $`F_v`$ throughout the present subsection, we shall use the notation $`\pi _1`$ rather than e.g. $`\pi _{1,v}`$).
Then $`\pi _1,\pi _2`$ are self-dual. We assume that $`\pi _2`$ is unramified and $`\pi _1`$ has conductor $`q_v`$, and denote by $`L(s,\pi _j)`$ the local $`L`$-factors.
Fix once and for all an additive unramified character $`\psi `$ of $`F_v`$. Let $`v:F_v^\times `$ be the valuation, put $`𝔬_{F_v}=\{xF_v^\times :v(x)0\}`$, and choose a uniformizer $`\varpi F_v^\times `$. Let $`𝔬_{F_v}^\times `$ be the multiplicative group of units in $`𝔬_{F_v}`$. Let $`d^\times x,dx`$ be Haar measures on $`F_v^\times ,F_v`$ respectively, assigning mass $`1`$ to $`𝔬_{F_v}^\times `$ and $`𝔬_{F_v}`$ respectively. For $`xF_v`$, put $`n(x)=\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right)`$. Also let $`w=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$ We choose a Whittaker model for $`\pi _1`$ transforming by the character $`n(x)\psi (x)`$, and a Whittaker model for $`\pi _2`$ transforming by the character $`n(x)\overline{\psi (x)}`$.
Let $`\mathrm{\Psi }_v`$ be the characteristic function of $`𝔬_{F_v}^2`$. Set $`W_1`$ to be the new vector in the Kirillov model of $`\pi _1`$, let $`W_2^{}`$ be the new vector in the Kirillov model of $`\pi _2`$, and set $`W_2=\pi _2(\left(\begin{array}{cc}1& 0\\ 0& \varpi \end{array}\right))W_2^{}`$. Then both $`W_1,W_2`$ are invariant by the subgroup
(11.11)
$$K_0=\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):a,b,d𝔬_{F_v},c\varpi 𝔬_{F_v}\}.$$
Moreover $`W_2`$ is invariant by $`n(\varpi ^1𝔬_{F_v})`$.
###### Lemma 11.6.
Notations being as in (11.6), let $`L(s,\pi _1\times \pi _2)`$ be the local $`L`$-factor. Then:
$$\frac{I_v(W_{1,v},W_{2,v},\mathrm{\Psi }_v)}{L(s,\pi _1\times \pi _2)}=\pm \frac{q_v^s}{q_v+1}$$
###### Proof.
We shall often use the following shorthand: if $`W`$ is a function in the Whittaker model of $`\pi \{\pi _1,\pi _2\}`$ and for $`zF_v^\times `$, we write $`W(z)`$ for $`W(\left(\begin{array}{cc}z& 0\\ 0& 1\end{array}\right))`$. Thus, for instance, $`W_2(z)=W_2^{}(z\varpi ^1)`$. The function $`zW(z)`$ belongs to the Kirillov model of $`\pi `$.
Let $`ϵ\{1,1\}`$ be the local root number of $`\pi _1`$ (it lies in $`\{1,1\}`$ since $`\pi _1`$ is self-dual). Then:
(11.12) $`{\displaystyle _{aF_v^\times }}W_1(a)|a|^{s1/2}d^\times a=L(s,\pi _1),`$
$`{\displaystyle _{aF_v^\times }}W_2(a)|a|^{s1/2}d^\times a=q_v^{(s1/2)}L(s,\pi _2),`$
as follows from defining properties of newforms and the fact $`W_2(z)=W_2^{}(z\varpi ^1)`$; moreover
(11.13) $`{\displaystyle _{aF_v^\times }}\pi _1(w)W_1(a)|a|^{s1/2}d^\times a=ϵq_v^{(s1/2)}L(s,\pi _1),`$
$`{\displaystyle _{aF_v^\times }}\pi _2(w)W_2(a)|a|^{s1/2}d^\times a=q_v^{(s1/2)}L(s,\pi _2),`$
as follows from local functional equation for the standard $`L`$-function on $`\mathrm{GL}(2)`$: see \[16, 2.18\]. <sup>23</sup><sup>23</sup>23 That is, $`_{aF^\times }W(a)|a|^{s1/2}d^\times a=\frac{L(s,\pi )}{ϵ(s,\pi )L(1s,\stackrel{~}{\pi })}_{F^\times }W(aw)|a|^{1/2s}\omega ^1(a)d^\times a`$. In particular, if $`\pi `$ is a representation with trivial central character, and $`\chi `$ a character of $`F^\times `$, $`_{aF^\times }W(a)\chi (a)d^\times a=\frac{L(\frac{1}{2},\pi \chi )}{ϵ(\frac{1}{2},\pi \chi )L(\frac{1}{2},\stackrel{~}{\pi }\overline{\chi })}_{F^\times }W(aw)\chi ^1(a)d^\times a.`$
Note moreover that $`W_1`$, $`W_2`$, $`\pi _1(w)W_1`$, and $`\pi _2(w)W_2`$ are all invariant by the maximal compact subgroup of the diagonal torus of $`\mathrm{GL}_2`$. Thus (11.12) and (11.13) completely determine their restriction to the diagonal torus; we now explicate this.
Choose $`\alpha `$ so that $`L(s,\pi _1)=(1\alpha q_v^s)^1`$. In fact, $`\alpha =ϵq_v^{1/2}`$, by \[16, Prop. 3.6\]. Choose $`\gamma _1,\gamma _2`$ so that $`L(s,\pi _2)=((1\gamma _1q_v^s)(1\gamma _2q_v^s))^1`$. Recalling the notational convention established in the paragraph prior to (11.12), we see:
(11.14) $`W_1(\varpi ^r)=\{\begin{array}{cc}\alpha ^rq_v^{r/2},r0\hfill & \\ 0,r<0\hfill & \end{array},\pi _1(w)W_1(\varpi ^r)=\{\begin{array}{cc}ϵ\alpha ^{r+1}q_v^{\frac{r+1}{2}},r1\hfill & \\ 0,r<1\hfill & \end{array}`$
$`W_2(\varpi ^r)=\{\begin{array}{cc}0,r0\hfill & \\ 1,r=1\hfill & \\ (\gamma _1^{r1}+\gamma _1^{r2}\gamma _2+\mathrm{}+\gamma _2^{r1})q_v^{\frac{r1}{2}},r2\hfill & \end{array}`$
$`\pi _2(w)W_2(\varpi ^r)=\{\begin{array}{cc}0,r<1\hfill & \\ 1,r=1\hfill & \\ (\gamma _1^{r+1}+\gamma _1^r\gamma _2+\mathrm{}+\gamma _2^{r+1})q_v^{\frac{r+1}{2}},r0\hfill & \end{array}`$
The local integral we wish to evaluate is the right hand side of (11.6). In the case at hand, with $`N,G,Z`$ denoting the $`F_v`$-points of the respective groups, we have:
(11.15)
$$I(s):=_{ZN\backslash G}W_1(g)W_2(g)|det(g)|^s\left(_t\mathrm{\Psi }_v((0,t)g)|t|^{2s}d^\times t\right)𝑑g$$
Using the Iwasawa decomposition, and recalling $`\mathrm{\Psi }_v`$ was the characteristic function of $`𝔬_v^2`$ one finds:
$$I(s)=(1q_v^{2s})^1_{A\times K_v}\pi _1(k)W_1(a)\pi _2(k)W_2(a)|a|^{s1}d^\times a𝑑k,$$
where the measure $`dk`$ is the Haar measure of total mass $`1`$, and $`d^\times a`$ assigns mass $`1`$ to $`AK_v`$.
The function $`k\pi _1(k)W_1(a)\pi _2(k)W_2(a)`$ is right invariant by $`K_0`$ (see (11.11) for definition) and left invariant by $`NK_v`$. There are two $`(NK_v,K_0)`$ double cosets in $`K_v`$, and we may therefore express $`I(s)`$ as a sum:
(11.16) $`(1q_v^{2s})I(s)={\displaystyle \frac{1}{q_v+1}}{\displaystyle _{F_v^\times }}W_1(a)W_2(a)|a|^{s1}d^\times a`$
$`+{\displaystyle \frac{q_v}{q_v+1}}{\displaystyle _{aF_v^\times }}\pi _1(w)W_1(a)\pi _2(w)W_2(a)|a|^{s1}d^\times a`$
To evaluate $`I(s)`$, we use (11.14). Noting that $`L(s,\pi _1\times \pi _2)=\frac{1}{(1\alpha \gamma _1q_v^s)(1\alpha \gamma _2q_v^s)}`$, an easy computation shows
$$I(s)=\frac{L(s,\pi _1\times \pi _2)}{(q_v+1)(1q_v^{2s})}\left(\alpha q_v^{1/2}q_v^{(s1)}+ϵq_vq_v^{s1}\right)$$
from where we obtain $`I(s)=ϵ\frac{q_v^s}{q_v+1}L(s,\pi _1\times \pi _2)`$. Note also that $`I(s)`$ satisfies the necessary functional equation. ∎
### 11.4. Hecke-Jacquet-Langlands integral representations for standard $`L`$-functions
Our goal here is to prove Prop. 6.1 and 6.2, used in the text. This amounts to explicit computations connected to Hecke-Jacquet-Langlands integral representations. Since, in the main text, we obtain subconvexity for $`\mathrm{GL}(1)`$ twists of $`\mathrm{GL}(2)`$ $`L`$-functions, with polynomial dependence in all parameters, we will have to be somewhat more precise than in the case of Rankin-Selberg $`L`$-functions.
Let $`\pi `$ be a cuspidal representation of $`\mathrm{GL}_2`$ over $`𝔸_F`$. Let $`\chi `$ be a unitary character of $`𝔸_F^\times /F^\times `$ of finite conductor $`𝔣`$. Put $`L_{unr}(s,\pi \times \chi )`$ to be the unramified part of the (finite) standard $`L`$-function:
$$L_{unr}(s,\pi \times \chi ):=\underset{v\mathrm{finite},\chi _v\mathrm{unramified}}{}L_v(s,\pi _v\times \chi _v).$$
Define $`\mu _z`$ as in (6.4), i.e. the measure on $`𝐗_{\mathrm{GL}(2)}`$ defined as
$$\mu _z(f)=_{|y|=z}f(a(y)n([𝔣]))\chi (y)d^\times y.$$
We refer to Sec. 2.3 and Sec. 2.5 for notation, as well as the start of Section 6 for a discussion of the meaning of $`\mu _z`$ in classical terms.
###### Lemma 11.7.
Let $`v`$ be a nonarchimedean place of $`F`$ with residue characteristic $`q_v`$, and $`\pi _v`$ an irreducible generic representation of $`\mathrm{GL}_2(F_v)`$. Let $`\psi _v`$ be an unramified additive character of $`F_v`$. Let $`\chi _v:F_v^\times `$ a multiplicative character of conductor $`r`$, $`W_v`$ be the new vector in the $`\psi _v`$-Whittaker model of $`\pi _v`$. Then
$$_{yF_v^\times }W_v(a(y)n(\varpi _v^r))\chi _v(y)|y|^{s1/2}d^\times y=\{\begin{array}{cc}L_v(s,\pi _v\times \chi _v),r=0.\hfill & \\ \theta ,r1,\hfill & \end{array}$$
where $`\theta `$ is a scalar of absolute value $`q_v^{r/2}(1q_v^1)^1`$.
###### Proof.
If $`r=0`$, then $`\chi _v`$ is unramified, the result follows immediately from the definition of the new vector. Otherwise, $`\chi _v`$ is ramified, and we rewrite the integral under consideration as
(11.17)
$$_{yF_v^\times }W_v(a(y))\psi _v(\varpi _v^ry)\chi _v(y)|y|^{s1/2}d^\times y.$$
Now $`W_v(a(y))`$ vanishes when $`v(y)<0`$ and it is $`𝔬_{F_v}^\times `$-invariant. The integral $`_{v(y)=k}\chi _v(y)\psi _v(\varpi _v^ry)d^\times y`$ is nonvanishing only when $`k=0`$. In that case, it is a Gauss sum with absolute value $`\frac{q_v^{r/2}}{(1q_v^1)}`$, where the factor $`(1q_v^1)^1`$ arises from the measure normalization (cf. Sec. 2.6) namely $`_{v(y)=0}d^\times y=1`$. The result follows. ∎
###### Lemma 11.8.
Let $`d,\beta 0`$. Then there exists $`\phi \pi `$ such that, with
(11.18)
$$\mathrm{\Phi }(s)=\mathrm{N}(𝔣)^{1/2}\frac{_z\mu _z(\phi )|z|^{s1/2}d^\times z}{L_{unr}(s,\pi \times \chi )}$$
then $`\mathrm{\Phi }(s)`$ is holomorphic and satisfies:
1. $`|\mathrm{\Phi }(s)|_{\mathrm{}(s),ϵ}\mathrm{N}(𝔣)^ϵ`$ and $`|\mathrm{\Phi }(\frac{1}{2})|_ϵ\mathrm{N}(𝔣)^ϵ`$.
2. $`\phi `$ is new at every finite place (i.e., for each finite prime $`𝔮`$ it is invariant by $`K_0[𝔮^{s_𝔮}]`$, where $`s_𝔮`$ is the local conductor of the local constituent $`\pi _𝔮`$).
3. The Sobolev norms of $`\phi `$ satisfy the bounds (conductor notation as in Sec. 2.12.2)
(11.19)
$$S_{2,d,\beta }(\phi )_ϵ\mathrm{Cond}_{\mathrm{}}(\pi )^{2d+ϵ}\mathrm{Cond}_f(\pi )^{\beta +ϵ}\mathrm{Cond}_{\mathrm{}}(\chi )^{1/2+2d}$$
###### Proof.
For each infinite place $`w`$ of $`F`$, denote by $`\mathrm{Cond}_w(\chi )`$ the contribution from $`w`$ to the Iwaniec-Sarnak analytic conductor of $`\chi `$ (see Sec. 2.12.2.)
The map $`\phi W_\phi =_{F\backslash 𝔸_F}e_F(x)\phi (n(x)g)`$ is an isomorphism between the space of $`\pi `$ and the Whittaker model of $`\pi `$. For each finite $`v`$, take $`W_v`$ to be a new vector in the Whittaker model of $`\pi _v`$. A point of caution is that $`e_F`$ may not be unramified on $`F_v`$; to be absolutely concrete, we set $`W_v(g)=W_{v,\mathrm{new}}(a(\varpi ^{d_v})g)`$, where $`W_{v,\mathrm{new}}`$ is the new vector in the Whittaker model of $`\pi _v`$ taken w.r.t an unramified additive character of $`F_v`$, and $`d_v=v(𝔡)`$ is the local valuation of the different.
Let us now choose $`W_v`$ at the infinite places. Let $`g_1`$ be a smooth positive function of compact support on $`F_v`$. Let $`\mathrm{deg}(v)=2`$ if $`v`$ is complex and $`\mathrm{deg}(v)=1`$ if $`v`$ is real. For $`\mathrm{}|v`$, define
$$W_v(y)=\mathrm{Cond}_v(\chi )g_1(\mathrm{Cond}_v(\chi )^{1/\mathrm{deg}(v)}(y1)).$$
This is possible by the theory of the Kirillov model; thus $`W_v`$ is a smooth (but not $`K_v`$-finite) vector. In words, if $`v`$ is real, the function $`W_v`$ is supported in a neighbourhood of the identity of size $`\mathrm{Cond}_v(\chi )^1`$ and takes values of size $`|\mathrm{Cond}_v(\chi )|`$ there; if $`v`$ is complex, a similar statement holds but now $`W_v`$ is supported in a disc around the identity with area $`\mathrm{Cond}_v(\chi )^1`$.
Then there exists $`\phi \pi `$ with $`W_\phi =_vW_v`$. By unfolding, it follows that for $`\mathrm{}(s)1`$
(11.20)
$$_z\mu _z(\phi )|z|^{s1/2}d^\times z=c_F\underset{v}{}_{yF_v^\times }W_v(a(y)n([𝔣]))|y|^{s1/2}\chi _v(y)d^\times y$$
Here $`c_F`$ is a constant depending only on $`F`$, arising from change of measure; it is entirely unimportant as we will be only interested in bounds.<sup>24</sup><sup>24</sup>24 The measure $`\mu _z`$ is normalized as a probability measure, whereas to unfold from $`𝔸_F^\times `$ to $`_vF_v^\times `$ we use the measures previously set up there (see Section 2.6).
By Lem. 11.7, with $`\mathrm{\Phi }(s)`$ as in the statement of the Lemma,
(11.21)
$$\mathrm{\Phi }(s)=c_F\theta ^{}\mathrm{N}(𝔡)^{s1/2}\underset{\mathrm{infinite}v}{}_{F_v^\times }W_v(a(y))|y|^{s1/2}\chi _v(y)$$
where $`|\theta ^{}|=_{𝔮|𝔣}(1\mathrm{N}(𝔮)^1)^1`$. For this choice of $`\phi `$, the second assertion if the Lemma is clear, and, if we choose the support of $`g_1`$ to be small enough, the first assertion also follows easily. <sup>25</sup><sup>25</sup>25For the assertion concerning the lower bound for $`|\mathrm{\Phi }(1/2)|`$, the point, in words, is that our choices are so that $`\chi _v`$ does not oscillate over the support of $`W_v`$, cf. Remark 2.2. Note how convenient it is, here and elsewhere, to use smooth vectors rather than $`K_{\mathrm{}}`$-finite vectors; one could not e.g. achieve $`W_v`$ of compact support with $`K_{\mathrm{}}`$-finite vectors.
(11.19) follows from Lem. 8.4, together with Lem. 11.3 and the upper bound for $`L`$-functions near $`1`$ due to Iwaniec. See \[15, Chapter 8\] for this bound. ∎
The previous Lemma shows that $`L(1/2,\pi \times \chi )`$ may be “well-approximated” by an appropriate period integral. Unfortunately, this period integral is against a measure of infinite mass, since $`𝔸_F^\times /F^\times `$ is of infinite volume. It is, therefore, convenient to know that the $`\mu _z`$-integral of (11.18) can be truncated to a compact range without affecting the answer too much. This is, roughly speaking, the geometric equivalent of the approximate functional equation in the classical theory, and is provided by the next Lemma. It says, roughly speaking, that the integral of (11.18) can be truncated to the range where $`z`$ is around $`\mathrm{N}(𝔣)^1`$.
###### Lemma 11.9.
Let notation be as in Lem. 11.8. Let $`g_+,g_{}`$ be positive smooth functions on $`_0`$ such that $`g_++g_{}=1`$, $`g_+(t)=1`$ for $`t2`$ and $`g_{}(t)=1`$ for all $`t1/2`$. Then
$$I_+:=_z\mu _z(\phi )g_+(z/T)d^\times z_{g_+,ϵ}(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\pi )\mathrm{Cond}(\chi ))^ϵ$$
$$I_{}:=_z\mu _z(\phi )g_{}(z/T)d^\times z_{g_{},ϵ}(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵ(\mathrm{Cond}_{\mathrm{}}(\chi )\mathrm{Cond}(\pi ))^{1+ϵ}$$
###### Proof.
Recall the definition of $`\mu _z`$ from (6.4). Put $`\widehat{g_\pm }(s)=g_\pm (x)x^{s1}𝑑x`$, the Mellin transform of $`g_\pm `$; then $`g_\pm `$ is holomorphic in $`\pm \mathrm{}(s)<0`$ and for any $`M0,\pm \sigma <0`$ the integral $`_{\mathrm{}(s)=\sigma }|\widehat{g_\pm }(s)|(1+|s|)^M𝑑s`$ is convergent. Then, for any $`\pm \sigma >0`$, we have, by the Plancherel formula on $`^\times `$, that:
$$_z\mu _z(\phi )g_\pm (z/T)d^\times z=\frac{1}{2\pi i}_{\mathrm{}(s)=\sigma }\left(\mu _z(\phi )|z|^sd^\times z\right)T^s\widehat{g_\pm }(s)𝑑s.$$
So for any $`M>0`$,
$$|I_\pm |_{\sigma ,g_\pm ,M}T^\sigma \mathrm{N}(𝔣)^{1/2}\underset{\mathrm{}(s)=1/2+\sigma }{sup}\frac{|L_{unr}(s)\mathrm{\Phi }(s)|}{(1+|s|)^M}$$
where $`\mathrm{\Phi }`$ is as in the previous Lemma.Take $`\sigma =1/2+\epsilon `$ in the $`+`$ case, $`1/2\epsilon `$ in the $``$ case.
Using Iwaniec’s bounds on $`L`$-functions near $`1`$ \[15, Chapter 8\] and the functional equation, we see that for sufficiently large $`M`$:
(11.22)
$$\begin{array}{c}\underset{\mathrm{}(s)=1+\epsilon }{sup}|L_{unr}(s,\pi \times \chi )|\mathrm{Cond}(\pi \chi )^\epsilon \hfill \\ \hfill \frac{sup_{\mathrm{}(s)=\epsilon }|L_{unr}(s,\pi \times \chi )|}{(1+|s|)^M}_M\mathrm{Cond}(\pi \chi )^{1/2+\epsilon }\underset{\chi _v\mathrm{ramified}\mathrm{finite}}{}\underset{\mathrm{}(s)=\epsilon }{sup}|L_v(s,\pi _v\times \chi _v)|^1\end{array}$$
For each $`v`$ where $`\chi _v`$ is ramified and $`L_v(s,\pi _v\times \chi _v)`$ is not identically $`1`$, the representation $`\pi _v`$ must also be ramified (i.e., not spherical). So one can bound the product on the second line on (11.22), using trivial bounds towards the Ramanujan conjecture, by $`\mathrm{Cond}(\pi )^{1/2+2\epsilon }`$. The fact that $`\mathrm{Cond}(\chi )=\mathrm{Cond}_{\mathrm{}}(\chi )\mathrm{N}(𝔣)`$, the bound $`\mathrm{Cond}(\pi \chi )\mathrm{Cond}(\pi )\mathrm{Cond}(\chi )^2`$, and the (easily verified) analogue of this bound of at archimedean places, allows one to conclude. ∎
We now address the analogue of the previous Lemmas when $`\pi `$ is noncuspidal. <sup>26</sup><sup>26</sup>26The content of the following Lemma, in classical language, is related to the following observation. Let $`\chi `$ be an even Dirichlet character mod $`q`$, and let $`E^{}(s,z)`$ to be the Eisenstein series of (10.6), and $`\overline{E}^{}(s,z):=E^{}(s,z)\xi (2s)y^s\xi (2(1s))y^{1s}`$, then $`\frac{1}{q}_0^{\mathrm{}}_{1xq1}\overline{E}^{}(s,\frac{x}{q}+iy)y^s^{}d^\times y`$ coincides, up to some harmless factor, with $`q^{1/2}\mathrm{\Lambda }(\chi ,s+s^{})\mathrm{\Lambda }(\chi ,1s+s^{}))`$, where $`\mathrm{\Lambda }(\chi ,s)`$ is the usual Dirichlet $`L`$-function completed to include the $`\mathrm{\Gamma }`$-factor at $`\mathrm{}`$. This particular expression is actually not quite suitable for our needs, because of the rapid decay of the $`\mathrm{\Gamma }`$-factor swamps information about the finite $`L`$-function, and in fact the Lemma uses (the equivalent of) a different test vector belonging to the automorphic representation underlying $`E^{}(z,s)`$.
###### Lemma 11.10.
Let $`s_0,s_0^{}`$. There is an absolute $`C>0`$ (i.e., depending only on $`F`$) and a Schwarz function $`\mathrm{\Psi }`$ (depending on $`\chi `$) so that if we put
$$\mathrm{\Phi }(s,s^{}):=\mathrm{N}(𝔣)^{1/2}\frac{_{y𝔸_F^\times /F^\times }\overline{E}_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)|y|^s^{}d^\times y}{L(\chi ,s+s^{})L(\chi ,1s+s^{})}$$
where $`\overline{E}`$ is defined as in (10.9), then the integral defining $`\mathrm{\Phi }`$ is absolutely convergent in a right half-plane $`\mathrm{}(s)1`$. Moreover, $`\mathrm{\Phi }`$ extends from $`\mathrm{}(s),\mathrm{}(s^{})1`$ to a holomorphic function on $`^2`$, satisfying
1. $`|\mathrm{\Phi }(1/2,0)|1`$ and $`|\mathrm{\Phi }(s,s^{})|C^{1+|\mathrm{}(s)|+|\mathrm{}(s^{})|}(1+|s|+|s^{}|)^C`$.
Moreover, given $`N>0`$ we have that
(11.23)
$$|\mathrm{\Phi }(s,s^{})|(1+|s|+|s^{}|)^N_{\mathrm{}(s),\mathrm{}(s^{}),N}\mathrm{Cond}_{\mathrm{}}(\chi )^N^{}$$
where $`N^{}`$ and the implicit constant may be taken to depend continuously on $`N,\mathrm{}(s),\mathrm{}(s^{})`$.
2. $`\mathrm{\Psi }`$, and so also $`E_\mathrm{\Psi }(s,g)`$ is invariant by $`K_{\mathrm{max}}`$;
3. Let $`h(\kappa )`$ be as in (10.18), and put $`E_h:=_{\mathrm{}(s)1}h(s)E_\mathrm{\Psi }(s,g)𝑑g`$. For each $`d,\beta `$ there is $`N`$ such that $`S_{\mathrm{},d,\beta }(E_h)_\kappa h_0\mathrm{Cond}_{\mathrm{}}(\chi )^N`$.
###### Proof.
We shall not explicitly address details of convergence. The manipulations that follow may be justified by similar reasoning to that of Lem. 10.6.
We now define a Schwarz function $`\mathrm{\Psi }_v`$ on $`F_v^2`$ for each place $`v`$. For each finite place $`v`$, let $`\mathrm{\Psi }_v`$ be the characteristic function of $`𝔬_v^2`$.
For infinite $`v`$, we will first define a Schwarz function $`\rho _v`$ on $`F_v`$, and then take $`\mathrm{\Psi }_v(x,y)=\rho _v(x)\widehat{\rho _v}(y)`$; here $`\widehat{\rho _v}`$ is the inverse Fourier transform of $`\rho _v`$, satisfying $`_{F_v}\widehat{\rho _v}(y)e_{F_v}(xy)𝑑y=\rho _v(x)`$.
Let $`g_1`$ be a smooth positive function of compact support on $`F_v`$. Let $`\mathrm{deg}(v)=2`$ if $`v`$ is complex and $`\mathrm{deg}(v)=1`$ if $`v`$ is real. For $`\mathrm{}|v`$, define
$$\rho _v(y)=\mathrm{Cond}_v(\chi )g_1(\mathrm{Cond}_v(\chi )^{1/\mathrm{deg}(v)}(y1)).$$
In words: in the real (resp. complex) case, $`\rho _v`$ is localized in a real (resp. complex) interval (resp. disc) around $`1`$, of length (resp. area) $`\mathrm{Cond}_v(\chi )^1`$. Now put $`\mathrm{\Psi }_v(x,y)=\rho _v(x)\widehat{\rho _v}(y)`$. The function $`\mathrm{\Psi }_v`$ is not compactly supported; however, it is of rapid decay. Indeed for each Schwarz norm $`𝒮`$, there is $`M>0`$ such that
(11.24)
$$𝒮(\mathrm{\Psi }_v)\mathrm{Cond}_v(\chi )^M$$
Define a Schwarz function on $`𝔸_F^2`$ via $`\mathrm{\Psi }(x,y)=_v\mathrm{\Psi }_v(x,y)`$. Define $`W_\mathrm{\Psi }(s,g)`$ as in Lem. 10.5 to be the Fourier coefficient of $`E_\mathrm{\Psi }(s,g)`$. The choice of $`\mathrm{\Psi }`$ and Lem. 10.5 shows that $`W_\mathrm{\Psi }(s,g)=_vW_v(g)`$, where, for each finite $`v`$, $`W_v`$ is given by Cor. 10.1, and satisfies
(11.25)
$$_{F_v^\times }W_v(a(y))|y|^s^{}d^\times y=q_v^{d_v(1+s^{}s)}L_v(||^s,s^{})L_v(||^{1s},s^{}),$$
For infinite $`v`$, $`W_v`$ satisfies (Lem. 10.5)
(11.26)
$$_{F_v^\times }W_v(a(y))\omega (y)|y|^s^{}d^\times y=_{F_v^\times }\rho _v(x)\omega (x)|x|^{s+s^{}}d^\times x_{F_v^\times }\rho _v(x)\omega (x)|x|^{1s+s^{}}d^\times x$$
By Fourier analysis and Lem. 10.3, $`\overline{E}_\mathrm{\Psi }(s,g)=_{\alpha F^\times }W_\mathrm{\Psi }(a(\alpha )g)`$. Thus, for $`\mathrm{}(s)1`$,
(11.27)
$$\begin{array}{c}_{y𝔸_F^\times /F^\times }\overline{E}_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)|y|^s^{}d^\times y=_{y𝔸_F^\times }W_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)|y|^s^{}d^\times y\hfill \\ \hfill =\underset{v}{}_{yF_v^\times }W_v(a(y)n([𝔣]))\chi (y)|y|^s^{}d^\times y\end{array}$$
For $`s1`$, we use (10.17), (11.26) and Lemma 11.7 to evaluate the local factors, obtaining:
(11.28)
$$\begin{array}{c}\mathrm{\Phi }(s,s^{})=\theta ^{}\mathrm{N}(𝔡)^{1+s^{}s}\underset{v\mathrm{infinite}}{}_{yF_v^\times }W_v(a(y))\chi (y)|y|^s^{}d^\times y\hfill \\ \hfill =\theta ^{}\mathrm{N}(𝔡)^{1+s^{}s}_{F_v^\times }\rho _v(x)\chi (x)|x|^{s+s^{}}d^\times x_{F_v^\times }\rho _v(x)\chi (x)|x|^{1s+s^{}}d^\times x\end{array}$$
where $`|\theta ^{}|=_{𝔮|𝔣}(1\mathrm{N}(𝔮)^1)^1`$. Now, by choice of $`\phi _v`$, the integral $`I_v(s):=_{yF_v^\times }\rho _v(y)\chi (y)|y|^sd^\times y`$ satisfies $`|I_v(1/2)|1`$ and $`|I_v(s)|(1+|s|)^CC^{1+|\mathrm{}(s)|}`$, at least when we choose the support of $`g_1`$ to be sufficiently small. It also satisfies $`|I_v(s)|(1+|s|)^N_{N,\mathrm{}(s)}\mathrm{Cond}_{\mathrm{}}(\chi )^N^{}`$, where $`N^{}`$ and the implicit constant may be taken to depend continuously on $`N,\mathrm{}(s)`$.
The corresponding facts (i.e. the first assertion of the Lemma) about $`\mathrm{\Phi }`$ follow immediately. The second assertion of the Lemma is immediate from our choice of $`\mathrm{\Psi }`$.
As for the third and final assertion, it follows from Rem. 10.3 and (11.24).
###### Lemma 11.11.
Let notations be as in the previous Lemma. Assume $`\chi `$ is ramified at at least one finite place. Let $`g_+,g_{}`$ be positive smooth functions on $`_0`$ such that $`g_++g_{}=1`$, $`g_+(t)=1`$ for $`t2`$ and $`g_{}(t)=1`$ for all $`t1/2`$.
Then
(11.29)
$$\mu _z(E_h)_{K,\mathrm{\Psi },h}\mathrm{min}(z^K,z^K)$$
for any $`K1`$.
Moreover, there is an absolute $`N>0`$ such that
$$I_+:=_z\mu _z(E_h)g_+(z/T)d^\times z(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵh_N$$
$$I_{}:=_z\mu _z(E_h)g_{}(z/T)d^\times z(\mathrm{N}(𝔣)T)^{1/2}(T\mathrm{Cond}(\chi ))^ϵ\mathrm{Cond}_{\mathrm{}}(\chi )^{1+ϵ}h_N$$
(Here the norms $`_N`$ are as in (10.18).)
###### Proof.
Again, we shall leave verification of convergence to the reader. Recall that, with the relevant measure on $`𝔸_F^\times /F^\times `$ having mass $`1`$:
(11.30)
$$\begin{array}{c}\mu _z(E_h)=_{y𝔸_F^\times /F^\times ,|y|=z}E_h(a(y)n[𝔣])\chi (y)d^\times y\hfill \\ \hfill =_{y𝔸_F^\times /F^\times ,|y|=z}_{\mathrm{}(s)1}h(s)E_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)d^\times y\\ \hfill =_{y𝔸_F^\times /F^\times ,|y|=z}_{\mathrm{}(s)1}h(s)\overline{E}_\mathrm{\Psi }(s,a(y)n([𝔣]))\chi (y)d^\times y\end{array}$$
Here, the last equality is justified by the fact that $`(E_\mathrm{\Psi }\overline{E}_\mathrm{\Psi })(s,a(y)n([𝔣]))`$ is invariant under $`yyy^{}`$, for $`y^{}_v𝔬_v^\times `$. On the other hand, $`\chi `$ is nontrivial on $`_v𝔬_v^\times `$, by assumption.
Combining (11.30) with Lem. 11.10, we have
(11.31)
$$_z\mu _z(E_h)|z|^s^{}d^\times z=c_F\mathrm{N}(𝔣)^{1/2}_{\mathrm{}(s)1}h(s)L(\chi ,1s+s^{})L(\chi ,s+s^{})\mathrm{\Phi }(s,s^{})𝑑s.$$
Here $`c_F`$ is an (unimportant) constant arising from measure normalization, as in (11.20).
The assertion (11.29) follows immediately from this, inverse Mellin transform, and analytic properties of the right-hand side.
Now proceed as in Lem. 11.9; it follows that (for any $`M`$)
$$|I_\pm |T^{(1/2+\epsilon )}\mathrm{N}(𝔣)^{1/2}\underset{\mathrm{}(s^{})=\pm (1/2+\epsilon )}{sup}(1+|s^{}|)^Mh(s)L(\chi ,1s+s^{})L(\chi ,s+s^{})\mathrm{\Phi }(s,s^{})𝑑s.$$
We deal with the case of $`I_{}`$. In that case, we take the inner integral to be over $`\mathrm{}(s)=1/2`$, and put $`s^{}=1/2\epsilon it^{}`$, and it will suffice to bound $`h(1/2+it)L(\chi ,\epsilon itit^{})L(\chi ,\epsilon +itit^{})(1+|t|+|t^{}|)^C`$. This is bounded, up to an implicit constant depending on $`\epsilon `$, by $`\mathrm{Cond}(\chi )^{1+2\epsilon }h_M^{}(1+|t^{}|)^C^{}`$ for sufficiently big $`M^{},C^{}`$, whence the result. ∎
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# Rates of Neutrino Absorption on Nucleons and the Reverse Processes in Strong Magnetic Fields
## I Introduction
The processes
$`\nu _e+n`$ $``$ $`e^{}+p,`$ (1a)
$`\overline{\nu }_e+p`$ $``$ $`e^++n,`$ (1b)
play important roles in supernovae. A supernova is initiated by the collapse of a stellar core, which leads to the formation of a protoneutron star. Nearly all the gravitational binding energy of the protoneutron star is emitted in $`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$, $`\overline{\nu }_\mu `$, $`\nu _\tau `$, and $`\overline{\nu }_\tau `$, some of which would interact to heat the material above the protoneutron star. The forward neutrino absorption processes in Eq. (1) provide the dominant heating mechanism, which is counteracted by cooling of the material through the reverse neutrino emission processes. In a prevalent paradigm Bethe and Wilson (1985), supernova explosion is determined by the competition between these heating and cooling processes. These processes also interconvert neutrons and protons, thereby setting the neutron-to-proton ratio of the material above the protoneutron star Qian and Woosley (1996). This ratio is a key parameter that governs the production of heavy elements during the ejection of this material Woosley and Hoffman (1992); Hoffman et al. (1996). Thus, accurate rates of the processes in Eq. (1) are important for understanding supernova dynamics and nucleosynthesis.
Observations and theoretical considerations indicate that protoneutron stars with magnetic fields of $`10^{16}`$ G may be formed. The rates of the above processes in such strong fields have been studied in the literature with various approximations Leinson and Pérez (1998); Roulet (1998); Lai and Qian (1998); Arras and Lai (1999); Gvozdev and Ognev (1999); Chandra et al. (2002); Bhattacharya and Pal (2004); Shinkevich and Studenikin (2004). In our previous work Duan and Qian (2004), we used the Landau wave functions of $`e^\pm `$ and derived a set of simple and consistent formulas to calculate the rates of the processes in Eq. (1) in the presence of strong magnetic fields. We also applied these rates to discuss the implications of such fields for supernova dynamics. However, all the calculations in the literature, including our previous work, were mostly carried out to $`𝒪(1)`$, the zeroth order in $`1/m_N`$ with $`m_N`$ being the nucleon mass. None of them included both the effects of nucleon recoil and weak magnetism, which are of $`𝒪(1/m_N)`$ and known to be important for the conditions in supernovae Horowitz (2002). For modeling the production of heavy elements during the ejection of the material from the protoneutron star, an accuracy of $`1`$% for the rates of the processes in Eq. (1) is required to determine precisely the neutron-to-proton ratio in the material Hoffman et al. (1996). To achieve such an accuracy, the $`𝒪(1/m_N)`$ effects on these rates must be taken into account. In this paper we recalculate these rates using the respective Landau wave functions of $`e^\pm `$ and protons and including the $`𝒪(1/m_N)`$ corrections from both nucleon recoil and weak magnetism. Our goal is to identify the important factors in computing accurate rates of the processes in Eq. (1) for application to supernova nucleosynthesis in the presence of strong magnetic fields.
This paper is organized as follows. The energies and the wave functions of the relevant particles in magnetic fields are discussed in Sec. II. The cross sections of the neutrino absorption processes in Eq. (1) and the differential reaction rates of the reverse neutrino emission processes are derived to $`𝒪(1/m_N)`$ in Sec. III. The rates of these processes in supernova environments with magnetic fields of $`10^{16}`$ G are calculated and discussed in Sec. IV. Conclusions are given in Sec. V.
## II Particle energies and wave functions in magnetic fields
The importance of magnetic field effects can be gauged from the energy scale
$$\sqrt{eB}=7.69\left(\frac{B}{10^{16}\text{G}}\right)^{1/2}\text{MeV},$$
(2)
where $`e`$ is the charge of $`e^+`$ and $`B`$ is the field strength. However, there is no detailed knowledge of magnetic fields in supernovae. Observations indicate that neutron stars may have $`B10^{15}\text{ G}`$ long after their birth in supernovae Kouveliotou et al. (1999); Gotthelf et al. (1999); Ibrahim et al. (2003). This suggests that at least $`B10^{15}\text{ G}`$ can be generated during the formation of some protoneutron stars. A recent theoretical model suggests that $`B10^{16}\text{ G}`$ may be produced near the surface of a protoneutron star Akiyama et al. (2003). An upper limit of $`B10^{18}\text{ G}`$ can be estimated for such a star by equating the magnetic energy to its gravitational binding energy Lai (2001). To explore the effects of strong magnetic fields on the rates of the processes in Eq. (1), we consider that $`B10^{16}\text{ G}`$ may exist in the region of interest to supernova nucleosynthesis, which lies well below $`10^7`$ cm from the protoneutron star Qian and Woosley (1996). For such fields, the associated energy scale is much smaller than the mass of the $`W`$ boson $`M_W=80`$ GeV. So there will be no change in the description of the weak interaction that is involved in the processes in Eq. (1). On the other hand, the energy scale for $`B10^{16}\text{ G}`$ is larger than the temperature ($`T1`$ MeV) of the material above the protoneutron star and comparable to the typical neutrino energy ($`E_\nu 10`$ MeV). Thus, magnetic field effects on energy levels of charged particles ($`e^\pm `$ and $`p`$) will be important. Furthermore, $`B10^{16}\text{ G}`$ will induce polarization of nucleon spin at the level of $`eB/m_NT10^2`$. This is significant due to parity violation of weak interaction and should also be taken into account.
We discuss the energy levels and the corresponding wave functions of all the relevant particles in this section. We assume a uniform magnetic field $`𝐁`$ in the positive z-direction, for which the vector potential is
$$𝐀=(\frac{1}{2}By,\frac{1}{2}Bx,0).$$
(3)
All the wave functions will be given in Dirac-Pauli representation.
### II.1 Electron and positron
The motion of $`e^\pm `$ along the $`z`$-axis is not affected by the magnetic field, but the motion in the $`xy`$-plane is quantized into Landau levels with energies (see, e.g., Ref. Landau and Lifshitz (1977))
$$E_e=\sqrt{m_e^2+k_{ez}^2+2n_eeB},$$
(4)
where $`m_e`$ is the rest mass of $`e^\pm `$, $`k_{ez}`$ is the $`z`$-component of the momentum, and $`n_e`$ is an integer quantum number (i.e., $`n_e=0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{}`$). For the $`e^\pm `$ in the initial states of the neutrino emission processes in Eq. (1), the relevant $`E_e`$ is of the order of the temperature $`T1`$ MeV for the material above the protoneutron star. It can be seen from Eqs. (2) and (4) that these $`e^\pm `$ predominantly occupy the ground Landau level ($`n_e=0`$) for $`B10^{16}\text{ G}`$. In comparison, the $`e^\pm `$ in the final states of the neutrino absorption processes typically have $`E_e`$ of the order of the neutrino energy $`E_\nu 10`$ MeV. These $`e^\pm `$ can occupy excited Landau levels ($`n_e1`$).
The wave function of $`e^{}`$ in cylindrical coordinates $`(\xi ,\varphi ,z)`$ is
$$(\psi _e^{})_{s_e}=\frac{e^{i(k_{ez}zE_et)}e^{i(n_er_e)\varphi }}{\sqrt{2\pi L/eB}}(U_e^{})_{s_e},$$
(5)
where $`s_e=1`$ and $`1`$ for spin up and down, respectively, $`r_e`$ is the quantum number labeling the center of gyromotion in the $`xy`$-plane, and $`L`$ is the linear size of the normalization volume. In Eq. (5), the spinor $`(U_e^{})_{s_e}`$ is
$`(U_e^{})_{s_e=1}`$ $`={\displaystyle \frac{1}{\sqrt{2E_e(E_e+m_e)}}}\left(\begin{array}{c}(m_e+E_e)e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ 0\\ k_{ez}e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ i\sqrt{2n_eeB}I_{n_e,r_e}(eB\xi ^2/2)\end{array}\right)`$ (10)
and
$`(U_e^{})_{s_e=1}`$ $`={\displaystyle \frac{1}{\sqrt{2E_e(E_e+m_e)}}}\left(\begin{array}{c}0\\ (m_e+E_e)I_{n_e,r_e}(eB\xi ^2/2)\\ i\sqrt{2n_eeB}e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ k_{ez}I_{n_e,r_e}(eB\xi ^2/2)\end{array}\right).`$ (15)
The special function $`I_{n,r}(\zeta )`$ in the above equations is defined in Ref. Sokolov and Temov (1968) and can be calculated using the method given in Appendix A.
The wave function of $`e^+`$ is
$$(\psi _{e^+})_{s_e}=\frac{e^{i(k_{ez}zE_et)}e^{i(n_er_e)\varphi }}{\sqrt{2\pi L/eB}}(U_{e^+})_{s_e},$$
(16)
where
$`(U_{e^+})_{s_e=1}`$ $`={\displaystyle \frac{1}{\sqrt{2E_e(E_e+m_e)}}}\left(\begin{array}{c}i\sqrt{2n_eeB}e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ k_{ez}I_{n_e,r_e}(eB\xi ^2/2)\\ 0\\ (m_e+E_e)I_{n_e,r_e}(eB\xi ^2/2)\end{array}\right)`$ (21)
and
$`(U_{e^+})_{s_e=1}`$ $`={\displaystyle \frac{1}{\sqrt{2E_e(E_e+m_e)}}}\left(\begin{array}{c}k_{ez}e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ i\sqrt{2n_eeB}I_{n_e,r_e}(eB\xi ^2/2)\\ (m_e+E_e)e^{i\varphi }I_{n_e1,r_e}(eB\xi ^2/2)\\ 0\end{array}\right).`$ (26)
Clearly, $`E_e`$ does not depend on the quantum number $`r_e`$ in the wave functions. This leads to a degeneracy factor
$$\underset{r_e}{}1=\frac{eBL^2}{2\pi }$$
(27)
for each Landau level of $`e^\pm `$ (see, e.g., Ref. Lai (2001)). Each level is further degenerate with respect to spin except for the ground level \[$`(\psi _e^{})_{s_e=1}=(\psi _{e^+})_{s_e=1}=0`$ for $`n_e=0`$\]. This introduces an additional spin degeneracy factor $`g_{n_e}`$, which is 1 for $`n_e=0`$ and 2 for $`n_e1`$.
### II.2 Proton
Protons are nonrelativistic in the supernova environment of interest. For nonrelativistic $`e^+`$ with the same charge and spin as protons, expansion of Eq. (4) to $`𝒪(1/m_e)`$ gives
$$E_{e,\text{NR}}=m_e+\frac{k_{ez}^2}{2m_e}+\frac{n_eeB}{m_e}.$$
(28)
The above equation already accounts for the contribution from the $`e^+`$ magnetic moment of $`e/2m_e`$. Unlike $`e^+`$, protons have an anomalous magnetic moment of $`\stackrel{~}{\mu }_p=1.79\mu _N`$ in addition to the value $`\mu _N=e/2m_p`$ expected for a spin-1/2 point particle of charge $`e`$ and mass $`m_p`$. Taking this into consideration, we obtain the energies of the proton Landau levels as
$$E_p=m_p+\frac{k_{pz}^2}{2m_p}+\frac{n_peB}{m_p}s_p\stackrel{~}{\mu }_pB,$$
(29)
where symbols have similar meanings to those for $`e^\pm `$. For $`B10^{16}\text{ G}`$, $`eB/m_p\stackrel{~}{\mu }_pB60`$ keV. The protons in the initial states of the $`\overline{\nu }_e`$ absorption and $`\nu _e`$ emission processes in Eq. (1) have $`E_pm_pT1`$ MeV, and therefore, can occupy many Landau levels. By the correspondence principle, the quantum effects of the magnetic field on these protons are insignificant. However, the protons in the final states of the $`\nu _e`$ absorption and $`\overline{\nu }_e`$ emission processes are less energetic, with typical recoil energies of $`E_pm_pE_\nu ^2/m_p100`$ keV and $`T^2/m_p1`$ keV, respectively. Thus, proper treatment of proton Landau levels is especially important for these processes.
The proton wave function can be written as
$$(\psi _p)_{s_p}=\frac{e^{i(k_{pz}zE_pt)}e^{i(r_pn_p)\varphi }}{\sqrt{2\pi L/eB}}(U_p)_{s_p},$$
(30)
where
$`(U_p)_{s_p=1}`$ $`=\left(\begin{array}{c}I_{n_p,r_p}(eB\xi ^2/2)\\ 0\\ (k_{pz}/2m_p)I_{n_p,r_p}(eB\xi ^2/2)\\ ie^{i\varphi }(\sqrt{2n_peB}/2m_p)I_{n_p1,r_p}(eB\xi ^2/2)\end{array}\right)`$ (35)
and
$`(U_p)_{s_p=1}`$ $`=\left(\begin{array}{c}0\\ e^{i\varphi }I_{n_p1,r_p}(eB\xi ^2/2)\\ i(\sqrt{2n_peB}/2m_p)I_{n_p,r_p}(eB\xi ^2/2)\\ e^{i\varphi }(k_{pz}/2m_p)I_{n_p1,r_p}(eB\xi ^2/2)\end{array}\right).`$ (40)
Note that each proton Landau level is also degenerate with respect to the quantum number $`r_p`$, but the spin degeneracy of the excited levels is lifted due to the contribution from the anomalous magnetic moment.
The proton wave function contains terms of the form
$$\mathrm{\Psi }_{n,r}\frac{e^{ik_zz}e^{i(rn)\varphi }}{\sqrt{2\pi L/eB}}I_{n,r}(eB\xi ^2/2),$$
(41)
which has the following properties:
$`\pi _+\mathrm{\Psi }_{n1,r}`$ $`(\pi _xi\pi _y)\mathrm{\Psi }_{n1,r}=i\sqrt{2neB}\mathrm{\Psi }_{n,r},`$ (42a)
and
$`\pi _{}\mathrm{\Psi }_{n,r}`$ $`(\pi _x+i\pi _y)\mathrm{\Psi }_{n,r}=i\sqrt{2neB}\mathrm{\Psi }_{n1,r}.`$ (42b)
The operator $`𝝅`$ in the above equations is defined as
$$𝝅i\mathbf{}e𝐀.$$
(43)
Equations (42a) and (42b) can be used to simplify the evaluation of the transition amplitudes for the processes in Eq. (1) Duan (2004).
### II.3 Neutron
Neutrons are also nonrelativistic in the supernova environment of interest, and their energy is
$`E_n`$ $`=m_n+{\displaystyle \frac{k_n^2}{2m_n}}s_n\mu _nB,`$ (44)
where $`\mu _n=1.91\mu _N`$ is the neutron magnetic moment. The corresponding wave function to $`𝒪(1/m_N)`$ is
$$(\psi _n)_{s_n}=\frac{e^{i(𝐤_n𝐱E_nt)}}{L^{3/2}}(U_n)_{s_n},$$
(45)
where
$`(U_n)_{s_n=1}`$ $`=\left(\begin{array}{c}1\\ 0\\ (k_n/2m_n)\mathrm{cos}\mathrm{\Theta }_n\\ (k_n/2m_n)\mathrm{sin}\mathrm{\Theta }_ne^{i\mathrm{\Phi }_n}\end{array}\right)`$ (50)
and
$`(U_n)_{s_n=1}`$ $`=\left(\begin{array}{c}0\\ 1\\ (k_n/2m_n)\mathrm{sin}\mathrm{\Theta }_ne^{i\mathrm{\Phi }_n}\\ (k_n/2m_n)\mathrm{cos}\mathrm{\Theta }_n\end{array}\right).`$ (55)
In the above equations, $`\mathrm{\Theta }_n`$ and $`\mathrm{\Phi }_n`$ are the polar and azimuthal angles of the neutron momentum $`𝐤_n`$ in spherical coordinates.
### II.4 Neutrinos
The neutrino energy is not affected by the magnetic field. For left-handed $`\nu _e`$ with momentum $`𝐤_\nu `$, the wave function is
$$\psi _{\nu _e}=\frac{e^{i(𝐤_\nu 𝐱E_\nu t)}}{L^{3/2}}U_\nu ,$$
(56)
where
$$U_\nu =\left(\begin{array}{c}\mathrm{sin}(\mathrm{\Theta }_\nu /2)\\ \mathrm{cos}(\mathrm{\Theta }_\nu /2)\\ \mathrm{sin}(\mathrm{\Theta }_\nu /2)\\ \mathrm{cos}(\mathrm{\Theta }_\nu /2)\end{array}\right).$$
(57)
The azimuthal angle of $`𝐤_\nu `$ is taken to be $`\mathrm{\Phi }_\nu =0`$ in the above equation. The wave function of right-handed $`\overline{\nu }_e`$ with the same momentum $`𝐤_\nu `$ is
$$\psi _{\overline{\nu }_e}=\frac{e^{i(𝐤_\nu 𝐱E_\nu t)}}{L^{3/2}}U_\nu ,$$
(58)
where $`U_\nu `$ is the same as in Eq. (57).
## III Cross sections and differential reaction rates
As discussed in Sec. II, $`B10^{16}\text{ G}`$ will not affect the weak interaction, which is still described by the effective four-fermion Lagrangian
$$_{\text{int}}=\frac{G_F\mathrm{cos}\theta _C}{\sqrt{2}}\left(N_\alpha ^{}L^\alpha +N^\alpha L_\alpha ^{}\right),$$
(59)
where $`G_F=(292.8\text{GeV})^2`$ is the Fermi constant, $`\theta _C`$ is the Cabbibo angle ($`\mathrm{cos}^2\theta _C=0.95`$), the leptonic charged current $`L^\alpha `$ is
$$L^\alpha =\overline{\psi }_\nu \gamma ^\alpha (1\gamma _5)\psi _e,$$
(60)
and the nucleonic current $`N^\alpha `$ is
$$N^\alpha =\overline{\psi }_p\left[f\gamma ^\alpha g\gamma ^\alpha \gamma _5+\frac{if_2}{2m_p}\sigma ^{\alpha \beta }\left(i\stackrel{}{𝒟}_\beta \right)\right]\psi _n.$$
(61)
In the above equations, $`\gamma ^\alpha `$, $`\gamma _5`$, and $`\sigma ^{\alpha \beta }`$ are the standard matrices describing fermionic transitions in Dirac-Pauli representation, and $`f=1`$, $`g=1.26`$, and $`f_2=3.7`$ are the nucleon form factors. \[A more up-to-date value of $`g`$ is 1.27 Eidelman and et al. (Particle Data Group). This value is recommended for calculating the rates of the processes in Eq. (1) for specific application to supernova nucleosynthesis.\] The term involving $`f_2`$ in Eq. (61) represents weak magnetism and must be included for calculations to $`𝒪(1/m_N)`$. The covariant derivative $`i\stackrel{}{𝒟}_\beta `$ in this term preserves the gauge invariance and operates according to
$$\overline{\psi }_pO\left(i\stackrel{}{𝒟}_\beta \right)\psi _n=\left[\left(i_\beta eA_\beta \right)\overline{\psi }_p\right]O\psi _n\overline{\psi }_pO\left(i_\beta \psi _n\right),$$
(62)
where $`O`$ is a constant matrix and $`A_\beta `$ corresponds to the electromagnetic field ($`A_0=0`$ here).
Based on the above description of the weak interaction, we derive below the cross sections of the neutrino absorption processes in Eq. (1) and the differential reaction rates of the reverse neutrino emission processes. We will include the magnetic field effects on particle energies and wave functions and focus on corrections of $`𝒪(1/m_N)`$ in both the transition amplitude and kinematics. Radiative corrections and the effect of the Coulomb field of the proton on the electron wave function are ignored for simplicity. (The Coulomb field will modify the Landau wave function of the electron, thus making the calculation much more complicated.) We propose an approximate treatment of these factors at the end of Sec. IV.1.
### III.1 Cross sections for neutrino absorption
We first derive the cross section of $`\nu _e+ne^{}+p`$ in detail. The transition matrix of this process is
$$𝒯_{\nu _en}=\frac{G_F\mathrm{cos}\theta _C}{\sqrt{2}}\overline{\psi }_p\left[f\gamma ^\alpha g\gamma ^\alpha \gamma _5+\frac{if_2}{2m_p}\sigma ^{\alpha \beta }\left(i\stackrel{}{𝒟}_\beta \right)\right]\psi _n\overline{\psi }_e^{}\gamma _\alpha (1\gamma _5)\psi _{\nu _e}\text{d}^4x.$$
(63)
With the wave functions given in Sec. II, Eq. (63) can be rewritten as
$$𝒯_{\nu _en}=\frac{G_F\mathrm{cos}\theta _C}{\sqrt{2}}\frac{eB}{2\pi L^4}2\pi \delta (E_e+E_pE_\nu E_n)2\pi \delta (k_{ez}+k_{pz}k_{\nu z}k_{nz})𝔐_{\nu _en}.$$
(64)
The amplitude $`𝔐_{\nu _en}`$ in Eq. (64) is
$`𝔐_{\nu _en}`$ $`={\displaystyle _0^{\mathrm{}}}\xi \text{d}\xi {\displaystyle _0^{2\pi }}e^{i𝐰_{}𝐱_{}}e^{i(n_er_en_p+r_p)\varphi }\{\overline{U}_p\gamma ^\alpha (fg\gamma _5)U_n\overline{U}_e^{}\gamma _\alpha (1\gamma _5)U_{\nu _e}`$
$`+{\displaystyle \frac{if_2}{2m_p}}[(X_p)_\beta ^{}\gamma ^0\sigma ^{\alpha \beta }U_n(k_n)_\beta \overline{U}_p\sigma ^{\alpha \beta }U_n]\overline{U}_e^{}\gamma _\alpha (1\gamma _5)U_{\nu _e}\}\text{d}\varphi ,`$ (65)
where $`𝐰=𝐤_n+𝐤_\nu `$ is the total momentum, the subscript $``$ denotes a vector in the $`xy`$-plane, and
$$(X_p)_\beta \left[\frac{e^{i(k_{pz}zE_pt)}e^{i(r_pn_p)\varphi }}{\sqrt{2\pi L/eB}}\right]^1\left(i_\beta eA_\beta \right)\psi _p.$$
(66)
Evaluation of $`(X_p)_\beta `$ for $`\beta =1`$ and 2 ($`x`$ and $`y`$) can be simplified by using Eqs. (42a) and (42b) Duan (2004).
The $`\delta `$-functions in Eq. (64) enforce conservation of energy and of momentum in the $`z`$-direction, for which both the neutron and the proton momenta must be taken into account in calculations to $`𝒪(1/m_N)`$. For $`\nu _e+ne^{}+p`$ occurring in the material above the protoneutron star, the neutron momentum is especially important as the typical value $`k_n\sqrt{2m_nT}=43(T/\text{MeV})^{1/2}\text{MeV}`$ is larger than the typical $`\nu _e`$ momentum $`k_\nu =E_\nu 10`$ MeV. To account for this, we average the cross section over the normalized thermal distribution function $`f_n(𝐤_n,s_n)`$ for the neutrons and obtain
$$\sigma _{\nu _en}^{(1,B)}=\underset{s_n}{}f_n(𝐤_n,s_n)\text{d}^3k_n\underset{s_p,n_p,r_p}{}\frac{L\text{d}k_{pz}}{2\pi }\underset{s_e,n_e,r_e}{}\frac{L\text{d}k_{ez}}{2\pi }\frac{1}{L^3L^3}\frac{|𝒯_{\nu _en}|^2}{\tau L^3},$$
(67)
where the superscript $`(1,B)`$ denotes the cross section to $`𝒪(1/m_N)`$ and in the presence of a magnetic field, $`\tau `$ is the duration of the interaction, and
$$f_n(𝐤_n,s_n)=\frac{e^{(E_nm_n)/T}}{\left(2\pi m_nT\right)^{3/2}\left(e^{\mu _nB/T}+e^{\mu _nB/T}\right)}.$$
(68)
The summation and integration in Eq. (67) must be treated as nested integrals. For example, the summation over $`s_n`$ and the integration over $`𝐤_n`$ apply to not only $`f_n(𝐤_n,s_n)`$ but also the subsequent terms that have implicit dependence on $`s_n`$ and $`𝐤_n`$. The summation and integration in the equations below should be interpreted similarly.
Using
$`{\displaystyle \delta (E_e+E_pE_\nu E_n)\delta (k_{ez}+k_{pz}k_{\nu z}k_{nz})\text{d}^3k_n}`$
$`={\displaystyle \text{d}\mathrm{\Phi }_n\delta (E_e+E_pE_\nu E_n)\text{d}\left(\frac{k_n^2}{2}\right)\delta (k_{ez}+k_{pz}k_{\nu z}k_{nz})\text{d}k_{nz}}`$
$`=m_n{\displaystyle \text{d}\mathrm{\Phi }_n}`$ (69)
to integrate over the neutron momentum, we can rewrite Eq. (67) as
$`\sigma _{\nu _en}^{(1,B)}`$ $`={\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{4\pi }}{\displaystyle \frac{m_neB}{(2\pi m_nT)^{3/2}}}{\displaystyle \frac{1}{e^{\mu _nB/T}+e^{\mu _nB/T}}}`$
$`\times {\displaystyle \underset{n_e=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_p=0}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\text{d}k_{ez}{\displaystyle \underset{s_n=\pm 1}{}}{\displaystyle \underset{s_p=\pm 1}{}}{\displaystyle \underset{K}{}}\text{d}k_{pz}{\displaystyle _0^{2\pi }}e^{(E_nm_n)/T}𝒲_{\nu _en}\text{d}\mathrm{\Phi }_n,`$ (70)
where
$$𝒲_{\nu _en}\frac{eB}{2\pi L^2}\underset{s_e,r_e,r_p}{}|𝔐_{\nu _en}|^2$$
(71)
is the reduced amplitude squared given explicitly in Appendix B. It follows from Eq. (69) that $`E_n`$ and $`k_{nz}`$ in the integrand in Eq. (70) are determined in terms of other quantities by conservation of energy and of momentum in the $`z`$-direction. The integration region $`K`$ of $`\text{d}k_{pz}`$ in Eq. (70) is also set by these conservation laws, which require
$`E_\nu +m_n+{\displaystyle \frac{k_n^2}{2m_n}}+{\displaystyle \frac{(k_{ez}+k_{pz}k_{\nu z})^2}{2m_n}}s_n\mu _nB`$ $`=\sqrt{m_e^2+k_{ez}^2+2n_eeB}`$
$`+m_p+{\displaystyle \frac{k_{pz}^2}{2m_p}}+{\displaystyle \frac{n_peB}{m_p}}s_p\stackrel{~}{\mu }_pB.`$ (72)
The above equation can be rearranged into the form
$$ak_{pz}^2+bk_{pz}+c=\frac{k_n^2}{2m_n}0,$$
(73)
where
$`a`$ $`={\displaystyle \frac{\mathrm{\Delta }}{2m_pm_n}},`$ (74a)
$`b`$ $`={\displaystyle \frac{k_{\nu z}k_{ez}}{m_n}},`$ (74b)
$`c`$ $`=\sqrt{m_e^2+k_{ez}^2+2n_eeB}E_\nu \mathrm{\Delta }{\displaystyle \frac{(k_{\nu z}k_{ez})^2}{2m_n}}+{\displaystyle \frac{n_peB}{m_p}}s_p\stackrel{~}{\mu }_pB+s_n\mu _nB,`$ (74c)
$`\mathrm{\Delta }`$ $`m_nm_p.`$ (74d)
Thus,
$$K=\{\begin{array}{cc}(\mathrm{},+\mathrm{}),\hfill & \text{if }b^24ac\text{,}\hfill \\ (\mathrm{},(k_{pz})_{}][(k_{pz})_+,+\mathrm{}),\hfill & \text{if }b^2>4ac\text{,}\hfill \end{array}$$
(75)
where
$$(k_{pz})_\pm =\frac{b\pm \sqrt{b^24ac}}{2a}.$$
(76)
For $`\overline{\nu }_e+pe^++n`$ occurring in the material above the protoneutron star, the cross section can be written as
$$\sigma _{\overline{\nu }_ep}^{(1,B)}=\overline{\underset{r_p}{}}\underset{s_p,n_p}{}f_p(k_{pz},n_p,s_p)\text{d}k_{pz}\underset{s_n}{}\frac{L^3\text{d}^3k_n}{(2\pi )^3}\underset{s_e,n_e,r_e}{}\frac{L\text{d}k_{ez}}{2\pi }\frac{1}{L^3L^3}\frac{|𝒯_{\overline{\nu }_ep}|^2}{\tau L^3},$$
(77)
where
$$\overline{\underset{r_p}{}}\left(\frac{eBL^2}{2\pi }\right)^1\underset{r_p}{}$$
(78)
and
$$f_p(k_{pz},n_p,s_p)=\frac{e^{(E_pm_p)/T}}{\sqrt{2\pi m_pT}}\frac{1e^{eB/m_pT}}{e^{\stackrel{~}{\mu }_pB/T}+e^{(\stackrel{~}{\mu }_pB/T)(eB/m_pT)}}$$
(79)
is the normalized thermal distribution function for the protons. Using again the integration over the neutron momentum to get rid of the $`\delta `$-functions, we obtain
$`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ $`={\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{8\pi ^2}}{\displaystyle \frac{m_n}{\sqrt{2\pi m_pT}}}{\displaystyle \frac{1e^{eB/m_pT}}{e^{\stackrel{~}{\mu }_pB/T}+e^{(\stackrel{~}{\mu }_pB/T)(eB/m_pT)}}}`$
$`\times {\displaystyle \underset{n_e=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_p=0}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\text{d}k_{ez}{\displaystyle \underset{s_n=\pm 1}{}}{\displaystyle \underset{s_p=\pm 1}{}}{\displaystyle \underset{K^{}}{}}\text{d}k_{pz}{\displaystyle _0^{2\pi }}e^{(E_pm_p)/T}𝒲_{\overline{\nu }_ep}\text{d}\mathrm{\Phi }_n.`$ (80)
The reduced amplitude squared $`𝒲_{\overline{\nu }_ep}`$ in the above equation can be obtained from $`𝒲_{\nu _en}`$ by making the substitution
$$\begin{array}{cc}\hfill (E_\nu ,𝐤_\nu )& (E_\nu ,𝐤_\nu ),\hfill \\ \hfill (E_e,k_{ez})& (E_e,k_{ez}).\hfill \end{array}$$
(81)
The integration region $`K^{}`$ of $`\text{d}k_{pz}`$ in Eq. (80) is determined from energy and momentum conservation as in Eq. (72) but with the above substitution implemented.
For application to supernova neutrinos, it is useful to further average the cross sections in Eqs. (70) and (80) over the relevant normalized neutrino energy spectra $`f_\nu (E_\nu )`$ to obtain
$$\sigma _{\nu N}=\sigma _{\nu N}f_\nu (E_\nu )\text{d}E_\nu ,$$
(82)
where $`\sigma _{\nu N}`$ stands for $`\sigma _{\nu _en}^{(1,B)}`$ or $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$. A typical form of $`f_\nu (E_\nu )`$ adopted in the literature is
$$f_\nu (E_\nu )=\frac{1}{T_\nu ^3F_2(\eta _\nu )}\frac{E_\nu ^2}{\mathrm{exp}\left[\left(E_\nu /T_\nu \right)\eta _\nu \right]+1},$$
(83)
where $`T_\nu `$ and $`\eta _\nu `$ are constant parameters and
$$F_n(\eta _\nu )_0^{\mathrm{}}\frac{x^n}{\mathrm{exp}\left(x\eta _\nu \right)+1}\text{d}x.$$
(84)
For the neutrino energy spectra in Eq. (83), the average neutrino energy is
$$E_\nu =\frac{F_3(\eta _\nu )}{F_2(\eta _\nu )}T_\nu .$$
(85)
### III.2 Differential reaction rates for neutrino emission
As can be seen from Eqs. (70), (80), and (82), the cross sections $`\sigma _{\nu _en}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ for the neutrino absorption processes $`\nu _e+ne^{}+p`$ and $`\overline{\nu }_e+pe^++n`$, respectively, have the same generic form
$$\text{d}E_\nu \underset{n_e=0}{\overset{\mathrm{}}{}}\underset{n_p=0}{\overset{\mathrm{}}{}}_{\mathrm{}}^+\mathrm{}\text{d}k_{ez}\underset{s_n=\pm 1}{}\underset{s_p=\pm 1}{}\underset{\stackrel{~}{K}}{}\text{d}k_{pz}_0^{2\pi }\text{d}\mathrm{\Phi }_n,$$
(86)
where $`\stackrel{~}{K}=K\text{ or }K^{}`$, and $``$ is the integrand involving the relevant amplitude squared and distribution functions. If we use the differential reaction rates with respect to $`\mathrm{cos}\mathrm{\Theta }_\nu `$ to describe the neutrino emission processes $`e^{}+p\nu _e+n`$ and $`e^++n\overline{\nu }_e+p`$, then these rates also have the generic form in Eq. (86). This follows from the symmetry between the forward and reverse processes. In particular, the transition amplitudes squared $`\left|𝒯_{e^{}p}\right|^2`$ and $`\left|𝒯_{e^+n}\right|^2`$ are identical to $`\left|𝒯_{\nu _en}\right|^2`$ and $`\left|𝒯_{\overline{\nu }_ep}\right|^2`$, respectively. By taking advantage of the symmetry between the neutrino absorption and emission processes, numerical computation of the cross sections for the former and the differential reaction rates for the latter is greatly simplified.
For $`e^{}+p\nu _e+n`$ occurring in the material above the protoneutron star, the differential reaction rate is
$`{\displaystyle \frac{\text{d}\lambda _{e^{}p}^{(1,B)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }}`$ $`=`$ $`\overline{{\displaystyle \underset{r_p}{}}}{\displaystyle \underset{s_p,n_p}{}}{\displaystyle f_p(k_{pz},n_p,s_p)\text{d}k_{pz}\underset{s_e,n_e,r_e}{}\frac{L\text{d}k_{ez}}{2\pi }\frac{1}{L^3}\frac{1}{e^{(E_e/T)\eta _e}+1}}`$ (87)
$`\times {\displaystyle }{\displaystyle \frac{L^3E_\nu ^2\text{d}E_\nu }{4\pi ^2}}{\displaystyle \underset{s_n}{}}{\displaystyle }{\displaystyle \frac{L^3\text{d}^3k_n}{(2\pi )^3}}{\displaystyle \frac{1}{L^3L^3}}{\displaystyle \frac{\left|𝒯_{\nu _en}\right|^2}{\tau L^3}},`$
where $`\eta _e`$ is the degeneracy parameter characterizing the Fermi-Dirac distribution function of the electrons. Integrating over the neutron momentum as in Eq. (69), we obtain
$`{\displaystyle \frac{\text{d}\lambda _{e^{}p}^{(1,B)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }}`$ $`=`$ $`{\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{32\pi ^4}}{\displaystyle \frac{m_n}{\sqrt{2\pi m_pT}}}{\displaystyle \frac{1e^{eB/m_pT}}{e^{\stackrel{~}{\mu }_pBT}+e^{(\stackrel{~}{\mu }_pB/T)(eB/m_pT)}}}{\displaystyle _0^{\mathrm{}}}E_\nu ^2\text{d}E_\nu `$ (88)
$`\times {\displaystyle \underset{n_e=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_p=0}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\text{d}k_{ez}{\displaystyle \underset{s_n=\pm 1}{}}{\displaystyle \underset{s_p=\pm 1}{}}{\displaystyle \underset{K}{}}\text{d}k_{pz}`$
$`\times {\displaystyle _0^{2\pi }}{\displaystyle \frac{e^{(E_pm_p)/T}}{e^{(E_e/T)\eta _e}+1}}𝒲_{\nu _en}\text{d}\mathrm{\Phi }_n.`$
Similarly, we obtain the differential reaction rate for $`e^++n\overline{\nu }_e+p`$ as
$`{\displaystyle \frac{\text{d}\lambda _{e^+n}^{(1,B)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }}`$ $`=`$ $`{\displaystyle \underset{s_n}{}}{\displaystyle f_n(𝐤_n,s_n)\text{d}^3k_n\underset{s_e,n_e,r_e}{}\frac{L\text{d}k_{ez}}{2\pi }\frac{1}{L^3}\frac{1}{e^{(E_e/T)+\eta _e}+1}}`$ (89)
$`\times {\displaystyle }{\displaystyle \frac{L^3E_\nu ^2\text{d}E_\nu }{4\pi ^2}}{\displaystyle \underset{s_p,n_p,r_p}{}}{\displaystyle }{\displaystyle \frac{L\text{d}k_{pz}}{2\pi }}{\displaystyle \frac{1}{L^3L^3}}{\displaystyle \frac{\left|𝒯_{\overline{\nu }_ep}\right|^2}{\tau L^3}}`$
$`=`$ $`{\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{16\pi ^3}}{\displaystyle \frac{m_neB}{(2\pi m_nT)^{3/2}}}{\displaystyle \frac{1}{e^{\mu _nB/T}+e^{\mu _nB/T}}}{\displaystyle _0^{\mathrm{}}}E_\nu ^2\text{d}E_\nu `$
$`\times {\displaystyle \underset{n_e=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_p=0}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\text{d}k_{ez}{\displaystyle \underset{s_n=\pm 1}{}}{\displaystyle \underset{s_p=\pm 1}{}}{\displaystyle \underset{K^{}}{}}\text{d}k_{pz}`$
$`\times {\displaystyle _0^{2\pi }}{\displaystyle \frac{e^{(E_nm_n)/T}}{e^{(E_e/T)+\eta _e}+1}}𝒲_{\overline{\nu }_ep}\text{d}\mathrm{\Phi }_n.`$
## IV Rates of Neutrino Processes in Supernovae
We now calculate the rates of the neutrino absorption and emission processes in Eq. (1) for the supernova environment near a protoneutron star that possesses a strong magnetic field. A wide range of heavy elements may be produced during the ejection of the material above the protoneutron star. As mentioned in the introduction, a key parameter governing this production is the neutron-to-proton ratio of the material Woosley and Hoffman (1992); Hoffman et al. (1996), which depends on the competition between the neutrino absorption and emission processes Qian and Woosley (1996). We will calculate the rates of these processes in the context of heavy element nucleosynthesis, for which the accuracy of these rates is especially important. In this context, the material above the protoneutron star is characterized by temperatures of $`T1`$ MeV, entropies of $`S100`$ (in units of Boltzmann constant per nucleon), and electron fractions of $`Y_e0.5`$. For these conditions, the nucleons in the material are nonrelativistic and nondegenerate while the $`e^\pm `$ are relativistic and have a small degeneracy parameter $`0<\eta _e1`$. The thermal distribution functions of the nucleons and $`e^\pm `$ have been given in Sec. III. The neutrinos emitted from the protoneutron star are not in thermal equilibrium with the overlying material and their energy distribution functions are taken to be of the form in Eq. (83). As discussed in Ref. Duan and Qian (2004), Pauli blocking for the final states of the neutrino processes above the protoneutron star is unimportant and will be ignored.
### IV.1 Neutrino absorption
At a radius $`R`$ above the protoneutron star, the rate of neutrino absorption per nucleon can be estimated as
$$\lambda _{\nu N}=\frac{L_\nu \sigma _{\nu N}}{4\pi R^2E_\nu }=49.7\left(\frac{L_\nu }{10^{51}\mathrm{erg}\mathrm{s}^1}\right)\left(\frac{10\mathrm{MeV}}{E_\nu }\right)\left(\frac{\sigma _{\nu N}}{10^{41}\mathrm{cm}^2}\right)\left(\frac{10\mathrm{km}}{R}\right)^2\mathrm{s}^1,$$
(90)
where $`L_\nu `$ is the neutrino luminosity and has a typical value of $`10^{51}`$ erg s<sup>-1</sup> in the supernova epoch of interest. The key quantity $`\sigma _{\nu N}`$ in the above equation is obtained by averaging $`\sigma _{\nu N}`$ over the neutrino energy spectrum. We first compare various approximations for $`\sigma _{\nu N}`$ as functions of the neutrino energy $`E_\nu `$.
The cross sections for neutrino absorption on nucleons in a magnetic field have been derived to $`𝒪(1/m_N)`$ as $`\sigma _{\nu N}^{(1,B)}`$ in Sec. III.1. To $`𝒪(1)`$, the zeroth order in $`1/m_N`$, the cross sections are Duan and Qian (2004)
$`\sigma _{\nu N}^{(0,B)}`$ $`=\sigma _{B,1}\left[1+2\chi {\displaystyle \frac{(f\pm g)g}{f^2+3g^2}}\mathrm{cos}\mathrm{\Theta }_\nu \right]`$
$`+\sigma _{B,2}\left[{\displaystyle \frac{f^2g^2}{f^2+3g^2}}\mathrm{cos}\mathrm{\Theta }_\nu +2\chi {\displaystyle \frac{(fg)}{f^2+3g^2}}\right],`$ (91)
where
$`\sigma _{B,1}`$ $`={\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{2\pi }}(f^2+3g^2)eB{\displaystyle \underset{n_e=0}{\overset{n_{e,\text{max}}}{}}}{\displaystyle \frac{g_{n_e}E_e^{(0)}}{\sqrt{(E_e^{(0)})^2m_e^22n_eeB}}},`$ (92)
$`\sigma _{B,2}`$ $`={\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{2\pi }}(f^2+3g^2)eB{\displaystyle \frac{E_e^{(0)}}{\sqrt{(E_e^{(0)})^2m_e^2}}},`$ (93)
$`E_e^{(0)}`$ $`=E_\nu \pm \mathrm{\Delta },`$ (94)
$`n_{e,\text{max}}`$ $`=\left[{\displaystyle \frac{(E_e^{(0)})^2m_e^2}{2eB}}\right]_{\text{int}}.`$ (95)
In the above equations and elsewhere in this subsection, the upper sign is for $`\nu _e+ne^{}+p`$ and the lower sign is for $`\overline{\nu }_e+pe^++n`$. In Eq. (91),
$$\chi =\frac{\mathrm{exp}(\mu B/T)\mathrm{exp}(\mu B/T)}{\mathrm{exp}(\mu B/T)+\mathrm{exp}(\mu B/T)}$$
(96)
is the polarization of nucleon spin, where $`\mu `$ is the magnetic moment of the relevant nucleon with $`\mu _p=2.79\mu _N`$ and $`\mu _n=1.91\mu _N`$. For the case of interest here, $`|\mu |B/T1`$, so
$$\chi =\frac{\mu B}{T}=3.15\times 10^2\left(\frac{\mu }{\mu _N}\right)\left(\frac{B}{10^{16}\mathrm{G}}\right)\left(\frac{\text{MeV}}{T}\right).$$
(97)
The term proportional to $`\sigma _{B,2}`$ in Eq. (91) arises because the ground Landau level of $`e^\pm `$ has only one spin state while any other level has two. In Eq. (92) for $`\sigma _{B,1}`$, the product of $`eB`$ and the sum gives the total phase space of the $`e^\pm `$ in the final state. In the limit $`n_{e,\text{max}}1`$, the summation of the Landau levels can be replaced by integration and $`\sigma _{B,1}`$ approaches
$$\sigma _{\nu N}^{(0)}=\frac{G_F^2\mathrm{cos}^2\theta _C}{\pi }E_e^{(0)}\sqrt{(E_e^{(0)})^2m_e^2},$$
(98)
which is the cross section to $`𝒪(1)`$ in the absence of any magnetic field. In the same limit, $`\sigma _{B,2}`$ is negligible compared with $`\sigma _{B,1}`$, so $`\sigma _{\nu N}^{(0,B)}`$ approaches $`\sigma _{\nu N}^{(0)}(1+ϵ_\chi )`$, where
$$ϵ_\chi =\chi \frac{2(f\pm g)g}{f^2+3g^2}\mathrm{cos}\mathrm{\Theta }_\nu $$
(99)
results from the polarization of the initial nucleon spin by the magnetic field.
For numerical examples of the cross sections, we take $`B=10^{16}\text{G}`$. The cross sections $`\sigma _{\nu _en}^{(0,B)}`$ for $`\mathrm{cos}\mathrm{\Theta }_\nu =1`$, 0, and 1 as functions of $`E_\nu `$ are shown as the dotted curves in Figs. 1a, 1b, and 1c, respectively. The angle-dependent terms in Eq. (91) for $`\sigma _{\overline{\nu }_ep}^{(0,B)}`$ are proportional to the difference between $`f`$ and $`g`$. As the numerical values of $`f`$ and $`g`$ are close, these terms are very small. So we only show the cross section $`\sigma _{\overline{\nu }_ep}^{(0,B)}`$ for $`\mathrm{cos}\mathrm{\Theta }_\nu =0`$ as the dotted curve in Fig. 1d. All the dotted curves in Fig. 1 have spikes superposed on a general trend. The varying heights of these spikes are artifacts of the plotting tool: all the spikes should have been infinitely high as they correspond to “resonances” at $`E_e^{(0)}=\sqrt{m_e^2+2n_eeB}`$, for which a new Landau level opens up. These singularities are integrable and do not give infinite probabilities of interaction in practice. For example, at a given $`E_\nu `$, the thermal motion of the absorbing nucleons will produce a range of $`E_e`$ and the cross section obtained from integration over this range will be finite. Thus, the spikes in $`\sigma _{\nu N}^{(0,B)}`$ will be smeared out by the thermal motion of the absorbing nucleons, which is similar to the Doppler broadening of the photon absorption lines in the solar light spectrum. The effects of such motion are of $`𝒪(1/m_N)`$ and have been taken into account by the cross sections $`\sigma _{\nu N}^{(1,B)}`$ derived in Sec. III.1. Using $`T=2`$ MeV for illustration, we show $`\sigma _{\nu N}^{(1,B)}`$ as the solid curves in Fig. 1. It can be seen that where spikes occur in $`\sigma _{\nu N}^{(0,B)}`$, there are only smooth bumps in $`\sigma _{\nu N}^{(1,B)}`$. Clearly, $`\sigma _{\nu N}^{(1,B)}`$ is more physical than $`\sigma _{\nu N}^{(0,B)}`$.
Two more aspects of Fig. 1 require discussion. First, the bumps in $`\sigma _{\nu N}^{(1,B)}`$ diminish as $`E_\nu `$ increases and become invisible for $`E_\nu \sqrt{eB}8`$ MeV. This is expected from the correspondence principle: when a number of Landau levels for $`e^\pm `$ and protons can be occupied, the quantum effects of the magnetic field are small. As noted in Sec. II.2, the absorbing proton in $`\overline{\nu }_e+pe^++n`$ can occupy many levels for $`T1`$ MeV. However, for the $`e^+`$ in $`\overline{\nu }_e+pe^++n`$ and the $`e^{}`$ and the proton in $`\nu _e+ne^{}+p`$, occupation of many levels requires $`E_\nu \sqrt{eB}8`$ MeV (see Secs. II.1 and II.2). Second, while the general trends of the dotted curves for $`\sigma _{\nu _en}^{(0,B)}`$ appear to follow the corresponding solid curves for $`\sigma _{\nu _en}^{(1,B)}`$, the general trend of the dotted curve for $`\sigma _{\overline{\nu }_ep}^{(0,B)}`$ deviates substantially from the corresponding solid curve for $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$. This concerns the effects of weak magnetism and recoil of the final-state nucleons, both of which are of $`𝒪(1/m_N)`$ and are taken into account by $`\sigma _{\nu N}^{(1,B)}`$ but not by $`\sigma _{\nu N}^{(0,B)}`$. Figure 1 shows that these effects give small corrections to $`\sigma _{\nu _en}^{(0,B)}`$ but much larger corrections to $`\sigma _{\overline{\nu }_ep}^{(0,B)}`$.
To better understand the effects of weak magnetism and recoil of the final-state nucleons, we make use of the correspondence principle. As noted above, the effects of Landau levels become negligible for $`E_\nu \sqrt{eB}8`$ MeV. In this case, the only surviving quantum effect of the magnetic field is polarization of the initial nucleon spin, which gives rise to a dependence on $`\mathrm{cos}\mathrm{\Theta }_\nu `$ for the cross sections due to parity violation of weak interaction. Thus, allowing for this surviving effect, we should recover the results for no magnetic field in the limit of high $`E_\nu `$. In the absence of any field, the cross sections $`\sigma _{\nu N}^{(1)}`$ to $`𝒪(1/m_N)`$ is (see, e.g., Refs. Vogel and Beacom (1999); Horowitz (2002))
$$\sigma _{\nu N}^{(1)}=\sigma _{\nu N}^{(0)}\left\{1\frac{2[f^22(f+f_2)g+5g^2]}{f^2+3g^2}\frac{E_\nu }{m_N}\right\},$$
(100)
where we have ignored terms like $`m_e^2/E_\nu ^2`$ and $`\mathrm{\Delta }/m_N`$. The above zero-field cross sections assume that the initial nucleon spin is unpolarized, and therefore, do not depend on $`\mathrm{cos}\mathrm{\Theta }_\nu `$. If a small polarization $`\chi `$ is artificially imposed, the modified cross sections should have an additional factor $`1+ϵ_\chi `$. The term proportional to $`E_\nu /m_N`$ in Eq. (100) represents the effects of weak magnetism and recoil of the final-state nucleons. The coefficient in this term is 1.01 for $`\nu _e+ne^{}+p`$ and $`7.21`$ for $`\overline{\nu }_e+pe^++n`$. Therefore, over the range $`E_\nu 10`$–50 MeV typical of supernova neutrinos, the correction due to the above effects is $`1`$–5% for the former reaction but amounts to $`7`$% to $`36`$% for the latter reaction. The importance of these corrections has been discussed in other contexts Horowitz and Li (1999); Vogel and Beacom (1999).
Using $`\chi _n=0.03`$ and $`\chi _p=0.04`$ corresponding to $`B=10^{16}`$ G and $`T=2`$ MeV, we show $`\sigma _{\nu N}^{(0)}(1+ϵ_\chi )`$ and $`\sigma _{\nu N}^{(1)}(1+ϵ_\chi )`$ as the short-dashed and dot-dashed curves, respectively, in Fig. 1. The small increase from $`\sigma _{\nu _en}^{(0)}`$ to $`\sigma _{\nu _en}^{(1)}`$ and the much larger decrease from $`\sigma _{\overline{\nu }_ep}^{(0)}`$ to $`\sigma _{\overline{\nu }_ep}^{(1)}`$ given in Eq. (100) can be seen from this figure. In addition, as expected from the correspondence principle, the general trends of the dotted curves for $`\sigma _{\nu _en}^{(0,B)}`$ closely follow the short-dashed curves for $`\sigma _{\nu _en}^{(0)}(1+ϵ_{\chi _n})`$ at $`E_\nu 20`$ MeV and the solid curves for $`\sigma _{\nu _en}^{(1,B)}`$ become indistinguishable from the dot-dashed curves for $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ in the same regime (see Figs. 1a–c). However, while the relation between the dotted curve for $`\sigma _{\overline{\nu }_ep}^{(0,B)}`$ and the short-dashed curve for $`\sigma _{\overline{\nu }_ep}^{(0)}(1+ϵ_{\chi _p})`$ is in accordance with the correspondence principle, the solid curve for $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ clearly stays above the dot-dashed curve for $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ at $`E_\nu 25`$ MeV (see Fig. 1d). This apparent violation of the correspondence principle for $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ is caused by the slightly different treatments of the reaction kinematics in calculating $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}`$.
We have used the transition amplitudes to $`𝒪(1/m_N)`$ in calculating both $`\sigma _{\nu N}^{(1,B)}`$ and $`\sigma _{\nu N}^{(1)}`$. However, we have treated the reaction kinematics exactly for $`\sigma _{\nu N}^{(1,B)}`$ \[assuming nonrelativistic nucleons, see Eqs. (7276)\] but only to $`𝒪(1/m_N)`$ for $`\sigma _{\nu N}^{(1)}`$. This difference does not affect the comparison between $`\sigma _{\nu _en}^{(1,B)}`$ and $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ as the total correction from weak magnetism and nucleon recoil is small to $`𝒪(1/m_N)`$ in this case, and the terms of orders higher than $`𝒪(1/m_N)`$ are even smaller. In contrast, the importance of the weak magnetism and recoil effects for $`\overline{\nu }_e+pe^++n`$ enables terms of orders higher than $`𝒪(1/m_N)`$ to give rather large corrections to the cross section. Such terms are included in $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ due to exact treatment of the reaction kinematics but not in $`\sigma _{\overline{\nu }_ep}^{(1)}`$. In Ref. Strumia and Vissani (2003), the zero-field cross sections for neutrino absorption on nucleons were derived to $`𝒪(1/m_N)`$ but with reaction kinematics treated exactly. We denote these cross sections as $`\sigma _{\nu N}^{(1)}`$. For consistency with the rest of the paper, we ignore radiative corrections and the effect of the Coulomb interaction between the final-state particles for $`\nu _e+ne^{}+p`$, both of which were taken into account in Ref. Strumia and Vissani (2003). It was shown in this reference that $`\sigma _{\overline{\nu }_ep}^{(1)}`$ is more accurate than $`\sigma _{\overline{\nu }_ep}^{(1)}`$. We show $`\sigma _{\nu N}^{(1)}(1+ϵ_\chi )`$ as the long-dashed curves in Fig. 1. It can be seen that the long-dashed curves for $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ are indistinguishable from the corresponding dot-dashed curves for $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ over the range of $`E_\nu `$ shown but the long-dashed curve for $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ lies significantly above the dot-dashed curve for $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ at $`E_\nu 25`$ MeV. In addition, the solid curves for $`\sigma _{\nu _en}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ settle down to the corresponding long-dashed curves for $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ at $`E_\nu 20`$ and 25 MeV, respectively. Thus, the cross sections $`\sigma _{\nu _en}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ calculated above are in full agreement with the correspondence principle.
We now calculate the average cross sections $`\sigma _{\nu N}`$ using the neutrino energy spectra in Eq. (83). We take $`\eta _{\nu _e}=\eta _{\overline{\nu }_e}=3`$, $`E_{\nu _e}=11\text{MeV}`$, and $`E_{\overline{\nu }_e}=16\text{MeV}`$. For these parameters, $`T_{\nu _e}=2.75\text{MeV}`$ and $`T_{\overline{\nu }_e}=4\text{MeV}`$. Adopting the same $`B`$, $`T`$, $`\chi _n`$, and $`\chi _p`$ as for Fig. 1, we give $`\sigma _{\nu N}^{(0)}(1+ϵ_\chi )`$, $`\sigma _{\nu N}^{(1)}(1+ϵ_\chi )`$, $`\sigma _{\nu N}^{(1)}(1+ϵ_\chi )`$, $`\sigma _{\nu N}^{(0,B)}`$, and $`\sigma _{\nu N}^{(1,B)}`$ for $`\mathrm{cos}\mathrm{\Theta }_\nu =1`$, 0, and 1, respectively, in Table 1. As discussed above, $`\sigma _{\nu N}^{(1)}`$ and $`\sigma _{\nu N}^{(1)}`$ differ from $`\sigma _{\nu N}^{(0)}`$ due to the effects of weak magnetism and recoil of the final-state nucleons. These effects slightly increase the cross sections for $`\nu _e+ne^{}+p`$ but substantially decrease those for $`\overline{\nu }_e+pe^++n`$. As can be seen from Table 1, $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ and $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ are only a few percent larger than $`\sigma _{\nu _en}^{(0)}(1+ϵ_{\chi _n})`$ but $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$ are $`20`$% smaller than $`\sigma _{\overline{\nu }_ep}^{(0)}(1+ϵ_{\chi _p})`$. The differences between $`\sigma _{\nu N}^{(1,B)}`$ and $`\sigma _{\nu N}^{(0,B)}`$ are similar. On the other hand, the effects of the magnetic field on the average cross sections are small for both $`\nu _e+ne^{}+p`$ and $`\overline{\nu }_e+pe^++n`$. For $`B=10^{16}`$ G assumed above, $`\sigma _{\nu _en}^{(1,B)}`$ is at most 4% smaller than $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ or $`\sigma _{\nu _en}^{(1)}(1+ϵ_{\chi _n})`$ while $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ is indistinguishable from $`\sigma _{\overline{\nu }_ep}^{(1)}(1+ϵ_{\chi _p})`$. This is because with $`E_{\nu _e}=11\text{MeV}`$ and $`E_{\overline{\nu }_e}=16\text{MeV}`$, the important energy range for determining the average cross sections has $`E_\nu >\sqrt{eB}8`$ MeV, for which the effects of Landau levels are small. As $`E_{\overline{\nu }_e}`$ is substantially larger than $`E_{\nu _e}`$, the magnetic field affects $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ even less than $`\sigma _{\nu _en}^{(1,B)}`$.
The results for absorption of supernova neutrinos on nucleons can be summarized as follows. Generally speaking, one can use $`\sigma _{\nu N}^{(1)}\sigma _{\nu N}^{(0)}`$ to estimate the corrections without magnetic fields, and use $`\sigma _{\nu N}^{(0,B)}\sigma _{\nu N}^{(0)}`$ to estimate the corrections due to magnetic fields. For $`B10^{16}`$ G, $`\sigma _{\nu _en}^{(1)}+\left[\sigma _{\nu _en}^{(0,B)}\sigma _{\nu _en}^{(0)}\right]`$ is a good estimate of $`\sigma _{\nu _en}^{(1,B)}`$ with an accuracy of $`1`$%. For the same field strength, the effects of magnetic fields are not important for $`\sigma _{\overline{\nu }_ep}`$, and $`\sigma _{\overline{\nu }_ep}^{(1)}`$ is a good estimate of $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$ with an accuracy of $`1`$%. Note that we have ignored radiative corrections (see, e.g., Kurylov et al. (2003)) and the effect of the Coulomb field of the proton on the electron wave function (see, e.g., Strumia and Vissani (2003)). These factors give corrections at the level of $`2`$% Strumia and Vissani (2003). To account for them, we suggest calculating $`\sigma _{\nu _en}^{(1)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}`$ as in Ref. Strumia and Vissani (2003) and use the results in the above estimates for $`\sigma _{\nu _en}^{(1,B)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1,B)}`$.
### IV.2 Neutrino emission
The differential reaction rates with respect to $`\mathrm{cos}\mathrm{\Theta }_\nu `$ for $`e^{}+p\nu _e+n`$ and $`e^++n\overline{\nu }_e+p`$ in a magnetic field have been derived to $`𝒪(1/m_N)`$ as $`\text{d}\lambda _{e^{}p}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ and $`\text{d}\lambda _{e^+n}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$, respectively, in Sec. III.2. To $`𝒪(1)`$, the zeroth order in $`1/m_N`$, the differential reaction rates are Duan and Qian (2004)
$$\frac{\text{d}\lambda _{eN}^{(0,B)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }=\frac{eB}{2\pi ^2}\underset{n_e}{}g_{n_e}_0^{\mathrm{}}\frac{\text{d}\mathrm{\Gamma }_{eN}^{(0,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }{\mathrm{exp}[(E_e/T)\eta _e]+1}\text{d}k_{ez},$$
(101)
where
$`{\displaystyle \frac{\text{d}\mathrm{\Gamma }_{eN}^{(0,B)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }}`$ $`={\displaystyle \frac{\mathrm{\Gamma }_{eN}^{(0)}}{2}}\left[1+2\chi {\displaystyle \frac{(fg)g}{f^2+3g^2}}\mathrm{cos}\mathrm{\Theta }_\nu \right]`$
$`+\delta _{n_e,0}{\displaystyle \frac{\mathrm{\Gamma }_{eN}^{(0)}}{2}}\left[{\displaystyle \frac{f^2g^2}{f^2+3g^2}}\mathrm{cos}\mathrm{\Theta }_\nu +2\chi {\displaystyle \frac{(f\pm g)g}{f^2+3g^2}}\right],`$ (102)
$`\mathrm{\Gamma }_{eN}^{(0)}`$ $`={\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{2\pi }}(f^2+3g^2)(E_e\mathrm{\Delta })^2.`$ (103)
In the above equations and elsewhere in this subsection, the upper sign is for $`e^{}+p\nu _e+n`$ and the lower sign is for $`e^++n\overline{\nu }_e+p`$. In Eq. (102), $`\delta _{n_e,0}`$ is the Kronecker delta. For comparison, in the absence of any magnetic field, the differential reaction rates to $`𝒪(1)`$ are
$$\frac{\text{d}\lambda _{eN}^{(0)}}{\text{d}\mathrm{cos}\mathrm{\Theta }_\nu }=\frac{\mathrm{\Gamma }_{eN}^{(0)}}{\mathrm{exp}[(E_e/T)\eta _e]+1}\frac{\text{d}^3k_e}{(2\pi )^3}.$$
(104)
We have also calculated $`\text{d}\lambda _{eN}^{(1)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ to $`𝒪(1/m_N)`$ using the prescription in Ref. Strumia and Vissani (2003). Note that both $`\text{d}\lambda _{eN}^{(0)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ and $`\text{d}\lambda _{eN}^{(1)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ are independent of $`\mathrm{cos}\mathrm{\Theta }_\nu `$.
To compare $`\text{d}\lambda _{eN}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ with $`\text{d}\lambda _{eN}^{(0,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$, we take $`B=10^{16}`$ G, $`T=2`$ MeV, and $`\eta _e=0`$. The differential reaction rates $`\text{d}\lambda _{eN}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ are numerically calculated for $`\mathrm{\Theta }_\nu =0`$, $`\pi /4`$, $`\pi /2`$, $`3\pi /4`$, and $`\pi `$ and shown as the filled circles with error bars in Fig. 2. Here and elsewhere in this subsection, the error bars for our results represent the accuracy of the numerical calculation. To very good approximation, the rates $`\text{d}\lambda _{eN}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ are linear functions of $`\mathrm{cos}\mathrm{\Theta }_\nu `$ as shown by the solid lines fitted to the numerical results in Fig. 2. The rates $`\text{d}\lambda _{eN}^{(0,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ as functions of $`\mathrm{cos}\mathrm{\Theta }_\nu `$ are shown as the dashed lines in the same figure. It can be seen that relative to $`\text{d}\lambda _{e^{}p}^{(0,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$, $`\text{d}\lambda _{e^{}p}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ is smaller for $`\mathrm{cos}\mathrm{\Theta }_\nu 0.15`$ but larger for $`\mathrm{cos}\mathrm{\Theta }_\nu <0.15`$ due to corrections of $`𝒪(1/m_N)`$. So we expect that when integrated over $`\mathrm{cos}\mathrm{\Theta }_\nu `$, the difference between $`\lambda _{e^{}p}^{(0,B)}`$ and $`\lambda _{e^{}p}^{(1,B)}`$ is small. In contrast, corrections of $`𝒪(1/m_N)`$ make $`\text{d}\lambda _{e^+n}^{(1,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ smaller than $`\text{d}\lambda _{e^+n}^{(0,B)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ for all values of $`\mathrm{cos}\mathrm{\Theta }_\nu `$. As discussed in the case of neutrino absorption, such corrections are due to the effects of weak magnetism and recoil of the final-state nucleons, which tend to affect the processes involving $`\overline{\nu }_e`$ more than those involving $`\nu _e`$. These corrections can also be seen by comparing the zero-field results $`\text{d}\lambda _{eN}^{(0)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$ and $`\text{d}\lambda _{eN}^{(1)}/\text{d}\mathrm{cos}\mathrm{\Theta }_\nu `$, which are shown as the dotted and dot-dashed lines, respectively, in Fig. 2. Note that for the parameters adopted above, the magnetic field decreases the rates for $`e^{}+p\nu _e+n`$ but increases those for $`e^++n\overline{\nu }_e+p`$.
To further explore the effects of the magnetic field on the rates of neutrino emission, we consider two representative sets of supernova conditions: $`(T,S,Y_e)=(2\mathrm{MeV},50,0.5)`$ and $`(1\mathrm{MeV},100,0.5)`$ for cases I and II, respectively. For each case, the electron degeneracy parameter $`\eta _e`$ can be obtained from the equations of state as discussed in Ref. Duan and Qian (2004). We calculate the rates $`\lambda _{eN}^{(1,B)}`$ for a number of values of $`B`$ ($`4\times 10^{15}`$$`1.6\times 10^{16}`$ G for case I and $`2\times 10^{15}`$$`10^{16}`$ G for case II) and show the results as the filled circles with error bars in Fig. 3. The corresponding rates $`\lambda _{eN}^{(0,B)}`$ as functions of $`B`$ and the zero-field results $`\lambda _{eN}^{(0)}`$ and $`\lambda _{eN}^{(1)}`$ are shown as the dashed curves and dotted and dot-dashed lines, respectively, in the same figure. As can be seen from Fig. 3, in the limit of small $`B`$, the dashed curves for $`\lambda _{eN}^{(0,B)}`$ agree with the dotted curves for $`\lambda _{eN}^{(0)}`$. The approach of $`\lambda _{e^+n}^{(1,B)}`$ (filled circles) to the zero-field limit $`\lambda _{e^+n}^{(1)}`$ (dot-dashed line) is also clearly demonstrated for case I. As a large number of Landau levels must be included in the calculation for small $`B`$, it becomes computationally prohibitive to demonstrate the behavior of $`\lambda _{eN}^{(1,B)}`$ in this limit to the fullest extent. Nevertheless, the relation between $`\lambda _{eN}^{(1,B)}`$ and $`\lambda _{eN}^{(0,B)}`$ for small $`B`$ clearly agrees with that between $`\lambda _{eN}^{(1)}`$ and $`\lambda _{eN}^{(0)}`$.
The dependences on $`B`$ for the rates of neutrino emission shown in Fig. 3 require discussion. The effects of the magnetic field on these rates have been noted for the specific case of $`B=10^{16}`$ G, $`T=2`$ MeV, and $`\eta _e=0`$ shown in Fig. 2. More generally, Fig. 3 shows that $`\lambda _{e^{}p}^{(0,B)}`$ and $`\lambda _{e^{}p}^{(1,B)}`$ first decrease with increasing $`B`$ to reach some minimum values and then increase with $`B`$. In contrast, $`\lambda _{e^+n}^{(0,B)}`$ and $`\lambda _{e^+n}^{(1,B)}`$ appear to increase monotonically with $`B`$. The above results can be understood by considering two different effects of the magnetic field on $`e^\pm `$. On the one hand, a stronger field confines more $`e^\pm `$ to the ground Landau level, thus reducing the average $`e^\pm `$ energy. This tends to decrease the rates of neutrino emission. On the other hand, a magnetic field changes the $`e^\pm `$ phase space according to
$$2\frac{\text{d}^3k_e}{(2\pi )^3}\frac{eB}{4\pi ^2}\underset{n_e}{}g_{n_e}_{\mathrm{}}^+\mathrm{}\text{d}k_{ez}.$$
(105)
Thus, the $`e^\pm `$ phase space increases with $`B`$, which tends to increase the rates of neutrino emission due to the increase in the number density of $`e^\pm `$. The competition between the above two factors then determines the dependences on $`B`$ for the rates of neutrino emission.
To show quantitatively the two effects of the magnetic field on $`e^\pm `$ discussed above, we compare the average energy $`E_e_B`$ and the number density $`(\rho _e)_B`$ of $`e^\pm `$ in a field with the corresponding quantities for no field, $`E_e`$ and $`\rho _e`$, respectively, for a wide range of $`B`$ in Fig. 4. As $`0<\eta _e1`$ for the supernova conditions represented by cases I and II, we take $`\eta _e=0`$ for simplicity. The major difference between these two cases lies in the temperature. The ratios $`E_e_B/E_e`$ as functions of $`B`$ for $`T=1`$ and 2 MeV (cases II and I) are shown as the solid and dashed curves, respectively, in Fig. 4a. The corresponding ratios $`(\rho _e)_B/\rho _e`$ are shown in Fig. 4b. For large $`B`$, it is appropriate to consider the limiting case where all $`e^\pm `$ are in the ground Landau level, and therefore, $`E_e_B/E_e`$ is a constant and $`(\rho _e)_B/\rho _e`$ increases linearly with $`B`$. These limits are shown as the dot-dashed and dotted lines for $`T=1`$ and 2 MeV, respectively, in Fig. 4. As can be seen from this figure, $`E_e_B/E_e`$ monotonically decreases with increasing $`B`$, eventually approaching the constant limit, while $`(\rho _e)_B/\rho _e`$ monotonically increases with $`B`$, eventually approaching the limiting linear trend. The combined result of the two effects is that $`\lambda _{eN}`$ decreases with increasing $`B`$ in weak field regime, and starts to increase in strong field regime after some turn-over point. From dimensional analysis, we expect the field at the turn-over point to be $`B_\text{c}E_{\text{eff}}^2/e`$ with $`E_{\text{eff}}`$ being some typical energy of the particles participating in the reaction. Because of the threshold, $`e^{}`$ participating in $`e^{}+p\nu _e+n`$ is more energetic than $`e^+`$ in $`e^++n\overline{\nu }_e+p`$. So $`B_\text{c}`$ is larger for panels (a) and (c) than for panels (b) and (d), respectively, in Fig. 3. The turn-over points correspond to $`B_\mathrm{c}2\times 10^{15}\mathrm{G}`$ in panel (b) and $`B_\mathrm{c}<10^{15}\mathrm{G}`$ in panel (d). However, the turn-over in these two panels is much weaker than that in panels (a) and (c) so that $`\lambda _{e^+n}^{(0,B)}`$ and $`\lambda _{e^+n}^{(1,B)}`$ appear to increase monotonically with $`B`$ for $`B10^{15}\mathrm{G}`$. In addition, because $`E_{\text{eff}}`$ is higher for higher $`T`$, $`B_\text{c}`$ is larger for panels (a) and (b) ($`T=2\mathrm{MeV}`$) than for panels (c) and (d) ($`T=1\mathrm{MeV}`$) in Fig. 3.
In summary, we note that the rates $`\lambda _{eN}`$ are sensitive to the temperature $`T`$ of the supernova environment regardless of $`B`$: lowering $`T`$ by a factor of two reduces $`\lambda _{e^{}p}`$ by factors of $`30`$–50 and $`\lambda _{e^+n}`$ by factors of $`6`$–20 (see Fig. 3). In contrast, the average cross sections $`\sigma _{\nu N}`$ only have minor dependence on $`T`$ (mainly through the polarization of nucleon spin) so that the rates $`\lambda _{\nu N}`$ essentially scale with the radius $`R`$ as $`\lambda _{\nu N}R^2`$. For the temperature profile in the supernova environment of interest, the rates $`\lambda _{\nu N}`$ dominate $`\lambda _{eN}`$. Therefore, so long as the former rates are calculated accurately, the latter can be estimated using $`\lambda _{eN}^{(0,B)}`$ to good approximation for $`B10^{16}`$ G.
## V Conclusions
In a previous paper Duan and Qian (2004), we calculated the rates of $`\nu _e+ne^{}+p`$ and $`\overline{\nu }_e+pe^++n`$ in supernova environments with strong magnetic fields assuming that the nucleon mass $`m_N`$ is infinite. We also applied these rates to discuss the implications of such fields for supernova dynamics. In the present paper, we have taken into account the effects of a finite $`m_N`$ and developed a numerical method for calculating the above rates to $`𝒪(1/m_N)`$ for similar environments. Rates with such an accuracy are required for application to supernova nucleosynthesis.
We have shown that our results have the correct behavior in the limit of high neutrino energy or small magnetic field. We find that for typical supernova $`\nu _e`$ energy distributions, magnetic fields of $`B10^{16}`$ G reduce the rate of $`\nu _e+ne^{}+p`$ while the $`𝒪(1/m_N)`$ corrections due to weak magnetism and nucleon recoil increase this rate. These two opposite effects tend to cancel. On the other hand, the reduction of the rate of $`\overline{\nu }_e+pe^++n`$ by the $`𝒪(1/m_N)`$ corrections dominates the magnetic field effects for $`B10^{16}`$ G and typical supernova $`\overline{\nu }_e`$ energy distributions. We also find that for typical supernova conditions relevant for heavy element nucleosynthesis, the rates of $`e^{}+p\nu _e+n`$ and $`e^++n\overline{\nu }_e+p`$ first decrease and then increase with increasing $`B`$. As it is extremely time consuming to numerically calculate to $`𝒪(1/m_N)`$ the rates for the above processes in strong magnetic fields, we recommend that for $`B10^{16}`$ G, the following approximations be implemented in models of supernova nucleosynthesis. For $`\nu _e+ne^{}+p`$, it is simple to calculate the average cross section including the magnetic field effects but no $`𝒪(1/m_N)`$ corrections \[$`\sigma _{\nu _en}^{(0,B)}`$ in Table 1\] or vice versa \[$`\sigma _{\nu _en}^{(1)}`$ in Table 1\]. By comparing the two with $`\sigma _{\nu _en}^{(0)}`$, one can estimate the effects of magnetic fields and the $`𝒪(1/m_N)`$ corrections, respectively. With these two kinds of corrections combined, $`\sigma _{\nu _en}^{(1)}+\left[\sigma _{\nu _en}^{(0,B)}\sigma _{\nu _en}^{(0)}\right]`$ agrees with the result of the full calculation at the level of $`1`$%. For $`\overline{\nu }_e+pe^++n`$, the magnetic field effects on the average cross section can be ignored but the $`𝒪(1/m_N)`$ corrections should be included with an exact treatment of the reaction kinematics \[$`\sigma _{\overline{\nu }_ep}^{(1)}`$ in Table 1\]. While we have ignored radiative corrections and the effect of the Coulomb field of the proton on the electron wave function, these factors can be included in $`\sigma _{\nu _en}^{(1)}`$ and $`\sigma _{\overline{\nu }_ep}^{(1)}`$ following Ref. Strumia and Vissani (2003). For $`e^{}+p\nu _e+n`$ and $`e^++n\overline{\nu }_e+p`$, the rates including the magnetic field effects but no $`𝒪(1/m_N)`$ corrections \[Eqs. (101)–(103)\] are sufficient.
In conclusion, we note that the cross sections of neutrino absorption on nucleons are relevant not only for supernova nucleosynthesis but also for determining the thermal decoupling of $`\nu _e`$ and $`\overline{\nu }_e`$ from the protoneutron star. For example, a decrease in $`\sigma _{\overline{\nu }_ep}`$ would enable $`\overline{\nu }_e`$ to emerge from deeper and hotter regions of the protoneutron star, thus increasing the average $`\overline{\nu }_e`$ energy. Accurate neutrino energy spectra are essential to models of supernova nucleosynthesis Qian and Woosley (1996). However, our results on the cross sections for neutrino absorption cannot be applied directly to the discussion of neutrino decoupling from a strongly-magnetized protoneutron star because the conditions (e.g., temperature and density) inside such a star are very different from those considered here. We also note that neutrino scattering on $`e^\pm `$ plays a significant role in supernova explosion Bethe (1990). Similar to the case of $`e^\pm `$ capture on nucleons, the rates of neutrino scattering on $`e^\pm `$ above the protoneutron star will be modified substantially by strong magnetic fields. These issues remain to be explored in detail by future studies.
###### Acknowledgements.
We would like to thank Arkady Vainshtein for helpful discussions. We also want to thank John Beacom for communications regarding the effects of weak magnetism and nucleon recoil on neutrino processes. This work was supported in part by DOE grant DE-FG02-87ER40328.
## Appendix A Special Function
The special function $`I_{n,r}(\zeta )`$ can be written as (Arras and Lai, 1999)
$$I_{n,r}(\zeta )=\sqrt{\frac{r!}{n!}}e^{\zeta /2}\zeta ^{(nr)/2}L_r^{nr}(\zeta ),$$
(106)
where $`L_n^\alpha (x)`$ is the generalized Laguerre polynomial defined as (Gradshteyn and Ryzhik, 1980)
$`L_n^\alpha (x)`$ $`=`$ $`{\displaystyle \frac{1}{n!}}e^xx^\alpha {\displaystyle \frac{\text{d}^n}{\text{d}x^n}}(e^xx^{n+\alpha })`$ (107a)
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}(1)^m\left(\begin{array}{c}n+\alpha \\ nm\end{array}\right){\displaystyle \frac{x^m}{m!}}.`$ (107d)
To calculate $`I_{n,r}(\zeta )`$ efficiently, we use its properties given below.
* *Mirror relation*
Based on the identity (Sokolov and Temov, 1968)
$$(1)^{nr}\zeta ^{(nr)}Q_n^{rn}(\zeta )=Q_r^{nr}(\zeta ),$$
(108)
where
$$Q_r^{nr}(\zeta )r!L_r^{nr}(\zeta ),$$
(109)
it is straightforward to show that
$$I_{n,r}(\zeta )=(1)^{nr}I_{r,n}(\zeta ).$$
(110)
* *Recursion relation*s
Using the recursion relation of the generalized Laguerre polynomial (Gradshteyn and Ryzhik, 1980)
$$L_n^{\alpha 1}(\zeta )=L_n^\alpha (\zeta )L_{n1}^\alpha (\zeta ),$$
(111)
one can show that
$`I_{n,r}(\zeta )`$ $`=`$ $`\sqrt{{\displaystyle \frac{r!}{n!}}}e^{\zeta /2}\zeta ^{(nr)/2}L_r^{nr}(\zeta )\left[L_r^{(n1)r}(\zeta )+L_{r1}^{(n1)(r1)}(\zeta )\right]`$ (112)
$`=`$ $`\sqrt{{\displaystyle \frac{\zeta }{n}}}I_{n1,r}(\zeta )+\sqrt{{\displaystyle \frac{r}{n}}}I_{n1,r1}(\zeta ).`$
Using this recursion relation and the mirror relation in Eq. (110), one can prove another recursion relation
$$I_{n,r}(\zeta )=\sqrt{\frac{\zeta }{r}}I_{n,r1}(\zeta )+\sqrt{\frac{n}{r}}I_{n1,r1}(\zeta ).$$
(113)
Starting from the definition
$$I_{0,0}(\zeta )=e^{\zeta /2},$$
(114)
one can use the recursion relation in Eq. (112) to obtain
$$I_{n,0}(\zeta )=\sqrt{\frac{\zeta ^n}{n!}}I_{0,0}(\zeta ).$$
(115)
Using the above result and the mirror relation in Eq. (110), one has
$$I_{0,r}(\zeta )=(1)^r\sqrt{\frac{\zeta ^n}{r!}}I_{0,0}(\zeta ).$$
(116)
The function $`I_{n,r}(\zeta )`$ with $`n>0`$ and $`r>0`$ can be calculated as follows:
1. Compute $`I_{n1,0}(\zeta )`$ and $`I_{n,0}(\zeta )`$ from Eq. (115). Set $`r^{}=1`$.
2. Compute $`I_{n,r^{}}(\zeta )`$ from the recursion relation in Eq. (113).
3. If $`r^{}=r`$, finish. Otherwise, compute $`I_{n1,r^{}}(\zeta )`$ from the recursion relation in Eq. (112).
4. Advance $`r^{}`$ by unity and return to step 2.
## Appendix B Reduced Amplitude squared
The reduced amplitude squared $`𝒲_{\nu _en}`$ for $`\nu _e+ne^{}+p`$ is defined in Eq. (71). The amplitude $`𝔐_{\nu _en}`$ \[Eq. (65)\] contained in $`𝒲_{\nu _en}`$ can be simplified using
$`{\displaystyle _0^{\mathrm{}}}\xi \text{d}\xi {\displaystyle _0^{2\pi }}e^{i𝐰_{}𝐱_{}i(n_er_en_p+r_p)\varphi }I_{n_p,r_p}(eB\xi ^2/2)I_{n_e,r_e}(eB\xi ^2/2)\text{d}\varphi `$ (117)
$`=`$ $`{\displaystyle \frac{2\pi }{eB}}i^{(n_er_en_p+r_p)}e^{i(n_er_en_p+r_p)\varphi _w}I_{n_e,n_p}(w_{}^2/2eB)I_{r_e,r_p}(w_{}^2/2eB),`$
where $`\varphi _w`$ is the azimuthal angle of $`𝐰_{}`$. The above result follows from Sokolov and Temov (1968); Gradshteyn and Ryzhik (1980); Duan and Qian (2004)
$$_0^{2\pi }e^{i𝐰_{}𝐱_{}i(nr)\varphi }\text{d}\varphi =2\pi i^{nr}e^{i(nr)\varphi _w}J_{nr}(w_{}\xi )$$
(118)
and
$$_0^{\mathrm{}}J_{(nr)(n^{}r^{})}(2\sqrt{u\zeta })I_{n^{},r^{}}(u)I_{n,r}(u)\text{d}u=I_{n,n^{}}(\zeta )I_{r,r^{}}(\zeta ),$$
(119)
where $`J_n(\zeta )`$ is the Bessel function.
Noting that Sokolov and Temov (1968)
$$\underset{r}{}I_{n,r}(\zeta )I_{n^{},r}(\zeta )=\delta _{n,n^{}}$$
(120)
and using Eqs. (27) and (117), we are able to derive the following explicit expressions of $`𝒲_{\nu _en}`$ with the help of Mathmatica<sup>®</sup>:
$`(𝒲_{\nu _en})_{s_p=1,s_n=1}`$ $`=`$ $`(f+g)^2(1+v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e,n_p}^2(w_{}^2/2eB)`$ (121a)
$`+(fg)^2(1v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e1,n_p}^2(w_{}^2/2eB)`$
$`+2(f^2g^2){\displaystyle \frac{\sqrt{2n_eeB}}{E_e}}\mathrm{cos}\varphi _w\mathrm{sin}\mathrm{\Theta }_\nu I_{n_e1,n_p}(w_{}^2/2eB)I_{n_e,n_p}(w_{}^2/2eB)`$
$`+{\displaystyle \frac{1}{m_N}}\{[(f+g)^2(1+v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )(k_{nz}+k_{pz})`$
$`(f+g)(2f+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`+f_2(f+g)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_x`$
$`+(f+g)(2f+f_2)(1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e,n_p}^2(w_{}^2/2eB)`$
$`+[(fg)^2(1v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )(k_{nz}+k_{pz})`$
$`+f_2(fg)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`(fg)(2f+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_x`$
$`f_2(fg)(1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e1,n_p}^2(w_{}^2/2eB)`$
$`+[f_2(f+g)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`+(fg)(2f+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`+2\left(f^2+g(f+f_2)+f(fg+f_2)\mathrm{cos}\mathrm{\Theta }_\nu \right)`$
$`\times \mathrm{cos}(\mathrm{\Phi }_n\varphi _w){\displaystyle \frac{\sqrt{2n_eeB}k_n}{E_e}}`$
$`(f+g)(2f+f_2)(1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}w_{}}{E_e}}`$
$`+f_2(fg)(1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}w_{}}{E_e}}]`$
$`\times I_{n_e1,n_p}(w_{}^2/2eB)I_{n_e,n_p}(w_{}^2/2eB)\},`$
$`(𝒲_{\nu _en})_{s_p=1,s_n=1}`$ $`=`$ $`4g^2(1+v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e,n_p}^2(w_{}^2/2eB)`$ (121b)
$`+{\displaystyle \frac{2}{m_N}}\{[g(f+g+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`+2g(f+f_2)(1+v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )(k_{ez}k_{\nu z})`$
$`+g(f+g+f_2)(1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e,n_p}^2(w_{}^2/2eB)`$
$`+g(fg+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_xI_{n_e,n_p1}^2(w_{}^2/2eB)`$
$`+\left[g(fg+f_2)k_n\mathrm{cos}(\mathrm{\Phi }_n\varphi _w)g(f+g+f_2)w_{}\right]`$
$`\times (1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}}{E_e}}I_{n_e1,n_p}(w_{}^2/2eB)I_{n_e,n_p}(w_{}^2/2eB)`$
$`g(fg+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`\times I_{n_e1,n_p1}(w_{}^2/2eB)I_{n_e,n_p1}(w_{}^2/2eB)\},`$
$`(𝒲_{\nu _en})_{s_p=1,s_n=1}`$ $`=`$ $`4g^2(1v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e1,n_p1}^2(w_{}^2/2eB)`$ (121c)
$`+{\displaystyle \frac{2}{m_N}}\{[g(f+g+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`2g(f+f_2)(1v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )(k_{ez}k_{\nu z})`$
$`+g(f+g+f_2)(1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e1,n_p1}^2(w_{}^2/2eB)`$
$`+g(fg+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_xI_{n_e1,n_p}^2(w_{}^2/2eB)`$
$`+\left[g(fg+f_2)k_n\mathrm{cos}(\mathrm{\Phi }_n\varphi _w)g(f+g+f_2)w_{}\right]`$
$`\times (1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}}{E_e}}I_{n_e1,n_p1}(w_{}^2/2eB)I_{n_e,n_p1}(w_{}^2/2eB)`$
$`g(fg+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`\times I_{n_e1,n_p}(w_{}^2/2eB)I_{n_e,n_p}(w_{}^2/2eB)\},`$
$`(𝒲_{\nu _en})_{s_p=1,s_n=1}`$ $`=`$ $`(f+g)^2(1v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e1,n_p1}^2(w_{}^2/2eB)`$ (121d)
$`+(fg)^2(1+v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )I_{n_e,n_p1}^2(w_{}^2/2eB)`$
$`+2(f^2g^2){\displaystyle \frac{\sqrt{2n_eeB}}{E_e}}\mathrm{cos}\varphi _w\mathrm{sin}\mathrm{\Theta }_\nu I_{n_e1,n_p1}(w_{}^2/2eB)I_{n_e,n_p1}(w_{}^2/2eB)`$
$`+{\displaystyle \frac{1}{m_N}}\{[(f+g)^2(1v_{ez})(1\mathrm{cos}\mathrm{\Theta }_\nu )(k_{nz}+k_{pz})`$
$`(f+g)(2f+f_2)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`+f_2(f+g)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_x`$
$`+(f+g)(2f+f_2)(1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e1,n_p1}^2(w_{}^2/2eB)`$
$`+[(fg)^2(1+v_{ez})(1+\mathrm{cos}\mathrm{\Theta }_\nu )(k_{nz}+k_{pz})`$
$`+f_2(fg)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu k_{nx}`$
$`(fg)(2f+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu w_x`$
$`f_2(fg)(1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{2n_eeB}{E_e}}]\times I_{n_e,n_p1}^2(w_{}^2/2eB)`$
$`+[f_2(f+g)(1v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`+(fg)(2f+f_2)(1+v_{ez})\mathrm{sin}\mathrm{\Theta }_\nu \mathrm{cos}\varphi _w\sqrt{2n_eeB}`$
$`2\left(f^2g(f+f_2)+f(fg+f_2)\mathrm{cos}\mathrm{\Theta }_\nu \right)`$
$`\times \mathrm{cos}(\mathrm{\Phi }_n\varphi _w){\displaystyle \frac{\sqrt{2n_eeB}k_n}{E_e}}`$
$`(f+g)(2f+f_2)(1\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}w_{}}{E_e}}`$
$`+f_2(fg)(1+\mathrm{cos}\mathrm{\Theta }_\nu ){\displaystyle \frac{\sqrt{2n_eeB}w_{}}{E_e}}]`$
$`\times I_{n_e1,n_p1}(w_{}^2/2eB)I_{n_e,n_p1}(w_{}^2/2eB)\}.`$
In the above equations, $`v_{ez}=k_{ez}/E_e`$. The reduced amplitude squared $`𝒲_{\overline{\nu }_ep}`$ for $`\overline{\nu }_e+pe^++n`$ can be obtained from these equations by making the substitution given in Eq. (81).
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# Defensive forecasting for linear protocols
## 1 Introduction
In we suggested a new methodology for designing forecasting strategies. Considering only the simplest case of binary forecasting, we showed that any constructive, in the sense explained below, law of probability can be translated into a forecasting strategy that satisfies this law. In this paper this result is extended to a general class of protocols including multi-class forecasting. In proposing this approach to forecasting we were inspired by and papers further developing , although our methods and formal results appear to be completely different.
Whereas the meta-theorem stated in is mathematically trivial, the generalization considered in this paper is less so, depending on the Schauder-Tikhonov fixed-point theorem. Our general meta-theorem is stated in §4 and proved in §4 and Appendix A. The general forecasting protocols covered by this result are introduced and discussed in §§23.
In we demonstrated the value of the meta-theorem by applying it to the strong law of large numbers, obtaining from it a kernel forecasting strategy which we called K29. The derivation, however, was informal, involving heuristic transitions to a limit, and this made it impossible to state formally any properties of K29. In this paper we deduce K29 in a much more direct way from the weak law of large numbers and state its properties. (For binary forecasting, this was also done in , and the reader might prefer to read that paper first.) The weak law of large numbers is stated and proved in §5, and K29 is derived and studied in §6.
We call the approach to forecasting using our meta-theorem “defensive forecasting”: Forecaster is trying to defend himself when playing against Skeptic. The justification of this approach given in this paper and in is K29’s properties of proper calibration and resolution. Another justification, in a sense the ultimate justification of any forecasts, is given in : defensive forecasts lead to good decisions; this result, however, is obtained in for rather simple decision problems requiring only binary forecasts, and its extensions will require this paper’s results or their generalizations.
The exposition of probability theory needed for this paper is given in . The standard exposition is based on Kolmogorov’s measure-theoretic axioms of probability, whereas states several key laws of probability in terms of a game between the forecaster, the reality, and a third player, the skeptic. The game-theoretic laws of probability in are constructive in that we explicitly construct computable winning strategies for the forecaster in various games of forecasting.
## 2 Forecasting as a game
Following and we consider the following general forecasting protocol:
Forecasting Game 1
Players: Reality, Forecaster, Skeptic
Parameters: $`𝐗`$ (*data space*), $`𝐘`$ (*observation space*), $`𝐅`$ (*Forecaster’s move space*), $`𝐒`$ (*Skeptic’s move space*), $`\lambda :𝐒\times 𝐅\times 𝐘`$ (*Skeptic’s gain function* and *Forecaster’s loss function*)
Protocol:
$`𝒦_0:=1`$.
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`f_n𝐅`$.
Skeptic announces $`s_n𝐒`$.
Reality announces $`y_n𝐘`$.
$`𝒦_n:=𝒦_{n1}+\lambda (s_n,f_n,y_n)`$.
END FOR
Restriction on Skeptic: Skeptic must choose the $`s_n`$ so that his capital is always nonnegative ($`𝒦_n0`$ for all $`n`$) no matter how the other players move.
This is a perfect-information protocol: the players move in the order indicated, and each player sees the other player’s moves as they are made. It specifies both an initial value for Skeptic’s capital ($`𝒦_0=1`$) and a lower bound on its subsequent values ($`𝒦_n0`$). We will say that $`x_n`$ are the *data*, $`y_n`$ are the *observations*, and $`f_n`$ are the *forecasts*. In applications, the datum $`x_n`$ will contain all available information deemed useful in forecasting $`y_n`$.
Book contains several results (game-theoretic versions of limit theorems of probability theory) of the following form: Skeptic has a strategy that guarantees that either a property of agreement between the forecasts $`f_n`$ and observations $`y_n`$ is satisfied or Skeptic becomes very rich (without risking bankruptcy, according to the protocol). All specific strategies considered in have computable versions. According to Brouwer’s principle (see, e.g., §1 of for a recent review of the relevant literature) they must be automatically continuous; in any case, their continuity can be checked directly. In we showed that, under a special choice of the players’ move spaces and Skeptic’s gain function $`\lambda `$, for any continuous strategy for Skeptic Forecaster has a strategy that guarantees that Skeptic’s capital never increases when he plays that strategy. Therefore, Forecaster has strategies that ensure various properties of agreement between the forecasts and the observations.
The purpose of this paper is to extend the result of to a wide class of Skeptic’s gain functions $`\lambda `$. But first we consider several important special cases of Forecasting Game 1.
### Binary forecasting
The simplest non-trivial case, considered in , is where $`𝐘=\{0,1\}`$, $`𝐅=[0,1]`$, $`𝐒=`$, and
$$\lambda (s_n,f_n,y_n)=s_n(y_nf_n).$$
(1)
Intuitively, Forecaster gives probability forecasts for $`y_n`$: $`f_n`$ is his subjective probability that $`y_n=1`$. The operational interpretation of $`f_n`$ is that it is the price that Forecaster charges for a ticket that will pay $`y_n`$ at the end of the $`n`$th round of the game; $`s_n`$ is the number (positive, zero, or negative) of such tickets that Skeptic chooses to buy.
### Bounded regression
This is the most straightforward extension of binary forecasting, considered in , §3.2. The move spaces are $`𝐘=𝐅=[A,B]`$, where $`A`$ and $`B`$ are two constants, and $`𝐒=`$; the gain function is, as before, (1). This protocol allows one to prove a strong law of large numbers (, Proposition 3.3) and a simple one-sided law of the iterated logarithm (, Corollary 5.1).
### Multi-class forecasting
Another extension of binary forecasting is the protocol where $`𝐘`$ is a finite set, $`𝐅`$ is the set of all probability distributions on $`𝐘`$, $`𝐒`$ is the set of all real-valued functions on $`𝐘`$, and
$$\lambda (s_n,f_n,y_n)=s_n(y_n)s_ndf_n.$$
The intuition behind Skeptic’s move $`s_n`$ is that Skeptic buys the ticket which pays $`s_n(y_n)`$ after $`y_n`$ is announced; he is charged $`s_ndf_n`$ for this ticket.
The binary forecasting protocol is “isomorphic” to the special case of this protocol where $`𝐘=\{0,1\}`$: Forecaster’s move $`f_n`$ in the binary forecasting protocol is represented by the probability distribution $`f_n^{}`$ on $`\{0,1\}`$ assigning weight $`f_n`$ to $`\{1\}`$ and Skeptic’s move $`s_n`$ in the binary forecasting protocol is represented by any function $`s_n^{}`$ on $`\{0,1\}`$ such that $`s_n^{}(1)s_n^{}(0)=s_n`$. The isomorphism between these two protocols follows from
$$\begin{array}{c}s_n^{}(y_n)s_n^{}df_n^{}=s_n^{}(y_n)s_n^{}(1)f_ns_n^{}(0)(1f_n)\hfill \\ \hfill =s_n^{}(y_n)s_n^{}(0)s_nf_n=s_n(y_nf_n)\end{array}$$
(remember that $`y_n\{0,1\}`$).
### Bounded mean-variance forecasting
In this protocol, $`𝐘=[A,B]`$, where $`A`$ and $`B`$ are again two constants, $`𝐅=𝐒=^2`$, and
$$\lambda (s_n,f_n,y_n)=\lambda ((M_n,V_n),(m_n,v_n),y_n)=M_n(y_nm_n)+V_n((y_nm_n)^2v_n).$$
Intuitively, Forecaster is asked to forecast $`y_n`$ with a number $`m_n`$ and also forecast the accuracy $`(y_nm_n)^2`$ of his first forecast with a number $`v_n`$. This protocol, although usually without the restriction $`y_n[A,B]`$, is used extensively in (e.g., in Chaps. 4 and 5).
An equivalent representation of this protocol is $`𝐘=\{(t,t^2)|t[A,B]\}`$, $`𝐅=𝐒=^2`$ and
$$\lambda (s_n,f_n,y_n)=\lambda ((s_n^{},s_n^{\prime \prime }),(f_n^{},f_n^{\prime \prime }),(t_n,t_n^2))=s_n^{}(t_nf_n^{})+s_n^{\prime \prime }(t_n^2f_n^{\prime \prime }).$$
The equivalence of the two representations can be seen as follows: Reality’s move $`(x_n,t_n)`$ in the first representation corresponds to $`(x_n,y_n)=(x_n,(t_n,t_n^2))`$ in the second representation, Forecaster’s move $`(m_n,v_n)`$ in the first representation corresponds to $`(f_n^{},f_n^{\prime \prime })=(m_n,v_n+m_n^2)`$ in the second representation, and Skeptic’s move $`(s_n^{},s_n^{\prime \prime })`$ in the second representation corresponds to $`(M_n,V_n)=(s_n^{}+2m_ns_n^{\prime \prime },s_n^{\prime \prime })`$ in the first representation. This establishes a bijection between Reality’s move spaces, a bijection between Forecaster’s move spaces, and a bijection between Skeptic’s move spaces in the two representations; Skeptic’s gains are also the same in the two representations:
$$\begin{array}{c}s_n^{}(t_nf_n^{})+s_n^{\prime \prime }(t_n^2f_n^{\prime \prime })\hfill \\ \hfill =s_n^{}(t_nm_n)+s_n^{\prime \prime }\left(\left((t_nm_n)^2+2(t_nm_n)m_n+m_n^2\right)\left(v_n+m_n^2\right)\right)\\ \hfill =(s_n^{}+2m_ns_n^{\prime \prime })(t_nm_n)+s_n^{\prime \prime }\left((t_nm_n)^2v_n\right).\end{array}$$
## 3 Linear protocol
Forecasting Game 1 is too general to derive results of the kind we are interested in. In this subsection we will introduce a narrower protocol which will still be wide enough to cover all special cases considered so far.
All move spaces are now subsets of a Hilbert space $`𝐋`$ (we allow $`𝐋`$ to be non-separable or finite-dimensional; in fact, in this paper we emphasize the case where $`𝐋=^m`$ for some positive integer $`m`$). The observation space is a non-empty pre-compact subset $`𝐘𝐋`$ (we say that a set is *pre-compact* if its closure is compact; if $`𝐋=^m`$, this is equivalent to it being bounded), Forecaster’s move space $`𝐅`$ is the whole of $`𝐋`$, and Skeptic’s move space $`𝐒`$ is also the whole of $`𝐋`$. Skeptic’s gain function is
$$\lambda (s_n,f_n,y_n)=s_n,y_nf_n_𝐋.$$
Therefore, we consider the following perfect-information game:
Forecasting Game 2
Players: Reality, Forecaster, Skeptic
Parameters: $`𝐗`$, $`𝐋`$ (Hilbert space), $`𝐘`$ (non-empty pre-compact subset of $`𝐋`$)
Protocol:
$`𝒦_0:=1`$.
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`f_n𝐋`$.
Skeptic announces $`s_n𝐋`$.
Reality announces $`y_n𝐘`$.
$`𝒦_n:=𝒦_{n1}+s_n,y_nf_n_𝐋`$. (2)
END FOR
Restriction on Skeptic: Skeptic must choose the $`s_n`$ so that his capital is always nonnegative no matter how the other players move.
Let us check that the specific protocols considered in the previous section are covered by this *linear protocol* (and for all those protocols $`𝐋`$ can be taken finite dimensional, $`𝐋=^m`$ for some $`m\{1,2,\mathrm{}\}`$). At first sight, even the binary forecasting protocol is not covered, as Forecaster’s move space is $`𝐅=[0,1]`$ rather than $``$. It is easy to see, however, that Forecaster’s move $`f_n\overline{\mathrm{co}}𝐘`$ outside the convex closure $`\overline{\mathrm{co}}𝐘`$ of the observation space (the convex closure $`\overline{\mathrm{co}}A`$ of a set $`A`$ is defined to be the intersection of all convex closed sets containing $`A`$) is always inadmissible, in the sense that there exists Skeptic’s reply $`s_n`$ making him arbitrarily rich regardless of Reality’s move, and so we can as well choose $`𝐅:=\overline{\mathrm{co}}𝐘`$. Indeed, suppose that $`f_n\overline{\mathrm{co}}𝐘`$ in the linear protocol. Since $`𝐘`$ is pre-compact, $`\overline{\mathrm{co}}𝐘`$ is compact (, Theorem 3.20(c)). By the Hahn-Banach theorem (, Theorem 3.4(b)), there exists a vector $`s_n𝐋`$ such that
$$\underset{y𝐘}{inf}s_n,yf_n_𝐋>0.$$
(It would have been sufficient for either $`\{f_n\}`$ or $`\overline{\mathrm{co}}𝐘`$ to be compact; in fact both are.) Skeptic’s move $`Cs_n`$ can make him as rich as he wishes as $`C`$ can be arbitrarily large. In what follows, we will usually assume that Forecaster’s move space is $`\overline{\mathrm{co}}𝐘`$ and use $`𝐅`$ as a shorthand for $`\overline{\mathrm{co}}𝐘`$.
Now it is obvious that the binary forecasting, bounded regression, and bounded mean-variance forecasting (in its second representation) protocols are special cases of the linear protocol (perhaps with $`𝐅=\overline{\mathrm{co}}𝐘`$). For the multi-class forecasting protocol, we should represent $`𝐘`$ as the vertices
$$y^1:=(1,0,0,\mathrm{},0),y^2:=(0,1,0,\mathrm{},0),\mathrm{},y^m:=(0,0,0,\mathrm{},1)$$
of the standard simplex in $`^m`$, where $`m`$ is the size of $`𝐘`$, represent the probability distributions $`f`$ on $`𝐘`$ as vectors $`(f\{y^1\},\mathrm{},f\{y^m\})`$ in $`^m`$, and represent the real-valued functions $`s`$ on $`𝐘`$ as vectors $`(s(y^1),\mathrm{},s(y^m))`$ in $`^m`$.
## 4 Meta-theorem
In this section we state the main mathematical result of this paper: for any continuous strategy for Skeptic there exists a strategy for Forecaster that does not allow Skeptic’s capital to grow, regardless of what Reality is doing. As in , we make Skeptic announce his strategy for each round at the outset of that round rather than announce his strategy for the whole game at the beginning of the game, and we drop all restrictions on Skeptic. Forecaster’s move space is restricted to $`𝐅=\overline{\mathrm{co}}𝐘`$. The resulting perfect-information game is:
Forecasting Game 3
Players: Reality, Forecaster, Skeptic
Parameters: $`𝐗`$, $`𝐋`$ (Hilbert space), $`𝐘𝐋`$ (non-empty and pre-compact)
Protocol:
$`𝒦_0`$ is set to a real number.
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Skeptic announces continuous $`S_n:\overline{\mathrm{co}}𝐘𝐋`$.
Forecaster announces $`f_n\overline{\mathrm{co}}𝐘`$.
Reality announces $`y_n𝐘`$.
$`𝒦_n:=𝒦_{n1}+S_n(f_n),y_nf_n_𝐋`$.
END FOR
###### Theorem 1
Forecaster has a strategy in Forecasting Game 3 that ensures $`𝒦_0𝒦_1𝒦_2\mathrm{}`$.
* Fix a round $`n`$ and Skeptic’s move $`S_n:𝐅𝐋`$ (we will refer to $`S_n`$ as a vector field in $`𝐅`$). Our task is to prove the existence of a point $`f_n𝐅`$ such that, for all $`y𝐘`$, $`S_n(f_n),yf_n_𝐋0`$.
If for some $`f𝐅`$ (we use $`A`$ to denote the boundary of $`A𝐋`$) the vector $`S_n(f)`$ is normal and directed exteriorly to $`𝐅`$ (in the sense that $`S_n(f),yf_𝐋0`$ for all $`y𝐅`$), we can take such $`f`$ as $`f_n`$. Therefore, we assume, without loss of generality, that $`S_n`$ is never normal and directed exteriorly on $`𝐅`$. Then by Lemma 1 in Appendix A there exists $`f`$ such that $`S_n(f)=0`$, and we can take such $`f`$ as $`f_n`$.
* Notice that Theorem 1 will not become weaker if the first move by Reality (choosing $`x_n`$) is removed from each round of the protocol.
## 5 A weak law of large numbers in Hilbert space
Unfortunately, the usual law of large numbers is not useful for the purpose of designing forecasting strategies (see the discussion in ). Therefore, we state a generalized law of large numbers; at the end of this section we will explain connections with the usual law of large numbers. In this section we consider Forecasting Game 2 without the requirement $`𝒦_0=1`$ and with the restriction on Skeptic dropped. If we fix a strategy for Skeptic and Skeptic’s initial capital $`𝒦_0`$ (not necessarily $`1`$ or even a positive number), $`𝒦_n`$ defined by (3) becomes a function of Reality’s and Forecaster’s moves. Such functions will be called *capital processes*.
Let $`\mathrm{\Phi }:𝐅\times 𝐗𝐇`$ (as usual, $`𝐅=\overline{\mathrm{co}}𝐘`$) be a *feature mapping* into a Hilbert space $`𝐇`$; $`𝐇`$ is called the *feature space*. The next theorem uses the notion of tensor product; for details, see Appendix B.
###### Theorem 2
The function
$$𝒦_n:=\underset{i=1}{\overset{n}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2\underset{i=1}{\overset{n}{}}y_if_i_𝐋^2\mathrm{\Phi }(f_i,x_i)_𝐇^2$$
(3)
is a capital process (not necessarily non-negative) of some strategy for Skeptic.
* We start by noticing that
$`𝒦_n𝒦_{n1}`$ $`={\displaystyle \underset{i=1}{\overset{n1}{}}}(y_if_i)\mathrm{\Phi }(f_i,x_i)+(y_nf_n)\mathrm{\Phi }(f_n,x_n)_{𝐋𝐇}^2`$
$`{\displaystyle \underset{i=1}{\overset{n1}{}}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2y_nf_n_𝐋^2\mathrm{\Phi }(f_n,x_n)_𝐇^2`$
$`=2{\displaystyle \underset{i=1}{\overset{n1}{}}}(y_if_i)\mathrm{\Phi }(f_i,x_i),(y_nf_n)\mathrm{\Phi }(f_n,x_n)_{𝐋𝐇}`$
$`=2{\displaystyle \underset{i=1}{\overset{n1}{}}}y_if_i,y_nf_n_𝐋\mathrm{\Phi }(f_i,x_i),\mathrm{\Phi }(f_n,x_n)_𝐇`$
(in the last two equalities we used (18) and (19) from Appendix B). Introducing the notation
$$𝐤((f,x),(f^{},x^{})):=\mathrm{\Phi }(f,x),\mathrm{\Phi }(f^{},x^{})_𝐇,$$
(4)
where $`(f,x),(f^{},x^{})𝐅\times 𝐗`$, we can rewrite the expression for $`𝒦_n𝒦_{n1}`$ as
$$2\underset{i=1}{\overset{n1}{}}𝐤((f_i,x_i),(f_n,x_n))(y_if_i),y_nf_n_𝐋.$$
Therefore, $`𝒦_n`$ is the capital process corresponding to Skeptic’s strategy
$$2\underset{i=1}{\overset{n1}{}}𝐤((f_i,x_i),(f_n,x_n))(y_if_i);$$
(5)
this completes the proof.
### More standard statements of the weak law
In the rest of this section we explain connections of Theorem 2 with more standard statements of the weak law of large numbers; in this part of the paper we will use some notions introduced in . The rest of the paper does not depend on this material, and the reader may wish to skip this subsection.
Let us assume that
$$𝐜_\mathrm{\Phi }:=\underset{(f,x)𝐅\times 𝐗}{sup}\mathrm{\Phi }(f,x)_𝐇<\mathrm{}.$$
We will use the notation $`diam(𝐘):=sup_{y,y^{}𝐘}yy^{}_𝐋`$; it is clear that $`diam(𝐘)<\mathrm{}`$. For any initial capital $`𝒦_0`$,
$$𝒦_n:=𝒦_0+\underset{i=1}{\overset{n}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2\underset{i=1}{\overset{n}{}}y_if_i_𝐋^2\mathrm{\Phi }(f_i,x_i)_𝐇^2$$
is the capital process of some strategy for Skeptic. Suppose a positive integer $`N`$ (the duration of the game, or the *horizon*) is given in advance and $`𝒦_0:=diam^2(𝐘)𝐜_\mathrm{\Phi }^2N`$. Then, in the game lasting $`N`$ rounds, $`𝒦_n`$ is never negative and
$$𝒦_N\underset{i=1}{\overset{N}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2.$$
If we do not believe that Skeptic can increase his capital $`1/\delta `$-fold for a small $`\delta >0`$ without risking bankruptcy, we should believe that
$$\underset{i=1}{\overset{N}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2diam^2(𝐘)𝐜_\mathrm{\Phi }^2N/\delta ,$$
which can be rewritten as
$$\frac{1}{N}\underset{i=1}{\overset{N}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}diam(𝐘)𝐜_\mathrm{\Phi }(N\delta )^{1/2}.$$
(6)
In the terminology of , *the game-theoretic lower probability of the event (6) is at least $`1\delta `$*.
The game-theoretic version of Bernoulli’s law of large numbers is a special case of (6) corresponding to $`\mathrm{\Phi }(f,x)=1`$, for all $`f`$ and $`x`$, $`𝐘=\{0,1\}`$, and $`|𝐗|=1`$ (the last two conditions mean that we are considering the binary forecasting protocol without the data); as usual, we assume that $`f_i`$ are chosen from $`\overline{\mathrm{co}}𝐘=[0,1]`$. As explained in , in combination with the measurability of Skeptic’s strategy guaranteeing (6), this implies that the measure-theoretic probability of the event (6) is at least $`1\delta `$, assuming that the $`y_i`$ are generated by a probability distribution and that each $`f_i`$ is the conditional probability that $`y_i=1`$ given $`y_1,\mathrm{},y_{i1}`$. This measure-theoretic result was proved by Kolmogorov in 1929 (see ) and is the origin of the name “K29 strategy”.
We will see in the next section that the feature-space version (6) of the weak law of large numbers is much more useful than the standard version for the purpose of forecasting.
## 6 The K29 strategy and its properties
According to Theorem 1, under the continuity assumption there is a strategy for Forecaster that does not allow $`𝒦_n`$ to grow, where $`𝒦_n`$ is defined by (3). Fortunately (but not unusually), this strategy depends on the feature mapping $`\mathrm{\Phi }`$ only via the corresponding *kernel* $`𝐤`$ defined by (4). The continuity assumption needed is that $`𝐤((f,x),(f^{},x^{}))`$ should be continuous in $`f`$; such kernels will be called *admissible*. According to (5), the corresponding forecasting strategy, which we will call the *K29 strategy* with parameter $`𝐤`$, is to output, on the $`n`$th round, a forecast $`f_n`$ satisfying
$$S(f_n):=\underset{i=1}{\overset{n1}{}}𝐤((f_i,x_i),(f_n,x_n))(y_if_i)=0$$
(or, if such $`f_n`$ does not exist, the forecast is chosen to be a point $`f_n𝐅`$ where $`S(f_n)`$ is normal and directed exteriorly to $`𝐅`$).
The protocol of this section is essentially that of Forecasting Game 3; as Skeptic ceases to be an active player, it simplifies to:
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`f_n\overline{\mathrm{co}}𝐘`$.
Reality announces $`y_n𝐘`$.
END FOR
###### Theorem 3
The K29 strategy guarantees that always
$$\underset{i=1}{\overset{n}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}diam(𝐘)𝐜_\mathrm{\Phi }\sqrt{n},$$
(7)
where $`𝐜_\mathrm{\Phi }:=sup_{(f,x)𝐅\times 𝐗}\mathrm{\Phi }(f,x)_𝐇`$ is assumed to be finite.
* The K29 strategy ensures that (3) never increases; therefore,
$$\underset{i=1}{\overset{n}{}}(y_if_i)\mathrm{\Phi }(f_i,x_i)_{𝐋𝐇}^2\underset{i=1}{\overset{n}{}}y_if_i_𝐋^2\mathrm{\Phi }(f_i,x_i)_𝐇^2diam^2(𝐘)𝐜_\mathrm{\Phi }^2n.$$
* The property (7) is a special case of (6) corresponding to $`\delta =1`$; we gave an independent derivation to make our exposition self-contained and to avoid the extra assumptions used in the derivation of (6), such as the horizon being finite and known in advance.
### K29 with reproducing kernel Hilbert spaces
A *reproducing kernel Hilbert space* (usually abbreviated to RKHS) is a function space $``$ on some set $`Z`$ such that all evaluation functionals $`FF(z)`$, $`zZ`$, are continuous. We will be interested in RKHS on the Cartesian product $`𝐅\times 𝐗`$.
By the Riesz-Fischer theorem, for each $`zZ`$ there exists a function $`𝐤_z`$ such that
$$F(z)=𝐤_z,F_{},F.$$
Let
$$𝐜_{}:=\underset{zZ}{sup}𝐤_z_{};$$
(8)
we will be interested in the case $`𝐜_{}<\mathrm{}`$.
The *kernel* of an RKHS $``$ on $`Z`$ is
$$𝐤(z,z^{}):=𝐤_z,𝐤_z^{}_{}$$
(9)
(equivalently, we could define $`𝐤(z,z^{})`$ as $`𝐤_z(z^{})`$ or as $`𝐤_z^{}(z)`$). It is clear that (9) is a special case of the generalization
$$𝐤(z,z^{}):=\mathrm{\Phi }(z),\mathrm{\Phi }(z^{})_𝐇$$
(10)
of (4). In fact, the functions $`𝐤`$ that can be represented as (10) are exactly the functions that can be represented as (9); they can be equivalently defined as symmetric positive definite functions on $`Z^2`$ (see for a list of references).
A long list of RKHS together with their kernels is given in , §7.4. We will only give one example: the Sobolev space $`𝒮`$ of absolutely continuous functions $`F`$ on $``$ with finite norm
$$F_𝒮:=\sqrt{_{\mathrm{}}^{\mathrm{}}F^2(z)dz+_{\mathrm{}}^{\mathrm{}}(F^{}(z))^2dz};$$
(11)
its kernel is
$$𝐤(z,z^{})=\frac{1}{2}\mathrm{exp}\left(\left|zz^{}\right|\right)$$
(see or , §7.4, Example 24). From the last equation we can see that $`𝐜_𝒮=1/\sqrt{2}`$.
The following is an easy corollary of Theorem 3.
###### Theorem 4
Let $``$ be an RKHS on $`𝐅\times 𝐗`$. The K29 strategy with parameter $`𝐤`$ (defined by (9)) ensures
$$\underset{i=1}{\overset{n}{}}F(f_i,x_i)(y_if_i)_𝐋diam(𝐘)𝐜_{}F_{}\sqrt{n}$$
(12)
for each function $`F`$, where $`𝐜_{}`$ is defined by (8).
* Let $`\mathrm{\Phi }:𝐅\times 𝐗𝐇:=`$ be defined by $`\mathrm{\Phi }(z):=𝐤_z`$. Theorem 3 then implies
$`{\displaystyle \underset{i=1}{\overset{n}{}}}F(f_i,x_i)(y_if_i)_𝐋`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}𝐤_{f_i,x_i},F_𝐇(y_if_i)_𝐋`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}\left((y_if_i)𝐤_{f_i,x_i}\right)F_𝐋`$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}(y_if_i)𝐤_{f_i,x_i}_𝐋F_{}`$
$`diam(𝐘)𝐜_{}F_{}\sqrt{n}`$
(the second equality follows from Lemma 2 and the first inequality from Lemma 3 in Appendix B).
### Calibration and resolution
Two important properties of a forecasting strategy are its calibration and resolution, which we introduce informally. Our discussion in this section extends the discussion in , §5, to the case of linear protocols (in particular, to the case of multi-class forecasting). Forecaster’s move space is assumed to be $`𝐅=\overline{\mathrm{co}}𝐘`$.
We say that the forecasts $`f_n`$ are *properly calibrated* if, for any $`f^{}𝐅`$,
$$\frac{_{i=1,\mathrm{},n:f_if^{}}y_i}{_{i=1,\mathrm{},n:f_if^{}}1}f^{}$$
provided $`_{i=1,\mathrm{},n:f_if^{}}1`$ is not too small. (We shorten $`(1/c)v`$ to $`v/c`$, where $`v`$ is a vector and $`c0`$ is a number.) Proper calibration is only a necessary but far from sufficient condition for good forecasts: for example, a forecaster who ignores the data $`x_n`$ can be perfectly calibrated, no matter how much useful information $`x_n`$ contain. (Cf. the discussion in .)
We say that the forecasts $`f_n`$ are *properly calibrated and resolved* if, for any $`(f^{},x^{})𝐅\times 𝐗`$,
$$\frac{_{i=1,\mathrm{},n:(f_i,x_i)(f^{},x^{})}y_i}{_{i=1,\mathrm{},n:(f_i,x_i)(f^{},x^{})}1}f^{}$$
(13)
provided $`_{i=1,\mathrm{},n:(f_i,x_i)(f^{},x^{})}1`$ is not too small.
Instead of “crisp” points $`(f^{},x^{})𝐅\times 𝐗`$ one may consider “fuzzy points” $`I:𝐅\times 𝐗[0,1]`$ such that $`I(f^{},x^{})=1`$ and $`I(f,x)=0`$ for all $`(f,x)`$ outside a small neighborhood of $`(f^{},x^{})`$. A standard choice would be something like $`I:=𝕀_E`$, where $`E𝐅\times 𝐗`$ is a small neighborhood of $`(f^{},x^{})`$ and $`𝕀_E`$ is its indicator function, but we will want $`I`$ to be continuous (it can, however, be arbitrarily close to $`𝕀_E`$).
Suppose $`𝐅^m`$ and $`𝐗^l`$ for some $`m,l\{1,2,\mathrm{}\}`$. Let $`(f^{},x^{})`$ be a point in $`𝐅\times 𝐗`$; consider a small box $`E:=_{i=1}^m[a_i,b_i]\times _{j=1}^l[c_j,d_j]`$ containing this point, $`E(f^{},x^{})`$. The indicator $`𝕀_E`$ of $`E`$ can be arbitrarily well approximated by the tensor product
$$I(f_1,\mathrm{},f_m,x_1,\mathrm{},x_l)=\underset{i=1}{\overset{m}{}}F_i(f_i)\underset{j=1}{\overset{l}{}}G_j(x_j)$$
of some functions $`F_i`$ and $`G_j`$ from the Sobolev class (11). Let $`I_{}`$ be the norm of $`I`$ in the tensor product $``$ of $`m+l`$ copies of $`𝒮`$ (see , §I.8, for an explicit description of tensor products of RKHS). We can rewrite (12) as
$$\frac{_{i=1}^nI(f_i,x_i)(y_if_i)}{_{i=1}^nI(f_i,x_i)}_𝐋2^{\frac{m+l}{2}}\frac{diam(𝐘)I_{}\sqrt{n}}{_{i=1}^nI(f_i,x_i)}$$
(14)
(assuming the denominator $`_{i=1}^nI(f_i,x_i)`$ is positive); therefore, we can expect proper calibration and resolution in the soft neighborhood $`I`$ of $`(f^{},x^{})`$ when
$$\underset{i=1}{\overset{n}{}}I(f_i,x_i)\sqrt{n}.$$
(15)
## 7 Further research
The main result of this paper is an existence theorem: we did not show how to compute Forecaster’s strategy ensuring $`𝒦_0𝒦_1\mathrm{}`$. (The latter was easy in the case of binary forecasting considered in .) It is important to develop computationally efficient ways to find zeros of vector fields, at least when $`𝐋=^m`$. There are several popular methods for finding zeros, such as the Newton-Raphson method (see, e.g., , Chap. 9), but it would be ideal to have efficient methods that are guaranteed to find a zero (or a near zero) in a prespecified time.
### Acknowledgments
This work was partially supported by MRC (grant S505/65), Royal Society, and the Superrobust Computation Project (Graduate School of Information Science and Technology, University of Tokyo). We are grateful to anonymous reviewers for their comments.
## Appendix A Zeros of vector fields
The following lemma is the main component of the proof of Theorem 1.
###### Lemma 1
Let $`𝐅`$ be a compact convex non-empty set in a Hilbert space $`𝐋`$ and $`S:𝐅𝐋`$ be a continuous vector field on $`𝐅`$. If at no point of the boundary $`𝐅`$ the vector field $`S`$ is normal and directed exteriorly to $`𝐅`$ then there exists $`f𝐅`$ such that $`S(f)=0`$.
* For each $`f𝐋`$ define $`\sigma (f)`$ to be the point of $`𝐅`$ closest to $`f`$. A standard argument (see, e.g., , Theorem 12.3) shows that such a point exists: if $`d:=inf\{yf_𝐋|y𝐅\}`$, we can take any sequence $`y_n𝐅`$ with $`y_nf_𝐋d`$ and apply the parallelogram law $`ab^2+a+b^2=2a^2+2b^2`$ to obtain
$`y_my_n_𝐋^2`$ $`=(y_mf)(y_nf)_𝐋^2`$
$`=2y_mf_𝐋^2+2y_nf_𝐋^2(y_mf)+(y_nf)_𝐋^2`$
$`=2y_mf_𝐋^2+2y_nf_𝐋^24{\displaystyle \frac{y_m+y_n}{2}}f_𝐋^2`$
$`2y_mf_𝐋^2+2y_nf_𝐋^24d^22d^2+2d^24d^2=0`$
as $`m,n\mathrm{}`$; since $`𝐋`$ is complete and $`𝐅`$ is closed, $`y_ny`$ for some $`y𝐅`$, and it is clear that $`yf_𝐋=d`$. A closest point is indeed unique: if $`y_1f_𝐋=y_2f_𝐋=d`$ and $`y_1y_2`$, the parallelogram law would give
$$\begin{array}{c}\frac{y_1+y_2}{2}f_𝐋^2=\frac{1}{4}(y_1f)+(y_2f)_𝐋^2\hfill \\ \hfill =\frac{1}{2}y_1f_𝐋^2+\frac{1}{2}y_2f_𝐋^2\frac{1}{4}(y_1f)(y_2f)_𝐋^2\\ \hfill =d^2\frac{1}{4}y_1y_2_𝐋^2<d^2.\end{array}$$
(16)
Therefore, the function $`\sigma (f)`$ is well-defined. It is also continuous: if $`f\sigma (f)_𝐋=d`$ and $`f_nf`$, then $`f\sigma (f_n)_𝐋d`$ and, analogously to (16),
$$\begin{array}{c}d^2\frac{\sigma (f)+\sigma (f_n)}{2}f_𝐋^2=\frac{1}{4}(\sigma (f)f)+(\sigma (f_n)f)_𝐋^2\hfill \\ \hfill =\frac{1}{2}\sigma (f)f_𝐋^2+\frac{1}{2}\sigma (f_n)f_𝐋^2\frac{1}{4}(\sigma (f)f)(\sigma (f_n)f)_𝐋^2\\ \hfill =d^2+o(1)\frac{1}{4}\sigma (f)\sigma (f_n)_𝐋^2;\end{array}$$
therefore, $`\sigma (f_n)\sigma (f)`$ in $`𝐋`$.
For each $`f𝐅`$, let $`\mathrm{\Sigma }(f):=\sigma (f+S(f))`$ be the point of $`𝐅`$ closest to $`f+S(f)`$; since both $`\sigma `$ and $`S`$ are continuous, $`\mathrm{\Sigma }`$ is continuous. By the Schauder-Tikhonov theorem (see, e.g., , Theorem 5.28) there is a point $`f𝐅`$ such that $`\mathrm{\Sigma }(f)=f`$. If $`f`$ is an interior point of $`𝐅`$, $`\sigma (f+S(f))=f`$ implies $`S(f)=0`$, and so the conclusion of the lemma holds. It remains to consider the case $`f𝐅`$; in fact, we will show that this case is impossible. There exists $`y𝐅`$ such that $`S(f),yf_𝐋>0`$ (otherwise, $`S`$ would have been normal and directed exteriorly to $`𝐅`$), and we find for $`t(0,1)`$:
$$\begin{array}{c}(f+S(f))((1t)f+ty)_𝐋^2=S(f)t(yf)_𝐋^2\hfill \\ \hfill =S(f)_𝐋^22tS(f),yf_𝐋+t^2yf_𝐋^2;\end{array}$$
for a small enough $`t`$ this gives
$$(f+S(f))((1t)f+ty)_𝐋^2<S(f)_𝐋^2,$$
a contradiction.
## Appendix B Tensor product
In this appendix we list several definitions and simple facts about tensor products of Hilbert spaces, in the form used in this paper.
The *tensor product* $`𝐋𝐇`$ of Hilbert spaces $`𝐋`$ and $`𝐇`$ is defined in, e.g., , §II.4. Briefly, the definition is as follows. The space $`𝐋𝐇`$ is the subset of the set of bilinear forms $`v(l^{},h^{})`$, $`l^{}𝐋`$ and $`h^{}𝐇`$, obtained as the completion of the set of all linear combinations of the bilinear forms $`lh`$, where $`l𝐋`$ and $`h𝐇`$, defined by
$$(lh)(l^{},h^{}):=l,l^{}_𝐋h,h^{}_𝐇;$$
(17)
the inner product in $`𝐋𝐇`$ is determined uniquely by setting
$$l_1h_1,l_2h_2_{𝐋𝐇}:=l_1,l_2_𝐋h_1,h_2_𝐇.$$
(18)
In particular, (18) implies
$$lh_{𝐋𝐇}=l_𝐋h_𝐇$$
(19)
for all $`l𝐋`$ and $`h𝐇`$.
If $`v𝐋𝐇`$ and $`h𝐇`$, we define the *product* $`vh𝐋`$ by the requirement
$$v(l^{},h)=vh,l^{}_𝐋,l^{}𝐋$$
(the validity of this definition follows from the Riesz-Fischer theorem: all bilinear forms in $`𝐋𝐇`$ are clearly continuous).
###### Lemma 2
For any $`l𝐋`$ and $`h_1,h_2𝐇`$,
$$(lh_1)h_2=h_1,h_2_𝐇l.$$
(20)
* It suffices to prove
$$(lh_1)h_2,l^{}_𝐋=h_1,h_2_𝐇l,l^{}_𝐋,$$
which, by definition, is equivalent to
$$(lh_1)(l^{},h_2)=h_1,h_2_𝐇l,l^{}_𝐋$$
and, therefore, true (cf. (17)).
The following lemma is an easy implication of the Cauchy-Schwarz inequality.
###### Lemma 3
For any $`v𝐋𝐇`$ and $`h𝐇`$,
$$vh_𝐋v_{𝐋𝐇}h_𝐇.$$
* We are required to prove, for all $`l^{}𝐋`$,
$$vh,l^{}_𝐋v_{𝐋𝐇}h_𝐇l^{}_𝐋,$$
i.e.,
$$v(l^{},h)v_{𝐋𝐇}h_𝐇l^{}_𝐋.$$
We can assume that $`v=lh^{}`$, for some $`l𝐋`$ and $`h^{}𝐇`$, in which case the last inequality immediately follows from (17), (19), and the Cauchy-Schwarz inequality.
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# Gravitational waves from remnants of ultraluminous X-ray sources
## 1 Introduction
The counterpart of ultraluminous X-ray sources (ULXs), which have X-ray luminosities $`L_X`$ larger than the Eddington luminosity of a stellar mass object of mass $`M`$, $`L_X>L_E=1.3\times 10^{39}\mathrm{erg}\mathrm{s}^1M/10M_{}`$, are not known. Most likely there is no universal engine for ULXs. Some may be powered by anisotropic radiation (e.g. King et al. 2001; Rappaport, Podsiadlowski & Pfahl 2005) or super Eddington luminosity (Begelman 2002) from stellar mass black holes.
A third possibility is offered by the hypothesised existence of intermediate mass black holes (IBHs; see Miller & Colbert 2004 for a review) of masses $`10^3M_{}M_{}10^5M_{}`$, which could radiate isotropically at a sub-Eddington rate to account for the observed luminosities. There has been some observational evidence that at least some ULXs are powered by IBHs. In some cases the mass-temperature relation (Miller et al. 2003; Miller, Fabian & Miller 2004) or quasi periodic oscillations (Fiorito & Titarchuk 2004; Liu et al. 2005) indicate a high accreting mass.
Numerical N-body and Monte Carlo simulations indicate that the stellar collision rate in young dense stellar clusters during core collapse can become very large, giving rise to a hierarchical merger. In that case an object of thousands of solar masses may form (Portegies Zwart et al. 1999, 2004; Gürkan, Freitag & Rasio 2004; Freitag, Gürkan & Rasio 2005a ; Freitag, Rasio, & Baumgardt 2005b ). The fate of such a massive star is unclear, but it might lead to the formation of an IBH. This scenario is supported by the fact that ULXs correlate positively with star formation (Swartz et al. 2004; Liu, Bregman & Irwin 2005) and that some ULXs are associated with stellar clusters (e.g. Zezas et al. 2002).
The young stars in the host cluster of the IBH have strong winds which blow out the gas from the cluster, and there is insufficient free gas available to power a ULX. However, the IBH can acquire a companion star by dynamical capture or tidal capture. Here we discuss the latter possibility. Hopman, Portegies Zwart & Alexander (2004) showed that tidal capture of a main sequence star of mass $`M_{}`$ and radius $`R_{}`$ can lead to circularization close to the tidal radius
$`r_t`$ $`=`$ $`R_{}\left({\displaystyle \frac{M_{}}{M_{}}}\right)^{1/3},`$ (1)
which is the distance from the IBH where the tidal forces equal the forces which keep the star bound. As the star evolves it starts to fill its Roche lobe causing it to lose mass to the IBH. This leads to high X-ray luminosities, provided that the donor is sufficiently massive ($`M_{}10M_{}`$, Hopman et al. 2004; Portegies Zwart, Dewi, & Maccarone 2004; Li 2004). Interestingly, for some ULXs an optical counterpart has been identified, indicating that these ULXs are binary systems (Liu, Bregman, & Seitzer 2004; Kuntz et al. 2005).
In this scenario a ULX may turn on for at most the life-time $`t_{}10\mathrm{Myr}`$ of the captured star. In addition to strong X-ray emission the star emits gravitational waves (GWs). Portegies Zwart (2004) discussed the possibility of observing X-rays and GWs simultaneously. Only if the captured star is sufficiently light ($`M_{}2M_{}`$), its tidal radius is small enough for it to emit GWs in the LISA band during the mass transfer phase (Portegies Zwart 2004). However, in order to account for the high ULX luminosities the companion star should be considerably more massive than $`2M_{}`$ (Hopman et al. 2004; Portegies Zwart et al. 2004; Li 2004). In young and dense star clusters mass segregation causes most massive stars to accumulate in the cluster centre, and the combination of a high mass and large size make these stars excellent candidates for tidal capture. When after time $`t_{}`$ the massive donor explodes it turns into a compact remnant (CR); a neutron star (NS) or a stellar mass black hole (SBH). From that moment the ULX, deprived of its source of gas, turns off.
Here we focus on the subsequent –post supernova– evolution of the (IBH, CR) binary. In-spiral of the CR into an IBH due to GW emission results in a strong signal in frequencies measurable by space detectors such as LISA, providing a wealth of information about the system (see Miller 2002 for a discussion). We show that LISA may be able to observe such (IBH, CR) binaries.
## 2 Stellar capture and binary evolution
When stars orbit IBHs on wide, but highly eccentric orbits, with periapse close to the tidal radius of the IBH, $`r_pr_t`$, orbital energy is invested in tidal distortions of the star, causing it to spiral in. Hopman et al. (2004) showed that if the IBH is less massive than $`M_{\mathrm{max}}10^5M_{}`$, the star may be able to cool down by radiating the excess energy efficiently, and survive the strong tidal forces. Eventually it circularises near the IBH. We first estimate the fraction $`f_{\mathrm{merge}}`$ of isolated binaries which merge as a result of GW emission within the age of the Universe, while accounting for angular momentum conservation during the mass transfer phase and an isotropic velocity kick. We then discuss the consequences of the gravitational interaction of cluster stars with the (IBH, CR) binary.
### 2.1 Evolution of an isolated IBH-star binary
After tidal circularization near the IBH, the star may fill its Roche lobe, possibly after a period of stellar evolution and expansion. Roche lobe overflow from a $`M_{}10M_{}`$ donor can then give rise to high luminosities for a time limited by the life-time $`t_{}`$ of the star (Hopman et al. 2004; Portegies Zwart et al. 2004; Li 2004).
The main-sequence star fills its Roche lobe at a distance $`2r_t`$ before the terminal age main sequence. We assume that the entire hydrogen envelope is transferred to the IBH, while the binary conserves angular momentum. The mass of the stellar core of the donor star is given by $`M_c=0.08M_{}(M_{}/M_{})^{1.4}`$ (Iben, Tutukov & Yungelson 1995), while the mass of the hydrogen envelope is $`M_H=M_{}M_c`$. Angular momentum conservation implies that the binary separation increases during the mass transfer phase.
The remaining helium star subsequently explodes in a super nova after a time $`t_{}`$, and forms a CR. In this event a star of $`10M_{}<M_{}<20M_{}`$ forms a NS of Chandrasekhar mass, $`M_{\mathrm{Ch}}=1.4M_{}`$. Stars more massive stars than $`20M_{}`$ collapse to SBHs; we assume that the mass distribution is that found by Fyer & Kalogera (2001).
The semi-major axes and eccentricities of the CR are determined by the post mass-transfer orbital elements while accounting for the mass lost in the super nova and for the velocity kick, imparted to the CR at the time of the explosion. The kick velocity is taken in a random direction with magnitude from the distribution of pulsars velocities, which is well fitted by a double Gaussian distribution
$`f(v)`$ $`=`$ $`4\pi v^2[{\displaystyle \frac{w}{(2\pi \sigma _1^2)^{3/2}}}\mathrm{exp}(v^2/2\sigma _1^2)`$ (2)
$`+{\displaystyle \frac{(1w)}{(2\pi \sigma _2^2)^{3/2}}}\mathrm{exp}(v^2/2\sigma _2^2)],`$
where $`w=0.4`$, $`\sigma _1=90\mathrm{kms}^1`$, and $`\sigma _2=500\mathrm{kms}^1`$ (Arzoumanian, Chernoff & Cordes 2002). For SBHs we adopt the same distribution, but with velocities smaller by a factor $`M_{\mathrm{Ch}}/M_{\mathrm{SBH}}`$ (White & van Paradijs 1996; Gualandris et al. 2005). The velocity kick leads to an increase of the eccentricity of binary and a change in its energy; the kicks are generally insufficient to ionise the binary systems (van den Heuvel et al. 2000).
The surviving binary loses orbital energy due to the emission of GWs. A CR of mass $`M_{\mathrm{CR}}`$ on an orbit with semi-major axis $`a`$ and eccentricity $`e`$ around an IBH of mass $`M_{}M_{\mathrm{CR}}`$ loses energy due to the emission of GWs at rate of
$`\dot{E}_{\mathrm{GW}}={\displaystyle \frac{32}{5}}{\displaystyle \frac{G^4}{c^5}}{\displaystyle \frac{M_{}^3M_{\mathrm{CR}}^2}{a^5}}f(e),`$ (3)
with
$`f(e)={\displaystyle \frac{1+\frac{73}{24}e^2+\frac{37}{96}e^4}{(1e^2)^{7/2}}}.`$ (4)
If the orbit is circular $`(e=0)`$ spiral-in occurs on a time-scale of (Peters 1964)
$`t_{\mathrm{merge}}`$ $`=`$ $`{\displaystyle \frac{5}{256}}{\displaystyle \frac{c^5}{G^3}}{\displaystyle \frac{a^4}{M_{}^2M_{\mathrm{CR}}}}`$ (5)
$`=`$ $`3.5\mathrm{Gyr}\left({\displaystyle \frac{a}{\mathrm{AU}}}\right)^4\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^2\left({\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}\right)^1.`$
For non-zero eccentricities the in-spiral time is shorter by a factor $`(1e)^{7/2}`$.
We perform binary population synthesis with the above described scenario to compute the fraction of high mass binaries that become potential LISA sources. The mass of the stellar companion was selected from the initial mass function $`dN/dM_{}M_{}^\alpha `$, where we assumed $`\alpha =2`$, consistent with the mass function in the core of clusters in which a runaway merger occurred (Portegies Zwart et al. 2004). We assume a minimal mass of $`10M_{}`$, since lighter donors cannot account for ULX luminosities (Hopman et al. 2004; Portegies Zwart et al. 2004), and a maximal donor mass of $`100M_{}`$.
LISA is likely to observe (IBH, CR) binaries in the phase before the merger, rather than the merger event itself. However, the CRs spent only a very short time in the LISA frequency band as compared to the Hubble time (see §3). The question whether these sources are detectable therefore depends on how many merge within a the Hubble time.
In figures (1) and (2) we show the results for an IBH of $`M_{}=3\times 10^3M_{}`$ and $`M_{}=10^4M_{}`$, respectively. Only a fraction of the star spirals in within a Hubble time. For $`M_{}=3\times 10^3M_{}`$ nearly all objects which spiral in are NSs, with only very few SBHs. For more massive IBHs, a significant fraction of SBHs also spirals in fast enough, and the total fraction of merging objects exceeds 10% for $`M_{}\stackrel{>}{}10^4M_{}`$. In figure (3) we show $`f_{\mathrm{merge}}`$, the fraction of stars which spiral in within the age of the Universe, as a function of $`M_{}`$. After the velocity kick, the CRs have eccentricities up to $`e0.9`$. While this decreases $`t_{\mathrm{merge}}`$, these eccentricities are much smaller than the $`e0.995`$ eccentricities found by Hopman & Alexander (2005) for direct GW capture. By the time the stars spiral in to orbital frequencies $`\nu >10^4`$s, which is in the LISA band, the orbits are close to circular.
### 2.2 Interactions of the binary with cluster stars
The two-body relaxation time $`t_r`$ of star clusters which host a runaway merger is short ($`t_r30\mathrm{Myr}`$, Portegies Zwart & McMillan 2002), and these evaporate on a time-scale of $`10^8`$ years if tidally limited. As long as the cluster has not yet evaporated the (IBH, CR) binary has interactions with cluster stars. Since the (IBH, CR) binary is “hard” (Heggie 1975), these interactions tend to decrease the orbital separation between the IBH and the CR (Miller 2002). In addition, scattering changes the eccentricity of the binary. The angular momentum vector $`𝐉`$ of the binary performs a random walk, and its magnitude $`J`$ samples angular momenta $`0JJ_m`$ on the cluster’s relaxation time-scale (Alexander & Hopman 2003; Hopman & Alexander 2005). Here $`J_m`$ is the maximum angular momentum of a binary of given energy. When $`J`$ decreases the orbit becomes more eccentric, and the in-spiral time decreases (see eq. ).
It is not straightforward to quantify the effect of gravitational two-body scattering, since the cluster is not in a steady state, and in particular the relaxation time in the core can vary wildly. For simplicity we therefore neglect the effect of scattering in the following discussion on the in-spiral rate and the number of observable LISA sources. We note, however, that scattering may significantly increase the number of GW sources, so that the following LISA detection rates should be regarded as lower limits.
## 3 Observable gravitational waves from IBHs
The dimensionless strain of the GWs emitted at a frequency $`\nu =10^3\nu _3\mathrm{s}^1`$ from source at a distance $`d=d_{\mathrm{Mpc}}`$ Mpc is
$`h=3.5\times 10^{21}\nu _3^{2/3}\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^{2/3}{\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}d_{\mathrm{Mpc}}^1`$ (6)
(e.g. Sigurdsson & Rees 1997). LISA is sensitive to frequencies in the range $`10^4\mathrm{Hz}\nu 1`$ Hz. At $`\nu =10^3`$ Hz, LISA can detect sources with strains larger than $`h=\widehat{h}_{23}10^{23}`$, where $`\widehat{h}_{23}1`$. This estimate is based on a 1 yr observation with signal-to-noise ratio S/N = 1 (see, e.g., http://www.srl.caltech.edu/lisa). In the following we assume for concreteness a GW source with frequency $`\nu =10^3`$ Hz; application to other frequencies is straightforward. Sources can be observed to distances up to
$`d_{\mathrm{max}}354\mathrm{Mpc}\widehat{h}_{23}\nu _3^{2/3}\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^{2/3}{\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}.`$ (7)
We assume that only ULXs with luminosities $`L_X>10^{40}\mathrm{ergs}^1`$ contain IBHs (Portegies Zwart et al. 2004). Presently the average number of $`>10^{40}\mathrm{ergs}^1`$ ULXs per galaxy is $`0.1`$ (Swartz et al. 2004). The star formation rate dropped by an order of magnitude since $`z2`$ (Madau, Pozzetti & Dickinson 1998). ULXs correlate with star formation; here we assume that the number of ULXs is proportional to the star formation rate, in which case the number of luminous ULXs per galaxy at an earlier time was $`N_{\mathrm{ULX}}1`$. Since the ULX lives for a time $`t_{}=10^7\mathrm{yr}t_7`$, this yields an ULX formation rate per galaxy of $`10^7\mathrm{yr}^1N_{\mathrm{ULX}}t_7^1`$. As was discussed in the previous section, only a small fraction $`f_{\mathrm{merge}}=0.1f_1`$ of the ULXs with $`M_{}<10^4M_{}`$ leave behind a remnant binary which spirals in within a Hubble time. The rate at which observable GW sources are produced per relic ULX is then
$`\mathrm{\Gamma }_{\mathrm{GW}}=10^8\mathrm{yr}^1f_1N_{\mathrm{ULX}}t_7^1.`$ (8)
At the point where the period $`P=2\pi a^{3/2}/(GM_{})^{1/2}`$ equals $`10^3\mathrm{s}`$, the in-spiral time is
$`t_L=155\mathrm{yr}\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^{2/3}\left({\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}\right)^1,`$ (9)
where we assumed that the orbit has circularised when the frequency is this high. We thus find that the mean number of sources emitting GWs with frequencies of $`\nu 10^3\mathrm{s}^1`$ per galaxy is $`𝒩_1=\mathrm{\Gamma }_{\mathrm{GW}}t_L`$, or
$`𝒩_1=10^6f_1t_7^1N_{\mathrm{ULX}}\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^{2/3}\left({\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}\right)^1.`$ (10)
The local galaxy density is estimated to be $`n_{\mathrm{gal}}3\times 10^2\mathrm{Mpc}^3`$ (Marinoni et al. 1999). If the maximal distance at which the GW source can be observed is $`d_{\mathrm{max}}`$, the number of LISA sources in the sky at any moment is given by $`𝒩_{\mathrm{tot}}=𝒩_1n_{\mathrm{gal}}4\pi d_{\mathrm{max}}^3/3`$, or, using equation (7),
$`𝒩_{\mathrm{tot}}=8.6f_1\widehat{h}_{23}^3t_7^1N_{\mathrm{ULX}}\left({\displaystyle \frac{M_{}}{3\times 10^3M_{}}}\right)^{4/3}\left({\displaystyle \frac{M_{\mathrm{CR}}}{10M_{}}}\right)^2.`$ (11)
## 4 Discussion
For IBHs of $`M_{}3\times 10^3M_{}`$, $`f_10.3`$ and almost no SBHs spiral in (figure ). The main contribution to LISA sources comes from NSs, i.e. $`M_{\mathrm{CR}}/10M_{}0.14`$.We estimate that the merger rate for NSs with IBHs is too low to be likely to be seen by LISA. However, in this estimate we ignored the effect of three-body scattering on the binary orbit. Though hard to quantise, its result may be a marginal detection rate for NSs. We also note that we estimated that SBHs have a velocity kick distribution similar to that of NSs, but with velocities smaller by the mass ratio $`M_{\mathrm{Ch}}/M_{\mathrm{SBH}}`$ of the objects. This estimate is rather uncertain, and may be too conservative. Larger kick velocities cause more SBHs to merge within a Hubble time, in which case the contribution of $`M_{}3\times 10^3M_{}`$ IBHs to LISA would increase.
For IBHs of $`M_{}3\times 10^3M_{}`$ we find that $`f_11`$ and for SBHs $`t_{\mathrm{merge}}`$ decreases to well within a Hubble time. In this case stellar mass black holes $`(M_{\mathrm{CR}}/10M_{}1)`$ give the most promising GW sources, with a detection rate of about 10 per year, in the case that three-body scattering is ignored. Including three-body scattering in this estimate may boost the detection rate with a sizable fraction.
Our estimate in equation (11) for the number of sources that LISA can detect is conservative. First, in §3 the (IBH, CR) binary is considered to be isolated. Indeed the host cluster eventually evaporates, but this is preceded by a phase during which the binary interacts with other cluster stars (§2.2). These interactions tend to harden the binary, and change its eccentricity. As a result $`t_{\mathrm{merge}}`$ decreases, and thus the number of potential LISA sources increases. It is not implausible that nearly all (IBH, CR) binaries merge within a Hubble time, in which case $`f_1=10`$. Second, the life-time of ULXs is probably significantly shorter than the main sequence life-time $`t_{}`$ of the star. A more realistic assumption would be that ULX luminosities are only achieved when the donor is near the terminal age main sequence, in which case $`t_70.1`$, boosting the predicted detection rate for LISA with an order of magnitude. In conclusion, the number of observable GW sources could easily be orders of magnitude larger than the expression (11) indicates.
Previous studies of the number of potential LISA sources from binaries with an extremely small mass ratio focused mainly on cases in which the orbital energy is dissipated by the GWs themselves (e.g. Hils & Bender 1995; Sigurdsson & Rees 1997; Ivanov 2002; Miller 2002; Freitag 2001, 2003; Alexander & Hopman 2003; Hopman & Alexander 2005). In that case the event rate is on the order of a few per Gyr, and the orbits are highly eccentric, with typical eccentricities as large as $`e0.995`$ for $`\nu =10^4\mathrm{s}^1`$ near an IBH of $`M_{}=10^3M_{}`$ (Hopman & Alexander 2005). In that case stars spiral in very quickly, and emit GWs in the LISA band only for a short time, in contrast to the here discussed tidal capture scenario.
GW sources originating from tidal capture sources have nearly circular orbits when they enter the LISA frequency band, leading to a distinctly different signal than the highly eccentric sources originating from direct capture (e.g. Barack & Cutler BC04 2004; Wen & Gair WG05 2005). LISA can determine both the mass of the IBH and the eccentricity of the orbit. Detection of GWs from an IBH will give proof of the existence of these objects. If the signal stems from an orbit with low eccentricity, this supports the scenario that ULXs are accreting IBHs in binary systems.
## Acknowledgments
We thank T. Alexander for discussions. We are grateful to the Dutch Royal Academy of Arts and Sciences (KNAW), The Dutch Organization for Scientific Research (NWO) and the Dutch Advanced School for Astronomy (NOVA).
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# Statistics of lattice points in thin annuli for generic lattices
## 1 Introduction
We consider a variant of the lattice points counting problem. Let $`\mathrm{\Lambda }^2`$ be a planar lattice, with $`det\mathrm{\Lambda }`$ the area of its fundamental cell. Let
$$N_\mathrm{\Lambda }(t)=\{x\mathrm{\Lambda }:|x|t\},$$
denote its counting function, that is, we are counting $`\mathrm{\Lambda }`$-points inside a disc of radius $`t`$.
As well known, as $`t\mathrm{}`$, $`N_\mathrm{\Lambda }(t)\frac{\pi }{det\mathrm{\Lambda }}t^2`$. Denoting the remainder or the error term
$$\mathrm{\Delta }_\mathrm{\Lambda }(t)=N_\mathrm{\Lambda }(t)\frac{\pi }{det\mathrm{\Lambda }}t^2,$$
it is a conjecture of Hardy that
$$|\mathrm{\Delta }_\mathrm{\Lambda }(t)|_ϵt^{1/2+ϵ}.$$
Another problem one could study is the statistical behavior of the value distribution of $`\mathrm{\Delta }_\mathrm{\Lambda }`$ normalized by $`\sqrt{t}`$, namely of
$$F_\mathrm{\Lambda }(t):=\frac{\mathrm{\Delta }_\mathrm{\Lambda }(t)}{\sqrt{t}}.$$
Heath-Brown \[HB\] shows that for the standard lattice $`\mathrm{\Lambda }=^2`$, the value distribution of $`F_\mathrm{\Lambda }`$, weakly converges to a non-Gaussian distribution with density $`p(x)`$. Bleher \[BL3\] established an analogue of this theorem for a more general setting, where in particular it implies a non-Gaussian limiting distribution of $`F_\mathrm{\Lambda }`$, for any lattice $`\mathrm{\Lambda }^2`$.
However, the object of our interest is slightly different. Rather than counting lattice points in the circle of varying radius $`t`$, we will do the same for annuli. More precisely, we define
$$N_\mathrm{\Lambda }(t,\rho ):=N_\mathrm{\Lambda }(t+\rho )N_\mathrm{\Lambda }(t),$$
that is, the number of $`\mathrm{\Lambda }`$-points inside the annulus of inner radius $`t`$ and width $`\rho `$. The ”expected” value is the area $`\frac{\pi }{det\mathrm{\Lambda }}(2t\rho +\rho ^2)`$, and the corresponding normalized remainder term is
$$S_\mathrm{\Lambda }(t,\rho ):=\frac{N_\mathrm{\Lambda }(t+\rho )N_\mathrm{\Lambda }(t)\frac{\pi }{det\mathrm{\Lambda }}(2t\rho +\rho ^2)}{\sqrt{t}}.$$
The statistics of $`S_\mathrm{\Lambda }(t,\rho )`$ vary depending to the size of $`\rho (t)`$. Of our particular interest is the intermediate or macroscopic regime. Here $`\rho 0`$, but $`\rho t\mathrm{}`$. A particular case of the conjecture of Bleher and Lebowitz \[BL4\] states that $`S_\mathrm{\Lambda }(t,\rho )`$ has a Gaussian distribution. In 2004 Hughes and Rudnick \[HR\] established the Gaussian distribution for the unit circle, under an additional assumption that $`\rho (t)t^ϵ`$ for every $`ϵ>0`$.
By a rotation and dilation (which does not essentially effect the counting function), we may assume, with no loss of generality, that $`\mathrm{\Lambda }`$ admits a basis one of whose elements is the vector $`(1,0)`$, that is $`\mathrm{\Lambda }=1,\alpha +i\beta `$ (we make the natural identification of $`i`$ with $`(0,\mathrm{\hspace{0.17em}1})`$). In a previous paper \[W\] we already dealt with the problem of investigating the statistical properties of the error term for rectangular lattice $`\mathrm{\Lambda }=1,i\beta `$. We established the limiting Gaussian distribution for the ”generic” case in this 1-parameter family.
Some of the work done in \[W\] extends quite naturally for the 2-parameter family of planar lattices $`1,\alpha +i\beta `$. That is, in the current work we will require algebraic independence of $`\alpha `$ and $`\beta `$ as well as the ”strong Diophantinity” of the pair $`(\alpha ,\beta )`$ (to be defined), rather than transcendence and strong Diopantinity of the aspect ratio of the ellipse, as in \[W\].
We say that a real number $`\xi `$ is strongly Diophantine, if for every fixed natural $`n`$, there exists $`K_1>0`$, such that for integers $`a_j`$ with $`\underset{j=0}{\overset{n}{}}a_j\alpha ^j0`$,
$$\left|\underset{j=0}{\overset{n}{}}a_j\xi ^j\right|_n\frac{1}{\left(\underset{0jn}{\mathrm{max}}|a_j|\right)^{K_1}}.$$
It was shown by Mahler \[MAH\], that this property holds for a ”generic” real number. We say that a pair of numbers $`(\alpha ,\beta )`$ is strongly Diophantine, if for every fixed natural $`n`$, there exists a number $`K_1>0`$, such that for every integral polynomial $`p(x,y)=\underset{i+jn}{}a_{i,j}x^iy^j`$ of degree $`n`$, we have
$$|p(\alpha ,\beta )|_n\frac{1}{\underset{i+jn}{\mathrm{max}}|a_{i,j}|^{K_1}},$$
whenever $`p(\alpha ,\beta )0`$. This holds for almost all real pairs $`(\alpha ,\beta )`$, see section 2.2.
###### Theorem 1.1.
Let $`\mathrm{\Lambda }=1,\alpha +i\beta `$ where $`(\alpha ,\beta )`$ is algebraically independent and strongly Diophantine pair of real numbers. Assume that $`\rho =\rho (T)0`$, but for every $`\delta >0`$, $`\rho T^\delta `$. Then for every interval $`𝒜`$,
$$\underset{T\mathrm{}}{lim}\frac{1}{T}meas\{t[T,\mathrm{\hspace{0.17em}2}T]:\frac{S_\mathrm{\Lambda }(t,\rho )}{\sigma }𝒜\}=\frac{1}{\sqrt{2\pi }}\underset{𝒜}{}e^{\frac{x^2}{2}}𝑑x,$$
(1)
where the variance is given by
$$\sigma ^2:=\frac{4\pi }{\beta }\rho .$$
(2)
##### Remark:
Note that the variance $`\sigma ^2`$ is $`\alpha `$-independent, since the determinant $`det(\mathrm{\Lambda })=\beta `$.
One of the features of a rectangular lattice is that it is quite easy to show that the number of so-called close pairs of lattice points or pairs of points lying within a narrow annulus is bounded by essentially its average (see lemma 5.2 of \[W\]). This particular feature of the rectangular lattices was exploited while reducing the computation of the moments to the ones of a smooth counting function (we call it ”unsmoothing”). In order to prove an analogous bound for a general lattice, we extend a result from Eskin, Margulis and Mozes \[EMM\] for our needs to obtain proposition 3.1. We believe that this proposition is of independent interest.
## 2 The distribution of $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}`$
In this section, we are interested in the distribution of the smooth version of $`S_\mathrm{\Lambda }(t,\rho )`$, denoted $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}(t)`$, where $`L:=\frac{1}{\rho }`$ and $`M`$ is the smoothing parameter. Just as in \[W\] and \[HR\],
$$\stackrel{~}{S}_{\mathrm{\Lambda },M,L}(t)=\frac{\stackrel{~}{N}_{\mathrm{\Lambda },M}(t+\frac{1}{L})\stackrel{~}{N}_{\mathrm{\Lambda },M}(t)\frac{\pi }{d}(\frac{2t}{L}+\frac{1}{L^2})}{\sqrt{t}},$$
(3)
where $`\stackrel{~}{N}_{\mathrm{\Lambda },M}`$ is the smooth version of $`N_\mathrm{\Lambda }`$, computed by means of convolution of the characteristic function of the unit ball with $`\psi `$, a smooth function with a compact support (see \[HR\] or \[W\] for details). We assume that for every $`\delta >0`$, $`L=L(T)=O(T^\delta )`$, which corresponds to the assumption of theorem 1.1 regarding $`\rho :=\frac{1}{L}`$.
Rather than drawing $`t`$ at random from $`[T,\mathrm{\hspace{0.17em}2}T]`$ with a uniform distribution, we prefer to work with smooth densities: introduce $`\omega 0`$, a smooth function of total mass unity, such that both $`\omega `$ and $`\widehat{\omega }`$ are rapidly decaying, namely
$$|\omega (t)|\frac{1}{(1+|t|)^A},|\widehat{\omega }(t)|\frac{1}{(1+|t|)^A},$$
for every $`A>0`$. Define the averaging operator
$$f_T=\frac{1}{T}\underset{\mathrm{}}{\overset{\mathrm{}}{}}f(t)\omega (\frac{t}{T})𝑑t,$$
and let $`_{\omega ,T}`$ be the associated probability measure:
$$_{\omega ,T}(f𝒜)=\frac{1}{T}\underset{\mathrm{}}{\overset{\mathrm{}}{}}1_𝒜(f(t))\omega (\frac{t}{T})𝑑t.$$
##### Remark:
In what follows, we will suppress the explicit dependency on $`T`$, whenever convenient.
###### Theorem 2.1.
Suppose that $`M(T)`$ and $`L(T)`$ are increasing to infinity with $`T`$, such that $`M=O(T^\delta )`$ for all $`\delta >0`$, and $`L/\sqrt{M}0`$. Then if $`(\alpha ,\beta )`$ is an algebraically independent strongly Diophantine pair, we have for $`\mathrm{\Lambda }=1,\alpha +i\beta `$,
$$\underset{T\mathrm{}}{lim}_{\omega ,T}\left\{\frac{\stackrel{~}{S}_{\mathrm{\Lambda },M,L}}{\sigma }𝒜\right\}=\frac{1}{\sqrt{2\pi }}\underset{𝒜}{}e^{\frac{x^2}{2}}𝑑x,$$
for any interval $`𝒜`$, where
$$\sigma ^2:=\frac{4\pi }{\beta L}.$$
(4)
##### Definition:
A tuple of real numbers $`(\alpha _1,\mathrm{},\alpha _n)^n`$ is called Diophantine, if there exists a number $`K>0`$, such that for every integer tuple $`\{a_i\}_{i=0}^n`$,
$$\left|a_0+\underset{i=1}{\overset{n}{}}a_i\alpha _i\right|\frac{1}{q^K},$$
(5)
where $`q=\underset{0in}{\mathrm{max}}|a_i|`$. Khintchine proved that almost all tuples in $`^n`$ are Diophantine (see, e.g. \[S\], pages 60-63).
Denote the dual lattice
$$\mathrm{\Lambda }^{}=1,\frac{\alpha }{\beta }+i\frac{1}{\beta }.$$
We assume for the rest of current section that the set of squared norms of $`\mathrm{\Lambda }^{}`$ satisfy the Diophantine property, which means that $`(\alpha ^2,\alpha \beta ,\beta ^2)`$ is a Diophantine triple of numbers. We may assume the Diophantinity of $`(\alpha ^2,\alpha \beta ,\beta ^2)`$, since theorem 1.1 (and theorem 2.1) assume $`(\alpha ,\beta )`$ is strongly Diophantine, which is obviously a stronger assumption.
We use the following approximation to $`\stackrel{~}{N}_{\mathrm{\Lambda },M}(t)`$ (see e.g \[W\], lemma 4.1):
###### Lemma 2.2.
As $`t\mathrm{}`$,
$$\stackrel{~}{N}_{\mathrm{\Lambda },M}(t)=\frac{\pi t^2}{\beta }\frac{\sqrt{t}}{\beta \pi }\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{\mathrm{cos}\left(2\pi t|\stackrel{}{k}|+\frac{\pi }{4}\right)}{|\stackrel{}{k}|^{\frac{3}{2}}}\widehat{\psi }\left(\frac{|\stackrel{}{k}|}{\sqrt{M}}\right)+O\left(\frac{1}{\sqrt{t}}\right),$$
(6)
where, again, $`\mathrm{\Lambda }^{}`$ is the dual lattice.
By the definition of $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}`$ in (3) and appropriately manipulating the sum in (6) we obtain the following
###### Corollary 2.3.
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{\mathrm{\Lambda },M,L}(t)& =\frac{2}{\beta \pi }\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{\mathrm{sin}\left(\frac{\pi |\stackrel{}{k}|}{L}\right)}{|\stackrel{}{k}|^{\frac{3}{2}}}\mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{k}|+\frac{\pi }{4}\right)\widehat{\psi }\left(\frac{|\stackrel{}{k}|}{\sqrt{M}}\right)\hfill \\ & +O\left(\frac{1}{\sqrt{t}}\right).\hfill \end{array}$$
(7)
One should note that $`\widehat{\psi }`$ being compactly supported means that the sum essentially truncates at $`|\stackrel{}{k}|\sqrt{M}`$.
Unlike the standard lattice, clearly there are no nontrivial multiplicities in $`\mathrm{\Lambda }`$, that is
###### Lemma 2.4.
Let $`\stackrel{}{a_j}=m_j+n_j(\alpha +i\beta )\mathrm{\Lambda }`$, $`j=1,\mathrm{\hspace{0.17em}2}`$, with an irrational $`\alpha `$ such that $`\gamma (\alpha )`$. Then if $`|\stackrel{}{a_1}|=|\stackrel{}{a_2}|`$, either $`n_1=n_2`$ and $`m_1=m_2`$ or $`n_1=n_2`$ and $`n_2=m_2`$.
###### Proof of theorem 2.1.
We will show that the moments of $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}`$ corresponding to the smooth probability space converge to the moments of the normal distribution with zero mean and variance which is given by theorem 2.1. This allows us to deduce that the distribution of $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}`$ converges to the normal distribution as $`T\mathrm{}`$, precisely in the sense of theorem 2.1.
First, we show that the mean is $`O(\frac{1}{\sqrt{T}})`$. Since $`\omega `$ is real,
$$\left|\mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{k}|+\frac{\pi }{4}\right)\right|=\left|\mathrm{}m\left\{\widehat{\omega }\left(T|\stackrel{}{k}|\right)e^{i\pi (\frac{|\stackrel{}{k}|}{L}+\frac{1}{4}}\right\}\right|\frac{1}{T^A|\stackrel{}{k}|^A}$$
for any $`A>0`$, where we have used the rapid decay of $`\widehat{\omega }`$. Thus
$$\left|\stackrel{~}{S}_{\mathrm{\Lambda },M,L}\right|\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{1}{T^A|\stackrel{}{k}|^{A+3/2}}+O\left(\frac{1}{\sqrt{T}}\right)O\left(\frac{1}{\sqrt{T}}\right),$$
due to the convergence of $`\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{1}{|\stackrel{}{k}|^{A+3/2}}`$, for $`A>\frac{1}{2}`$
Now define
$$_{\mathrm{\Lambda },m}:=\left(\frac{2}{\beta \pi }\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{\mathrm{sin}\left(\frac{\pi |\stackrel{}{k}|}{L}\right)}{|\stackrel{}{k}|^{\frac{3}{2}}}\mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{k}|+\frac{\pi }{4}\right)\widehat{\psi }\left(\frac{|\stackrel{}{k}|}{\sqrt{M}}\right)\right)^m$$
(8)
Then from (7), the binomial formula and the Cauchy-Schwartz inequality,
$$\left(\stackrel{~}{S}_{\mathrm{\Lambda },M,L}\right)^m=_{\mathrm{\Lambda },m}+O\left(\underset{j=1}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)\frac{\sqrt{_{2m2j}}}{T^{j/2}}\right)$$
Proposition 2.5 together with proposition 2.8 allow us to deduce the result of theorem 2.1 for an algebraically independent strongly Diophantine $`(\xi ,\eta ):=(\frac{\alpha }{\beta },\frac{1}{\beta })`$. Clearly, $`(\alpha ,\beta )`$ being algebraically independent and strongly Diophantine is sufficient. ∎
### 2.1 The variance
The computation of the variance is done in two steps. First, we reduce the main contribution to the diagonal terms, using the assumption on the pair $`(\alpha ,\beta )`$ (i.e. $`(\alpha ^2,\alpha \beta ,\beta ^2)`$ is Diophantine). Then we compute the contribution of the diagonal terms. We sketch these steps, since they are very close to the corresponding one \[W\].
Suppose that the triple $`(\alpha ^2,\alpha \beta ,\beta ^2)`$ satisfies (5).
###### Proposition 2.5.
If $`M=O\left(T^{1/(K+1/2+\delta )}\right)`$ for fixed $`\delta >0`$, then the variance of $`\stackrel{~}{S}_{\mathrm{\Lambda },M,L}`$ is asymptotic to
$$\sigma ^2:=\frac{4}{\beta ^2\pi ^2}\underset{\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}}{}\frac{\mathrm{sin}^2\left(\frac{\pi |\stackrel{}{k}|}{L}\right)}{|\stackrel{}{k}|^3}\widehat{\psi }^2\left(\frac{|\stackrel{}{k}|}{\sqrt{M}}\right)$$
If $`L\mathrm{}`$, but $`L/\sqrt{M}0`$, then
$$\sigma ^2\frac{4\pi }{\beta L}$$
(9)
###### Proof.
Expanding out (8), we have
$$\begin{array}{cc}\hfill _{\mathrm{\Lambda },\mathrm{\hspace{0.17em}2}}=\frac{4}{\beta ^2\pi ^2}& \underset{\stackrel{}{k},\stackrel{}{l}\mathrm{\Lambda }^{}\{0\}}{}\frac{\mathrm{sin}\left(\frac{\pi |\stackrel{}{k}|}{L}\right)\mathrm{sin}\left(\frac{\pi |\stackrel{}{l}|}{L}\right)\widehat{\psi }\left(\frac{|\stackrel{}{k}|}{\sqrt{M}}\right)\widehat{\psi }\left(\frac{|\stackrel{}{l}|}{\sqrt{M}}\right)}{|\stackrel{}{k}|^{\frac{3}{2}}|\stackrel{}{l}|^{\frac{3}{2}}}\hfill \\ & \times \mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{k}|+\frac{\pi }{4}\right)\mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{l}|+\frac{\pi }{4}\right)\hfill \end{array}$$
(10)
It is easy to check that the average of the second line of the previous equation is:
$$\begin{array}{cc}\hfill \frac{1}{4}[& \widehat{\omega }\left(T(|\stackrel{}{k}||\stackrel{}{l}|)\right)e^{i\pi (1/L)(|\stackrel{}{l}||\stackrel{}{k}|)}+\hfill \\ & \widehat{\omega }\left(T(|\stackrel{}{l}||\stackrel{}{k}|)\right)e^{i\pi (1/L)(|\stackrel{}{k}||\stackrel{}{l}|)}+\hfill \\ & \widehat{\omega }\left(T(|\stackrel{}{k}|+|\stackrel{}{l}|)\right)e^{i\pi (1/2+(1/L)(|\stackrel{}{k}|+|\stackrel{}{l}|))}\hfill \\ & \widehat{\omega }(T(|\stackrel{}{k}|+|\stackrel{}{l}|))e^{i\pi (1/2+(1/L)(|\stackrel{}{k}|+|\stackrel{}{l}|))}]\hfill \end{array}$$
(11)
Recall that the support condition on $`\widehat{\psi }`$ means that $`\stackrel{}{k}`$ and $`\stackrel{}{l}`$ are both constrained to be of length $`O(\sqrt{M})`$. Thus the off-diagonal contribution (that is for $`|\stackrel{}{k}||\stackrel{}{l}|`$ ) of the first two lines of (11) is
$$\underset{\begin{array}{c}\stackrel{}{k},\stackrel{}{l}\mathrm{\Lambda }^{}\{0\}\\ |\stackrel{}{k}|,|\stackrel{}{k^{}}|\sqrt{M}\end{array}}{}\frac{M^{A(K+1/2)}}{T^A}\frac{M^{A(K+1/2)+2}}{T^A}T^B,$$
for every $`B>0`$, by Diophantinity of $`(\alpha ,\alpha \beta ,\beta ^2)`$.
Obviously, the contribution to (10) of the two last lines of (11) is negligible both in the diagonal and off-diagonal cases, justifying the diagonal approximation of (10) in the first statement of the proposition, and we omit the rest of the proof. ∎
### 2.2 The higher moments
In order to compute the higher moments we will prove that the main contribution comes from the so-called diagonal terms (to be explained later). In order to bound the contribution of the off-diagonal terms, we will use the following theorem, which is a consequence of the work of Kleinbock and Margulis \[KM\]. The contribution of the diagonal terms is computed exactly in the same manner it was done in \[W\], and so we will omit it here.
###### Theorem 2.6.
Let an integer $`n`$ be given. Then almost all pairs of real numbers $`(\xi ,\eta )^2`$ satisfy the following property: there exists a number $`K_1`$ such that for every integer polynomial of $`2`$ variables $`p(x,y)=\underset{i+jn}{}a_{i,j}x^iy^j`$ with degree $`n`$, we have
$$\left|p(\xi ,\eta )\right|h^{K_1},$$
where $`h=\underset{i+jn}{\mathrm{max}}|a_{i,j}|`$ is the height of $`p`$. The constant involved in the $`\mathrm{"}\mathrm{"}`$ notation depends only on $`\xi ,\eta `$ and $`n`$.
We will remark that theorem A in \[KM\] is much more general when the result we are using. As a matter of fact, we have the inequality
$$\left|b_0+b_1f_1(x)+\mathrm{}+b_nf_n(x)\right|_ϵ\frac{1}{h^{n+ϵ}}$$
with $`b_i`$ and
$$h:=\underset{0in}{\mathrm{max}}|b_i|.$$
The inequality above holds for every $`ϵ>0`$ for a wide class of functions $`f_i:U`$, for almost all $`xU`$, where $`U^m`$ is an open subset. Here we use this inequality for the monomials.
##### Definition:
We call the pairs $`(\xi ,\eta )`$ which satisfy for all natural $`n`$ the property of theorem 2.6, strongly Diophantine. Thus theorem 2.6 states that almost all real pairs of numbers are strongly Diophantine.
##### Remark:
Simon Kristensen \[KR\] has recently shown, that the set of all pairs $`(\xi ,\eta )^2`$ which fail to be strongly Diophantine has Hausdorff dimension $`1`$.
Obviously, strong Diophantinity of $`(\xi ,\eta )`$ implies Diophantinity of any $`n`$-tuple of real numbers which consists of any set of monomials in $`\xi `$ and $`\eta `$. Moreover, $`(\xi ,\eta )`$ is strongly Diophantine iff $`(\frac{\alpha }{\beta },\frac{1}{\beta })`$ is such.
We have the following analogue of lemma 4.7 in \[W\], which will eventually allow us to exploit the strong Diophantinity of $`(\alpha ,\beta )`$.
###### Lemma 2.7.
If $`(\xi ,\eta )`$ is strongly Diophantine, then it satisfies the following property: for any fixed natural $`m`$, there exists $`K`$, such that if
$$z_j=a_j^2+b_j^2\xi ^2+2a_jb_j\xi +b_j^2\eta ^2M,$$
and $`ϵ_j=\pm 1`$ for $`j=1,\mathrm{},m`$, with integral $`a_j,b_j`$ and if $`\underset{j=1}{\overset{m}{}}ϵ_j\sqrt{z_j}0`$, then
$$\left|\underset{j=1}{\overset{m}{}}ϵ_j\sqrt{z_j}\right|M^K,$$
(12)
where the constant involved in the $`\mathrm{"}\mathrm{"}`$ notation depends only on $`\eta `$ and $`m`$.
The proof is essentially the same as the one of lemma 4.7 from \[W\], considering the product $`Q`$ of numbers of the form $`\underset{j=1}{\overset{m}{}}\delta _j\sqrt{z_j}`$ over all possible signs $`\delta _j`$. Here we use the Diophantinity of the real tuple $`(\xi ,\eta )`$ rather than of a single real number.
###### Proposition 2.8.
Let $`m`$ be given. Suppose that $`\mathrm{\Lambda }=1,\alpha +i\beta `$, such that the pair $`(\xi ,\eta ):=(\frac{\alpha }{\beta },\frac{1}{\beta })`$ is algebraically independent strongly Diophantine, which satisfy the property of lemma 2.7 for the given $`m`$, with $`K=K_m`$. Then if $`=O\left(T^{\frac{1\delta }{K_m}}\right)`$ for some $`\delta >0`$, and if $`L\mathrm{}`$ such that $`L/\sqrt{M}0`$, the following holds:
$$\frac{_{\mathrm{\Lambda },m}}{\sigma ^m}=\{\begin{array}{cc}\frac{m!}{2^{m/2}\left(\frac{m}{2}\right)!}+O\left(\frac{\mathrm{log}L}{L}\right),\hfill & m\text{ is even}\hfill \\ O\left(\frac{\mathrm{log}L}{L}\right),\hfill & m\text{ is odd}\hfill \end{array}$$
###### Proof.
Expanding out (8), we have
$$\begin{array}{cc}\hfill _{\mathrm{\Lambda },m}=\frac{2^m}{\beta ^m\pi ^m}\underset{\stackrel{}{k_1},\mathrm{},\stackrel{}{k_m}\mathrm{\Lambda }^{}\{0\}}{}& \underset{j=1}{\overset{m}{}}\frac{\mathrm{sin}\left(\frac{\pi |\stackrel{}{k_j}|}{L}\right)\widehat{\psi }\left(\frac{|\stackrel{}{k_j}|}{\sqrt{M}}\right)}{|\stackrel{}{k_j}|^{\frac{3}{2}}}\hfill \\ & \times \underset{j=1}{\overset{m}{}}\mathrm{sin}\left(2\pi \left(t+\frac{1}{2L}\right)|\stackrel{}{k_1}|+\frac{\pi }{4}\right)\hfill \end{array}$$
(13)
Now,
$$\begin{array}{cc}\hfill \underset{j=1}{\overset{m}{}}& \mathrm{sin}(2\pi (t+\frac{1}{2L})|\stackrel{}{k_1}|+\frac{\pi }{4})\hfill \\ & =\underset{ϵ_j=\pm 1}{}\frac{\underset{j=1}{\overset{m}{}}ϵ_j}{2^mi^m}\widehat{\omega }\left(T\underset{j=1}{\overset{m}{}}ϵ_j|\stackrel{}{k_j}|\right)e^{\pi i\underset{j=1}{\overset{m}{}}ϵ_j\left((1/L)|\stackrel{}{k_j}|+1/4\right)}\hfill \end{array}$$
We call a term of the summation in (13) with $`\underset{j=1}{\overset{m}{}}ϵ_j|\stackrel{}{k_j}|=0`$ diagonal, and off-diagonal otherwise. Due to lemma 2.7, the contribution of the off-diagonal terms is:
$$\underset{\stackrel{}{k_1},\mathrm{},\stackrel{}{k_m}\mathrm{\Lambda }^{}\{0\}}{}\left(\frac{T}{M^{K_m}}\right)^AM^mT^{A\delta },$$
for every $`A>0`$, by the rapid decay of $`\widehat{\omega }`$ and our assumption regarding $`M`$.
Since $`m`$ is constant, this allows us to reduce the sum to the diagonal terms. In order to be able to sum over all the diagonal terms we need the following analogue of a well-known theorem due to Besicovitch \[BS\] about incommensurability of square roots of integers.
###### Proposition 2.9.
Suppose that $`\xi `$ and $`\eta `$ are algebraically independent, and
$$z_j=a_j^2+2a_jb_j\xi +b_j^2(\xi ^2+\eta ),$$
(14)
such that $`(a_j,b_j)_+^2`$ are all different primitive vectors, for $`1jm`$. Then $`\{\sqrt{z_j}\}_{j=1}^m`$ are linearly independent over $``$.
The last proposition is an immediate consequence of a theorem proved in the appendix of \[BL2\].
Computing the contribution of the diagonal terms is done literally the same way it was done in \[W\] and thus it is omitted here. In order to be able to sum over the diagonal terms, we use here proposition 2.9, rather than proposition 3.2 in \[W\].
## 3 Bounding the number of close pairs of lattice points
Roughly speaking, we say that a pair of lattice points, $`n`$ and $`n^{}`$ is close, if $`\left||n||n^{}|\right|`$ is small. We would like to show that this phenomenon is rare. This is closely related to the Oppenheim conjecture, as $`|n|^2|n^{}|^2`$ is a quadratic form on the coefficients of $`n`$ and $`n^{}`$. In order to establish a quantative result, we use a technique developed in a paper by Eskin, Margulis and Mozes \[EMM\].
### 3.1 Statement of the results
The ultimate goal of this section is to establish the following
###### Proposition 3.1.
Let $`\mathrm{\Lambda }`$ be a lattice and denote
$$A(R,\delta ):=\{(\stackrel{}{k},\stackrel{}{l})\mathrm{\Lambda }:R|\stackrel{}{k}|^22R,|\stackrel{}{k}|^2|\stackrel{}{l}|^2|\stackrel{}{k}|^2+\delta \}.$$
(15)
Then if $`\delta >1`$, such that $`\delta =o(R)`$, we have
$$\mathrm{\#}A(R,\delta )R\delta \mathrm{log}R$$
In order to prove this result, we note that evaluating the size of $`A(R,\delta )`$ is equivalent to counting integer points $`\stackrel{}{v}^4`$ with $`T\stackrel{}{v}2T`$ such that
$$0Q_1(v)\delta ,$$
where $`Q_1`$ is a quadratic form of signature $`(2,\mathrm{\hspace{0.17em}2})`$, given explicitly by
$$Q_1(\stackrel{}{v})=(v_1+v_2\alpha )^2+(v_2\beta )^2(v_3+v_4\alpha )^2(v_4\beta )^2.$$
(16)
For a fixed $`\delta >0`$ and a large $`R`$, this situation was considered extensively by Eskin, Margulis and Mozes \[EMM\]. We will examine how the constants involved in their result depend on $`\delta `$, and find out that there is a linear dependency, which is what we essentially need. The author wishes to thank Alex Eskin for his assistance with this matter.
##### Remark:
For our purposes we need a weaker result:
$$\mathrm{\#}A(R,\delta )_ϵR\delta R^ϵ,$$
for every $`ϵ>0`$. If $`\mathrm{\Lambda }`$ is a rectangular lattice (i.e. $`\alpha =0`$), then this result follows from properties of the divisor function (see e.g. \[BL\], lemma 3.2).
Theorem 2.3 in \[EMM\] considers a more general setting than proposition 3.1. We state here theorem 2.3 from \[EMM\] (see theorem 3.2). It follows from theorem 3.3 from \[EMM\], which will be stated as well (see theorem 3.3). Then we give an outline of the proof of theorem 2.3 of \[EMM\], and inspect the dependency on $`\delta `$ of the constants involved.
### 3.2 Theorems 2.3 and 3.3 from \[EMM\]
Let $`\mathrm{\Delta }`$ be a lattice in $`^n`$. We say that a subspace $`L^n`$ is $`\mathrm{\Delta }`$-rational, if $`L\mathrm{\Delta }`$ is a lattice in $`L`$. We need the following definitions:
##### Definitions:
$$\alpha _i(\mathrm{\Delta }):=sup\left\{\frac{1}{d_\mathrm{\Delta }(L)}\right|L\text{ is a }\mathrm{\Delta }\text{rational subspace of dimension }i\},$$
where
$$d_\mathrm{\Delta }(L):=vol(L/(L\mathrm{\Delta })).$$
Also
$$\alpha (\mathrm{\Delta }):=\underset{0in}{\mathrm{max}}\alpha _i(\mathrm{\Delta }).$$
Since the space of unimodular lattices is canonically isomorphic to
$`SL(n,)/SL(n,)`$, the notation $`\alpha (g)`$ makes sense for $`gG:=SL(n,)`$.
For a bounded function $`f:^n`$, with $`|f|M`$, which vanishes outside a ball $`B(0,R)`$, define $`\stackrel{~}{f}:SL(n,)`$ by the following formula:
$$\stackrel{~}{f}(g):=\underset{v^n}{}f(gv).$$
Lemma 3.1 in \[S2\] implies that
$$\stackrel{~}{f}(g)<c\alpha (g),$$
(17)
where $`c=c(f)`$ is an explicit constant constant
$$c(f)=c_0M\mathrm{max}(1,R^n),$$
for some constant $`c_0=c_0(n)`$, independent on f. In section 3.4 we prove a stronger result, assuming some additional information about the support of $`f`$.
Let $`Q_0`$ be a quadratic form defined by
$$Q_0(\stackrel{}{v})=2v_1v_n+\underset{i=2}{\overset{p}{}}v_i^2\underset{i=p+1}{\overset{n1}{}}v_i^2.$$
Since
$$v_1v_n=\frac{(v_1+v_n)^2(v_1v_n)^2}{2},$$
$`Q_0`$ is of signature $`p,q`$. Obviously, $`G:=SL(n,)`$ acts on the space of quadratic forms of signature $`(p,q)`$, and discriminant $`\pm 1`$, $`𝒪=𝒪(p,q)`$ by:
$$Q^g(v):=Q(gv).$$
Moreover, by the well known classification of quadratic forms, $`𝒪`$ is the orbit of $`Q_0`$ under this action.
In our case the signature is $`(p,q)=(2,\mathrm{\hspace{0.17em}2})`$ and $`n=4`$. We fix an element $`h_1G`$ with $`Q^{h_1}=Q_1`$, where $`Q_1`$ is given by (16). There exists a constant $`\tau >0`$, such that for every $`v^4`$,
$$\tau ^1vh_1v\tau v.$$
(18)
We may assume, with no loss of generality that $`\tau 1`$.
Let $`H:=Stab_{Q_0}(G)`$. Then the natural mophism $`H\backslash G𝒪(p,q)`$ is a homeomorphism. Define a $`1`$-parameter family $`a_tG`$ by:
$$a_te_i=\{\begin{array}{cc}e^te_1,\hfill & i=1\hfill \\ e_i,\hfill & i=2,\mathrm{},n1\hfill \\ e^te_n,\hfill & i=n\hfill \end{array}.$$
Clearly, $`a_tH`$. Furthermore, let $`\widehat{K}`$ be the subgroup of $`G`$ consisting of orthogonal matrices, and denote $`K:=H\widehat{K}`$.
Let $`(a,b)^2`$ be given and let $`Q:^n`$ be any quadratic form. The object of our interest is:
$$V_{(a,b)}()=V_{(a,b)}^Q()=\{x^n:a<Q(x)<b\}.$$
Theorem 2.3 states, in our case:
###### Theorem 3.2 (Theorem 2.3 from \[EMM\]).
Let $`\mathrm{\Omega }=\{v^4|v<\nu (v/v)\}`$, where $`\nu `$ is a nonnegative continuous function on $`S^3`$. Then we have:
$$\mathrm{\#}V_{(a,b)}^{Q_1}()T\mathrm{\Omega }<cT^2\mathrm{log}T,$$
where the constant $`c`$ depends only on $`(a,b)`$.
The proof of theorem 3.2 relies on theorem 3.3 from \[EMM\], and we give here a particular case of this theorem
###### Theorem 3.3 (Theorem 3.3 from \[EMM\]).
For any (fixed) lattice $`\mathrm{\Delta }`$ in $`^4`$,
$$\underset{t>1}{sup}\frac{1}{t}\underset{K}{}\alpha (a_tk\mathrm{\Delta })𝑑m(k)<\mathrm{},$$
where the upper bound is universal.
### 3.3 Outline of the proof of theorem 3.2:
##### Step 1:
Define
$$J_f(r,\zeta )=\frac{1}{r^2}\underset{^2}{}f(r,x_2,x_3,x_4)𝑑x_2𝑑x_3,$$
(19)
where
$$x_4=\frac{\zeta x_2^2+x_3^2}{2r}$$
Lemma 3.6 in \[EMM\] states that $`J_f`$ is approximable by means of an integral over the compact subgroup K. More precisely, there is some constant $`C>0`$, such that for every $`ϵ>0`$,
$$\left|Ce^{2t}\underset{K}{}f(a_tkv)\nu (k^1e_1)𝑑m(k)J_f(ve^t,Q_0(v))\nu (\frac{v}{v})\right|<ϵ$$
(20)
with $`e^t,v>T_0`$ for some $`T_0>0`$.
##### Step 2:
Choose a continuous nonnegative function $`f`$ on $`_+^4=\{x_1>0\}`$ which vanishes outside a compact set so that
$$J_f(r,\zeta )1+ϵ$$
on $`[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [a,b]`$. We will show later, how one can choose $`f`$.
##### Step 3:
Denote $`T=e^t`$, and suppose that $`Tv2T`$ and $`aQ_0(h_1v)b`$. Then by (18), $`J_f(h_1vT^1,Q_0(h_1v))1+ϵ`$, and by (20), for a sufficiently large $`t`$,
$$CT^2\underset{K}{}f(a_tkh_1v)𝑑m(k)1,$$
(21)
for $`Tv2T`$ and
$$aQ_0^x(v)b.$$
(22)
##### Step 4:
Summing (21) over all $`v^4`$ with (22) and $`Tv2T`$, we obtain:
$$\begin{array}{cc}\hfill \mathrm{\#}V_{(a,b)}()[T,\mathrm{\hspace{0.17em}2}T]S^3& \underset{v^n}{}CT^2\underset{K}{}f(a_tkh_1v)𝑑m(k)\hfill \\ & =CT^2\underset{K}{}\stackrel{~}{f}(a_tkh_1)𝑑m(k)\hfill \end{array}$$
(23)
using the nonnegativity of $`f`$.
##### Step 5:
By (17), (23) is
$$Cc(f)T^2\underset{K}{}\alpha (a_tkh_1)𝑑m(k).$$
##### Step 6:
The result of theorem 2.3 is obtained by using theorem 3.3 on the last expression.
### 3.4 $`\delta `$-dependency:
In this section we assume that $`(a,b)=(0,\delta )`$, which suits the definition of the set $`A(R,\delta )`$, (15). One should notice that there only $`3`$ $`\delta `$-dependent steps:
$``$ Choosing $`f`$ in step 2, such that $`J_f1+ϵ`$ on $`[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [0,\delta ]`$. We will construct a family of functions $`f_\delta `$ with an universal bound $`|f_\delta |M`$, such that $`f_\delta `$ vanishes outside of a compact set which is only slightly larger than
$$V(\delta )=[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [1,1]^2\times [0,\frac{\delta \tau }{2}].$$
(24)
This is done in section 3.4.1.
$``$ The dependency of $`T_0`$ of step 3, so that the usage of lemma 3.6 in \[EMM\] is legitimate. For this purpose we will have to examine the proof of this lemma. This is done in section 3.4.2.
$``$ The constant $`c`$ in (17). We would like to establish a linear dependency on $`\delta `$. This is straightforward, once we are able to control the number of integral points in a domain defined by (24). This is done in section 3.4.3.
#### 3.4.1 Choosing $`f_\delta `$:
##### Notation:
For a set $`U^n`$, and $`ϵ>0`$, denote
$$U_ϵ:=\{x^n:\underset{1in}{\mathrm{max}}|x_iy_i|ϵ,\text{for some }yU\}.$$
Choose a nonnegative continuous function $`f_0`$, on $`_+^4`$, which vanishes outside a compact set, such that its support, $`E_{f_0}`$, slightly exceeds the set $`V(1)`$. More precisely, $`V(1)E_{f_0}V(1)_{\delta _0}`$ for some $`\delta _0>0`$. By the uniform continuity of $`f`$, there are $`ϵ_0,\delta _0>0`$, such that if $`\underset{1i4}{\mathrm{max}}|x_ix_i^0|\delta _0`$, then $`f(x)>ϵ_0`$, for every $`x^0=(x_1^0,0,\mathrm{\hspace{0.17em}0},x_4^0)V(1)`$.
Thus for $`(r,\zeta )[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [0,\delta ]`$, the contribution of $`[\delta _0,\delta _0]^2`$ to $`J_{f_0}`$ is $`ϵ_0(2\delta _0)^2`$. Multiplying $`f_0`$ by a suitable factor, and by the linearity of $`J_{f_0}`$, we may assume that this contribution is at least $`1+ϵ`$.
Now define $`f_\delta (x_1,\mathrm{},x_4):=f_0(x_1,x_2,x_3,\frac{x_4}{\delta })`$. We have for $`\delta 1`$
$$\frac{\zeta x_2^2+x_3^2}{2r\delta }=\frac{\zeta /2r}{\delta }\frac{(x_2/\sqrt{\delta })^2}{2r}+\frac{(x_3/\sqrt{\delta })^2}{2r}.$$
Thus for $`\delta 1`$, if $`(r,\zeta )[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [0,\delta ]`$ and for $`i=2,\mathrm{\hspace{0.17em}3}`$, $`|x_i|<\delta _0`$, $`f_\delta `$ satisfies:
$$f_\delta (r,x_2,x_3,x_4)>ϵ_0,$$
and therefore the contribution of this domain to $`J_{f_\delta }`$ is
$$ϵ_0(2\delta )^21+ϵ$$
by our assumption.
By the construction, the family $`\{f_\delta \}`$ has a universal upper bound $`M`$ which is the one of $`f_0`$.
#### 3.4.2 How large is $`T_0`$
The proof of lemma 3.6 from \[EMM\] works well along the same lines, as long as
$$f(a_tx)0$$
(25)
implies that for $`t\mathrm{}`$, $`x/x`$ converges to $`e_1=(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0})`$. Now, since $`a_t`$ preserves $`x_1x_4`$, (25) implies for the particular choice of $`f=f_\delta `$ in section 3.4.1:
$$|x_1x_4|=O(\delta );x_1T.$$
Thus
$$x=x_1+O(\frac{\delta }{T})+O(1),$$
and so, as long as $`\delta =o(T)`$, $`x/x`$ indeed converges to $`e_1`$.
#### 3.4.3 Bounding integral points in $`V_\delta `$:
###### Lemma 3.4.
Let $`V(\delta )`$ defined by
$$V(\delta )=[\tau ^1,\mathrm{\hspace{0.17em}2}\tau ]\times [1,1]^{n2}\times [0,\frac{\delta \beta }{2}].$$
(26)
for some constant $`\tau `$ and $`n3`$. Let $`gSL(n,)`$ and denote
$$N(g,\delta ):=\mathrm{\#}V(\delta )g^n.$$
Then for $`\delta 1`$,
$$\left|N(g,\delta )\frac{2^{n2}(2\tau \tau ^1)\delta }{detg}\right|c_5\delta \underset{i=1}{\overset{n1}{}}\frac{1}{vol(L_i/(g^nL_i)}$$
for some $`g`$-rational subspaces $`L_i`$ of $`^4`$ of dimension $`i`$, where $`c_5=c_5(n)`$ depends only on $`n`$.
A direct consequence of lemma 3.4 is the following
###### Corollary 3.5.
Let $`f:^n`$ be a nonnegative function which vanishes outside a compact set $`E`$. Suppose that $`EV_ϵ(\delta )`$ for some $`ϵ>0`$. Then for $`\delta 1`$, (17) is satisfied with
$$c(f)=c_3M\delta ,$$
where the constant $`c_3`$ depends on $`n`$ only.
In order to prove lemma 3.4, we shall need the following:
###### Lemma 3.6.
Let $`\mathrm{\Lambda }^n`$ be a $`m`$-dimensional lattice, and let
$$A_t=\left(\begin{array}{c}1\\ & 1\\ & & \mathrm{}\\ & & & t\end{array}\right)$$
(27)
an $`n`$-dimensional linear transformation. Then for $`t>0`$ we have
$$detA_t\mathrm{\Lambda }tdet\mathrm{\Lambda }.$$
(28)
###### Proof.
We may assume that $`m<n`$, since if $`m=n`$, we obviously have an equality. Let $`v_1,\mathrm{},v_m`$ the basis of $`\mathrm{\Lambda }`$ and denote for every $`i`$, $`u_i^{n1}`$ the vector, which consists of first $`n1`$ coordinates of $`v_i`$. Also, let $`x_i`$ be the last coordinate of $`v_i`$. By switching vectors, if necessary, we may assume $`x_10`$. We consider the function
$$f(t):=(detA_t\mathrm{\Lambda })^2,$$
as a function of $`t`$.
Obviously,
$$f(t)=det(<u_i,u_j>+x_ix_jt^2)_{1i,jm}.$$
Substracting $`\frac{x_i}{x_1}`$ times the first row from any other, we obtain:
$$f(t)=\left|\begin{array}{c}<u_1,u_j>+x_1x_jt^2\\ <u_2,u_j>\frac{x_2}{x_1}<u_1,u_j>\\ \mathrm{}\\ <u_m,u_j>\frac{x_m}{x_1}<u_1,u_j>\end{array}\right|,$$
and by the multilinearity property of the determinant, $`f`$ is a linear function of $`t^2`$. Write
$$f(t)=a(t^21)+bt^2.$$
Thus
$$b=f(1);a=f(0),$$
and so $`b=det\mathrm{\Lambda }`$, and $`a=det<u_i,u_j>0`$, being minus the determinant of a Gram matrix. Therefore,
$$(detA_t\mathrm{\Lambda })^2t^2det\mathrm{\Lambda }=a(t^21)0$$
for $`t1`$, implying (27). ∎
###### Proof of lemma 3.4.
We will prove the lemma, assuming $`\beta =2`$. However, it implies the result of the lemma for any $`\beta `$, affecting only $`c_5`$. Let $`\delta >0`$. Trivially,
$$N(g,\delta )=N(g_0,\mathrm{\hspace{0.17em}1}),$$
where $`g_0=A_\delta ^1g`$ with $`A_\delta `$ given by (27). Let $`\lambda _1\lambda _2\mathrm{}\lambda _n`$ be the successive minima of $`g_0`$, and pick linearly independent lattice points $`v_1,\mathrm{},v_n`$ with $`v_i=\lambda _i`$. Denote $`M_i`$ the linear space spanned by $`v_1,\mathrm{},v_i`$ and the lattice $`\mathrm{\Lambda }_i=g_0^nM_i`$.
First, assume that $`\lambda _n\sqrt{\tau ^2+(n1)}=:r`$. Now, by Gauss’ argument,
$$\left|N(g_0,\mathrm{\hspace{0.17em}1})\frac{2^{n1}(2\tau \tau ^1)\delta }{detg}\right|\frac{1}{detg_0}vol(\mathrm{\Sigma }),$$
where
$$\mathrm{\Sigma }:=\{x:dist(x,V(1))n\lambda _n\}.$$
Now, for $`\lambda _nr`$,
$$vol(\mathrm{\Sigma })\lambda _n,$$
where the constant implied in the $`\mathrm{`}\mathrm{`}\mathrm{`}\mathrm{`}`$-notation depends on $`n`$ only (this is obvious for $`\lambda _n\frac{1}{2n}`$, and trivial otherwise, since for $`\lambda _nr`$, $`vol(\mathrm{\Sigma })=O(1)`$). Thus,
$$\begin{array}{cc}\hfill |& N(g_0,\mathrm{\hspace{0.17em}1})\frac{2^{n1}(2\tau \tau ^1)\delta }{detg}|\frac{\lambda _n}{detg_0}\frac{1}{det\mathrm{\Lambda }_{n1}}\hfill \\ & =\frac{1}{vol(M_{n1}/M_{n1}g_0^n)}\frac{\delta }{vol(A_\delta M_{n1}/A_\delta M_{n1}g^n)}\hfill \end{array}$$
Next, suppose that $`\lambda _n>r`$. Then,
$$V(\delta )g_0^nV(\delta )\mathrm{\Lambda }_{n1}.$$
Thus, by the induction hypothesis, the number of such points is:
$$\begin{array}{cc}\hfill & c_4\underset{i=0}{\overset{k1}{}}\frac{1}{det(\mathrm{\Lambda }_i)}=\underset{i=0}{\overset{k1}{}}\frac{1}{vol(M_i/M_ig_0^n)}\hfill \\ & \delta \underset{i=0}{\overset{k1}{}}\frac{1}{vol(A_\delta M_i/A_\delta _ig^n)}.\hfill \end{array}$$
Since $`\lambda _n>r`$, we have
$$\frac{1}{detg}=\frac{1}{\lambda _n}\frac{1}{detg/\lambda _n}\frac{1}{detg/\lambda _n}\frac{1}{\lambda _1\mathrm{}\lambda _{n1}},$$
and we’re done by defining $`L_i:=A_\delta M_i`$. ∎
## 4 Unsmoothing
### 4.1 An asymptotic formula for $`N_\mathrm{\Lambda }`$
We need an asymptotic formula for the sharp counting function $`N_\mathrm{\Lambda }`$. Unlike the case of the standard lattice, $`^2`$, in order to have a good control over the error terms we should use some Diophantine properties of the lattice we are working with. We adapt the following notations:
Let $`\mathrm{\Lambda }`$ be a lattice and $`t>0`$ a real variable. Denote the set of squared norms of $`\mathrm{\Lambda }`$ by
$$SN_\mathrm{\Lambda }=\{|\stackrel{}{n}|^2:n\mathrm{\Lambda }\}.$$
Suppose we have a function $`\delta _\mathrm{\Lambda }:SN_\mathrm{\Lambda }`$, such that given $`\stackrel{}{k}\mathrm{\Lambda }`$, there are no vectors $`\stackrel{}{n}\mathrm{\Lambda }`$ with $`0<||\stackrel{}{n}|^2|\stackrel{}{k}|^2|<\delta _\mathrm{\Lambda }(|\stackrel{}{k}|^2)`$. That is,
$$\mathrm{\Lambda }\{\stackrel{}{n}\mathrm{\Lambda }:|\stackrel{}{k}|^2\delta _\mathrm{\Lambda }(|\stackrel{}{k}|^2)<|\stackrel{}{n}|^2<|\stackrel{}{k}|^2+\delta _\mathrm{\Lambda }(|\stackrel{}{k}|^2)\}=A_{|\stackrel{}{k}|},$$
where
$$A_y:=\{\stackrel{}{n}\mathrm{\Lambda }:|\stackrel{}{n}|=y\}.$$
Extend $`\delta _\mathrm{\Lambda }`$ to $``$ by defining $`\delta _\mathrm{\Lambda }(x):=\delta _\mathrm{\Lambda }(|\stackrel{}{k}|^2)`$, where $`\stackrel{}{k}\mathrm{\Lambda }`$ minimizes $`|x|\stackrel{}{k}|^2|`$ (in the case there is any ambiguity, that is if $`x=\frac{|\stackrel{}{n_1}|^2+|\stackrel{}{n_2}|^2}{2}`$ for vectors $`\stackrel{}{n_1},\stackrel{}{n_2}\mathrm{\Lambda }`$ with consecutive increasing norms, choose $`\stackrel{}{k}:=\stackrel{}{n_1}`$). We have the following lemma:
###### Lemma 4.1.
For every $`a>0,c>1`$,
$$\begin{array}{cc}\hfill N_\mathrm{\Lambda }(t)& =\frac{\pi }{\beta }t^2\frac{\sqrt{t}}{\beta \pi }\underset{\begin{array}{c}\stackrel{}{k}\mathrm{\Lambda }^{}\{0\}\\ |\stackrel{}{k}|\sqrt{N}\end{array}}{}\frac{\mathrm{cos}\left(2\pi t|\stackrel{}{k}|+\frac{\pi }{4}\right)}{|\stackrel{}{k}|^{\frac{3}{2}}}+O(N^a)\hfill \\ & +O\left(\frac{t^{2c1}}{\sqrt{N}}\right)+O(\frac{t}{\sqrt{N}}(\mathrm{log}t+\mathrm{log}(\delta _\mathrm{\Lambda }(t^2)))\hfill \\ & +O\left(\mathrm{log}N+\mathrm{log}(\delta _\mathrm{\Lambda }^{}(t^2))\right)\hfill \end{array}$$
As a typical example of such a function, $`\delta _\mathrm{\Lambda }`$, for $`\mathrm{\Lambda }=1,\alpha +i\beta `$, with a Diophantine $`(\alpha ,\alpha ^2,\gamma ^2)`$, we may choose $`\delta _\mathrm{\Lambda }(y)=\frac{c}{y^K}`$, where $`c`$ is a constant. In this example, if $`\mathrm{\Lambda }\stackrel{}{k}=(a,b)`$, then by lemma 2.4, $`A_{|\stackrel{}{k}|}=\pm (a,b)`$, provided that $`\gamma `$ is irrational.
The proof of this lemma is essentially the same as the one of lemma 5.1 in \[W\], starting from
$$𝒵_\mathrm{\Lambda }(s):=\frac{1}{4}\underset{\stackrel{}{k}\mathrm{\Lambda }0}{}\frac{1}{|\stackrel{}{k}|^{2s}}=\underset{(m,n)_+^20}{}\frac{1}{\left((m+n\alpha )^2+(\beta n)^2\right)^s}$$
###### Proposition 4.2.
Let a lattice $`\mathrm{\Lambda }=1,\alpha +i\beta `$ with a Diophantine triple of numbers $`(\alpha ^2,\alpha \beta ,\beta ^2)`$ be given. Suppose that $`L\mathrm{}`$ as $`T\mathrm{}`$ and choose $`M`$, such that $`L/\sqrt{M}0`$, but $`M=O\left(T^\delta \right)`$ for every $`\delta >0`$ as $`T\mathrm{}`$. Suppose furthermore, that $`M=O(L^{s_0})`$ for some (fixed) $`s_0>0`$. Then
$$\left|S_\mathrm{\Lambda }(t,\rho )\stackrel{~}{S}_{\mathrm{\Lambda },M,L}(t)\right|^2\frac{1}{\sqrt{M}}$$
The proof of proposition 4.2 proceeds along the same lines as the one of proposition 6.1 in \[W\], using again an asymptotic formula for the sharp counting function, given by lemma 4.1. The only difference is that here we use proposition 3.1 rather than lemma 6.2 from \[W\].
Once we have proposition 4.2 in our hands, the proof of our main result, namely, theorem 1.1 proceeds along the same lines as the one of theorem 1.1 in \[W\].
##### Acknowledgement.
This work was supported in part by the EC TMR network Mathematical Aspects of Quantum Chaos, EC-contract no HPRN-CT-2000-00103 and the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities. This work was carried out as part of the author’s PHD thesis at Tel Aviv University, under the supervision of prof. Zeév Rudnick. The author wishes to thank Alex Eskin for his help. A substantial part of this work was done during the author’s visit to the university of Bristol.
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# Red giant branch stars as probes of stellar populations. II. Properties of the newly discovered globular cluster GLIMPSE-C01
## 1 Introduction
The all-sky infrared surveys carried out during the recent years have brought the discovery of a number of new clusters, hidden by the dust extinction in the plane of the Milky Way. They usually suffer A<sub>V</sub>$``$10-20 mag of extinction, making them invisible in the optical wavebands. The vast majority of these objects appear to be a few million years old (Ivanov et al. iva02 (2002), iva05 (2005), Borissova et al. bor03 (2003), bor05 (2005)) but a few have proved to be analogues of “classical” globular clusters (Hurt et al. hur00 (2000), Kobulnicky et al. kob05 (2005), Carraro car05 (2005)).
The globular cluster GLIMPSE-C01 was discovered (Kobulnicky et al. kob05 (2005)) during the Galactic Legacy Infrared Mid-Plane Survey Extraordinaire (hereafter GLIMPSE; Benjamin et al. ben03 (2003))). The survey is mapping the Galactic plane from $``$l$``$ = 10$`\mathrm{°}`$ to 65$`\mathrm{°}`$ and $``$b$``$$``$1$`\mathrm{°}`$ with the Infrared-Array Camera (IRAC; Fazio et al. faz04 (2004)) on Spitzer Space Telescope in 3.6, 4.5, 5.8, and 8.0 micron bands. It was also identified independently by Simpson & Cotera (sim04 (2004)) who cross-correlated ASCA X-ray and IRAS sources with the 2 MASS.
The isochrone comparison indicated that GLIMPSE-C01 is indeed a globular cluster, at least a few Gyr old. The near infrared color-magnitude diagram (CMD) suggested that the object suffers A<sub>V</sub>=15$`\pm `$3 mag of extinction, and the analysis of the line-of-sight <sup>13</sup>CO yielded a distance of 3.1-5.2 Kpc, based on a cinematic model of the Milky Way. The question of the cluster abundance remained open, as well as the need to verify the distance.
Here we report properties – metallicity, extinction, and distance – derived from deep near infrared photometry of GLIMPSE-C01. We also estimate the possible number of undiscovered globular clusters in the central region of the Milky Way.
## 2 Observations and data reduction
The $`JHK_S`$ imaging observations of GLIMPSE-C01 were carried out in Nov 2004 under non-photometric condition and 1 arcsec seeing with the SofI (Son of ISAAC) at the NTT. The instrument is equipped with Hawaii HgCdTe 1024$`\times `$1024 detector, with pixel scale of 0.288 arcsec px<sup>-1</sup>.
We took 12 images in each filter, in jittering mode with 3 arcmin jitter box size to ensure that on each image the cluster is located on a different place on the array. Each image was the average of 3 frames of 10 sec for $`J`$, 13 frames of 3 sec for $`H`$, and 20 frames of 3 sec for $`K_S`$, comprising total integration times of 6, 7.8, and 12 min, respectively. Individual images were sky-subtracted, flat-fielded, aligned, and combined into a single image.
The stellar photometry of the final images was carried out with ALLSTAR in DAOPHOT II (Stetson ste93 (1993)). The typical photometric errors vary from 0.03 mag for the K<sub>S</sub>$``$10 mag stars to 0.10 mag for K<sub>S</sub>$``$13 and 0.15 mag for K<sub>S</sub>$``$16 mag. The photometry of the the faint stars is affected by confusion, especially near the cluster center but our metallicity estimates is based on stars with K<sub>S</sub>$``$13 mag.
The photometric calibration was performed by comparing our instrumental magnitudes with the 2 MASS measurements of 67 stars, covering the color ranges 0.61$``$$`J`$$``$$`K_S`$$``$5.98 and 0.21$``$$`H`$$``$$`K_S`$$``$2.32 mag, and magnitude range 8.33$``$$`K_S`$$``$13.13 mag. The transformation equations are:
$$JK_S=(0.981\pm 0.005)\times (jk)+(3.208\pm 0.007),$$
(1)
$$HK_S=(0.971\pm 0.016)\times (hk)+(1.820\pm 0.014),$$
(2)
and
$$K_S=(1.008\pm 0.007)\times k+(0.0126\pm 0.006)\times (JK_S)(10.096\pm 0.131),$$
(3)
with r.m.s.=0.047, 0.062, and 0.044 mag, respectively. Here $`k`$, $`j`$$``$$`k`$, $`h`$$``$$`k`$ are the instrumental magnitudes and colors and $`K_S`$, $`J`$$``$$`K_S`$, $`H`$$``$$`K_S`$ are the magnitudes and colors in the 2 MASS system. The standard error values are given after each coefficient.
A true-color composite of GLIMPSE-C01 is shown in Figure 1. The image indicates substantial variations of the extinction in the field. Dust structures are present even across the cluster face, suggesting that the completeness of the photometry varies significantly. The crowding also contributes to this effect, as it can be seen from the luminosity function of stars in the field of GLIMPSE-C01 (Figure 2): it is clear that the completeness of the photometry in the cluster area (long-dashed line) is at least a magnitude shallowed than that in the field (dotted line).
## 3 Properties of GLIMPSE-C01
### 3.1 Structural parameters
The structural parameters of GLIMPSE-C01 were determined using the iterative star counts method of King (kin62 (1962)) after randomly removing the field stars (see Sec. 3.3). The radial profile was built only from the stars in the South-Eastern half of the GLIMPSE-C01 to exclude the effects of the variable extinction across the cluster face. The best fit to the radial surface brightness profile with a single-mass King’s model (Figure 3) yielded core radius r<sub>c</sub>=0.78 arcmin, tidal radius r<sub>t</sub>=27 arcmin, and central concentration c=1.54. The inhomogeneous foreground extinction – significant even in the near-infrared – makes it impossible to confirm the suggestion of Kobulnicky et al. (kob05 (2005)) that GLIMPSE-C01 is elongated but we certainly can not exclude such a possibility.
### 3.2 Distance and extinction
The CMD of GLIMPSE-C01 is plotted in Figure 4. The left panel contains all stars with $`J`$ and $`K_S`$ photometry, and shows the presence of a few sequences. The nature of these sequences become apparent from the other two panels: the cluster (middle panel) shows a red giant branch (hereafter RGB) at $`J`$$``$$`K_S`$$``$3-4 mag, and the field (right panel) is dominated by main sequence at $`J`$$``$$`K_S`$$``$1-2 mag and a group of red stars with $`J`$$``$$`K_S`$$``$3 mag, possible just highly reddened background objects.
An inspection of the CMD reveals the presence of a clump of stars in the red giant at $`K_S`$$``$13.0 mag. It is also noticeable on the cluster luminosity function (Figure 2). Taking into account the width of the structure along the $`K_S`$ axis, we assign an uncertainty of 0.25 mag (equal to the bin size used to build the luminosity function) to the apparent magnitude of the clump. If this structure is indeed the red clump, it has an absolute magnitude of $`M_K`$=$``$1.61$`\pm `$0.03 mag (Alves alv00 (2000)), yielding a distance modulus $`(m`$$``$$`M)_K`$$``$14.6$`\pm `$0.3 mag. Note that this calibration doesn’t take into account any metallicity effects (see for discussion Salaris & Girardi sal02 (2002)).
The distance to the cluster can also be measured using the RGB tip (i.e. Ivanov & Borissova iva02 (2002)) but this method suffers from strong metallicity dependence and it is hampered by the small number of stars at the upper end of the RGB. Nevertheless, we carried out this test as a consistency check. If the tip is located at $`K_S`$$``$8.7 mag, and assuming $`M_K`$=$``$6.1 mag for \[Fe/H\]=$``$1.5 (Ivanov & Borissova iva02 (2002)) we obtain $`(m`$$``$$`M)_K`$$``$14.8 mag, in agreement with the estimate given above.
The color-color diagram of the cluster field is shown in Figure 5. Our data confirm the visual extinction A<sub>V</sub>$``$15$`\pm `$3 mag derived by Kobulnicky et al. (kob05 (2005)). Throughout this paper we used the reddening law of Rieke and Lebofsky (rie85 (1985)), giving A<sub>K</sub>$``$1.7$`\pm `$0.3 mag. This yields reddening-free distance modulus of $`(m`$$``$$`M)_0`$$``$12.9$`\pm `$0.4 mag, or D$``$3.8$`\pm `$0.7 kpc, in agreement with Kobulnicky et al. (kob05 (2005)). The CMD shows a notable spread of colors among the fainter stars that might be contributed to differential extinction and to contamination from extended emission and background sources.
### 3.3 Metal abundance
The RGB slope allows to derive abundances of globular clusters (i.e. De Costa & Armandroff dac90 (1990)) because of the metallicity-dependent opacities in cool stars. The RGB slope is independent of the reddening, which is an important advantage in studies of heavily obscured clusters such as GLIMPSE-C01. Here we applied the analysis and the calibrations, described in Ivanov & Borissova (iva02a (2002)).
The first step was to remove the fore- and background stars, which constituted $``$8-9% of the stars in the designated cluster area – a circle with 40 arcsec diameter. We used a Monte-Carlo technique: the CMD of the cluster was divided into rectangular bins. Similarly, the CMD of a circular annulus, centered on the cluster, with an inner radius 60 arcsec, and the same are as the cluster region. Then, from each bin of the cluster CMD we removed randomly a number of star, equal to the number of stars in the corresponding bin of the “field” CMD, producing a single realization. Next, we carried out a least square fit to the RGB (Ivanov & Borissova iva02a (2002); see also Appendix A in this work):
$$(JK_S)=ZP+Slope\times K_S$$
(4)
on the corrected diagram to derive the slope and the zero point. The RG locus was defined after inspecting the CMD: 8.5$``$$`K`$<sub>S</sub>$``$13.0 mag, and 2.5$``$$`J`$$``$$`K`$<sub>S</sub>$``$4.0 mag. The bin-size was 0.5 mag along both axis. We removed the 10-$`\sigma `$ outliers, and repeated the fitting. Typically, the fitting coefficients obtained in the two iterations were statistically indistinguishable. Note that the spread of colors along the RGB is dominated by differential extinction, rather than measurement errors. Total of 10<sup>4</sup> realizations were obtained, and we averaged the fitting results. The distributions of the RGB slope and the zero point of the fits are shown in Figure 6.
Finally, we derived the metal abundance from the calibration of Ivanov & Borissova (iva02a (2002)), and obtained \[Fe/H\]=$``$1.61$`\pm `$0.14 in the scale of Zinn (as implemented in Harris har96 (1996)), \[Fe/H\]=$``$1.44$`\pm `$0.12 in the scale of Caretta & Gratton (car97 (1997)), and \[Fe/H\]=$``$1.12$`\pm `$0.12 in the scale of Ferraro et al. (fer99 (1999)). The calibrations of Valenti, Ferraro, & Origlia (val04 (2004)) wield \[Fe/H\]=$``$1.14$`\pm `$0.16 and \[Fe/H\]=$``$0.97$`\pm `$0.15, for the last two scales. Both values are systematically lower than ours but the differences are within 1-2 $`\sigma `$.
These measurements rely on the important assumption that GLIMPSE-C01 has the same age as the other Milky Way globular clusters, used to derive the metallicity versus the RGB slope calibrations. The comparison of the cluster CMD with stellar isochrones (Kobulnicky et al. kob05 (2005)) seems to indicate an age of at least 8 Gyr but this issue can not be addressed until more accurate age estimate becomes available.
## 4 The case of the missing clusters
The discoveries of GLIMPSE-C01 and Whitting-1 prompted us to reinvestigate the question if there are any more undiscovered globular clusters in the Milky Way. The total number of Milky Way globulars estimated by van den Bergh (van98 (1998)) is 160$`\pm `$20, slightly above the currently known $``$150 objects. Barbuy, Bica, & Ortolani (bar98 (1998); see their Figure 3 and 4) argued that some missing clusters are probably located in the general direction of the Galactic Center. Here we adopt a modification of their method. The XYZ Galactic coordinates used in this section are taken from Feb 2003 edition of the Milky Way Globular Cluster Catalog (Harris har96 (1996)) with their face value. A possible caveat is the variation of the ratio of the total extinction to the reddening R<sub>V</sub>. For example, a change or R<sub>V</sub> from 3 to 3.3 introduces 10% change in the distances, or almost 1 kpc at the distance of the Galactic Center. The averaging over different dust properties along the line of sight helps to minimize the related distance uncertainties but the variations of R<sub>V</sub> toward the bulge definitely need further investigation.
First, we assumed that all missing clusters are located in the Galactic plane and close to the Galactic Center. Therefore, we consider only the region with $``$Z$``$$``$0.5 kpc and R<sub>GC</sub>$``$3.0 kpc. These constraints reflect the spatial distribution of the obscuring material. They make our estimate of the missing clusters only a lower limit, as it was demonstrated by the recent discovery of the off-the-plane cluster Whiting-1.
An observer at the Galactic Center should detect an equal number of globulars toward the Sun, in the antisolar direction, and in the directions perpendicular to the line, connecting the Sun and the Galactic Center, i.e. the globular cluster distribution along any direction in the Galactic plane should be flat. However, it show some structure (Figure 7): there are more clusters in the directions toward the Sun and in the opposite direction than in the direction perpendicular to the line connected the Galactic Center and the Sun. Ten additional clusters are necessary to flatten the histogram (indicated by the shaded area), setting a lower limit of the missing clusters.
This estimate has to be treated with caution for a number of reasons. For example, it is sensitive to the adopted globular cluster distance, and may contain systematic errors related to the adopted globular cluster scale (and the Galactic Center distance). In addition, while the location of the clusters on the sky are known very well, their distances contain statistical observational errors, artificially extending the distribution along the line of sight, similarly to the “finger of God” effect known to the extragalactic astronomers. This would move some of the clusters from the middle two bins to the outer bins, increasing estimate.
Finally, the spatial location of the clusters may be affected by the presence of a bar or a triaxial bulge in our Galaxy. The major axis of the bar has an angle of only about 36$`\pm `$10 deg from the line of sight toward the Galactic center (Weinberg wei92 (1992)), with the near end lying in the first quadrant. The major axis of the bar is close to the direction in which the cluster distribution is elongated.
Our result suggests that the searches for hidden clusters carried out so far are biased toward the regions closer to the Galactic Center so the further away from the Galactic Center is a region, the less likely it is to get attention. This implies that the future searches has to modify their strategy, including regions in the Galactic plane, as far as 7-15 deg from the Galactic Center.
## 5 Summary
We report deep near infrared photometry of the newly discovered Galactic globular Cluster GLIMPSE-C01, and we derived for the first time the metal abundance of this object from the slope of the RGB. The cluster appears metal-poor, with \[Fe/H\]=$``$1.6 in the scale of Zinn. We confirm the distance and reddening estimates of Kobulnicky et al. (kob05 (2005)), placing the cluster at D$``$3.7$`\pm `$0.8 kpc, behind A<sub>V</sub>$``$15 mag of visual extinction.
Finally, we estimate the number of the missing clusters in the central region of the Milky Way. Based on the location of the known clusters, and assuming radial symmetry of the cluster distribution around the Galactic center, we conclude that the Milky Way contains at least 10$`\pm `$3 undiscovered globular clusters.
###### Acknowledgements.
This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. This research has made use of the SIMBAD database, operated at CDS, Strasbourg, France. The authors thank Dr. Ortolani for the useful comments.
## Appendix A Red Giant Slope – Metallicity calibration
The equations given bellow superseded Eqns. 1-4 from Ivanov & Borissova (iva02 (2002)) which contained an error.
$$(JK_S)_0=RGB_{ZP}+RGB_{Sl}\times M_{K_S}$$
(5)
$$RGB_{Sl}=a_0^{Sl}+a_1^{Sl}\times [\mathrm{Fe}/\mathrm{H}]$$
(6)
$$RGB_{ZP}=a_0^{ZP}+a_1^{ZP}\times [\mathrm{Fe}/\mathrm{H}]$$
(7)
$$[\mathrm{Fe}/\mathrm{H}]=[(JK_S)_0a_0^{Sl}\times M_{K_S}a_0^{ZP}]/[a_1^{Sp}\times M_{K_S}+a_1^{ZP}]$$
(8)
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# Teleportation via multi-qubit channels
## 1 Introduction
In recent times entanglement has come to be recognised as one of the major distinguishing features between quantum systems and classical systems, where it is now seen as being as fundamental as the uncertainty principle. This point of view has arisen due to the realisation that entanglement is a resource to be exploited in the processing of quantum information through processes such as teleportation , dense-coding and quantum cryptography . It has also opened new perspectives in other areas such as condensed matter physics, due the to emerging understanding of the relationship between entanglement and quantum critical phenomena . As a consequence there has been an intense level of activity in characterising entanglement and studying its properties.
At the level of bi-partite systems entanglement is well understood and can be quantified . From studies of three-qubit systems it was realised that different categories of entanglement exist in multi-qubit systems, with the specific example of three-way entanglement shown to be essentially different from bi-partite entanglement through the examples of the Greenberger-Horne-Zeilinger (GHZ) and W states. Now, a clear picture of three-qubit entanglement has emerged with the demonstration of three different types of entanglement existing in the three-qubit case, which are characterised by five generically independent invariants . Though the above results for three-qubit systems can in principle be generalised to arbitrary multi-qubit systems, it is technically challenging to undertake. Despite many studies of specific types of multi-qubit entanglement (e.g. ), a complete description remains elusive.
Our aim in this work is to investigate entanglement in multi-qubit systems through a study of one of its applications, viz. teleportation. The protocol for the procedure is as follows, with a schematic representation shown in Fig. (1). An unknown qubit state is held by a client (Alice), and is to be teleported to a recipient (Bob). Alice and Bob share a quantum channel which is some state of $`2`$ qubits, so the entire system consists of $`2+1`$ qubits. The channel is distributed in such a way that Alice may access $`21`$ qubits of the channel while Bob has access to a single qubit of the channel. Alice is to perform $``$ Bell state measurements on the $`2`$-qubit subsystem which is comprised of her unknown state and $`21`$ qubits of the channel. The consequence of this measurement is that Bob is left with a single qubit which is not entangled with the remainder of the system. From the results of the measurements Alice is to send classical information to Bob. Upon receiving this information, and some knowledge of the channel, Bob determines a local unitary operation called a correction gate which he applies to his qubit. Any channel for which this procedure exactly reproduces the client state for Bob (i.e. the teleportation is effected with perfect fidelity) we will call a perfect channel. One of the aims of this work is to determine the complete set of perfect channels for this protocol. We mention that this protocol is not tight in the sense of , and consequently does not belong to the classification of teleportation schemes given therein. On the other hand it does bear similarity to the quantum repeater described in , with the major difference being that we employ Bell measurements whereas local measurements are used in .
It is well known that teleportation can be performed with perfect fidelity across a 2-qubit channel when the channel is one of the four Bell states . This is achieved by making a Bell state measurement and then sending two bits of information to the recipient via a classical channel, which is then used to determine the correction gate. It is thus clear that teleportation can also be achieved with perfect fidelity using a channel which is a product of $``$ Bell states, by performing $``$ successive Bell state measurements. However this is not the most general solution to the problem we have described above, and our analysis below shows some surprising results. The first is that there exist four orthogonal perfect channel subspaces, the direct sum of which is the entire Hilbert space of channels. Also, despite the fact that $``$ Bell measurements need to be performed by Alice to implement the teleportation, only two bits of classical information need to be sent from Alice to Bob for him to determine the correction gate.
Because perfect channels fall into one of only four subspaces of the channel state space, an interesting question to consider is whether the ground states of common many-body systems fall into these classes. This is indeed the case. For example, our results indicate that all spin singlet states are perfect channels, and so the ground state of the antiferromagnetic Heisenberg model, for a number of different lattices, is a perfect channel, as is the ground state of the one-dimensional Majumdar–Ghosh model . It is also true that the ground state of the model of Affleck, Kennedy, Lieb and Tasaki (AKLT) is a perfect channel, under an equivalent protocol . In identifying the perfect channel states we determine a teleportation-order parameter which provides a measure of the effectiveness of an arbitrary channel. The teleportation-order parameter has close connection with string-order , as discussed in in relation to the AKLT model, and is also an example of the string operators discussed in . We mention however that the results of our analysis are independent of the dimension and topology of the lattice on which the qubits are arranged.
We will also show that this analysis extends to formulate a teleportation protocol for the case of 3-qubit channels, and that there is a generalisation for qudits. Finally, we will discuss some implications of these results towards understanding entanglement in multi-qubit systems. The results presented here describe in detail the mathematical aspects which underly the results reported in .
## 2 Teleportation via two-qubit channels
In this section we recall teleportation across a channel of two qubits which exists in one of the four Bell states . While this phenomenon is now well known, the notational conventions we adopt, which are convenient for the following sections, are not standard.
Let $`|+,|`$ denote the standard basis for a qubit space $`V`$ such that $`|j`$ is an eigenvector of the Pauli matrix $`\sigma ^z`$ with eigenvalue $`j`$. Throughout, we will label $`\pm 1`$ simply by $`\pm `$. A natural basis for two coupled qubits is $`|j,k|j|k`$. In making a basis change to the Bell states we define
$`|+:+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,++|,\right)`$
$`|+:\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+|,+\right)`$
$`|:+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+|,\right)`$
$`|:\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,|,+\right)`$
such that we can write
$`|j:k\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,k+j|,\overline{k}\right)`$ (1)
where we adopt the notation $`\overline{k}=k`$. It is known that each of the Bell basis states are related by a local unitary transformation, which we express as
$`|j:k\}=(IX_{pq}^{jk})|p:q\}`$
where
$`X_{jk}^{jk}=U^0`$
$`X_{++}^+=X_+^{++}=X_{}^+=X_+^{}=U^1`$
$`X_{++}^+=X_+^{++}=X_{}^+=X_+^{}=U^2`$
$`X_{++}^{}=X_+^+=X_{}^{++}=X_+^+=U^3,`$ (2)
and the unitary operators $`U^i`$ are given by
$`U^0`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$
$`U^1`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$
$`U^2`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$
$`U^3`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$
(Note that the $`X_{pq}^{jk}`$ could just as easily have been defined in terms of Pauli matrices. We generally prefer to not use Pauli matrix notation, as this eliminates $`\sqrt{1}`$ terms which would otherwise appear in many subsequent formulae.) The above is just a statement of the fact that the Bell states are equivalent: two states are said to be equivalent if they are equal up to a tensor product of local unitary transformations. Equivalent states have identical entanglement properties. Likewise, we say that two subspaces $`Y,Z`$ with the same (finite) dimension are equivalent if and only if for a fixed transformation, each $`yY`$ is equivalent to some $`zZ`$. Moreover two operators are equivalent if they are similar by a transformation which is a tensor product of local unitary transformations.
Define $`\nu :\left\{\pm 1\right\}_2`$ by
$`\nu (+1)=0,`$ $`\nu (1)=1`$
which satisfies $`\nu (ab)=\nu (a)+\nu (b)`$. We can express the relations (2) as
$`X_{pq}^{jk}`$ $`=`$ $`\delta _{pq}^{jk}(U^1)^{\nu (kq)}(U^2)^{\nu (jp)}`$ (3)
where $`\delta _{pq}^{jk}=\pm 1`$ can be read off from (2). The operators $`X_{pq}^{jk}`$ satisfy the following properties:
$`X_{pq}^{jk}`$ $`=`$ $`\epsilon _{pq}^{jk}X_{jk}^{pq}=(X_{jk}^{pq})^{}`$ (4)
$`X_{pq}^{jk}X_{ab}^{pq}`$ $`=`$ $`X_{ab}^{jk}`$ (5)
$`X_{pq}^{jk}X_{cd}^{ab}`$ $`=`$ $`\epsilon _{pq}^{jk}\epsilon _{jb}^{pq}X_{pq}^{jb}X_{cd}^{ak}`$ (6)
$`X_{pq}^{jk}X_{cd}^{ab}`$ $`=`$ $`\delta _{jk}^{ak}\delta _{jb}^{ab}\epsilon _{pq}^{jk}\epsilon _{ak}^{pq}X_{pq}^{ak}X_{cd}^{jb}`$ (7)
$`X_{pq}^{jk}X_{cd}^{ab}`$ $`=`$ $`\delta _{pq}^{jk}\delta _{cd}^{ab}\delta _{(pc)(qd)}^{(ja)(kb)}\epsilon _{pd}^{jb}X_{(pc)(qd)}^{(ja)(kb)}`$ (8)
where $`\epsilon _{pq}^{jk}=\epsilon _{jk}^{pq}`$ is defined by
$$\epsilon _{pq}^{jk}=\{\begin{array}{cc}1\hfill & \text{ for }jp\text{ and }kq\hfill \\ 1\hfill & \text{ otherwise }\hfill \end{array}$$
(9)
Property (4) is deduced by inspection, while (5) follows from the definition of the $`X_{pq}^{jk}`$. To show (6), first observe that it is true if $`k=b`$. Assuming $`kb`$ and using (4,5) we find
$`X_{pq}^{jk}X_{cd}^{ab}`$ $`=`$ $`\epsilon _{pq}^{jk}X_{jk}^{pq}X_{cd}^{ab}`$
$`=`$ $`\epsilon _{pq}^{jk}X_{jb}^{pq}X_{jk}^{jb}X_{ak}^{ab}X_{cd}^{ak}`$
$`=`$ $`\epsilon _{pq}^{jk}X_{jb}^{pq}U^1U^1X_{cd}^{ak}`$
$`=`$ $`\epsilon _{pq}^{jk}\epsilon _{jb}^{pq}X_{pq}^{jb}X_{cd}^{ak}.`$
Property (7) is proved similarly. To show (8) we calculate
$`X_{pq}^{jk}X_{cd}^{ab}`$ $`=`$ $`\delta _{pq}^{jk}\delta _{cd}^{ab}(U^1)^{\nu (kq)}(U^2)^{\nu (jp)}(U^1)^{\nu (bd)}(U^2)^{\nu (ac)}`$
$`=`$ $`\delta _{pq}^{jk}\delta _{cd}^{ab}\epsilon _{pd}^{jb}(U^1)^{\nu (kqbd)}(U^2)^{\nu (jpac)}`$
$`=`$ $`\delta _{pq}^{jk}\delta _{cd}^{ab}\delta _{(pc)(qd)}^{(ja)(kb)}\epsilon _{pd}^{jb}X_{(pc)(qd)}^{(ja)(kb)}.`$
We also find that
$`U^1U^1|j:k\}`$ $`=`$ $`j|j:k\}`$
$`U^2U^2|j:k\}`$ $`=`$ $`k|j:k\}`$
so the eigenvalues of $`U^1U^1,U^2U^2`$ provide good quantum numbers to label the basis states. (It is easily checked that $`U^1U^1`$ and $`U^2U^2`$ commute.) A measurement which is represented by the action of these two operators is called a Bell measurement, where $`(j:k)`$ denotes the measurement outcome.
Let $`|v=\alpha |++\beta |`$ be arbitrary such that
$$|\alpha |^2+|\beta |^2=1$$
and $`\alpha `$ and $`\beta `$ are completely unknown. We call $`|v`$ the client state. The state which will be used to teleport the client state will be called the channel. When the channel is one of the Bell basis states $`|j:k\}`$ we look to rewrite the total state $`|v|j:k\}`$ as a linear combination of states where the first two qubits are expressed in the Bell basis, i.e.
$`2(II(X_{++}^{jk})^1)|v|j:k\}`$
$`=2|v|+:+\}`$
$`=\sqrt{2}\left(\alpha |+,+,++\beta |,+,++\alpha |+,,+\beta |,,\right)`$
$`=|+:+\}(\alpha |++\beta |)+|+:\}(\beta |++\alpha |)`$
$`+|:+\}(\alpha |+\beta |)+|:\}(\beta |++\alpha |)`$
$`=|+:+\}|v+|+:\}U^1|v+|:+\}U^2|v+|:\}U^3|v`$
$`=|+:+\}X_{++}^{++}|v+|+:\}X_{++}^+|v`$
$`+|:+\}X_{++}^+|v+|:\}X_{++}^{}|v`$
$`=|+:+\}X_{++}^{++}|v+|+:\}X_+^{++}|v`$
$`+|:+\}X_+^{++}|v+|:\}\epsilon _{++}^{}X_{}^{++}|v`$
and so
$`2|v|j:k\}`$ $`=`$ $`|+:+\}X_{++}^{jk}|v+|+:\}X_+^{jk}|v`$
$`+|:+\}X_+^{jk}|v+|:\}\epsilon _{++}^{}X_{}^{jk}|v`$
$`=`$ $`|+:+\}\epsilon _{++}^{++}X_{++}^{jk}|v+|+:\}\epsilon _{++}^+X_+^{jk}|v`$
$`+|:+\}\epsilon _{++}^+X_+^{jk}|v+|:\}\epsilon _{++}^{}X_{}^{jk}|v.`$
This last expression can be expressed in a compact form:
$`|v|j:k\}={\displaystyle \frac{1}{2}}{\displaystyle \underset{p,q}{}}|p:q\}\stackrel{~}{X}_{pq}^{jk}|v`$ (10)
where $`\stackrel{~}{X}_{pq}^{jk}=\epsilon _{++}^{pq}X_{pq}^{jk}`$. Thus, when a Bell measurement is made on the first and second qubits by Alice, the system is projected onto a state
$$|p:q\}\stackrel{~}{X}_{pq}^{jk}|v.$$
Note that the probabilities for measuring each of the four possible states are equal. The result of the measurement may be communicated to Bob using only two bits of classical information. This, together with knowledge of which channel was used, is sufficient information for Bob to determine the correction gate $`\stackrel{~}{X}_{jk}^{pq}`$, to be implemented in order to recover the client state. Thus the client state has been teleported from Alice to Bob via the channel and classical communication.
## 3 Teleportation via multi-qubit channels
### 3.1 Singlet channels: an example of teleportation via multi-qubit channels with perfect fidelity
Our goal is to extend the above construction to the multi-qubit channel case. Here, we will first look at the case when the channel is a $`U(2)`$ singlet. The Hilbert space for an $`L`$-qubit system is given by the tensor product of the local qubit spaces $`V`$;
$$V^LV^L=V_1V_2\mathrm{}V_L.$$
Throughout we take $`L`$ to be even and define $`=L/2`$. Recall that the action of the Lie group $`U(2)`$ on the space of a single qubit space $`V`$ is represented by the set of all $`2\times 2`$ unitary matrices. Given any such matrix $`AU(2)`$, the action extends to the space of $`L`$ qubits through
$$AA^L.$$
A $`U(2)`$ singlet is any state $`|\mathrm{\Psi }V^L`$ such that for all $`AU(2)`$
$`A^L|\mathrm{\Psi }=\mathrm{exp}(i\theta )|\mathrm{\Psi }`$ (11)
for some real $`\theta `$. An example of a singlet is given by the Bell state $`|:\}`$.
Let $`P^{pq}`$ denote the projection onto the Bell state $`|p,q\}`$. Each projection can be related to the projection onto the $`U(2)`$ singlet state $`|:\}`$ through
$`P^{pq}=(IX_{}^{pq})P^{}(IX_{pq}^{}).`$
Now since $`|:\}`$ is a $`U(2)`$ singlet then $`P^{}`$ is an invariant operator in the sense that
$`(AA)P^{}(A^1A^1)=P^{}AU(2),`$
that is, the action of $`U(2)`$ commutes with $`P^{}`$. It follows that $`P^{}I^{L2}`$ is an invariant operator on the $`L`$-fold space $`V^L`$. An important result we will use subsequently is Schur’s lemma, which asserts that any invariant operator maps an irreducible $`U(2)`$ invariant space to an isomorphic space by a scalar multiple .
Let $`P_r^{pq}`$ be the projector $`P^{pq}`$ acting on the $`r`$th and $`(r+1)`$ qubits of the tensor product space and let $`(X_{pq}^{jk})_r`$ be $`X_{pq}^{jk}`$ acting on the $`r`$th space. Let $`|v`$ again be an arbitrary client state, and let the channel $`|\mathrm{\Psi }V^L`$ be an arbitrary singlet state. We denote the space to which the client state belongs by $`V_0`$. The initial state of the total system is thus
$$|v^{(0)}=|v|\mathrm{\Psi }.$$
Now we employ Schur’s lemma, which in particular means that
$$P_0^{}|v|\mathrm{\Psi }=\chi |:\}|v^{(1)}$$
for some scalar $`\chi `$, where $`|v^{(1)}`$ is some state in $`W^{(1)}=V_2V_3\mathrm{}V_L`$ which is isomorphic to $`|v`$. In other words, if we decompose $`W^{(1)}`$ into $`U(2)`$ spaces then $`|v^{(1)}`$ belongs to a doublet.
Starting with $`|v^{(0)}`$, a Bell measurement is made on $`V_0V_1`$, which is denoted by the projection $`P_0^{pq}`$ where $`(p:q)`$ is the result of the measurement. With reference to the above discussions and notational conventions this means we may write
$`P_0^{pq}|v^{(0)}`$ $`=`$ $`P_0^{pq}|v|\mathrm{\Psi }`$
$`=`$ $`(X_{}^{pq})_1P_0^{}(X_{pq}^{})_1|v|\mathrm{\Psi }`$
$`=`$ $`\left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{}^{pq})_r\right)P_0^{}\left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{pq}^{})_r\right)|v|\mathrm{\Psi }`$
$`=`$ $`\left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{}^{pq})_r\right)P_0^{}|v\left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{pq}^{})_r\right)|\mathrm{\Psi }`$
$`=`$ $`e^{i\theta }\left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{}^{pq})_r\right)P_0^{}|v|\mathrm{\Psi }(\mathrm{since}|\mathrm{\Psi }\mathrm{is}\mathrm{a}\mathrm{singlet})`$
$`=`$ $`e^{i\theta }\chi \left({\displaystyle \underset{r=1}{\overset{L}{}}}(X_{}^{pq})_r\right)|:\}|v^{(1)}(\mathrm{by}\mathrm{Schur}^{}\mathrm{s}\mathrm{lemma})`$
$`=`$ $`e^{i\theta }\chi (X_{}^{pq})_1|:\}\left({\displaystyle \underset{r=2}{\overset{L}{}}}(X_{}^{pq})_r\right)|v^{(1)}`$
$`=`$ $`e^{i\theta }\chi |p:q\}\left({\displaystyle \underset{r=2}{\overset{L}{}}}(X_{}^{pq})_r\right)|v^{(1)}.`$
This procedure can be iterated by taking $`l`$ consecutive Bell measurements to give
$$P_0^{p_1q_1}P_2^{p_2q_2}\mathrm{}.P_{2l2}^{p_lq_l}|v^{(0)}=\gamma |p_1:q_1\}|p_2:q_2\}\mathrm{}|p_l:q_l\}\left(\underset{t=l}{\overset{1}{}}\left(\underset{r=2l}{\overset{L}{}}(X_{}^{p_tq_t})_r\right)\right)|v^{(l)}$$
where $`\gamma `$ is a scalar. In particular
$$P_0^{p_1q_1}P_2^{p_2q_2}\mathrm{}.P_{L2}^{p_{}q_{}}|v^{(0)}=\gamma |p_1:q_1\}|p_2:q_2\}\mathrm{}|p_{}:q_{}\}\left(\underset{t=}{\overset{1}{}}(X_{}^{p_tq_t})_L\right)|v^{()}.$$
Note the notation employed means for any operator $`\mathrm{\Xi }`$
$`{\displaystyle \underset{j=k}{\overset{1}{}}}\mathrm{\Xi }_j=\mathrm{\Xi }_k\mathrm{}.\mathrm{\Xi }_2\mathrm{\Xi }_1.`$
In each case $`|v^{(l)}W^{(l)}=V_{2l}\mathrm{}.V_L`$ is isomorphic to $`|v`$ due to Schur’s lemma. However, since $`|v^{()}W^{()}=V_L`$ and $`V_L`$ is an irreducible $`U(2)`$ space, we must have
$`|v^{()}=|v_L.`$
After the $``$ Bell basis measurements are made by Alice, Bob needs to apply the correction gate
$`D={\displaystyle \underset{t=1}{\overset{}{}}}X_{p_tq_t}^{}`$
to the $`L`$th qubit in order to recover the client state. In view of (8) we see that
$`D=\kappa X_{pq}^{jk}`$ (12)
where $`j=k=(1)^{}`$, $`p=_{t=1}^{}p_t`$ and $`q=_{t=1}^{}q_t`$. Note that $`\kappa =\pm 1`$ in (12) is a function of all indices $`p_t`$ and $`q_t`$ and can in principle be determined through (8). However its value is inconsequential, as it will only alter the corrected state by a phase. Throughout, whenever such a phase arises we will generically denote it by $`\kappa `$. For ease of notation, we will not explicitly state its dependence on particular indices, although this should be clear from the context.
After Alice has performed the Bell measurements, she need only send two bits of classical information, viz. $`p`$ and $`q`$, to Bob in order for him to determine the correction gate. This is a case where teleportation occurs with perfect fidelity, and shows that all singlet states are perfect channels. Next we look to extend this result to cover the most general possibilities.
### 3.2 A basis of perfect multi-qubit channels
Our first step to classifying the perfect channels is to establish that there exists a basis for the $`L`$-qubit Hilbert space $`V^L`$ in which each basis state is a perfect channel. Since the Bell states provide a basis for $`VV`$ it immediately follows that the set of all vectors of the form
$`|\stackrel{}{j}:\stackrel{}{k}\}=|j_1:k_1\}\mathrm{}|j_{}:k_{}\}`$ (13)
forms a basis for $`V^L`$. Through repeated use of (10) we arrive at
$`|v|\stackrel{}{j}:\stackrel{}{k}\}={\displaystyle \frac{1}{2^{}}}{\displaystyle \underset{\stackrel{}{p},\stackrel{}{q}}{}}|\stackrel{}{p}:\stackrel{}{q}\}\stackrel{~}{X}_{p_{}q_{}}^{j_{}k_L}\mathrm{}\stackrel{~}{X}_{p_1q_1}^{j_1k_1}|v`$ (14)
where the sum is taken over all possible values of $`\stackrel{}{p}`$ and $`\stackrel{}{q}`$. By Alice making pairwise Bell measurements on the first $`L`$ spaces, a projection is made to a state
$`|\stackrel{}{p}:\stackrel{}{q}\}\stackrel{~}{X}_{p_{}q_{}}^{j_{}k_{}}\mathrm{}\stackrel{~}{X}_{p_1q_1}^{j_1k_1}|v.`$ (15)
Given a basis vector $`|\stackrel{}{j}:\stackrel{}{k}\}`$ we say that it belongs to the Bell class $`[j:k],j,k=\pm `$ if
$`{\displaystyle \underset{i=1}{\overset{}{}}}j_i=j,{\displaystyle \underset{i=1}{\overset{}{}}}k_i=k.`$
There are four distinct Bell classes. Given an arbitrary vector
$`|\mathrm{\Phi }={\displaystyle \underset{\stackrel{}{j},\stackrel{}{k}}{}}\mathrm{\Gamma }_{\stackrel{}{j},\stackrel{}{k}}|\stackrel{}{j}:\stackrel{}{k}\}`$ (16)
we say that $`|\mathrm{\Phi }`$ belongs to the Bell class $`[j:k]`$ if, for $`\mathrm{\Gamma }_{\stackrel{}{j}\stackrel{}{k}}`$ non-zero, then $`|\stackrel{}{j}:\stackrel{}{k}\}[j:k]`$. In other words, $`|\mathrm{\Phi }`$ belongs to the Bell class $`[j:k]`$ if it is a linear combination of basis vectors (13) of Bell class $`[j:k]`$. It is clear that the notion of Bell classes leads to a vector space decomposition
$$V^L=V_{[+:+]}^LV_{[+:]}^LV_{[:+]}^LV_{[:]}^L$$
and we refer to each $`V_{[j:k]}^L`$ as a Bell subspace. In view of (14) we arrive at
$`|v|\mathrm{\Phi }={\displaystyle \frac{1}{2^{}}}{\displaystyle \underset{\stackrel{}{p},\stackrel{}{q},\stackrel{}{j},\stackrel{}{k}}{}}\mathrm{\Gamma }_{\stackrel{}{j}\stackrel{}{k}}|\stackrel{}{p}:\stackrel{}{q}\}\stackrel{~}{X}_{p_{}q_{}}^{j_{}k_L}\mathrm{}\stackrel{~}{X}_{p_1q_1}^{j_1k_1}|v`$ (17)
where the sum is taken over all possible values of $`\stackrel{}{j},\stackrel{}{k},\stackrel{}{p}`$ and $`\stackrel{}{q}`$. Making pairwise measurements on the first $`L`$ spaces then projects out a state
$`𝒩{\displaystyle \underset{\stackrel{}{j},\stackrel{}{k}}{}}\mathrm{\Gamma }_{\stackrel{}{j}\stackrel{}{k}}|\stackrel{}{p}:\stackrel{}{q}\}\stackrel{~}{X}_{p_{}q_{}}^{j_{}k_{}}\mathrm{}\stackrel{~}{X}_{p_1q_1}^{j_1k_1}|v`$ (18)
where $`𝒩`$ is a normalisation factor. Now suppose $`|\mathrm{\Phi }V_{[j:k]}^L`$. Again appealing to (8) and using the fact that $`\mathrm{\Gamma }_{\stackrel{}{j}\stackrel{}{k}}=0`$ for $`|\stackrel{}{j}:\stackrel{}{k}\}[j:k]`$, then up to a phase we can express (18) as
$`|\stackrel{}{p}:\stackrel{}{q}\}\stackrel{~}{X}_{pq}^{jk}|v`$
where $`[j:k]`$ is the Bell class of the channel, and $`[p:q]`$ is the Bell class of the measurement (more precisely, the Bell class of the measured tensor product of Bell states). As in the case of singlet channels, Alice again just needs to communicate the Bell class of her measurement (i.e. two bits of classical information) to Bob for him to determine the correction gate. Here we assume, as in the case of teleportation across a single Bell state, that the Bell class of the channel is known to Bob.
A characteristic of the states of Bell subspaces is that they are simultaneous eigenvectors of the operators
$`\mathrm{{\rm Y}}^1`$ $`=`$ $`{\displaystyle \underset{p=1}{\overset{L}{}}}U_p^1,\mathrm{{\rm Y}}^2={\displaystyle \underset{q=1}{\overset{L}{}}}U_q^2,`$
and therefore an eigenstate of the product
$`\mathrm{{\rm Y}}^3`$ $`=`$ $`\left({\displaystyle \underset{p=1}{\overset{L}{}}}U_p^1\right)\left({\displaystyle \underset{q=1}{\overset{L}{}}}U_q^2\right)={\displaystyle \underset{p=1}{\overset{L}{}}}U_p^3.`$
Note also that
$`[\mathrm{{\rm Y}}^\alpha ,\mathrm{{\rm Y}}^\beta ]=0\alpha ,\beta =1,2,3.`$
It is apparent that each space $`V_{[j:k]}^L`$ is a stabiliser space for the set of operators $`\{j\mathrm{{\rm Y}}^1,k\mathrm{{\rm Y}}^2\}`$. In fact our protocol can be re-expressed as a multi-qubit generalisation of the stabiliser description of teleportation given in . Another result that can be immediately deduced from the above is that any perfect channel $`|\phi V_{[j:k]}^L`$ is maximally locally disordered; i.e.,
$`\phi \left|U_p^\alpha \right|\phi =0\alpha =1,2,3,p=1,\mathrm{},L.`$ (19)
The result follows from the fact that $`|\phi `$ is an eigenstate of each $`\mathrm{{\rm Y}}^\alpha `$, the $`\mathrm{{\rm Y}}^\alpha `$ are self-adjoint, and
$`\mathrm{{\rm Y}}^\alpha U_p^\beta =U_p^\beta \mathrm{{\rm Y}}^\alpha \mathrm{for}\alpha \beta .`$ (20)
Note that (20) also shows that the Bell subspaces are equivalent.
Let $`𝒫_{mn}`$ denote the permutation operator which permutes the $`m`$th and $`n`$th qubits of the tensor product space $`V^L`$. These operators provide a representation of the symmetric group. Since the $`𝒫_{mn}`$ commute with the $`\mathrm{{\rm Y}}^\alpha `$, it follows that each of the subspaces $`V_{[j:k]}^L`$ is invariant under the action of the symmetric group. Thus given any perfect channel, it can be used by Alice and Bob to achieve unit fidelity teleportation independent of which qubit of the channel is Bob’s, how Alice chooses to pair the qubits in making the Bell basis measurements, and the order in which she performs the measurements. In particular, it is not necessary that her first measurement involves the client state.
We mention here that unlike the $`=1`$ case, the probabilities for the measurements that may be made by Alice are not necessarily equal in the case of general $``$. An illustration of this fact is given by the example in the Appendix. It is true however that the probability a measurement made by Alice is of the Bell class $`[j:k]`$ is always 1/4, independent of $`j,k`$ or $`L`$. To show this we first construct projection operators $`P_{[j:k]}`$ onto the subspaces $`V_{[j:k]}^L`$ by
$`P_{[j:k]}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(I+j\mathrm{{\rm Y}}^1)(I+k\mathrm{{\rm Y}}^2)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(I+j\mathrm{{\rm Y}}^1+k\mathrm{{\rm Y}}^2+jk\mathrm{{\rm Y}}^3\right).`$
Given any perfect channel $`|\phi _{[p:q]}`$ of Bell class $`[p:q]`$ and client state $`|v`$ the density matrix is
$`\rho =|vv||\phi _{[p:q]}\phi _{[p:q]}|.`$
The probability $`𝒫_{[j:k]}`$ that Alice makes a measurement of Bell class $`[j:k]`$ is given by
$`𝒫_{[j:k]}`$ $`=`$ $`\mathrm{tr}[(P_{[j:k]}I)\rho ]`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{tr}[\rho ]+\mathrm{tr}[j(\mathrm{{\rm Y}}^1I)\rho ]+\mathrm{tr}[k(\mathrm{{\rm Y}}^2I)\rho ]+\mathrm{tr}[jk(\mathrm{{\rm Y}}^3I)\rho ]\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1+\mathrm{tr}[j(\mathrm{{\rm Y}}^1U^1)\rho U_L^1]+\mathrm{tr}[k(\mathrm{{\rm Y}}^2U^2)\rho U_L^2]\mathrm{tr}[jk(\mathrm{{\rm Y}}^3U^3)\rho U_L^3]\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1+\mathrm{tr}[j(U^1\mathrm{{\rm Y}}^1)\rho U_L^1]+\mathrm{tr}[k(U^2\mathrm{{\rm Y}}^2)\rho U_L^2]\mathrm{tr}[jk(U^3\mathrm{{\rm Y}}^3)\rho U_L^3]\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1+\mathrm{tr}[jpU_1^1\rho U_L^1]+\mathrm{tr}[kqU_1^2\rho U_L^2]\mathrm{tr}[jkpqU_1^3\rho U_L^3]\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}(1+jpv|U_1^1|v\phi _{[p:q]}|U_L^1|\phi _{[p:q]}`$
$`+kqv|U_1^2|v\phi _{[p:q]}|U_L^2|\phi _{[p:q]}jkpqv|U_1^3|v\phi _{[p:q]}|U_L^3|\phi _{[p:q]})`$
$`=`$ $`{\displaystyle \frac{1}{4}}`$
where the last line follows from (19).
It is of interest to consider how the above results relate to common physical models. It was shown earlier that all singlet states are perfect channels, and since they form a subspace (for fixed $`L`$), they must belong to the same Bell class. The ground state of the antiferromagnetic Heisenberg model is a singlet, as is the ground state of the one-dimensional Majumdar-Ghosh model , so each is a perfect channel. In one-dimension the Heisenberg model is gapless, so any physical realisation would be susceptible to errors arising from thermal fluctuations. One way to reduce errors is to use a gapped system, which is the case for the Heisenberg model on a two-leg ladder lattice as well as the one-dimensional Majumdar-Ghosh model. For the Heisenberg model on the two-dimensional Kagome lattice the system is gapless, but the elementary gapless excitations are also singlets . In this instance error due to thermal fluctuation is again reduced since all singlet states belong to the same Bell class. The existence of such subspaces for which all states provide perfect fidelity teleportation is reminiscent of decoherence free subspaces used to encode logical qubits which are immune to decoherence effects .
### 3.3 Cluster states
Cluster states were introduced in as examples of multi-qubit states with maximal connectedness and high persistency of entanglement. The utilisation of these states for one-way quantum computation has been studied in . We will indicate here how each of the one-dimensional cluster states is equivalent to a particular Bell class state.
The one-dimensional cluster states may be defined as the $`L`$-qubit states
$`|\varphi ^{(L)}={\displaystyle \frac{1}{2^{}}}\left[{\displaystyle \underset{j=1}{\overset{L1}{}}}\left(|+_j+|_jU_{j+1}^2\right)\right]\left(|+_L+|_L\right).`$ (21)
Consider the set of operators $`K_j`$ defined by
$`K_1`$ $`=`$ $`U_1^1U_2^2,`$
$`K_j`$ $`=`$ $`U_{j1}^2U_j^1U_{j+1}^2,j=2,\mathrm{},L1,`$
$`K_L`$ $`=`$ $`U_{L1}^2U_L^1.`$ (22)
It is straightforward to verify these operators satisfy
$`[K_j,K_l]`$ $`=`$ $`0j,l=1,\mathrm{},L,`$
$`K_j|\varphi ^{(L)}`$ $`=`$ $`|\varphi ^{(L)}j=1,\mathrm{},L.`$
Define the operators $`G_1`$ and $`G_2`$ by
$`G_1`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{/2}{}}}(K_{4j3}K_{4j}),`$
$`G_2`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{/2}{}}}(K_{4j2}K_{4j1})`$
for $`/2`$ even and
$`G_1`$ $`=`$ $`K_{21}{\displaystyle \underset{j=1}{\overset{(1)/2}{}}}(K_{4j3}K_{4j}),`$
$`G_2`$ $`=`$ $`K_2{\displaystyle \underset{j=1}{\overset{(1)/2}{}}}(K_{4j2}K_{4j1})`$
for when $`/2`$ is odd. The operators $`G_1`$ and $`G_2`$ necessarily commute and moreover
$`G^1|\varphi ^{(L)}`$ $`=`$ $`|\varphi ^{(L)},`$
$`G^2|\varphi ^{(L)}`$ $`=`$ $`|\varphi ^{(L)}.`$
It can be shown that $`G^1,G^2`$ are equivalent to $`\mathrm{{\rm Y}}^1,\mathrm{{\rm Y}}^2.`$ We will not give a detailed proof, but rather illustrate some examples. When $`L=6`$ we have
$`G^1`$ $`=`$ $`U_1^1U_2^2U_3^2U_4^1U_5^2U_4^2U_5^1U_6^2`$
$`=`$ $`U_1^1U_2^2U_3^2U_4^3U_5^3U_6^2,`$
$`G^2`$ $`=`$ $`U_1^2U_2^1U_3^2U_2^2U_3^1U_4^2U_5^2U_6^1`$
$`=`$ $`U_1^2U_2^3U_3^3U_4^2U_5^2U_6^1`$
whilst in the case $`L=8`$ we have
$`G^1`$ $`=`$ $`U_1^1U_2^2U_3^2U_4^1U_5^2U_4^2U_5^1U_6^2U_7^2U_8^1`$
$`=`$ $`U_1^1U_2^2U_3^2U_4^3U_5^3U_6^2U_7^2U_8^1,`$
$`G^2`$ $`=`$ $`U_1^2U_2^1U_3^2U_2^2U_3^1U_4^2U_5^2U_6^1U_7^2U_6^2U_7^1U_8^2`$
$`=`$ $`U_1^2U_2^3U_3^3U_4^2U_5^2U_6^3U_7^3U_8^2.`$
For these instances the equivalence of $`G^1,G^2`$ to $`\mathrm{{\rm Y}}^1,\mathrm{{\rm Y}}^2`$ can be deduced by inspection. The result holds true not only for all linear cluster states, but can be generalised to cluster states defined on arbitrary $`d`$-dimensional lattices as defined in . However, the proof is made tedious by the fact that the definition of the $`K_j`$ depends on the choice of cluster in each case, so we omit any details.
The fact that each cluster state is equivalent to some perfect channel means that it has exactly the same entanglement properties as that channel. However, teleportation under our protocol using a cluster state will generally fail because the choice of measurement basis is not optimal. Mathematically, this is because cluster states do not belong to the stabiliser space of $`\{\mathrm{{\rm Y}}^1,\mathrm{{\rm Y}}^2\}`$. This serves to remind that while entanglement is necessary to achieve perfect fidelity teleportation, it is just as necessary that the entanglement be ordered with respect to a choice of measurement basis. In the next section we will construct a teleportation-order parameter which quantifies this order, and in turn the efficiency of a channel.
## 4 Teleportation-order
The manner in which we will construct a teleportation-order parameter is motivated by works studying the AKLT model and the role of the string-order parameter . The starting point for this study is the concept of localisable entanglement.
### 4.1 Localisable entanglement
Localisable entanglement is defined as the maximal possible entanglement that can be localised between two qubits (or more generally qudits), by an optimal choice of measurements on all other qubits of the system. The concept of localisable entanglement we follow is somewhat looser than that of in that we do not impose that the measurements are local, but rather are Bell measurements. Here we will show that each of the basis states $`|\stackrel{}{j}:\stackrel{}{k}\}`$ has maximal localisable entanglement between any two qubits with respect to any choice of Bell state measurements on all the other qubits of the system. Once we have established this fact, we then show that the same result holds for all states within a Bell subspace.
Below,
$$|\stackrel{}{j}:\stackrel{}{k}\}^{}$$
is defined to be such that
$$|\stackrel{}{j}:\stackrel{}{k}\}=|j_1:k_1\}|\stackrel{}{j}:\stackrel{}{k}\}^{}.$$
Now if $`|\stackrel{}{j}:\stackrel{}{k}\}`$ belongs to the Bell class $`[j:k]`$ then $`|\stackrel{}{j}:\stackrel{}{k}\}^{}`$ belongs to the Bell class $`[(jj_1):(kk_1)]`$. Next we appeal to (1), which permits us to write
$`|\stackrel{}{j}:\stackrel{}{k}\}`$ $`=`$ $`|j_1:k_1\}|\stackrel{}{j}:\stackrel{}{k}\}^{}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\right|+,k_1+j_1|,\overline{k_1})|\stackrel{}{j}:\stackrel{}{k}\}^{}`$
$`=`$ $`{\displaystyle \frac{\kappa }{2^{1/2}}}|+{\displaystyle \underset{\stackrel{}{p},\stackrel{}{q}}{}}|\stackrel{}{p}:\stackrel{}{q}\}X_{pq}^{(jj_1)(kk_1)}|k_1`$
$`+{\displaystyle \frac{j_1\kappa }{2^{1/2}}}|{\displaystyle \underset{\stackrel{}{p},\stackrel{}{q}}{}}|\stackrel{}{p}:\stackrel{}{q}\}X_{pq}^{(jj_1)(kk_1)}|\overline{k_1}`$
Making pairwise measurements on the interior qubits then projects out a state
$$\frac{1}{\sqrt{2}}\left(\right|+|\stackrel{}{p}:\stackrel{}{q}\}X_{pq}^{(jj_1)(kk_1)}|k_1+j_1||\stackrel{}{p}:\stackrel{}{q}\}X_{pq}^{(jj_1)(kk_1)}|\overline{k_1})$$
It is clear that the two end qubits are disentangled from the rest of the system by this process, and together form the state
$$(IX_{pq}^{(jj_1)(kk_1)})|j_1:k_1\}$$
which is one of the Bell states. Using the properties (4,8) we may rewrite this as
$`(IX_{pq}^{(jj_1)(kk_1)})|j_1:k_1\}`$ $`=`$ $`\kappa ^{}(IX_{pq}^{jk}X_{++}^{j_1k_1})|j_1:k_1\}`$
$`=`$ $`\kappa ^{\prime \prime }(IX_{pq}^{jk}X_{j_1k_1}^{++})|j_1:k_1\}`$
$`=`$ $`\kappa ^{\prime \prime }(IX_{pq}^{jk})|+:+\}`$
This final expression only depends on the Bell class of the channel and the Bell class of the measurement, so it extends to linear combinations of states from the same Bell class. It then follows that for any state from a fixed Bell class, keeping in mind the the subspace associated with each Bell class is invariant under the symmetric group, any sequence of Bell measurements on $`L2`$ qubits will leave the remaining two qubits in a Bell state. Depending on the context, it can be said that each Bell class state has maximal localisable entanglement under Bell measurements, or maximal entanglement length. These concepts have been discussed in in relation to the spin-1 AKLT model with spin-1/2 boundary sites. A significant feature of this model is that the system is gapped with finite-range spin correlations, yet has maximal entanglement length.
Following the notational conventions of , the Hamiltonian for the AKLT model with spin-1/2 boundaries reads
$`H=h_{1,2}+h_{+1,}+{\displaystyle \underset{j=2}{\overset{1}{}}}H_{j,j+1}`$ (23)
where
$`h_{j,k}`$ $`=`$ $`{\displaystyle \frac{2}{3}}(I+\stackrel{}{s}_j.\stackrel{}{S}_k)`$
$`H_{j,k}`$ $`=`$ $`\stackrel{}{S}_j.\stackrel{}{S}_k+{\displaystyle \frac{1}{3}}(\stackrel{}{S}_j.\stackrel{}{S}_k)^2.`$
Above, $`\stackrel{}{S}`$ is the vector spin-1 operator and $`\stackrel{}{s}`$ is the vector spin-1/2 operator. The ground state for the system can be constructed exactly using $`22`$ virtual qubits to represent the $`1`$ local spin-1 spaces . Let
$`P^t=P^{++}+P^++P^+\mathrm{End}(VV)`$
denote the projection onto the triplet space contained in $`VV`$, and let $`P_r^t`$ denote this operator acting on the $`r`$th and $`(r+1)`$th qubits of $`V^L`$. The Hilbert space of states for (23) is the image $`WV^L`$ of the operator
$`=P_2^tP_4^t\mathrm{}.P_{L2}^t,`$ (24)
and the ground state is given by
$`|AKLT=|\stackrel{}{p}:\stackrel{}{q}\}`$
where $`p_j=q_k=1j,k=1,\mathrm{}.,L`$; i.e. the ground state is the projection of a product of virtual 2-qubit singlet states into $`W`$. The ground state is a perfect channel under a protocol which employs a Bell measurement on the client state and one boundary spin, followed by a sequence of local measurements on the spin-1 sites . In this procedure the client state is teleported to the other boundary site. The fact that this protocol works with perfect fidelity can be understood through the string-order parameter .
For any state $`|\vartheta W`$ the string-order parameter $`𝒮(\vartheta )`$ is defined as
$`𝒮(\vartheta )=4\vartheta |s_1^z[_{k=2}^{}\mathrm{exp}(i\pi S_k^z)]s_{+1}^z|\vartheta `$ (25)
which takes values between $`1`$ and 1. It can be checked that for the ground state
$`𝒮(AKLT)=1.`$
We may extend the domain of the local operators $`\stackrel{}{S}_k`$ to act on the direct sum of the triplet and singlet spaces, and represent each local spin-($`10`$) space by the full tensor product $`VV`$ of two virtual qubits. It is then found that
$$\mathrm{exp}(i\pi S_k^z)=U^2U^2.$$
Therefore the expectation value (25) is precisely the expectation value of $`\mathrm{{\rm Y}}^2`$ up to a phase factor of $`(1)^{}`$. We thus see that the expectation value of $`\mathrm{{\rm Y}}^2`$ restricted to states in $`W`$ is equivalent to the string-order parameter for the case of the AKLT model.
### 4.2 Teleportation-order parameter
By analogy with the string-order parameter, for any state $`|\mathrm{\Psi }V^L`$ we define the teleportation-order parameter $`\stackrel{}{𝒯}^3`$ to be
$`\stackrel{}{𝒯}(\mathrm{\Psi })={\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle \underset{j=1}{\overset{3}{}}}\mathrm{\Psi }\left|\mathrm{{\rm Y}}^j\right|\mathrm{\Psi }\stackrel{}{e}_j`$
where the $`\{\stackrel{}{e}_j:j=1,2,3\}`$ denotes a set of orthonormal vectors for $`^3`$. Given an arbitrary $`|\mathrm{\Psi }V^L`$ we can make the decomposition into a linear combination of representatives from each Bell subspace:
$$|\mathrm{\Psi }=\underset{j,k}{}c_{[j:k]}|\mathrm{\Psi }_{[j:k]}$$
where $`|\mathrm{\Psi }_{[j:k]}V_{[j:k]}^L`$ is assumed to be normalised so that $`_{j,k}|c_{[j:k]}|^2=1`$. We can then determine that
$`c_{[j:k]}|\mathrm{\Psi }_{[j:k]}`$ $`=`$ $`P_{[j:k]}|\mathrm{\Psi }`$ (26)
$`=`$ $`{\displaystyle \frac{1}{4}}\left(|\mathrm{\Psi }+j\mathrm{{\rm Y}}^1|\mathrm{\Psi }+k\mathrm{{\rm Y}}^2|\mathrm{\Psi }+jk\mathrm{{\rm Y}}^3|\mathrm{\Psi }\right)`$
which in turn gives
$`|c_{[j:k]}|^2={\displaystyle \frac{1}{4}}\left(1+\mathrm{\Omega }_{[j:k]}\right)`$ (27)
where we have defined
$`\mathrm{\Omega }_{[j:k]}=j\mathrm{\Psi }\left|\mathrm{{\rm Y}}^1\right|\mathrm{\Psi }+k\mathrm{\Psi }\left|\mathrm{{\rm Y}}^2\right|\mathrm{\Psi }+jk\mathrm{\Psi }\left|\mathrm{{\rm Y}}^3\right|\mathrm{\Psi }.`$ (28)
Inverting these relations yields
$`\mathrm{\Psi }\left|\mathrm{{\rm Y}}^1\right|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{j,k}{}}j|c_{[j:k]}|^2`$
$`\mathrm{\Psi }\left|\mathrm{{\rm Y}}^2\right|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{j,k}{}}k|c_{[j:k]}|^2`$
$`\mathrm{\Psi }\left|\mathrm{{\rm Y}}^3\right|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{j,k}{}}jk|c_{[j:k]}|^2.`$
For any channel we define the efficiency of teleportation $``$ through
$`(\mathrm{\Psi })`$ $`=`$ $`|\stackrel{}{𝒯}(\mathrm{\Psi })|^2`$
$`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \underset{j=1}{\overset{3}{}}}\mathrm{\Psi }|\mathrm{{\rm Y}}^j|\mathrm{\Psi }^2`$
$`=`$ $`{\displaystyle \frac{1}{3}}\left(4{\displaystyle \underset{j,k}{}}|c_{[j:k]}|^41\right).`$
Since
$`(\mathrm{{\rm Y}}^j)^2=I^L`$
then for any state $`|\mathrm{\Psi }`$
$`1\mathrm{\Psi }\left|\mathrm{{\rm Y}}^j\right|\mathrm{\Psi }1,j=1,\mathrm{},L`$ (29)
and so the efficiency takes values in the range $`01`$. When $`=0`$ we see from (27,28) that the state is an equally weighted linear combination of states from each Bell subspace. At the other extreme, when $`=1`$ it indicates that the state belongs to a Bell subspace.
Next we will show that for any product state the efficiency is bounded:
$$01/3.$$
Let $`|v_jV`$ be arbitrary and let
$$|w=|w_1|w_2\mathrm{}|w_{}.$$
where
$$|w_j=|v_{(2j1)}|v_{2j}.$$
It is an exercise to show that
$`v_j|U^1|v_j^2+v_j|U^2|v_j^2v_j|U^3|v_j^2=1j=1,\mathrm{},L`$ (30)
and that
$`1v_j|U^1|v_j^20,`$
$`1v_j|U^2|v_j^20,`$
$`1v_j|U^3|v_j^20.`$
We can then deduce that
$`w_j|U^1U^1|w_j^2+w_j|U^2U^2|w_j^2+w_j|U^3U^3|w_j^2`$
$`=v_{(2j1)}|U^1|v_{(2j1)}^2v_{2j}|U^1|v_{2j}^2+v_{(2j1)}|U^2|v_{(2j1)}^2v_{2j}|U^2|v_{2j}^2`$
$`+v_{(2j1)}|U^3|v_{(2j1)}^2v_{2j}|U^3|v_{2j}^2`$
$`v_{2j}|U^1|v_{2j}^2+v_{2j}|U_{2j}^2|v_{2j}^2v_{2j}|U^3|v_{2j}^2=1`$
and moreover
$`1w_j|U^1U^1|w_j^20,`$
$`1w_j|U^2U^2|w_j^20,`$
$`1w_j|U^3U^3|w_j^20.`$
Proceeding analogously it follows by an inductive argument that
$$\underset{j=1}{\overset{3}{}}w|\mathrm{{\rm Y}}^j|w^21$$
and thus $``$ is bounded above by 1/3 for product states. Alternatively, if $`>1/3`$ for a state $`|\mathrm{\Psi }`$ then it must certainly be entangled, so $``$ provides a generalised notion of an entanglement witness .
### 4.3 Fidelity
The above analysis identified channels for which teleportation is achieved with perfect fidelity; viz. those channels which lie in a Bell subspace. In practise, there may be some error in the channel which leads to a loss of fidelity. Below we discuss how such a loss of fidelity may be quantified.
Without loss of generality, since the Bell subspaces are equivalent, let us assume that Alice and Bob believe the channel to lie in $`V_{[+:+]}^L`$. The protocol requires that Alice makes $``$ Bell measurements on her subsystem which is comprised of the client qubit and $`L1`$ qubits of the channel. In the case of perfect channels, we have shown that the protocol is independent of the way in which she pairs the qubits, nor the order in which she makes the measurements. This is also true for the case of non-perfect channels, which can be seen from eq. (17). So we may simplify the problem by assuming that Alice first makes $`1`$ Bell measurements on qubits which are contained within the channel, and then the final measurement involving one channel qubit and the client qubit. In view of our earlier discussion on localisable entanglement, the first $`1`$ measurements project each of the component states $`|\mathrm{\Psi }_{[j:k]}`$ onto a product of Bell states tensored with a Bell state shared by Alice and Bob. Suppose that the result of Alice’s first $`1`$ Bell measurements is of class $`[r:s]`$. The Bell class of each of the component Bell states of the total state shared between Alice and Bob after this measurement can be determined from the Bell class of the measurement, as in Sect. 4.1. Thus, after Alice’s first $`1`$ Bell measurements, the channel shared by Alice and Bob is found to be of the form
$`|\mathrm{\Theta }|\mathrm{\Psi }`$
where $`|\mathrm{\Theta }`$ is a state of Bell class $`[r:s]`$, onto which Alice has projected as a result of her measurement, and
$`|\mathrm{\Psi }={\displaystyle \underset{j,k}{}}C_{[j:k]}|j:k\}`$
with
$`C_{[j:k]}=e^{i\theta _{[j:k]}}c_{[rj:sk]}.`$
The above phase factors $`\theta _{[j:k]}`$ are unknown, because Alice’s measurement results do not determine the overall phases of the components of the remaining shared state. Now the problem has been reduced to the investiagation of teleportation across the 2-qubit channel $`|\mathrm{\Psi }`$ which Alice, as a result of her measurements, believes to be the Bell state $`|r:s\}`$.
From (10) we may write
$`|v|\mathrm{\Psi }={\displaystyle \frac{1}{2}}{\displaystyle \underset{p,q}{}}|p:q\}\stackrel{~}{X}_{pq}|v`$ (31)
where
$`\stackrel{~}{X}_{pq}={\displaystyle \underset{j,k}{}}C_{[j:k]}\stackrel{~}{X}_{pq}^{jk}.`$
Note that $`\stackrel{~}{X}_{pq}`$ is not necessarily a unitary matrix. The probability $`𝒫_{[pr:qs]}`$ of Alice’s final measurement result being $`(p:q)`$, thus making her overall measurement of class $`[pr:qs]`$, is given by
$`𝒫_{[pr:qs]}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left|v\left|\stackrel{~}{X}_{pq}^{}\stackrel{~}{X}_{pq}\right|v\right|.`$
Now suppose that Alice’s final measurement result is $`(p:q)`$, so she communicates to Bob that the total measurement class is $`[pr:qs]`$. Upon receiving this information, he would apply the correction gate $`(\stackrel{~}{X}_{(pr)(qs)}^{++})^{}`$ in attempting to recover the client state. We define the fidelity $`_{(pr)(qs)}^{++}`$ of this attempted teleportation as the square of the magnitude of the overlap between the client state and the state Bob has obtained; i.e
$`_{(pr)(qs)}^{++}`$ $`=`$ $`{\displaystyle \frac{\left|v\left|(\stackrel{~}{X}_{++}^{(pr)(qs)})^{}\stackrel{~}{X}_{pq}\right|v\right|^2}{\left|v\left|\stackrel{~}{X}_{pq}^{}\stackrel{~}{X}_{pq}\right|v\right|}}`$ (32)
$`=`$ $`{\displaystyle \frac{\left|v\left|(\stackrel{~}{X}_{pq}^{rs})^{}\stackrel{~}{X}_{pq}\right|v\right|^2}{\left|v\left|\stackrel{~}{X}_{pq}^{}\stackrel{~}{X}_{pq}\right|v\right|}}`$
where we have used (4,8).
We return to eq. (31). Alice believes that $`|\mathrm{\Psi }`$ is the Bell state $`|r:s\}`$ (up to an overall phase which we hereafter ignore) so we write for real $`\theta `$
$`|\mathrm{\Psi }`$ $`=`$ $`(IR(\theta ,\widehat{n}))|r:s\}`$
with
$`R(\theta ,\widehat{n})`$ $`=`$ $`\mathrm{cos}(\theta /2)Ii\mathrm{sin}(\theta /2)(n_1U^1+n_2U^2+in_3U^3)`$
so that
$`\stackrel{~}{X}_{pq}=R(\theta ,\widehat{n})\stackrel{~}{X}_{pq}^{rs}.`$ (33)
Using (2) we then find
$`|\mathrm{\Psi }`$ $`=`$ $`(\mathrm{cos}(\theta /2)Ii\mathrm{sin}(\theta /2)(n_1X_{rs}^{r\overline{s}}+(1)^sn_2X_{rs}^{\overline{r}s}+(1)^sin_3X_{rs}^{\overline{r}\overline{s}}))|r:s\}`$
such that we can idenitfy
$`C_{[r:s]}`$ $`=`$ $`\mathrm{cos}(\theta /2)`$
$`C_{[r:\overline{s}]}`$ $`=`$ $`in_1\mathrm{sin}(\theta /2)`$
$`C_{[\overline{r}:s]}`$ $`=`$ $`(1)^sin_2\mathrm{sin}(\theta /2)`$
$`C_{[\overline{r}:\overline{s}]}`$ $`=`$ $`(1)^sn_3\mathrm{sin}(\theta /2).`$
Normalisation of $`|\mathrm{\Psi }`$ requires that $`\widehat{n}=(n_1,n_2,n_3)`$ is a unit complex vector. In the case that $`\widehat{n}`$ is real then $`R(\theta ,\widehat{n})`$ is a unitary matrix corresponding to a rotation of the Bloch sphere by an angle $`\theta `$ about an axis determined by $`\widehat{n}`$. However $`R(\theta ,\widehat{n})`$ is not unitary for a generic complex unit vector $`\widehat{n}`$. Substituting (33) into (32) gives
$`_{(pr)(qs)}^{++}`$ $`=`$ $`{\displaystyle \frac{\left|v\left|(\stackrel{~}{X}_{pq}^{rs})^{}R(\theta ,\widehat{n})\stackrel{~}{X}_{pq}^{rs}\right|v\right|^2}{\left|v\left|(\stackrel{~}{X}_{pq}^{rs})^{}R^{}(\theta ,\widehat{n})R(\theta ,\widehat{n})\stackrel{~}{X}_{pq}^{rs}\right|v\right|}}.`$
The minimum fidelity is
$`\mathrm{min}\left(_{(pr)(qs)}^{++}\right)`$ $`=`$ $`\underset{|v}{\mathrm{min}}{\displaystyle \frac{\left|v\left|(\stackrel{~}{X}_{pq}^{rs})^{}R(\theta ,\widehat{n})\stackrel{~}{X}_{pq}^{rs}\right|v\right|^2}{\left|v\left|(\stackrel{~}{X}_{pq}^{rs})^{}R^{}(\theta ,\widehat{n})R(\theta ,\widehat{n})\stackrel{~}{X}_{pq}^{rs}\right|v\right|}}`$ (34)
$`=`$ $`\underset{|v}{\mathrm{min}}{\displaystyle \frac{\left|v\left|R(\theta ,\widehat{n})\right|v\right|^2}{\left|v\left|R^{}(\theta ,\widehat{n})R(\theta ,\widehat{n})\right|v\right|}}`$
$``$ $`2\mathrm{cos}^2(\theta /2)1`$
where the above inequality holds for all $`\widehat{n}`$. The proof of this result is given in Appendix B. We now have
$`\mathrm{min}\left(_{(pr)(qs)}^{++}\right)`$ $``$ $`2\mathrm{cos}^2(\theta /2)1`$
$`=`$ $`2|C_{[r:s]}|^21`$
$`=`$ $`2|c_{[+:+]}|^21`$
$`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Omega }_{[+:+]}1),`$
so that generally
$`_{pq}^{jk}`$ $``$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Omega }_{[j:k]}1).`$
That is the quantitity $`\mathrm{\Omega }_{[j:k]}`$, which is simply a linear combination of the components of the teleportation-order parameter, provides a lower bound on the fidelity. The results of simulations for four-qubit channels are shown in Fig. 2 (also in ).
## 5 Teleportation via three-qubit channels
In the above, teleportation was only investigated for an even number of channel qubits. It leaves open the problem of devising a teleportation protocol when the number of channel qubits is odd. Here we won’t address this in a general context, but we will show that a protocol does exist for teleportation via three-qubit channels.
We begin by defining an orthonormal basis for three-qubit states which generalises the Bell basis. Let
$`|+:+:+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+,++|,,\right)`$
$`|+:+:\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+,+|,,+\right)`$
$`|+::+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,,++|,+,\right)`$
$`|+::\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,,+|,+,+\right)`$
$`|:+:+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+,+|,,\right)`$
$`|:+:\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,+,|,,+\right)`$
$`|::+\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,,+|,+,\right)`$
$`|::\}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|+,,|,+,+\right).`$
Any basis state can be conveniently expressed as
$`|j:k:l\}={\displaystyle \frac{1}{\sqrt{2}}}\left(\right|+,k,l+j|,\overline{k},\overline{l})`$ (35)
and satisfies
$`\mathrm{\Lambda }^1|j:k:l\}`$ $`=`$ $`j|j:k:l\}`$
$`\mathrm{\Lambda }^2|j:k:l\}`$ $`=`$ $`k|j:k:l\}`$
$`\mathrm{\Lambda }^3|j:k:l\}`$ $`=`$ $`l|j:k:l\}`$
where
$`\mathrm{\Lambda }^1`$ $`=`$ $`U^1U^1U^1`$
$`\mathrm{\Lambda }^2`$ $`=`$ $`U^2U^2I`$
$`\mathrm{\Lambda }^3`$ $`=`$ $`U^2IU^2.`$
Above, the $`\mathrm{\Lambda }^j`$ are mutually commuting self-adjoint operators, so their actions represent a simultaneous measurement. We call such a measurement a three-qubit Bell measurement where $`(j:k:l)`$ is the measurement outcome. Defining $`Y_{pqr}^{jkl}\mathrm{End}(VV)`$ by the relation
$`|j:k:l\}=IY_{pqr}^{jkl}|p:q:r\}`$
we find that
$`Y_{pqr}^{jkl}=Z_q^kX_{pr}^{jl}`$
where
$`Z_j^j=I,Z_{\overline{j}}^j=U^1`$
and the $`X_{pr}^{jl}`$ are as before.
Now consider for an arbitrary client state $`|v=(\alpha |++\beta |)/\sqrt{2}`$
$`2|v|+:+:+\}`$ $`=`$ $`\sqrt{2}\left(\alpha \right|+,+,+,++\beta |,+,+,+`$
$`+\alpha |+,,,+\beta |,,,)`$
$`=`$ $`|+:+:+\}(\alpha |++\beta |)`$
$`+|+::\}(\beta |++\alpha |)`$
$`+|:+:+\}(\alpha |+\beta |)`$
$`+|::\}(\beta |++\alpha |)`$
$`=`$ $`|+:+:+\}|v+|+::\}U^1|v`$
$`+|:+:+\}U^2|v+|::\}U^3|v.`$
By applying the operators $`U^j`$ to the last qubit of the space we deduce the following relations
$`2|v|+::\}`$ $`=`$ $`|+:+:\}U^1|v+|+::+\}|v`$
$`+|:+:\}U^3|v|::+\}U^2|v,`$
$`2|v|:+:+\}`$ $`=`$ $`|+:+:+\}U^2|v|+::\}U^3|v`$
$`+|:+:+\}|v+|::\}U^1|v,`$
$`2|v|::\}`$ $`=`$ $`|+:+:\}U^3|v|+::+\}U^2|v`$
$`+|:+:\}U^1|v+|::+\}|v.`$
Applying the operator $`U^1`$ to the second last qubit of the space in the above four cases gives
$`2|v|+::+\}`$ $`=`$ $`|+:+:\}|v+|+::+\}U^1|v`$
$`+|:+:\}U^2|v+|::+\}U^3|v,`$
$`2|v|+:+:\}`$ $`=`$ $`|+:+:+\}U^1|v+|+::\}|v`$
$`+|:+:+\}U^3|v|::\}U^2|v,`$
$`2|v|::+\}`$ $`=`$ $`|+:+:\}U^2|v|+::+\}U^3|v`$
$`+|:+:\}|v+|::+\}U^1|v,`$
$`2|v|:+:\}`$ $`=`$ $`|+:+:+\}U^3|v|+::\}U^2|v`$
$`+|:+:+\}U^1|v+|::\}|v.`$
The eight relations above can all be expressed in a unified way:
$`|v|j:k:l\}={\displaystyle \frac{1}{2}}{\displaystyle \underset{p,q}{}}|p:q:(kq)\}\stackrel{~}{X}_{pq}^{jl}|v,`$ (36)
which is a three-qubit channel generalisation of (10). By the same argument as in the two-qubit channel case, we conclude that each $`|j:k:l\}`$ is a perfect channel for teleportation. After Alice performs a three-qubit Bell measurement which projects the system onto a state
$`|p:q:(kq)\}\stackrel{~}{X}_{pq}^{jl}|v,`$
two bits of classical information (i.e. $`p`$ and $`q`$) need to be transmitted to Bob for him to determine the correction gate.
Because the teleporation protocol requires that only two bits of classical information be sent to Bob, that part of the Bell measurement represented by $`\mathrm{\Lambda }_2`$ becomes redundant. Hence, in analogy with the case where the channel is a state of a system with an even number of qubits, will can still define four Bell classes of channels which give rise to the Hilbert space decomposition
$`V^3=V_{[+:+]}^3V_{[+:]}^3V_{[:+]}^3V_{[:]}`$
where the class indices $`[j:k]`$ are the eigenvalues of the operators $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_3`$. Indeed, we can take a generalised form of (36)
$`|v{\displaystyle \underset{k}{}}\alpha _k|j:k:l\}={\displaystyle \frac{1}{2}}{\displaystyle \underset{p,q,k}{}}\alpha _k|p:q:(kq)\}\stackrel{~}{X}_{pq}^{jl}|v,`$ (37)
and perform a reduced three-qubit Bell measurement which is represented by $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_3`$. Above, $`\alpha _1,\alpha _2`$ are arbitrary up to the constraint of normalisation. The consequence of the reduced Bell measurement is that it leaves the system in a state of the form
$`{\displaystyle \underset{k}{}}\alpha _k|p:q:(kq)\}\stackrel{~}{X}_{pq}^{jl}|v,`$
and once again teleportation can be achieved with perfect fidelity. What this result tells us is that the state of the third qubit of the system (i.e. the second qubit of the channel) is of no consequence in this teleportation protocol. In fact, by a suitable choice of $`\alpha _1`$ and $`\alpha _2`$ the channel factorises into a tensor product of a Bell state for the first and third qubits of the channel and a disentangled qubit state for the second qubit. Thus this protocol for teleportation via a three-qubit channel is essentially a two-qubit channel protocol as the third qubit can be made redundant.
This raises the question of whether the entanglement of a three-qubit channel can be used to achieve more efficient teleportation than a two-qubit channel. Specifically, can two qubits be teleported via a three-qubit channel? Within the protocol considered here this is not the case. Let $`|v`$ and $`|w`$ be arbitrary qubit states. Using (36) we calculate
$`|v|w|j:k:l\}`$ $`=`$ $`|v\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{p,q}{}}\right|p:q:(kq)\}\stackrel{~}{X}_{pq}^{jl}|w)`$
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{p,q,r,s}{}}|r:s:(sq)\}\stackrel{~}{X}_{rs}^{p(kq)}|v\stackrel{~}{X}_{pq}^{jl}|w`$
Now make the change of variable $`t=sq`$:
$`|v|w|j:k:l\}`$
$`={\displaystyle \frac{1}{4}}{\displaystyle \underset{p,r,s,t}{}}|r:s:t\}\stackrel{~}{X}_{rs}^{p(kst)}|v\stackrel{~}{X}_{p(st)}^{jl}|w`$
$`={\displaystyle \frac{1}{\sqrt{8}}}{\displaystyle \underset{r,s,t}{}}|r:s:t\}{\displaystyle \underset{p}{}}{\displaystyle \frac{1}{\sqrt{2}}}(\stackrel{~}{X}_{rs}^{p(kst)}\stackrel{~}{X}_{p(st)}^{jl})\left(\right|v|w)`$
$`={\displaystyle \frac{1}{\sqrt{8}}}{\displaystyle \underset{r,s,t}{}}|r:s:t\}(\stackrel{~}{X}_{rs}^{+(kst)}\stackrel{~}{X}_{+(st)}^{jl})\mathrm{\Theta }\left(\right|v|w)`$
where $`\mathrm{\Theta }=\left(II+\kappa \stackrel{~}{X}_{++}^+\stackrel{~}{X}_+^{++}\right)/\sqrt{2}=\left(II+\kappa U^2U^2\right)/\sqrt{2}`$. Because $`\mathrm{\Theta }`$ is not invertible, there is no possibility to effect two-qubit teleportation in this manner.
## 6 A qudit generalisation
As discussed in the original work , it is also possible to teleport qudit states (see also ). Let $`\{|l:l=0,\mathrm{},d1\}`$ denote a set of orthonormal basis states for a qudit. For $`\omega `$ a fixed primitive $`d`$th root of unity, we introduce the permutation and phase matrices defined by
$`P|l`$ $`=`$ $`|l+1,`$
$`Q|l`$ $`=`$ $`\omega ^l|l`$
and set
$`R^{kj}=P^kQ^j.`$
Throughout, the state labels are taken modulo $`d`$ so for example $`|d|0`$. A qudit generalisation of Bell states is given by
$`|j:k\}`$ $`=`$ $`(IU^{kj})|0:0\}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{d}}}{\displaystyle \underset{l=0}{\overset{d1}{}}}\omega ^{jl}|l|l+k`$
where
$`|0:0\}={\displaystyle \frac{1}{\sqrt{d}}}{\displaystyle \underset{l=0}{\overset{d1}{}}}|l|l.`$
The generalised Bell states provide a basis which allows us to write
$`|j|k={\displaystyle \frac{1}{\sqrt{d}}}{\displaystyle \underset{l=0}{\overset{d1}{}}}\omega ^{jl}|l:kj\}.`$ (39)
Letting $`|v=_{j=0}^{d1}\alpha _j|j`$ denote an arbitrary qudit state we find by using (39)
$`|v|0:0\}={\displaystyle \frac{1}{d}}{\displaystyle \underset{p,q=0}{\overset{d1}{}}}|p:q\}R^{q\overline{p}}|v`$
where as before $`\overline{p}=p`$. From (6) it follows that
$`|v|j:k\}={\displaystyle \frac{1}{d}}{\displaystyle \underset{p,q=0}{\overset{d1}{}}}|p:q\}\stackrel{~}{X}_{pq}^{jk}|v`$ (40)
with
$`\stackrel{~}{X}_{pq}^{jk}=R^{kj}R^{q\overline{p}}.`$
It is apparent that (40) is a qudit generalisation of (10). As the operators $`R^{jk}`$ generate a group, since $`QP=\omega PQ`$, so do the $`\stackrel{~}{X}_{pq}^{jk}`$. It follows that our analysis for qubit systems generalises to qudit systems, with the main finding being that for the $`L`$-qudit case there are $`d^2`$ Bell subspaces of perfect channels. The Bell subspaces are equivalent, with each having dimension $`d^{L2}`$.
## 7 Conclusion
To conclude, we discuss some aspects of our results in the context of 4-qubit channels. The 16-dimensional Hilbert space $`V^4`$ decomposes into four Bell subspaces $`V_{[j:k]}^4`$, each of dimension four. Since these subspaces are all equivalent, we can focus on the space $`V_{[+:+]}^4`$. This space is precisely the space $`G_{abcd}`$ of , the generic equivalence class of 4-qubit states representing the orbits arising from stochastic local operations and classical communication (see also ). It contains the 4-qubit case of the celebrated Greenberger-Horne-Zeilinger states, for which it has been argued are the only states exhibiting essential multi-partite entanglement . It also contains the state of Higuchi and Sudbery , which has the largest known average 2-qubit bi-partite entanglement in a 4-qubit system (see also ). This state, together with its complex conjugate state, provides a basis for the space of singlets contained in $`V^4`$. Further, there are three states in $`V_{[+:+]}^4`$ which are equivalent to the three 4-qubit cluster states, known to have maximal connectedness and high persistency of entanglement . All the above mentioned states, by representing different forms of multi-partite entanglement, are not equivalent in the sense that they are not related by local unitary transformations. They are however all entirely equivalent for the purpose of teleportation under our prescribed protocol, since they all belong to $`V_{[+:+]}^4`$. This highlights the fact that the entanglement needed to implement this protocol is of a specific type, which depends on each qubit being maximally entangled with the rest of the system (maximal local disorder). Other forms of entanglement the channel might possess are irrelevant. We do emphasise though that for the channel to be effective, this entanglement has to be ordered with respect to the prescribed measurement basis, which is quantified by the teleportation-order parameter.
Lastly we mention that, besides the one described here, there are many possible teleportation protocols which generalise the original work of ; e.g., see . It would be useful in future work to identify a correspondence between any given teleportation protocol, and a teleportation-order parameter which signifies when a channel can be used to implement the protocol and effect teleportation with full fidelity. Furthermore, the possibilities for performing teleportation without a shared reference frame, following the ideas developed in , also warrant investigation.
Acknowledgements \- This work was supported by the Australian Research Council. We thank Steve Bartlett and Michael Nielsen for helpful advice, and we are indebted to Dominic Berry for providing the proof in Appendix B of the inequality (34).
Appendix A
Here we show by example that not all measurement outcomes are equally likely. Let the channel be
$`|\mathrm{\Psi }=\mathrm{cos}(\varphi )|+:\}|:+\}+\mathrm{sin}(\varphi )|:+\}|+:\}`$
which is of Bell class $`[:]`$. Using (6,7,17) we may write
$`4|v|\mathrm{\Psi }`$ $`=`$ $`\mathrm{cos}(\varphi ){\displaystyle \underset{p_1,p_2,q_1,q_2}{}}|p_1p_2:q_1q_2\}\stackrel{~}{X}_{p_2q_2}^+\stackrel{~}{X}_{p_1q_1}^+|v`$
$`+\mathrm{sin}(\varphi ){\displaystyle \underset{p_1,p_2,q_1,q_2}{}}|p_1p_2:q_1q_2\}\stackrel{~}{X}_{p_2q_2}^+\stackrel{~}{X}_{p_1q_1}^+|v`$
$`=`$ $`\mathrm{cos}(\varphi ){\displaystyle \underset{p_1,p_2,q_1,q_2}{}}\epsilon _{++}^{p_1q_1}\epsilon _+^{p_2q_2}|p_1p_2:q_1q_2\}X_{p_2q_2}^{++}X_{p_1q_1}^{}|v`$
$`+\mathrm{sin}(\varphi ){\displaystyle \underset{p_1,p_2,q_1,q_2}{}}\epsilon _{++}^{p_1q_1}\epsilon _+^{p_2q_2}|p_1p_2:q_1q_2\}X_{p_2q_2}^{++}X_{p_1q_1}^{}|v`$
$`=`$ $`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|++:++\}U^3|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|++:+\}U^2|v`$
$`+(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|++:+\}U^2|v`$
$`(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|++:\}U^3|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|+:++\}U^1|v`$
$`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|+:+\}U^0|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|+:+\}U^0|v`$
$`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|+:\}U^1|v`$
$`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|+:++\}U^1|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|+:+\}U^0|v`$
$`+(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|+:+\}U^0|v`$
$`(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|+:\}U^1|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|:++\}U^3|v`$
$`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|:+\}U^2|v`$
$`+(\mathrm{cos}(\varphi )+\mathrm{sin}(\varphi ))|:+\}U^2|v`$
$`(\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ))|:\}U^3|v.`$
From the above it can be seen that for any measurement projecting onto a state
$`|p_1p_2:q_1q_2\}U^j|v`$
the probability is either
$`{\displaystyle \frac{1}{16}}\left(1+\mathrm{sin}(2\varphi )\right)`$
or
$`{\displaystyle \frac{1}{16}}\left(1\mathrm{sin}(2\varphi )\right)`$
so not all measurement outcomes are equally likely. It is also easily checked for this case that the probability a measurement is of the Bell class $`[j:k]`$ is 1/4, independent of $`j`$ and $`k`$, which is consistent with our earlier result.
Appendix B
Here we prove the inequality (34), viz.
$`\underset{|v}{\mathrm{min}}{\displaystyle \frac{\left|v\left|R(\theta ,\widehat{n})\right|v\right|^2}{\left|v\left|R^{}(\theta ,\widehat{n})R(\theta ,\widehat{n})\right|v\right|}}`$ $``$ $`2\mathrm{cos}^2(\theta /2)1.`$
Suppose that $`|w`$ is a vector for which the miminum is achieved. We can then perform a unitary transformation such that $`|w`$ is transformed into the state $`|+`$. Under such a unitary transformation, the operator $`R(\theta ,\widehat{n})`$ is transformed into $`R(\theta ,\widehat{m})`$ for some unit complex vector $`\widehat{m}`$. Importantly, the variable $`\theta `$ is the same for both operators. Setting
$`\mathrm{\Delta }(\theta ,\widehat{m})={\displaystyle \frac{|+\left|R(\theta ,\widehat{m})\right|+|^2}{\left|+\left|R^{}(\theta ,\widehat{m})R(\theta ,\widehat{m})\right|+\right|}}`$
we now need to show that, for all $`\widehat{m}`$,
$`\mathrm{\Delta }(\theta ,\widehat{m})`$ $``$ $`2\mathrm{cos}^2(\theta /2)1.`$
We may express a generic $`R(\theta ,\widehat{m})`$ as
$`R(\theta ,\widehat{m})=\mathrm{cos}(\theta /2)I+\mathrm{sin}(\theta /2)\left(\begin{array}{cc}a& c\\ b& a\end{array}\right)`$
where $`a,b,c`$ are complex parameters subject to the normalisation constraint
$`|a|^2+{\displaystyle \frac{1}{2}}(|b|^2+|c|^2)=1.`$
Without loss of generality we may impose $`0\theta <\pi `$. In terms of these parameters we have
$`\mathrm{\Delta }(\theta ,\widehat{m})={\displaystyle \frac{|\mathrm{cos}(\theta /2)+a\mathrm{sin}(\theta /2)|^2}{|\mathrm{cos}(\theta /2)+a\mathrm{sin}(\theta /2)|^2+|b\mathrm{sin}(\theta /2)|^2}}.`$ (41)
For any fixed $`a`$, (41) is minimised by maximising $`|b|^2`$. We thus choose $`c=0`$ leading to
$`\mathrm{\Delta }(\theta ,\widehat{m})`$ $`=`$ $`{\displaystyle \frac{|\mathrm{cos}(\theta /2)+a\mathrm{sin}(\theta /2)|^2}{|\mathrm{cos}(\theta /2)+a\mathrm{sin}(\theta /2)|^2+2(1|a|^2)\mathrm{sin}^2(\theta /2)}}`$ (42)
$`=`$ $`{\displaystyle \frac{\mathrm{cos}^2(\theta /2)+|a|^2\mathrm{sin}^2(\theta /2)+2\mathrm{}(a)\mathrm{sin}(\theta /2)\mathrm{cos}(\theta /2)}{1+(1|a|^2)\mathrm{sin}^2(\theta /2)+2\mathrm{}(a)\mathrm{sin}(\theta /2)\mathrm{cos}(\theta /2)}}`$
where $`\mathrm{}(a)`$ denotes the real part of $`a`$. The above expression is minimised when $`a`$ is real and given by
$`a={\displaystyle \frac{\sqrt{14\mathrm{cos}^2(\theta /2)\mathrm{sin}^2(\theta /2)}1}{2\mathrm{sin}(\theta /2)\mathrm{cos}(\theta /2)}}={\displaystyle \frac{\mathrm{cos}(\theta )1}{\mathrm{sin}(\theta )}}.`$ (43)
Finally, substituting (43) into (42) gives
$`\underset{\widehat{m}}{\mathrm{min}}\left(\mathrm{\Delta }(\theta ,\widehat{m})\right)=\mathrm{cos}(\theta )=2\mathrm{cos}^2(\theta /2)1`$
as required.
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# Contents
## 1 Introduction
In 1983, Berry made the surprising discovery that a quantum system adiabatically transported round a closed circuit $`𝒞`$ in the space of external parameters acquires, besides the familiar dynamical phase, a non-integrable phase depending only on the geometry of the circuit $`𝒞`$ . This Berry phase, which had been overlooked for more than half a century, provides us a very deep insight on the geometric structure of quantum mechanics and gives rise to various observable effects. The concept of the Berry has now become a central unifying concept in quantum mechanics, with applications in fields ranging from chemistry to condensed matter physics .
The aim of the present lecture is to give an elementary introduction to the Berry phase, and to discuss its various implications in the field of magnetism, where it plays an increasingly important role. The reader is referred to specialized textbooks for a more comprehensive presentation of this topic.
## 2 Parallel transport in geometry
The importance of the Berry phase stems from the fact that it reveals the intimate geometrical structure underlying quantum mechanics. It is therefore appropriate to start with an introduction of the fundamental concept of *parallel transport* in a purely geometrical context; we follow here the discussion given by Berry in Ref. .
This is best illustrated by means of a simple example. Consider a surface $`\mathrm{\Sigma }`$ (e.g., a plane, a sphere, a cone, etc.) and a vector constrained to lie everywhere in the plane tangent to the surface. Next, we wish to transport the vector on the surface, *without rotating it around the axis normal to the surface*, as illustrated in Fig. 1. We are interested, in particular in the case, in which the arrow is transported round a closed circuit $`𝒞(1231)`$. We may encounter two different situations: (i) if the surface is flat, as on Fig. 1(a), then the arrow always remains parallel to its original orientation, and therefore is unchanged after completion of the circuit $`𝒞`$; (ii) if, however, the surface $`\mathrm{\Sigma }`$ is curved as on Fig. 1(b,c), the arrow, being constrained to lie in the local tangent plane, cannot remain parallel to its original orientation, and after completion of the circuit $`𝒞`$, it is clearly seen to have been rotated by an angle $`\theta (𝒞)`$, a phenomenon referred to as *anholonomy*.
Let us now formalize this procedure. The arrow is represented by a tangent unit vector $`𝐞^1`$, transported along a circuit $`𝒞\{𝐫(t)|t=0T\}`$ on the surface. Defining $`𝐧(𝐫)`$ as the unit vector normal to the surface at point $`𝐫`$, we define a second tangent unit vector $`𝐞^2𝐧\times 𝐞^1`$, which is also parallel transported on the surface along $`𝒞`$. The 3 unit vectors $`(𝐧,𝐞^1,𝐞^2)`$ form an orthonormal reference frame. As $`𝐞^1`$ and $`𝐞^2`$ are transported, they have to rotate with an angular velocity $`𝝎`$ (to be determined) if the surface is not flat, i.e., the equation of motion of $`𝐞^1`$ and $`𝐞^2`$ is
$$\dot{𝐞^r}=𝝎\times 𝐞^r(r=1,2),$$
(1)
where the overdot indicates the time derivative. One can easily see that in order to fulfill the requirements that $`𝐞^1`$ and $`𝐞^2`$ remain tangent unit vectors (i.e., $`𝐞^r𝐧=0`$, $`(r=1,2)`$) and never rotate around $`𝐧`$ (i.e., $`𝝎𝐧=0`$), the angular velocity has to be given by
$$𝝎=𝐧\times \dot{𝐧}.$$
(2)
The law of parallel transport is therefore
$$\dot{𝐞}^r=(𝐧\times \dot{𝐧})\times 𝐞^r=(𝐞^r\dot{𝐧})𝐧.$$
(3)
This law can be expressed in a form more suitable for generalization to the case of quantum mechanics, by defining the complex unit vector
$$\mathit{\varphi }\frac{𝐞^1+\mathrm{i}𝐞^2}{\sqrt{2}},$$
(4)
with
$$\mathit{\varphi }^{}\mathit{\varphi }=1.$$
(5)
The law of parallel transport now reads
$$\mathit{\varphi }^{}\dot{\mathit{\varphi }}=0.$$
(6)
In order to express the rotation of the unit vectors $`(𝐞^1,𝐞^2)`$ as they move around $`𝒞`$, we need to chose a *fixed* local orthonormal frame $`(𝐧(𝐫),𝐭^1(𝐫),𝐭^2(𝐫))`$on the surface. The normal unit vector $`𝐧(𝐫)`$ is of course uniquely determined by the surface, but we have an infinity of possible choices for $`𝐭^1(𝐫)`$ (we simply impose that it is a smooth function of $`𝐫`$), which corresponds to a gauge freedom; once we have made a choice for $`𝐭^1(𝐫)`$, then $`𝐭^2(𝐫)`$ is of course uniquely determined. We next define the complex unit vector
$$𝐮(𝐫)\frac{𝐭^1(𝐫)+\mathrm{i}𝐭^2(𝐫)}{\sqrt{2}},$$
(7)
with, of course,
$$𝐮^{}(𝐫)𝐮(𝐫)=1.$$
(8)
The relation between the parallel transported frame and the fixed one is expressed as
$$\mathit{\varphi }(t)=\mathrm{exp}[\mathrm{i}\theta (t)]𝐮\left(𝐫(t)\right),$$
(9)
where $`\theta (t)`$ is the angle by which $`(𝐭^1,𝐭^2)`$ must be rotated to coincide with $`(𝐞^1,𝐞^2)`$. We obtain the equation satisfied by $`\theta (t)`$ by inserting the above definition in the equation of parallel transport (6), and obtain
$$0=\mathit{\varphi }^{}\dot{\mathit{\varphi }}=\mathrm{i}\dot{\theta }𝐮^{}𝐮+𝐮^{}\dot{𝐮}.$$
(10)
Since $`𝐮^{}𝐮=1`$ and then $`𝐮^{}\dot{𝐮}`$ is imaginary, we get
$$\dot{\theta }=\mathrm{Im}(𝐮^{}\dot{𝐮}),$$
(11)
so that
$`\theta (𝒞)`$ $`=`$ $`\mathrm{Im}{\displaystyle _𝒞}𝐮^{}d𝐮`$ (12)
$`=`$ $`{\displaystyle _𝒞}𝐭^2d𝐭^1.`$ (13)
If we choose a coordinate system $`(X_1,X_2)`$ on our surface $`\mathrm{\Sigma }`$, and define the vector field $`𝐀(𝐫)`$ (usually called a *connection*) on $`\mathrm{\Sigma }`$ as
$$A_i(𝐗)\mathrm{Im}\left[u_j^{}(𝐗)\frac{u_j(𝐗)}{X_i}\right],$$
(14)
where we have used Einstein’convention of summation over repeated indices, we get
$$\theta (𝒞)=_𝒞𝐀(𝐗)d𝐗,$$
(15)
which constitutes the *1-form* expression of the anholonomy angle $`\theta (𝒞)`$. The connection $`𝐀(𝐗)`$ depends on our particular gauge choice for $`𝐭^1(𝐗)`$: if we make a new choice $`𝐭_{}^{1}{}_{}{}^{}(𝐗)`$ which is brought in coincidence with $`𝐭^1(𝐗)`$ by a rotation of angle $`\mu (𝐗)`$, i.e., if we make the gauge transformation
$$𝐮(𝐗)𝐮^{}(𝐗)\mathrm{exp}\left(\mathrm{i}\mu (𝐗)\right)𝐮(𝐗),$$
(16)
we obtain a new connection
$$A_i^{}(𝐗)\mathrm{Im}\left[u_j^{}(𝐗)\frac{u_j(𝐗)}{X_i}\right]=A_i(𝐗)\frac{\mu (𝐗)}{X_i}.$$
(17)
However, since
$$_𝒞\mu (𝐫)d𝐫=_𝒞d\mu (𝐫)=0,$$
(18)
we can see that the expression (15) for the anholonomy angle $`\theta (𝒞)`$ is indeed gauge invariant, as it should.
A more intuitive understanding of the anholonomy angle may be obtained if we use Stokes’ theorem to express it as a surface integral. In doing so, however, we should pay attention to the possible existence of holes in the surface $`\mathrm{\Sigma }`$. If this is the case, $`\mathrm{\Sigma }`$ is said to be non simply connected. An example is sketched on Fig. 2, where the surface $`\mathrm{\Sigma }`$ has 2 holes limited by the contours $`𝒞_1`$ and $`𝒞_2`$ (hatched areas on Fig. 2). Applying Stokes theorem, we then obtain
$$\theta (𝒞)=_𝒮B(𝐗)dX_1dX_2+\underset{i}{}N_i(𝒞)\theta (𝒞_i).$$
(19)
where the surface $`𝒮`$ is the subset of the surface $`\mathrm{\Sigma }`$ limited by the circuit (dotted area on Fig. 2), $`𝒞`$, $`N_i(𝒞`$ is the winding number of circuit $`𝒞`$ around the hole $`i`$ (i.e., the difference between the number of turns in counterclockwise and clockwise directions),
$$\theta (𝒞_i)_{𝒞_i}𝐀(𝐗)d𝐗$$
(20)
is the anholonomy angle of circuit $`𝒞_i`$ and
$`B(𝐗)`$ $``$ $`\left({\displaystyle \frac{A_2}{X_1}}{\displaystyle \frac{A_1}{X_2}}\right)`$ (21)
$`=`$ $`\mathrm{Im}\left[{\displaystyle \frac{𝐮^{}}{X_1}}{\displaystyle \frac{𝐮}{X_2}}{\displaystyle \frac{𝐮^{}}{X_2}}{\displaystyle \frac{𝐮}{X_1}}\right].`$
Equation (19) constitutes the *2-form* expression of the anholonomy angle $`\theta (𝒞)`$. One can see immediately that, unlike the connection $`𝐀(𝐗)`$, the quantity $`B(𝐗)`$ is gauge invariant. The geometrical meaning of $`B(𝐗)`$ stems from its relation to the *Gaussian curvature* $`K`$ of $`\mathrm{\Sigma }`$ at point $`𝐗`$, i.e.,
$$B(𝐗)\mathrm{d}X_1\mathrm{d}X_2=K\mathrm{d}S\frac{\mathrm{d}S}{R_1(𝐗)R_2(𝐗)},$$
(22)
where $`R_1(𝐗)`$ and $`R_2(𝐗)`$ are the principal curvature radii at point $`𝐗`$. In the case of the sphere, this is easily checked by explicit calculation, taking the usual spherical angles $`(\theta ,\phi )`$ as variables $`(X_1,X_2)`$. Since the Gaussian curvature is related to the solid angle $`\mathrm{\Omega }`$ of spanned by the normal unit vector $`𝐧`$ by
$$B=\frac{\mathrm{d}^2\mathrm{\Omega }}{\mathrm{d}X_1\mathrm{d}X_2}$$
(23)
we finally get
$$\theta (𝒞)\underset{i}{}N_i(𝒞)\theta (𝒞_i)=_𝒮\frac{\mathrm{d}^2\mathrm{\Omega }}{\mathrm{d}X_1\mathrm{d}X_2}dX_1dX_2=_𝒮\mathrm{d}^2\mathrm{\Omega }=\mathrm{\Omega }(𝒮),$$
(24)
where $`\mathrm{\Omega }(𝒮)`$ is the solid angle described by the normal vector $`𝐧`$ on the surface $`𝒮`$. That the above results hold not only for a sphere, but for any surface can be understood easily from the following argument: Eq. (3) shows that the trajectory of the parallel transported tangent vectors is entirely determined by the trajectory of the normal unit vector $`𝐧`$ along $`𝒞`$. We can therefore map the trajectory $`𝒞`$ on the surface $`\mathrm{\Sigma }`$ to a trajectory $`𝒞^{}`$ on the sphere of unit radius $`S^2`$, by mapping each point of $`\mathrm{\Sigma }`$ onto the point of $`S^2`$ with the same normal vector $`𝐧`$. This implies that we can restrict ourselves to studying the case of parallel transport on $`S^2`$ and obtain conclusions valid for parallel transport on any surface $`\mathrm{\Sigma }`$.
Let us examine these results for the examples sketched on Fig. 1. For the case of the plane, the anholonomy of course trivially vanishes. For the sphere, the anholonomy angle is given by the solid angle $`\mathrm{\Omega }(𝒮)`$, and is therefore a *geometric* property of the circuit $`𝒞`$; this can be easily checked by making the following experiment: take your pen in you left hand, and raise your arm above you head, the pen pointing in front of you; then rotate your arm until it is horizontal in front of you, without twisting your hand; then rotate it by 90<sup>o</sup> to your left; finally rotate your arm back to the vertical (pay attention to never twist your hand in whole process); the pen is now pointing to your left, i.e., it has rotated by $`4\pi /8=\pi /2`$. For the case of the cone, the Gaussian curvature vanishes everywhere (a cone can be fabricated by rolling a sheet of paper), so that the the anholonomy angle is in fact a topological property of the circuit $`𝒞`$, being given by the winding number of the circuit $`𝒞`$ around the cone (multiplied by the solid angle of the cone).
## 3 Parallel transport in classical mechanics: Foucault’s pendulum and the gyroscope
Let us now consider the famous experiment of Foucault’s pendulum that demonstrated the earth’s rotation. If the pendulum trajectory is originally planar (swinging oscillation), the vertical component of the angular momentum vanishes. Since forces exerted on the pendulum (gravity and wire tension) produce a vanishing vertical torque, the vertical component of the angular momentum has to be conserved. The absence of any vertical torque imposes that the swing plane has to follow a law of parallel transport as the direction of gravity slowly changes due to the earth’s rotation. Therefore, within one day it rotates by an angle equal to the solid angle described by the vertical $`2\pi (1\mathrm{cos}\theta )`$, where $`\theta `$ is the colatitude.
The parallel transport may also affect the phase of the periodic motion of the Foucault pendulum or the rotation phase of a gyroscope, but also the phase of their periodic motion. Let us consider a gyroscope whose rotation axis is constrained to remain parallel to the axis $`𝐧`$; let us now move the rotation axis $`𝐧`$ round a closed circuit $`𝒞`$. The rotation angle of the gyroscope will be the sum of the *dynamic rotation angle* $`\omega t`$ and of a *geometric anholonomy angle* $`\theta (𝒞)`$ equal to the solid angle described by the rotation axis. Thus if we have two synchronous gyroscopes and perform different circuits with the rotation axes, they will eventually be dephased with respect to each other, an effect that could easily be observed by stroboscopy. This geometric anholonomy angle is known as Hannay’s angle . If the Foucault pendulum is given a conical oscillation, instead of a planar swing, then we have exactly the same situation as described above for the gyroscope, and the rotation angle will have an anholonomy excess angle given by the solid angle described by the vertical. Thus, two identical Foucault pendula (i.e., of same length) with circular oscillations in opposite directions will have slightly different oscillation frequencies, and will progressively get dephased with respect to each other. The swinging motion of the usual Foucault may be view as the superposition of circular motions in opposite direction, so that the rotation of the swinging plane may be viewed as resulting from the above mentioned frequency shift.
## 4 Parallel transport in quantum mechanics: the Berry phase
Let us now consider a quantum mechanical system described by a Hamiltonian controlled by a set of external parameters $`(R_1,R_2,\mathrm{})`$, which we describe collectively as a vector $`𝐑`$ in some abstract parameter space. Physically, the external parameters may be magnetic or electric fields, etc. For each value $`𝐑`$ of the external parameters, the Hamiltonian $`H(𝐑)`$ has eigenvalues $`E_n(𝐑)`$ and eigenvectors $`|n(𝐑)`$ satisfying the independent Schrödinger equation, i.e.,
$$H(𝐑)|n(𝐑)=E_n(𝐑)|n(𝐑)$$
(25)
The eigenvectors $`|n(𝐑)`$ are defined up to an arbitrary phase, and there is *a priori* no particular phase relation between eigenstates corresponding to different values of the parameter $`𝐑`$. We make a particular choice for the phase of the eigenstates, simply requiring that $`|n(𝐑)`$ varies smoothly with $`𝐑`$ in the region of interest. It may happen that the eigenstates we have chosen are not single valued functions of $`𝐑`$. If this happens, special care must be given to this point.
Let us perform an adiabatic closed circuit $`𝒞\left\{𝐑(t)\right|t=0T\}`$ in the parameter space. The adiabatic theorem tells us that if the rate of variation of the external parameters is low enough, a system that is initially in the $`n`$th stationary state $`|n`$ (assumed non-degenerate) of the Hamiltonian will remain continuously in the state $`|n`$. The condition of adiabaticity is that the stationary state under consideration remains non-degenerate, and the rate of variation of the Hamiltonian is low enough to make the probability of transition to another state $`|m`$ vanishingly small, i.e.,
$$\mathrm{}|m|\dot{H}|n||E_mE_n|^2mn.$$
(26)
Then of course, if one performs a closed adiabatic circuit $`𝒞`$, the system has to return to its original state.
Berry asked the following question: what will be the phase of the state after completion of the circuit $`𝒞`$ ? It may be difficult at first sight to realize that this question may be of any interest. Indeed, the expectation value of any observable quantity $`A`$,
$$A\psi |A|\psi $$
(27)
does not depend of the phase of $`|\psi `$. This lack of interest is certainly the main reason why the Berry phase was (almost <sup>1</sup><sup>1</sup>1Some early precursor work on effects related to the Berry phase include notably Pancharatnam’s work on optical polarization , Aharonov and Bohm’s work on the phase due to the electromagnetic potential vector , and Mead and Truhlar’s work on the molecular Aharonov-Bohm effect in the Born-Oppenheimer theory of molecular vibrations .) completely overlooked for more than half a century of quantum mechanics.
So, following Berry, taking
$$|\psi (t=0)|n(𝐑(t=0))$$
(28)
we express the state $`|\psi (t)`$ at a latter time $`t`$ as
$$|\psi (t)\mathrm{exp}\left[\frac{\mathrm{i}}{\mathrm{}}_0^tdt^{}E_n(𝐫(t^{}))\right]|\varphi _n(t),$$
(29)
i.e., we introduce an auxiliary wavefunction $`|\varphi _n(t)`$ with a zero dynamical phase. Using the time-dependent Schrödinger equation,
$$\mathrm{i}\mathrm{}|\dot{\psi }(t)=H(t)|\psi (t),$$
(30)
and projecting it on $`\psi (t)|`$, we get
$`0`$ $`=`$ $`\psi (t)|\left(H(t)\mathrm{i}\mathrm{}{\displaystyle \frac{}{t}}\right)|\psi (t)`$ (31)
$`=`$ $`\varphi _n(t)|\dot{\varphi }_n(t),`$
where we have used the relation
$$\psi (t)|H(t)|\psi (t)=E_n(t),$$
(32)
which follows from the adiabatic theorem. Equation (31) shows that the wavefunction $`|\varphi _n(t)`$ obeys a quantum mechanical analogue of the law of parallel transport (6).
In complete analogy with the problem of parallel transport on a surface, we now express the parallel transported state $`|\varphi _n(t)`$ in terms of the fixed eigenstates $`|n(𝐑)`$ as
$$|\varphi _n(t)\mathrm{exp}((\mathrm{i}\gamma _n(t))|n(𝐑),$$
(33)
where the phase $`\gamma _n(t)`$ plays the same role as the angle $`\theta (t)`$ for the problem of parallel transport on a surface. We then immediately get the equation of motion of $`\gamma _n(t)`$, i.e.,
$$\dot{\gamma }_n(t)=\mathrm{i}n|\dot{n}=\mathrm{Im}n(𝐑(t))|\frac{\mathrm{d}}{\mathrm{d}t}n(𝐑(t)),$$
(34)
which is analogous to Eq. (11).
Finally, the answer to the question originally asked by Berry is
$$|\psi (T)=\mathrm{exp}\left[\mathrm{i}(\delta _n+\gamma _n(𝒞))\right]|\psi (0),$$
(35)
where
$$\delta _n\frac{1}{\mathrm{}}_0^TE_n(𝐑(t))dt$$
(36)
is the dynamical phase, and
$$\gamma _n(𝒞))\mathrm{Im}[_𝒞n(𝐑)|_𝐑|n(𝐑)\mathrm{d}𝐑]\alpha _n(𝒞)$$
(37)
is the Berry phase. The last term in the latter equation arises when the states $`|n(𝐑)`$ are not a single-valued function of $`𝐑`$ in the region of interest of the parameter space<sup>2</sup><sup>2</sup>2This term was absent in Berry’s original paper , because the basis states $`|n(𝐑)`$ were assumed to be single-valued., and is given by
$$\alpha _n(𝒞)=\mathrm{i}\mathrm{ln}\left[n(𝐑(0))|n(𝐑(T))\right].$$
(38)
We shall omit this term below, and consider only the case of single-valued basis states.
We note the very close analogy between the result obtained for quantum and classical systems. The dynamical phase of a quantum system is analogous to the rotation angle $`\omega T`$ in classical mechanics, whereas the Berry phase is analogous to Hannay’s angle (they both arise from the anholonomy of parallel transport).
Defining the connection $`𝐀^n(𝐑)`$ as
$$𝐀^n(𝐑)\mathrm{Im}\left[n(𝐑)|_𝐑n(𝐑)\right],$$
(39)
we re-express the Berry phase as
$$\gamma _n(𝒞))_𝒞𝐀^n(𝐑)\mathrm{d}𝐑,$$
(40)
which constitutes the 1-form expression of the Berry phase. The latter clearly depends only on the geometry of the circuit $`𝒞`$. The connection $`𝐀^n(𝐑)`$ is not gauge invariant: if we make a new choice for the phase of the reference state, i.e.,
$$|n(𝐑)^{}=\mathrm{exp}(i\mu (𝐑))|n(𝐑),$$
(41)
with a single-valued function $`\mu (𝐑)`$, we obtain a different connection
$$𝐀_{}^{n}{}_{}{}^{}(𝐑)=𝐀^n(𝐑)+_𝐑\mu (𝐑).$$
(42)
However, the Berry phase $`\gamma _n(𝒞)`$ is gauge invariant, as it should.
As for the geometric parallel transport on surfaces, we may obtain a gauge invariant and more transparent expression by transforming the above result to a surface integral by using Stokes’ theorem. Here too, we have to pay attention to the existence of holes in the parameter space: if the parameter space is multiply connected, and if the circuit $`𝒞`$ cannot be continuously deformed to a point<sup>3</sup><sup>3</sup>3A circuit that can be continuously deformed to a point is said to be *homotopic* to a point., we must take into account terms associated with the winding of $`𝒞`$ around holes of the parameter space.
The formulation of the Berry phase as a surface integral in a form that is independent of a particular choice of coordinates of the parameter space generally requires to use the mathematical formalism of differential forms , which is beyond the scope of the present lecture. We can nevertheless obtain a useful result without resorting to any advanced mathematics if we make a suitable choice of coordinates of the parameter space. Let us choose a surface $`𝒮`$ in the parameter space which is bound by the circuit $`𝒞`$, and a parameterization $`(R_1,R_2)`$ of the surface $`𝒮`$. Using Stokes’ theorem we then get
$$\gamma _n(𝒞)=_𝒮B^n(𝐑)dR_1dR_2+\underset{i}{}N_i(𝒞)\gamma _n(𝒞_i),$$
(43)
where $`𝒞_i`$ are the circuits bounding the holes of the parameter space and $`N_i`$ the corresponding winding numbers of the circuit $`𝒞`$ around them, and where
$`B^n(𝐑)`$ $``$ $`\left(_{R_1}A_2^n_{R_2}A_1^n\right)`$ (44)
$`=`$ $`\mathrm{Im}\left[_{R_1}n(𝐑)|_{R_2}n(𝐑)_{R_2}n(𝐑)|_{R_1}n(𝐑)\right]`$
is the *Berry curvature*. In the case where the parameter space is three-dimensional, then we can use the familiar language of vector calculus, as in electrodynamics, and Stokes’ theorem yields
$$\gamma _n(𝒞)=_𝒮𝐁^b(𝐑)𝐧dS+\underset{i}{}N_i(𝒞)\gamma _n(𝒞_i),$$
(45)
$`𝐁^n(𝐑)`$ $``$ $`\times 𝐀^n(𝐑)`$ (46)
$`=`$ $`\mathrm{Im}\left[n(𝐑)|\times |n(𝐑)\right]`$
$`=`$ $`\mathrm{Im}{\displaystyle \underset{mn}{}}n(𝐑)|m(𝐑)\times m(𝐑)|n(𝐑).`$ (47)
Making use of the relation
$$m|n=\frac{m|H|n}{E_nE_m},$$
(48)
one eventually get
$$𝐁^n(𝐑)=\mathrm{Im}\underset{mn}{}\frac{n(𝐑)|H(𝐑)|m(𝐑)\times m(𝐑)|H(𝐑)|n(𝐑)}{\left(E_m(𝐑)E_n(𝐑)\right)^2}.$$
(49)
Obviously, the Berry curvature is gauge invariant. As the notation suggests, the Berry curvature $`𝐁^n`$ plays the role of a magnetic field in the space of parameters, whose vector potential is the Berry connection $`𝐀^n`$.
The energy denominator in Eq. (49) shows that if the circuit $`𝒞`$ lies in a region of the parameter space that is close to a point $`𝐑^{}`$ of two-fold degeneracy involving the two states labelled $`+`$ and $``$, the corresponding Berry connections $`𝐁^+`$ and $`𝐁^+`$ are dominated by the term involving the denominator $`(E_+E)^2`$ and the contribution involving other states can be neglected. So, to first order in $`𝐑𝐑^{}`$, one has
$$𝐁_+(𝐑)=𝐁_{}(𝐑)=\mathrm{Im}\frac{+(𝐑)|H(𝐑^{})|(𝐑)\times (𝐑)|H(𝐑^{})|+(𝐑)}{\left(E_+(𝐑)E_{}(𝐑)\right)^2}.$$
(50)
The general form of the Hamiltonian $`H(𝐑)`$ of a two-level system is (without loss of generality, we may take $`𝐑^{}=0`$)
$$H(𝐑)\frac{1}{2}\left(\begin{array}{cc}Z& X\mathrm{i}Y\\ X+\mathrm{i}Y& Z\end{array}\right),$$
(51)
with eigenvalues
$$E_+(𝐑)=E_{}(𝐑)=\frac{1}{2}R.$$
(52)
This illustrates a theorem due to von Neumann and Wigner , stating that it is necessary to adjust 3 independent parameters in order to obtain a two-fold degeneracy from an Hermitian matrix. The gradient of the Hamiltonian is
$$H=\frac{1}{2}𝝈,$$
(53)
where $`𝝈`$ is the vector matrix whose components are the familiar Pauli matrices. Simple algebra then yields
$$𝐁_+=𝐁_{}=\frac{𝐑}{R^3}.$$
(54)
The above Berry curvature $`𝐁_\pm `$ is the magnetic field in parameter space generated by a Dirac magnetic monopole of strength $`1/2`$. Thus, the Berry phase $`\gamma _\pm (𝒞)`$ of a circuit $`𝒞`$ is given by the flux of the monopole through the surface $`𝒮`$ subtended by the circuit $`𝒞`$, which, by Gauss’ theorem, is nothing else as $`\mathrm{\Omega }(𝒞)`$, where $`\mathrm{\Omega }(𝒞)`$ is the solid angle described by $`𝐑`$ along the circuit $`𝒞`$.
The corresponding vector potential (or Berry connection) $`𝐀_\pm `$ (not calculated here), has an essential singularity along a line (Dirac string) ending at the origin, and carrying a ”flux” of magnitude $`\pm 2\pi `$. The position of the Dirac string can be moved (but not removed!) by a gauge transformation, as sketched on Fig. 3. If the Dirac string happen to cross the cross the surface $`𝒮`$, the Berry phase remains unchanged (modulo $`2\pi `$), so that the result is indeed gauge invariant.
## 5 Examples of Berry phase
### 5.1 Spin in a magnetic field
As a first example, we consider the case of a single spin (of magnitude $`S`$) in a magnetic field, which is both the most immediate application of the formal theory presented above and one of the most frequent case encountered in experimentally relevant situations. The Hamiltonian considered is
$$H(𝐛)𝐛𝐒,$$
(55)
with the magnetic field $`𝐛`$ being the external parameters. The eigenvalues are
$$E_n(𝐛)=nb,$$
(56)
with $`2n`$ integer and $`SnS`$. For $`𝐛=0`$, the $`2S+1`$ eigenstates are degenerate, so the circuit $`𝒞`$ has to avoid the origin. The Berry connection can be calculated using Eq. (49) and well known properties of the spin operators, and one gets
$$𝐁^n(𝐛)=n\frac{𝐛}{b^3},$$
(57)
which is the ”magnetic field” (in parameter space) of a monopole of strength $`n`$ located at the origin. The Berry phase is thus
$$\gamma _n(𝒞)=n\mathrm{\Omega }(𝒞),$$
(58)
where $`\mathrm{\Omega }(𝒞)`$ is the solid angle described by the field $`𝐛`$ along the circuit $`𝒞`$. For $`S=1/2`$, this of course reduces to the result obtained above for the two-level problem. Note that the Berry phase $`\gamma _n(𝒞)`$ depends only on the quantum number $`n`$ (projection of $`𝐒`$ on $`𝐛`$) and not on the magnitude $`S`$ of the spin. Note also, that while $`H(𝐛)`$ is the most general Hamiltonian for a spin $`S=1/2`$, this is not the case for a spin $`S1`$; in the latter case we are restricting ourselves here to a subspace of the full parameter space. If a more general Hamiltonian and a wider parameter space is considered, the simple result obtained above would not hold any more.
### 5.2 Aharonov-Bohm effect
Another example which is a great interest, both conceptually and experimentally, is the well known Aharonov-Bohm effect . We follow here the presentation of the Aharonov-Bohm effect given by Berry .
Let us consider the situation depicted in Fig. 4, namely a magnetic field confined in a tube with flux $`\mathrm{\Phi }`$ and a box located at $`𝐑`$ in which particles of charge $`q`$ are confined. The magnetic field is vanishes everywhere outside the flux tube, and in particular inside the box. Let $`𝐀(𝐫)`$ be the corresponding vector potential. The latter generally does not vanish in the regions of vanishing field (unless the flux $`\mathrm{\Phi }`$ is a multiple of the flux quantum $`\mathrm{\Phi }_0h/e`$.
Let the Hamiltonian describing the particles in the box be $`H(𝐩,𝐫𝐑)`$; the corresponding wave functions, for a vanishing vector potential, are of the form $`\psi _n(𝐫𝐑)`$, with energies $`E_n`$ independent of $`𝐑`$. When the flux is non-zero, we can chose as basis states $`|n(𝐑)`$, satisfying
$$H(𝐩q𝐀(𝐫),𝐫𝐑)|n(𝐑)=E_n|n(𝐑),$$
(59)
whose solutions are given by
$$𝐫|n(𝐑)=\mathrm{exp}\left[\frac{\mathrm{i}q}{\mathrm{}}_𝐑^^𝐫d𝐫^{}𝐀(𝐫^{})\right]\psi _n(𝐫𝐑),$$
(60)
where the integral is performed along a path contained in the box. The energies $`E_n`$ are independent of the vector potential, because it is always possible to find a gauge transformation that would make it zero in the box (but not everywhere in space!).
The Hamiltonian depends on the position $`𝐑`$ of the box via the vector potential. Thus our parameter space, in this example, is nothing else than the real space, with exclusion of the region of the flux tube. If we transport the box around a closed circuit $`𝒞`$, the Berry phase will be given by
$$\gamma _n(𝒞))_𝒞𝐀^n(𝐑)\mathrm{d}𝐑,$$
(61)
with the Berry connection
$`𝐀^n(𝐑)`$ $``$ $`\mathrm{Im}\left[n(𝐑)|_𝐑n(𝐑)\right]`$ (62)
$`=`$ $`\mathrm{Im}{\displaystyle \mathrm{d}^3𝐫\psi _n^{}(𝐫𝐑)\left[\frac{\mathrm{i}q}{\mathrm{}}𝐀(𝐑)\psi _n(𝐫𝐑)+_𝐑\psi _n(𝐫𝐑)\right]}`$
$`=`$ $`{\displaystyle \frac{q}{\mathrm{}}}𝐀(𝐑).`$
The Berry curvature $`𝐁^n(𝐑)=\times 𝐀^n(𝐑)=(q/\mathrm{})𝐁(𝐑)`$ is just given by the magnetic field, and vanishes everywhere outside the flux tube. But because the tube region is excluded from the allowed parameter space, the latter is multiply connected, and the Berry phase is purely topological, being given by the winding number $`N(𝒞)`$ of the circuit $`𝒞`$ around the flux tube, and by the flux $`\mathrm{\Phi }`$:
$$\gamma _n(𝒞)=2\pi N(𝒞)\frac{q}{h}\mathrm{\Phi }.$$
(63)
The Aharonov-Bohm effect was confirmed experimentally by electron holography by Tonomura *et al.* in a configuration where the magnetic truly vanishes, and plays an outstanding role in the physics of mesoscopic systems, where it gives rise to conductance oscillations and to persistent currents in mesoscopic metallic rings threaded by a magnetic flux .
## 6 Experimental observations of the Berry phase for a single spin
Let us now discuss how the Berry phase could be detected experimentally. As already mentioned, this is not immediately clear since the expectation value of any observable would be independent of the phase of the system. As always when considering phases, some kind of interference has to be observed. There various ways in which this can be done.
* Berry original proposal was the following: a monoenergetic polarized beam of particles in the spin state $`n`$ along the magnetic field $`𝐛`$ is split in two beams. For one of the beams, the field $`𝐛`$ is kept constant in magnitude and direction, whereas in the second beam, the magnitude of $`𝐛`$ is kept constant and its direction slowly varied along a circuit $`𝒞`$ subtending a solid angle $`\mathrm{\Omega }`$. The two beams are then recombined to interfere, and the intensity is monitored as a function of the solid angle $`\mathrm{\Omega }`$. Since the dynamical phase is the same for both beams, the phase difference between the two beams is given purely by the Berry phase (plus a propagation factor is determined by the phase shift for $`\mathrm{\Omega }=0`$. Although conceptually possible, it seems unlikely that such an experiment would be feasible in practice. In particular, it would be extremely difficult to ensure that the difference between the dynamical phases of the two beams be smaller that the Berry phase one wants to detect, unless some physical principle enforces it. This kind of experiment may be said to be of type ”one state – two Hamiltonians”. This kind of experiment, being based on interferences is truly quantum mechanical.
* An alternative approach, more amenable to experimental test is to prepare the system into a superposition of two states, i.e.,
$$|\psi (t=0)=\alpha |n(𝐑(t=0))+\beta |m(𝐑(t=0)),$$
(64)
with $`m=n1`$ and $`|\alpha |^2+|\beta |^2=1`$, for example by polarizing it along a direction perpendicular to the field $`𝐛`$. The orientation of the transverse component of the spin is given by the angle $`\theta (t=0)\mathrm{arg}(\beta )\mathrm{arg}(\alpha )`$. The spin of course precesses at around $`𝐛`$ at the Larmor frequency $`\omega _L=b/\mathrm{}`$. After completion of the circuit $`𝒞`$, the system state has evolved to
$$|\psi (T)=\alpha \mathrm{exp}[\mathrm{i}(\delta _n+\gamma _n(𝒞))]|n(𝐑(t=0))+\beta \mathrm{exp}[\mathrm{i}(\delta _m+\gamma _m(𝒞))]|m(𝐑(t=0)),$$
(65)
and the polarization angle has evolved to $`\theta (T)=\theta (t=0)+\mathrm{\Delta }\theta `$ with
$`\mathrm{\Delta }\theta `$ $`=`$ $`\mathrm{\Delta }\theta _{\mathrm{dyn}}+\mathrm{\Delta }\theta _\mathrm{B},`$ (66)
$`\mathrm{\Delta }\theta _{\mathrm{dyn}}`$ $``$ $`\delta _m\delta _n=\omega _LT,`$ (67)
$`\mathrm{\Delta }\theta _\mathrm{B}`$ $``$ $`\gamma _m(𝒞))\gamma _n(𝒞)).`$ (68)
Here the angle $`\mathrm{\Delta }\theta _{\mathrm{dyn}}`$ gives the polarization rotation due to the Larmor precession (dynamic phase), while $`\mathrm{\Delta }\theta _\mathrm{B}`$ is the polarization rotation due to the Berry phase accumulated along the circuit $`𝒞`$. Thus by investigating how the polarization varies as the circuit $`𝒞`$ is modified, the Berry phase can be detected. Such an experiment may be said to be of the type ”two states – one Hamiltonian”. Note that this type of experiment can be interpreted in purely classical terms (it bears a clear analogy to the rotation of swinging plane of the Foucault pendulum); this is related to the fact that only Berry phase differences between two states, and not the absolute Berry phase of a given state is detected.
* A further possibility consists in repeating the circuit $`𝒞`$ in a periodic manner. Thus, the Berry phase is accumulated linearly in time, just as the dynamical phase, and leads to an apparent energy shift for the state $`n`$,
$$\mathrm{\Delta }E_n=\frac{\mathrm{}}{T}\gamma _n(𝒞),$$
(69)
which gives rise to an observable shift of the transition between to levels $`n`$ and $`m`$. Such an experiment is of type ”two states – one Hamiltonian”, too. It can also be interpreted in classical terms and has close analogy to the period shift of a Foucault pendulum with circular oscillation.
### 6.1 Berry phase of neutrons
The Berry phase has been observed for neutrons by Bitter and Dubbers , who used the experimental shown in Fig. 5. A slow ($`v500\mathrm{m}.\mathrm{s}^1`$), monochromatic, beam neutrons polarized ($`P0.97`$) along an axis perpendicular to the beam axis $`z`$ is injected in a cylinder with a helical magnetic field with longitudinal component $`B_z`$ and transverse component $`B_1`$making a right-handed turn of $`2\pi `$. Depending on the values of $`B_z`$ and $`B_1`$, various values of the solid angle $`\mathrm{\Omega }`$ may be achieved.
After having traversed the cylinder, the polarization of the beam is measured, from which the Berry phase can be extracted. The comparison of the measured Berry phase (or more precisely the difference of Berry phase between states $`S_z=+1/2`$ and $`S_z=1/2`$) and of the solid angle is shown in Fig. 5. The observation is in good agreement with the theoretical prediction.
### 6.2 Berry phase of photons
The photon is a particle of spin $`S=1`$ and can thus experience a Berry phase. The particularity of the photon is that, being massless, only the states $`S_z=\pm 1`$ occur and that the quantization axis is fixed by the direction of the wave vector $`𝐤`$. The wavevector therefore plays the role of a magnetic field for the photon .
If the latter is constrained to make a closed circuit $`𝒞`$ of solid angle $`\mathrm{\Omega }`$, then a Berry phase of $`\pm \mathrm{\Omega }`$ is expected for the two circular polarizations, respectively. If a monochromatic, linearly polarized, optical wave
$$|\chi =\frac{|++|}{\sqrt{2}},$$
(70)
where $`|+`$ and $`|`$ represent, respectively the two circular polarization modes, is injected in a single mode optical fiber, whose axis describes a helix of solid angle $`\mathrm{\Omega }`$, then the emerging wave will be (omitting the dynamical phase)
$$|\chi ^{}=\frac{\mathrm{exp}(\mathrm{i}\gamma _+)|++\mathrm{exp}(\mathrm{i}\gamma _{})|}{\sqrt{2}},$$
(71)
with $`\gamma _+=\gamma _{}=\mathrm{\Omega }`$, which yields
$$|\chi ^{}|\chi |^2=\mathrm{cos}^2\mathrm{\Omega }.$$
(72)
By Malus’ law, this means that the polarization has rotated by an angle $`\mathrm{\Omega }`$ (the sense of rotation, when looking into the output of the fiber is counterclockwise, i.e., dextrorotatory, for a left-handed helix) .
The experiment carried out by Tomita and Chiao is shown in Fig. 6, and shows a very good agreement between theoretical prediction and experimental results. Note that this kind of experiment can also be explained entirely from classical electrodynamics considerations .
### 6.3 Berry phase effects in nuclear magnetic resonance
Nuclear spins interact very weakly with each other and with their environment and therefore offer constitute systems that ideally suited to test the Berry phase of a single spin. The experiment described below has been performed on protons ($`S=1/2`$) by Suter *et al.* following a proposal of Moody *et al.* .
As in a typical nuclear magnetic resonance (NMR) experiment, the spins are subject to a large, static, orienting field parallel to the $`z`$ axis, and to a weak, transverse, field rotating around $`𝐁_0`$ at angular frequency $`\omega _{\mathrm{RF}}`$. For convenience, we express here energies and magnetic fields in units of angular frequencies, i.e., the Hamiltonian, expressed in the laboratory frame, reads
$$H(t)=\omega _LS_z\omega _1[S_x\mathrm{cos}(\omega _{\mathrm{RF}}t)+S_y\mathrm{sin}(\omega _{\mathrm{RF}}t)].$$
(73)
The measured signal is the transverse magnetization $`S_x(t)+\mathrm{i},S_y(t)`$. In the present case, it is of convenient to perform a transformation from the laboratory frame to a detector frame, rotating at angular frequency $`\omega _D`$. In practice, this is achieved by beating the measured signal against a reference signal of angular frequency $`\omega _D`$. In the detector frame, the Hamiltonian now reads
$$H^{}(t)=(\omega _L\omega _D)S_z\omega _1[S_x\mathrm{cos}((\omega _{\mathrm{RF}}\omega _D)t)+S_y\mathrm{sin}((\omega _{\mathrm{RF}}\omega _D)t)].$$
(74)
Let us define
$$\omega _{\mathrm{eff}}\sqrt{(\omega _L\omega _D)^2+\omega _1^2},$$
(75)
which is the magnitude of the total field in the detector frame, at angle $`\theta \mathrm{arcsin}(\omega _1/\omega _{\mathrm{eff}})`$ from the $`z`$ axis and precessing around the $`z`$ axis at angular frequency $`\omega _{\mathrm{RF}}\omega _D`$. In the adiabatic limit, i.e., if $`|\omega _{\mathrm{RF}}\omega _D|\omega _{\mathrm{eff}})`$, the adiabatic eigenstates have an energy $`\omega _n=n\omega _{\mathrm{eff}}`$ ($`n=\pm 1/2`$). For each cycle of the effective field around the $`z`$ axis, the state $`n`$ acquires a Berry phase $`\gamma _n=n2\pi (1\mathrm{cos}\theta )`$. Thus if we prepare the system in a superposition of the two states $`n=\pm 1/2`$, the Fourier spectrum of the transverse magnetization will have a component of angular frequency $`\omega _{\mathrm{eff}}+(\omega _{\mathrm{RF}}\omega _D)2\pi (1\mathrm{cos}\theta )`$, where the last term arising from the Berry phase, as shown in Fig. 7.
## 7 Berry phase for itinerant electrons in a solid
### 7.1 General formulation
We now want to discuss the Berry phase of electrons in solids. Let us consider (non-interacting) electrons subject to a scalar potential and to a Zeeman (or exchange) field, whose direction is spatially nonuniform. As they move through the solid, the electrons experience, in their proper reference frame, a Zeeman field whose direction changes with time. If this change is slow enough, the electron spin has to follow it adiabatically and therefore accumulates a Berry phase .
Before discussing the resulting physical consequences, let us formulate the problem more precisely. I follow here the discussion of Ref. . The corresponding Hamiltonian is
$$H=\frac{\mathrm{}^2}{2m}\frac{^2}{𝐫^2}+V(𝐫)\mathrm{\Delta }(𝐫)𝐧(𝐫)𝝈.$$
(76)
We use a gauge transformation $`T(𝐫)`$, which makes the quantization axis oriented along vector $`𝐧(𝐫)`$ at each point. It transforms the last term in the above equation as $`T^{}(𝐫)\left[𝝈𝐧(𝐫)\right]T(𝐫)=\sigma _z`$, corresponding to a local rotation of the quantization axis from $`z`$ axis to the axis along $`𝐧(𝐫)`$. The transformed Hamiltonian describes the electrons moving in a (spinor) gauge potential $`𝐀(𝐫)`$,
$$H^{}T^{}HT=\frac{\mathrm{}^2}{2m}\left(\frac{}{𝐫}\frac{ie}{\mathrm{}c}𝐀(𝐫)\right)^2+V(𝐫)\mathrm{\Delta }(𝐫)\sigma _z,$$
(77)
where $`A_i(𝐫)=2\pi i\varphi _0T^{}(𝐫)_iT(𝐫)`$, $`\varphi _0=hc/\left|e\right|`$ is the flux quantum. For convenience, we have defined the gauge potential $`𝐀(𝐫)`$ to have the same dimension as the electromagnetic vector potential. The components of $`𝐀(𝐫)`$ can be found easily using an explicit form of $`T(𝐫)`$.
The above Hamiltonian with the matrix $`𝐀(𝐫)`$ contains terms coupling the two spin states. If the rate at which the spin-quantization axis varies is slow enough as compared to the Larmor precession frequency (as seen from the moving electron’s frame), the spin will follow adiabatically the local spin-quantization axis, and these spin-flip terms can be neglected. The variation rate of the spin quantization axis is (for an electron at the Fermi level) $`v_F/\xi `$, where $`\xi `$ is the characteristic length of variation of the spin-quantization axis, so that the adiabaticity condition is
$$\frac{\mathrm{}v_F}{\xi }\mathrm{\Delta }.$$
(78)
With this approximation we now obtain
$$H^{}\frac{\mathrm{}^2}{2m}\left(\frac{}{𝐫}\frac{ie}{\mathrm{}c}𝐚(𝐫)\sigma _z\right)^2+V(𝐫)+V^{}(𝐫)\mathrm{\Delta }(𝐫)\sigma _z,$$
(79)
where
$$a_i(𝐫)\frac{\pi \varphi _0\left(n_x_in_yn_y_in_x\right)}{1+n_z}$$
(80)
is an effective vector potential arising from the Berry phase, and
$$V^{}(𝐫)=(\mathrm{}^2/8m)\underset{i,\mu }{}(_in_\mu )^2$$
(81)
is an effective scalar potential. Since the Hamiltonian is diagonal in spin, we can treat the two spin subbands separately. For each of the spin subband, we have mapped the original problem on the simpler one of a spinless electron moving in effective scalar and vector potentials. The effective vector potential in turn gives rise to an effective magnetic field $`𝐛\times 𝐚`$, whose components are expressed in terms of the magnetization texture as
$$B_i=\frac{\varphi _0}{4\pi }\epsilon _{ijk}\epsilon _{\mu \nu \lambda }n_\mu (_jn_\nu )(_kn_\lambda ).$$
(82)
The physical consequences of the effective vector potential and effective magnetic field are exactly the same as those of the familiar electro-magnetic counterparts, and can be classified in two different classes:
1. local effects such as the Lorentz force. These effects are classical in origin (see the illuminating discussion on this point given by Aharonov and Stern ), and therefore do not require phase coherence.
2. non-local interference effects, such as Aharonov-Bohm-like effects and persistent currents. These effects are intrinsically quantum mechanical, and require phase coherence.
### 7.2 Anomalous Hall effect due to the Berry phase in a textured ferromagnet
The Hall effect consists in the appearance of a voltage transverse to the current in a conducting system. As it is antisymmetric with respect to time reversal, it may appear only in the presence of a term in the Hamiltonian breaking time reversal invariance. Until recently, two different mechanism were recognize to give rise to the Hall effect:
1. the electro-magnetic Lorentz force due to a usual magnetic field; in the classical regime (normal Hall effect), this is well described by the Drude theory; in the quantum limit, the spectacular quantum Hall effect is obtained.
2. in absence of an external magnetic field, time-reversal invariance is also broken if the system exhibits magnetic order. However, this fact is not enough to induce a Hall effect if the magnetic order is uniform and spin-orbit coupling is absent or negligible, because the system is then invariant by under a global rotation of $`\pi `$ of the spins, which is equivalent to time reversal in this case (uniform magnetization). Therefore Hall effect arises only as a consequence of simultaneous presence of spontaneous magnetic order and spin-orbit coupling. This mechanism is called the anomalous Hall effect.
Recently, however, it was realized that a third mechanism could give rise to the Hall effect in ferromagnetic, in absence of an external magnetic field, and of the spin-orbit coupling, if the magnetization is non-uniform and exhibits a non-trivial texture .
The central idea is the following: as we have discussed above that if the spin-orbit coupling is negligible, the system is invariant under a global rotation of the magnetization. However, if the magnetization is non-uniform, a rotation of $`\pi `$ is generally not equivalent to a time-reversal, so that Hall effect may arise. At the microscopic level, the origin of the Hall effect in such a case is the effective Lorentz force due to the Berry phase of the magnetization texture.
We shall discuss below as an example the results of Taguchi *et al.* . In this work the authors investigated the compound Nd<sub>2</sub>Mo<sub>2</sub>O<sub>7</sub>, which has the pyrochlore structure shown on Fig. 8. Due to their large spin-orbit coupling the Nd 4f moments of a given tetrahedron adopt the “two-in, two-out” structure shown on Fig. 8B. The Mo 4d moments which are ferromagnetically coupled to each other and antiferromagnetically to the the Nd moments therefore adopt and non-collinear umbrella structure, whose chirality gives rise to the anomalous Hall effect: as electrons moves on a triangle face of a tetrahedron, they acquire a Berry phase, and experience the associated Lorentz force. This mechanism is dominant at low temperature, where other mechanisms due to the spin-orbit coupling (giving contributions to the Hall resistivity linear or quadratic in the longitudinal resistivity) are strongly suppressed. That the origin of the anomalous Hall effect is the Berry phase due to the texture is further indicated by noting that the application of a magnetic field reduces the solid angle of the umbrella texture and hence the Berry phase and the associated Hall effect.
One should note, however, that the average effective magnetic field due to the Berry phase is zero on the present case. This can be understood by noting that the Mo planes perpendicular to (111) axes are kagomé lattices, with a Berry phase of $`+\phi `$ on the triangles and $`2\phi `$ on the hexagons. Since there are two triangles and one hexagon per unit cell the effective magnetic field due to the Berry phase is zero on average. However, since the circuits with positive (triangles) and negative (hexagons) Berry phase are inequivalent, time reversal invariance is still broken, and a non zero Hall effect may (and generally does) result . However, because the effective field due to the Berry phase vanishes on average, the resulting Hall effect is considerably weaker than that one would obtain for a non-zero net effective field.
One should also notice that, as in the example discussed above, spin-orbit coupling always has to be invoked in order to obtain a spin-texture with a definite chirality and, hence, a non-zero Hall effect . This is, however, quite different from the mechanisms of Hall effect in which the spin-orbit coupling directly influences the motion of the conduction electrons.
Until recently, all cases of anomalous Hall due to a chiral spin texture discussed in the literature considered a texture at the atomic level. Recently, it has been proposed that a mesoscopic scale magnetic field texture can be produced by using magnetic nanostructures . For example, one can consider an array of magnetic nanocylinders (all magnetized in the same direction along their axis), as shown in Fig. 11 (left panel), in order to generate a non-uniform dipolar stray field in a two-dimensional electron (or hole) gas placed just underneath. Such structure can be fabricated by electrolytic techniques, as shown on Fig. 11 (right panel).
The Fig. 12 shows the $`z`$-component of the dipolar stray field (left panel) and topological field generated by the Berry phase (right). It is noteworthy that for a lattice of Fe cylinders with a pitch of 100 nm, the dipolar field at a distance 20 nm underneath (the average value of which is always zero) has a maximum absolute value of about 2 kG, whereas the topological field has an non-zero average value of about 5 kG, with local values ranging between $`5`$ and $`+15`$ kG. Thus, one sees that a rather weak dipolar stray field averaging to zero generates in the two-dimensional electron gas a much stronger topological field with a non-zero average!
The topological field has a number of interesting properties:
* First, one can show that the total flux of the topological field through a unit cell is always an integer multiple of the flux quantum. To see this, let us consider the Berry phase corresponding to circuit going around a unit cell (e.g., along the path ABCDA on Fig. 13 (left panel). Because of the translational periodicity of the system, the local field points along the same direction at the four corners A, B, C, D of the unit cell, and thus correspond to the same point (represented by the solid dot) on the sphere of unit radius on which the corresponding path for the field direction takes place; the path ABC is shown, and one sees immediately that, because of the translational periodicity, the path CDA is exactly the same, described in the reversed direction, so that the Berry phase corresponding to the path ABCDA is zero modulo $`2\pi `$ which implies the quantization of the total flux in units of $`\varphi _0`$. This constitutes an example of topological quantization.
* Next, let us try to understand what is the actual integer value taken by the total topological flux. The above reasoning does tell us anything about it, because we don’t know how many times we are wrapping the sphere when integrating over the unit cell. To know this, we may consider the closed lines where the $`z`$-component of the stray field vanishes. This lines corresponds to paths going an integer number of times around the “equator” of the unit sphere. Such lines are shown on Fig. 12 (left panel, solid lines). One can then see that each round trip around the equator contributes to one flux quantum. In counting this, one should be careful in getting the sign correctly. In the case shown in Fig. 12 for a vanishing external field, one obtains that the total topological flux per unit cell is $`+\varphi _0`$.
* Finally, one can see that the lines of zero $`z`$-component of the stray field will change if one applies an external magnetic field. For example, in the case considered here, under application of a uniform external field along the magnetization direction of the nanocylinders, the regions of positive $`z`$-field will expand, so that the topology of the lines of zero $`z`$-field will eventually change. This is shown by the dashed lines on the left panel of Fig.12 ($`B_{\mathrm{ext}}/4\pi M_s=0.058`$). From the above discussion, this implies a change in the topological flux per unit cell, which now takes the value $`2\varphi _0`$ (the factor 2 is because there are 2 lines per unit cell, and the factor $`1`$ is because the circulation is reversed.
The properties discussed above indicate that by application an a rather small uniform external field, one can change the average value of the topological field. This will give rise to changes in the anomalous Hall effect. As a first approximation, one can estimate the Hall effect by using the Drude model and by neglecting the spacial fluctuations of the dipolar and topological fields. In this case, the Hall effect is just given by the familiar Drude formula, with an effective field equal to the sum of the external magnetic field and the average topological field. If the two spin subbands contribute, one simply has to sum the contributions of the two subbands, taking into account the fact that the topological field is of opposite sign for the two spin subbands. For the case discussed here, and assuming that only one spin subband is occupied, the resulting behavior of the Hall effect is shown on Fig. 14. For the chosen parameters, he values of the critical fields are $`B_12`$ kG, $`B_20.9`$ kG, and $`B_31.3`$ kG. The uniform slope comes from the normal Hall effect (Lorentz force of the external magnetic field), whereas the remaining non-monotonous contribution arises from the Lorentz force due to the topological field of the Berry phase. Such a characteristic non-monotonous behavior would constitute a signature of the Berry phase contribution of the Hall effect and can be tested experimentally. One should point out, however, that in the vicinity of the critical fields where the topological flux abruptly changes, the adiabaticity condition cannot be well satisfied, so that, in practice, a rounded curve would be obtained.
In order to identify a system for which the above predictions can be satisfied, we have to look for a system with a large Zeeman splitting, in order to satisfy as well as possible the condition of adiabaticity. We propose to use II VI dilute magnetic semiconductors (DMS) which exhibit giant Zeeman splitting; p-type Mn doped DMS are best suited since the exchange constants for holes are much larger than for electrons . For a detailed discussion see Ref. .
Before closing this section on the anomalous Hall effect, I wish to point that, while the Berry phase allowed to identified a new mechanism of Hall effect arising from the chirality of the spin texture, in absence of an external field and of the spin-orbit coupling, as we have discussed above, it also allows to give a modern interpretation to the previously known mechanisms for the Hall effect (normal Hall effect, classical or quantized; anomalous Hall effect), due either to an external magnetic field, or to the spin-orbit coupling. In this case, one deals with Berry phase in momentum space instead of real space as discussed above. For a detailed discussion, see Refs. .
### 7.3 Interference effects due to the Berry phase in an Aharonov-Bohm ring
As mentioned earlier, the Berry phase accumulated by electrons moving in a non-trivial magnetic texture can give rise to interference effects, of which the archetype is the Aharonov-Bohm effect. It has been proposed by Loss *et al.* and by Stern that a metallic ring subject to a textured magnetic field (or magnetization) as depicted in Fig. 15 would yield a Berry phase for an electron moving around the ring, and hence a dependence of the conductance of the ring (when connected to current leads) upon the solid angle described by the magnetization , as well as to persistent charge and spin currents (for a non-connected ring) .
So far, it has not been possible to test experimentally this prediction in the configuration described above (i.e., by using a textured magnetic field or magnetization). However, several authors have indicated that a similar Berry phase may be obtained by using the combined effect of the Zeeman coupling to a uniform magnetic field (parallel to the ring axis) and of Rashba-type spin-orbit coupling , as described by the following Hamiltonian
$$H=\frac{𝐩^2}{2m}B\sigma _z+\alpha (𝐩\times 𝝈)\widehat{𝐳}+V(𝐫).$$
(83)
In the above equation, the third term gives the Rashba spin-orbit coupling, while the last one confines the electron to the ring. The Rashba effect acts as an effective magnetic field perpendicular to the plane and to the direction of motion.
There is, however, a crucial difference with respect to a true magnetic field, namely that the Rashba term is invariant under time reversal (unlike a true magnetic field), which is manifest from the fact the associated effective field changes sign as the motion is reversed. Because of this, there is no phase shift from the Berry phase between the paths going through the upper and lower arms of the ring, as sketched on Fig. 16, unlike what would be expected for a real magnetic field as in Fig. 15. Therefore, Aharonov-Bohm like interferences in this configuration have to involve paths that wind completely the ring a different number of times. The associated Berry phase will of course be superimposed to the usual Aharonov-Bohm phase and therefore modify the Aharonov-Bohm oscillations of the ring conductance versus magnetic field. Such observations, indicating the presence of the Berry phase, have been made by various groups . The results of Mopurgo *et al.* are shown in Fig. 17, where the signature of the Berry phase is given by the splitting of the $`e/h`$ peak in the average Fourier spectrum.
## 8 Further effects of Berry phase in magnetism
In closing these lecture notes, I wish to briefly mention some further developments and applications of the concept of Berry phase in magnetism.
In the previous section, we mentioned that the Berry phase can give rise to interference phenomena for interfering paths in real space. There may be also interferences associated with different paths in spin space as well; this plays an important role in the theory of tunnelling of magnetization in large spin molecular magnets . The situation is sketched in Fig. 18. Depending on the value of the spin, and on the solid angle between the 2 tunnelling trajectories from state A to state B, interference due to the Berry phase take place; in special cases, the interferences are destructive and tunnelling becomes forbidden. This gives rise to very spectacular parity effect that have been observed experimentally . For a detailed review of this topic, see Ref. .
The Berry phase plays a ubiquitous role in quantum mechanical problems where one wants to treat the dynamics of some “slow” degrees of freedom, after having “integrated out” the “fast” degrees of freedom. An application of this concept in magnetism concerns the adiabatic spin-wave dynamics of itinerant magnets. Here the fast degrees of freedom are the electron degrees of freedom giving rise to charge fluctuations and longitudinal spin fluctuations, whereas the slow degrees of freedom are the transverse spin fluctuations, i.e., the long wavelength magnons. This was pioneered by Wen and Zee , and later on further developed by Niu *et al.* . They obtained an equation of motion for the spins which is controlled by the Berry phase; in the case of localized systems, this reduces to the Landau-Lifshitz equation, but contains non-local contributions in the case of strongly delocalized systems.
For a spin $`S`$ coupled to a slowly moving magnetic field, we have seen that the Berry phase is given by the solid angle described by the field. For the case of a spin $`S=1/2`$, this situation constitutes the most general case, however, for larger spins $`S1`$, the most general Hamiltonian may contain more further contributions, such as anisotropy terms, so that the parameter space is much larger and richer than for a spin $`1/2`$. It is therefore of interest to investigate the Berry phase in this more general context. Recently, this has been investigating, by considering more specifically the Berry phase associated with global rotations of anisotropic spin systems . This study reveals that beside the familiar solid angle term, these is also a topological term, related to a winding number giving the number of rotations of the systems around its magnetization axis. This is relevant to any spin system of spin $`S1`$, in particular to magnons ($`S=1`$), holes in semiconductors ($`S=3/2`$), etc. A general theory of the Berry phase of magnons is in preparation . Interesting spin-wave interference phenomena have recently be obtained by Hertel *et al.* from micromagnetic simulations .
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# I Introduction
## I Introduction
Hadronic matter under extreme conditions has attracted a lot of interest in recent years. On the one hand, many theoretical and experimental efforts have been devoted to the discussion of heavy ion collisions where the temperature is high. On the other hand, the physics of neutron stars has become a hot topic which connects astrophysics with high density nuclear physics. In 1934, Baade and Zwicky suggested that neutron stars could be formed in supernovae. The first theoretical calculation of a neutron star was performed by Oppenheimer and Volkoff , and independently by Tolman . Observing a range of masses and radii of neutron stars will reveal the equations of state (EOS) of dense hadronic matter. Determination of the EOS of neutron stars has been an important goal for more than two decades. Six double neutron-star binaries are known so far, and all of them have masses in the surprisingly narrow range $`1.36\pm 0.08M_{\mathrm{sun}}`$ . A number of early theoretical investigations on neutron stars were based on the non-relativistic Skyrme framework . Since the Walecka model was proposed and applied to study the properties of nuclear matter, the relativistic mean field approach has been widely used in the determination of the masses and radii of neutron stars. These models lead to different predictions for neutron star masses and radii . For a recent review see Ref. . Though models with maximum neutron star masses considerably smaller than $`1.4M_{\mathrm{sun}}`$ are simply ruled out, the constraint on EOS of nuclear matter (for example, the density dependence of pressure of hadronic system) has certainly not been established from the existing observations.
In the process of neutron star formation, $`\beta `$-equilibrium can be achieved. As a consequence, hyperons will exist in neutron stars, especially in stars with high baryon density. These hyperons will affect the EOS of hadronic matter. As a result, the mass-radius relationship of strange hadronic stars will be quite different from that of pure neutron stars. The simplest way to discuss the effects of hyperons is to study strange hadronic stars including only $`\mathrm{\Lambda }`$ hyperons. This is due to the fact that $`\mathrm{\Lambda }`$ is the lightest hyperon and the $`\mathrm{\Lambda }`$-N interaction is known better than other hyperon-nucleon interactions. However, one must also consider hyperons with negative charge in neutron stars because the negatively charged hyperons can substitute for electrons. There have been many discussions of strange hadronic stars including $`\mathrm{\Lambda }`$ hyperons, $`\mathrm{\Lambda }`$ and $`\mathrm{\Sigma }^{}`$ or even the whole baryon octet -.
At high baryon density, the overlap effects of baryons are very important and the quark degrees of freedom within baryons should be considered. There are some phenomenological models based on the quark degrees of freedom, such as the quark meson coupling model , the cloudy bag model , the quark mean field model and the NJL model . Several years ago, a chiral $`SU(3)`$ quark mean field model was proposed . In this model, quarks are confined within baryons by an effective potential. The quark-meson interaction and meson self-interaction are based on $`SU(3)`$ chiral symmetry. Through the mechanism of spontaneous symmetry breaking the resulting constituent quarks and mesons (except for the pseudoscalars) obtain masses. The introduction of an explicit symmetry breaking term in the meson self-interaction generates the masses of the pseudoscalar mesons which satisfy the partially conserved axial-vector current (PCAC) relations. The explicit symmetry breaking term in the quark-meson interaction gives reasonable hyperon potentials in hadronic matter. This chiral $`SU(3)`$ quark mean field model has been applied to investigate nuclear matter , strange hadronic matter , finite nuclei, hypernuclei , and quark matter . Recently, we improved the chiral $`SU(3)`$ quark mean field model by using the linear definition of effective baryon mass . This new treatment is applied to study the liquid-gas phase transition of asymmetric nuclear system and strange hadronic matter . By and large the results are in reasonable agreement with existing experimental data.
In this paper, we will study the neutron star and strange star in the chiral $`SU(3)`$ quark mean field model. The paper is organized in the following way. In section II, we briefly introduce the model. In section III, we apply this model to investigate the neutron star and strange hadronic star. The numerical results are discussed in section IV. We summarize the main results in section V.
## II The model
Our considerations are based on the chiral $`SU(3)`$ quark mean field model (for details see Refs. ), which contains quarks and mesons as the basic degrees of freedom. In the chiral limit, the quark field $`\mathrm{\Psi }`$ can be split into left and right-handed parts $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$: $`\mathrm{\Psi }=\mathrm{\Psi }_L+\mathrm{\Psi }_R`$. Under $`SU(3)`$$`{}_{L}{}^{}\times `$ $`SU(3)`$<sub>R</sub> they transform as
$$\mathrm{\Psi }_L\mathrm{\Psi }_L^{}=L\mathrm{\Psi }_L,\mathrm{\Psi }_R\mathrm{\Psi }_R^{}=R\mathrm{\Psi }_R.$$
(1)
The spin-0 mesons are written in the compact form
$$\genfrac{}{}{0pt}{}{M}{M^{}}=\mathrm{\Sigma }\pm i\mathrm{\Pi }=\frac{1}{\sqrt{2}}\underset{a=0}{\overset{8}{}}\left(s^a\pm ip^a\right)\lambda ^a,$$
(2)
where $`s^a`$ and $`p^a`$ are the nonets of scalar and pseudoscalar mesons, respectively, $`\lambda ^a(a=1,\mathrm{},8)`$ are the Gell-Mann matrices, and $`\lambda ^0=\sqrt{\frac{2}{3}}I`$. The alternatives indicated by the plus and minus signs correspond to $`M`$ and $`M^{}`$, respectively. Under chiral $`SU(3)`$ transformations, $`M`$ and $`M^{}`$ transform as $`MM^{}=LMR^{}`$ and $`M^{}M^{^{}}=RM^{}L^{}`$. The spin-1 mesons are arranged in a similar way as
$$\genfrac{}{}{0pt}{}{l_\mu }{r_\mu }=\frac{1}{2}\left(V_\mu \pm A_\mu \right)=\frac{1}{2\sqrt{2}}\underset{a=0}{\overset{8}{}}\left(v_\mu ^a\pm a_\mu ^a\right)\lambda ^a$$
(3)
with the transformation properties: $`l_\mu l_\mu ^{}=Ll_\mu L^{}`$, $`r_\mu r_\mu ^{}=Rr_\mu R^{}`$. The matrices $`\mathrm{\Sigma }`$, $`\mathrm{\Pi }`$, $`V_\mu `$ and $`A_\mu `$ can be written in a form where the physical states are explicit. For the scalar and vector nonets, we have the expressions
$`\mathrm{\Sigma }={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{a=0}{\overset{8}{}}}s^a\lambda ^a=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\left(\sigma +a_0^0\right)\hfill & a_0^+& \hfill K^+\\ a_0^{}\hfill & \frac{1}{\sqrt{2}}\left(\sigma a_0^0\right)& \hfill K^0\\ K^{}\hfill & \overline{K}^0& \hfill \zeta \end{array}\right),`$ (7)
$`V_\mu ={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{a=0}{\overset{8}{}}}v_\mu ^a\lambda ^a=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\left(\omega _\mu +\rho _\mu ^0\right)\hfill & \rho _\mu ^+& \hfill K_\mu ^+\\ \rho _\mu ^{}\hfill & \frac{1}{\sqrt{2}}\left(\omega _\mu \rho _\mu ^0\right)& \hfill K_\mu ^0\\ K_\mu ^{}\hfill & \overline{K}_\mu ^0& \hfill \varphi _\mu \end{array}\right).`$ (11)
Pseudoscalar and pseudovector nonet mesons can be written in a similar fashion.
The total effective Lagrangian is written:
$`_{\mathrm{eff}}=_0+_{qM}+_{\mathrm{\Sigma }\mathrm{\Sigma }}+_{VV}+_{\chi SB}+_{\mathrm{\Delta }m_s}+_h,+_c,`$ (12)
where $`_0=i\overline{\mathrm{\Psi }}\gamma ^\mu _\mu \mathrm{\Psi }`$ is the free part for massless quarks. The quark-meson interaction $`_{qM}`$ can be written in a chiral $`SU(3)`$ invariant way as
$`_{qM}=g_s\left(\overline{\mathrm{\Psi }}_LM\mathrm{\Psi }_R+\overline{\mathrm{\Psi }}_RM^{}\mathrm{\Psi }_L\right)g_v\left(\overline{\mathrm{\Psi }}_L\gamma ^\mu l_\mu \mathrm{\Psi }_L+\overline{\mathrm{\Psi }}_R\gamma ^\mu r_\mu \mathrm{\Psi }_R\right)`$ (13)
$`={\displaystyle \frac{g_s}{\sqrt{2}}}\overline{\mathrm{\Psi }}\left({\displaystyle \underset{a=0}{\overset{8}{}}}s_a\lambda _a+i\gamma ^5{\displaystyle \underset{a=0}{\overset{8}{}}}p_a\lambda _a\right)\mathrm{\Psi }{\displaystyle \frac{g_v}{2\sqrt{2}}}\overline{\mathrm{\Psi }}\left(\gamma ^\mu {\displaystyle \underset{a=0}{\overset{8}{}}}v_\mu ^a\lambda _a\gamma ^\mu \gamma ^5{\displaystyle \underset{a=0}{\overset{8}{}}}a_\mu ^a\lambda _a\right)\mathrm{\Psi }.`$ (14)
From the quark-meson interaction, the coupling constants between scalar mesons, vector mesons and quarks have the following relations:
$`{\displaystyle \frac{g_s}{\sqrt{2}}}`$ $`=`$ $`g_{a_0}^u=g_{a_0}^d=g_\sigma ^u=g_\sigma ^d=\mathrm{}={\displaystyle \frac{1}{\sqrt{2}}}g_\zeta ^s,g_{a_0}^s=g_\sigma ^s=g_\zeta ^u=g_\zeta ^d=0,`$ (15)
$`{\displaystyle \frac{g_v}{2\sqrt{2}}}`$ $`=`$ $`g_\rho ^u=g_\rho ^d=g_\omega ^u=g_\omega ^d=\mathrm{}={\displaystyle \frac{1}{\sqrt{2}}}g_\varphi ^s,g_\omega ^s=g_\rho ^s=g_\varphi ^u=g_\varphi ^d=0.`$ (16)
In the mean field approximation, the chiral-invariant scalar meson $`_{\mathrm{\Sigma }\mathrm{\Sigma }}`$ and vector meson $`_{VV}`$ self-interaction terms are written as
$`_{\mathrm{\Sigma }\mathrm{\Sigma }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}k_0\chi ^2\left(\sigma ^2+\zeta ^2\right)+k_1\left(\sigma ^2+\zeta ^2\right)^2+k_2\left({\displaystyle \frac{\sigma ^4}{2}}+\zeta ^4\right)+k_3\chi \sigma ^2\zeta `$ (18)
$`k_4\chi ^4{\displaystyle \frac{1}{4}}\chi ^4\mathrm{ln}{\displaystyle \frac{\chi ^4}{\chi _0^4}}+{\displaystyle \frac{\delta }{3}}\chi ^4\mathrm{ln}{\displaystyle \frac{\sigma ^2\zeta }{\sigma _0^2\zeta _0}},`$
$`_{VV}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\chi ^2}{\chi _0^2}}\left(m_\omega ^2\omega ^2+m_\rho ^2\rho ^2+m_\varphi ^2\varphi ^2\right)+g_4\left(\omega ^4+6\omega ^2\rho ^2+\rho ^4+2\varphi ^4\right),`$ (19)
where $`\delta =6/33`$; $`\sigma _0`$ and $`\zeta _0`$ are the vacuum expectation values of the corresponding mean fields $`\sigma `$, $`\zeta `$ which are expressed as
$$\sigma _0=F_\pi ,\zeta _0=\frac{1}{\sqrt{2}}(F_\pi 2F_K).$$
(20)
The vacuum value $`\chi _0`$ is about 280 MeV in our numerical calculation. The Lagrangian $`_{\chi SB}`$ generates the nonvanishing masses of pseudoscalar mesons
$$_{\chi SB}=\frac{\chi ^2}{\chi _0^2}\left[m_\pi ^2F_\pi \sigma +\left(\sqrt{2}m_K^2F_K\frac{m_\pi ^2}{\sqrt{2}}F_\pi \right)\zeta \right],$$
(21)
leading to a nonvanishing divergence of the axial currents which in turn satisfy the partial conserved axial-vector current (PCAC) relations for $`\pi `$ and $`K`$ mesons. Pseudoscalar, scalar mesons and also the dilaton field $`\chi `$ obtain mass terms by spontaneous breaking of chiral symmetry in the Lagrangian of Eq. (18). The masses of $`u`$, $`d`$ and $`s`$ quarks are generated by the vacuum expectation values of the two scalar mesons $`\sigma `$ and $`\zeta `$. To obtain the correct constituent mass of the strange quark, an additional mass term has to be added:
$`_{\mathrm{\Delta }m_s}=\mathrm{\Delta }m_s\overline{q}Sq`$ (22)
where $`S=\frac{1}{3}\left(I\lambda _8\sqrt{3}\right)=\mathrm{diag}(0,0,1)`$ is the strangeness quark matrix. Based on these mechanisms, the quark constituent masses are finally given by
$`m_u=m_d={\displaystyle \frac{g_s}{\sqrt{2}}}\sigma _0\text{and}m_s=g_s\zeta _0+\mathrm{\Delta }m_s.`$ (23)
The parameters $`g_s=4.76`$ and $`\mathrm{\Delta }m_s=29`$ MeV are chosen to yield the constituent quark masses $`m_q=313`$ MeV and $`m_s=490`$ MeV. In order to obtain reasonable hyperon potentials in hadronic matter, we include an additional coupling between strange quarks and the scalar mesons $`\sigma `$ and $`\zeta `$ . This term is expressed as
$`_h=[h_1(\sigma \sigma _0)+h_2(\zeta \zeta _0)]\overline{s}s.`$ (24)
Therefore, the strange quark scalar-coupling constants are modified and do not exactly satisfy Eq. (15). The hyperon potentials were listed in our previous paper . In the quark mean field model, quarks are confined in baryons by the Lagrangian $`_c=\overline{\mathrm{\Psi }}\chi _c\mathrm{\Psi }`$ (with $`\chi _c`$ given in Eq. (25), below). We note that this confining term is not chiral invariant. Possible extensions of the model which would restore chiral symmetry in this term have been discussed in Ref. .
The Dirac equation for a quark field $`\mathrm{\Psi }_{ij}`$ under the additional influence of the meson mean fields is given by
$$\left[i\stackrel{}{\alpha }\stackrel{}{}+\beta \chi _c(r)+\beta m_i^{}\right]\mathrm{\Psi }_{ij}=e_i^{}\mathrm{\Psi }_{ij},$$
(25)
where $`\stackrel{}{\alpha }=\gamma ^0\stackrel{}{\gamma }`$ , $`\beta =\gamma ^0`$ , the subscripts $`i`$ and $`j`$ denote the quark $`i`$ ($`i=u,d,s`$) in a baryon of type $`j`$ ($`j=N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }`$) ; $`\chi _c(r)`$ is a confinement potential, i.e. a static potential providing the confinement of quarks by meson mean-field configurations. In the numerical calculations, we choose $`\chi _c(r)=\frac{1}{4}k_cr^2`$, where $`k_c=1`$ (GeV fm$`{}_{}{}^{2})`$, which yields baryon radii (in the absence of the pion cloud ) around 0.6 fm. The quark mass $`m_i^{}`$ and energy $`e_i^{}`$ are defined as
$$m_i^{}=g_\sigma ^i\sigma g_\zeta ^i\zeta +m_{i0}$$
(26)
and
$$e_i^{}=e_ig_\omega ^i\omega g_\rho ^i\rho g_\varphi ^i\varphi ,$$
(27)
where $`e_i`$ is the energy of the quark under the influence of the meson mean fields. Here $`m_{i0}=0`$ for $`i=u,d`$ (nonstrange quark) and $`m_{i0}=\mathrm{\Delta }m_s`$ for $`i=s`$ (strange quark). The effective baryon mass can be written as
$`M_j^{}={\displaystyle \underset{i}{}}n_{ij}e_i^{}E_j^0,`$ (28)
where $`n_{ij}`$ is the number of quarks with flavor $`\mathrm{`}\mathrm{`}i\mathrm{"}`$ in a baryon with flavor $`j`$, with $`j=N\{p,n\},\mathrm{\Sigma }\{\mathrm{\Sigma }^\pm ,\mathrm{\Sigma }^0\},\mathrm{\Xi }\{\mathrm{\Xi }^0,\mathrm{\Xi }^{}\},\mathrm{\Lambda }`$ and $`E_j^0`$ was found to be only very weakly dependent on the external field strength. We therefore use Eq. (28), with $`E_j^0`$ a constant, independent of the density, which is adjusted to give a best fit to the free baryon masses. Compared with the earlier square root ansätz, here we use the linear definition of effective baryon mass. As we have explained in Ref. the linear definition of effective baryon mass has been derived using a symmetric relativistic approach , while to the best of our knowledge, no equivalent derivation exists for the square root case.
## III hadronic system
Based on the previously defined quark mean field model the thermodynamical potential for the study of hadronic systems is written as
$$\mathrm{\Omega }=\underset{j=N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }}{}\frac{2k_BT}{(2\pi )^3}_0^{\mathrm{}}d^3k\left\{\mathrm{ln}\left(1+e^{(E_j^{}(k)\nu _j)/k_BT}\right)+\mathrm{ln}\left(1+e^{(E_j^{}(k)+\nu _j)/k_BT}\right)\right\}_M,$$
(29)
where $`E_j^{}(k)=\sqrt{M_j^2+k^2}`$ and $`M_j^{}`$ is the effective baryon mass. The quantity $`\nu _j`$ is related to the usual chemical potential $`\mu _j`$ by $`\nu _j=\mu _jg_\omega ^j\omega g_\rho ^j\rho g_\varphi ^j\varphi `$. The mesonic Lagrangian
$$_M=_{\mathrm{\Sigma }\mathrm{\Sigma }}+_{VV}+_{\chi SB}$$
(30)
describes the interaction between mesons which includes the scalar meson self-interaction $`_{\mathrm{\Sigma }\mathrm{\Sigma }}`$, the vector meson self-interaction $`_{VV}`$ and the explicit chiral symmetry breaking term $`_{\chi SB}`$ defined previously in Eqs. (18), (19) and (21). The Lagrangian $`_M`$ involves scalar ($`\sigma `$, $`\zeta `$ and $`\chi `$) and vector ($`\omega `$, $`\rho `$ and $`\varphi `$) mesons. The interactions between quarks and scalar mesons result in the effective baryon masses $`M_j^{}`$. The interactions between quarks and vector mesons generate the baryon-vector meson interaction terms. The energy per volume and the pressure of the system can be derived as $`\epsilon =\mathrm{\Omega }\frac{1}{T}\frac{\mathrm{\Omega }}{T}+\nu _j\rho _j`$ and $`p=\mathrm{\Omega }`$, where $`\rho _j`$ is the density of baryon $`j`$. At zero temperature, $`\mathrm{\Omega }`$ can be expressed as
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \underset{j=N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }}{}}{\displaystyle \frac{1}{24\pi ^2}}\{\nu _j[\nu _j^2M_j^2]^{1/2}[2\nu _j^25M_j^2]`$ (31)
$`+`$ $`3M_j^4\mathrm{ln}\left[{\displaystyle \frac{\nu _j+\left(\nu _j^2M_j^2\right)^{1/2}}{M_j^{}}}\right]\}_M,`$ (32)
The equations for mesons $`\varphi _i`$ can be obtained by the formula $`\frac{\mathrm{\Omega }}{\varphi _i}=0`$. Therefore, the equations for $`\sigma `$, $`\zeta `$ and $`\chi `$ are
$`k_0\chi ^2\sigma 4k_1\left(\sigma ^2+\zeta ^2\right)\sigma 2k_2\sigma ^32k_3\chi \sigma \zeta {\displaystyle \frac{2\delta }{3\sigma }}\chi ^4+{\displaystyle \frac{\chi ^2}{\chi _0^2}}m_\pi ^2F_\pi `$ (33)
$`\left({\displaystyle \frac{\chi }{\chi _0}}\right)^2m_\omega \omega ^2{\displaystyle \frac{m_\omega }{\sigma }}\left({\displaystyle \frac{\chi }{\chi _0}}\right)^2m_\rho \rho ^2{\displaystyle \frac{m_\rho }{\sigma }}+{\displaystyle \underset{j=N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }}{}}{\displaystyle \frac{M_j^{}}{\sigma }}<\overline{\psi _j}\psi _j>`$ $`=`$ $`0,`$ (34)
$`k_0\chi ^2\zeta 4k_1\left(\sigma ^2+\zeta ^2\right)\zeta 4k_2\zeta ^3k_3\chi \sigma ^2{\displaystyle \frac{\delta }{3\zeta }}\chi ^4+{\displaystyle \frac{\chi ^2}{\chi _0^2}}\left(\sqrt{2}m_k^2F_k{\displaystyle \frac{1}{\sqrt{2}}}m_\pi ^2F_\pi \right)`$ (35)
$`\left({\displaystyle \frac{\chi }{\chi _0}}\right)^2m_\varphi \varphi ^2{\displaystyle \frac{m_\varphi }{\zeta }}+{\displaystyle \underset{j=\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }}{}}{\displaystyle \frac{M_j^{}}{\zeta }}<\overline{\psi _j}\psi _j>=0,`$ (36)
$`k_0\chi \left(\sigma ^2+\zeta ^2\right)k_3\sigma ^2\zeta +\left(4k_4+1+4ln{\displaystyle \frac{\chi }{\chi _0}}{\displaystyle \frac{4\delta }{3}}ln{\displaystyle \frac{\sigma ^2\zeta }{\sigma _0^2\zeta _0}}\right)\chi ^3`$ (37)
$`+{\displaystyle \frac{2\chi }{\chi _0^2}}\left[m_\pi ^2F_\pi \sigma +\left(\sqrt{2}m_k^2F_k{\displaystyle \frac{1}{\sqrt{2}}}m_\pi ^2F_\pi \right)\zeta \right]{\displaystyle \frac{\chi }{\chi _0^2}}(m_\omega ^2\omega ^2+m_\rho ^2\rho ^2+m_\varphi ^2\varphi ^2)`$ $`=`$ $`0,`$ (38)
where $`<\overline{\psi }_j\psi _j>`$ is expressed as
$`<\overline{\psi }_j\psi _j>`$ $`=`$ $`{\displaystyle \frac{M_j^{}}{\pi ^2}}{\displaystyle _0^{k_{F_j}}}𝑑k{\displaystyle \frac{k^2}{\sqrt{M_j^2+k^2}}}`$ (39)
$`=`$ $`{\displaystyle \frac{M_j^3}{2\pi ^2}}\left[{\displaystyle \frac{k_{F_j}}{M_j^{}}}\sqrt{1+{\displaystyle \frac{k_{F_j}^2}{M_j^2}}}\mathrm{ln}\left({\displaystyle \frac{k_{F_j}}{M_j^{}}}+\sqrt{1+{\displaystyle \frac{k_{F_j}^2}{M_j^2}}}\right)\right],`$ (40)
with $`k_{F_j}=\sqrt{\nu _j^2M_j^2}`$.
For the $`\beta `$-equilibrium, the chemical potentials for the baryons satisfy the following equations:
$$\mu _\mathrm{\Lambda }=\mu _{\mathrm{\Sigma }^0}=\mu _{\mathrm{\Xi }^0}=\mu _n=\mu _p+\mu _e=\mu _p+\mu _\mu ,$$
(41)
$$\mu _{\mathrm{\Sigma }^+}=\mu _p,$$
(42)
$$\mu _\mathrm{\Sigma }^{}=\mu _\mathrm{\Xi }^{}=\mu _n+\mu _e.$$
(43)
There are only two independent chemical potentials which are determined by the total baryon density and neutral charge:
$$\rho _B=\rho _p+\rho _n+\rho _\mathrm{\Lambda }+\rho _{\mathrm{\Sigma }^+}+\rho _{\mathrm{\Sigma }^0}+\rho _\mathrm{\Sigma }^{}+\rho _{\mathrm{\Xi }^0}+\rho _\mathrm{\Xi }^{},$$
(44)
$$\rho _p+\rho _{\mathrm{\Sigma }^+}\rho _\mathrm{\Sigma }^{}\rho _\mathrm{\Xi }^{}\rho _e\rho _\mu =0.$$
(45)
In order to get the mass-radius relation, one has to resolve the Tolman-Oppenheimer-Volkoff (TOV) equation:
$$\frac{dp}{dr}=\frac{\left[p(r)+\epsilon (r)\right]\left[M(r)+4\pi r^3p(r)\right]}{r(r2M(r))},$$
(46)
where
$$M(r)=4\pi _0^r\epsilon (r)r^2𝑑r.$$
(47)
With the equations of state, the functions, such as $`M(r)`$, $`\rho (r)`$ and $`p(r)`$, etc. can be obtained.
## IV Numerical results
The parameters of this model were determined by the meson masses in vacuum and the saturation properties of nuclear matter which were listed in the table 1 of Ref. . The improved linear definition of effective baryon mass is chosen in our numerical calculations. We first discuss the equations of state of neutron matter and strange hadronic matter which are needed for the calculation of neutron stars. For pure neutron stars, there are only neutrons present. For strange hadronic stars, with increasing baryon density, other kinds of baryons will appear. In Fig. 1 we show the fractions of octet baryons versus density with $`\beta `$-equilibrium. With the increasing of baryon density, the neutron fraction decreases slowly from 1. If the density is lower than about 0.19 fm<sup>-3</sup>, the fraction of electrons is the same as that of protons which makes the system charge neutral. The muon appears when the density is in the range 0.19 - 0.98 fm<sup>-3</sup>. The maximum fractions of muons and electrons appear at $`\rho _B`$ 0.4 fm<sup>-3</sup>. Their fractions decrease with the increasing fractions of hyperons. When the density is larger than about 0.4 fm<sup>-3</sup>, the $`\mathrm{\Sigma }^{}`$ hyperons appear and the fraction of neutrons decreases faster. After the density is larger than about 0.57 fm<sup>-3</sup>, $`\mathrm{\Lambda }`$ hyperons start to appear. The fraction of $`\mathrm{\Sigma }^{}`$ hyperons decreases with the increasing density after $`\mathrm{\Xi }^{}`$ hyperons appear where the density is about 0.84 fm<sup>-3</sup>.
The density dependence of the effective baryon masses and scalar mean fields are shown in Fig. 2. The $`\sigma `$ field decreases quickly with the increasing baryon density when the density is small, $`\rho _B<0.4`$ fm<sup>-3</sup>. This is because at small baryon density, the nucleon is dominant and there are no hyperons. With the increasing of density, more and more hyperons appear. As a result, the $`\zeta `$ field decreases quickly. At a broad range of densities, the value of $`\chi `$ changes little.
In Fig. 3, the pressure versus baryon density is shown. The dashed and solid lines are for the pure neutron star and the strange hadronic star with $`\beta `$-equilibrium, respectively. When the density is low, the two curves are close to each other. With the increasing of baryon density, the contributions of protons and hyperons are not negligible. The inclusion of hyperons will soften the equation of state of hadronic matter. As a result, at a given baryon density the pressure of strange hadronic matter is smaller than the corresponding pressure of pure neutron matter. The pressure $`p`$ versus energy density $`\epsilon `$ is shown in Fig. 4. Again, one can see that the equation of state of strange hadronic matter is softer than that of pure neutron matter.
We now study neutron stars with the obtained EOS. By solving the TOV equation, the baryon density versus radius can be obtained which is shown in Fig. 5. The central densities $`\rho _c`$ are chosen to be $`3\rho _0`$ and $`5\rho _0`$ where $`\rho _0`$ (0.16 fm<sup>-3</sup>) is the saturation density of symmetric nuclear matter. The dashed and solid lines are for pure neutron stars and strange hadronic stars with $`\beta `$-equilibrium, respectively. With the increasing radius, the density of strange hadronic stars decreases a little faster than that of pure neutron stars which results in a smaller radius. The radii of stars are not sensitive to their the central density. For example, for $`\rho _c`$ of $`3\rho _0`$ and $`5\rho _0`$, the radii are both around 11-12 km.
We plot the star mass ratio $`M/M_{\mathrm{sun}}`$ versus central baryon density in Fig. 6. The maximum mass of pure neutron stars is about 1.8$`M_{\mathrm{sun}}`$ with a central density 1.05 fm<sup>-3</sup>. After the central density is larger than 1.05 fm<sup>-3</sup>, the star will become unstable. The maximum mass changes to $`1.45M_{\mathrm{sun}}`$ when hyperons are included. In the range $`3\rho _0<\rho _c<6\rho _0`$, the masses of pure neutron stars and strange hadronic stars are $`1.48M_{\mathrm{sun}}<M<1.8M_{\mathrm{sun}}`$ and $`1.23M_{\mathrm{sun}}<M<1.45M_{\mathrm{sun}}`$, respectively. Our results are reasonable compared with the observation of the six known stars with masses in the range $`1.36\pm 0.08M_{\mathrm{sun}}`$, since the “neutron star” is in fact a strange hadronic star with $`\beta `$-equilibrium. We should also keep in mind that there are some heavy stars reported in recent years. For PSR J0437-4715, the mass is found to be $`1.58\pm 0.18M_{\mathrm{sun}}`$ . For Vale X-1, Cygnus X-2 and 4U 1820-30, their masses are determined to be $`1.87_{0.17}^{+0.23}M_{\mathrm{sun}}`$ , $`1.8\pm 0.4M_{\mathrm{sun}}`$ and $`2.3M_{\mathrm{sun}}`$ . The rotation of a star can increase its mass by $`10\%`$ . Therefore, the calculated maximum mass of strange hadronic stars can be as large as $`1.6M_{\mathrm{sun}}`$. If the heavy stars such as 4U 1820-30 are confirmed, the strange hadronic star would be ruled out if this model is a good description of Nature. It is possible to increase the maximum star mass by making the EOS stiffer at higher densities. Whether the inclusion of a quark core in the strange star will result in a large-maximum mass is an interesting topic.
In Fig. 7, the masses of stars versus their radii are shown. For pure neutron stars, when their masses are in the range $`0.5M_{\mathrm{sun}}<M<1.8M_{\mathrm{sun}}`$, their radii are about 11.0-12.3 km. For the strange hadronic stars, when the masses are in the range $`0.5M_{\mathrm{sun}}<M<1.45M_{\mathrm{sun}}`$, the radii are about 10.7-11.7 km. With the same mass ($`M>0.2M_{\mathrm{sun}}`$), strange hadronic stars have smaller radii compared with pure neutron stars. Because the size of neutron stars is small, it is very difficult to observe and measure their radii directly. Different indirect methods lead to different values of radii with large errors. For example, for the RX J1856-3754, the radius varies from 5 km to 15 km with a mass of $`1.4M_{\mathrm{sun}}`$ . More accurate values are needed to obtain a more strict constraint on the EOS of hadronic matter.
## V summary
We have investigated pure neutron stars and strange hadronic stars in the chiral $`SU(3)`$ quark mean field model. The $`\mathrm{\Lambda }`$, $`\mathrm{\Sigma }`$ and $`\mathrm{\Xi }`$ hyperons are included in the model. The proton and hyperon contributions to the system are important at high baryon density when $`\beta `$-equilibrium is achieved, and soften the EOS of hadronic matter. The maximum pure neutron star mass is about $`M=1.8M_{\mathrm{sun}}`$ with a corresponding radius $`R=11.0`$ km and central density $`\rho _c=1.05`$ fm<sup>-3</sup>. For the strange hadronic stars, the maximum masses are about $`1.45M_{\mathrm{sun}}`$ and the corresponding radii and central density are $`R=10.9`$ km and $`\rho _c=1.0`$ fm<sup>-3</sup>. When the central densities are between $`3\rho _0`$ and $`6\rho _0`$, the masses of stars are in the range $`1.23M_{\mathrm{sun}}<M<1.45M_{\mathrm{sun}}`$ (strange hadronic stars) and $`1.48M_{\mathrm{sun}}<M<1.8M_{\mathrm{sun}}`$ (pure neutron stars). If the masses of stars are larger than $`0.5M_{\mathrm{sun}}`$, the typical values of radii are 10.7-11.7 km (strange hadronic stars) km and 11.0-12.3 km (pure neutron stars).
Our results are reasonable compared with astrophysical observations where the six known neutron stars have masses in the narrow range $`1.36\pm 0.08M_{\mathrm{sun}}`$. Accurate values of radii for neutron stars are needed to get a more strict constraint on the EOS of hadronic matter. As for the heavy stars, for example 4U 1820-30, if its mass $`M2.3M_{\mathrm{sun}}`$ is confirmed, then strange hadronic stars are obviously ruled out if the model explored herein is a good description of Nature. It is therefore of interest to see whether including quark degrees of freedom can lead to this large mass.
## Acknowledgements
P.W. thanks the Theory Group at Jefferson Lab for their kind hospitality. This work was supported by the Australian Research Council and by DOE contract DE-AC05-84ER40150, under which SURA operates Jefferson Laboratory.
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# Mass Flows in Cometary UCHII Regions
## 1 Introduction
Ultracompact HII (UCHII) regions form when massive OB stars ignite inside dense molecular clouds and high energy UV photons ionize the surrounding neutral material. They have small sizes ($`0.1`$ pc), high electron densities ($`10^4`$ cm<sup>-3</sup>), and high emission measures ($`EM=n_in_edl10^7`$pc cm<sup>-6</sup>) (Wood & Churchwell, 1989b). The ionization increases the temperature ($`10^4`$K) and number density ($`H_22p^++2e^{}`$) causing the gas pressure inside the regions to be at least two orders of magnitude higher than in the surrounding molecular gas. UCHII regions should then expand at approximately the speed of sound in the ionized gas ($`10`$ km s<sup>-1</sup>) until they reach pressure equilibrium with the surrounding material or break out of their parent molecular clouds. Because of this expansion, the density and emission measure both should drop with time. Once the emission measure drops below $`10^7`$ pc cm<sup>-6</sup>, these HII regions will no longer be considered ultracompact.
A longstanding puzzle regarding UCHII regions involves their numbers and their lifetimes. If HII regions expand at 10 km s<sup>-1</sup>, they should remain ultracompact for only $`10^4`$ yr. Since this is $`<`$1% of the lifetime of an OB star, the number of UCHII regions should be $`<`$1% the number of OB stars. Radio interferometry observations confirm that the ratio of IRAS far-infrared (FIR) flux densities provides an efficient way to select embedded OB stars from the IRAS data (Wood & Churchwell, 1989b; Churchwell, 1990; Garay et al., 1993; Kurtz et al., 1994; Miralles et al., 1994). However, counts of IRAS FIR sources with the spectral characteristics of UCHII regions and the optically visible O stars in the solar neighborhood show that OB stars spend $`1020\%`$ of their main-sequence lifetime embedded in molecular clouds (Wood & Churchwell, 1989a). Apparently, the lifetimes of UCHII regions are an order of magnitude larger than predicted by the classical pressure-driven spherical expansion model.
The variety of observed UCHII morphologies, including spherical, shell-like, cometary, and irregular, also needs to be explained. An understanding of the kinematics of the ionized gas will help provide an explanation of both the morphologies and the lifetime of UCHII regions. It has been suggested that long lifetimes could result from some kind of containment or gas replenishment mechanism (Hollenbach et al., 1994; Dyson et al., 1995; Redman et al., 1996; Williams et al., 1996; Redman et al., 1998). Various mechanisms would have different effects on the morphologies and the kinematics of UCHII region, suggesting that observations could distinguish between these mechanisms.
The well-organized overall appearance of cometary UCHII regions make them good objects in which to study UCHII region kinematics and morphology. One way to explain these objects is with a bow shock model, which suggests the cometary structure of some UCHII regions is the result of the supersonic motion of OB stars with high speed winds through molecular clouds (Hughes & Viner, 1976; Wood & Churchwell, 1989b; van Buren et al., 1990; Mac Low et al., 1991; van Buren & Mac Low, 1992). Swept-up ambient gas and stellar wind material accumulate where the ram pressures associated with these two mass flows balance, resulting in a shell-like structure and a surface flow. This model may not explain the lifetime of all UCHII regions, but could provide a part of the solution to this problem.
Because ultracompact HII regions are formed inside dense molecular clouds, the extinction toward these objects usually is very high, with a typical line of sight extinction of $`A_v=3050`$ (Hanson et al., 2002). Thus, optical techniques used in diffuse HII region observations are inappropriate for studies of UCHII regions. Low extinction at infrared and radio wavelengths makes these wavelength bands more suitable for studies of UCHII regions. Hydrogen radio recombination line observations toward UCHII regions have been carried out extensively (Garay et al., 1985; Kim & Koo, 2001; Araya et al., 2002). These observations reveal that thermal motions cannot account for the linewidths of many UCHII regions. However, significant thermal broadening makes it hard to study the velocity structure in these regions using hydrogen recombination lines because of hydrogen’s low atomic mass (Jaffe & Martín-Pintado, 1999; Sewilo et al., 2004; De Pree et al., 2004, and references therein). The thermal line width of hydrogen recombination lines is $`\mathrm{\Delta }\upsilon _{FWHM}21.4`$km s<sup>-1</sup> for gas with a temperature of $`10^4`$K. This large thermal linewidth means that heavier ions are better probes of bulk motion; the thermal broadening is only about $`4.8`$km s<sup>-1</sup> for $`Ne^+`$ ions at the same temperature. Mid-infrared ionic fine-structure lines have been used to probe the structure and excitation of HII regions (Beck et al., 1981; Lacy et al., 1982; Takahashi et al., 2000), but generally with too low spectral resolution to study gas motions. To study kinematics and distinguish organized motions from turbulence and thermal broadening of heavy ions, one needs $`5`$km s<sup>-1</sup> velocity resolution.
We have begun a program of high spectral resolution observations of mid-infrared fine structure line emission from UCHII regions. By mapping a small sample of UCHII regions of different morphologies at high spatial resolution in \[Ne II\]$`\lambda 12.8\mu `$m, \[Ar III\]$`\lambda 9.0\mu `$m, \[S III\]$`\lambda 18.7\mu `$m and \[S IV\]$`\lambda 10.5\mu `$m, we can study the kinematics and physical conditions in UCHII regions in order to better understand massive star formation and the relationship between these stars and their surrounding environments. Among these lines, the \[Ne II\] line is particularly bright and is present in HII regions with a broad range of excitation, so it is well suited for kinematic analysis. We presented high spectral resolution observations of one cometary UCHII region, Mon R2 IRS1 (Mon R2 hereafter), in Jaffe et al. (2003). In the current paper, we present \[Ne II\] observations of another cometary UCHII region, G29.96 -0.02 (G29.96 hereafter). We also examine bow shock models in some detail and compare them to the observations of both G29.96 and Mon R2.
In section 2, we describe our mapping and data reduction methods. We present our \[Ne II\] line emission observations toward G29.96 in section 3. In section 4, we describe the kinematics of the bow shock model and our relaxation method for simulating the formation of the surface flow. We also introduce an additional acceleration, due to the pressure gradient in the ionized gas along the flow. In section bf 5, we compare the model predictions with the \[Ne II\] observations of G29.96 and Mon R2. Finally, we discuss the existing problems of models for cometary UCHII regions in section 6.
## 2 Observations
### 2.1 Instrument
The \[Ne II\] observations of G29.96 were carried out with TEXES (the Texas Echelon Cross Echelle Spectrograph, Lacy et al., 2002) on the NASA 3 meter Infrared Telescope Facility (IRTF) on Mauna Kea, in June, 2001. TEXES is a high resolution (R $``$ 100,000) spectrograph operating at mid-infrared (5-25 $`\mu `$m) wavelengths. The slit, with a $`1.4^{\prime \prime }`$ width and a $`11.5^{\prime \prime }`$ length, was oriented north-south on the sky. With this slit, we achieved a spectral resolution of $`4`$km s<sup>-1</sup> or R = 75,000 at $`12.8\mu `$m. Each pixel along the slit was about $`0.36^{\prime \prime }`$ on the sky, which is a little smaller than half of the diffraction limit ($`0.88^{\prime \prime }`$) of the IRTF at $`12.8\mu `$m. To map the regions of interest we stepped the telescope west to east across the objects without chopping. Multiple, partially overlapping scans were needed to cover the entire region. For G29.96, each scan had a length of $`20^{\prime \prime }`$ with a step size of $`0.4^{\prime \prime }`$. For Mon R2, each scan was $`45^{\prime \prime }`$ long with $`0.7^{\prime \prime }`$ steps (Jaffe et al., 2003). Nine overlapping scans were made on G29.96, and 20 scans were made on Mon R2. Because the scans overlapped, the integration time at each point in the map is approximately 18 seconds for G29.96 and 8 seconds for Mon R2.
### 2.2 Data Reduction
The spectral-spatial datacubes produced from the scans were first reduced with a custom Fortran reduction program that performs general purpose corrections, such as the correction of optical distortions, flat fielding, bad pixel masking and cosmic ray spike removal (Lacy et al., 2002). Radiometric calibration was obtained from measurements of an ambient temperature blackbody before each set of scans, giving intensities in units of erg cm<sup>-2</sup> s<sup>-1</sup> sr<sup>-1</sup> (cm$`{}_{}{}^{1})^1`$. The uncertainties in the intensity are mostly systematic, and are probably $`\pm 20\%`$. The wavelength calibration was obtained from sky emission lines, and is accurate to $`1`$km s<sup>-1</sup>. IDL scripts were used to do the sky background subtraction, multiple scan cross-correlation and combining, and datacube manipulation. Positions safely off the object at both ends of each scan allow us to interpolate the sky emission at each step position. In doing the sky background subtraction, we assumed that the sky background varies linearly with time during the course of a scan, based on the brightness of pixels off the object. After sky emission subtraction, multiple scans of the same region were cross-correlated, shifted, and added to make a complete datacube for the object.
## 3 Results
We show the integrated \[Ne II\] line emission map for G29.96 in Fig. 1. The overall cometary shape is apparent and is symmetric about an axis at a position angle $`70^{}`$ east of north. The brightest emission forms an arc perpendicular to the symmetry axis. To the west of the arc, the line emission drops rapidly from the peak value to the background level within $`4^{\prime \prime }`$. On the east side of the bright arc, fainter emission extends over $`13^{\prime \prime }`$ to the east-northeast. The change of the emission level is much more gradual and less uniform. The east edge of the region is much fuzzier and broken by a faint “lane” into two parts. The southern portion extends the curve of the emission arc out $`6^{\prime \prime }`$ and ends at a clump of ionized gas of size $`3^{\prime \prime }\times 3^{\prime \prime }`$. From our velocity channel maps (Fig. 2), we can see that it is the continuous extension of the arc.
Velocity channel maps of G29.96 (Fig. 2) show additional details. A cavity is present at medium and lower velocity channels, from $`82`$km s<sup>-1</sup> to $`94`$km s<sup>-1</sup> . The fainter structures on the east side of the region and the brighter arc form an almost complete ring around the cavity, although the symmetry axis of the elliptical ring is at a somewhat larger position angle ($`120^{}`$) than that of the bright arc ($`70^{}`$). A similar ring is seen in channel maps of Mon R2 (Jaffe et al., 2003), except that the major axis of the ring in Mon R2 is nearly perpendicular to the symmetry axis of its bright arc. A compact emission peak centered at $`V_{LSR}108`$ km s<sup>-1</sup> can be seen $`5^{\prime \prime }`$ northeast of the peak of the arc in G29.96. The \[Ne II\] line is generally broad, reaching $`40`$km s<sup>-1</sup> near the bright arc. The emission is typically bright near the velocity of the ambient cloud ($`V_{LSR}98`$km s<sup>-1</sup>, obtained from single dish observations of rotational transitions of CS molecules, Martín-Hernández et al., 2003, and references therein), but the line is significantly blueshifted ahead of the emission peak. The line center reaches a $`V_{LSR}82`$km s<sup>-1</sup> at the west edge of the region
In Fig. 3, we plot \[Ne II\] line contours on the top of the VLA 2 cm continuum image of G29.96 (Fey et al., 1995). We first aligned the two maps, assuming they are emitted by the same ionized gas. We then convolved the radio map with our beam. The close resemblance of the maps after convolution suggests that the much sharper edges seen in the VLA image are the result of its higher spatial resolution ($`<0.56^{\prime \prime }`$). Comparing a cross-cut along the symmetry axis (Fig. 4) confirms that \[Ne II\] and radio free-free flux are both proportional to the emission measure of the region. Radio recombination line observations at $`0.62^{\prime \prime }`$ and 4 km s<sup>-1</sup> resolution showed similar arc-like structure (Wood & Churchwell, 1991). They also showed that the H76$`\alpha `$ line is broader in front of the radio continuum emission arc, where the velocity gradient is also the highest.
Martín-Hernández et al. (2003) obtained a narrow-band H<sub>2</sub> 1-0 S(1) filter map and K band ($`2.072.19\mu `$m) spectra of the region with a resolution of R=8,000 along a $`120^{\prime \prime }`$ long slit along the symmetry axis. Their map shows many local near-infrared emission clumps together with the proposed ionizing star of the region lying $`2^{\prime \prime }.3`$ northeast ($`p.a.=64^{}`$) of the emission peak. We mark the star position with an asterisk in our map according to this offset (Fig. 1). Their $`Br\gamma `$ linewidth ranges between $`42`$km s<sup>-1</sup> and $`62`$km s<sup>-1</sup> and shows significant blueshift in front of the emission arc, which is consistent with our \[Ne II\] observations. The fact that $`Br\gamma `$ line is systematically broader than \[Ne II\] along the slit supports the idea that \[Ne II\] is a better tracer of the ionized gas motions.
## 4 Bow Shock and Relaxation - A Physical Model for Parabolic Flow
In this section, we present an extension of the analytic solution for bow shocks using a numerical relaxation technique. We then use the resulting gas flow as a template with which we can compare the observed kinematics of G29.96 and Mon R2.
### 4.1 Analytical Solution
Bow shock models have been explored to explain the formation of cometary UCHII regions by several authors (van Buren et al., 1990; Mac Low et al., 1991; van Buren & Mac Low, 1992; Wilkin, 1996). They describe a situation where a star with a strong stellar wind and high UV luminosity moves inside a dense molecular cloud. In the frame of reference of the star, the stellar wind and ambient material collide and create a stationary shock region approximately paraboloidal in shape in front of the star. The standoff distance $`r_0`$ is determined by pressure balance at the apex of the shell:
$`\rho _wv_w^2+n_wkT_w=\rho _av_{}^2+n_akT_a,`$ (1)
where $`\rho _w=\frac{\dot{M}_w}{4\pi r_0^2v_w}`$ is the stellar wind density at the distance $`r_0`$ and $`\rho _a=n_a\mu `$ is the density of the ambient medium. The corresponding number densities are $`n_w`$ and $`n_a`$. Where the stellar wind and ambient medium collide near head-on, the ram pressures are much greater than the gas thermal pressure, thus the standoff distance can be expressed as:
$`r_0=\sqrt{{\displaystyle \frac{\dot{M}_wv_w}{4\pi \rho _av_{}^2}}}`$ (2)
by neglecting gas thermal pressure (Wilkin, 1996). The gas pressure becomes non-negligible far downstream from the star where the shell surface is roughly parallel to the stellar motion, resulting in reduced ram pressure from the ambient medium.
For physical and mathematical simplicity, most models assume that the stellar wind and ambient medium material are well-coupled inside the shock region, and radiative cooling is very efficient on the time scale of the ionized gas recombination time. The resulting shell is “momentum supported”, rather than “energy supported”, and relatively thin compared to its scale. Mass and momentum are conserved and transported within the shell. This momentum conserving assumption provides the possibility of investigating the kinematics analytically. Wilkin (1996) derived the formula for the shell’s shape by applying mass and momentum conservation and neglecting the gas thermal pressure on both sides of the shell, as well as the tangential acceleration caused by pressure gradients within the shell:
$`r(\theta )`$ $`=`$ $`r_0\mathrm{csc}\theta \sqrt{3(1\theta \mathrm{cot}\theta )},`$ (3)
where $`r(0)=r_0`$ is the standoff distance, and $`\theta `$ is the angle between the point on the shell and the apex. He also derived expressions for the surface density and the tangential velocity of the shell.
### 4.2 Relaxation
The analytical method has no numerical uncertainties, but necessarily leaves out some details. Once we include more terms in the momentum equation, it can no longer be integrated analytically. Assuming a bow shock is a steady state configuration, we can apply a relaxation method to the mathematical description of the system. Using the relaxation method allows us to include more physical processes than the analytical method of Wilkin (1996), while being easier to formulate and less time-consuming to calculate than a full hydrodynamic approach.
We assumed a cylindrically symmetric bow shock shaped as an approximately paraboloidal shell. A grid with fixed angular size is created on the surface of the shell and each cell has the same angular size along both the azimuthal and the polar directions, $`\delta \theta =\delta \varphi =0.5^{}`$. Due to the symmetry of such a geometry, we only need to calculate quantities in one strip of cells along the polar direction. We assume that the gas in a cell moves as a single fluid. That is, we assume rapid cooling behind the shock where the ambient and stellar wind material mix and form a uniform shell, and we neglect shear motion between front and back sides of a cell. Mass flow, gas velocity, and position calculated from the analytical formula (Wilkin, 1996) are assigned to each cell as initial conditions (which were shown not to affect the results). Then, we let the program iteratively adjust these parameters by applying mass and momentum conservation.
During relaxation, the physical and geometric properties of each cell are recalculated to replace the previous values in each step until the whole system achieves convergence. In each step, the mass flow from one cell (i) into a neighboring cell (i+1) is given by:
$`\delta m^i=\delta m^{i1}+\delta m_w^i+\delta m_a^i`$ (4)
where $`\delta m^{i1}`$ is the mass flowing out of the previous cell. The mass flows from the stellar wind and the ambient medium are $`\delta m_w^i`$ and $`\delta m_a^i`$. They are given by:
$`\delta m_w^i`$ $`=`$ $`{\displaystyle \frac{\delta \mathrm{\Omega }}{4\pi }}\dot{M}_w\tau `$ (5)
$`\delta m_a^i`$ $`=`$ $`\rho _a(𝐯_a𝐒^i)\tau `$ (6)
$`\delta \mathrm{\Omega }`$ is the solid angle subtended by the cell element, $`\dot{M}_w`$ is the mass loss rate of the stellar wind, $`𝐒^i`$ is the outward pointing normal vector with magnitude equal the surface area of the cell $`i`$, $`𝐯_a`$ is the velocity vector of the ambient material, and $`\tau `$ is the time step of the iteration. Under the same condition, we also have the formula for total momentum flow from the cell:
$$\delta 𝐩^i=\delta 𝐩^{i1}+\delta m_w^i𝐯_w^i+\delta m_a^i𝐯_a^i+(P_w^iP_a^i)𝐒^i\tau +\mathrm{\Phi }_P$$
(7)
where $`P_a^i`$ and $`P_w^i`$ are the gas pressures of the external media. The orientation of the normal is perpendicular to the direction of tangential velocity, $`𝐯^i=\delta 𝐩^i/\delta m^i`$. $`\mathrm{\Phi }_P`$ is the extra pressure term resulting from the tangential pressure gradient inside the shell. We include it in the relaxation calculation after testing our model for the case without pressure acceleration. The distance of the cell from the star is calculated from the law of sines:
$$r^{i+1}=\frac{\mathrm{sin}(\beta ^i)}{\mathrm{sin}(\beta ^i\delta \theta )}r^i,$$
(8)
where $`\delta \theta `$ is the angular separation between cells and $`\beta `$ is the angle formed by the radius and the velocity vector. Because we assumed that the system will eventually reach a steady state, we did not try to solve a non-steady problem. Our calculation simply proceeds and finds a converging solution.
The shape of the shell and the tangential velocity along the shell, calculated by the relaxation method neglecting the contribution of the gas pressure gradient, are shown (dashed line) in Fig. 5. For this illustration, we set the star’s speed equal to $`20`$km s<sup>-1</sup>. For comparison, the analytical solution (Wilkin, 1996) is also shown (dotted line). We also show an overall shape of the calculated shell in Fig. 6. In this later calculation, we include gas thermal pressure contributions from both sides of the shell by assuming a temperature of $`50`$K for the ambient medium and $`10^4`$K for the free flowing stellar wind, which gives a more accurate but only slightly different solution for the tail portion of the shell than the analytical solution. We can see that the discrepancy between our iterated result and the analytical solution (Wilkin, 1996) is smaller before we add pressure acceleration to the relaxation solution. Calculations show that, in the rest frame of the star, gas starts to move with zero tangential speed from the apex of the paraboloidal like shell and accelerates until it reaches the star’s travelling speed, $`20`$km s<sup>-1</sup>, at the end of the shell. This acceleration is due to the accumulation of mass and momentum, which are provided by swept-up ambient material and the stellar wind. In later sections, we will show that this simple picture may not be totally correct since the effect of the pressure gradient within the shell can be significant.
Fig. 5 also shows the surface density and the particle number density along the shell. The surface density increases, almost linearly, from its value at the apex of the shell. The particle number density, derived from the shape of the shell and the normal pressure components, drops rapidly from its value at the apex. It drops by a factor of four by the position where the shell passes the star. In steady state, gas in the shell should be in appoximate pressure equilibrium with the average of the normal components of the pressures on the two sides. The difference between the two normal components causes a pressure gradient across the shell which causes its centripetal accelaration. From the pressure in the ionized gas, we can calculate the density: $`n=\frac{P_{ram}+P_{gas}}{kT}`$, where we take the ionized shell’s temperature as $`10^4`$K.
### 4.3 Pressure Gradient Acceleration
Most existing solutions, both analytical (Wilkin, 1996) and numerical (Mac Low et al., 1991), neglect the effect of the gas pressure gradient along the compressed shell. The exception is Comeron & Kaper (1998), who make a numerical hydrodynamic calculation, but concentrate on the case of a runaway O star moving through the diffuse ISM. In addition to the centripetal acceleration caused by the unbalanced external pressures, the variation in the pressure along the shell results in a pressure gradient that accelerates the gas along the shell. The momentum deposited in a cell by this effect is given by:
$`\mathrm{\Phi }_P={\displaystyle \frac{P^{i1}P^{i+1}}{2}}A^i\tau =P^iV^i\tau `$ (9)
Where $`P^i`$ and $`P^i`$ are the pressure and the pressure gradient in the $`i`$th cell. $`A^i`$ and $`V^i`$ are the cross section and the volume of the cell. We derive the thickness of each cell from the emission measure, using the ionization-recombination equilibrium equation, and the density, calculated from the balance between gas pressure and ram pressure assuming stellar parameters appropriate for G29.96 (section 4.4.1). We then assume that the thickness of the given cell is constant. If the central star is unable to ionize the whole shell, the thickness of the shell is slightly greater than that of the ionized layer, with the neutral gas in a thin, dense region just outside of the HII region. The thickness increases from the apex of the shell to the tail region because the density drops along the shell. The results of the relaxation are shown with a solid line in Fig. 5. The analytical solution is shown as a dotted line for comparison.
From the plots in Fig. 5, we can see that the change to the shape of the shell caused by including the pressure gradient is very small. Thus the change in the ionized gas number density is also small, because it is calculated from the ram pressure normal component, which directly relates to the shape of the shell. In the region close to the apex, this change can be neglected. A larger change is seen in the tangential velocity plot. For the parameters assumed, which includes a partially ionized shell, the velocity near the head of the bow shock increases so that the gas speed is $`1`$km s<sup>-1</sup> higher than without the pressure gradient acceleration at $`\theta 1radian`$. Without the pressure gradient, the maximum velocity that gas can achieve is the speed of the star. With the additional acceleration provided by the pressure gradient, the gas in the shell can reach higher velocities. The tangential speed finally drops as the shell picks up more mass from the ambient medium and the pressure gradient acceleration effect decreases toward the end of the shell due to the drop of stellar wind ram pressure. With the parameters we have chosen, the maximum speed that gas can achieve under the pressure gradient acceleration is $`1.3`$km s<sup>-1</sup> higher than the speed of the star. Up to this point, we have assumed that the ionized gas and the swept-up neutral gas are coupled, so they share the momentum carried by the stellar wind and ambient medium. When including the effect of a pressure gradient, we assume that the force resulting from this effect is also shared by ionized and neutral gases. In fact, most of this force should be exerted on the ionized gas only, since a pressure gradient causes a force per unit volume and the ionized gas fills most of the volume. If the ionized gas can slip past the neutral layer, it should reach a higher velocity than what we present in this work. The resulting speed should be similar to that in a fully ionized shell situation.
If the optical depth of the shell is small for ionizing radiation, the ionizing photons will ionize the whole shell or penetrate the shell and ionize the ambient medium beyond it. In this “ionizing photon leaking” situation, the pressure gradient acceleration effect will reach its maximum because the force is proportional to the thickness of the ionized shell. Our calculations show that the gas speed in the shell can reach 1.25 times the star speed in the range 10 - 20 km$`s^1`$, and there is no drop in speed for angles larger than 2.4 radians seen in Fig. 5(b).
In this work, we only show the results of the single layer bow shock model. Including multiple layers to our model is a logical next step. Furthermore, we leave out the pressure gradient in the hot post-shock gas on the inner side of the shell from our model. Since it is so hot ($`10^7`$ K), it should cool very slowly. It is even harder to include the contribution from this part of HII regions in the model.
The pressure gradient accelerates gas inside the shell, thus decreasing the surface density. It does not affect our flux maps because the brightness of each cell in our model depends only on the number of ionizing photons, as long as the shell is only partially ionized. However, the change in the surface density will affect the thickness of the neutral part of the shell, thus affecting the appearance of corresponding molecular flux maps. A kinematic study of the neutral component associated with ionized bow shock structures might prove useful.
If pressure gradient acceleration is ignored, the spatial and kinematic structure of a modeled bow shock depend on two parameters: the speed of the star and the stand-off distance. The dependence of the stand-off distance on the stellar mass-loss rate, stellar wind terminal speed, stellar speed, and the ambient medium density is given in Eq. 2. As is discussed above, inclusion of the effect of the pressure gradient in the ionized gas has little effect on the spatial distribution, but changes the velocities much like an increase of the stellar speed would. A set of models for UCHII regions in G29.96 and Mon R2 is given in Table 1, where plausible values, based on other observations (see references in the table), are given for unconstrained parameters.
## 5 Comparison with Observations
Our relaxation models of bow shocks provide the shape of the shell and the Doppler shift at each point on the shell. We calculated the \[Ne II\] emission at each point on the shell, assuming ionization-recombination balance in the shell, taking into account the angle at which the ionizing radiation hits the shell. We ignore any dust that might be present in the region.
The model ionized shell was tilted relative to the line of sight to improve the agreement with the data, sampled and remapped onto the sky with a spatial grid $`0.2^{\prime \prime }\times 0.2^{\prime \prime }`$, and then convolved with TEXES’s velocity resolution of $`4`$km s<sup>-1</sup> and a model of TEXES’s beam of half-maximum radius $`r0.8^{\prime \prime }`$. For G29.96, we convolved with an additional turbulent linewidth of $`10.0`$km s<sup>-1</sup> to improve the agreement with the observations. In Mon R2, we convolved the data with the the narrowest line width $`8.8`$km s<sup>-1</sup> (Jaffe et al., 2003). Our observing sampling grid is much coarser than our modeling grids in the region corresponding to our observations, so the error caused by the quantization of the modeling grids is small.
### 5.1 G29.96 -0.02
Four model parameters are constrained with varying uncertainties by our observations: the standoff distance ($`r_0`$), the tilt of the symmetry axis from our line-of-sight, the speed of the star relative to the ambient medium and LSR velocity of the star (or equivalently of the molecular cloud). In addition, the stellar ionizing luminosity and spectral type, and the neon abundance combine to determine the \[Ne II\] brightness. Other parameters, notably those from which the standoff distance can be calculated, can be determined from other observations of our objects or estimated from typical O-star properties. Parameters determined in this way are given in Table 1.
We rotate the observed flux datacube $`26^{}`$ (assuming p.a.=$`64^{}`$, Martín-Hernández et al., 2003) counter-clockwise in order to orient the axis of the shell horizontally in the figures. We also position the observed flux map so that the position of the ionizing star matches that in our model. A cross at (0, 0) in the map (Fig. 1) indicates the position of the ionizing star. The spatial distribution of the \[Ne II\] emission is reproduced better if the motion of the star in the model is within $`45^{}`$ of the plane of the sky, but a bigger angle tends to give a better fit in position-velocity diagrams. We chose to tilt the shell $`50^{}`$ away from us in our model.
As discussed in section 4.3, we assume that the ionized gas and the neutral gas in the shell are moving together, and we pick parameters so that the shell is only partially ionized. A stellar speed of $`v_{}\mathrm{\hspace{0.17em}20}`$km s<sup>-1</sup> is picked to fit the observed range of velocities. An LSR velocity of the star of $`V_{,LSR}=104`$km s<sup>-1</sup>, or $`V_{amb,LSR}=89`$km s<sup>-1</sup> is needed to fit the velocity offset observed. However, observations of the nearby molecular material (Churchwell et al., 1990; Cesaroni et al., 1992; Afflerbach et al., 1994; Olmi & Cesaroni, 1999; Lumsden & Hoare, 1999; Martín-Hernández et al., 2003) give $`V_{amb,LSR}=92100`$km s<sup>-1</sup>, with most numbers near $`98`$km s<sup>-1</sup>. A stellar speed of $`v_{}\mathrm{\hspace{0.17em}10}`$km s<sup>-1</sup> would allow better agreement between the model velocity offset and the molecular observations, but would not fit the observed \[Ne II\] velocity range unless the pressure gradient acceleration is larger than in our model.
The simulated \[Ne II\] flux map and P-V diagrams on cuts through various positions in the model are shown with corresponding diagrams from observed data in Fig. 7-8s. The accelerating, paraboloidal like flow produced in the bow shock model successfully matches many global features in the observed data. The model produces the limb brightening at the head region, although the fit would be improved by using a smaller tilt away from the observer. The curvature of the bright ridge is also fit well. Cuts perpendicular to the shell axis show similar central velocity shifts, spatial and spectral ranges, and overall shape, including a “$`>`$”, in both the observations and the model (Fig. 7). From cut to cut, the curvature of the “$`>`$” changes because the motion of the shell with the star gives more redshift to gas closer to the apex than to down-stream gas. Our line-of-sight passes through the ionized shell twice for positions near the symmetry axis. Given the orientation of the UCHII region, the far side is closer to the apex, contains higher density gas, and dominates the line emission. The near side is farther from the apex and has lower density gas. At the edges of the cuts, we also see more lower density gas. The “$`>`$” forms because gas closer to the apex is less blue-shifted, compared with gas farther from the apex. Because the cuts are perpendicular to the symmetry axis, the model P-V diagrams are necessarily symmetric. Asymmetries in the observations are also small. In cuts parallel to the symmetry axis (Fig. 8), the P-V diagrams show a “7” like pattern, with a curved leg and a rather flat top “arm”, indicating a large velocity gradient at the head and less velocity change towards the tail. The line is also broader at the head region. These are predicted by the model too. In our velocity plot (Fig. 5), we can see that more than $`80\%`$ of the change in the tangential velocity happens in the first $`\pi /2radians`$. The line is broader because the scale length for velocity change is smaller in the head region. The remarkable similarity of the models and the observations is a clear indication for a large-scale paraboloidal like flow in this UCHII region.
### 5.2 Mon R2 IRS1
Our observations of the compact HII region in Mon R2 IRS1 are presented in a previous paper (Jaffe et al., 2003). With a $`24^{\prime \prime }`$ diameter shell and a bright southeast ridge, \[Ne II\] emission in Mon R2 shows complex and broad velocity structure. We speculated that we were looking at a kinematic pattern in which material flows from the bottom to the rim of a bowl-like feature. Here, we try using the bow shock model to interpret the observations, since the object also has a cometary appearance. In order to show the similarities between the kinematics of a shell-like flow structure and the observed data, we choose parameters so that the pressure gradient acceleration is negligible. The results of two models with different standoff distances are shown in Figs. 9-14. In the first model, we try to match the ionizing star position to that in near-infrared observations that are good to $`<1^{\prime \prime }`$ (Yao et al., 1997). In the second model (Table 1), we use a bigger standoff distance, which gives a better fit to our \[Ne II\] observations. In both cases, we tilt the shell $`20^{}`$ toward the observer in the simulation for a good fit to the observed P-V diagrams. As in the G29.96 -0.02 case, we derive a $`V_{amb,LSR}`$ different from the value we found in the literature for Mon R2. An additional $`8`$ km s<sup>-1</sup> redshift is needed to shift the center of the line to the value in the observational data. It is interesting to note that Mon R2 requires a redshift and is tilted toward us while G29.96 is tilted away and requires additional blueshift. This may indicate the pressure gradient may have a bigger effect on the gas acceleration than we calculate and that a more sophisticated model might fit the data with a smaller stellar velocity.
Qualitatively, the model agrees with the data. The observed morphology appears cometary, with a bright arc and a fainter tail, although the tail ends more abruptly than in the model. The P-V diagrams show two peaks over most of the region, as is expected where our line of sight passes through the front and back sides of a shell.
If the predominant motion were due to expansion of the shell, rather than flow along the surface, the velocity splitting would increase only gradually moving from the edge of the shell into the center. As in G29.96, the motions in Mon R2 are consistent with gas flow along a shell. However, there are a few facts which make Mon R2 HII region hard to fit into the bow shock model. First, there is a compact broad-lined region near or just inside of the apex of the observed shell, which is not predicted by the model. It is most apparent in Fig. 11 and Fig. 14, where it produces the central ridge in the P-V diagrams. This component is also shown in our previous paper (Jaffe et al., 2003, Fig. 7). This source may be a result of the shock front overtaking a dense clump in the molecular cloud, or, it could result from an instability in the front, as is seen in the hydrodynamic models of Comeron & Kaper (1998). The second difficulty is that a larger standoff distance is required to make the curvature of the shell the same as in the data. In the first cut of Figs. 9 and 12, where our observed P-V diagrams only show one component, our model P-V diagrams show two components. In addition, the shell appears to be closed on the back side. The outermost contours on the images are nearly circular, although the rim is much brighter at the bottom of the map than at the top. The model predicts that the shell is always open toward the tail.
In our figures 9,10 for both the model and the data, cuts parallel to and near the shell axis show two components. At the ridge of the shell-like region, these two components are connected by a broad line with a width up to $`40`$km s<sup>-1</sup>. The relative strengths and the separation of the two components changes with the position of cuts. Normally, the blue-shifted component is seen farther toward the tail because the tilt of the shell puts the denser part of the near side of the shell farther from the head than the far side. The single broad line in the P-V diagram of the cut made in the center of the ridge is probably caused by the additional source there.
The observed P-V diagrams of the cuts perpendicular to the shell axis (Fig. 11,14) demonstrate the gradual change of a circular structure, which is typical for cuts across a rotationally symmetric shell. If we neglect the broad-line component at the peak of the flux map, our model and observed P-V diagrams agree well. Compared to the observations, the model with a smaller standoff distance predicts a smaller spatial span (Fig. 11). Once we increase the standoff distance (Fig. 14), we no longer have this problem. A density gradient along the symmetry axis of the shell might be able to explain the curvature of the shell, although there is at present no concrete evidence for such a gradient.
## 6 Discussion
We have shown that the overall velocity structure behind the cometary shape of UCHII regions in G29.96 -0.02 and Mon R2 is formed when ionized gas flows along a paraboloidal like surface. Using the bow-shock model, we can reproduce this kind of structure. We found that the model provides a good qualitative explanation of our observations of G29.96 and Mon R2. Both sources have morphologies similar to that predicted by the model, at least near the apex of the shell. The observed line profiles, as seen in the P-V diagrams, are generally double-peaked, indicating that the gas is swept up into a thin shell, so each line of sight passes through two surfaces with different Doppler shifts. The two Doppler components reach a maximum separation near the shell apex, indicating that the gas accelerates along the shell, and that the dominant motion is along the shell rather than radial.
There are quantitative differences between the model and the data, especially in the case of Mon R2, as discussed above, indicating that a complete description of these sources will have to be more complicated than our simple model. Nevertheless, we think that the bow shock model must be essentially valid. Now we consider whether this model resolves any of the problems associated with the study of UCHII regions, and how the model might be improved.
The current bow shock model assumes that the stellar wind material and the ambient medium mix well and cool efficiently after they collide, so their mass and momentum contribute to the swept-up layer and the velocity of the ionized gas is that of the single swept-up layer. However, this momentum-conserving thin shell assumption is not well justified. Behind the shock, stellar wind material is collisionally heated and ionized. This gas will have a temperature $`>10^7`$K and does not cool efficiently, keeping the layer thick. This high temperature layer forms between the stellar wind and the photonionized gas layer and makes the mass and momentum exchange between the two layers difficult. The thin shell assumption is invalid in this situation. The photonionized layer is also evaporated through conduction with the extremly hot shocked stellar wind gas, which makes predicting the kinematics more difficult. There are many instabilities which can inflate the shell to make the thin shell assumption invalid. Possible candidates include the transverse acceleration instability (TAI) (Dgani et al., 1996) caused by the acceleration of the flow normal to the surface of the shell and the non-linear thin shell instability (NTSI) (Vishniac, 1994; Hueckstaedt, 2003) caused directly by the collision of isothermal flows. These instabilities will disturb the shell and create sub-structures with scales comparable to the thickness of the bow shock. Under such conditions, the momentum conservation assumption will not hold.
Among the other models proposed to explain cometary UCHII regions, the champagne model is the most interesting. The model applies when a massive star is found in a region with a large density gradient, such as the edge of a dense molecular cloud. The resulting HII region will expand supersonically away from the high-density region, or simply break out of the edge of the molecular cloud and cause the ionized gas to stream out of the opening in response to a large pressure gradient. The classic champagne model without a stellar wind (Tenorio-Tagle, 1979; Bodenheimer et al., 1979; Yorke et al., 1983) can explain the cometary shape of HII regions but, because the ionized gas fills up the bubble and the main pressure gradient is along the density gradient, the champagne model has difficulty accounting for limb brightening and for the line profiles observed. Adding a stellar wind to the champagne model will probably produce limb brightening (Comeron, 1997). Without ram pressure of the external gas, the champagne model with a stellar wind tends to produce a bigger shell in the clouds with the same density and temperature as the clouds we investigated with the bow shock model.
The main differences between the bow shock and the champagne flow are seen in the kinematic properties of the ionized gas (Garay & Lizano, 1999). First, the bow shock model predicts that the velocity gradient is steeper in the head than in the tail, which can be seen easily in our P-V diagrams, whereas the champagne model expects the largest gradient in the tail, where the fractional pressure gradient also reaches the highest value (Comeron, 1997). Second, the champagne model predicts that the line widths are broader in the tail region because the gas has higher speed there and the gas motion is less parallel, while the bow shock predicts broader line widths near the apex because the gas gains more momentum there and moves cylindrical-symmetrically along the surface (in the frame of reference of the star). Finally, ionized gas near the apex of the cometary structure is at rest with respect to the ambient molecular gas in the champagne model, whereas it moves with the star and/or the shock front in the bow shock model. In all of these respects, our observations of G29.96 and Mon R2 are in better agreement with the bow shock model.
A central problem about UCHII regions is their lifetimes. The number ratio of UCHII regions and OB stars requires that UCHII regions should have an average lifetime $`20\%`$ of that of OB stars (Wood & Churchwell, 1989a). Any model of UCHII regions has to be able to account for this. We have observed over a dozen of UCHII regions. They all have large velocity range, usually over $`20`$ km s<sup>-1</sup>. In the bow shock model, the velocity range of the line is directly connected to the star speed relative to the ambient medium. If $`20\%`$ of UCHII regions have a bow shock-like structure, as shown in Wood & Churchwell (1989a, b) and Kurtz et al. (1994), the same percentage of OB stars should move supersonically through molecular clouds. Although evolved OB stars can be accelerated up to $`200`$km s<sup>-1</sup> through the association ejection or supernova explosions (Blaauw, 1993) and form bow shock-like structures around them (van Buren et al., 1995; Noriega-Crespo et al., 1997; Kaper et al., 1997), high speed OB stars are rare in OB associations. The velocity dispersion of OB association is generally small, only a few km s<sup>-1</sup> (Jones & Walker, 1988; Tian et al., 1996). Even if the molecular gas initially made the potential a bit deeper, stellar speeds of $`1520`$ km s<sup>-1</sup> are extremely improbable as part of a normal distribution. Including pressure gradient acceleration helps to reduce this requirement. Our calculations suggest that the needed speed is still too high to solve the problem. This disagreement indicates that the stellar speeds may be less than in our model. Due to the simplicity, our model could underestimate the acceleration of the ionized gas along the shock front. A full hydrodynamical treatment would help determine the true effect of a pressure gradient.
We are grateful to Alan Fey and Ed Churchwell for letting us use their VLA 2 cm data (Fey et al., 1995; Wood & Churchwell, 1989b). We thank Gregory Shields for help on the ionization model. We also need to thank NASA IRTF staff for their help on observations. This work was supported by NSF grant AST-0205518, NSF grant AST-0307497 and NASA grant NNG04GG92G. Thomas Greathouse is currently supported by the Lunar and Planetary Institute, which is operated by the Universities Space Research Association under NASA CAN-NCC5-679.
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# Optical sum in Nearly Antiferromagnetic Fermi Liquid Model
## I Introduction
The high $`T_c`$ oxides fall in the category of highly correlated systems. A manifestation of this fact is that in the underdoped regime there exists a pseudogap.r1 ; r2 ; r3 ; r4 ; r5 Precisely how it is to be described remains controversial.r5 ; r6 ; r7 ; r23a ; r8 ; r9 ; r10 ; r11 ; r12 ; r13 In certain theories it is closely related to superconducting correlations,r5 ; r6 ; r7 ; r23a and the superconducting transition temperature $`T_c`$ is the temperature at which phase coherence is lost. In other theories the pseudogap has its origin in completely different correlations and is a manifestation of a competing order such as a $`d`$-density wave (DDW).r8 ; r9 ; r10 ; r11 ; r12 ; r13 In any case interactions among charge carriers play an important role in such systems and cannot be ignored in any realistic approach to their properties even around optimum doping which is the case of interest here.
While many aspects of the superconducting state can be understood qualitatively on the basis of extensions of BCS theory, particularly around optimum doping, the search for essential differencesr14 has remained an important avenue of investigation. In particular the idea of kinetic energy as opposed to potential energy driven superconductivity (i.e.: the kinetic energy is reduced at the transition temperature) and its relation to the OS has recently been given serious consideration in theoryr14 ; r15 ; r16 ; r17 ; r18 ; r19 ; r20 ; r21 ; r22 ; r23 and experimentally.r24 ; r25 ; r26 ; r27 ; r28 ; r29 An issue of importance is the relation of the kinetic energy (KE) to the optical sum (OS). A recent paperr28 has provided some insight into this relationship and has given a comparison with experiment for the temperature variation of the OS in the normal state, its change at $`T_c`$, and its further evolution in the superconducting state. Correlations beyond BCS pairing were not considered, however. There have also been several recent studies of the temperature dependence of the OS, experimentally for the system LA<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, Ref. r42a, , and, theoretically, using other models.r42b ; r42c ; r42d
Within a BCS model, the condensation is potential energy driven and, in fact, the KE increases because the probability of occupation of the state k $`(n_𝐤)`$ becomes smeared at the Fermi energy by the opening of the gap. This translates into a decrease in the OS as compared with its normal state value. But this is opposite to the behavior obtained experimentally in Ref. r26, . However, as mentioned, the theoretical discussion of Ref. r28, is based on a non interacting model in the normal state and includes pairing correlations in the superconducting state only at the level of BCS. More sophisticated formulations of the theory of superconductivity could give different results. Also, interactions can significantly modify the results even in the normal state, as has been demonstrated recently by Knigavko et al.r30 in a simplified model in which the charge carriers are coupled to a single Einstein mode. While only the normal state was considered, it was found that boson hardening (softening) results in an increase (decrease) in the OS and that interactions play an important role in determining temperature dependences.
In this paper we study a tight binding model with an emphasis on the effect interactions can have on the OS, particularly on its temperature dependence and its relationship to the KE. Here, the interactions among the charge carriers are treated in the Nearly Antiferromagnetic Fermi Liquid (NAFFL) model.r31 A review of its main properties and successes is given in Ref. r31, . The model is phenomenological and falls into the general class of boson exchange models where the interaction between electrons proceeds through the exchange of spin fluctuations.r32 The imaginary part of the spin susceptibility replaces the phonon propagator of the classic electron-phonon Eliashberg theory. The spin susceptibility could be calculated from microscopic theory. More usually, however, it is fit to experimental data, specifically to NMR in Ref. r32, . The basic idea is that the doped metallic cuprates are near an antiferromagnetic phase boundary and that coupling to spin fluctuations is therefore strong. While there is no consensus as to the validity of such a model when applied to the oxides the NAFFL model has been widely discussed and has had considerable successes particularly in correlating superconducting properties,r33 ; r34 ; r35 ; r35a ; r35b ; r35c and has been used to give a detailed description of the optical propertiesr35 ; r36 ; r37 ; r38 ; r39 ; r40 ; r43a in the high $`T_c`$ oxides at optimum doping. In any case, the NAFFL modelr41 provides a convenient and specific framework within which we can study the effect of correlations on the OS and on related properties.
In Sec. II we summarize the basic equations that are needed to compute the optical sum integral as well as the KE. They are a set of three coupled generalized Eliashberg equations written for any momentum k in the two dimensional CuO<sub>2</sub> Brillouin zone. They involve renormalized Matsubara frequencies, the renormalized quasiparticle energies as well as the superconducting energy gap. The interaction between the charge carriers is mediated by the exchange of spin fluctuations and involve the spin susceptibility. Fast Fourier transforms (FFT) provide solutions which give a $`d`$-wave gap as observed in experiments. In Sec. III we apply our solutions to evaluate the probability of occupation of the state of momentum k and spin $`\sigma `$ $`(n_{𝐤,\sigma })`$ from which the OS and the KE follow. Their temperature dependence is studied. Comparison with the non interacting case is made and it is concluded that interactions can profoundly modify results. Different behaviors can result depending on the choice of microscopic parameters. In Sec. IV we consider explicitly the superconducting state. Here, again, generalized Eliashberg equations give variations with $`T`$ which are considerably different from earlier BCS results. When the electron-exchange boson interaction is taken as temperature independent and independent of state, the OS is lower in the superconducting than in the normal state. Nevertheless, it can keep increasing with decreasing temperature in contrast to BCS where it was found to decrease. This increase can be traced to the underlying temperature dependence of the OS which depends on details of the band structure and interactions involved, for example on the model spin susceptibility. Further, if we take account of the low energy gaping of the spin susceptibility which is brought about by the superconducting transition (a process which is not operative in the normal state) the KE can be further decreased and there is an additional increase in the OS which can effectively increase faster than it does in the normal state. In this case the KE in the superconducting state with low energy gaping of the spin susceptibility can be less than in the normal state without low energy gaping. An important conclusion of our work is that the observation of a faster increase in the OS with decreasing temperature in the superconducting state than in the normal state cannot unambiguously be taken to be an indication of kinetic energy driven superconductivity in contrast to a recent claim by van der Marel et al.r28 In Sec. V we provide a brief conclusion.
We use units in which $`\mathrm{}=c=1`$ throughout this paper.
## II Formalism
In the NAFFL model the interaction between holes proceeds through the exchange of spin fluctuations and the spin susceptibility $`\chi (𝐪,\omega )`$ plays a central role. The three Eliashberg equations for the renormalized frequencies $`\stackrel{~}{\omega }(𝐤,i\omega _n)`$, the energy renormalization $`\xi (𝐤,i\omega _n)`$ and the pairing energy $`\varphi (𝐤,i\omega _n)`$ as a function of momentum k in the two dimensional CuO<sub>2</sub> Brillouin zone, and of fermionic Matsubara frequencies $`i\omega _n=i\pi T(2n+1),n=0,\pm 1,\pm 2,\mathrm{}`$ and the temperature $`T`$ arer35a ; r35b ; r35c
$`\stackrel{~}{\omega }(𝐤,i\omega _n)`$ $`=`$ $`\omega _n+T{\displaystyle \underset{m}{}}{\displaystyle \underset{𝐤^{}}{}}\lambda _{SF}(𝐤𝐤^{},i\omega _ni\omega _m)\mathrm{¸}`$ (1a)
$`\times {\displaystyle \frac{\stackrel{~}{\omega }(𝐤^{},i\omega _m)}{D(𝐤^{},i\omega _m)}},`$
$`\xi (𝐤,i\omega _n)`$ $`=`$ $`T{\displaystyle \underset{m}{}}{\displaystyle \underset{𝐤^{}}{}}\lambda _{SF}(𝐤𝐤^{},i\omega _ni\omega _m)`$ (1b)
$`\times {\displaystyle \frac{ϵ_𝐤^{}+\xi (𝐤^{},i\omega _m)}{D(𝐤^{},i\omega _m)}},`$
$`\varphi (𝐤,i\omega _n)`$ $`=`$ $`T{\displaystyle \underset{m}{}}{\displaystyle \underset{𝐤^{}}{}}\lambda _{SF}(𝐤𝐤^{},i\omega _ni\omega _m)`$ (1c)
$`\times {\displaystyle \frac{\varphi (𝐤^{},i\omega _m)}{D(𝐤^{},i\omega _m)}}.`$
In Eqs. (1) $`D(𝐤,i\omega _n)`$ is given by
$$D(𝐤,i\omega _n)=\stackrel{~}{\omega }^2(𝐤,i\omega _n)+\left[ϵ_𝐤+\xi (𝐤,i\omega _n)\right]^2+\varphi ^2(𝐤,i\omega _n),$$
(2)
and $`ϵ_𝐤`$ is the charge carrier dispersion relation. In a tight binding model without the inclusion of the coupling to the spin fluctuations it is given by:
$`ϵ_𝐤`$ $`=`$ $`2\{t[\mathrm{cos}(ak_x)+\mathrm{cos}(ak_y)]`$ (3)
$`2t^{}\mathrm{cos}(ak_x)\mathrm{cos}(ak_y)\}\mu ^{}.`$
Here, $`a`$ is the lattice parameter in the copper oxide plane, $`t`$ the nearest neighbor hopping, $`t^{}`$ the next nearest neighbor hopping, and $`\mu ^{}`$ the chemical potential. We will discuss within this context, two tight binding models with the parameters given in Table 1. The corresponding Fermi surface is presented in Fig. 1a for model A and in Fig. 1b for model B. These figures define also the various points in the Brillouin zone. The dotted lines indicate the antiferromagnetic Brillouin zone boundary. For model A, the Fermi surface crosses the anti-ferromagnetic
Brillouin zone around $`X`$ and symmetry related points. These points are ‘hot spots’ for which Fermi surface to Fermi surface transitions are possible with momentum transfer $`(\pi /a,\pi /a)`$ (nesting property). Model A has been used previously in Ref. r28, to discuss the OS and KE in Bi<sub>2</sub>Sr<sub>2</sub>Cu<sub>2</sub>O<sub>8+δ</sub> (BSCCO) assuming no interactions between the electrons in the normal state. Therefore, it is natural to employ this same model in the present study which extends the previous calculations to include exchange of spin fluctuations between charge carriers in the NAFFL model. The second model, Model B, was chosen to contrast with the first. It has no hot spots and its Fermi surface is closed around the $`\mathrm{\Gamma }`$-point in contrast to Model A which has a Fermi surface which is closed around the $`M`$-point. It also has a smaller value of $`t`$ which, on its own, would imply a smaller absolute value of KE. These differences in band parameters lead, as we shall see, to some differences in KE and optical sum at $`T=0`$ in the normal state.
In the phenomenological model of Millis et al.r31 ; r32 and Monthoux et al.r41 (MMP-model) the magnetic susceptibility was fit to NMR data and its imaginary part takes on the form
$$\mathrm{}\mathrm{m}\chi _{MMP}(𝐪,\omega )=\frac{\chi _𝐐(\omega /\omega _{SF})}{\left[1+\zeta ^2(𝐪𝐐)^2\right]^2+(\omega /\omega _{SF})^2},$$
(4)
where $`\chi _𝐐`$ is the static susceptibility, Q is the commensurate antiferromagnetic wave vector $`(\pi /a,\pi /a)`$ in the upper right hand quadrant of the CuO<sub>2</sub>-plane Brillouin zone and symmetry related points. $`\zeta `$
is the magnetic coherence length, and $`\omega _{SF}`$ a characteristic spin fluctuation frequency. We set $`\zeta =2.5a`$ throughout this paper and various values for $`\omega _{SF}`$ are investigated.
The kernel $`\lambda _{SF}(𝐪,i\nu _{nm})`$ in Eqs. (1) with momentum transfer $`𝐪=𝐤𝐤^{}`$ and the bosonic Matsubara frequency $`i\nu _{nm}=i\omega _ni\omega _m`$ is given as
$$\lambda _{SF}(𝐪,i\nu _n)=\frac{g^2\chi _𝐐}{1+\zeta ^2(𝐪𝐐)^2+(|\nu _n|/\omega _{SF})},$$
(5)
with $`g^2\chi _𝐐`$ adjusted to get the desired value of the critical temperature $`T_c`$ for a certain value of $`\omega _{SF}`$ from the solution of the linearized Eqs. (1). This defines the model.
The aim of this paper is to investigate the effect interactions have on the OS defined as
$$\pi e^2I_\sigma =\underset{\mathrm{\Omega }}{\overset{\mathrm{\Omega }}{}}𝑑\omega \mathrm{}\mathrm{e}\sigma _{xx}(\omega )=\frac{\pi e^2}{V}\underset{𝐤,\sigma }{}n_{𝐤,\sigma }\frac{^2ϵ_𝐤}{k_x^2},$$
(6)
where $`e`$ is the charge on the electron, $`V`$ the volume, and $`n_{𝐤,\sigma }`$ is the probability of occupation of a state of momentum k and spin $`\sigma `$. Finally, $`\sigma _{xx}(\omega )`$ is the optical conductivity. The integral in Eq. (6) is to be taken over the single band with $`\mathrm{\Omega }`$, the upper limit in the integral of Eq. (6), large enough to include all possible transitions in that band. We are also interested in the relationship between the OS and the kinetic energy. By definition
$$I_{\mathrm{KE}}=H_{\mathrm{KE}}=\frac{a^2}{V}\underset{𝐤,\sigma }{}n_{𝐤,\sigma }ϵ_𝐤.$$
(7)
We will see that, to a good approximation $`I_\sigma `$ and $`I_{\mathrm{KE}}`$ are nearly proportional to each other. Thus,
$$\rho _L\frac{1}{\pi e^2}\underset{\mathrm{\Omega }}{\overset{\mathrm{\Omega }}{}}𝑑\omega \mathrm{}\mathrm{e}\sigma _{xx}(\omega )\frac{1}{2}H_{\mathrm{KE}}.$$
(8)
holds approximately. (An equal sign would be appropriate if the dispersion relation (3) contained only nearest neighbor interaction, i.e.: $`t^{}=0`$.) Here, $`\rho _L`$ is the experimentally determined value of the OS $`(I_\sigma )`$.
Interactions have a profound effect on the probability of occupation of the state $`|𝐤,\sigma `$ which would be one or zero for occupied and unoccupied states respectively in the non interacting case. In Fig. 2 we show results for $`n_𝐤`$, along certain selected
directions in the CuO<sub>2</sub> Brillouin zone. In all cases Model A of Table 1 is used for the electronic dispersion (3) with $`\mu ^{}`$ adjusted to the required filling which is is defined as
$$n=\frac{1}{2}\underset{𝐤}{}\underset{n0}{}\frac{ϵ_𝐤+\xi (𝐤,i\omega _n)}{\stackrel{~}{\omega }^2(𝐤,i\omega _n)+\left[ϵ_𝐤+\xi (𝐤,i\omega _n)\right]^2+\varphi ^2(𝐤,i\omega _n)},$$
(9)
and the charge carrier spin fluctuation strength $`g^2\chi _𝐐`$ is adjusted to get a $`T_c=90`$K for the superconducting state. In the top frame of Fig. 2 we show $`n_𝐤`$ in the non interacting case as we go from $`\mathrm{\Gamma }`$ to $`X`$ and from $`X`$ to $`M`$ in the Brillouin zone with the Fermi surface defined as the value of k at which $`n_𝐤`$ jumps from one to zero. It is obvious from Fig. 1a (dashed line) that the path $`\overline{XM}`$ crosses the Fermi surface. A second crossing of the Fermi surface can be observed along the path from $`M`$ to $`\mathrm{\Gamma }`$.
The center frame of Fig. 2 shows $`n_𝐤`$ when interactions are taken into account. (The corresponding Fermi surface is shown as the solid line in Fig. 1a.) We see a drastic difference in the value of $`n_𝐤`$ which is now of the order 0.9 at the $`\mathrm{\Gamma }`$ point indicating that the effect of the interaction is very significant even in the center of the Brillouin zone. Also, $`n_𝐤`$ is of the order 0.1 outside the non interacting Fermi surface where $`n_𝐤0`$ for the non interacting case. The solid line applies to the normal state at $`T=20`$K and the dashed line to the superconducting state at the same temperature. On comparing these two cases we see that the transition to the superconducting state depletes even further the non interacting sea in the sense that it further reduces $`n_𝐤`$ at $`𝐤`$s below the Fermi surface and correspondingly increases the probability of occupation of states right above the Fermi surface beyond the effect of interactions in the normal state. It is to be noted, however, that the onset of superconductivity results in rather modest changes in $`n_𝐤`$ as compared to the difference between interactions and no interactions in the normal state.
The bottom frame of Fig. 2 gives results in the normal state but compares two temperatures, namely $`T=20`$K (solid line) and $`T=150`$K (dashed line). Comparison with the middle frame shows that increasing the temperature has roughly the same qualitative effect on $`n_𝐤`$ as does the transition to the superconducting state. In both cases, the KE given by Eq. (7) increases because the states of lower $`ϵ_𝐤`$ get depleted while states with higher $`ϵ_𝐤`$ are occupied with increasing probability. This will also hold for the OS according to Eq. (6) which will now depend on interactions and on temperature.
## III Results for the optical sum in the normal state
In Fig. 3 we show results for the
optical sum $`(I_\sigma )`$ Eq. (6) and compare with the kinetic energy $`(I_{\mathrm{KE}})`$, Eq. (7). The top frame is based on Model A and the bottom frame on Model B of Table 1. The solid squares and circles are $`I_{\mathrm{KE}}/2`$ and $`I_\sigma `$ respectively in the free tight binding case, i.e.: no interactions, plotted as a function of the square of the temperature for the normal state. (A $`128\times 128`$ sampling of the $`𝐤`$-space, $`ak_x,ak_y[0,\pi ]`$, was used but going to a $`256\times 256`$ sampling did not influence the results.) In both models variation with $`T`$ over the range $`0`$ to $`200`$K is small (of order 1% for Model A and 2% for Model B) as was also found in the work of Molegraaf et al.r26 Also, the two integrals $`(I_\sigma `$, $`I_{\mathrm{KE}}/2)`$ track each other closely even though they are not equal in magnitude. (They would be equal for $`t^{}=0`$. We tried other Fermi surfaces, even one with perfect nesting, i.e.: $`t^{}=0`$ and $`n=0.5`$, and found no qualitative changes.) These results are for comparison with results indicated by up/down-triangles which include interactions in the NAFFL model. The magnitude of both, the optical sum integral and the kinetic energy has been changed considerably by the interactions although the order remains the same, i.e.: $`I_{\mathrm{KE}}/2`$ is greater than $`I_\sigma `$ and, again, they track each other. More importantly, for the discussion here, the temperature variation has been changed. Both integrals now show variation of the order 8 to $`9\%`$. Also the dependence on $`T^2`$ is not linear at small values of $`T^2`$. Model A shows similar behavior for small values of $`\omega _{SF}`$. It is clear that any estimate based on the independent particle tight binding model is unreliable. It is, however, possible to chose specific parameters in the MMP-model which show variations in the interacting case that are much closer to the non interacting case. This is illustrated in the top frame of Fig. 3. Here we used Model A of Table 1. Again, results with and without interaction are compared and both show little temperature variation. To get this we used $`\omega _{SF}=82`$meV in our MMP form of Eq. (4) without a change in the magnetic coherence length. Both results, with and without interaction, show little variation with temperature. What this shows is that the magnitude as well as the temperature variation of KE and of OS depends significantly on the parameters used to characterize their electronic structure, particularly the spin susceptibility.
We have done additional calculations for Model A with $`\omega _{SF}=40,\mathrm{\hspace{0.17em}20},`$ and $`10`$meV. In all cases the change in KE due to interactions at $`T=0`$ increases with decreasing values of $`\omega _{SF}`$. In particular, it changes at $`T=0`$ by 5.3% when compared with the non interacting case, for $`\omega _{SF}=82`$meV and by 15.8% for $`\omega _{SF}=10`$meV. The corresponding temperature changes from $`T=0`$ to $`T=200`$K are roughly a factor of 5 smaller, more precisely, they are 0.7% and 3.4% respectively. Thus, a change in KE at $`T=0`$ due to interactions also implies a corresponding change in temperature dependence with both changes tracking each other. For Model B the change in KE due to interactions is 34% for $`\omega _{SF}=10`$meV with a 8.7% increase in KE from $`T=0`$ to $`T=200`$K. These variations are about a factor of two larger than for the equivalent case of Model A with a comparable value of $`\omega _{SF}`$. Despite the fact that the two models represent very different band structures $`I_\sigma `$ and $`I_{KE}/2`$ show essentially the same qualitative features in their temperature and $`\omega _{SF}`$ dependence. However, the quantitative differences are important.
## IV Results for the optical sum in the superconducting state
Results for the superconducting state are illustrated in Fig. 4
which has two frames. The top frame applies to the band structure Model A of Table 1 and is for $`\omega _{SF}=82`$meV as in the top frame of Fig. 3. We have also included 2% impurities in the unitary limit but this serves mainly to illustrate that impurities introduce no qualitative differences into our results. We see that, as we expect, superconductivity reduces the optical integral (open triangles) as compared with its normal state (solid triangles) value at the same temperature. This reduction is small. For the top frame which shows the least temperature dependence, the KE integral shows a reduction of about 0.25% below its normal state value which can be compared with the results shown in the bottom frame of Fig. 6 of Ref. r29, where the difference is 0.2% in their BCS calculations. On the other hand, in the bottom frame of our Fig. 4 for Model B the reduction is about 0.8% (four times larger). This shows that the Eliashberg results depend on band structure as well as on the details of the interactions involved, in particular on the value of $`\omega _{SF}`$. In the BCS limit the increase in KE normalized to the absolute value of the condensation energy is given by the formula $`\left[\mathrm{ln}\left(\frac{\omega _D}{T_c}\right)0.38\right]`$ for both $`s`$\- and $`d`$-wave superconductors. Here $`\omega _D`$ is the Debye energy. This shows a strong dependence on $`\omega _D/T_c`$. The formula itself, however, is valid only for $`\omega _D/T_c1`$ and cannot be used to understand our Eliashberg results. The NAFFL model includes interactions which, as we have seen, change importantly the probability of occupation $`n_𝐤`$ and consequently the optical integral as well as the kinetic energy. For the parameters of Model B and $`\omega _{SF}=10`$meV we find that $`I_\sigma `$ and $`I_{\mathrm{KE}}/2`$ can keep increasing with decreasing temperature in the superconducting state (bottom frame of Fig. 4). The open squares and triangles (superconducting state) are below their solid counterparts (normal state) but still keep growing as the temperatures is reduced. This does not indicate an exotic mechanism but comes directly from a generalization of Eliashberg theory that includes anisotropy in the band structure and, more importantly, the interaction due to coupling to the spin fluctuations.
In Fig. 5 we show additional results where the OS is seen to
increase even more rapidly, with reduced temperature below the onset of superconductivity than it does in the normal state above $`T_c`$. What is shown is $`\rho _L`$ which is $`I_\sigma `$ or $`I_{\mathrm{KE}}/2`$ scaled to agree with experiment as discussed below; either, $`I_\sigma `$ or $`I_{\mathrm{KE}}/2`$, will do since these differ mainly by a different scaling factor. Only the OS is considered but the KE integral follows the same trend and therefore the optical measurement can again be used to get information on KE and its variation with temperature.
While in obtaining Fig. 5 we applied Model A which was used by van der Marel et al.r28 to describe their optimally doped and underdoped BSCCO samples, we now vary, in contrast to the top frame of Fig. 4, the value of $`\omega _{SF}`$ used in the MMP-model for the spin susceptibility Eq. (4). Results are presented for $`\omega _{SF}=20`$meV (dashed line), $`10`$meV (dotted line), $`13`$meV (solid squares) for the normal state, and solid triangles for the superconducting state. Also shown as the thick solid line are the experimental results of Ref. r26, for their optimally doped sample. We have scaled our theoretical results to agree with experiment at $`T=120`$K. We first note that varying $`\omega _{SF}`$ in the normal state can strongly influence the temperature dependence obtained for $`\rho _L`$ (in meV). The value of $`\omega _{SF}=13`$meV was chosen from a best fit in the region $`120\mathrm{K}T200`$K. The scaling factor required to get agreement with this data is approximately 2. (When interactions are neglected, as in Ref. r28, , the scaling factor is approximately 1.5.) To reduce this discrepancy, the value of $`t`$ would need to be increased but this would also decrease the sensitivity of $`I_\sigma (T)`$ to temperature variations and, thus, $`\omega _{SF}`$ would need to be adjusted as well. To prove this we employ results of local density approximation calculations by Markievicz et al.r42 which suggest significantly bigger values for $`t`$ in BSCCO with a Fermi surface which is little different from the one presented in Fig. 2a. Using the dispersion relation, Eq. (3), and the parameter values of Table I of Ref. r42, we find $`I_\sigma (T=0)=267.34`$meV for the non-interacting system, well above the value of $`\rho _L(T=0)171.53`$meV as has been extrapolated from the normal state experimental data of Ref. r26, . In order to reproduce the experimentally observed temperature dependence $`\rho _L(T)/\rho _L(T=0)`$ we have to introduce interactions and a value $`\omega _{SF}=8`$meV is found to give excellent agreement of $`I_\sigma (T)/I_\sigma (T=0)`$ with experiment. In this case $`I_\sigma (T=0)=230.34`$meV, still well above the experimental value and a down-scaling of $`I_\sigma (T)`$ by a factor of 0.745 is required to achieve agreement with experiment. Ultimately, a dispersion relation somewhere between Model A and the one reported by Markievicz et al.r42 and an $`\omega _{SF}`$ between $`8`$ and $`13`$meV will result in a $`I_\sigma (T)`$ from theory which agrees with the experimental $`\rho _L(T)`$ without scaling. However, our main aim is not to treat a specific case but to understand better the role interactions between the charge carriers can play in the OS. Interactions introduce a new energy scale into the problem, namely $`\omega _{SF}`$ for the NAFFL model. This energy scale is additional to the chemical potential or the hopping parameter $`t`$. With the values of the microscopic parameters associated with the NAFFL model just described, we proceed to compute the OS for temperatures at and below $`T_c`$. The solid squares give the continuation of the normal state curve and are presented for comparison with the solid triangles which are the equivalent results in the superconducting state. Again, superconducting state results fall below the normal state ones but they are seen to, nevertheless, increase with decreasing $`T`$. This occurs even if an Eliashberg formulation is used which represents a generalization of BCS theory and, in that sense, is not exotic. The mechanism is the exchange of antiferromagnetic spin fluctuations. On comparing the top frame of Fig. 4 with Fig. 5, we note that whether or not the superconducting state results keep increasing with decreasing temperature is, for a given band structure, governed by the value of $`\omega _{SF}`$.
In Fig. 5 we show additional results, open squares for a model normal state and open triangles for the superconducting state. Now, there is a further dramatic increase in the OS both in the normal and the superconducting state as compared to its value in the normal state at $`T=T_c=90`$K. A detailed explanation of how these results were arrived at is required. In their analysis of optical data Carbotte et al.r34 found that the spin fluctuations themselves are modified when the superconducting state sets in. To carry out their analysis these authors used a simplified version of our Eliashberg Eqs. (1) which follows when the sum over k is changed to an energy integral as well as an angular average and the energy integral is done analytically in a constant density of states approximation for an infinite band model with interactions pinned to the Fermi surface. Here we have been more realistic but what is important for us in the work of Ref. r34, is that they find that the spin fluctuation spectrum is gaped at low energies, or at the very least loses intensity and a spin resonance or peak forms at higher energy. This readjustment in the spin susceptibility is not unexpected and is a characteristic that should be seen in any electronic mechanism for superconductivity.r37 ; r38 ; r39 ; r40 ; r43a ; r41 ; r44 ; r45 ; r46 ; r47 ; r48 Details are not important for the present discussion beyond the fact that some adjustment of the spin susceptibility $`\mathrm{}\mathrm{m}\chi (𝐪,\omega )`$ at small $`\omega `$ is expected, which weakens the inelastic scattering. Here we simply use the same low $`\omega `$ $`[\omega _c(T)]`$ cutoff applied to Eq. (4) which was determined by E. Schachinger et al.,r43a through consideration of microwave data. Another approach would be to calculate the low energy gaping of the spin susceptibility from first principles but this would go beyond the scope of this work and would introduce additional uncertainties. The temperature dependence of $`\omega _c(T)`$ follows the temperature dependence of the superconducting gap with a maximum value of $`24`$meV. Application of this cutoff in otherwise standard Eliashberg calculations based on our Eqs. (1) yield the open triangles (superconducting state) and open squares (normal state) of Fig. 5. The physics underlying these curves has been made clear in a simple model recently studied by Knigavko et al.r30 These authors studied a model in which the charge carriers are coupled to a single Einstein mode of unspecified origin. What they found was that stiffening of this mode decreases the kinetic energy and hence increases the OS. This is precisely the same mechanism that is operative in Fig. 5. By applying a low frequency cutoff to the spin fluctuations in our MMP-model we are decreasing the KE. This decrease in KE, present in the underlying normal state below $`T_c`$, compensates for the increase in KE intrinsic to the superconducting transition which results from the opening of the superconducting gap. We note, however, that in our formulation the OS, at any given temperature, is always below (although not very much) its normal state value at this same temperature calculated with the spectrum with a low frequency cutoff (open triangles). But this cutoff is only operative in the superconducting state and is responsible for making the open triangles fall above the solid squares. The kinetic energy in the superconducting state with low frequency cutoff is now less than the normal state kinetic energy without cutoff. Including the feedback mechanism of the formation of the superconducting state on the spin susceptibility itself has the net effect, at zero temperature (where it is largest), of changing the sign of the KE contribution to the condensation energy from that in BCS.
## V Conclusion
We have used the Nearly Antiferromagnetic Fermi Liquid model to study the effect of interactions on the optical sum and on the kinetic energy in tight binding bands. Comparison of normal state results with equivalent results when interactions are neglected showed that temperature variations can be strongly affected by details of the microscopic parameters involved in the spin fluctuation exchange mechanism. Behaviors are possible which can be quite different from the non interacting independent particle model. Comparison with normal state experimental data proves that the tight binding model of non interacting particles is certainly not adequate to describe properly the temperature dependence of the optical sum. (This has also been observed by Benfatto et al.r42b ) Taking into account interactions between the charge carriers makes the tight binding model a viable model for the analysis of the temperature dependence of the normal state optical sum. This was demonstrated for the particular case of BSCCO and a particle interaction modeled on the NAFFL. Other models, like the one presented by Toschi et al.r42c are also capable to reproduce the temperature dependence $`I_\sigma (T)/I_\sigma (T=0)`$ but they lack agreement with the value of the optical sum at zero temperature.
When superconductivity is considered within an Eliashberg formalism, the superconducting gap has $`d`$-wave symmetry as a function of momentum in the two dimensional Brillouin zone. The optical sum is found to decrease with decreasing temperature for some range of parameters characterizing the spin susceptibility but can also increase. This increase cannot necessarily be interpreted as kinetic energy driven superconductivity. In fact, in all cases considered, the optical sum is always lower, at a given temperature in the superconducting state, than it is in the corresponding normal state but, in some cases not by much. Correspondingly, the kinetic energy is increased in the superconducting state. What makes the optical sum and KE integral continue to go up (in some cases) with decreasing temperature is the fact that the interactions themselves introduce a temperature dependence in the underlying normal state.
The results just described were obtained for a fixed (i.e.: temperature independent) value of the spin susceptibility. If we consider the possibility that the spin fluctuation spectrum may itself be modifiedr33 ; r34 ; r35 ; r44 ; r45 ; r46 ; r47 ; r48 by the onset of superconductivity and by temperature, even larger increases in the optical sum with decreasing temperature can be obtained. It is widely recognized that a generic feature of an electronic mechanism of superconductivity is the possible gaping of the excitation spectrum itself at small energies due to the opening of the superconducting gap. This leads to a weakening of interactions at small $`\omega `$ and to the so called collapse of the inelastic scattering rater44 ; r45 ; r46 ; r47 ; r48 which manifests itself as a large peak in the temperature dependence of the microwave conductivity. The weakening of the interaction in the superconducting state through the opening of a low energy gap in the spin susceptibility corresponds to a decrease in KE in the superconducting state which can, in the case considered, more than compensate for the intrinsic increase that accompanies the formation of Cooper pairs and, consequently, the OS rises with a larger slope in the superconducting state than in the normal state just above $`T_c`$. At $`T_c`$ there is no low energy gaping of the spin susceptibility and, therefore, the mechanism for KE reduction just described is not operative. In this sense our model does not describe KE driven superconductivity.
## Acknowledgment
Research supported by the Natural Sciences and Engineering Research Council of Canada (NSERC) and by the Canadian Institute for Advanced Research (CIAR).
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# Hadronic high-energy gamma-ray emission from the microquasar LS I +61 303
## 1 Introduction
LS I $`+\mathrm{61\hspace{0.17em}303}`$ is a Be/X-ray binary that presents unusually strong and variable radio emission (Gregory & Taylor 1978). The X-ray emission is weaker than in other objects of the same class (e.g. Greiner & Rau 2001) and shows a modulation with the radio period (Paredes et al. 1997). The most recent determination of the orbital parameters (Casares et al. 2005) indicates that the eccentricity of the system is $`0.72\pm 0.15`$ and that the orbital inclination is $`30^{}\pm 20^{}`$. The best determination of the orbital period ($`P=26.4960\pm 0.0028`$) comes from radio data (Gregory 2002). The primary star is a B0 V with a dense equatorial wind. Its distance is $`2`$ kpc. The X-ray/radio ourtbursts are triggered 2.5-4 days after the periastron passage of the compact object, usually thought to be a neutron star. These outbursts can last until well beyond the apastron passage.
Recently, Massi et al. (2001) have detected the existence of relativistic radio jets in LS I $`+\mathrm{61\hspace{0.17em}303}`$, which makes of it a member of the microquasar class. Microquasars are thought to be potential gamma-ray sources (Paredes et al. 2000, Kaufman Bernadó et al. 2002, Bosch-Ramon et al. 2005a) and, in fact, LS I $`+\mathrm{61\hspace{0.17em}303}`$ has long been associated with a gamma-ray source. First with the COS-B source CG135+01, and later on with 3EG J0241+6103 (Gregory & Taylor 1978, Kniffen et al 1997). The gamma-ray emission is clearly variable (Tavani et al. 1998) and has been recently shown that the peak of the gamma-ray lightcurve is consistent with the periastron passage (Massi 2004), contrary to what happens with the radio/X-ray emission, which peaks after the passage.
The matter content of microquasars jets is unknown, although in the case of SS 433 iron X-ray line observations have proved the presence of ions in the jets (Kotani et al. 1994, 1996; Migliari et al. 2002). In the present paper we will assume that relativistic protons are part of the content of the observed jets in LS I $`+\mathrm{61\hspace{0.17em}303}`$ and we will develop a simple model for the high-energy gamma-ray production in this system, with specific predictions for Cherenkov telescopes like MAGIC. We emphasize that our model is not opposed, but rather complementary to pure leptonic models as those presented by Bosch-Ramon & Paredes (2004) and Bosch-Ramon et al. (2005a), since the leptonic contribution might dominate at lower gamma-ray energies and after the periastron passage. In the next section we will describe the basic features of the model, and then we will present the calculations and results.
## 2 General picture
A hadronic model for the gamma-ray emission in microquasars with early-type companions has been already developed by Romero et al. (2003). This model, however, is limited to the simple case of a massive star with a spherically symmetric wind and a compact object in a circular orbit. Here we will consider a B-type primary with a wind that forms a circumstellar outflowing disk of density $`\rho _\mathrm{w}(r)=\rho _0(r/R_{})^n`$ (Gregory & Neish, 2002). The continuity equation implies a wind velocity of the type $`v_\mathrm{w}=v_0(r/R_{})^{n2}`$. We will consider that the wind remains mainly near to the equatorial plane, confined in a disk with half-opening angle $`\varphi =15^{}`$, with $`n=3.2`$, $`\rho _0=10^{11}`$ g cm<sup>-3</sup>, and $`v_0=5`$ km s<sup>-1</sup> (Martí and Paredes 1995).
The modeled properties of the system will be expressed in terms of the orbital phase $`\psi `$ ($`\psi =0.23`$ at the periastron passage according to the latest determination by Casares et al. 2005) which is related to the separation between the stars by $`r(\psi )=a(1e^2)/[1e\mathrm{cos}(2\pi (\psi +0.73))]`$, where $`a`$ is the semi-major axis of the orbit and $`e`$ the eccentricity. The wind accretion rate onto the compact object of mass $`M_\mathrm{c}`$ can be estimated as:
$$\dot{M}_\mathrm{c}=\frac{4\pi (GM_\mathrm{c})^2\rho _w(r)}{v_{\mathrm{rel}}^3},$$
(1)
where $`v_{\mathrm{rel}}`$ is the relative velocity between the neutron star (moving in a Keplerian orbit) and the circumstellar wind, assumed to be flowing radially on the equatorial plane.
Following the basic assumption of the jet-disk symbiosis model (Falcke & Biermann 1995) we will assume that the accretion rate is coupled to the kinetic jet power by:
$$Q_\mathrm{j}=q_\mathrm{j}\dot{M}_\mathrm{c}c^2,$$
(2)
where $`q_\mathrm{j}0.1`$ is the coupling constant. Most of this power will consist of cool protons that are ejected with a macroscopic Lorentz factor $`\mathrm{\Gamma }1.25`$ (Massi et al. 2001). Only a small fraction $`q_\mathrm{j}^{\mathrm{rel}}10^3`$ is in the form of relativistic hadrons. The relativistic jet is confined by the pressure of the cold particles ($`P_{\mathrm{cold}}>P_{\mathrm{rel}}`$), which expand laterally at the local sound speed.
The jet axis, $`z`$, will be assumed normal to the orbital plane. The jet will be conical, with a radius $`R_\mathrm{j}(z)=z(R_0/z_0)`$, where $`z_0`$ is the injection point and $`R_0`$ is the initial radius of the jet. We will adopt $`z_0=10^7`$ cm and $`R_0=z_0/10`$ as reasonable values (see Romero et al. 2003 and Bosch-Ramon et al. 2005a, who deals with similar jets for additional details). The relativistic proton spectrum will be a power law $`N_p^{}(E_p^{})=K_pE_{}^{}{}_{p}{}^{\alpha }`$, valid for $`E_{p}^{}{}_{}{}^{\mathrm{min}}E_p^{}E_{p}^{}{}_{}{}^{\mathrm{max}}`$ (in the jet frame). The corresponding relativistic proton flux will be $`J_p^{}(E_p^{})=(c/4\pi )N_p^{}(E_p^{})`$. Since the jet expands in a conical way, the proton flux evolves with $`z`$ as:
$$J_p^{}(E_p^{})=\frac{c}{4\pi }K_0\left(\frac{z_0}{z}\right)^2E_{p}^{}{}_{}{}^{\alpha },$$
(3)
where it is implicit the assumption of the conservation of the number of particles (see Ghisellini et al. 1985), and a prime refers to the jet frame. Using relativistic invariants, it can be shown that the proton flux, in the lab (observer) frame, becomes (e.g. Purmohammad & Samimi 2001):
$$J_p(E_p,\theta )=\frac{cK_0}{4\pi }\left(\frac{z_0}{z}\right)^2\frac{\mathrm{\Gamma }^{\alpha +1}\left(E_p\beta _\mathrm{b}\sqrt{E_p^2m_p^2c^4}\mathrm{cos}\theta \right)^\alpha }{\left[\mathrm{sin}^2\theta +\mathrm{\Gamma }^2\left(\mathrm{cos}\theta \frac{\beta _\mathrm{b}E_p}{\sqrt{E_p^2m_p^2c^4}}\right)^2\right]^{1/2}}.$$
(4)
In this expression, $`\mathrm{\Gamma }`$ is the jet Lorentz factor, $`\theta `$ is the angle subtended by the proton velocity direction (which will be roughly the same as that of the emerging photon) and the jet axis (notice that then $`\theta \theta _{\mathrm{obs}}`$), and $`\beta _\mathrm{b}`$ is the bulk velocity in units of $`c`$. We will make all calculations in the lab frame, where the cross sections for $`pp`$ interactions have suitable parametrizations.
The number density $`n_{0}^{}{}_{}{}^{}`$ of particles flowing in the jet at $`R_0`$, and the normalization constant $`K_0`$ can be determined as in Romero et al. (2003). In the numerical calculations of the next section we have considered $`E_{}^{}{}_{p}{}^{\mathrm{max}}=100`$ TeV, $`E_{}^{}{}_{p}{}^{\mathrm{min}}=1`$ GeV, $`\mathrm{\Gamma }=1.25`$, and, $`\alpha =2.2`$ (see the list of the assumed parameters in Table 1). The assumed maximum energy is consistent with the jet size and shock acceleration with an efficiency $`0.010.1`$.
The matter from the wind can penetrate the jet from the side, diffusing into it as long as the particle gyro-radius is smaller than the radius of the jet. This imposes a constraint onto the value of the magnetic field in the jet: $`B_{\mathrm{jet}}E_k/(eR_0)`$, where $`E_k=m_pv_{\mathrm{rel}}^2/2`$. For the periastron passage ($`E_k`$ maximum) results $`B_{\mathrm{jet}}\mathrm{2.8\hspace{0.17em}10}^6`$ G, which is surely satisfied. However, some effects, like shock formation on the boundary layers, could prevent some particles from entering into the jet. Given our ignorance of the microphysics involved, we adopt a parameter $`f_\mathrm{p}`$ that takes into account particle rejection from the boundary in a phenomenological way. In a conservative approach, we will adopt $`f_\mathrm{p}0.1`$ .
Some of the particles entering the jet flow would be immediately accelerated to the jet velocity (by Coloumb interactions or wave-particle interactions). As a consequence, the jet should be slowed down during its motion through the equatorial wind. However, it is a fact that the jet survives this interaction since it is seen at radio wavelengths far beyond the wind region, up to distances of $`400`$ AU (Massi et al. 2001, 2004). Since the bulk velocity seems not to be very high (Massi et al. 2001) and hence its change does not affect seriously the calculations of the gamma-ray emissivity, we will neglect, in what follows, the effects of a macroscopic deceleration. The reader interested in the case of the hadronic gamma-ray emission of a jet slowed down to rest by the effects of the wind and the resulting standing shock wave as the major source of radiation is referred to the recent treatment presented by Romero & Orellana (2005).
In Figure 1 we show a sketch of the general situation and in Figure 2 we show the orbit of the system and the corresponding phases.
## 3 Gamma-ray emission
Relativistic protons in the jet will interact with target protons in the wind through the reaction channel $`p+pp+p+\xi _{\pi ^0}\pi ^0+\xi _{\pi ^\pm }(\pi ^++\pi ^{})`$, where $`\xi _\pi `$ is the corresponding multiplicity. Then pion decay chains will lead to gamma-ray and neutrino emission. The differential gamma-ray emissivity from $`\pi ^0`$-decays can be expressed as (e.g. Aharonian & Atoyan 1996):
$$q_\gamma (E_\gamma ,\theta )=4\pi \sigma _{pp}(E_p)\frac{2Z_{p\pi ^0}^{(\alpha )}}{\alpha }J_p(E_\gamma ,\theta )\eta _\mathrm{A},$$
(5)
where $`Z_{p\pi ^0}^{(\alpha )}`$ is the so-called spectrum-weighted moment of the inclusive cross-section (see, for instance, Gaisser 1990). $`J_p(E_\gamma )`$ is the proton flux distribution (4) evaluated at $`E=E_\gamma `$. The cross section $`\sigma _{pp}(E_p)`$ for inelastic $`pp`$ interactions at energy $`E_p6\xi _{\pi ^0}E_\gamma /K`$, where $`K0.5`$ is the inelasticity coefficient and $`\xi _{\pi ^0}=1.1(E_p/\mathrm{GeV})^{1/4}`$, can be represented for $`E_p1`$ GeV by
$$\sigma _{pp}(E_p)30\times [0.95+0.06\mathrm{log}(E_p/\mathrm{GeV})](\mathrm{mb}).$$
Finally, the parameter $`\eta _\mathrm{A}`$ takes into account the contribution from different nuclei in the wind and in the jet (for standard composition of cosmic rays and interstellar medium $`\eta _\mathrm{A}1.4`$).
The spectral energy distribution is:
$$L_\gamma (E_\gamma ,\theta )=E_\gamma ^2_Vn(\stackrel{}{r^{}})q_\gamma (E_\gamma ,\theta )d^3\stackrel{}{r^{}},$$
(6)
where $`V`$ is the interaction volume between the jet and the circumstellar disk. The particle density of the wind that penetrates the jet is $`n(r)f_\mathrm{p}\rho _w(r)/m_p`$.
In our calculations, we adopt a viewing angle of $`\theta =30^{}`$ in accordance with the average value given by Casares et al. (2005). In Figure 3 we show a 3-D plot that shows the evolution of the gamma-ray spectral energy distribution as a function of the orbital phase. Other two plots in this figure show cuts at both the periastron and apastron, and the luminosity evolution with the orbital phase at 100 GeV. In both cases we show the unabsorbed (dashed lines) and the absorbed (continuum lines) curves. This absorption is discussed in the next section.
At the periastron passage the unattenuated luminosity is $`10^{33}`$ erg s<sup>-1</sup>. We can make a simple order-of-magnitude estimate of this value. The accretion rate at the periastron is $`3\times 10^{17}`$ g s<sup>-1</sup>. This means that the total power in relativistic protons should be $`Q_\mathrm{j}^{\mathrm{rel}}=10^3\dot{M}_\mathrm{c}c^22.8\times 10^{35}`$ erg s<sup>-1</sup>. The density of the stellar wind at the injection point of the jet is $`n4\times 10^{11}`$ cm<sup>-2</sup> and the cross section for protons of $`E_p1`$ TeV, $`\sigma _{pp}34`$ mb. Hence, the mean free path of the protons results $`\lambda _{pp}8.3\times 10^{13}`$ cm. The thickness of the region of the disk traversed by the jet is $`\mathrm{\Delta }zr_{\mathrm{perias}}\mathrm{tan}15^{}4.4\times 10^{11}`$ cm. Consequently, we can approximate the gamma-ray luminosity by:
$$L_\gamma =2f_\pi Q_\mathrm{j}^{\mathrm{rel}}\left(1e^{\mathrm{\Delta }z/\lambda _{pp}}\right),$$
(7)
where $`f_\pi 0.2`$ is the fraction of the energy of the leading proton that goes into neutral pions and hence into gamma-rays. With a simple substitution into Eq. (7) we get $`L_\gamma 6.6\times 10^{32}`$ erg s<sup>-1</sup>, in good agreement with the detailed numerical calculations presented in Fig. 3.
## 4 Opacity
The optical depth for a photon with energy $`E_\gamma `$, which in this case depends upon the direction observed, can be estimated as
$$\tau (\rho ,E_\gamma )=_{E_{\mathrm{min}}(E_\gamma )}^{\mathrm{}}_\rho ^{\mathrm{}}n_{\mathrm{ph}}(E_{\mathrm{ph}},\rho ^{})\sigma _{e^{}e^+}(E_{\mathrm{ph}},E_\gamma )𝑑\rho ^{}𝑑E_{\mathrm{ph}},$$
(8)
where $`E_{\mathrm{ph}}`$ is the energy of the ambient photons, $`n_{\mathrm{ph}}(E_{\mathrm{ph}},\rho )`$ is their density at a distance $`\rho `$ from the neutron star, and $`\sigma _{e^{}e^+}(E_{\mathrm{ph}},E_\gamma )`$ is the photon-photon pair creation cross section given by:
$$\sigma _{e^+e^{}}(E_{\mathrm{ph}},E_\gamma )=\frac{\pi r_0^2}{2}(1\xi ^2)\left[2\xi (\xi ^22)+(3\xi ^4)\mathrm{ln}\left(\frac{1+\xi }{1\xi }\right)\right],$$
(9)
where $`r_0`$ is the classical radius of the electron and
$$\xi =\left[1\frac{(m_ec^2)^2}{E_{\mathrm{ph}}E_\gamma }\right]^{1/2}.$$
(10)
In Eq. (8), $`E_{\mathrm{min}}`$ is the threshold energy for pair creation in the ambient photon field. This field can be considered as formed by two components, one from the Be star and the other from the hot accreting matter impacting onto the neutron star: $`n_{\mathrm{ph}}=n_{\mathrm{ph},1}+n_{\mathrm{ph},2}`$. Here,
$$n_{\mathrm{ph},1}(E_{\mathrm{ph}},\rho )=\left(\frac{\pi B(E_{\mathrm{ph}})}{hcE_{\mathrm{ph}}}\right)\frac{R_{}^2}{\rho ^2+r^22\rho r\mathrm{sin}\theta },$$
(11)
is the black body emission from the star, with
$$B(E_{\mathrm{ph}})=\frac{2E_{\mathrm{ph}}^3}{(hc)^2(e^{E_{\mathrm{ph}}/kT_{\mathrm{eff}}}1)}$$
(12)
and $`T_{\mathrm{eff}}=22500`$ K (Martí & Paredes 1995). The separation $`r`$ between the stars is again variable with the phase angle $`\psi `$.
The emission from the heated matter can be approximated by a Bremsstrahlung spectrum:
$$n_{\mathrm{ph},2}(E_{\mathrm{ph}},\rho )=\frac{L_XE_{\mathrm{ph}}^2}{4\pi c\rho ^2e^{E_{\mathrm{ph}}/E_{\mathrm{cut}\mathrm{off}}}}\text{ for }E_{\mathrm{ph}}1\text{ keV},$$
(13)
where $`L_X`$ is the total luminosity in hard X-rays and $`E_{\mathrm{cut}\mathrm{off}}100`$ keV. The photon index of the hard X-rays is taken to be within the range published by Greiner & Rau (2001), which was observationally determined. $`L_X`$ is also constrained by observations, being $`L_X10^{34}`$ erg s<sup>-1</sup> (Paredes et al. 1997). Notice that no bump due to a putative accretion disk has been observed at X-rays, so we neglect this contribution.
As an example, Figure 4 shows the dependence of the optical depth $`\tau `$ with the energy of the $`\gamma `$-rays and its variation along the $`z`$ axis for the observer at $`\theta _{\mathrm{obs}}=30^{}`$. From detailed versions of this plot, we find that for photons of $`E_\gamma =100`$ GeV significant absorption occurs mostly between $`\psi =0.1`$ and $`\psi =0.5`$. The optical depth remains well below the unity along the whole orbit for photons of energies $`E_\gamma 30`$ GeV and $`E_\gamma 2`$ TeV.
## 5 Secondary electron-positron pairs and synchrotron emission
Secondary pairs are produced by the decays of charged pions and muons, as well as by photon-photon interactions. The main reactions that lead to charged pions are:
$`p+p`$ $``$ $`p+p+\xi _{\pi ^0}\pi ^0+\xi _{\pi ^\pm }(\pi ^++\pi ^{})`$ (14)
$`p+p`$ $``$ $`p+n+\pi ^++X`$ (15)
$`p+p`$ $``$ $`2n+2\pi ^++X`$ (16)
where $`n`$ is a neutron, $`X`$ stands for anything (neutral) else, and the charged pion multiplicity is $`\xi _{\pi ^\pm }2(E_p/\mathrm{GeV})^{1/4}`$. The neutrons have a proper lifetime of $`886\pm 1`$ s and since they move at ultrarelativistic speed can escape from the source, decaying at considerable distances (Eichler & Wiita 1978). On the contrary, pions decay into the jet trough $`\pi ^\pm \mu ^\pm +\nu `$ and $`\mu ^\pm e^\pm +\nu +\overline{\nu }`$. For an injection proton spectrum given by Eq. (3) with $`\alpha =2.2`$, we have that the pion spectrum (in the jet’s system) will be a power-law $`J_{\pi ^\pm }^{}(E_{\pi ^\pm }^{})=K_{\pi ^\pm }E_{}^{}{}_{\pi ^\pm }{}^{\alpha _\pi }`$, with $`\alpha _\pi 2.3`$. The electron-positron distribution mimics this power law (Ginzburg & Syrovatskii 1964, Dermer 1986):
$$J_{}^{}{}_{e^\pm }{}^{}(E_{}^{}{}_{e^\pm }{}^{})=K_{\pi e^\pm }E_{}^{}{}_{e^\pm }{}^{\alpha _\pm },$$
(17)
with
$$K_{\pi e^\pm }=\left(\frac{m_\mu }{m_e}\right)^{\alpha _\pm 1}\frac{2(\alpha _\pm +5)}{\alpha _\pm (\alpha _\pm +2)(\alpha _\pm +3)}K_{\pi ^\pm },$$
(18)
and $`\alpha _\pm =\alpha _\pi `$.
The energy density of pion-generated pairs along the jet at the periastron passage can be calculated as:
$$w_{\pi e^\pm }=(4\pi /c)E_{}^{}{}_{e^\pm }{}^{}J_{}^{}{}_{e^\pm }{}^{}(E_{}^{}{}_{e^\pm }{}^{})𝑑E_{}^{}{}_{e^\pm }{}^{},$$
(19)
where $`J_{}^{}{}_{e^\pm }{}^{}(E_{}^{}{}_{e^\pm }{}^{})`$ takes into account all the contributions from $`z_0`$ to $`z_{\mathrm{max}}`$. Integrating we get $`w_{\pi e^\pm }3\times 10^9`$ erg cm<sup>-3</sup>.
We can compare the energy density of pairs from the charged pion decays with that of the pairs produced by direct gamma-ray absorption. The total luminosity of these pairs is:
$$L_{e^\pm }=L_\gamma ^0(1e^\tau ).$$
(20)
Then, using the opacity calculated in the previous section, the pair energy density results
$$w_{\gamma \gamma e^\pm }\frac{L_{e^\pm }}{4\pi R_0^2c}.$$
(21)
At the periastron passage, we get $`w_{\gamma \gamma e^\pm }3.7\times 10^9`$ erg cm<sup>-3</sup>. Hence, the pair injection from the photon-photon annihilation is similar to that of pion decay. In what follows we will evaluate the spectrum of these particles using the approximation derived by Aharonian et al. (1983), which is in excellent agreement with the more detailed calculations (exact to 2nd order QED) presented by B$`\ddot{\mathrm{o}}`$ttcher & Schlickeiser (1997).
The differential pair injection rate is given by (B$`\ddot{\mathrm{o}}`$ttcher & Schlickeiser 1997):
$`\dot{n}_{e^\pm }(\gamma )`$ $`=`$ $`{\displaystyle \frac{3}{32}}c\sigma __\mathrm{T}{\displaystyle _\gamma ^{\mathrm{}}}dϵ_\gamma {\displaystyle \frac{N_\gamma (ϵ_\gamma )}{ϵ_\gamma ^3}}{\displaystyle _{\frac{ϵ_\gamma }{4\gamma (ϵ_\gamma \gamma )}}^{\mathrm{}}}dϵ_{\mathrm{ph}}{\displaystyle \frac{n_{\mathrm{ph}}(ϵ_{\mathrm{ph}})}{ϵ_{\mathrm{ph}}^2}}\times `$ (22)
$`[{\displaystyle \frac{4ϵ_\gamma ^2}{\gamma (ϵ_\gamma \gamma )}}\mathrm{ln}\left({\displaystyle \frac{4ϵ_{\mathrm{ph}}\gamma (ϵ_\gamma \gamma )}{ϵ_\gamma }}\right)8ϵ_\gamma ϵ_{\mathrm{ph}}+{\displaystyle \frac{2(2ϵ_\gamma ϵ_{\mathrm{ph}}1)ϵ_\gamma ^2}{\gamma (ϵ_\gamma \gamma )}}`$
$`(1{\displaystyle \frac{1}{ϵ_\gamma ϵ_{\mathrm{ph}}}}){\displaystyle \frac{ϵ_\gamma ^4}{\gamma ^2(ϵ_\gamma \gamma )^2}}],`$
where $`\gamma =E_{e^\pm }/m_\mathrm{e}c^2`$, $`ϵ_\gamma =E_\gamma /m_\mathrm{e}c^2`$, and $`ϵ_{\mathrm{ph}}=E_{\mathrm{ph}}/m_\mathrm{e}c^2`$. A numerical integration yields a pair spectrum that can be well fitted by a power law $`N_{e^\pm }E_{e^\pm }^{1.9}`$. The proportionality constant $`K_{\gamma \gamma e^\pm }`$ can be obtained from the absorbed gamma-ray luminosity.
The presence of a magnetic field in the jet will imply that all these secondary pairs will produce synchrotron emission. Following Bosch-Ramon et al. (2005b) we assume that the magnetic field is entangled to cold protons in such a way it has random directions and hence the synchrotron emission is isotropic in the jet’s frame. To calculate the synchrotron luminosity we estimate the specific emission ($`j_ϵ(z)`$) and absorption ($`k_ϵ(z)`$) coefficients from the secondary particle distribution (see Pacholczyk 1970 for the detailed formulae), in such a way that:
$$\frac{dL_ϵ(z)}{dz}=2\pi R_\mathrm{j}\frac{j_ϵ(z)}{k_ϵ(z)}\times [1\mathrm{exp}(l_\mathrm{j}k_ϵ(z))],$$
(23)
where to simplify the notation we are not using now primes to indicate that the calculation is in the jet’s frame. In Eq. (23) $`l_\mathrm{j}R_\mathrm{j}`$ is the typical size of the synchrotron emitting plasma and $`ϵ`$ is the photon energy in units of $`m_\mathrm{e}c^2`$. Integrating over the jet length we get the spectral energy distribution as:
$$L_{\mathrm{syn}}^{\mathrm{obs}}=ϵ_{z_0}^{z_{\mathrm{max}}}\delta ^2\frac{dL_ϵ}{dz}𝑑z,$$
(24)
where $`\delta `$ is the Doppler boosting factor defined as:
$$\delta =\frac{1}{\mathrm{\Gamma }(1\beta _\mathrm{b}\mathrm{cos}\theta _{\mathrm{obs}})}.$$
(25)
To calculate the specific emission $`j_ϵ(z)`$ we adopt different values of the magnetic field at $`z_0`$: $`B_0=1`$, 10, and 100 Gauss (Bosch-Ramon & Paredes 2004). In Figure 5 we show the spectral energy distribution of the synchrotron radiation of all secondary pairs for the 3 different values of $`B_0`$. The radio emission is quite negligible in comparison to the observed values, which at the minimum imply a luminosity of $`10^{31}`$ erg s<sup>-1</sup> (e.g. Ribó et al. 2005).
## 6 Discussion
The predicted gamma-ray luminosity is clearly at its maximum during the periastron passage, when the neutron star travels through the densest parts of the wind. This is in accordance with the fact noticed by Massi (2004) that the peaks of the EGRET flux are coincident with the periastron and not with the radio maxima. The radio outbursts are the result of particle injection in the jet that occurs after some relaxation time from the periastron passage, when the accretion rate is increased (Paredes et al. 1991). Any purely leptonic model for the gamma-ray emission would have to explain why the radio and gamma-ray peaks are not observed in similar orbital phases.
Other specific feature of the gamma-ray emission predicted by our model is the presence of a local, secondary maximum at $`\psi 0.65`$ when the accretion rate, given by (1), has also a local maximum due to the fact that the wind velocity is roughly parallel to the neutron star orbital velocity, hence reducing $`v_{\mathrm{rel}}`$ and increasing $`\dot{M}_\mathrm{c}`$, as noticed by Martí & Paredes (1995).
The effects of the opacity of the ambient photon fields to gamma-ray propagation produces a “valley” in the spectral energy distribution, between a few tens of GeV and a few TeV, with a local minimum at around 100 GeV, during the periastron passage. The predicted luminosity is within the detection possibilities of an instrument like MAGIC, which, integrating over several periastron passages, could build up a SED which can be compared with that presented in Fig. 3. Upper limits obtained with the Whipple telescope (Hall et al. 2003, Fegan et al. 2005) are indicated in the figure. The source is too weak for the sensitivity of this instrument according to our model.
## 7 Concluding remarks
We have presented a hadronic model for the high-energy gamma-ray production in the microquasar LS I $`+\mathrm{61\hspace{0.17em}303}`$. The model is based on the interaction of a mildly relativistic jet with a small content of relativistic hadrons with the dense equatorial disk of the companion B0 V star. Gamma-rays are the result of the decay of neutral pions produced by $`pp`$ collisions. Charged pion decay will lead to neutrino production, that will be discussed elsewhere. The model takes into account the opacity of the ambient photon fields to the propagation of the gamma-rays. The predictions include a peak of gamma-ray flux in the periastron passage, with a secondary maximum at phase $`\psi 0.65`$. The spectral energy distribution presents a minimum around 100 GeV due to absorption. The spectral features should be detectable by an instrument like MAGIC through exposures $`50`$ hr, integrated along different periastron passages.
## Acknowledgments
We thank J.M. Paredes and V. Bosch-Ramon for careful readings of the manuscript and comments. The latter gave us useful support on calculations for the secondary emission. We also thank constructive suggestions by an anonymous referee. This work has been supported by the Argentinian agencies CONICET and ANPCyT (PICT 03-13291). HRC thanks support from FUNCAP and CNPq (Brazil).
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# Dynamics of fluctuation of the top location of a sandpile
## 1 Introduction
Dense systems of granular materials exhibit solid-like and fluid-like behaviors , and many researches are devoted on them. For the static phase, the spatial distribution of stress on the bottom of a sandpile was measured in experiments, and it is known that its functional form depends on the history of the formation of the sandpile . For the fluid phase, steady flows and avalanches are investigated from several perspectives . Fluidized phase appears in a localized layer near the surface, and the particle velocity under the layer obeys an exponential function of the depth from the surface . It is known for granular flows in a pipe that the temporal power spectra of particle density obey power laws .
In systems of a sandpile, feeding particles at small feed rate, the surface of a sandpile is kept in solid states except when avalanches occur intermittently, and continuous flows appear as the feed rate increases. However, even in the case that the feed rate is rather large, it is infrequent that the whole surface of a sandpile is kept in the fluid phase, and the states of the surface change temporally or spatially . The time evolution of a sandpile is caused by complicated interactions between avalanches and the shape of the surface.
A sandpile typically becomes mountain shaped with a top as it grows . Because particles on the surface run down from the top to the foot of a sandpile, the top plays the role of a singular point in an average flow of particles. In the formation process of a sandpile, the top location moves with time and determines the global flow on the surface. In this paper, we focus on the fluctuation of the top location to characterize the long time evolution of a sandpile. We carry out two-dimensional numerical simulations and measure the power spectrum of the top location. We find that the power spectrum obeys a power function, and that the exponent depends on the feeding rate.
The organization of the remainder of this paper is as follows. In §2, we explain the method of our numerical simulations. In §3, we investigate the power spectrum of the top location and relations between its exponent and the feed rate of particles. In §4, we discuss the origin of this power law. Finally, we draw conclusions in §5.
## 2 Method Of Simulations
In order to simulate polydispersed circular particles, we adopt a two-dimensional Discrete Elements Method (DEM) . In DEM, we assume that particles are cohesionless, and the mass density of a particle per unit area, $`\rho _0`$, is constant. The linear spring model is adopted to describe the repulsion of two particles in contact, and the viscous forces and Coulomb slip are assumed to act on them. We assume that the $`i`$th particle has the radius $`r_i`$, the mass $`m_i=\pi r_i^2\rho _0`$ and the momentum of inertia $`I_i=m_ir_i^2/2`$. The center of mass and the angular velocity of the $`i`$th particle represent $`𝐱_i`$ and $`\omega _i`$, respectively. The equations of motion of the $`i`$th particle are given by
$`\begin{array}{ccc}\hfill m_i\ddot{𝐱}_i& =& _j\left(F_n^{ij}𝐧_{ij}+F_t^{ij}𝐭_{ij}\right)+m_i𝐠,\hfill \\ \hfill I_i\dot{\omega }_i& =& r_i_jF_t^{ij},\hfill \end{array}`$ (3)
where $`𝐧_{ij}`$ and $`𝐭_{ij}`$ represent the normal and the tangential unit vectors at the contact point between the $`j`$th particle and the $`i`$th particle, and $`𝐠`$ is the acceleration of gravity, as depicted in Fig. 1. The normal contact force $`F_n^{ij}`$ is defined by
$$F_n^{ij}=\stackrel{~}{F}_n^{ij}\mathrm{\Theta }(\stackrel{~}{F}_n^{ij}),$$
(4)
and
$$\stackrel{~}{F}_n^{ij}=k_nm_r^{ij}\left(r_i+r_j|𝐱_i𝐱_j|\right)\eta m_r^{ij}𝐧_{ij}\left(\dot{𝐱}_i\dot{𝐱}_j\right),$$
(5)
where $`\mathrm{\Theta }(x)`$ is the Heaviside function, and $`m_r^{ij}`$ is the reduced mass $`m_r^{ij}\frac{m_im_j}{m_i+m_j}`$. $`\mathrm{\Theta }(x)`$ is introduced so that the contact force is repulsive. The tangential contact force $`F_t^{ij}`$ is defined by
$$F_t^{ij}=k_tm_r^{ij}u_t^{ij}.$$
(6)
$`u_t^{ij}`$ is obtained by integrating the equation
$$\dot{u}_t^{ij}=\left((\dot{𝐱}_i\dot{𝐱}_j)𝐭_{ij}+r_i\omega _i+r_j\omega _j\right)\mathrm{\Theta }\left(\mu |F_n^{ij}||F_t^{ij}|\right),$$
(7)
where $`u_t^{ij}`$ is zero when the particles are not in contact (i. e. for $`|𝐱_j𝐱_i|>r_j+r_i`$). Here, Coulomb frictional coefficient $`\mu `$ is assumed to be $`0.5`$. We express the maximum diameter and the maximum weight of a particle as $`d`$ and $`m`$. The distribution function of the diameters is uniform in the range between $`0.8d`$ and $`d`$. We assume that the spring constants are $`k_n=1.0\times 10^4mg/d`$ in the normal direction, and $`k_t=2.0\times 10^3mg/d`$ in the tangential direction, and that the viscosity is $`\eta =1.0\times 10^2\sqrt{g/d}`$. The coefficient of restitution in our model is about $`0.2`$ for a head-on particles collision. We adopt the second-order Adams-Bashforth method for time-integration with the time interval $`\mathrm{\Delta }t=1.0\times 10^3\sqrt{d/g}`$.
We assume that a particle is in a sandpile if the particle contacts other particles, and the top location of the sandpile is defined as the center of mass of the highest particle in the sandpile (Fig. 2). The floor under a sandpile is a horizontal array of $`80`$ fixed particles with diameter $`d`$. We introduce the $`x`$ coordinate along the floor and define the origin at the center of the floor.
To investigate the fluctuation of the top in the formation process of a sandpile, we drop particles with the time interval $`T`$ to it. The particles are released at a position just above the center of the floor and its height is $`H`$ from the top location as shown in Fig. 2. We first make a sandpile grow until it covers the floor and use this sandpile as an initial state. Because particles run off the edges of the floor with finite length, the size of a sandpile is maintained almost constant.
After the time series of the top location $`x_{top}(t)`$ reaches a statistical stationary state, we calculate its power spectrum $`S(f)`$ with respect to frequency $`f`$. To calculate the power spectrum $`S(f)`$ from the time series, we divide it into $`M`$ time series with a time interval $`T^{(s)}`$, the $`m`$th power spectrum $`S_m(f_j)`$ is defined by
$`S_m(f_j){\displaystyle \frac{1}{N}}|{\displaystyle \underset{n=1}{\overset{N}{}}}x_{top,n}^{(m)}e^{\frac{2\pi in}{N}j}|^2`$ (8)
where $`x_{top,n}^{(m)}x_{top}((m+\frac{n}{N})T^{(s)})`$ and $`f_j\frac{j}{T^{(s)}}`$.
We introduce the power spectrum $`S(f)`$ as the average of the $`M`$ power spectra,
$`S(f_j){\displaystyle \frac{1}{M}}{\displaystyle \underset{m=0}{\overset{M1}{}}}S_m(f_j).`$ (9)
## 3 Results
We measure $`x_{top}`$ for various values of the time interval $`T`$ and the height $`H`$. We change $`H`$ in the range $`20dH110d`$. If $`H`$ is sufficiently large beyond this range, the impact of a dropped particle is large and collapses the top shape of a sandpile into a caldera . Figures 3(a) and 3(b) shows the time series of $`x_{top}(t)`$ obtained from simulations with $`T=2\sqrt{d/g}`$ (a) and $`T=80\sqrt{d/g}`$ (b). The top fluctuates frequently in the case of small $`T`$, on the other hand, in the case of large $`T`$, the top almost stays for long time in comparison with $`T`$ because the motion of particles induced by the impact of a fed particle ceases before the next particle is dropped. Figure 4 is the power spectra of the time series, $`S(f)`$, which is calculated using eqs. (8) and (9) with $`T^{(s)}=N=10000`$ and $`M=10`$. We find that $`S(f)`$ behaves as a power-law, and its exponent depends on the value of $`T`$. At $`T=2\sqrt{d/g}`$, $`S(f)`$ approximately obeys $`1/f`$ law. At $`T=80\sqrt{d/g}`$, the exponents of $`S(f)`$ are smaller than $`1`$.
We investigate the dependence of the exponent of $`S(f)`$, $`\alpha `$, on $`T`$ and $`H`$. For the range $`5/T^{(s)}f1/(2T)`$, we calculate $`\alpha `$ from the double logarithmic plot of $`S(f)`$ in the least square method. $`\alpha `$ is insensitive to $`H`$ as shown in Fig. 5(a). In contrast, Fig. 5(b) indicates that $`\alpha `$ strongly depends on $`T`$ for small $`T`$ and approaches $`1`$ as $`T`$ decreases. As $`\alpha `$ approaches $`1`$, the power spectrum is approximated as a power function with a high degree of accuracy as indicated with the error bars. In the range $`10\sqrt{d/g}<T<60\sqrt{d/g}`$, $`\alpha `$ is approximately a constant $`1.43\pm 0.03`$, although $`\alpha `$ decreases as $`T`$ increases beyond $`60\sqrt{d/g}`$. In the case of large $`T`$, the error bars in Figs. 5 are large because the range of frequency used for fitting is small. In the region of higher frequency than $`1/T`$, the power function with the same exponent is not best fit with $`S(f)`$. If we fit $`S(f)`$ with a power function in this range of high frequency, its exponent changes from that indicated in Figs. 5, as shown with the data of $`T=80\sqrt{d/g}`$ in Fig. 4.
We mainly focus on the case of small $`T`$ because we are, in particular, interested in the case that $`\alpha `$ is close to $`1`$. The displacement of the top location $`x_{top}`$ is caused by avalanches, and avalanches occur on either slope at almost all times in this case. Although the instantaneous magnitude of avalanches is characterized by the kinetic energy of the particles, $`S(f)`$ and the power spectrum of the kinetic energy differ in functional form as mentioned below. Eliminating the narrow region $`dxd`$ in the center of a sandpile, we divide the sandpile into the left part and the right part with respect to $`x=0`$. We measure the kinetic energies of the left and right parts, $`K_l(t)`$ and $`K_r(t)`$, respectively. Using the same definition in eqs. (8) and (9), the power spectra of the time series of $`K_l(t)`$ and $`K_r(t)`$ are calculated with $`T^{(s)}=1.0\times 10^4\sqrt{d/g}`$, $`N=1.0\times 10^5`$ and $`M=10`$. For $`T=2\sqrt{d/g}`$ and $`H=20d`$, the power spectra of both $`K_l(t)`$ and $`x_{top}(t)`$ are shown in Fig. 6. The power spectrum of $`K_r(t)`$ is similar to that of $`K_l(t)`$. Because the power spectra of $`K_l(t)`$ and $`K_r(t)`$ are Lorentzian-like, avalanches seem to occur at random.
From the results of numerical simulations, it is found to be rare that avalanches occur simultaneously on the left and right slopes of a sandpile. We refer the states that avalanches occur on the left and right slopes as left mode and right mode, respectively. To investigate switchings between the both mode, we define the binarized time series,
$$K(t)=\{\begin{array}{cc}+1,\hfill & \text{for }K_l(t)<K_r(t),\hfill \\ 1,\hfill & \text{for }K_l(t)K_r(t).\hfill \end{array}$$
(10)
The sign of $`K(t)`$ represents the side on which avalanches occur mainly at time $`t`$. The switchings are well defined by $`K(t)`$ in the case that $`T`$ is small. However, as $`T`$ increases, it is difficult to define the switchings because avalanches occur at intervals, and the time intervals between avalanches are comparable to the time scale of switchings. We find that the time series of $`K(t)`$ is similar to $`x_{top}(t)`$ for small $`T`$. The power spectrum of the time series $`K(t)`$ is shown in Fig. 7 for $`T=2\sqrt{d/g}`$ and $`H=20d`$. The power spectrum of $`K(t)`$ is approximated as a power law with the exponent of $`1`$ in the long time scale. The exponent is approximately the same as that of the top location. Investigating the conditional probability of $`K(t)=1`$ for a given $`x_{top}`$, this probability increases with the value of $`x_{top}`$. Therefore, in each mode, the top location $`x_{top}(t)`$ is mainly in the opposite side on which avalanches occur. Thus the fluctuation of $`x_{top}`$ corresponds to the switching between the two modes, but not to the fluctuation of the magnitude of avalanches.
## 4 Discussion
For the binarized time series such as $`K(t)`$ defined by eq. (10), it is known from an analytical theory that its spectrum is expressed as a power function if the waiting time has a power-law distribution and each interval is independent . Here, the waiting time $`\tau `$ is defined as a time interval between neighboring switchings in the binarized time series. We assume that the probability density of $`\tau `$, $`p(\tau )`$, is the abrupt-cutoff power law,
$`p(\tau )=\{\begin{array}{cc}ct^D,\hfill & \text{for }a<t<b,\hfill \\ 0,\hfill & \text{otherwise},\hfill \end{array}`$ (11)
where the constants $`a`$ and $`b`$ are sufficiently small and large respectively, and $`c`$ is the normalization constant. In the range of $`1/bf1/a`$, the power spectrum of this binarized time series, $`S_b(f)`$, is given approximately by
$`S_b(f)f^{(D+3)}`$ (12)
for $`3<D<1`$. In the case of small $`T`$, $`p(\tau )`$ is expected to be a power function with $`D2`$ because the exponent of the power spectrum of $`K(t)`$ is approximately $`1`$ in our simulations. However there are few intervals with waiting time longer than $`\tau 100\sqrt{d/g}`$ in $`K(t)`$ because small noises chop up long intervals. Therefore applying the median filter of a time width of $`60\sqrt{d/g}`$ to $`K(t)`$, we calculate the distribution of waiting time for the coarse-grained time series, $`p(\tau ^{^{}})`$. We find that $`p(\tau ^{^{}})`$ decays approximately as a power function with $`D2`$ as shown in Fig. 8.
The plateau with $`\alpha 1.4`$ appears in a wide range of $`T`$ as shown in Fig. 5(a), although the switchings can not be well-defined by $`K(t)`$ as $`T`$ increases. It is possible that the dynamics of the top location has some relations with the density fluctuation of flows. In experiments of granular flow in vertical pipes filled with fluid, the exponents $`1`$, $`4/3`$ and $`3/2`$ are reported for the temporal power spectra of density , and the exponents $`1.4`$ and $`3/2`$ appear in traffic flows . We note that these exponents are close to $`1.4`$.
As $`T`$ becomes sufficiently large, we believe that $`\alpha `$ approaches $`2`$. The top location stays at the same place for long time in comparison with $`T`$, and its displacements are caused by impulsive force of avalanches. We infer that the top location moves like as a random walk, and its power spectrum is Lorentzian-like, which decays as $`f^2`$.
## 5 Conclusions
Carrying out 2-D DEM simulations, we have investigated the fluctuation of the top location of a sandpile that is caused by avalanches and piling up particles. We have found that the power spectra of the time series of the top location $`x_{top}(t)`$ behave as power functions in the range of long time scale. The exponent of the power spectrum, $`\alpha `$, depends on the time interval $`T`$ at which particles are fed to the sandpile. $`\alpha `$ is close to $`1`$ for small $`T`$ and decreases through a plateau with $`\alpha 1.4`$ as $`T`$ increases. In the case of $`\alpha 1`$, avalanches occur mainly either on the left or right side slopes, and the states of the sandpile switch intermittently between the left and right modes. The power spectrum of the top location is approximately the same as that of the binarized time series defined from the switchings. In our simulations, the distribution of waiting time of the switchings obeys a power function with the exponent $`D2`$ in this case. The relation between $`D`$ and $`\alpha `$ is consistent with the equation $`\alpha +D=3`$ proposed in the analytical theories .
## Acknowledgment
C. Urabe thanks H. Hayakawa, H. Tomita, S. Takesue and S. Kitsunezaki for fruitful discussion and N. Fuchikami for discussion on the binarized time series . The author thanks H. Nakao for providing information on the paper of S. B. Lowen and M. C. Teich . The author appreciates H. Hayakawa and S. Kitsunezaki for their critical reading. This work is partially supported by Grant-in-Aids for Japan Space Forum and Scientific Research (Grant No. 15540393) of the Ministry of Education, Culture, Sports, Science and Technology, Japan.
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# SALT: a Spectral Adaptive Light curve Template for Type Ia Supernovae
## 1 Introduction
Type Ia supernovae (SNe Ia) are a powerful tool for studying the evolution of the luminosity distance as a function of redshift and for subsequently constraining the cosmological parameters. SNe Ia are indeed very luminous and “standardizable” candles, and have lead to the discovery of the acceleration of the Universe (Riess et al. 1998; Perlmutter et al. 1999).
Although often described as a homogeneous class of objects, SNe Ia exhibit variability in light curve shapes, colors, intrinsic luminosity, and spectral features. Finding correlations among SN Ia observables is motivated by improving the estimation of their intrinsic luminosity on an event-by-event basis, in order to reduce the scatter in luminosity distance estimates. The main correlations observed in photometric measurements are:
* a width-luminosity (or brighter-slower) relation, which expresses the fact that brighter supernovae have a slower decline rate than fainter ones (Pskovskii 1977; Phillips 1993; Riess et al. 1995; Hamuy et al. 1996b; Perlmutter et al. 1997).
* a brighter-bluer relation, which was made explicit in Tripp (1998); Tripp & Branch (1999); Parodi et al. (2000), and assumed to be due to extinction by dust in other works (Riess et al. 1996a, 1998; Perlmutter et al. 1999; Tonry et al. 2003; Knop et al. 2003; Barris et al. 2004).
The case of the brighter-bluer relation is interesting. Even if authors fully agree neither on the origin of the effect nor on the strength of the correlation, most, if not all, recent attempts to build a SN Ia Hubble diagram have made use of color in their distance estimator<sup>1</sup><sup>1</sup>1In Perlmutter et al. (1999), it is checked that $`BV`$ rest-frame colors of nearby and distant type agree on average, and the color measurement is not used event per event.. Riess et al. (1996b) summarizes previous work on the subject and proposes a way to reconcile divergent interpretations of data by taking into account the correlation between light curve shape and color. New methods using color information have also been recently proposed to estimate luminosity distances (see for example Wang et al. 2003 and Wang et al. 2005). We will bring to the debate our own estimations of the correlation strength and its wavelength dependence.
Cosmological measurements using SNe Ia are based on comparing nearby and distant objects. In order to reduce the sources of systematic uncertainties, it is important that all distances are derived using the same procedure, especially when considering large samples such as those being collected by the ESSENCE<sup>2</sup><sup>2</sup>2http://www.ctio.noao.edu/~wsne or SNLS<sup>3</sup><sup>3</sup>3http://cfht.hawaii.edu/SNLS projects. To analyze these data sets, the following constraints have to be taken into account:
* High redshift objects often lack late-time photometric data (or have one of too poor quality) which makes it impracticable to estimate color from late-time data as proposed by Lira et al. (1998).
* Above redshifts $`z0.8`$, rest-frame $`U`$-band measurements have to be used because of the limitation of silicon detectors. Incorporating rest-frame $`U`$-band measurements on the same footing as the commonly used rest-frame $`B`$\- and $`V`$-band data is hence highly desirable.
* Applying cuts on light curve parameters should be prohibited since the resolution on these parameters follows the photometric resolutions and hence degrades with redshift, therefore biasing the event sample. In particular, cuts on colors to eliminate reddened events are to be avoided.
Various techniques have been used to estimates distances. Phillips (1993) and Perlmutter et al. (1997) fit one band at a time, and derive distances from light curve parameters measured in different bands (usually $`B`$ and $`V`$ bands), sometimes relying on late time measurements to measure host galaxy extinction. A recent and refined version of this approach can be found in Wang et al. (2005), and in Wang et al. (2003) it is proposed to fit $`BV`$ as a function of $`B`$. These methods do not make use of rest-frame $`U`$-band measurements. Nugent et al. (2002) discuss in detail the problems associated with SN Ia $`U`$-band photometry. They mention in particular the apparent large intrinsic variations of the UV luminosity among similar supernovae as well as our lack of understanding of the SN Ia UV photometry –principally due the uncertainties which affect the large extinction corrections which have to be applied to the data. Nevertheless, it is important to try and incorporate rest-frame $`U`$-band data to estimate distances.
An extension of Perlmutter et al. (1997) to rest-frame $`U`$-band data was developed in Knop et al. (2003), but at the expense of adding a large distance systematic (and probably statistical) uncertainty. Similarly, the MLCS method (Riess et al. 1995, 1996a) has been extended to include $`U`$-band measurements under the name of MLCS2k2 (Jha 2002) and used in Riess et al. (2004), but also at the expense of a worsened distance resolution (Jha 2002). We will show that our model predicts light curves for any band located between rest-frame $`U`$ and $`R`$ bands and that we are able to get a similar or better distance resolution with both rest-frame ($`B,V`$) and ($`U,B`$) band pairs.
We propose here to parameterize the light curve model with a minimal parameter set: a luminosity parameter, a decline rate parameter and a single color parameter. Our approach will be to build a phenomenological model of the expected SN flux, continuously varying with phase, wavelength, decline rate and color, in order to capture all these features at once. This approach offers several practical advantages which make it easily applicable to high-redshift SNe Ia currently measured in large projects. First, the k-corrections are built into the model and not applied to the data. This allows one to propagate all the uncertainties directly from the measurement errors. More importantly, when needed, we make use of the SN rest-frame $`U`$-band fluxes to estimate the supernova distances.
In Sect. 2 we describe the semi-analytic model used. We then describe, in Sect. 3, how the coefficients of the model are determined by an iterative training based on a set of well-sampled nearby SNe Ia taken from the literature. We also highlight some properties of the resulting model. At this stage, the aim is to model multi-color light curves and not to estimate luminosity distances. The model is tested in Sect. 4 with an independent set of SNe Ia in the Hubble flow. A luminosity distance estimate is then constructed from the fitted parameters of the light curve model. It is used to build Hubble diagrams successively from $`(B,V)`$ and $`(U,B)`$ light curve pairs. In order to assess the precision of this approach, we compare distance estimates of the same events obtained from the $`(U,B)`$ and $`(B,V)`$ band pairs and with other distance estimators.
## 2 The light curve model
### 2.1 Model definitions
As already mentioned, we choose to parameterize light curves (more precisely light curve pairs or light curve triplets when available) using a single luminosity, a single shape parameter and a single color. The choice among possible implementations is largely arbitrary. We choose the following parameters, which enable comparisons to be made with previous works:
* $`𝐟_\mathrm{𝟎}`$ : a global intensity parameter which varies with redshift like the inverse of the luminosity distance squared,
* $`𝐬`$ : a time stretch factor as the decline rate indicator (Perlmutter et al. 1997). In Goldhaber et al. (2001), this parameter is shown to apply to the rising part of the light curve as well. However, while the stretch paradigm describes well the bright part of the $`B`$ light curve, it does poorly at late time. It also fails to capture the shape variations in the other bands. This is why our model uses the stretch parameter as an index rather than the stretch paradigm itself. As described below, by construction our model follows exactly the stretch paradigm in the $`B`$ band.
* $`𝐜`$ = $`(BV)_{max}+0.057`$, where $`(BV)_{max}`$ is measured at B maximum, and $`0.057`$ is the chosen reference color (Vega magnitudes) of a SN Ia.
* $`t_{max}^B`$ : the date of maximum in the rest-frame $`B`$ band
With these definitions, the expected counting rate $`f_{SN}`$ in a given pass-band $`T`$, of a supernova at redshift $`z`$, and at a phase $`p(tt_{max}^B)/(1+z)`$, can be written:
$`f_{SN}(p,z,T)=f_0(1+z){\displaystyle \varphi (p,\lambda ,s,c)\frac{\lambda }{hc}T(\lambda (1+z))𝑑\lambda }`$
$`\varphi (p,\lambda ,s,c)`$ is a model of the SN Ia energy luminosity per unit wavelength. It may vary with the supernova stretch and color. Note that the potential extinction by dust in the host galaxy is not explicit in the equation. Instead, we choose to incorporate it in the model $`\varphi (p,\lambda ,s,c)`$ as discussed below.
Building an average spectral template $`\varphi `$ as a function of phase, wavelength, color and stretch from observations is complicated because of the inhomogeneity and incompleteness of published data. Although some spectral features have been correlated with stretch, we do not have yet a complete knowledge of spectral diversity as a function of phase. In most of the approaches to SN Ia light curve fitting, the limited knowledge of spectral variability impacts on the accuracy of cross-filter k-corrections, defined as expected ratios of fluxes in different bands at the same phase (Kim et al. 1996; Nugent et al. 2002).
In Nugent et al. (2002), it was shown that the variation in k-corrections from one supernova to another depends primarily on the supernova color, and to a lesser extent on spectral features<sup>4</sup><sup>4</sup>4Most of the SN Ia photometric reductions transform instrumental magnitudes into standard magnitudes using color equations, derived from standard stars observations. This assumes that color rather than spectral features dominates cross-filter corrections.. Hence we neglected the variability of those spectral features in the modeling of light curves.
In order to implement stretch dependent light curve shapes and colors we therefore used the following approximation:
$`f_{SN}(p_s,z,T)=f_0(1+z){\displaystyle \varphi (p_s,\lambda )\frac{\lambda }{hc}T(\lambda (1+z))𝑑\lambda }`$
$`\times \mathrm{exp}\left[0.4\mathrm{ln}(10)\times 𝒦(p_s,\lambda _T,s,c)\right]`$ (1)
where $`p_sp/s`$ is a stretch-corrected phase. This functional form defines the light curve model.
In equation 1, $`\varphi `$ no longer depends explicitly on $`s`$ and $`c`$, and $`𝒦(p_s,\lambda ,s,c)`$ is a smooth “correction” function of our four variables. $`\lambda _T`$ is the central wavelength of the filter $`T`$. $`𝒦`$ enables one to implement light curve shape variations that are more complicated than simple dilation of the time scale, along with stretch dependent colors. As described below, $`𝒦`$ varies smoothly with $`\lambda `$; this justifies placing it outside the integral over wavelength. We also considered keeping $`𝒦(p_s,\lambda _T,s,c)`$ inside the integral. It changed the results of the fits by less than 1%, while the computing time was multiplied by a factor of 10.
Equation (1) implements k-corrections for an average SN Ia, conforming to the common practice. The approach proposed here will be easy to adapt when constructing a stretch-dependent spectral template becomes possible.
The k-corrections are usually applied to data. Here, they are incorporated into the model. This offers a few practical advantages: first, light curves can be generated for arbitrary pass-bands (within the spectral coverage of $`\varphi `$ and $`𝒦`$); light curves in the observed pass-bands can be directly fitted to the data. Second, the light curve parameter uncertainties extracted from the fit incorporate all uncertainties propagated from measurement uncertainties; for example, uncertainties introduced in the k-corrections by the possibly poor determination of the date of maximum and/or color are propagated into the parameters uncertainties.
The functions $`\varphi `$ and $`𝒦`$ define the model. Once they are determined, one can fit the supernova photometric data points, measured in a minimum of two pass-bands, to estimate $`f_0`$, $`s`$, $`c`$, and a date of $`B`$ maximum light, which is a nuisance parameter. With only one passband, $`c`$ must be held fixed.
For $`\varphi (p_s,\lambda )`$, we use a template spectrum assembled by P. Nugent (Nugent et al. 2002 and private communication) smoothed along the phase (time) axis, and normalized as a function of phase to the $`B`$-band light curve template “Parab -18” of Goldhaber et al. (2001). Any smooth variation of the template $`\varphi `$ with phase or wavelength is irrelevant at this level, since it may be changed by the function $`𝒦(p_s,\lambda ,s,c)`$. The only important quantities in $`\varphi (p_s,\lambda )`$ are the SN Ia spectral features, intended to be realistic on average.
The empirical correction function $`𝒦`$ is implemented as a sum of two polynomials:
$$𝒦(p_s,\lambda ,s,c)=𝒦_s(p_s,\lambda ,s)+𝒦_c(\lambda ,c)$$
(2)
where we explicitly separate the corrections associated with the parameters $`s`$ and $`c`$ to clarify their interpretation. $`𝒦_s(p_s,\lambda ,s)`$ modifies the shape of light curves and absorbs any stretch–color relation except for the $`(BV)_{max}`$ color. Indeed we want $`c`$ to describe exactly the $`(BV)_{max}`$ color. $`𝒦_c(\lambda ,c)`$ is then a color correction as a function of wavelength and color.
In order to remove all degeneracies among coefficients, the model must fulfill the following constraints:
$`𝒦_s(p_s<35,s,\lambda _B)`$ $`=`$ $`0`$
$`𝒦_s(p_s=0,s,\lambda _V)`$ $`=`$ $`0`$
$`𝒦_c(\lambda ,c=0)=𝒦_c(\lambda _B,c)`$ $`=`$ $`0`$
$`𝒦_c(\lambda _B,c)𝒦_c(\lambda _V,c)`$ $`=`$ $`c`$
where $`\lambda _B`$ and $`\lambda _V`$ refer to the mean wavelengths of the B and V filters. The first constraint ensures that the parameter $`s`$ actually defines the stretch in $`B`$ band, since the template is not modified for $`\lambda =\lambda _B`$. The remaining constraints define $`c`$ as the $`BV`$ color (relative to the spectral template) at maximum $`B`$, for all stretches. Note that other colors may (and in fact will) depend on $`s`$ at fixed $`c`$. $`f_0`$ describes the actual peak flux in $`B`$. It is in no way “corrected” for the brighter-slower relation.
$`𝒦_s`$ is implemented as a polynomial of degree $`D_p`$ in phase, $`D_\lambda `$ in $`\lambda `$ and $`D_s`$ in stretch respectively. Similarly, $`𝒦_c`$ is a polynomial of degree $`D_p`$ in phase and $`D_c`$ in color. In this study, we chose for the degrees of polynomials $`(D_p,D_\lambda ,D_s,D_c)=(4,3,1,1)`$. The large degree in phase permits a detailed adaptation of the model light curve shapes to the actual data. The other degrees correspond to the minimal number of coefficients (a degree 3 in wavelength is chosen to adapt colors in $`UBVR`$, independently). The $`𝒦`$ polynomials are then defined by 48 coefficients, but the constraints reduce the number of independent coefficients to 34. There are 32 free coefficients for $`𝒦_s`$ and 2 free coefficients for $`𝒦_c`$. We will call “training” the determination of these coefficients from measurements of nearby SNe Ia.
### 2.2 Normalization of transmissions
The value of $`f_0`$ depends on the normalization of the spectral template $`\varphi `$. This is not an issue when comparing nearby and distant supernovae to measure cosmological parameters, since only flux ratios matter: all objects have to be fitted using the same model. For $`H_0`$ measurements, one would need to normalize $`(\varphi )`$ using SNe Ia at known distances, but we will not perform this here.
More important is the apparent dependence of the value of $`f_0`$ with respect to the normalization of the transmission $`T`$ (equation 1). Since we aim at fitting multi-color light curves, observed with different instruments, with a single set of parameters $`(f_0,s,c)`$, it is mandatory to remove this dependence.
All photometric data are expressed in units of the integrated flux $`f_{ref}`$ (deduced from the zero point) of a known standard spectrum $`\varphi _{ref}`$ (e.g. Vega). While the functional form of $`T(\lambda )`$ is determined by optical transmission measurements its normalization can be determined from $`f_{ref}`$ and $`\varphi _{ref}`$ via the relation
$$\varphi _{ref}(\lambda )\frac{\lambda }{hc}T(\lambda )𝑑\lambda =f_{ref}$$
(3)
The filter transmissions are usually published as series of numbers (as in Bessell 1990), and this leads to ambiguities: for example, Suntzeff et al. (1999) writes ”Note that Bessell (1990) defines the sensitivity function as the product of the quantum efficiency of the detector+telescope, the filter transmission curve, the atmospheric extinction, and a linearly increasing function of wavelength.”. The ambiguity is whether the transmission refers to signal per unit energy, or signal per photon flux. Attributing dimensions to transmissions solves the ambiguity.
There have also been concerns about the relative weights of the different wavelengths in the integrals over wavelength (Nugent et al. 2002), namely whether one should sum photons or energies. If one aims at reproducing the instrument response to an arbitrary spectral energy density, the mathematical integrals should mimic the physical integration process of the instrument (Fukugita et al. 1996). For example, if one considers a CCD-based observing system, the effective transmission $`T(\lambda )`$ will be proportional to a number of photo-electrons per photon. So one should not integrate photon counts nor energies, but charge on the detector (or ADC counts).
## 3 Training the model
Since $`f_0`$ describes the observed luminosity in the $`B`$ band (because of the constraints applied to the corrections), the model only incorporates stretch-shape and stretch-color relations, but no correlation involving luminosity. This option was chosen in order to allow us to train the model with objects at unknown distances, in particular the nearby objects in the sample of Jha (2002) measured in the $`U`$ band. If one offsets all magnitudes of each training object by an arbitrary amount, possibly different for each object, the resulting model will not change. One could then consider incorporating high redshift objects into the training, but we choose not to do it here.
### 3.1 The nearby SN Ia Sample
The model was trained and tested using a sample of published nearby supernova light curves. We collected 122 SNe Ia for which $`B`$\- and $`V`$-band light curves are available in the literature, including data from Hamuy et al. (1996b); Riess et al. (1996b) and Jha (2002) for a total of 94 objects, and 28 additional supernovae collected from various sources (see the caption of table 1).
Objects were then selected based on two main criteria. First, we kept supernovae with at least two measurements before the maximum in the $`B`$ or the $`V`$ band. This is necessary to ensure that the date of maximum is well defined and that the measurements can safely be used as a function of phase. Out of the whole sample, 56 SNe satisfied this criterion. Then under-luminous peculiar supernovae: SN 1991bg, SN 1998bp, SN 1998de and SN 1999by (Howell 2001; Li et al. 2001b), and the peculiar objects SN 2000cx and SN 2002er (Li et al. 2001a; Pignata et al. 2004) were rejected from our sample. There are a number of reasons for this choice; the most important one being the spectral difference between normal and SN 1991bg-like SNe Ia. Since our model is built using a spectral template describing the features of the average normal SN Ia, it is not well suited for describing very different objects, such as SN 1991bg-like events. We note that, this is not a constraint for cosmological studies since, so far, no under-luminous SNe have been found in distant SN searches. Furthermore, if present at high-redshift, these objects would easily be identified both from spectroscopy and from their $`BV`$ color evolution (see for instance Phillips et al. 1999; Garnavich et al. 2004).
SN 1991T-like events were kept in the sample. This ”sub-class” of over-luminous SNe Ia is spectroscopically identified by the presence of unusually weak absorption lines during the pre-maximum phase. These are more difficult to identify than their under-luminous counterparts in low signal-to-noise spectra, such as those usually available for high redshift SNe. Moreover, the $`BV`$ color evolution is not particularly different from the one of normal SNe. Ideally, one should use a different spectral template for fitting this kind of events. We estimate the errors introduced in the k-corrections using a standard SN Ia spectral template to be of a few percent in the worst cases, and well within the final dispersion in the peak luminosity of SNe Ia. We further note that, the uncertainty due to using a non optimal spectral template is analogous to the uncertainty in the cross filter $`k`$-corrections used in standard methods.
The resulting sample of 50 SNe was then split into two sets: a training sample and a test sample. The training sample was used to adjust the coefficients of the polynomials of the model. It contains all the supernovae with redshifts smaller than 0.015 (not in the Hubble flow) and 6 supernovae at redshifts above 0.015, for which with $`U`$-band data was available, in order to improve the model in this wavelength region. The training sample contains the 34 supernovae listed in table 1. The test sample contains 26 supernovae (table 2). Note that the two samples are not completely independent. Indeed, they share ten supernovae with $`U`$-band light curves, since such events are rather scarce.
The data was not pre processed in any way prior to fitting. To account for the Milky Way extinction, we incorporate it into the instrument transmission, using the law from Cardelli et al. (1989) with a color excess $`E(BV)`$ obtained from Schlegel et al. (1998) dust maps at the position of the object to fit.
### 3.2 Training the Model
All the published nearby supernova magnitudes are expressed in the Johnson-Cousins $`UBVR`$ system. In Equation (1), we use models of the instrument transmissions as a function of wavelength. We adopted the transmission functions published by Bessell (1990), and interpreted them as $`\lambda T(\lambda )`$ (see equation 1), i.e. counts per unit energy, following a footnote of Suntzeff et al. (1999).
Training the model consists in determining the $`𝒦(p_s,\lambda ,s,c)=𝒦_s(p_s,\lambda ,s)+𝒦_c(\lambda ,c)`$ correction function (Eq. 2) using the training sample data in the $`UBVR`$ bands. The $`R`$ band was introduced to avoid relying on extrapolation for measurements beyond rest-frame $`V`$ band. We start with a first guess:
$`𝒦_s(p_s,\lambda ,s)`$ $`=`$ $`0`$
$`𝒦_c(\lambda ,c)`$ $`=`$ $`c\times (\lambda \lambda _B)/(\lambda _V\lambda _B)`$
and we use an iterative algorithm which can be sketched as follows:
1. Fit the light curves using the current determination of $`𝒦`$.
2. Fit $`\delta 𝒦(p_s,\lambda ,s,c)`$, an instance of the $`𝒦`$ function, on the light curve residuals. During this step, identify and remove the outliers data points.
3. $`𝒦𝒦+\delta 𝒦`$.
4. GOTO step 1, until $`\delta 𝒦`$ becomes negligible.
All the available data, i.e. the $`UBVR`$ residuals, up to large phases were used to determine the $`𝒦`$ function (step 2). However, not all the data points were used to fit the light curves (step 1): the $`R`$-band data was not used and only $`UBV`$ photometric points with phases ranging from $`15`$ days to $`+35`$ days were used since we are mostly interested in describing the rest-frame $`UBV`$ central part of the supernova light curves.
### 3.3 Results of the training
The fit converged after four iterations. 2480 measurement points were fitted, and 39 were discarded as outliers (at the 3 $`\sigma `$ level). Compared to the number of free coefficients of the model, we can safely conclude that the model is not over-trained. The standard deviations of the residuals to the model in $`UBVR`$ are respectively of $`0.09,0.09,0.06,0.07`$ magnitudes. Figure (1) shows the final $`U`$, $`B`$, $`V`$ and $`R`$ templates obtained at the end of the process as a function of stretch. By construction, the rest-frame $`B`$ and $`V`$-band magnitudes at maximum do not vary with stretch. We find a strong dependence of $`(UB)_{max}`$ with stretch ($`\delta (UB)_{max}\delta s`$, compatible with Jha 2002), which is an essential feature for the model to reproduce in order to estimate a reliable color in the wavelength range between $`U`$ and $`B`$. The model also manages to reproduce a a stretch-dependent secondary shoulder in the $`R`$ band. We also notice that the $`R`$-band light curves cross each-other for different values of the stretch. This reproduces well the fact that the brighter slower relation is weaker in the redder pass-bands, as noted by Phillips et al. (1999). The residuals to the light curve fit are shown figure 2.
Figure 3 represents the color correction $`𝒦_c(\lambda ,c)`$ for $`c=0.1`$ compared to the dust extinction law from Cardelli et al. (1989). Interestingly enough, the law we obtain follows pretty well that of Cardelli in the $`R`$ band but not in the $`U`$ band where we get a stronger dependence on $`c`$. Also shown (shaded area) is the uncertainty on $`𝒦_c`$ derived from the $`\chi ^2`$ increment (normalized to the number of supernovae) of the fit to the light curve residuals. The Cardelli law in the $`U`$ band is at 3.7 standard deviation from the best fit value.
As a consequence, we deduce that the relation between E(B-V) and E(V-R) are very similar (i.e. indistinguishable) to the ones expected from reddening by dust. This similarity, noted by Riess et al. (1996b), does not prove however that $`c`$ can be interpreted as reddening by dust; an additional requirement for this hypothesis to be valid would be that the peak $`B`$-band magnitude increases with $`c`$ by a value of $`R_B\times c`$. We will see that this is not the case in the next section. Let us also emphasize that the stretch dependent part of the $`UB`$ and $`VR`$ colors are included in the stretch dependent term $`𝒦_s(p_s,\lambda ,s)`$, and not in the color curve of figure 3.
## 4 Performance study
Once the correction function $`𝒦(p_s,\lambda ,s,c)`$ is determined, we can fit the model on the sample of nearby SN Ia light curves listed in table 2. This allows us to perform various consistency checks, in order to make sure that the model describes well the $`UBV`$ photometry of SNe Ia.
A first test consists in checking the ability of the model to reproduce the shape and color features of the independent set of SNe Ia. This is demonstrated on figure 2 which shows the residuals to the fits of light curves of SNe from the test sample.
### 4.1 Distance estimate
The global intensity parameter $`f_0`$ is proportional to $`d_L(z)^2`$ and inversely proportional to the normalization of $`\varphi `$. One can define a rest-frame $`B`$ magnitude $`m_B^{}`$ (Perlmutter et al. 1997) which removes this artificial dependence on the model normalization,
$$m_B^{}=2.5\mathrm{log}_{10}\frac{f_{SN}(0,z,T_B^{})}{(1+z)f_{ref}(T_B)}$$
(4)
where $`f_{SN}`$ and $`f_{ref}`$ are respectively defined by equations 1 and 3, $`T_B`$ is the transmission of the $`B`$ filter and
$`T_B^{}(\lambda )=T_B(\lambda /(1+z))`$ is a redshifted $`B`$ transmission. One can check that $`m_B^{}`$ varies as $`5\mathrm{log}_{10}d_L(z)`$ with redshift and that $`m_B^{}m_B`$ for $`z1`$, where $`m_B`$ is the conventional $`B`$ magnitude.
We incorporate the Hubble parameter dependence of $`d_L`$ in a constant parameter $`M_B^{70}=M_B5\mathrm{log}_{10}\left(h_{70}\right)`$, which is the average absolute magnitude of a SN Ia with $`s=1`$ and $`c=0`$, for a value of the Hubble parameter of 70 km.s<sup>-1</sup>.Mpc<sup>-1</sup>.
As mentioned in the introduction, the peak luminosity of SNe Ia is correlated to stretch and color, so we may build a distance estimator that accounts for those correlations and as a result reduces the dispersion. Following Tripp (1998), we adopt linear corrections of coefficients $`\alpha `$ and $`\beta `$ respectively for stretch and color. The distance estimator is then
$$m_B^{}M_B^{70}43.16+\alpha (s1)\beta c$$
(5)
Its expectation value for a supernova at redshift $`z`$ is $`5\mathrm{log}_{10}\left(d_L(z)H_0c^1\right)`$. Since our goal is here to test the distance estimator rather than actually perform a cosmological fit, we impose the “concordance” cosmological parameters ($`\mathrm{\Omega }_M=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$) when fitting $`M_B^{70}`$, $`\alpha `$ and $`\beta `$ and apply the method to build low-z Hubble diagrams using successively $`(B,V)`$ only and $`(U,B)`$ only light curves of supernovae from the test sample (table 2). In general, rest-frame $`(U,B,V)`$ light curve triplets should be used, when available, to determine the light curve parameters.
#### 4.1.1 Hubble diagram in $`BV`$ only
Using $`B`$\- and $`V`$-band only light curves of supernovae with redshifts larger than 0.015 from the test sample, we obtain: $`M_B^{70}=19.41\pm 0.04`$, $`\alpha =1.56\pm 0.25`$ and $`\beta =2.19\pm 0.33`$. The standard deviation of residuals is $`0.16\pm 0.03`$ <sup>5</sup><sup>5</sup>5Note that this number takes into account the number of parameters in the fit. The measured RMS value is $`0.14\pm 0.03`$.. Uncertainties on $`m_B,s,c`$ along with their covariance were included in the fit<sup>6</sup><sup>6</sup>6The uncertainties on the distance estimate formally depend on $`\alpha `$ and $`\beta `$, and increase with them. As a consequence, the $`\chi ^2`$ minimum is biased toward large values of these parameters. We therefore computed the uncertainties with the initial values, and use the result of the fit at the final iteration., we also considered an uncertainty on redshifts due to peculiar velocities of 300 km.s<sup>-1</sup>; an additional “intrinsic” dispersion of 0.13 is needed in order to get a $`\chi ^2`$ per degree of freedom of 1. The resulting Hubble diagram as well as the observed brighter–slower and brighter–bluer relations are shown figure 4.
Our approach to estimating distances easily compares to the one adopted in Tripp (1998): the main differences are the light curve model and the brighter-slower parameterization. When we fit the same SNe sample (The Calán-Tololo sample from Hamuy et al. 1996a), a value of $`\alpha =1.04\pm 0.24`$ and $`\beta =2.08\pm 0.27`$ are obtained, which compare well to $`\alpha =0.88`$<sup>7</sup><sup>7</sup>7$`b=0.52`$ translates to $`\alpha 0.88`$ when using stretch and the first order relation $`(\mathrm{\Delta }M_{15}1.1)1.7(1s)`$., $`\beta =2.09`$ of Tripp (1998), based on peak luminosity, color, and decline rate estimates from Hamuy et al. (1996a). We conclude that our model correctly reproduces basic parameter estimations of previous works.
Concerning the interpretation of the brighter-bluer correlation, we find a value of $`\beta `$ which is incompatible with $`R_B=4.1`$, expected for extinction by dust analogous to the observed law in the Milky Way. The value we find is compatible with those found in previous works (see Tripp (1998) and references therein). However, as stressed in Riess et al. (1996b), the color excess (or deficit) at maximum should not be interpreted as entirely due to extinction but be corrected for the part of this excess that is correlated with stretch . We measure a stretch-color slope of about 0.2, similar to the relation proposed in Phillips et al. (1999)<sup>8</sup><sup>8</sup>8The proposed relation is $`\frac{dc}{d\mathrm{\Delta }M_{15}}=0.114\pm 0.037`$. With the approximate relation $`\frac{d\mathrm{\Delta }M_{15}}{ds_B}1.7`$ (at $`s_B`$ = 1), we expect $`dc/ds0.2`$. and can redefine our parameters to account for this correlation:
$`c^{}`$ $`=`$ $`c+0.2(s1)`$
$`s^{}`$ $`=`$ $`s`$
so that $`s^{}`$ and $`c^{}`$ are uncorrelated. The correlation coefficients then become $`\alpha ^{}=\alpha +0.2\beta `$ and $`\beta ^{}=\beta `$, which means that redefining the color excess to explicitly assign to stretch the color variations correlated to stretch does not change the brighter-bluer correlation strength.
The two-parameter (stretch and color) pragmatic approach we followed can accommodate both reddening by dust and any intrinsic color effect dependent or not on stretch. One may reasonably assume that reddening by dust and stretch independent intrinsic colors mix (as proposed in Nobili et al. 2003), and that disentangling the contributions would improve the distance resolution. Our distance indicator is however independent of the interpretation of the color variations. Since the low value of $`\beta `$ may indicate that some intrinsic effect plays a role, we did not interpret color as only due to reddening by dust and hence accepted negative $`c`$ values as such.
#### 4.1.2 Hubble diagram in $`UB`$ only
We applied the same procedure as in the previous section to fit the $`U`$ and $`B`$-band light curves of the sub-sample of table 2 for which $`U`$-band measurements are available and redshifts larger than 0.015 (9 supernovae). We obtain $`M_B^{70}=19.37\pm 0.05`$, $`\alpha =0.8\pm 0.4`$, $`\beta =3.6\pm 0.6`$, and the standard deviation of residuals is $`0.16\pm 0.05`$. As expected, these results are consistent with the fit using $`B`$ and $`V`$, as shown by the confidence contours for $`\alpha `$ and $`\beta `$ fitted using either $`UB`$ or $`BV`$ light curves shown figure 5. Note the covariance between the estimated values of $`\alpha `$ and $`\beta `$, particularly in the $`U+B`$ band case, which simply reflects the correlation between the parameters $`c`$ and $`s`$ in the test sample.
One of the SN present in this sample (1996bo) show a large statistical uncertainty due to its limited $`U`$-band photometry. Removing this point from the fit has the effect of bringing down the $`U,B`$ contour to the point that it contains almost all of the $`B,V`$ contour, bringing the two determinations of $`\alpha `$ and $`\beta `$ closer to each other.
Figure 6 presents the residuals to the Hubble diagram as a function of redshift, stretch and color using the values of $`M_B^{},\alpha ,\beta `$ fitted with $`B`$\- and $`V`$-band light curves in the previous section. Note the large uncertainty affecting 1996bo, which appears as a $`2`$ sigma outlier in the 3 plots. Fitting the Hubble diagram with the values of $`\alpha `$ and $`\beta `$ obtained with $`B`$\- and $`V`$-band light curves, the standard deviation of residuals is $`0.20\pm 0.05`$.
Comparisons of the fitted values of $`s`$ and $`c`$ fitted using either $`UB`$ or $`BV`$ data are shown figure 7. The error bars only reflect the propagation of photometric errors and do not account for any intrinsic dispersion. There is no significant bias between the two estimates.
### 4.2 Comparison with other luminosity distance estimators
We compare our method with various estimators of SN Ia distances by examining the dispersions about the Hubble line of a sample of nearby supernovae. The comparison is made on distances measured in $`B`$ and $`V`$ bands only for which we obtain a dispersion of $`0.16\pm 0.03`$. Our $`U`$\- plus $`B`$-band distances estimates could not be compared with other estimate due to the lack of published distances measured in these 2 bands.
As shown in many papers (e.g. Wang et al. 2003 and Wang et al. 2005), testing luminosity distance indicators on “low-extinction sample” greatly improves their performance. As said in the introduction, cutting on color estimate measured with a redshift dependent accuracy is a source of systematic errors which we need to avoid for cosmological applications. We hence compare our resolutions to “full-sample” resolutions, for distance indicators involving rest-frame $`U`$, $`B`$ and $`V`$ bands.
The MLCS method was originally presented in Riess et al. (1996a), quoting a distance resolution of 0.12. Its latest development is presented in Jha (2002), with a distance resolution of 0.18, and used in Riess et al. (2004). From the latter, we collected the distance measurements to the 20 objects in common with our test sample, and compute a Hubble diagram dispersion of 0.24 (0.22 with the low-redshift “golden sample”) ) to be compared with our value of 0.16 when measured on the same 20 events.
The CMAGIC method of Wang et al. (2003) finds a weighted dispersion of 0.08 for a sub-sample of SNe with $`B_{max}V_{max}<0.05`$. With a weaker cut on color, $`B_{max}V_{max}<0.5`$, the dispersion rises to about 0.15, which is consistent with our result. Similarly, the $`\mathrm{\Delta }C_{12}`$ method presented in Wang et al. (2005), also present an exquisite distance resolution of 0.07 in the $`V`$ band based on a low-extinction subsample. However, when considering the full sample, the distance resolution degrades to 0.18. A more detailed comparison with CMAGIC and $`\mathrm{\Delta }C_{12}`$ would require comparing the distances of the same sample of SNe (due to the limited statistics) but their distances are not published. Note also that in these papers, the test and training samples are not separated.
In summary, the method we propose gives a dispersion on distances measured using $`B`$\- and $`V`$-band data only comparable or lower than obtained with other methods while also providing, for the first time, comparable dispersion values for distances measured using $`U`$\- and $`B`$-band data only.
## 5 Conclusion
We have proposed a new method to fit broadband light curves of Type Ia supernovae. It allows us to determine simultaneously the SN Ia rest-frame $`B`$ magnitude at maximum, stretch and color excess (or deficit) using any measured multi-color light curve within the wavelength range of rest-frame $`UBV`$ bands. This technique is particularly well-suited to the treatment of high-redshift SNe Ia for which limited coverage is obtained in both wavelength and phase.
The k-corrections, which allow the observer to transform the observed magnitudes into the standard rest-frame magnitudes, are built-in; the model includes the dependence on stretch and color of the spectrum template needed to estimate those corrections. In particular, the well-known correlation between $`(UB)_{max}`$ and stretch is reproduced.
The $`(BV)`$ and $`(VR)`$ stretch-independent colors we obtain are extremely similar to the ones expected from reddening by dust. The $`(UB)`$ color departs from this law. We find a relation between $`(BV)`$ color and observed $`B`$ luminosity incompatible with $`R_B=4.1`$, at more than 3 standard deviations, even when accounting for the stretch–color correlation.
We have tested this fitting procedure on an independent sample of SNe Ia. Alternatively using $`B`$\- and $`V`$-band data and $`U`$\- and $`B`$-band data, we are able to retrieve consistent parameters and hence build Hubble diagrams with both sets of data. The dispersions about the Hubble line were found to be $`0.16\pm 0.03`$ and $`0.16\pm 0.05`$ in the $`B`$ plus $`V`$ and $`U`$ plus $`B`$ bands only, respectively. This method is particularly well adapted to reliably measure SN Ia distances in the full redshift range and in particular beyond redshift $`z0.8`$ for which rest-frame $`V`$-band measurements are often not available.
###### Acknowledgements.
It is always a pleasure to acknowledge the stimulating discussions among the FROGS (FRench Observing Group of Supernovae), especially with G. Garavini, J. Rich and R. Taillet whom we would also like to thank for their critical reading of the manuscript. We thank P. Nugent for providing us with his template spectrum time series.
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# Spin-1/2 particles moving on a 2D lattice with nearest-neighbor interactions can realize an autonomous quantum computer
## 1 Introduction
To understand which sets of quantum control operations are sufficient for quantum computing is still an important issue of research. Whereas the standard model of the quantum computer is based on one- and two-qubit unitaries there is meanwhile a large number of alternative proposals, e.g. computing by measurements only or adiabatic computing . The latter model encodes a computational problem into an interaction such that the ground state of the Hamiltonian indicates the solution. In order to understand the power of adiabatic computing it is particularly interesting to know to what extent the transition into the ground state of simple interactions can already be a non-trivial computation process. Since it has been shown that every problem in the complexity class QMA can be encoded in a nearest-neighbor interaction of qubits located on a 2D lattice it is clear that, indeed, relatively simple Hamiltonians are sufficient for adiabatic computing. This highlights the computational power of simple Hamiltonians from one point of view. Another aspect of the “computational power of Hamiltonians” has been studied in models where the natural time evolution is interspersed by external control operations (see e.g. ). The intention of the present paper is to understand to what extent models with simple interactions can already be autonomous devices with full quantum computing power in the sense that quantum computation reduces to the following protocol: (1) prepare an initial state in the computational basis which contains the data, (2) wait for a sufficiently long time, and (3) measure a sufficiently large set of qubits in the computational basis. In we have constructed a Hamiltonian satisfying these conditions with the additional feature that the readout need not necessarily be performed within a specific time interval; in our “ergodic quantum computer” the result is also present in the time average of the dynamics after some simple post-selection with high success probability. Even though this property may be of minor practical relevance, we considered it as a necessary feature of an autonomous computer since one would otherwise require a clock as an additional device. Here we construct an interaction that is even simpler than the one constructed in since the latter requires 10-qubit interactions.
The motivation to consider autonomous quantum computers is twofold. First, it could trigger new ideas how to reduce the set of necessary control operations in current implementation proposals by using the “natural power” of the interactions. In addressing this issue, the ergodic model defines only an extreme case; realistic perspectives for quantum computing could arise by combining it with more conventional approaches. Second it is an interesting fundamental issue in the thermodynamics of computation how to realize computation in closed physical systems. Benioff, Feynman, and Margolus have already presented such Hamiltonian computers . Margolus’ asynchronous Hamiltonian cellular automaton (CA) has the appealing feature to be lattice-symmetric. In other words, it is driven by a finite-range interaction among qubits<sup>1</sup><sup>1</sup>1In it is furthermore argued that one of this clock Hamiltonians for CAs in one dimension is quite close to real solid state interactions. located on a 2D lattice such that the total Hamiltonian is invariant under discrete translations. The synchronization of the CA is realized by a kind of spin wave propagating along the lattice and triggering the update of the cells according to some computationally universal update rules which are not specified any further in . However, its clock wave has to start in an uncertain position in order to obtain a well localized momentum distribution with mainly positive momenta. A localized wave front would also propagate backwards and would therefore not trigger a correct computation. In we have chosen the same synchronization scheme but we start with a localized wave front since we allow only for preparations of basis states. The fact that the dispersion of the wave front leads to completely undefined computation steps is irrelevant since the time average of the dynamics encodes the correct result. The feature of our ergodic model to show the correct computation result also in the time average was hence only a nice byproduct of the fact that we must use a concept which works with a strange kind of clock: The latter starts with a well-defined time but then it counts backwards and forward with completely undefined counting speed. Note that the time average can also be considered as the “generalized final state” of our computer since a final output state in the usual sense cannot exist for finite dimensional Hamiltonian models.
To ask for a simple finite range Hamiltonian that is universal for quantum computation is in some sense similar to asking for a simple computationally universal quantum cellular automaton (for some proposals see ) with the only difference of considering update rules which change the state only in an “infinitesimal” way. However, local interactions in lattices have typically the property to spread the information into increasingly large regions, whereas it is possible to construct update rules for cellular automata that work by propagating the information forward column by column . This apparent difference between discrete and continuous time dynamics can already be explained with the translation operator on $`n`$ qubits: the cyclic shift is a unitary that can be achieved by local update rules , whereas the Hamiltonian obtained from the logarithm of the shift contains interactions between distant qubits. Hamiltonians with finite interaction length will typically spread the information over the lattice. Therefore, one has to be more modest and demand that the correct result can only be present with high probability.
The paper is organized as follows. In section 2 we explain the hardware of our model consisting of a chain of atoms on a 2D square lattice. We construct a nearest-neighbor interaction which leads via a simple effective Hamiltonian to a coupled random walk that will later trigger the implementation of gates acting on the atom spins. In section 3 we explain how to arrange spin-spin interactions among the atoms such that they implement logical gates based on holonomic quantum computing. In section 4 we describe the complete Hamiltonian of our computing device. In section 5 we show that the time evolution of the effective Hamiltonian is solvable since it can be transformed into a quasifree evolution of fermions. Based on this solution, we estimate the time required for the computation process. In section 6 we briefly sketch some ideas on the realization. Section 7 describes how to construct the Hamiltonian such that it implements a universal quantum cellular automaton. This is to obtain programmable hardware. In the appendix we present an alternative model for the atom propagation which is caused by an even simpler Hamiltonian since the nearest-neighbor interactions in the model presented in the main part involves also diagonal neighbors in the lattice. In contrast, the model in the appendix uses only nearest neighbors in a strict sense. However, the corresponding Hamiltonian dynamics will not be solvable. We will therefore consider an incoherent analogue and show that the corresponding classical random walk would trigger the implementation of gates in the desired way. For the coherent model, we can only conjecture that an appropriate propagation is achieved.
## 2 The synchronization Hamiltonian
Now we present a model of a 2D lattice with a coupled propagation of a chain of atoms which will later be the synchronization mechanism of the computer. While propagating along the horizontal direction (see Fig. 1) the spin of the atoms (which represent the logical states) will be subjected to spatially inhomogeneous interactions that implement the desired gates. It is important that the chain does not tear since the spin-spin interactions to be described later will only be active between adjacent atoms. Furthermore, it is important for the synchronization that the chain remains connected. This will become clear in section 3.
We will now describe the synchronization Hamiltonian in detail and show that it generates a diagonalizable quantum walk. The lattice and the initial atom configuration are shown in Fig. 1. The lattice sites are given by $`n\times n`$ black fields of the modified chess-board in Fig. 1.
The $`2n1`$ atoms are only allowed to move forward or backward along the horizontal direction, i.e., along the rows. Each atom can only hop from a black site to either the next or the previous black site in the row. Furthermore, the atoms can only move such that the chain does not tear, i.e., atoms in adjacent rows are always diagonal neighbors. Fig. 2 shows a possible configuration.
Now we describe the Hamiltonian that allows only collective motions of the atoms respecting these rules. The grid of black fields is labeled by $`(i,j)`$ with $`i,j=1,\mathrm{},n`$. The sites $`(1,j)`$ with $`j=1,\mathrm{},n`$ indicate the diagonal line leading from the leftmost field to the uppermost one and sites $`(i,1)`$ with $`i=1,\mathrm{},n`$ indicate the diagonal line from the leftmost field to the lower-most one. We consider the Hilbert space of all atom configurations as a subspace of an $`n^2`$-qubit register of the form
$$_q:=(^2)^{n^2},$$
where qubit $`(i,j)`$ corresponds to lattice site $`(i,j)`$ and the basis states $`|0`$ and $`|1`$ indicate that there is no atom or that there is an atom, respectively, at position $`(i,j)`$. A subspace of $`_q`$ is the clock Hilbert space $`_c`$ spanned by all allowed configurations. They correspond to those basis states in $`_q`$ for which there is exactly one atom in each row and for which the atom chain is connected. First we introduce interactions between site $`(i,j)`$ and $`(i+1,j+1)`$ which generates independent hopping of atoms along the rows:
$$K:=\underset{i,j=1}{\overset{n1}{}}(a_{i,j}a_{i+1,j+1}^{}+h.c.).$$
(1)
To achieve that the atom configuration remains a connected chain we introduce a strong attractive force between diagonal neighbors:
$$H_{pot}:=E\underset{<(i,j),(k,l)>}{}N_{i,j}N_{k,l}+E_0\mathbf{\hspace{0.17em}1},$$
(2)
where $`E^+`$ is assumed to be much larger than $`1`$ and $`N_{i,j}=a_{i,j}^{}a_{i,j}`$ is the projection on the states with an atom in position $`i,j`$. The sum runs over all unordered pairs of sites $`(i,j),(k,l)`$ with $`|ik|+|jl|=1`$. The physically irrelevant term $`E_0`$ is (for purely technical reasons) chosen such that the initial configuration has zero energy. The energy is minimal if all attractive interactions are active. Otherwise, if the chain is not connected, the potential energy is at least $`E`$. We define the “synchronization Hamiltonian”, which leads to the desired coupled motion of atoms, by
$$\stackrel{~}{H}:=K+H_{pot}.$$
(3)
In order to analyze the dynamics generated by $`\stackrel{~}{H}`$ we first argue that it can be replaced with an effective clock Hamiltonian $`H_{eff}`$ having $`_c`$ as an invariant subspace. The initial atom configuration defines a ground state of $`H_{pot}`$. Due to $`E1`$ the comparably small perturbation $`K`$ cannot circumvent the gap to one of the first excited states of $`H_{pot}`$. If $`P`$ denotes the spectral projection of $`H_{pot}`$ corresponding to the ground state energy $`0`$ we will therefore expect a dynamical evolution according to
$$H_{eff}:=PKP.$$
(4)
Note that not all states in the image of $`P`$ are allowed clock states since the image of $`P`$ contains also states with more than one atom in a row. However, these ground states could only be reached from an allowed state if $`K`$ contained hopping terms along columns.
The following lemma (proven in the appendix) justifies more formally that we may analyze the dynamics according to $`H_{eff}`$ instead of $`\stackrel{~}{H}`$:
###### Lemma 1
For $`E>n^6`$ and $`n10`$ the norm distance between the effective Hamiltonian in equation (4) and the restriction $`\stackrel{~}{H}_{E/2}`$ of $`\stackrel{~}{H}`$ to the eigenspace corresponding to eigenvalues smaller than $`E`$ satisfies
$$\stackrel{~}{H}_{E/2}H_{eff}9\frac{n^3}{\sqrt{E}}.$$
If we increase the attractive part of the interaction proportional to $`n^6`$ if $`n`$ increases the time evolution induced by $`H_{eff}`$ is a good approximation for the true one. If the intention of this article was to propose a physical implementation scheme this would be a rather bad scaling, leading possibly to hard practical problems. However, from the computer science point of view, it is essential that the increase of energy is only polynomial in $`n`$ and we will thus obtain an efficient model of computation.
We will now give a more explicit form of $`H_{eff}`$. To be more precise, we will describe an operator whose restriction to $`_c`$ coincides with $`H_{eff}`$. Let $`|c`$ denote some basis state in $`_c`$, i.e., $`c`$ describes an allowed atom configuration. Then
$$H_{eff}|c=PKP|c=PK|c=P\underset{1=i,j}{\overset{n1}{}}(a_{i,j}a_{i+1,j+1}^{}+h.c.)|c$$
(5)
is a superposition of all atom configurations with connected chain that can be reached from $`c`$ by moving one atom forward or backwards in its diagonal row. In other words, each hopping term $`a_{i,j}a_{i+1,j+1}^{}`$ is only active if the sites $`i,j+1`$ and $`i+1,j`$ are occupied. We have therefore a conditional hopping given by
$$H_{eff}=\underset{i,j=1}{\overset{n1}{}}a_{i,j}a_{i+1,j+1}^{}N_{i,j+1}N_{i+1,j}+h.c.|__c,$$
where the symbol $`_c`$ at the right indicates the restriction to $`_c`$. Using a slightly more intuitive notation, which respects the orientation of Fig. 1, the summands of $`H_{eff}`$ can be denoted by
$$\begin{array}{ccccc}& & N& & \\ & & & & \\ a& & & & a^{}\\ & & & & \\ & & N& & \end{array}+h.c.,$$
(6)
where this $`4`$-local operator acts on $`4`$ black sites enclosing a common white site.
We are interested in the dynamical evolution of the quantum state defined by the initial atom configuration. Since it is a ground state of the dominating Hamiltonian $`H`$ the major part of the state vector is contained in that spectral subspace of $`\stackrel{~}{H}`$ which corresponds to energy values not greater than $`E/2`$. Using the bound of Lemma 1 we may then estimate the error in the dynamical evolution caused by replacing the true evolution with the evolution according to $`H_{eff}`$.
###### Theorem 1
Let $`|\psi _c`$ be an arbitrary allowed atom configuration. Then we have
$$(e^{i\stackrel{~}{H}t}e^{iH_{eff}t})|\psi ϵt+2n\sqrt{\frac{2}{E}},$$
where
$$ϵ:=9\frac{n^3}{\sqrt{E}}$$
as in Lemma 1.
Proof: Due to $`H_{pot}|\psi =0`$ the average energy of $`|\psi `$ satisfies
$$\psi |\stackrel{~}{H}|\psi =\psi |K|\psi Kn^2,$$
where the last inequality is given by counting the number of terms in eq. (1). Let $`Q`$ be the spectral projection of $`\stackrel{~}{H}`$ corresponding to all eigenvalues less than $`E`$. Then we have
$$\psi |Q\stackrel{~}{H}Q|\psi +\psi |(\mathrm{𝟏}Q)\stackrel{~}{H}(\mathrm{𝟏}Q)|\psi =\psi |\stackrel{~}{H}|\psi n^2.$$
Using
$$\psi |(\mathrm{𝟏}Q)\stackrel{~}{H}(\mathrm{𝟏}Q)|\psi E(\mathrm{𝟏}Q)|\psi ^2,$$
and
$$\psi |Q\stackrel{~}{H}Q|\psi =\psi |K|\psi Kn^2$$
we conclude
$$(\mathrm{𝟏}Q)|\psi ^2\frac{2n^2}{E}.$$
(7)
We have
$`(e^{i\stackrel{~}{H}t}e^{iH_{eff}t})|\psi `$ $``$ $`(e^{i\stackrel{~}{H}_{E/2}t}e^{iH_{eff}t})Q|\psi `$ (8)
$`+(e^{i\stackrel{~}{H}t}e^{iH_{eff}t})(\mathrm{𝟏}Q)|\psi `$
$``$ $`ϵt+2(\mathrm{𝟏}Q)|\psi ,`$
where we have used $`\mathrm{exp}(iA)\mathrm{exp}(iB)AB`$ for two self-adjoint matrices $`A,B`$. The statement follows then by replacing the last term in eq. (8) with the right hand side of ineq. (7). $`\mathrm{}`$
One could use the condition $`t\sqrt{E}/(9n^3)`$ of Theorem 1 to determine the time scale for which $`H_{eff}`$ is a good approximation when $`E`$ is given. We can also assume that the time scale is given and we have to determine the interaction strength $`E`$. After all, we have to ensure that the approximation is valid for all $`tt_0`$ where $`t_0`$ is some a priori given upper bound on the time required by the clock wave to pass the relevant region on the lattice, i.e., the region containing the spin-spin interactions which will be described in the next section.
The effective Hamiltonian leads to a quantum walk on the space of atom configurations which will be diagonalized in section 5. In the next section we will explain how this walk may trigger the implementation of gates.
## 3 Holonomic implementation of logical gates
Now we consider not only the positions of the atoms but also their inner degree of freedom, e.g. their spin. We replace the qubit at each lattice site with a qutrit and extend hence the space $`_q`$ to
$$:=(^3)^{n^2}.$$
The basis states $`|0,|,|`$ indicate the absence of the atom or its two possible spin states, respectively. Before we describe how to add a spin-spin interaction Hamiltonian to $`\stackrel{~}{H}`$ which leads to the implementation of gates when the atoms are moving along the rows, we should first mention two obvious approaches to imprint gates by interactions and why they are not suitable for our goal.
In and the gates are directly coupled (as additional tensor product components) to the synchronization Hamiltonian. Certainly we could define the hopping term in eq. (1) between column $`j`$ and $`j+1`$ in any desired row such that the spin is not conserved during the atom propagation but subjected to some rotation. This would provide us with one-qubit gates. Such a one-qubit gate would clearly be inverted when the atom moves back to column $`j`$ again. The conditional state of the considered qubit, given that the atom is found on the right of column $`j`$ would hence be subjected to the desired transformation, no matter whether the atom has traveled back and forth several times as is is typical for a random walk.
To imprint two-qubit gates into the Hamiltonian is more difficult. The obvious method to couple their implementation with some coupled motion of two adjacent atoms would require more non-local interactions than only two-particle terms. The second obvious method would be to complete the clock Hamiltonian by an additive spin-spin interaction term between some adjacent sites in the same column. This leads to the following problems. First the atoms do not stay there for a well-defined time. Some part of the wave packet moves already and one part stays. Second the atoms travel back and forth and pass the interaction region several times. Both effects would in general entangle atom position and logical spin states in an uncontrollable way.
A solution to the latter problem is holonomic quantum computing . The idea of this approach, in the usual setting, is that some time-dependent Hamiltonian $`G(t)`$ is adiabatically changed along a closed loop such that, on an appropriate degenerate subspace, the overall effect is a unitary which depends only on the loop (described by representing $`G(t)`$ in some parameter space) and not on the speed of the change of $`G(t)`$. The unitary in the end remains the same even if the vector $`G(t)`$ has been moving back and forth on this loop in the parameter space.
We apply this concept to our model where we achieve the time-dependence of the spin-spin interaction by the motion of the atoms when they pass spatially inhomogeneous interactions. We imprint a spin-spin interaction (varying slowly along the rows) such that is describes a closed loop after the atoms have passed a certain region. Then the spin state of the atoms does not depend on the time required for the passing, neither does it depend on the number of times the atoms were traveling back and forth before the whole region has been passed. This ensures that the entanglement between the position degree of freedom and the spin is lost as soon as the region has been passed by the whole atom chain. Before we explain how to make use of this idea in detail, we rephrase a result in section II of Ref. , which is useful for us, as a lemma:
###### Lemma 2
(Holonomic Gates)
Given a family of Hamiltonians $`G(t)`$ on a finite dimensional Hilbert space $``$ by
$$G(t)=e^{iXt}G_0e^{iXt}\text{ with }t[0,T],$$
where $`X`$ is a self-adjoint operator on $``$ such that the family $`G(t)`$ is a closed loop, i.e., $`G(0)=G(T)`$. Let the change of $`G(t)`$ be sufficiently slowly to consider it as an adiabatic change, i.e., for each entry $`G_{i,j}(t)`$ of $`G(t)`$ we have
$$\left|\frac{d}{dt}G_{i,j}(t)/G_{i,j}(t)\right|\mathrm{\Delta },$$
where $`\mathrm{\Delta }`$ is the smallest spectral gap of $`G(0)`$.
Let $`U_T`$ be the time evolution generated by the family $`G(t)`$ after the closed loop. Let $`𝒦`$ be an eigenspace of $`G_0`$ with eigenvalue $`\lambda `$. Then $`U_T`$ keeps $`𝒦`$ invariant and its restriction is given by
$$e^{i\lambda T}e^{iQXQT},$$
(9)
where $`Q`$ is the projection onto $`𝒦`$. The unitary (9) is also called a non-abelian geometric phase.
Since the holonomic concept implements the well defined unitary
$$\mathrm{exp}(iQXQT)$$
only on a subspace $`𝒦`$ of the full Hilbert space our $`2n1`$ spin particles cannot be used for $`2n1`$ logical qubits. Instead, the atom in row $`2i1`$ will encode a logical qubit together with the atom in row $`2i`$. Moreover, we can only use rows being not too far from the row in the middle since the others are too short to imprint interactions that implement gates. Two adjacent spin particles encode one qubit with logical states $`|0_l,|1_l`$ via
$$|0_l:=||\text{ and }|1_l:=||.$$
We will call their span $`𝒞`$ the “code space”.
### One-qubit gates
To explain how to implement one-qubit gates we first restrict our attention to two adjacent rows. We add interactions as follows. The whole region which carries the gate consists of a stripe of width $`2`$ as shown in Fig. 3. We call this region “interaction stripe”. Inside this stripe, we add a “gate Hamiltonian” to $`\stackrel{~}{H}`$ which consists of spin-spin interactions $`V_j`$ between certain diagonal neighbors in adjacent rows (as indicated by the edges labeled with $`V_1,\mathrm{},V_l`$ in Fig. 3). Each $`V_j`$ is a pair-interaction between two sites.
We choose the following interactions:
$$V_j:=e^{i\frac{2\pi (j1)}{l1}X}(\sigma _zN+N\sigma _z)e^{i\frac{2\pi (j1)}{l1}X},$$
(10)
with $`j=1,\mathrm{},l`$. Here $`X`$ is some self-adjoint operator on $`^3^3`$ specified later and
$$N:=||+||$$
is the particle number operator in analogy to the operator $`N_{i,j}`$ in eq. (2). The Pauli matrix $`\sigma _z`$ has to be read as formally acting on the space $`^3`$ of the corresponding qutrit, even though it is zero for the state $`|0`$ and acts therefore only on the spin states. It is important to note that the degenerate eigenspace of $`V_0`$ and $`V_{l1}`$ coincides with the code space $`𝒞`$. If the spins are initialized such that their state is in $`𝒞`$ before they pass the interaction region they will at every time instant remain in the degenerate subspace of the interaction that is active (provided that the change between different $`V_j`$ is adiabatic). This will implement the “non-abelian phase” $`\mathrm{exp}(iQXQT)`$ of Lemma 2.
The length $`l`$ of the interaction stripe is chosen such that the change of interactions from $`j`$ to $`j+1`$ can be considered as approximating a continuous change and furthermore as adiabatic when the atom chain is propagating in horizontal direction. To sketch how to choose $`l`$ we define a typical time scale $`\tau `$ on which an atom jumps from one site to its neighbor. Then we have to ensure that each entry $`V_j^{ik}`$ of $`V_j`$ satisfies
$$\left|\frac{V_{j+1}^{ik}V_j^{ik}}{\tau V_j^{ik}}\right|\mathrm{\Delta },$$
(11)
where $`\mathrm{\Delta }`$ is the smallest energy gap of $`V_1`$. Here, the time $`\tau `$ is defined as a dimensionless quantity of the order $`1`$ since the hopping terms $`a_{i,j}a_{i+1,j+1}^{}+h.c.`$ in eq. (1) appears with constant one. Physical dimensions are irrelevant since inequality (11) is invariant with respect to a simultaneous rescaling of all $`V_j`$ and the hopping terms in eq. (1). We obtain therefore the dimensionless condition $`l1`$.
Now we may set
$$X:=\sigma _x(\mathrm{cos}\varphi \sigma _x+\mathrm{sin}\varphi \sigma _z),$$
(12)
with arbitrary angle $`\varphi `$, or
$$X:=\sigma _y(\mathrm{cos}\varphi \sigma _x+\mathrm{sin}\varphi \sigma _z),$$
(13)
where we have again slightly abused notation since we did not explicitly indicate the embedding of the 2-qubit operators in eqs. (12) and (13) (acting on the 4 possible spin configurations) into the two qutrit space. The following lemma shows which gates are implemented by the above choices for $`X`$.
###### Lemma 3
Let two atoms pass a region where they are subjected to the Hamiltonians $`V_1,\mathrm{},V_l`$ in eq. (10) with $`X`$ as in eq. (12) or eq. (13). Assume furthermore that $`l`$ is sufficiently large to consider the change of interactions as adiabatic. Then their encoded qubit is, up to an irrelevant global phase, subjected to the transformation $`\mathrm{exp}(2\pi i\mathrm{cos}\varphi \sigma _x)`$ or $`\mathrm{exp}(2\pi i\mathrm{cos}\varphi \sigma _y),`$ respectively.
Proof: On the space $`^2^2`$ spanned by the $`4`$ spin states the interaction changes according to
$$V_j:=e^{i\frac{2\pi (j1)}{l1}X}(\sigma _z\mathrm{𝟏}+\mathrm{𝟏}\sigma _z)e^{i\frac{2\pi (j1)}{l1}X},$$
with $`j=1,\mathrm{},l`$. Due to Lemma 2 the code space $`𝒞`$ is then subjected to $`\mathrm{exp}(i\mathrm{\hspace{0.17em}2}\pi A)`$ with $`A:=P_𝒞XP_𝒞`$, where $`P_𝒞`$ is the projection onto $`𝒞`$. For both choices for $`X`$, the relevant operator $`A`$ consists then only of transitions between $`|`$ and $`|`$, which corresponds to a $`\sigma _x`$ if the transition amplitudes are $`1`$ (as it is the case given by eq. (12)) and to $`\sigma _y`$ if they are $`i`$ and $`i`$ (as in the case of eq. 13). $`\mathrm{}`$
Lemma 3 shows that we can generate arbitrary one-qubit transformations by concatenating interactions with $`X`$ as in eq. (12) and (13) using appropriate angles $`\varphi `$.
### Two-qubit gates
For the implementation of logical two-qubit gates we consider interaction stripes that consist of $`3`$ adjacent rows (see Fig. 4), where row $`1`$ and $`2`$ belong to the first logical qubit and row $`3`$ is part of the second logical qubit. We define interactions $`V`$ and $`U_j`$ such that $`V`$ connects row $`1`$ and $`2`$ and $`U_j`$ connects row $`2`$ and $`3`$ as shown in Fig. 4. We will refer to the atoms in row $`1,2,3`$ as atoms $`1,2,3`$, respectively.
The interaction between row 1 and the row 2 is always $`V`$ in the whole interaction stripe. Row 2 interacts with row 3 according to $`U_1,U_2\mathrm{},U_l`$.
As soon as atom $`2`$ has entered the interaction stripe (shown in Fig. 4 for $`l=7`$) it is simultaneously subjected to interactions $`V`$ and $`U_1`$. Moreover, $`V`$ and $`U_l`$ are turned off simultaneously as soon as atom $`2`$ leaves the interaction stripe. The total interaction that is active on the system consisting of the spins of atoms 1,2,3 is therefore either
$$V\mathrm{𝟏}+\mathrm{𝟏}U_j\text{ with }j=1,\mathrm{},l,$$
when atom $`2`$ is inside the interaction stripe, or $`0`$ otherwise.
We choose
$$V:=\sigma _zN,$$
(14)
where the first tensor component refers to row 1 in Fig. 4 and the second to row 2. The operator $`N`$ has no effect on the spin of the atom in row $`2`$. It only ensures that the $`\sigma _z`$-Hamiltonian is not switched on before the atom in row $`2`$ has entered the left-most black field in row 3 in Fig. 4.
The interaction $`U_j`$ changes according to
$$U_j:=\sigma _z||+(e^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}\sigma _ze^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}})||,$$
(15)
where $`\stackrel{~}{X}`$ is an operator that acts on the one-qutrit space. It will be specified later. On the space $`^2^2^2`$ spanned by all $`8`$ possible spin states, we have
$$V\mathrm{𝟏}+\mathrm{𝟏}U_j=\sigma _z\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\sigma _z||+\mathrm{𝟏}(e^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}\sigma _ze^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}})||.$$
The idea is that the adiabatic change of the degenerate Hamiltonian on row $`1`$ and $`2`$ is controlled by the spin of the atom in row $`3`$, i.e., by the logical state of the second qubit. Whenever the state in row $`3`$ is $`|`$, rows $`1`$ and $`2`$ are subjected to the Hamiltonians
$$\sigma _z\mathrm{𝟏}+\mathrm{𝟏}(e^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}\sigma _ze^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}),$$
otherwise they are subjected to a constant Hamiltonian. We find:
###### Lemma 4
Let $`V`$ and $`U_j`$ be as in eq. (14) and eq. (15), respectively, such that they connect row $`1`$ and $`2`$ as well as $`2`$ and $`3`$, respectively. Choose the corresponding operator $`\stackrel{~}{X}`$ as
$$\stackrel{~}{X}:=cos\varphi \sigma _x+\mathrm{sin}\varphi \sigma _z,$$
with arbitrary angle $`\varphi `$. Consider the logical 2-qubit space $`𝒞𝒞`$, where the first tensor component refers to the atoms in rows $`1`$ and $`2`$ and the second to those in $`3`$ and the additional row $`4`$ (which is not depicted in Fig. 4). Then the adiabatic change of the interaction $`V\mathrm{𝟏}+\mathrm{𝟏}U_j`$ from $`j=1`$ to $`j=l`$ implements a controlled-$`\mathrm{exp}(i\mathrm{sin}\varphi \sigma _z)`$ transformation on $`𝒞𝒞`$ with the lower qubit as control qubit and the upper as target (in the tensor product in eq. (15) this correponds to the right and the left side, respectively).
Proof: Let the second qubit be in the logical $`0`$ state $`|0_l`$. This means that the spins of atom $`3`$ and $`4`$ are in the state $`|`$. However, for the interaction only the state of atom $`3`$ matters. Since it is $`|`$, atoms $`1`$ and $`2`$ are constantly subjected to the spin Hamiltonian
$$\sigma _z\mathrm{𝟏}+\mathrm{𝟏}\sigma _z.$$
(16)
Since the kernel of this operator coincides with the code space the logical state is not affected at all.
Assume now that the second qubit is in the state $`|1`$ and atom $`3`$ is therefore in the state $`|`$. Then atom $`1`$ and $`2`$ are subjected to the Hamiltonians
$$\sigma _z\mathrm{𝟏}+\mathrm{𝟏}(e^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}\sigma _ze^{i\frac{2\pi (j1)}{l1}\stackrel{~}{X}}),$$
with $`j=1,\mathrm{},l`$. Due to Lemma 2 the unitary that is implemented on the code space after such an adiabatic change is given by $`\mathrm{exp}(i2\pi A)`$ with
$$A:=P_𝒞(\mathrm{𝟏}X)P_𝒞=\mathrm{sin}\varphi P_𝒞(\mathrm{𝟏}\sigma _z)P_𝒞.$$
Describing this operator with respect to the logical basis $`|0_l,|1_l`$ we obtain the operator $`\mathrm{sin}\varphi \sigma _z`$. This shows that the unitary $`\mathrm{exp}(i2\pi \mathrm{sin}\varphi \sigma _z)`$ is implemented after the atoms have passed the interaction stripe. Since we have also shown that this unitary is only implemented if the state of the second qubit is $`|1_l`$, we have know shown that the net effect is a controlled-$`\mathrm{exp}(i2\pi \mathrm{sin}\varphi \sigma _z)`$ gate. $`\mathrm{}`$.
We can imprint interactions like those above between any adjacent logical qubits as well as we may obtain one-qubit rotations as in Lemma 3 by appropriate interactions. It is known that conditional phase gates together with the set of one-qubit rotations allow for universal quantum computing . Hence we are able to design interactions for arbitrary quantum circuits.
## 4 The complete Hamiltonian
We obtain the complete Hamiltonian $`\stackrel{~}{H}`$ of our autonomous computing device as a sum of the synchronization Hamiltonian and the spatially inhomogeneous spin-spin interaction that implement holonomic transformations. It reads:
$$\widehat{H}:=H_{pot}+K+\underset{<i,j,i^{},j^{}>}{}W_{i,j,i^{},j^{}},$$
where $`H_{pot}`$ is the strong attractive force defined in eq. 2 and $`K`$ describes the hopping terms
$$K:=\underset{i,j=1}{\overset{n1}{}}a_{i,j,}a_{i+1,j+1,}^{}+a_{i,j,}a_{i+1,j+1,}^{}+h.c,$$
where we have, in slight abuse of notation, adapted the definition of eq. (1) for spin-less creation and annihilation operators into those for spin 1/2 particles. For instance, $`a_{i,j,}^{}`$ creates a particle with spin down. Similarly, the term $`H_{pot}`$ is adapted to spin 1/2 particles in the sense that the particle number operator $`N_{i,j}`$ is given by
$$N_{i,j}=a_{i,j,}^{}a_{i,j,}+a_{i,j,}^{}a_{i,j,}.$$
The spin-spin interactions $`W_{i,j,i^{},j^{}}`$ are only non-zero if site $`(i,j)`$ is adjacent to site $`(i^{},j^{})`$ and if the pair of sites belongs to some interaction stripe. Inside these stripes, they are given by the interactions $`V,V_j,U_j`$ described in the preceding section. We assume that all these rectangles (”interaction stripes”) are confined to a square of length $`k`$, i.e., the spin-spin interaction is only non-zero for $`i,j,i^{},j^{}k`$ for some $`k<n`$. We call this region the circuit region (see Fig. 5). As soon as all atoms have passed it, the whole circuit is implemented. The complement of the interaction region will be called output region since we will read out the result of the computation there.
## 5 Diagonalization of the quantum walk
Now we show that the Hamiltonian $`\widehat{H}`$ leads to a propagation of atoms that implements the gates in correct order. First of all, we observe that the gates are irrelevant for the quantum walk. Different particle configurations correspond to mutually orthogonal vectors in the Hilbert space. Whether or not the particles are subjected to an additional change of their inner degree of freedom is irrelevant for the dynamics as long as we consider the adiabatic limit where the effect of the spin-spin interaction is only that it implements unitary gates on the spin states. Hence we will analyze the dynamics generated by $`H_{eff}`$ instead of the dynamics generated by $`\widehat{H}`$ to derive the relevant probabilities for finding the atom chain at certain positions. After recalling that the adiabatic change of the spin-spin interactions has implemented the desired gates given that the atom chain has passed the interaction stripes we have then shown that the circuit has been implemented.
To diagonalize the quantum walk we show that it is isomorphic to a walk of free fermions propagating on a one-dimensional chain. On the left of Figs. 1 and 2 we have characterized the configurations by binary words of length $`2n2=:2m`$. The symbols $`0,1`$ as $`j`$th digit indicate whether the atom in row $`j+1`$ is behind the atom on row $`j`$ or in front of it, respectively. The initial configuration is hence characterized by $`m=n1`$ symbols $`0`$ followed by $`m`$ symbols $`1`$. Since the atom on the very top and that one on the very bottom are fixed, the number of symbols $`1`$ remains constant and the vector space spanned by the possible atom configurations is therefore the subspace $`_m`$ of $`(^2)^{2m}`$ spanned by words with Hamming weight $`m`$. With respect to such a representation, the effective clock Hamiltonian $`H_{eff}`$, when restricted to $`_c`$, consists of operators which replace some pattern $`10`$ by $`01`$ or vice versa. It can be written as
$$H_s:=\underset{j=1}{\overset{2m1}{}}\mathrm{𝟏}^{j1}|01||10|\mathrm{𝟏}^{2mj1}+h.c.=\underset{j=1}{\overset{2m1}{}}b_j^{}b_{j+1}+h.c.,$$
(17)
where $`b`$ and $`b^{}`$ are fermion annihilation and creation operators, respectively. They are defined by
$$b_j:=\sigma _z^{j1}|01|\mathrm{𝟏}^{2mj},$$
and satisfy the canonical anti-commutation relation
$$b_i^{}b_j+b_jb_i^{}=\delta _{ij}\mathbf{\hspace{0.17em}1},$$
where we have chosen the convention $`\sigma _z|1=|1`$ and $`\sigma _z|0=|0`$.
Observe that $`H_s`$ is the XY-Hamiltonian which is well-known in solid states physics and generates a quasi-free evolution of fermions. Since the so-called Jordan Wigner transformation, describing the fermion interpretation formally, is standard , we rephrase it only briefly. The subspace $`_m`$ can be reinterpreted as the space of $`m`$ fermions moving without interactions in a $`2m`$-dimensional state space. Then $`_m`$ is simply the antisymmetric tensor product
$$_m\underset{m}{\underset{}{(^{2m}\mathrm{}^{2m})^{}}}.$$
The correspondence between these two pictures can be described easily: Given a basis state in the qubit register by a binary word with symbols $`1`$ at the positions $`j_1,\mathrm{},j_m`$. This state corresponds in the fermionic picture to
$$(|j_1|j_2\mathrm{}|j_m)^{}$$
where $`|j_i`$ denotes here the $`j`$th canonical basis vector in $`^{2m}`$.
The restriction of $`H_s`$ to $`_m`$ is in the fermionic picture given by
$`H_s`$ $`=`$ $`(S+S^{})\mathrm{𝟏}\mathrm{}\mathrm{𝟏}`$
$`+`$ $`\mathrm{𝟏}(S+S^{})\mathrm{𝟏}\mathrm{}\mathrm{𝟏}`$
$`+`$ $`\mathrm{}`$
$`+`$ $`\mathrm{𝟏}\mathrm{}\mathrm{𝟏}(S+S^{}),`$
where $`S`$ is the linear (non-unitary) shift acting on the basis states of $`^{2m}`$ via
$$S|j:=|j+1$$
for $`j<2m`$ and $`S|2m=0`$. This is because the term $`_jb_j^{}b_{j+1}`$ shifts the fermion by one site. In other words, the time evolution in each tensor component in eq. (5) is generated by the Hamiltonian $`S+S^{}`$, i.e., the adjacency matrix of the linear chain. This implies that the time evolution transfers annihilators and creators to linear combinations of annihilators and creators, respectively, when considered in the Heisenberg picture . Let
$$U_t:=\mathrm{exp}(i(S+S^{})t)$$
(18)
be the time evolution of one particle on a quantum walk on a chain and $`u_{jl;t}`$ its entries. Then the time evolution of the creation operator on site $`j`$ is
$$b_j^{}\underset{l2m}{}u_{jl;t}b_l^{},$$
(19)
i.e., it evolves into an operator creating a fermion which is in a superposition of different sites. For the annihilation operator at site $`j`$ we obtain
$$b_j\underset{l2m}{}\overline{u}_{jl;t}b_l.$$
(20)
To analyze the time evolution of relevant observables we will only make use of this formulation of the dynamics.
#### Time required for passing the circuit region
The running time of our computer is given by the length of the time one has to wait until all atoms will have passed the interaction region $`R`$, i.e., a square of length $`k<n`$, almost with certainty. The dynamical evolution $`\mathrm{exp}(iH_st)`$ is clearly quasiperiodic because the Hilbert space $`_m`$ is finite dimensional. Therefore the atoms will always return to $`R`$. However, the essential idea of the ergodic quantum computer is that the probability of finding it outside of $`R`$ is high if one measures at a random time instant. Later, we will therefore also consider the time average, before we estimate the time needed by the atoms to pass $`R`$ for the first time. We found:
###### Theorem 2
Let the circuit region be a square of length $`k`$. Then there is a time $`tO(k)`$ such that the probability of finding all atoms outside the circuit region $`R`$ is at least $`112/k`$ given that the size $`nk`$ of the whole lattice is sufficiently large.
The proof can be found in the appendix. Note that the probability of finding some atom in $`R`$ could even be made smaller if the factor $`4`$ in eq. (26) in the proof would be replaced with some larger number. We conclude that the implementation time for our imprinted circuit is in $`O(k)`$.
#### Computing the time average state
Now we want to show that also at a random time instant the probability of finding all atoms in the output region (i.e., the probability of finding at least $`k`$ fermions on the left half) is close to $`1`$. We found:
###### Theorem 3
Let the circuit region be a square of length $`k=m/4=(n1)/4`$. Then the probability of finding all atoms outside the circuit region tends to $`1`$ for $`n\mathrm{}`$.
This theorem is also proven in the appendix. One should maybe mention the following limitation of our model: in a lattice with finite length $`n`$ the adiabatic approximation underlying the holonomic implementation is only true up to some error. Therefore the implemented transformation, given that the atom chain has passed the circuit region, does not exactly coincide with the desired quantum circuit. Instead, it depends to some extent on the time needed to pass the region. In the limit of averaging the time evolution over an infinite time interval this error will increase. Roughly speaking (and expressed in a too classical language), the atom chain has then traveled back and forth many times and the error can increase since passing the circuit region backward would not necessarily invert the implemented circuit. However, it is a matter of the time scale on which the time average is taken whether this error is relevant, i.e., the degree to which the adiabatic approximation is justified determines the time scale on which the time average is described by our analysis.
## 6 Remarks on the realization
To judge whether it is likely that systems with Hamiltonians as above could be found in real systems would go beyond the scope of this article. Maybe one should rather ask which modifications are possible for our models that would make them more feasible. So far, the required spin-spin interactions are quite specific.
To show that the required diagonal hopping is not a priori unphysical one could think of electrons on quantum dot arrays arranged as in Fig. 6. The hopping along the rows does not require tunneling between distant dots even though it is the diagonal direction with respect to the square lattice. Spin-spin interactions would then only be needed in the direction of the dashed square lattice.
Attractive interactions between particles on adjacent dots could, for instance, be achieved if the particles alternate between electrons and holes from row to row. Certainly, it will be difficult to find a system that combines the attractive interaction and the hopping with the spatially inhomogeneous spin-spin interactions.
## 7 Imprinting a cellular automaton
We have shown that nearest-neighbor interactions among qutrits located on a 2D square lattice can be designed in such a way that their autonomous time evolution simulates any desired quantum circuit. The Hamiltonian given here is significantly simpler than that one given in our previous paper since the latter contained $`10`$-qubit interactions. However, the advantage of the Hamiltonian in is that it is programmable. The Hamiltonian given here contains the gates to be implemented already as hardware. In order to make a programmable quantum computer one could proceed as follows. Consider a universal quantum cellular automaton in one dimension consisting of cells $`C_j`$ with $`j`$. Each cell contains a quantum system with Hilbert space $`^d`$. Using the so-called Margolus partition scheme , one time step of the CA consists of two parts. The first one is given by identical copies of a unitary $`U`$ acting on all pairs $`(C_{2j1},C_{2j})`$, the second one by identical copies<sup>2</sup><sup>2</sup>2To avoid technical problems with infinite tensor products of the form $`\mathrm{}VV\mathrm{}`$ one could use a $`C^{}`$-algebraic description and work with automorphisms on so-called quasilocal algebras . of $`V`$ acting on all pairs $`(C_{2j},C_{2j+1})`$. In our scheme, we may represent each cell by the subspace of an appropriate number of adjacent logical qubits. Then we cover some part of our 2D lattice with a pattern of interaction regions that implement unitaries $`U`$ and $`V`$ as shown in Fig. 7.
The number of pairs of columns in this pattern determines the number of time steps of the simulated CA. The translation symmetry of such a crystal would then, however, be described by huge unit cells. Nevertheless, our construction shows that translation invariant finite-range interactions among finite dimensional quantum systems can in principle implement a universally programmable quantum computer.
## 8 Conclusions
We have shown that nearest-neighbor interactions on 2-dimensional lattices can induce dynamical evolutions that are powerful enough to represent a programmable quantum computer. Even though we have not described a realistic implementation scheme we have given interactions which are not a priori unphysical. The main reason why it may be difficult to construct systems having exactly the Hamiltonians considered here is that the spatial homogeneity of the interactions follow a rather sophisticated law.
To find even simpler Hamiltonians which can perform quantum computing in a closed physical system is an interesting challenge for further research.
## Acknowledgments
Part of this work was done during a visit of Hans Briegel’s group at IQOQI, whose hospitality is gratefully acknowledged. The model was inspired by discussions with Robert Raussendorf about cellular automata and with several people at IQOQI about optical lattices and quantum dots, especially Gregor Thalhammer and Alessio Recati. Thanks also to Pawel Wocjan for encouraging remarks and helpful discussions.
This work was partly funded by the Landesstiftung Baden-Württemberg, project AZ1.1422.01.
## Appendix A Alternative model avoiding diagonal interactions
The clocking scheme of our autonomous computer presented so far is somewhat sophisticated in the sense that the hopping terms connect lattice sites lying diagonal with respect to the grid given by the possible atom positions. In other words, the rows defining the direction of atom propagation were diagonal with respect to the square lattice structure. It would be more natural to have tunneling along rows or columns of the square lattice itself.
We will therefore describe an alternative model for the propagation of an atom chain which has, however, the disadvantage, that we can only conjecture that the atoms move in forward direction with sufficient speed and will appear outside the interaction region with reasonable probability. We can, however, prove, that the corresponding classical random walk would really implement the computation.
The whole lattice consists of $`2n\times m`$ sites. For reasons that will become clear later, we will color the sites such that we obtain a pattern similar to a chess board. As indicated in Fig. 8 the $`2n`$ atoms start in the first column. In contrast to our first model, the white sites are used, too.
The rules for the propagation of atoms along the rows are as follows. An atom on a white site is allowed to move forward or backwards by one column if the atoms in the two adjacent rows are currently in the same column. An atom on a black site is only allowed to move one column forward or backwards if it is finally located in between two adjacent atoms in the same column. In the first column, only the atoms on the white sites are allowed to move forward. A possible configuration obtained by applying these rules is shown on the left of Fig. 8.
By induction, we will argue that the configuration of atoms has always the following properties:
1. Atoms located in adjacent rows are either in the same column or in adjacent columns.
2. For every atom located on a white site there are always atoms on the two adjacent black sites in the same column. The only exception are particles on white sites in the upper-most and lower-most rows which have certainly only one black neighbor.
It is clear that these conditions are satisfied by the initial state. In order to prove that they are preserved we observe that an atom that is located on a black site can only move to a white site if the latter is adjacent to two occupied black sites. But then it will satisfy condition (2). When an atom on a white site moves to a black site, condition (1) is certainly preserved since (2) was true for its initial position. The atoms being on black sites on the boundary can only move to column $`j`$ when the atom in the row below is already in column $`j`$.
Now we present a Hamiltonian that induces these propagation rules. We start with
$$K:=\underset{k=1}{\overset{2n}{}}\underset{l=1}{\overset{m1}{}}a_{k,l}a_{k,l+1}^{}+h.c.,$$
which is a usual hopping term (as it appears in Hubbard models ) annihilating the atom at a certain position and creating it at the right or left neighboring site in the same row. Physically, one could think of an optical lattice with periodic potentials generated by standing waves which result from the superposition of counter-propagating waves . To allow tunneling in horizontal direction and to avoid it along the vertical axis one could use high potential walls between rows and low potential walls between columns. This can be achieved by a superposition of high amplitude waves in vertical direction with low amplitude waves in horizontal direction.
Now we modify the lattice such that there are two types of minima in the potential. The lower potential corresponds to the black sites and the higher potential to the white sites. By superposition of two lattices with wave length $`\lambda `$ and $`2\lambda `$ one could also generate a chess-board like alternating potential (this is called a “superlattice” in ).
We describe the alternating potential by introducing an additional term
$$H_1:=2E\underset{(k,l)B}{}N_{k,l},$$
where $`B`$ denotes the set of black sites. Furthermore we introduce an attractive interaction between atoms in adjacent rows which is only active between pairs of atoms that are in the same column. The strength of the attractive interaction has strength $`E`$, i.e., we add a term
$$H_2:=E\underset{k=1}{\overset{2n1}{}}\underset{l=1}{\overset{m}{}}N_{k,l}N_{k+1,l}.$$
This ensures that all the allowed atom configurations have the same energy: when an atom sitting on a white site moves to a black one the attractive interaction is switched off, this compensates the energy gap between the black and the white site. The same is true for an atom on a black field that moves to a white field since the attractive interaction is then switched on. For the upper-most and the lower-most rows we need to decrease the energy gap between black and white sites because there is only one attractive interaction compensating it. We achieve this by adding a potential
$$H_3:=E\left(\underset{(1,l)B}{}N_{1,l}+\underset{(2n,l)B}{}N_{2n,l}\right).$$
Given the Hamiltonian $`H_1+H_2+H_3`$, one checks easily that an allowed forward or backward motion of an atom leads to a configuration with the same energy. Moreover, one can see that every step of an atom that violates the propagation rules would lead to a state with different energy. We observe furthermore that forbidden motions would always lead to a configuration with higher energy and never to lower energy. This additional feature provides the scheme with some thermodynamical stability.
Now we add the small perturbation $`K`$ ($`E`$ is again thought to be large) to the potentials and obtain
$$\stackrel{~}{H}:=H_1+H_2+H_3+K.$$
In analogy to the discussion of our first model we obtain an effective Hamiltonian that is given by
$$H_{eff}:=P\stackrel{~}{H}P,$$
where $`P`$ is projecting onto the subspace of states having the same energy as the initial atom configuration.
The effective Hamiltonian modifies the hopping terms such that they are multiplied by an additional number operator checking the position of the atoms in adjacent rows. In other words, each term of the form $`a_{j,k}^{}a_{j,k+1}`$ as well as those of the form $`a_{j,k}a_{j,k+1}^{}`$ are only active if the two black fields are occupied which are either the current two neighbors of the considered atom (if the latter is located on a white site) or the future two neighbors (if it moves to a white site). Explicitly, we get therefore
$$\begin{array}{ccc}N& & \\ & & \\ a& & a^{}\\ & & \\ N& & \end{array}+\begin{array}{ccc}& & N\\ & & \\ a^{}& & a\\ & & \\ & & N\end{array}+h.c..$$
The operators $`N`$ of these $`4`$-local terms act on black sites in the same column.
To implement one-qubit gates in this scheme, we imprint interactions as shown in Fig. 9. The interactions $`V_1,\mathrm{},V_l`$ are chosen as in Section 3. A decisive difference to our first model is that the hopping of the atoms does not turn $`V_j`$ on immediately after $`V_{j1}`$ was turned off. Instead, the atoms pass always a configuration where all interactions are turned off. However, this is irrelevant for the holonomic scheme because the statement of Lemma 2 remains certainly true if the evolution is interspersed by time intervals where $`G(t)`$ is completely switched off.
To implement two-qubit gates, we design the spin-spin interactions as depicted in Fig. 10. Now we have described how interactions can be imprinted that would implement the desired gates given that the whole atom chain passes indeed the circuit region.
However, as already stated, it seems to be difficult to derive explicit formulas for the dynamical evolution. We will therefore only consider the corresponding classical random walk and argue that one will have at least probability $`1/2`$ to find all particles outside the circuit region in the time average. Here the circuit region is defined by the first $`k`$ columns of the chess-board where $`k`$ is chosen such that the region contains all interaction stripes. The computation result is then obtained by measuring the spin of the particles after they have left the circuit region.
The classical random walk would describe the physical dynamics in the limit of strongly incoherent atom hopping. The incoherent model is, of course inconsistent with our intention to study computation in a closed physical system. The motivation to study the classical walk is that the result gives us some small hope that the coherent walk behaves also as it is required for computation even though it may clearly be very different from the classical random walk.
Let $`C`$ be the set of allowed atom configurations on the $`2n\times m`$-chess-board and $`G`$ be the graph with nodes $`C`$. Two nodes $`c_1,c_2C`$ are adjacent if $`c_2`$ can be reached from $`c_1`$ by an allowed step of an atom. The probability distribution on the set of possible atom configurations at some time instant $`t`$ is described by a vector $`p(t):=(p_1(t),\mathrm{},p_l(t))`$ having the $`l`$ nodes of $`G`$ as indices. Then a continuous classical random walk on $`G`$ is described by
$$\frac{d}{dt}p(t)=Lp(t),$$
where $`L`$ is the graph Laplacian. Its entries are $`L_{ij}:=1`$ for $`ij`$ if $`(i,j)`$ is adjacent, and $`L_{j,j}=d_j`$ where $`d_j`$ is the degree of node $`j`$, i.e., its number of neighbors. Since $`G`$ is connected, $`\lambda _0:=0`$ is a non-degenerate eigenvalue of $`L`$ and there is hence a unique stationary distribution. The latter is clearly invariant if we replace each atom configuration with the configuration obtained by reflecting it at the vertical symmetry axis of the chess-board. The latter is defined by a line between column $`2/m`$ and $`1+2/m`$ if $`m`$, where we have assumed that $`m`$ is even for simplicity.
Therefore the probability of finding at least one atom of the chain in the right half is at least $`1/2`$. Since the chain does not tear this implies that the column index of every atom is at least $`m/22n`$. Provided that $`k<m/22n`$ we find with probability at least $`1/2`$ all atoms outside the circuit region.
## Appendix B Proof of Lemma 1:
We have a dominating Hamiltonian $`H_{pot}`$ and a weak perturbation $`K`$. In order to use a perturbation theorem of we first have to rephrase some notations used in .
Given a Hamiltonian $`H_{pot}`$ and a small perturbation $`K`$, both acting on some finite dimensional Hilbert space $``$. Let $`H_{pot}`$ have a spectral gap $`\mathrm{\Delta }`$ such that no eigenvalues lie between $`\lambda _{}=\lambda _{}\mathrm{\Delta }/2`$ and $`\lambda _+=\lambda _{}+\mathrm{\Delta }/2`$ for some $`\lambda _{}`$. Let $`_{}`$ be the eigenspace of $`H_{pot}`$ corresponding to the eigenvalues below $`\lambda _{}`$ and denote its complement by $`_+`$. Denote the corresponding spectral projections by $`P_\pm `$. For an arbitrary operator $`X`$, write $`X_\pm `$ for $`P_\pm XP_{}`$ and $`X_+`$ for $`P_+XP_+`$. We define the self-energy operator $`\mathrm{\Sigma }_{}(z)`$ for real-valued $`z`$ by
$$\mathrm{\Sigma }_{}(z):=H_{pot}+K_{}+K_+G_+(\mathrm{𝟏}K_{++}G_+)^1K_+,$$
where $`G_+`$, called the unperturbed Greens function (resolvent) in the physics literature, is defined by
$$G_+^1(z):=z\mathbf{\hspace{0.17em}1}_+H_{pot+}.$$
Using the above notations, we rephrase the following theorem which can be found as Theorem 8 in .
###### Theorem 4
Set $`\stackrel{~}{H}=H_{pot}+K`$ with $`K\mathrm{\Delta }/2`$. Let there be an effective Hamiltonian $`H_{eff}`$ with spectral width $`w_{eff}`$. We assume that $`H_{eff}=P_{}H_{eff}P_{}`$. Let $`0<ϵ<\mathrm{\Delta }`$ and furthermore
1. there is an $`r`$ such that
$$w_{eff}+2ϵr\lambda _{},$$
2. for all $`z`$ such that $`|z|r`$, $`\mathrm{\Sigma }_{}(z)H_{eff}ϵ`$.
Then the restriction $`\stackrel{~}{H}_{<\lambda _{}}`$ of $`\stackrel{~}{H}`$ to the eigenspaces with eigenvalues smaller than $`\lambda _{}`$ satisfies
$`\stackrel{~}{H}_{<\lambda _{}}H_{eff}`$ $``$ $`3(w_{eff}+ϵ)\sqrt{{\displaystyle \frac{K}{\lambda _{}+\mathrm{\Delta }/2w_{eff}ϵ}}}`$
$`+`$ $`{\displaystyle \frac{r^2ϵ}{(rw_{eff}ϵ)(rw_{eff}2ϵ)}}.`$
We set $`\lambda _{}:=E/2`$ and $`\mathrm{\Delta }:=E`$. Then the condition $`K\mathrm{\Delta }/2`$ is true since the expression in eq. (1) is a sum of less than $`n^2`$ terms of norm $`1`$ and we have hence $`K<n^2<E`$ due to the assumption $`En^6`$. For the same reason, the spectral width $`w_{eff}`$ of $`H_{eff}`$ is less than $`n^2`$. The operator $`P_{}`$ projects onto $`_c`$ and condition $`H_{eff}=P_{}H_{eff}P_{}`$ required by the theorem is satisfied by definition (see eq. (4)). We will now derive a bound on
$$\mathrm{\Sigma }_{}(z)H_{eff}$$
in order to define an appropriate $`ϵ`$ and a corresponding constant $`r`$.
Due to $`K_{}=P_{}KP_{}=H_{eff}`$ and $`H_{pot}=0`$ we have
$$\mathrm{\Sigma }_{}(z)H_{eff}=K_+G_+(\mathrm{𝟏}K_{++}G_+)^1K_+$$
The eigenvalues of $`H_{pot+}`$ are $`E,2E,3E,\mathrm{}`$ according to the number of inactive interactions. Therefore the norm of
$$G_+(z)=(z\mathbf{\hspace{0.17em}1}_++H_{pot+})^1$$
is at most $`2/E`$ for all $`z`$ with $`|z|E/2`$. With
$$(\mathrm{𝟏}K_{++}G_+)^1=\underset{n0}{}(K_{++}G_+)^n,$$
we have
$$(\mathrm{𝟏}K_{++}G_+)^1\underset{n0}{}\left(\frac{2K}{E}\right)^n=\frac{1}{12K/E}.$$
Hence we get
$$\mathrm{\Sigma }_{}(z)H_{eff}K^2G_+\frac{1}{12K/E}.$$
For $`E4K`$ we obtain with $`G_+2/E`$
$$\mathrm{\Sigma }_{}(z)H_{eff}K^2\frac{4}{E}4\frac{n^4}{E}=:ϵ.$$
We may choose $`r:=\sqrt{E}`$ in order to fulfill $`w_{eff}+2ϵr\lambda _{}`$. We may now use inequality (4) and obtain
$`\stackrel{~}{H}_{E/2}H_{eff}`$ $``$ $`3(n^2+1)\sqrt{{\displaystyle \frac{n^2}{En^21}}}`$
$`+`$ $`{\displaystyle \frac{4n^4}{(\sqrt{E}n^21)(\sqrt{E}n^22)}},`$
where we have inserted the following results and definitions from above: $`w_{eff}<n^2`$, $`\lambda _{}=E/2`$, $`\mathrm{\Delta }=E`$, $`r=\sqrt{E}`$. We have inserted $`ϵ=4n^4/E`$ only in the numerator of the second fraction and replaced it with $`1`$ (as an upper bound) at the other places. Using $`En^21>E/2`$ for $`E>n^6`$ and $`n10`$ and $`\sqrt{E}n^22>\sqrt{E}/2`$ we have
$$\stackrel{~}{H}_{E/2}H_{eff}4n^2\sqrt{\frac{2n^2}{E}}+\frac{16n^4}{E}9\frac{n^3}{\sqrt{E}},$$
where we have used $`16n^4/E<n^6/E<n^3/\sqrt{E}`$. Thus, we have obtained the desired upper bound on the norm distance between $`\stackrel{~}{H}_{E/2}`$ and $`H_{eff}`$.
## Appendix C Proof of Theorem 2
Let
$$|I:=|0\mathrm{}01\mathrm{}1_m$$
(22)
be the initial configuration when denoted in the qubit picture. Recall the beginning of section 5, where we have described the correspondence between clock states and binary words (as indicated by the sequence of symbols $`1`$ and $`0`$ at the left of Figure 1).
Then all atoms are outside of $`R`$ if and only if at least $`k`$ symbols $`1`$ have traveled from the right half of the chain to the left half, i.e., at least $`k`$ fermions are contained in the left half of the interval of length $`2m`$. We have therefore to solve a mixing problem of a “discrete free fermion gas” where all particles start in the right interval. First we show that after the time $`O(k)`$ it is likely that at least $`k`$ symbols $`1`$ can be found on the left side. We define the observable
$$𝒩:=\underset{j=1}{\overset{m}{}}P_j,$$
where $`P_j:=b_j^{}b_j`$ is the projector on the upper state of qubit $`j`$. The observable $`𝒩`$ counts the number of symbols $`1`$ on the left side. In Figure 2 this corresponds to the number of symbols $`1`$ above the row in the middle. We will estimate the probability that less than $`k`$ symbols $`1`$ have moved to the left by the Chebyshev inequality. It states that for any random variable $`X`$ we have
$$P(|XE(X)|ϵ)\frac{V(X)}{ϵ^2},$$
(23)
where $`E(X)`$ and $`V(X)`$ denote the expectation value and the variance of $`X`$, respectively. In the sequel we will consider the probability distribution of $`𝒩`$ as the distribution of such a classical random variable.
Let $`|I_t`$ be the time evolved state after time $`t`$. Then the expectation value of $`𝒩`$ after the time $`t`$ is given by
$$E_t(𝒩):=\underset{j=1}{\overset{m}{}}I_t|P_j|I_t=\underset{j=1}{\overset{m}{}}I_t|b_j^{}b_j|I_t.$$
(24)
Using the dynamics (19) and (20) we get
$`I_t|P_j|I_t`$ $`=`$ $`I_t|b_j^{}b_j|I_t={\displaystyle \underset{l=1}{\overset{2m}{}}}u_{jl;t}\overline{u}_{jl;t}I|b_l^{}b_l|I`$
$`=`$ $`{\displaystyle \underset{l=m+1}{\overset{2m}{}}}u_{jl;t}\overline{u}_{jl;t}={\displaystyle \underset{l=m+1}{\overset{2m}{}}}|u_{jl;t}|^2,`$
where we have used that $`I|b_l^{}b_l|I`$ is $`1`$ or $`0`$ depending on whether there the state $`|I`$ contains the symbol $`1`$ at this position. Each term $`|u_{jl;t}|^2`$ is the probability for a particle starting at site $`l`$ to be found at site $`j`$ when measured after the time $`t`$ in a single particle quantum walk on a linear chain of length $`2m`$. Since the time evolution of such a walk has been discussed in detail in the literature , we will only describe the implications for our model.
Consider a particle starting at site $`l`$ with
$$m+1lm+4k.$$
(26)
In the notation of eq. (22) such a fermion corresponds to one of the $`4k`$ leftmost symbols $`1`$. The assumption $`n4k`$ and hence $`m4k`$ ensures that the time interval considered in the sequel is sufficiently small compared to the time to reach the boundaries of the lattice. Then reflections at the boundaries can be neglected. We have to wait only the time $`O(k)`$ in order to achieve that the width of the wave function of a particle starting at a definite position is much larger than $`4k`$ (see ) but still smaller than the size $`n`$ of the whole computer. Then the probability of finding it on the left half is larger than $`1/3`$. Recalling that this holds true for each $`l`$ satisfying eq. (26) and that the number of expected fermions ($``$ symbols “$`1`$”) in the left half is given by the sum in eq. (24), we may choose a time instant $`tO(k)`$ such that the expectation value satisfies exactly
$$E_t(𝒩)=\frac{4}{3}k.$$
In order to estimate the variance
$$V_t(𝒩)=E_t(𝒩^2)(E_t(𝒩))^2,$$
(27)
we observe that eq. (24) implies
$$(E_t(𝒩))^2=\underset{i,jm}{}I_t|P_i|I_tI_t|P_j|I_t.$$
(28)
Moreover, we have
$$E_t(𝒩^2)=\underset{i,jm}{}I_t|P_iP_j|I_t=\underset{i,jm,ij}{}I_t|P_iP_j|I_t+E_t(𝒩),$$
(29)
where we have used
$$\underset{i=1}{\overset{m}{}}I_t|P_iP_i|I_t=\underset{i=1}{\overset{m}{}}I_t|P_i|I_t=E_t(𝒩),$$
due to eq. (24). We rewrite one summand of the first term on the right of eq. (29) as
$$I_t|P_iP_j|I_t=\underset{l,s,r,p=1}{\overset{2m}{}}u_{il;t}\overline{u}_{is;t}u_{jr;t}\overline{u}_{jp;t}I|b_l^{}b_sb_r^{}b_p|I.$$
(30)
The inner product can only be nonzero if annihilators meet creators, i.e., if either $`l=s`$ and $`r=p`$ or $`l=p`$ and $`s=r`$ or if all indices coincide. In the first case (including the third) the term is only non-vanishing for $`l=s>m`$ and $`r=p>m`$ since $`b_l^{}b_l`$ is the projection $`|11|`$ on qubit $`l`$. In the second case we must have $`l=p>m`$ and $`s=rm`$ since $`b_sb_s^{}`$ is the projection $`|00|`$ on qubit $`s`$. Hence eq. (30) becomes
$$\underset{m<l,s2m}{}u_{il;t}\overline{u}_{il;t}u_{js;t}\overline{u}_{js;t}+\underset{m<l2m,sm}{}u_{il;t}\overline{u}_{jl;t}u_{js;t}\overline{u}_{is;t}.$$
(31)
The first term coincides with $`I_t|P_i|I_tI_t|P_j|I_t`$ by eq. (C). We rewrite the second term as
$`{\displaystyle \underset{ml2m,sm}{}}u_{il;t}\overline{u}_{jl;t}u_{js;t}\overline{u}_{is;t}`$
$`=`$ $`{\displaystyle \underset{1l2m,sm}{}}u_{il;t}\overline{u}_{jl;t}u_{js;t}\overline{u}_{is}{\displaystyle \underset{1lm}{}}u_{il;t}\overline{u}_{jl;t}{\displaystyle \underset{1sm}{}}u_{js;t}\overline{u}_{is;t}`$
$`=`$ $`|{\displaystyle \underset{1lm}{}}u_{il;t}\overline{u}_{jl;t}|^2,`$
where the last equality is due to
$$\underset{1l2m}{}u_{il;t}\overline{u}_{jl;t}=\delta _{ij},$$
because $`U_t`$ is unitary (see eq. 18). Hence we have found
$$I_t|P_iP_j|I_tI_t|P_i|I_tI_t|P_j|I_tij.$$
(32)
Combining (27) and (32) with (28) straightforward computation shows
$$V_t(𝒩)\left(E_t(𝒩)\underset{im}{}(I_t|P_i|I_t)^2\right)E_t(𝒩).$$
(33)
Recall that we have chosen the time instant such that $`E_t(𝒩)=4k/3`$ and hence $`V_t(𝒩)4k/3`$ by ineq. (33). Assume we would find less than $`k`$ symbols $`1`$ on the left half. Then the random variable defined by $`𝒩`$-measurements would deviate at least $`k/3`$ from its expectation value. Hence we can apply eq. (23) with $`V_t(𝒩)4k/3`$ and $`ϵ^2=k^2/9`$ and find that the probability of finding less than $`k`$ symbols in the left half is at most $`12/k`$. This completes the proof.
## Appendix D Proof of Theorem 3
In analogy to the proof of Theorem 2 we will compute the expectation value and the variance of $`𝒩`$ in the time average state and then use the Chebyshev inequality (23).
The time average expectation value of the fermion number on the left half is given by averaging eq. (24) over all $`t`$:
$$E(𝒩):=\underset{T\mathrm{}}{lim}\frac{1}{T}_0^TE_t(𝒩)𝑑t=\underset{l=m+1}{\overset{2m}{}}\left(\underset{j=1}{\overset{m}{}}\underset{T\mathrm{}}{lim}\frac{1}{T}_0^T|u_{jl;t}|^2𝑑t\right),$$
(34)
where the last equality follows from eq. (C). Recall that we interpret each summand for $`l=m+1,\mathrm{},2m`$ as the probability of finding a particle (that has started at site $`l`$) in the left half of the chain when measured at a random time instant. We will show that it is close to $`1/2`$ up to an error in $`O(1/\sqrt{m})`$. It is known that the Hamiltonian $`S+S^{}`$ on $`^{2m}`$ generating the walk (i.e., the adjacency matrix of the “path graph $`P_{2m}`$”) has $`2m`$ different eigenvalues
$$\lambda _r=2\mathrm{cos}\frac{r\pi }{2m+1}r=1,\mathrm{},2m.$$
The eigenspaces of $`S+S^{}`$, are therefore one-dimensional and the $`r`$th eigenvector is given by
$$|e_r:=\sqrt{c}\underset{j=1}{\overset{2m}{}}\mathrm{cos}\left((j\frac{1}{2})(r1)\frac{\pi }{2m}\right)|j,$$
where the normalization factor is $`c=1/(2m)`$ for $`r=1`$ and $`c=1/m`$ for $`r1`$.
Using these eigenvectors, we may write the time average density matrix of a particle that has started at position $`l`$ as
$$\underset{r=1}{\overset{2m}{}}|e_re_r|ll|e_re_r|.$$
By evaluating the probability to be at position $`j`$ using this state we may compute the time average in eq. (34) and obtain
$$\underset{T\mathrm{}}{lim}\frac{1}{T}_0^T|u_{jl;t}|^2𝑑t=\underset{r=1}{\overset{2m}{}}|j|e_r|^2|e_r|l|^2.$$
(35)
Each eigenvector defines a probability distribution $`p_r`$ on $`\{1,\mathrm{},2m\}`$ by
$$p_r(j):=|j|e_r|^2=\frac{c}{2}\left[1+\mathrm{cos}\left((j\frac{1}{2})(r1)\frac{\pi }{m}\right)\right].$$
(36)
The spatial probability oscillations are given by waves with frequencies $`\nu _r:=(r1)\pi /m`$. We will refer to those frequencies $`\nu `$ that satisfy
$$2\pi \frac{1}{\sqrt{m}}\nu 2\pi (1\frac{1}{\sqrt{m}})$$
as high frequencies (they are neither close to $`0`$ nor close to $`2\pi `$) and to the others as low frequencies. For a high frequency wave, the probability distribution $`jp_r(j)`$ leads almost to equal probabilities for both halfs. This is seen from
$$\underset{j=1}{\overset{m}{}}p_r(j)=\frac{1}{2}+\frac{1}{m}\underset{j=1}{\overset{m}{}}\mathrm{cos}(j\nu _r\frac{\nu _r}{2}).$$
The sum over the oscillating terms can be bounded from above by observing
$$|\underset{1lm}{}\mathrm{cos}(l\nu \frac{\nu }{2})||\underset{1lm}{}e^{il\nu }|\frac{2}{|1e^{i\nu }|}O(\sqrt{m}),$$
where we have used that $`\nu `$ differs from both $`0`$ and $`2\pi `$ by $`\mathrm{\Omega }(1/\sqrt{m})`$ and the last inequality follows directly from the geometric sum formula
$$\underset{j=0}{\overset{n}{}}q^j=\frac{1+q^{n+1}}{1q}.$$
Using eq. (35) and eq. (36), the term in the bracket in eq. (34) can be written as
$$\underset{j=1}{\overset{m}{}}\underset{r=1}{\overset{2m}{}}p_r(j)|e_r|l|^2.$$
(37)
Neglecting the low frequency terms $`p_r`$ in this sum can only cause an error in $`O(1/\sqrt{m})`$ due to $`|e_r|j|^2O(1/m)`$ taking into account that there are $`O(\sqrt{m})`$ such terms. Hence the probability of finding a particle that has started at any site $`l`$ in the left half is close to $`1/2`$ up to an error in $`O(1/\sqrt{m})`$. By summation over all $`l=m+1,\mathrm{},2m`$ we find $`E(𝒩)m/2O(m/\sqrt{m})`$.
To derive bounds on the variance of $`𝒩`$ in the time average state we use eq. (29) and observe
$`E(𝒩^2)`$ $`=`$ $`{\displaystyle \underset{1i,jm,ij}{}}\text{average}\left(I_t|P_iP_j|I_t\right)+E(𝒩)`$ (38)
$``$ $`{\displaystyle \underset{1i,jm,ij}{}}\text{average}\left(I_t|P_i|I_tI_t|P_j|I_t\right)+E(𝒩),`$ (39)
where the last inequality is due to ineq. (32). We may rewrite the term in the big bracket in equation (39) using eq. (C) into
$$I_t|P_i|I_tI_t|P_jI_t=\underset{l=m+1}{\overset{2m}{}}|u_{il;t}|^2\underset{l=m+1}{\overset{2m}{}}|u_{jl;t}|^2.$$
(40)
Recall that the coefficients $`u_{il;t}`$ are the matrix entries of a unitary that describes a random walk on a linear chain. Note furthermore that the corresponding unitary group $`U_t`$ is generated by a real-symmetric Hamiltonian and we have therefore $`u_{il;t}=u_{li:t}`$. Thus, we are allowed to interpret $`|u_{il:t}|^2`$ as the probability for finding a particle at position $`l`$ that has started at position $`i`$, which interchanges the original roles of $`i`$ and $`l`$. This will be convenient in the sequel. The product of the sums on the right hand side of eq. (40) can either be interpreted as arising from two independent walks of two particles or the walk of one particle in two dimensions. Explicitly, we consider a dynamical evolution in $`^{2m}^{2m}`$ generated by the Hamiltonian
$$(S+S^{})\mathrm{𝟏}+\mathrm{𝟏}(S+S^{}).$$
(41)
Since the Hamiltonian (41) is the adjacency graph of the square lattice
$$\{1,\mathrm{},2m\}\times \{1,\mathrm{},2m\},$$
we can consider the right hand of eq. (40) as the probability of finding a particle starting at $`|i,j`$ in a quantum walk on the square lattice in the “target quadrant” $`\{m,\mathrm{},2m\}\times \{m,\mathrm{},2m\}`$. We want to prove that we have sufficient mixing in order to obtain $`1/4`$ for the time average of term (40) up to an error in $`O(1/\sqrt{m})`$. We proceed similarly as for the one-dimensional walk with the essential difference that the spectrum of the Hamiltonian (41) is degenerate since $`|e_r,e_p:=|e_r|e_p`$ and $`|e_p,e_r`$ have the same eigenvalues. We denote the projection onto their span by $`P_{r,p}`$. The rank of $`P_{r,p}`$ is $`2`$ for $`rp`$ and $`1`$ for $`r=p`$. Using these projections, we may compute the time average of (40) for fixed $`i,j`$ by
$$\underset{l,s=m+1}{\overset{2m}{}}\underset{r,p=1}{\overset{2m}{}}l,s|P_{r,p}|i,ji,j|P_{r,p}|l,s=\underset{l,s=n+1}{\overset{2m}{}}\underset{r,p=1}{\overset{2m}{}}q_{r,p}(l,s)d_{r,p},$$
(42)
where we have defined the coefficients
$$d_{r,p}:=\underset{l,s=1}{\overset{2m}{}}|l,s|P_{r,p}|i,j|^2=i,j|P_{r,p}|i,j$$
and probability measures $`q_{r,p}`$ on the square lattice by
$$q_{r,p}(l,s):=\frac{1}{d_{r,p}}|l,s|P_{r,p}|i,j|^2.$$
Defining the frequencies $`\nu _r:=(r1)\pi /n`$ and $`\nu _p:=(p1)\pi /n`$, the state vector given by normalizing $`P_{r,p}|i,j`$ for a given $`r,p`$ is a superposition of cosine wave functions with frequency vector $`(\nu _r,\nu _p)`$ with another wave having the vector $`(\nu _p,\nu _r)`$. The corresponding probability distributions $`q_{r,p}`$ contains then terms with frequencies $`2\nu _r,2\nu _p,\nu _r\nu _p,\nu _r+\nu _p`$. We refer to a term as low frequency term whenever at least one of these frequencies is low. For each “high frequency distribution” $`q_{r,p}`$ the probability for the target quadrant is $`1/4`$ up to an error in $`O(1/\sqrt{m})`$. This is seen in straightforward analogy to the corresponding argument for the one-dimensional walk. Moreover, the contribution of the low frequency terms in (42) can be bounded from above by the number of low frequency terms (which is in $`O(m\sqrt{m})`$) times the maximal coefficient $`d_{l,s}`$ of each term (which is in $`O(1/m^2)`$ due to $`|i,j|e_r,e_p|O(1/m)`$). We conclude that the total contribution of all low frequency terms to (42) is in $`O(1/\sqrt{m})`$. Hence the total sum in eq. (42) is $`1/4`$ up to an error in $`O(1/\sqrt{m})`$. We conclude
$$\underset{1i,jm,ij}{}\text{average}\left(I_t|P_i|I_tI_t|P_j|I_t\right)=\frac{m^2}{4}+O(m/\sqrt{m}).$$
Using eq. (39) we conclude
$$E(𝒩^2)=m^2/4+E(𝒩)+O(m/\sqrt{m})$$
and hence
$$E(𝒩^2)=(E(𝒩))^2+O(m).$$
Hence the variance $`E(𝒩^2)(E(𝒩))^2`$ is in $`O(m)`$ and the probability that a measurement of $`𝒩`$ leads to a result with less than $`m/4`$ converges to zero with $`O(1/m)`$ by the Chebyshev inequality. Therefore, the probability of finding less than $`m/4`$ fermions on the left half (i.e., the atom chain has left the circuit region) tends to zero, too.
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# Higher Order Slow-Roll Predictions for Inflation
## Abstract
We study the WKB approximation beyond leading order for cosmological perturbations during inflation. To first order in the slow-roll parameters, we show that an improved WKB approximation leads to analytical results agreeing to within $`0.1\%`$ with the standard slow-roll results. Moreover, the leading WKB approximation to second order in the slow-roll parameters leads to analytical predictions in qualitative agreement with those obtained by the Green’s function method.
Introduction. It is nowadays common to state that we are in the era of precision cosmology. The implication of present and future data sets will be able to discriminate among different inflationary models infla . For this reason, the comparison of inflationary models with observations requires theoretical advances in the predictions for the power spectrum of primordial perturbations beyond the lowest order in the slow-roll parameters worked out by Stewart and Lyth SL . Such slow-roll parameters quantify the deviation from an exactly exponential expansion during inflation and are related for a canonical inflaton to the derivative of the potential. Among the many different conventions, we find it convenient to use the hierarchy of horizon flow functions $`ϵ_i`$ terrero , defined as $`ϵ_{i+1}=\dot{ϵ}_i/\left(Hϵ_i\right)`$, with $`ϵ_1=\dot{H}/H^2`$, $`H=\dot{a}/a`$ the Hubble parameter, $`a`$ the scale factor and dots denote derivatives respect to the cosmic time.
The search for deviations from a simple power-law parameterization of the primordial power spectrum in the first year of WMAP data has begun peiris , with however no statistically conclusive evidence for a significant deviation. Stronger conclusions about any deviations will be made possible by the better resolution of the Planck satellite planck . In the context of inflation, the deviations from power-law spectra considered so far correspond to some of the predictions beyond the first order in the slow-roll parameters and are expected in general.
The purpose of this letter is to show how the primordial power spectrum of cosmological perturbations generated in a single field inflationary model can be analytically predicted by the WKB method with sufficient accuracy. We shall also show how two different analytical approximations, such as the Green’s function method (GFM henceforth) gongstewart and the WKB method MS ; WKB1 used here, agree to second order on the polynomial structure of the results in the horizon flow functions during inflation, but differ in the numerical coefficients of the second order terms (the same polynomial structure in the spectral indices is found by the uniform approximation HHJMP ; HHHJMP ; LA ).
Cosmological perturbations. Let us begin by recalling that scalar (density) and tensor (gravitational wave) fluctuations on a Robertson-Walker background are given respectively by $`\mu =\mu _\mathrm{S}aQ`$ ($`Q`$ is the Mukhanov variable mukh ) and $`\mu =\mu _\mathrm{T}ah`$ ($`h`$ is the amplitude of the two polarizations of gravitational waves gris ; staro ). The functions $`\mu `$ must satisfy the one-dimensional Schrödinger-like equation
$`\left[{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}\eta ^2}}+\mathrm{\Omega }^2(k,\eta )\right]\mu =0,`$ (1)
together with the initial condition
$`\underset{\frac{k}{aH}+\mathrm{}}{lim}\mu (k,\eta ){\displaystyle \frac{\mathrm{e}^{ik\eta }}{\sqrt{2k}}}.`$ (2)
In the above $`\eta `$ is the conformal time, $`k`$ is the wave-number, and
$`\mathrm{\Omega }^2(k,\eta )k^2{\displaystyle \frac{z^{\prime \prime }}{z}},`$ (3)
where $`z=z_\mathrm{S}a^2\varphi ^{}/H`$ for scalar and $`z=z_\mathrm{T}a`$ for tensor perturbations ($`\varphi `$ is the homogenous inflaton and primes denote derivatives with respect to $`\eta `$). The dimensionless power spectra of scalar and tensor fluctuations are then given by
$`𝒫_\zeta {\displaystyle \frac{k^3}{2\pi ^2}}\left|{\displaystyle \frac{\mu _\mathrm{S}}{z_\mathrm{S}}}\right|^2,𝒫_h{\displaystyle \frac{4k^3}{\pi ^2}}\left|{\displaystyle \frac{\mu _\mathrm{T}}{z_\mathrm{T}}}\right|^2`$ (4a)
and the spectral indices and runnings by
$`n_\mathrm{S}1{\displaystyle \frac{\mathrm{d}\mathrm{ln}𝒫_\zeta }{\mathrm{d}\mathrm{ln}k}}|_{k=k_{}},n_\mathrm{T}{\displaystyle \frac{\mathrm{d}\mathrm{ln}𝒫_h}{\mathrm{d}\mathrm{ln}k}}|_{k=k_{}}`$ (4b)
$`\alpha _\mathrm{S}{\displaystyle \frac{\mathrm{d}^2\mathrm{ln}𝒫_\zeta }{(\mathrm{d}\mathrm{ln}k)^2}}|_{k=k_{}},\alpha _\mathrm{T}{\displaystyle \frac{\mathrm{d}^2\mathrm{ln}𝒫_h}{(\mathrm{d}\mathrm{ln}k)^2}}|_{k=k_{}}`$ (4c)
where $`k_{}`$ is an arbitrary pivot scale. We also define the tensor-to-scalar ratio
$`R{\displaystyle \frac{𝒫_h}{𝒫_\zeta }}|_{k=k_{}}.`$ (4d)
Slow-roll and WKB approximations. One of the very few cases for which the equations for cosmological perturbations can be integrated exactly is that of power-law inflation PL ; LS , where the inflaton is a canonical scalar field with an exponential potential. In such a case, $`ϵ_1`$ is constant in time and the hierarchy of the horizon flow functions is therefore truncated with $`ϵ_i=0`$ for $`i2`$. The original slow-roll approximation corresponds to considering both $`ϵ_1`$ and $`ϵ_2`$ constant in time for any potential SL .
Although apparently inappropriate for studying physical problems such as the hydrogen atom langer (or cosmological perturbations MS ), the WKB method can be cleverly applied after suitable redefinitions of the wave-function (for Fourier modes $`k`$) and variables (with the corresponding “frequency” $`\mathrm{\Omega }`$ replaced by a new expression $`\omega `$, as given in detail in Refs. MS ; WKB1 ). This improved WKB method applied to cosmological perturbations with a linear turning point in $`\omega `$ does not however predict amplitudes with a sufficient accuracy to lowest order MS . We have shown in Ref. WKB1 how the prediction for the amplitude may be improved by using a next-to-leading WKB approximation. Our method involved an adiabatic expansion and the numerical evaluation of the higher-order coefficients which yielded the spectra for a given inflationary model (and were compared with the exact spectra of power-law inflation WKB1 ).
The power spectra to next-to-leading WKB order, in the adiabatic expansion, are given by
$`𝒫_\zeta ={\displaystyle \frac{H^2}{\pi ϵ_1m_{\mathrm{Pl}}^2}}\left({\displaystyle \frac{k}{aH}}\right)^3{\displaystyle \frac{\mathrm{e}^{2\xi _{\mathrm{II},\mathrm{S}}}\left(1+g_{(1)\mathrm{S}}^{\mathrm{AD}}\right)}{\left(1ϵ_1\right)\omega _{\mathrm{II},\mathrm{S}}}}`$
(5)
$`𝒫_h={\displaystyle \frac{16H^2}{\pi m_{\mathrm{Pl}}^2}}\left({\displaystyle \frac{k}{aH}}\right)^3{\displaystyle \frac{\mathrm{e}^{2\xi _{\mathrm{II},\mathrm{T}}}\left(1+g_{(1)\mathrm{T}}^{\mathrm{AD}}\right)}{\left(1ϵ_1\right)\omega _{\mathrm{II},\mathrm{T}}}},`$
where $`m_{\mathrm{Pl}}`$ is the Planck mass and all quantities are evaluated in the super-horizon limit. A crucial part of our method is the evaluation of $`\xi _{II}`$ and $`g_{(1)}^{\mathrm{AD}}`$. In particular, in this work we will estimate analytically next-to-leading WKB corrections to $`𝒪(ϵ_i)`$ and show that the leading WKB order can also give predictions to $`𝒪(ϵ_i^2)`$.
For the precise definition of $`g_{(1)}^{\mathrm{AD}}`$ and its evaluation we refer to Eq. (61b) and Section V of Ref. WKB1 . The quantity $`\xi _{II}`$ was also defined in Eq. (26b) of the same reference, and can in general be written as
$`\xi _{\mathrm{II}}(\eta _\mathrm{f},\eta _0)={\displaystyle _{\eta _\mathrm{f}}^{\eta _0}}\sqrt{A^2(\eta )k^2\eta ^2}{\displaystyle \frac{\mathrm{d}\eta }{\eta }},`$ (6)
with $`\eta _0`$ the time for which the integrand vanishes, $`\eta _\mathrm{f}`$ the super-horizon limit and $`A^2(\eta )`$ contains all the dependence on the horizon flow functions $`ϵ_i(\eta )`$. In a manner similar to repeated integration by parts, we can obtain general expressions valid for every $`A^2(\eta )`$ which contain terms easy to evaluate explicitly and new integrals of sufficiently high order in the $`ϵ_i(\eta )`$ so that they can be neglected. For example, on employing such a process once, we obtain the identity
$`\xi _{\mathrm{II}}(\eta _\mathrm{f},\eta _0)=\sqrt{A^2(\eta _\mathrm{f})k^2\eta _\mathrm{f}^2}`$
$`{\displaystyle \frac{A(\eta _\mathrm{f})}{2}}\mathrm{ln}\left[{\displaystyle \frac{A(\eta _\mathrm{f})\sqrt{A^2(\eta _\mathrm{f})k^2\eta _\mathrm{f}^2}}{A(\eta _\mathrm{f})+\sqrt{A^2(\eta _\mathrm{f})k^2\eta _\mathrm{f}^2}}}\right]`$
$`{\displaystyle _{\eta _\mathrm{f}}^{\eta _0}}{\displaystyle \frac{A^{}(\eta )}{2}}\mathrm{ln}\left[{\displaystyle \frac{A(\eta )\sqrt{A^2(\eta )k^2\eta ^2}}{A(\eta )+\sqrt{A^2(\eta )k^2\eta ^2}}}\right]d\eta .`$ (7)
Since the integral in the right hand side contains $`A^{}`$, it is of (at least) one order higher in the $`ϵ_i`$ than the remaining terms, which can be calculated explicitly. More details, omitted here for the sake of brevity, will be given in a forthcoming paper sequel .
Next-to-leading WKB order and first slow-roll order. The adiabatic corrections $`g_{(1)}^{\mathrm{AD}}`$ in Eqs. (5) to leading order in the horizon flow functions are given by
$`g_{(1)\mathrm{S}}^{\mathrm{AD}}={\displaystyle \frac{37}{324}}{\displaystyle \frac{19}{243}}\left(ϵ_1+{\displaystyle \frac{1}{2}}ϵ_2\right)`$
(8)
$`g_{(1)\mathrm{T}}^{\mathrm{AD}}={\displaystyle \frac{37}{324}}{\displaystyle \frac{19}{243}}ϵ_1.`$
We can now write the expressions for the scalar and tensor spectra to next-to-leading WKB order (indicated by the subscript WKB$``$) and first slow-roll order (indicated by the superscript (1))
$`𝒫_{\zeta ,\mathrm{WKB}}^{(1)}={\displaystyle \frac{H^2}{\pi ϵ_1m_{\mathrm{Pl}}^2}}A_{\mathrm{WKB}}\left[12\left(D_{\mathrm{WKB}}+1\right)ϵ_1D_{\mathrm{WKB}}ϵ_2\left(2ϵ_1+ϵ_2\right)\mathrm{ln}\left({\displaystyle \frac{k}{k_{}}}\right)\right]`$
(9a)
$`𝒫_{h,\mathrm{WKB}}^{(1)}={\displaystyle \frac{16H^2}{\pi m_{\mathrm{Pl}}^2}}A_{\mathrm{WKB}}\left[12\left(D_{\mathrm{WKB}}+1\right)ϵ_12ϵ_1\mathrm{ln}\left({\displaystyle \frac{k}{k_{}}}\right)\right],`$
where $`A_{\mathrm{WKB}}=361/18e^30.999`$ and $`D_{\mathrm{WKB}}\frac{7}{19}\mathrm{ln}30.7302`$. The slow-roll approximation SL predicts, for the corresponding quantities, $`A_{\mathrm{SR}}=1`$ and $`D_{\mathrm{SR}}=C\gamma _E+\mathrm{ln}220.7296`$ (where $`\gamma _E`$ is the Euler-Mascheroni constant). Thus, the next-to-leading WKB order gives an error of about $`0.1\%`$ for the estimate of the amplitude and one of about $`0.08\%`$ on the coefficient $`C`$. This shows that analytical results obtained by the WKB method have reached the same accuracy as the standard slow-roll approximation. We also obtain the same slow-roll spectral indices and $`\alpha `$-runnings,
$`n_{\mathrm{S},\mathrm{WKB}}^{(1)}1=2ϵ_1ϵ_2,n_{\mathrm{T},\mathrm{WKB}}^{(1)}=2ϵ_1`$ (9b)
$`\alpha _{\mathrm{S},\mathrm{WKB}}^{(1)}=\alpha _{\mathrm{T},\mathrm{WKB}}^{(1)}=0,`$ (9c)
on using respectively Eqs. (4b) and (4c). From Eq. (4d) the tensor-to-scalar ratio becomes
$`R_{\mathrm{WKB}}^{(1)}=16ϵ_1\left(1+D_{\mathrm{WKB}}ϵ_2\right).`$ (9d)
Leading WKB order and second slow-roll order. We would now like to increase our accuracy in the slow-roll parameters to second order, while keeping the WKB approximation to leading order (a further improved treatment to second order in the slow-roll parameters is in progress sequel ). On setting $`g_{(1)}^{\mathrm{AD}}=0`$ in Eqs. (5), we can write the expressions for the scalar and tensor perturbations to leading WKB order (indicated by the subscript WKB), and second slow-roll order (indicated with the superscript (2)) as
$`𝒫_{\zeta ,\mathrm{WKB}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{H^2}{\pi ϵ_1m_{\mathrm{Pl}}^2}}A_{\mathrm{WKB}}\{12(D_{\mathrm{WKB}}+1)ϵ_1D_{\mathrm{WKB}}ϵ_2+(2D_{\mathrm{WKB}}^2+2D_{\mathrm{WKB}}{\displaystyle \frac{1}{9}})ϵ_1^2`$
$`+\left(D_{\mathrm{WKB}}^2D_{\mathrm{WKB}}+{\displaystyle \frac{\pi ^2}{12}}{\displaystyle \frac{20}{9}}\right)ϵ_1ϵ_2+\left({\displaystyle \frac{1}{2}}D_{\mathrm{WKB}}^2+{\displaystyle \frac{2}{9}}\right)ϵ_2^2+\left({\displaystyle \frac{1}{2}}D_{\mathrm{WKB}}^2+{\displaystyle \frac{\pi ^2}{24}}{\displaystyle \frac{1}{18}}\right)ϵ_2ϵ_3`$
$`+\left[2ϵ_1ϵ_2+2\left(2D_{\mathrm{WKB}}+1\right)ϵ_1^2+\left(2D_{\mathrm{WKB}}1\right)ϵ_1ϵ_2+D_{\mathrm{WKB}}ϵ_2^2D_{\mathrm{WKB}}ϵ_2ϵ_3\right]\mathrm{ln}\left({\displaystyle \frac{k}{k_{}}}\right)`$
$`+{\displaystyle \frac{1}{2}}(4ϵ_1^2+2ϵ_1ϵ_2+ϵ_2^2ϵ_2ϵ_3)\mathrm{ln}^2\left({\displaystyle \frac{k}{k_{}}}\right)\}`$
$`𝒫_{h,\mathrm{WKB}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{16H^2}{\pi m_{\mathrm{Pl}}^2}}A_{\mathrm{WKB}}\{12(D_{\mathrm{WKB}}+1)ϵ_1+(2D_{\mathrm{WKB}}^2+2D_{\mathrm{WKB}}{\displaystyle \frac{1}{9}})ϵ_1^2`$
$`+\left(D_{\mathrm{WKB}}^22D_{\mathrm{WKB}}+{\displaystyle \frac{\pi ^2}{12}}{\displaystyle \frac{19}{9}}\right)ϵ_1ϵ_2+\left[2ϵ_1+2\left(2D_{\mathrm{WKB}}+1\right)ϵ_1^22\left(D_{\mathrm{WKB}}+1\right)ϵ_1ϵ_2\right]\mathrm{ln}\left({\displaystyle \frac{k}{k_{}}}\right)`$
$`+{\displaystyle \frac{1}{2}}(4ϵ_1^22ϵ_1ϵ_2)\mathrm{ln}^2\left({\displaystyle \frac{k}{k_{}}}\right)\},`$
where $`A_{\mathrm{WKB}}=18/e^30.896`$ and $`D_{\mathrm{WKB}}\frac{1}{3}\mathrm{ln}30.7653`$. The spectral indices (4b) are then given by
$`n_{\mathrm{S},\mathrm{WKB}}^{(2)}1`$ $`=`$ $`2ϵ_1ϵ_22ϵ_1^2\left(2D_{\mathrm{WKB}}+3\right)ϵ_1ϵ_2`$ (10b)
$`D_{\mathrm{WKB}}ϵ_2ϵ_3`$
$`n_{\mathrm{T},\mathrm{WKB}}^{(2)}`$ $`=`$ $`2ϵ_12ϵ_1^22\left(D_{\mathrm{WKB}}+1\right)ϵ_1ϵ_2`$
and their runnings (4c) by
$`\alpha _{\mathrm{S},\mathrm{WKB}}^{(2)}=2ϵ_1ϵ_2ϵ_2ϵ_3,\alpha _{\mathrm{T},\mathrm{WKB}}^{(2)}=2ϵ_1ϵ_2.`$ (10c)
We note that the second order correction $`2ϵ_1^2`$ to $`n_\mathrm{S}`$ and $`n_\mathrm{T}`$ agrees with the exact spectral index for power-law inflation <sup>1</sup><sup>1</sup>1For power-law inflation $`ϵ_1=1/p`$, $`n_S=n_T`$ and the exact spectral index is $`n_T=12/(p1)`$. For large $`p`$, $`n_T=12/p2/p^2+𝒪(p^3)`$. See LA for the most recent result based on the uniform approximation which contains the correct term $`2ϵ_1^2`$ in $`n_T`$.. The tensor-to-scalar ratio (4d) becomes
$`R_{\mathrm{WKB}}^{(2)}`$ $`=`$ $`16ϵ_1[1+D_{\mathrm{WKB}}ϵ_2+(D_{\mathrm{WKB}}+{\displaystyle \frac{1}{9}})ϵ_1ϵ_2`$ (10d)
$`+\left({\displaystyle \frac{1}{2}}D_{\mathrm{WKB}}^2{\displaystyle \frac{2}{9}}\right)ϵ_2^2`$
$`+({\displaystyle \frac{1}{2}}D_{\mathrm{WKB}}^2{\displaystyle \frac{\pi ^2}{24}}+{\displaystyle \frac{1}{18}})ϵ_2ϵ_3].`$
Of course, Eqs. (LABEL:P\_SlowRoll\_0\_2order)–(10c) give the same results as obtained in Ref. MS , to first order in the horizon flow functions. As already found in MS , the leading WKB order gives an error of about $`10\%`$ for the estimate of the amplitude and one of about $`5\%`$ on $`C`$ with respect to standard slow-roll results. The runnings $`\alpha _\mathrm{S}`$ and $`\alpha _\mathrm{T}`$ are predicted to be of $`𝒪(ϵ_i^2)`$ KT , and in agreement with those obtained by the GFM gongstewart ; LLMS . It is notworthy that our power spectra and spectral indices have the same polynomial structure in the $`ϵ_i`$ as those of the latter references. For the spectral indices $`n_\mathrm{S}`$ and $`n_\mathrm{T}`$ an analogous structure is also confirmed by the uniform approximation HHJMP ; HHHJMP ; LA <sup>2</sup><sup>2</sup>2An agreement on the polynomial structure, but not in the numerical coefficients, also holds for the spectral indices $`n_\mathrm{S}`$ and $`n_\mathrm{T}`$, which are the present predictions of the uniform approximation HHJMP ; HHHJMP ; LA ..
Conclusions. We have shown that inflationary theoretical predictions have now reached expressions to second order in the slow-roll parameters confirmed by two completely different approximation schemes, such as the GFM and WKB. As can be seen from Fig. 1, for some inflationary models, second order slow-roll corrections to the power spectrum are necessary to perform a correct comparison between theoretical predictions and observational data. The figure also shows that not only the runnings are important, but also the $`𝒪(ϵ_i^2)`$ terms in the spectral indices.
The different predictions of the two methods for the numerical coefficients in front of the $`𝒪(ϵ_i^2)`$ terms are at most of the order of $`5\%`$ for the spectral indices and of $`10\%`$ in the amplitudes. For slow-roll parameters $`ϵ_i0.1`$, this leads to an accuracy of about $`0.5\%`$ in the theoretical predictions for the tensor-to-scalar ratio (10d). A similar accuracy (of $`0.1\%`$) has also been reached by the Boltzmann codes cmb . The predictions in the numerical coefficients and in the amplitude can be further improved by employing the next-to-leading WKB method sequel while including terms to second order in the slow-roll parameters.
###### Acknowledgements.
We would like to thank Salman Habib, Katrin Heitmann and Jerome Martin for discussions and comments.
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# Distributed Selfish Load BalancingPreliminary version in proceedings of ACM-SIAM Symposium on Discrete Algorithms (SODA) 2006. This work was partially supported by the EPSRC grants “Discontinuous Behaviour in the Complexity of Randomized Algorithms” and “Algorithmics of Network-sharing Games”, and by the Natural Sciences and Engineering Research Council of Canada (NSERC) discovery grant 250284-2002.
## 1 Introduction
Suppose that a consumer learns the price she would be charged by some domestic power supplier other than the one she is currently using. It is plausible that if the alternative price is lower than the price she is currently paying, then there is some possibility that she will switch to the new power supplier. Furthermore, she is more likely to switch if the ratio of current price to new price is large. If there is only a small saving, then it becomes unattractive to make the switch, since an influx of new business (oneself and other consumers) may drive up the price of the new power supplier and make it no longer competitive.
We study a simple mathematical model of the above natural rule, in the context of a load balancing (or task allocation) scenario that has received a lot of recent attention. We assume the presence of many individual users who may assign their tasks to chosen resources. The users are selfish in the sense that they attempt to optimize their own situation, i.e., try to assign their tasks to minimally loaded resources, without trying to optimize the global situation. In general, a Nash equilibrium among a set of selfish users is a state in which no user has the incentive to change her current decision. In our setting, this corresponds to no user having an incentive to reallocate their task to some other resource. An $`ϵ`$-Nash equilibrium is a standard notion of an approximate Nash equilibrium, and is a state in which no user can reduce her cost by a multiplicative factor of less than $`1ϵ`$ by changing action. Here we do not focus on the quality of equilibria, but rather on the (perhaps more algorithmic) question of convergence time to such a state.
We assume a strongly distributed and concurrent setting, i.e., there is no centralized control mechanism whatsoever, and all users may choose to reallocate their tasks at the same time. Thus, we do not (and cannot) use the Elementary Step System (discussed in more detail in the next section), where the assumption is that at most one user may reallocate her task at any given stage.
Throughout we let $`m`$ denote the number of tasks (in the above discussion, customers) and $`n`$ the number of resources (power suppliers). As hinted in the above discussion, we assume that typically $`mn`$. In a single time step (or round) each task does the following. Let $`i`$ be the resource currently being used by the task. Select $`j`$ uniformly at random from $`\{1,\mathrm{},n\}`$ and find the load of resource $`j`$. Let $`X_i`$ and $`X_j`$ be the loads of resources $`i`$ and $`j`$ respectively. If $`X_j<X_i`$, migrate from $`i`$ to $`j`$ with a probability of $`1X_j/X_i`$; the transition from round $`t`$ to round $`t+1`$ is given in Figure 1. Notice that if we had unconditional migrations, i.e., without an additional coin flip (move only with probability $`1X_j(t)/X_i(t)`$), then this may lead to an unstable system; consider for example the case $`m=2`$ with initially most tasks assigned to one of the resources: the overload would oscillate between the two resources, with a load ratio tending towards 2:1. (This observation about the risk of oscillation has also been made in similar contexts in , and we will not elaborate on it further.)
It can easily be seen that, if all tasks use the above policy, then the expected load of every resource at the next step is $`m/n`$:
###### Observation 1.
Regardless of the load distribution at time step $`t`$, the expected load of every resource at the next step is $`m/n`$.
###### Proof.
To see this, assume that the loads $`X_i(t)`$ are arranged in descending order so that $`X_j(t)X_{j+1}(t)`$ and note that
$$𝔼[X_i(t+1)]=X_i(t)+\underset{\mathrm{}=1}{\overset{i1}{}}\frac{1}{n}X_{\mathrm{}}(t)\left(1\frac{X_i(t)}{X_{\mathrm{}}(t)}\right)\underset{\mathrm{}=i+1}{\overset{n}{}}\frac{1}{n}X_i(t)\left(1\frac{X_{\mathrm{}}(t)}{X_i(t)}\right)$$
$$=X_i(t)+\frac{1}{n}\underset{\mathrm{}=1}{\overset{i1}{}}(X_{\mathrm{}}(t)X_i(t))\frac{1}{n}\underset{\mathrm{}=i+1}{\overset{n}{}}(X_i(t)X_{\mathrm{}}(t))$$
$$=X_i(t)+\frac{1}{n}\underset{\mathrm{}=1}{\overset{n}{}}(X_{\mathrm{}}(t)X_i(t))=\frac{1}{n}\underset{\mathrm{}=1}{\overset{n}{}}X_{\mathrm{}}(t)=\frac{m}{n}.$$
This provides a compelling motivation for the policy, which is that as a result, no task has an incentive to deviate unilaterally from this policy. This implies that in the terminology of it is a Nash rerouting policy. It is also a simple regret-minimizing policy in the sense of since the average cost of resources used by an agent is no higher than the best choice of a single resource to be used repeatedly. Although the above rule is very natural and has the nice properties described above, we show that it may take a long time to converge to a perfectly balanced allocation of tasks to resources. We address this problem as follows. Define a neutral move to be a task migration from a resource with load $`\mathrm{}`$ at time $`t`$ to a resource with load $`\mathrm{}1`$ at time $`t`$ (so, if no other task migrates, then the cost to the migrating task is unchanged.) We consider a modification in which neutral moves are specifically disallowed (see Figure 2). That seemingly-minor change ensures fast convergence from an almost balanced state to a perfectly-balanced state. To summarize, here are the most important features of the modified protocol:
* We do not need any global information whatsoever (apart from the number of available resources); in particular, a task does not need to know the total number of tasks in the system. Also, it is strongly distributed and concurrent. If additional tasks were to enter the system, it would rapidly converge once again, with no outside intervention.
* A migrating task needs to query the load of only one other resource (thus, doing a constant amount of work in each round).
* When a task finds a resource with a significantly smaller load (that is, a load that is smaller by at least two), the migration policy is exactly the same as that used by the Nash rerouting policy of Figure 1, so the incentive is to use that probability.
* When a task finds a resource with a load that is smaller by exactly one unit, the migration policy is sufficiently close to the Nash rerouting policy that the difference in expected load is at most one, and there is little incentive to deviate.
* The protocol is simple (as well as provably efficient) enough to convince users to actually stick to it.
### 1.1 Related Work
We are studying a simple kind of congestion game. In their general form, congestion games specify a set of agents, and a set of resources, and for each agent, a set of allowed strategies, where a strategy is the selection of a subset of the resources (in this paper, any singleton subset is allowed). The cost of a resource is a non-decreasing function of the number of agents using it, and the cost for an agent is the sum of the costs of resources it uses. A classical result due to Rosenthal is that pure Nash equilibria (NE) always exist for congestion games, and this is shown by exhibiting a potential function; they are a type of potential game . The potential function also establishes that pure NE can be found via sequences of “better-response” moves, in which agents repeatedly switch to lower-cost strategies. The potential function we use later in this paper is the one of , modulo a linear re-scaling.
These results do not show how to find Nash equilibrium efficiently, the problem being that in the worst case, sequences of these self-improving moves may be exponentially-long. The following questions arise: when can NE be found by any efficient algorithm, and if so, whether it can be found via an algorithm that purports to be a realistic model of agents’ behavior. Regarding the first of these questions, the answer is no in the general setting (the problem is PLS-complete for general congestion games , see also ). PLS-completeness (introduced in ) is a generally-accepted criterion for intractability of computational problems in which we seek a local optimum of a given objective function.
However, due to the basic fact of that pure NE are sure to result from a sufficiently long better-response sequence, many algorithms for finding them are based on such sequences. An important sub-class is the Elementary step system (ESS), proposed in Orda et al. , which consists of best-response moves (where a migrating agent switches not to any improved choice, but to one that is optimal at the time of migration). For matroid games (a class of congestion games that includes the ones we consider here), Ackermann et al. show that best-response sequences must have length polynomial in the number of players, resources, and maximal rank of the matroids. In this paper we consider the special case of singleton congestion games (where players’ strategies are always single resources, thus the ranks of the matroids is 1). For these games, Ieong et al. give polynomial bounds for best-response and better-response sequences. Chien and Sinclair study a version of the ESS in the context of approximate Nash equilibria, and show that in some cases the $`ϵ`$-Nash dynamics may find an $`ϵ`$-NE where finding an exact NE is PLS-complete. Mirrokni and Vetta study the convergence rate of the ESS to solutions, and the quality of the approximation after limited iterations.
While best- and better-response dynamics are a plausible model of selfish behaviour, the associated algorithms typically require that migrations be done one-by-one, and another common assumption is that best (not better) responses are always selected. This means that to some extent, agents are being assumed to be governed by a centralized algorithm that finds a NE, and raises the question of what sort of distributed algorithms can do so, especially if agents have limited information about the state of the system (and so may not be able to find best responses). That issue is of central importance to us in this paper. Goldberg studied situations where simple better-response approaches can be realised as weakly distributed algorithms (where each agent looks for moves independently of the others, but it is assumed that moves take place consecutively, not simultaneously). In a strongly distributed setting (as we study here), where moves may occur simultaneously, we need to address the possibility that a change of strategy may increase an agent’s cost. It may happen that after a best response has been identified, it is not optimal at the time it is executed. Even-Dar and Mansour consider concurrent, independent rerouting decisions where tasks are allowed to migrate from overloaded to underloaded resources. Their rerouting process terminates in expected $`O(\mathrm{log}\mathrm{log}m+\mathrm{log}n)`$ rounds when the system reaches a Nash equilibrium. Note that their convergence rate as a function of the number $`n`$ of resources is faster than the one we obtain in this paper. The reason is that is requires agents to have a certain amount of global knowledge. A task is required to know whether its resource is overloaded (having above-average load) and tasks on underloaded resources do not migrate at all. Our rerouting policy does not require that agents know anything other that their current resource load, and the load of a randomly-chosen alternative. Even-Dar and Mansour also present a general framework that can be used to show a logarithmic convergence rate for a wide class of rerouting strategies. Our protocol does not fall into that class, since we do not require migrations to occur only from overloaded resources. Note that our lower bound is linear in $`n`$ (thus, more than logarithmic).
Distributed algorithms have been studied in the Wardrop setting (the limit of infinitely many agents), for which recent work has also extensively studied the coordination ratio . Fischer et al. investigate convergence to Wardrop equilibria for games where agents select paths through a shared network to route their traffic. (Singleton games correspond to a network of parallel links.) Their re-routing strategies are slightly different to ours — they assume that in each round, an agent queries a path with probability proportional to the traffic on that path. Here we assume paths (individual elements of a set of parallel links) are queried uniformly at random, so that agents can be assumed to have minimal knowledge. As in this paper, the probability of switching to a better path depends on the latency difference, and care has to be taken to avoid oscillation. Also in the Wardrop setting, Blum et al. show that approximate NE is the outcome of regret-minimizing rerouting strategies, in which an agent’s cost, averaged over time, should approximate the cost of the best individual link available to that agent.
Certain generalisations of singleton games have also been considered. These generalisations are not strictly congestion games according to the standard definition we gave above, but many ideas carry over. One version introduced by Koutsoupias and Papadimitriou has been studied extensively in different contexts (for example ). In this generalisation, each task may have a numerical weight (sometimes called traffic, or demand) and each resource has a speed (or capacity). The cost of using a resource is the total weight of tasks using it, divided by its speed. Even-Dar et al. give a generalized version of the potential function of that applies to these games, and which was subsequently used in . For these games however, it seems harder to find polynomial-length best-response sequences. Feldman et al. show how a sequence of steps may lead to NE, under the weaker condition that the maximal cost experienced by agents must not increase, but individual steps need not necessarily be “selfish”. They also note that poorly-chosen better-response moves may lead to an exponential convergence rate. Another generalisation of singleton games is player-specific cost functions , which allow different agents to have different cost functions for the same resource. In this setting there is no potential function and better-response dynamics may cycle, although it remains the case that pure NE always exist.
Our rerouting strategy is also related to reallocation processes for balls into bins games. The goal of a balls into bins game is to allocate $`m`$ balls as evenly as possible into $`n`$ bins. It is well-known that a fairly even distribution can be achieved if every ball is allowed to randomly choose $`d`$ bins and then the ball is allocated to the least loaded amongst the chosen bins (see for an overview). Czumaj et al. consider such an allocation where each ball initially chooses two bins. They show that, in a polynomial number of steps, the reallocation process ends up in a state with maximum load at most $`m/n+1`$. Sanders et al. show that a maximum load of $`m/n+1`$ is optimal if every ball is restricted to two random choices.
In conclusion, this paper sits at one end of a spectrum in which we study a very simple load-balancing game, but we seek solutions in a very adverse setting in which agents have, at any point in time, a minimal amount of information about the state of their environment, and carry out actions simultaneously in a strongly distributed sense.
### 1.2 Overview of our results
Section 3 deals with upper bounds on convergence time. The main result, Theorem 2, is that the protocol of Figure 2 converges to a Nash equilibrium within expected time $`O(\mathrm{log}\mathrm{log}m+n^4)`$.
The proof of Theorem 2 shows that the system becomes approximately balanced very rapidly. Specifically, Corollary 12 shows that if $`nm^{1/3}`$, then for all $`ϵ`$, either version of the distributed protocol (with or without neutral moves allowed) attains an $`ϵ`$-Nash equilibrium (where all load ratios are within $`[1ϵ,1+ϵ]`$; we use $`ϵ`$ to denote a multiplicative factor as in ) in expected $`O(\mathrm{log}\mathrm{log}m)`$ rounds. The rest of Section 3 analyses the protocol of Figure 2. It is shown that within an additional $`O(n^4)`$ rounds the system becomes optimally balanced.
In Section 4, we provide two lower bound results. The first one, Theorem 19, shows that the first protocol (of Figure 1, including moves that do not necessarily yield a strict improvement for an individual task but allow for simply “neutral” moves as well, results in exponential (in $`n`$) expected convergence time. Finally, in Theorem 20 we provide a general lower bound (regardless of which of the two protocols is being used) on the expected convergence time of $`\mathrm{\Omega }(\mathrm{log}\mathrm{log}m)`$. This lower bound matches the upper bound as a function of $`m`$.
## 2 Notation
There are $`m`$ tasks and $`n`$ resources. An assignment of tasks to resources is represented as a vector $`(x_1,\mathrm{},x_n)`$ in which $`x_i`$ denotes the number of tasks that are assigned to resource $`i`$. In the remainder of this paper, $`[n]`$ denotes $`\{1,\mathrm{},n\}`$. The assignment is a Nash equilibrium if for all $`i[n]`$ and $`j[n]`$, $`|x_ix_j|1`$. We study a distributed process for constructing a Nash equilibrium. The states of the process, $`X(0),X(1),\mathrm{}`$, are assignments. The transition from state $`X(t)=(X_1(t),\mathrm{},X_n(t))`$ to state $`X(t+1)`$ is given by the greedy distributed protocol in Figure 2.
Note that if $`X(t)`$ is a Nash equilibrium, then $`X(t+1)=X(t)`$ so the assignment stops changing. Here is a formal description of the transition from a state $`X(t)=x`$. Independently, for every $`i[n]`$, let $`(Y_{i,1}(x),\mathrm{},Y_{i,n}(x))`$ be a random variable drawn from a multinomial distribution with the constraint $`_{j=1}^nY_{i,j}(x)=x_i`$. ($`Y_{ij}`$ represents the number of migrations from $`i`$ to $`j`$ in a round.) The corresponding probabilities $`(p_{i,1}(x),\mathrm{},p_{i,n}(x))`$ are given by
$$p_{i,j}(x)=\{\begin{array}{cc}\frac{1}{n}\left(1\frac{x_j}{x_i}\right)\hfill & \text{if }x_i>x_j+1\text{,}\hfill \\ 0\hfill & \text{if }ij\text{ but }x_ix_j+1\text{,}\hfill \\ 1_{ji}p_{i,j}(x)\hfill & \text{if }i=j.\hfill \end{array}$$
Then $`X_i(t+1)=_{\mathrm{}=1}^nY_{\mathrm{},i}(x)`$.
For any assignment $`x=(x_1,\mathrm{},x_n)`$, let $`\overline{x}=\frac{1}{n}_{i=1}^nx_i`$. We define the potential function $`\mathrm{\Phi }(x)=_{i=1}^n(x_i\overline{x})^2`$. Note that $`\mathrm{\Phi }(x)=_{i=1}^nx_i^2n\overline{x}^2`$, and that a single selfish move reduces the potential.
## 3 Upper bound on convergence time
Our main result is the following.
###### Theorem 2.
Let $`T`$ be the number of rounds taken by the protocol of Figure 2 to reach a Nash equilibrium for the first time. Then $`𝔼[T]=O(\mathrm{log}\mathrm{log}m+n^4)`$.
The proof of this theorem proceeds as follows. First (Lemma 7) we give an upper bound on $`𝔼[\mathrm{\Phi }(X(t))]`$ which implies (Corollary 11) that there is a $`\tau =O(\mathrm{log}\mathrm{log}m)`$ such that, with high probability, $`\mathrm{\Phi }(X(\tau ))=O(n)`$. We also show (Observation 6 and Corollary 15) that $`\mathrm{\Phi }(X(t))`$ is a super-martingale and (Lemma 16) that it has enough variance. Using these facts, we obtain the upper bound on the convergence time.
Definition: Let $`S_i(x)=\{jx_j<x_i1\}`$. $`S_i(x)`$ is the set of resources that are significantly smaller than resource $`i`$ in state $`x`$ (in the sense that their loads are at least two tasks smaller than the load of resource $`i`$). Similarly, let $`L_i(x)=\{jx_j>x_i+1\}`$ and let $`d_i(x)=\frac{1}{n}_{j:|x_ix_j|1}(x_ix_j)`$.
###### Observation 3.
$`𝔼[X_i(t+1)X(t)=x]=\overline{x}+d_i(x)`$.
###### Proof.
$`𝔼[X_i(t+1)X(t)=x]`$ $`={\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}𝔼[Y_{\mathrm{},i}(x)]={\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}x_{\mathrm{}}p_{\mathrm{},i}(x)`$
$`={\displaystyle \underset{\mathrm{}L_i(x)}{}}x_{\mathrm{}}{\displaystyle \frac{1}{n}}\left(1{\displaystyle \frac{x_i}{x_{\mathrm{}}}}\right)+x_i\left(1{\displaystyle \underset{jS_i(x)}{}}{\displaystyle \frac{1}{n}}\left(1{\displaystyle \frac{x_j}{x_i}}\right)\right)`$
$`=x_i+{\displaystyle \frac{1}{n}}\left({\displaystyle \underset{\mathrm{}L_i(x)}{}}(x_{\mathrm{}}x_i){\displaystyle \underset{jS_i(x)}{}}(x_ix_j)\right)=x_i+{\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)S_i(x)}{}}(x_{\mathrm{}}x_i)`$
$`=x_i+{\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}(x_{\mathrm{}}x_i){\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)S_i(x)}{}}(x_{\mathrm{}}x_i)`$
$`=\overline{x}{\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)S_i(x)}{}}(x_{\mathrm{}}x_i)`$
$`=\overline{x}+{\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)S_i(x)}{}}(x_ix_{\mathrm{}}).`$
###### Observation 4.
$`_{i=1}^n(𝔼[X_i(t+1)X(t)=x])^2=n\overline{x}^2+_{i=1}^nd_i(x)^2`$.
###### Proof.
Using Observation 3,
$$\underset{i=1}{\overset{n}{}}(𝔼[X_i(t+1)X(t)=x])^2=\underset{i=1}{\overset{n}{}}(\overline{x}+d_i(x))^2=n\overline{x}^2+2\overline{x}\underset{i=1}{\overset{n}{}}d_i(x)+\underset{i=1}{\overset{n}{}}d_i(x)^2,$$
and the second term is zero since $`d_i(x)=𝔼[X_i(t+1)X(t)=x]\overline{x}`$. ∎
###### Observation 5.
$`\mathrm{var}[X_i(t+1)X(t)=x]\frac{1}{n}_{\mathrm{}L_i(x)}\left(x_{\mathrm{}}x_i\right)+\frac{1}{n}_{jS_i(x)}\left(x_ix_j\right)`$.
###### Proof.
$`\mathrm{var}(X_i(t+1)X(t)=x)`$ $`={\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}\mathrm{var}(Y_{\mathrm{},i}(x))={\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}x_{\mathrm{}}p_{\mathrm{},i}(x)(1p_{\mathrm{},i}(x))`$
$`={\displaystyle \underset{\mathrm{}L_i(x)}{}}x_{\mathrm{}}{\displaystyle \frac{1}{n}}\left(1{\displaystyle \frac{x_i}{x_{\mathrm{}}}}\right)(1p_{\mathrm{},i}(x))+x_ip_{i,i}(x)\left({\displaystyle \underset{jS_i(x)}{}}{\displaystyle \frac{1}{n}}\left(1{\displaystyle \frac{x_j}{x_i}}\right)\right)`$
$`={\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)}{}}(x_{\mathrm{}}x_i)(1p_{\mathrm{},i}(x))+p_{i,i}(x){\displaystyle \frac{1}{n}}{\displaystyle \underset{jS_i(x)}{}}(x_ix_j)`$
$`{\displaystyle \frac{1}{n}}{\displaystyle \underset{\mathrm{}L_i(x)}{}}(x_{\mathrm{}}x_i)+{\displaystyle \frac{1}{n}}{\displaystyle \underset{jS_i(x)}{}}(x_ix_j).`$
Definition: For any assignment $`x`$, let $`s_i(x)=|\{jx_j=x_i1\}|`$ and $`l_i(x)=|\{jx_j=x_i+1\}|`$. Let $`u_1(x)=_{i=1}^n_{j[n]:|x_ix_j|>1}|x_ix_j|`$ and $`u_2(x)=_{i=1}^n(s_i(x)l_i(x))^2`$. Let $`u(x)=u_1(x)/n+u_2(x)/n^2`$. We will show that $`u(x)`$ is on upper bound on the expected potential after one step, starting from state $`x`$. The quantity $`u_1(x)`$ corresponds to the contribution arising from the sum of the variances of the individual loads and $`u_2(x)`$ corresponds to the rest.
###### Observation 6.
$`𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]u(x)`$.
###### Proof.
$`𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]+n\overline{x}^2`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}𝔼[X_i(t+1)^2X(t)=x]`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}(𝔼[X_i(t+1)X(t)=x])^2+{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{var}(X_i(t+1)X(t)=x).`$
Using Observations 4 and 5, this is at most $`n\overline{x}^2+_{i=1}^nd_i(x)^2+u_1(x)/n`$. But
$$d_i(x)=\frac{1}{n}\underset{j:|x_ix_j|1}{}(x_ix_j)=\frac{1}{n}(s_i(x)\mathrm{}_i(x)),$$
so the result follows. ∎
###### Lemma 7.
$`𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]n+2n^{1/2}\mathrm{\Phi }(x)^{1/2}`$.
###### Proof.
In the proof of Observation 6, we established that $`𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]_{i=1}^nd_i(x)^2+u_1(x)/n`$. Upper-bounding $`u_1(x)`$ and using $`d_i(x)1`$, we have
$$𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]n+\frac{1}{n}\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n}{}}|x_ix_j|,$$
and since $`|x_ix_j||x_i\overline{x}|+|x_j\overline{x}|`$, this is at most $`n+2_{i=1}^n|x_i\overline{x}|`$. By Cauchy-Schwarz, $`(_i|x_i\overline{x}|1)^2_i|x_i\overline{x}|^2_i1`$ so
$$𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]n+2(n\underset{i=1}{\overset{n}{}}|x_i\overline{x}|^2)^{1/2}.$$
###### Corollary 8.
$`𝔼[\mathrm{\Phi }(X(t+1))]n+2n^{1/2}(𝔼[\mathrm{\Phi }(X(t))])^{1/2}`$.
###### Proof.
Using Lemma 7, $`𝔼[\mathrm{\Phi }(X(t+1))]n+2n^{1/2}𝔼[f^{1/2}]`$ where $`f`$ denotes the random variable $`\mathrm{\Phi }(X(t))`$. By Jensen’s inequality $`𝔼[f^{1/2}](𝔼[f])^{1/2}`$ since the square-root function is concave, so we get $`𝔼[\mathrm{\Phi }(X(t+1))]n+2n^{1/2}(𝔼[f])^{1/2}`$. ∎
###### Lemma 9.
Either there is a $`t^{}<t`$ s.t. $`𝔼[\mathrm{\Phi }(X(t^{}))]18n`$ or $`𝔼[\mathrm{\Phi }(X(t))]9^{12^t}n^{12^t}\mathrm{\Phi }(X(0))^{2^t}`$.
###### Proof.
The proof is by induction on $`t`$. The base case is $`t=0`$. For the inductive step, note that $`12^t=_{k=1}^t2^k`$. Suppose that for all $`t^{}<t`$, $`𝔼[\mathrm{\Phi }(X(t^{}))]>18n`$ (otherwise we are finished). Then by Corollary 8,
$$𝔼[\mathrm{\Phi }(X(t))]=n+2n^{1/2}(𝔼[\mathrm{\Phi }(X(t1))])^{1/2}3n^{1/2}(𝔼[\mathrm{\Phi }(X(t1))])^{1/2}.$$
Applying the inductive hypothesis,
$$𝔼[\mathrm{\Phi }(X(t))]3n^{1/2}(3^{2(12^{(t1)})}n^{12^{(t1)}}\mathrm{\Phi }(X(0))^{2^{(t1)}})^{1/2}.$$
###### Corollary 10.
There is a $`\tau \mathrm{lg}\mathrm{lg}\mathrm{\Phi }(X(0))`$ such that $`𝔼[\mathrm{\Phi }(X(\tau ))]18n`$.
###### Proof.
Take $`t=\mathrm{lg}\mathrm{lg}\mathrm{\Phi }(X(0))`$. Either there is a $`\tau <t`$ with $`𝔼[\mathrm{\Phi }(X(\tau ))]18n`$ or, by the lemma,
$$𝔼[\mathrm{\Phi }(X(t))]9n\mathrm{\Phi }(X(0))^{2^t}18n.$$
###### Corollary 11.
There is a $`\tau \mathrm{lg}\mathrm{lg}\mathrm{\Phi }(X(0))`$ such that $`\mathrm{Pr}(\mathrm{\Phi }(X(\tau ))>720n)1/40`$.
###### Proof.
Consider the (non-negative) random variable $`Y=\mathrm{\Phi }(X(\tau ))`$ where $`\tau `$ is the quantity from Corollary 10. Markov’s inequality says that for any $`a>0`$, $`\mathrm{Pr}(Ya)𝔼[Y]/a`$. Now use Corollary 10 with $`a=720n`$. ∎
###### Corollary 12.
For all $`ϵ>0`$, provided that $`n<m^{1/3}`$, the expected time to reach $`ϵ`$-Nash equilibrium is $`O(\mathrm{log}\mathrm{log}m)`$.
###### Proof.
Since the bound is asymptotic as a function of $`m`$ for fixed $`ϵ`$, we can assume without loss of generality that $`m>(60/ϵ)^2`$ and that $`ϵm/(2n)`$ is an integer. We show that for any starting assignment $`X(0)`$, there exists $`\tau \mathrm{log}\mathrm{log}(m^2)`$ such that $`\mathrm{Pr}(X(\tau )\mathrm{is}ϵ\mathrm{Nash})`$ $`>\frac{39}{40}`$. This implies the statement of the result since the number of blocks of $`\tau `$ steps needed to reach an $`ϵ`$-Nash equilibrium is at most
$$1+\left(\frac{1}{40}\right)+\left(\frac{1}{40}\right)^2+\mathrm{}=\frac{40}{39}<2.$$
Suppose assignment $`x`$ is not $`ϵ`$-Nash. If $`X(t)=x`$ there exist resources $`i,j`$ with $`X_i(t)X_j(t)>ϵm/n`$. We use the following notation. Let $`\mathrm{\Delta }=ϵm/(2n)`$. Let $`\beta =X_i(t)X_j(t)2\mathrm{\Delta }`$. Note $`\beta >0`$. If $`X(t+1)`$ is obtained from $`X(t)`$ by transferring $`\mathrm{\Delta }`$ tasks from $`i`$ to $`j`$, then
$`\mathrm{\Phi }(X(t))\mathrm{\Phi }(X(t+1))`$ $`=X_i(t)^2+X_j(t)^2X_i(t+1)^2X_j(t+1)^2`$
$`=(2\mathrm{\Delta }+\beta +X_j(t))^2+X_j(t)^2(\mathrm{\Delta }+\beta +X_j(t))^2(\mathrm{\Delta }+X_j(t))^2`$
$`=2\mathrm{\Delta }(\mathrm{\Delta }+\beta +X_j(t))+\mathrm{\Delta }^2\left(2\mathrm{\Delta }X_j(t)+\mathrm{\Delta }^2\right)`$
$`=2\mathrm{\Delta }(\mathrm{\Delta }+\beta )\mathrm{\Delta }^2=(ϵm/2n)^2.`$
It follows that $`\mathrm{\Phi }(X(t))(ϵm/2n)^2`$.
From Corollary 11, $`\mathrm{Pr}(\mathrm{\Phi }(X(\tau ))<720n)>\frac{39}{40}`$, for $`\tau =\mathrm{log}\mathrm{log}(\mathrm{\Phi }(0))=O(\mathrm{log}\mathrm{log}m)`$.
An assignment $`X(\tau )`$ with $`\mathrm{\Phi }(X(\tau ))720n`$ must be $`ϵ`$-Nash if $`(ϵm/2n)^2>720n`$. Note that $`m>n^3`$ and $`m>(60/ϵ)^2`$. Hence, from $`ϵ^2(60/ϵ)^2n^3>4.720.n^3`$, we can deduce $`ϵ^2m^2>4.720.n^3`$, hence $`(ϵm/2n)^2>720n`$. ∎
Corollary 11 tells us that $`\mathrm{\Phi }(X(\tau ))`$ is likely to be $`O(n)`$. We want to show that $`\mathrm{\Phi }(X(t))`$ quickly gets even smaller (all the way to a Nash equilibrium) and to this end, we show that $`\mathrm{\Phi }(X(t))`$ is a super-martingale. By Observation 6, it suffices to show $`u(x)\mathrm{\Phi }(x)`$, and we proceed with this. In the following, we shall consider the cases $`|x_i\overline{x}|<2.5`$ for all $`i[n]`$ (Lemma 13) and $`i[n]:|x_i\overline{x}|2.5`$ (Lemma 14) separately.
###### Lemma 13.
Suppose that assignment $`x=(x_1,\mathrm{},x_n)`$ satisfies $`|x_i\overline{x}|<2.5`$ for all $`i[n]`$. Then $`u(x)\mathrm{\Phi }(x)`$.
###### Proof.
For all $`i[n]`$ and $`j[n]`$ we have $`|x_ix_j||x_i\overline{x}|+|x_j\overline{x}|<5`$. Let $`z=\mathrm{min}_ix_i`$ so every $`x_i\{z,\mathrm{},z+4\}`$. Let $`n_i=|\{jx_j=z+i\}|`$. Then
$$n^2\mathrm{\Phi }(x)=n^2\underset{i=1}{\overset{n}{}}x_i^2n\left(\underset{i=1}{\overset{n}{}}x_i\right)^2=n^2\left(\underset{j=0}{\overset{4}{}}n_j(z+j)^2\right)\left(\underset{j=0}{\overset{4}{}}n_j(z+j)\right)^2.$$
Also, $`n^2u(x)=nu_1(x)+u_2(x)`$, where
$$u_1(x)=n_0(2n_2+3n_3+4n_4)+n_1(2n_3+3n_4)+n_2(2n_0+2n_4)+n_3(3n_0+2n_1)+n_4(4n_0+3n_1+2n_2)$$
and
$$u_2(x)=n_0n_1^2+n_1(n_0n_2)^2+n_2(n_1n_3)^2+n_3(n_2n_4)^2+n_4n_3^2.$$
Plugging in these expressions and simplifying, we get
$`n^2\mathrm{\Phi }(x)n^2u(x)=`$ $`4n_0n_1n_2+3n_0^2n_3+4n_0n_1n_3+4n_0n_2n_3+4n_1n_2n_3+3n_0n_3^2+8n_0^2n_4+12n_0n_1n_4`$
$`+3n_1^2n_4+8n_0n_2n_4+4n_1n_2n_4+12n_0n_3n_4+4n_1n_3n_4+4n_2n_3n_4+8n_0n_4^2+3n_1n_4^2,`$
which is clearly non-negative since all coefficients are positive. ∎
###### Lemma 14.
Suppose that assignment $`x=(x_1,\mathrm{},x_n)`$ satisfies $`|x_n\overline{x}|2.5`$ and, for all $`i[n]`$, $`|x_i\overline{x}||x_n\overline{x}|`$. Let $`w=(w_1,\mathrm{},w_{n1})`$ be the assignment with $`w_i=x_i`$ for $`i[n1]`$. Then $`\mathrm{\Phi }(x)u(x)\mathrm{\Phi }(w)u(w)`$, that is, the lower bound on the potential drop for $`x`$ is at least as big as that for $`w`$.
###### Proof.
Let $`k=|x_n\overline{x}|`$. We will show
1. $`\mathrm{\Phi }(x)\mathrm{\Phi }(w)k^2`$, and
2. $`u(x)u(w)2k+1`$.
Then
$$\mathrm{\Phi }(x)u(x)(\mathrm{\Phi }(w)u(w))k^2(2k+1),$$
which is non-negative since $`k2.51+\sqrt{2}`$.
First, we prove (1). Let $`f(z)=_{i=1}^{n1}(x_iz)^2`$. Note that the derivative of $`f(z)`$ is
$$f^{}(z)=2(n1)z2\underset{i=1}{\overset{n1}{}}x_i=2(n1)z2(n1)\overline{w}.$$
Furthermore the second derivative is $`f^{\prime \prime }(z)=2(n1)0`$. Thus, $`f(z)`$ is minimized at $`z=\overline{w}`$. Now note that
$$\mathrm{\Phi }(x)\mathrm{\Phi }(w)=k^2+\underset{i=1}{\overset{n1}{}}(x_i\overline{x})^2\underset{i=1}{\overset{n1}{}}(x_i\overline{w})^2k^2.$$
Now we finish the proof by proving (2). Assume first that $`x_n=\overline{x}+k`$. Then
$$u_1(x)u_1(w)=2\underset{i[n]:|x_ix_n|>1}{}|x_ix_n|2\underset{i=1}{\overset{n}{}}|x_ix_n|=2\underset{i=1}{\overset{n}{}}(x_nx_i)=2nk.$$
Let $`z_j=|\{\mathrm{}x_{\mathrm{}}=j\}|`$. Clearly $`z_j=0`$ for $`j>x_n`$. Let $`\xi =x_n2k`$. For $`\mathrm{}[n]`$ we have $`x_{\mathrm{}}\overline{x}k=x_n2k`$ so $`z_j=0`$ for $`j<\xi `$. Now $`u_2(x)=_{j=\xi }^{x_n}z_j(z_{j1}z_{j+1})^2`$. The representation of $`w`$ in terms of $`z_j`$s is the same as the representation of $`x`$ except that $`z_{x_n}`$ is reduced by one. Therefore,
$`u_2(x)u_2(w)`$ $`=z_{x_n1}((z_{x_n2}z_{x_n})^2(z_{x_n2}z_{x_n}+1)^2)++(z_{x_n1}z_{x_n+1})^2`$
$`=z_{x_n1}(2z_{x_n2}+2z_{x_n}+z_{x_n1}1)z_{x_n1}(2z_{x_n}+z_{x_n1}).`$
But since $`z_{x_n}nz_{x_n1}`$, the upper bound on the right-hand side is at most
$$z_{x_n1}(2n2z_{x_n1}+z_{x_n1})=2z_{x_n1}(nz_{x_n1}/2),$$
which is at most $`n^2`$ since the right-hand side is maximized at $`z_{x_n1}=n`$. To finish the proof of (2), use the definition of $`u`$ to deduce that
$$u(x)u(w)\frac{u_1(x)u_1(w)}{n}+\frac{u_2(x)u_2(w)}{n^2}.$$
The proof of (2) when $`x_n=\overline{x}k`$ is similar. ∎
###### Corollary 15.
For any assignment $`x=(x_1,\mathrm{},x_n)`$, $`\mathrm{\Phi }(x)u(x)0`$.
###### Proof.
The proof is by induction on $`n`$. The base case, $`n=1`$, follows from Lemma 13. Suppose $`n>1`$. Neither $`\mathrm{\Phi }(x)`$ nor $`u(x)`$ depends upon the order of the components in $`x`$, so assume without loss of generality that $`|x_i\overline{x}||x_n\overline{x}|`$ for all $`i`$. If $`|x_n\overline{x}|<2.5`$ then apply Lemma 13. Otherwise, use Lemma 14 to find an assignment $`w=(w_1,\mathrm{},w_{n1})`$ such that $`\mathrm{\Phi }(x)u(x)\mathrm{\Phi }(w)u(w)`$. By the inductive hypothesis, $`\mathrm{\Phi }(w)u(w)0`$. ∎
Together, Observation 6 and Corollary 15 tell us that $`𝔼[\mathrm{\Phi }(X(t+1))X(t)=x]\mathrm{\Phi }(x)`$. The next lemma will be used to give a lower bound on the variance of the process. Let $`V=0.4n^2`$.
###### Lemma 16.
Suppose that $`X(t)=x`$ and that $`x`$ is not a Nash equilibrium. Then
$$\mathrm{Pr}(\mathrm{\Phi }(X(t+1))\mathrm{\Phi }(x)X(t)=x)V.$$
###### Proof.
Choose $`s`$ and $`\mathrm{}`$ such that for all $`i[n]`$, $`x_sx_ix_{\mathrm{}}`$. Since $`x`$ is not a Nash equilibrium, $`x_{\mathrm{}}>x_s+1`$. Assuming $`X(t)=x`$, consider the following experiment for choosing $`X(t+1)`$.
The intuition behind the experiment is as follows. We wish to show that the transition from $`X(t)`$ to $`X(t+1)`$ has some variance in the sense that $`\mathrm{\Phi }(X(t+1))`$ is sufficiently likely to differ from $`\mathrm{\Phi }(X(t))`$. To do this, we single out a “least loaded” resource $`s`$ and a “most loaded” resource $`\mathrm{}`$ as above. In the transition from $`X(t)`$ to $`X(t+1)`$ we make transitions from resources other than resource $`\mathrm{}`$ in the usual way. We pay special attention to transitions from resource $`\mathrm{}`$ (and particular attention to transitions from resource $`\mathrm{}`$ which could either go to resource $`s`$ or stay at resource $`\mathrm{}`$). It helps to be very precise about how the random decisions involving tasks that start at resource $`\mathrm{}`$ are made. In particular, for each task $`b`$ that starts at resource $`\mathrm{}`$, we first make a decision about whether $`b`$ would *accept* the transition from resource $`\mathrm{}`$ to resource $`s`$ *if $`b`$ happened to choose resource $`s`$*. Then we make the decision about which resource task $`b`$ should choose. Of course, we can’t cheat and we have to sample from the original required distribution. Here are the details.
Independently, for every $`i\mathrm{}`$, choose $`(Y_{i,1}(x),\mathrm{},Y_{i,n}(x))`$ from the multinomial distribution described in Section 2. (In the informal description above, this corresponds to making transitions from resources other than resource $`\mathrm{}`$ in the usual way.) Now, for every task $`bx_{\mathrm{}}`$, let $`z_b=1`$ with probability $`1x_s/x_{\mathrm{}}`$ and $`z_b=0`$ otherwise. (In the informal description above, this corresponds to deciding whether $`b`$ would *accept* the transition to $`s`$ if resource $`s`$ were (later) chosen.) Let $`x_{\mathrm{}}^+`$ be the number of tasks $`b`$ with $`z_b=1`$ and let $`x_{\mathrm{}}^{}`$ be the number of tasks $`b`$ with $`z_b=0`$. Choose $`(Y_{\mathrm{},1}^+(x),\mathrm{},Y_{\mathrm{},n}^+(x))`$ from a multinomial distribution with the constraint $`_{j=1}^nY_{\mathrm{},j}^+(x)=x_{\mathrm{}}^+`$ and probabilities given by
$$p_{\mathrm{},j}^+(x)=\{\begin{array}{cc}\frac{1}{n}\hfill & \text{if }j=s\text{,}\hfill \\ \frac{1}{n}\left(1\frac{x_j}{x_{\mathrm{}}}\right)\hfill & \text{if }js\text{ and }x_{\mathrm{}}>x_j+1\text{,}\hfill \\ 0\hfill & \text{if }\mathrm{}j\text{ but }x_{\mathrm{}}x_j+1\text{,}\hfill \\ 1_j\mathrm{}p_{\mathrm{},j}(x)\hfill & \text{if }\mathrm{}=j.\hfill \end{array}$$
Similarly, choose $`(Y_{\mathrm{},1}^{}(x),\mathrm{},Y_{\mathrm{},n}^{}(x))`$ from a multinomial distribution with the constraint $`_{j=1}^nY_{\mathrm{},j}^{}(x)=x_{\mathrm{}}^{}`$ and probabilities given by
$$p_{\mathrm{},j}^{}(x)=\{\begin{array}{cc}0\hfill & \text{if }j=s\text{,}\hfill \\ \frac{1}{n}\left(1\frac{x_j}{x_{\mathrm{}}}\right)\hfill & \text{if }js\text{ and }x_{\mathrm{}}>x_j+1\text{,}\hfill \\ 0\hfill & \text{if }\mathrm{}j\text{ but }x_{\mathrm{}}x_j+1\text{,}\hfill \\ 1_j\mathrm{}p_{\mathrm{},j}(x)\hfill & \text{if }\mathrm{}=j.\hfill \end{array}$$
For all $`j`$, let $`Y_{\mathrm{},j}(x)=Y_{\mathrm{},j}^+(x)+Y_{\mathrm{},j}^{}(x)`$. Informally, the $`p_{\mathrm{},j}^+`$ transition probabilities are set up so that packets which decided that they would accept a transition to $`s`$ behave appropriately and the $`p_{\mathrm{},j}^{}`$ transition probabilities are set up so that packets which decided that they would *not* accept a transition to $`s`$ behave appropriately. By combining the probabilities, we see that $`X(t+1)`$ is chosen from the correct distribution in this way.
Now, consider the transition from $`x`$ to $`X(t+1)`$. Condition on the choice for $`(Y_{i,1}(x),\mathrm{},Y_{i,n}(x))`$ for all $`i\mathrm{}`$. Suppose $`x_{\mathrm{}}^+>2`$. Condition on the choice for $`(Y_{\mathrm{},1}^{}(x),\mathrm{},Y_{\mathrm{},n}^{}(x))`$. Flip a coin for each of the first $`x_b^+2`$ tasks with $`z_b=1`$ to determine which of $`Y_{\mathrm{},1}^+(x),\mathrm{},Y_{\mathrm{},n}^+(x)`$ the task contributes to. Condition on these choices. Consider the following options:
1. Let $`x_1`$ be the resulting value of $`X(t+1)`$ when we add both of the last two tasks to $`Y_\mathrm{},\mathrm{}^+(x)`$.
2. Let $`x_2`$ be the resulting value of $`X(t+1)`$ when we add one of the last two tasks to $`Y_\mathrm{},\mathrm{}^+(x)`$ and the other to $`Y_{\mathrm{},s}^+(x)`$.
3. Let $`x_3`$ be the resulting value of $`X(t+1)`$ when we add both of the last two tasks to $`Y_{s,s}^+(x)`$.
Note that, given the conditioning, each of these choices occurs with probability at least $`n^2`$. Also, $`\mathrm{\Phi }(x_1)`$, $`\mathrm{\Phi }(x_2)`$ and $`\mathrm{\Phi }(x_3)`$ are not all the same. Thus, $`\mathrm{Pr}(\mathrm{\Phi }(X(t+1)\mathrm{\Phi }(x)X(t)=x,x_{\mathrm{}}^+>2)n^2`$. Also,
$$\mathrm{Pr}(x_{\mathrm{}}^+>2)=1\left(\frac{x_s}{x_{\mathrm{}}}\right)^x_{\mathrm{}}x_{\mathrm{}}\left(1\frac{x_s}{x_{\mathrm{}}}\right)\left(\frac{x_s}{x_{\mathrm{}}}\right)^{x_{\mathrm{}}1}.$$
Since the derivative with respect to $`x_s`$ is negative, this is minimized by taking $`x_s`$ as large as possible, namely $`x_{\mathrm{}}2`$, so $`\mathrm{Pr}(x_{\mathrm{}}^+>2)17e^20.4`$, and the result follows. ∎
In order to finish our proof of convergence, we need the following observation about $`\mathrm{\Phi }(x)`$.
###### Observation 17.
For any assignment $`x`$, $`\mathrm{\Phi }(x)m^2`$. Let $`r=mmodn`$. Then $`\mathrm{\Phi }(x)r(1r/n)`$, with equality if and only if $`x`$ is a Nash equilibrium.
###### Proof.
Suppose that in assignment $`x`$ there are resources $`i`$ and $`j`$ such that $`x_ix_j2`$. Let $`x^{}`$ be the assignment constructed from $`x`$ by transferring a task from resource $`i`$ to resource $`j`$. Then
$`\mathrm{\Phi }(x)\mathrm{\Phi }(x^{})`$ $`=x_i^2x_{i}^{}{}_{}{}^{2}+x_j^2x_{j}^{}{}_{}{}^{2}=x_i^2(x_i^22x_i+1)+x_j^2(x_j^2+2x_j+1)`$
$`=2x_i2x_j2=2(x_ix_j)2>0.`$
Now suppose that, in some assignment $`x^{}`$, resources $`i`$ and $`j`$ satisfy $`x_i^{}x_j^{}>0`$. Let $`x`$ be the assignment constructed from $`x^{}`$ by transferring a task from resource $`j`$ to resource $`i`$. Since $`(x_i^{}+1)(x_j^{}1)2`$, the above argument gives $`\mathrm{\Phi }(x)>\mathrm{\Phi }(x^{})`$. We conclude that an assignment $`x`$ with maximum $`\mathrm{\Phi }(x)`$ must have all of the tasks in the same resource, with $`\mathrm{\Phi }(x)=m^2`$.
Furthermore, an assignment $`x`$ with minimum $`\mathrm{\Phi }(x)`$ must have $`|x_ix_j|1`$ for all $`i,j`$. In this case there must be $`r`$ resources with loads of $`q+1`$ and $`nr`$ resources with loads of $`q`$, where $`m=qn+r`$. So
$$\mathrm{\Phi }(x)=r(q+1\overline{x})^2+(nr)(q\overline{x})^2=r\left(1\frac{r}{n}\right)^2+(nr)\left(\frac{r}{n}\right)^2=r\left(1\frac{r}{n}\right).$$
Note that $`x`$ is a Nash assignment if and only if $`|x_ix_j|1`$ for all $`i`$ and $`j`$. ∎
Combining Observation 17 and Corollary 11 we find that there is a $`\tau \mathrm{lg}\mathrm{lg}m^2`$ such that $`\mathrm{Pr}(\mathrm{\Phi }(X(\tau ))>720n)1/40`$. Let $`B=7200n+\frac{m^2}{n}\frac{m^2}{n}`$. Let $`t^{}=\tau +10B^2/V`$.
###### Lemma 18.
Given any starting state $`X(0)=x`$, the probability that $`X(t^{})`$ is a Nash equilibrium is at least $`3/4`$.
###### Proof.
The proof is based on a standard martingale argument, see . Suppose that $`\mathrm{\Phi }(X(\tau ))720n`$. Let $`W_t=\mathrm{\Phi }(X(t+\tau ))r(1r/n)`$ and let $`D_t=\mathrm{min}(W_t,B)`$. Note that $`D_0720n`$. Together, Observation 6 and Corollary 15 tell us that $`W_t`$ is a supermartingale. This implies that $`D_t`$ is also a supermartingale since
$$𝔼[D_{t+1}D_t=x<B]𝔼[W_{t+1}W_t=x<B]W_t=D_t,$$
and
$$𝔼[D_{t+1}D_t=B]B=D_t.$$
Together, Lemma 16 and Observation 17 tell us that if $`x>0`$, $`\mathrm{Pr}(W_{t+1}W_tW_t=x)V`$. Thus, if $`0<x<B`$,
$`\mathrm{Pr}(D_{t+1}D_tD_t=x)`$ $`=\mathrm{Pr}(\mathrm{min}(W_{t+1},B)W_tW_t=x)`$
$`\mathrm{Pr}(W_{t+1}W_tBW_tW_t=x)=\mathrm{Pr}(W_{t+1}W_tW_t=x)V.`$
Since $`D_{t+1}D_t`$ is an integer, $`𝔼[(D_{t+1}D_t)^20<D_t<B]V`$. Let $`T`$ be the first time at which either (a) $`D_t=0`$ (i.e., $`X(t+\tau )`$ is a Nash equilibrium), or (b) $`D_t=B`$. Note that $`T`$ is a stopping time. Define $`Z_t=(BD_t)^2Vt`$, and observe that $`Z_{tT}`$ is a sub-martingale, where $`tT`$ denotes the minimum of $`t`$ and $`T`$. Let $`p`$ be the probability that (a) occurs. By the optional stopping theorem $`𝔼[D_T]D_0`$, so $`(1p)B=𝔼[D_T]D_0`$ and $`p1D_0/B\frac{9}{10}`$. Also, by the optional stopping theorem
$$pB^2V𝔼[T]=𝔼[(BD_T)^2]V𝔼[T]=𝔼[Z_T]Z_0=(BD_0)^2>0,$$
so $`𝔼[T]pB^2/V`$. Conditioning on (a) occurring, it follows that $`𝔼[TD_T=0]B^2/V`$. Hence $`\mathrm{Pr}(T>10B^2/VD_T=0)\frac{1}{10}`$. So, if we now run for $`10B^2/V`$ steps, then the probability that we do not reach a Nash equilibrium is at most $`\frac{1}{40}+2\frac{1}{10}<1/4`$. ∎
Now we can give the proof of Theorem 2.
###### Proof.
Subdivide time into intervals of $`t^{}`$ steps. The probability that the process has not reached a Nash equilibrium before the $`(j+1)`$st interval is at most $`(1/4)^j`$. ∎
## 4 Lower Bounds
In this section we prove the lower-bound results stated in the introduction. We will use the following Chernoff bound which can be found, for example, in . Let $`N1`$ and let $`p_i[0,1]`$ for $`i=1,\mathrm{},N`$. Let $`X_1,X_2,\mathrm{},X_N`$ be independent Bernoulli random variables with $`\mathrm{Pr}(X_i=1)=p_i`$ for $`i=1,\mathrm{},N`$ and let $`X=X_1+\mathrm{}+X_N`$. Then we have $`𝔼[X]=_{i=1}^Np_i`$ and for $`0ϵ1`$,
$$\mathrm{Pr}(X(1ϵ)𝔼[X])\mathrm{exp}\left(\frac{ϵ^2𝔼[X]}{3}\right).$$
(1)
The following theorem gives an exponential lower bound for the expected convergence time of the process in Figure 1.
###### Theorem 19.
Let $`X(t)`$ be the process in Figure 1 with $`m=n`$. Let $`X(0)`$ be the assignment given by $`X(0)=(n,0,\mathrm{},0)`$. Let $`T`$ be the first time at which $`X(t)`$ is a Nash equilibrium. Then $`𝔼[T]=\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$.
###### Proof.
For an assignment $`x`$, let $`n_0(x)`$ denote the number of resources $`i`$ with $`x_i=0`$. Thus, $`n_0(X(0))=n1`$. The (unique) Nash equilibrium $`x`$ assigns one task to each resource, so $`n_0(x)=0`$. Let $`k=\sqrt{n}`$. We will show that for any assignment $`x`$ with $`n_0(x)k`$,
$$\mathrm{Pr}(n_0(X(t))<kX(t1)=x)\mathrm{exp}(\mathrm{\Theta }(\sqrt{n})).$$
This implies the result.
Suppose $`X(t1)=x`$ with $`n_0(x)k`$. For convenience, let $`n_0`$ denote $`n_0(x)`$. Let $`x^{}`$ denote $`X(t)`$, and let $`n_0^{}`$ denote $`n_0(x^{})`$. We will show that, with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$, $`n_0^{}k`$. During the course of the proof, we will assume, where necessary, that $`n`$ is sufficiently large. This is without loss of generality given the $`\mathrm{\Theta }`$ notation in the statement of the result.
#### Case 1
$`n_0>8k`$.
Consider the protocol in Figure 1. Let $`U=\{bx_{j_b}=0\}`$. $`𝔼[|U|]=n_0`$, so by the Chernoff bound (Equation (1)), $`\mathrm{Pr}(|U|\frac{n_0}{2}+\frac{3n_0}{8})\mathrm{Pr}(|U|\frac{8}{9}n_0)=\mathrm{exp}\left(\mathrm{\Theta }(\sqrt{n})\right).`$ Thus, $`|U|n_0/2+3n_0/8`$ with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$. Suppose this is the case. Partition $`U`$ into $`U_1`$ and $`U_2`$ with $`|U_1|=n_0/2`$. Let $`W=_{bU_1}\{j_b\}`$. First, suppose $`|W|\frac{3}{8}n_0`$. In that case
$$|\{jx_j^{}>0\}|n|U_1|+\frac{3}{8}n_0=nn_0/2+\frac{3}{8}n_0nk,$$
so $`n_0^{}k`$. Otherwise, let $`U^{}=\{bU_2j_bW\}`$.
$$𝔼[|U^{}|]=|U_2|\frac{|W|}{n_0}\frac{9}{64}n_0>\frac{9}{8}k,$$
so by the Chernoff bound (1), $`\mathrm{Pr}(|U^{}|k)=\mathrm{Pr}(|U^{}|(1\frac{1}{9})𝔼[|U^{}|])=\mathrm{exp}\left(\mathrm{\Theta }(\sqrt{n})\right)`$, recalling that $`k=\sqrt{n}`$. Thus $`|U^{}|k`$ with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$, which implies $`n_0^{}k`$.
#### Case 2
$`kn_08k`$.
Consider the protocol in Figure 1. Let $`L`$ be the set of “loners” defined by $`L=\{ix_i=1\}`$ and let $`\mathrm{}=|L|`$. The number of resources $`i`$ with $`x_i>1`$ is $`nn_0\mathrm{}`$ and this is at most half as many as the number of tasks assigned to such resources (which is $`n\mathrm{}`$), so $`\mathrm{}n2n_0`$. Let $`U=\{bi_bL\text{ and }x_{j_b}=0\}`$. $`𝔼[|U|]=\mathrm{}\frac{n_0}{n}\frac{(n2n_0)n_0}{n}=\mathrm{\Theta }(\sqrt{n})`$, so by the Chernoff bound (1), $`\mathrm{Pr}(|U|2\frac{1}{4}\mathrm{}\frac{n_0}{n})\mathrm{Pr}(|U|\frac{2}{3}𝔼[|U|])\mathrm{exp}\left(\mathrm{\Theta }(\sqrt{n})\right)`$. Thus, $`|U|2\frac{1}{4}\mathrm{}\frac{n_0}{n}`$ with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$. Suppose this is the case. Let $`U_1`$ and $`U_2`$ be disjoint subsets of $`U`$ of size $`\frac{1}{4}\mathrm{}\frac{n_0}{n}`$. Order tasks in $`U`$ arbitrarily and let $`S=\{bU\text{for some }b^{}U\text{ with }b^{}<b\text{}j_b^{}=j_b\text{.}\}`$. (Note that $`|S|`$ does not depend on the ordering.) Let $`W=_{bU_1}\{j_b\}`$.
Note that if $`|W|\frac{1}{5}\mathrm{}\frac{n_0}{n}`$ then $`|S|\frac{1}{20}\mathrm{}\frac{n_0}{n}>\frac{n_0}{40}\left(\frac{\mathrm{}}{n}\right)^2`$. Otherwise, let $`U^{}=\{bU_2j_bW\}`$.
$$𝔼[|U^{}|]=|U_2|\frac{|W|}{n_0}\frac{n_0}{20}\left(\frac{\mathrm{}}{n}\right)^2,$$
so, by the Chernoff bound (1), $`\mathrm{Pr}(|U^{}|\frac{1}{2}\frac{n_0}{20}\left(\frac{\mathrm{}}{n}\right)^2)\mathrm{exp}\left(\mathrm{\Theta }(\sqrt{n})\right)`$ (recall that $`n_0\left(\frac{\mathrm{}}{n}\right)^2n_0\left(\frac{n2n_0}{n}\right)^2k\left(\frac{n16k}{n}\right)^2=\mathrm{\Theta }(\sqrt{n})`$), and thus $`|U^{}|\frac{n_0}{40}\left(\frac{\mathrm{}}{n}\right)^2`$ with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$, so $`|S|\frac{n_0}{40}\left(\frac{\mathrm{}}{n}\right)^2.`$
Suppose then that $`|S|\frac{n_0}{40}\left(\frac{\mathrm{}}{n}\right)^2.`$ Assuming that $`n`$ is sufficiently large, $`|S|k/41`$. Let $`B_0=_{bU}\{j_b\}`$ and $`B_1=_{bLU}\{i_b\}`$. Note that every resource in $`B_0B_1`$ is used in $`x^{}`$ for some task $`bL`$. Thus, $`|B_0B_1|\mathrm{}|S|`$. Let $`R=\{ix_i=0\}LB_0B_1`$. Then $`|R|n_0+\mathrm{}(\mathrm{}|S|)n_0+|S|(1+\frac{1}{41})k`$.
Let $`T=\{bi_bL,j_bR\}`$. $`𝔼[T]=(n\mathrm{})\frac{|R|}{n}`$ and
$$\mathrm{Pr}\left(T\frac{|R|}{100}\right)\left(\genfrac{}{}{0pt}{}{n\mathrm{}}{\frac{|R|}{100}}\right)\left(\frac{|R|}{n}\right)^{|R|/100}\left(\frac{2n_0e100}{n}\right)^{|R|/100},$$
so with probability at least $`1\mathrm{exp}(\mathrm{\Theta }(\sqrt{n}))`$, $`T<|R|/100`$. In that case, $`n_0^{}|R|(1\frac{1}{100})k`$. ∎
The following theorem provides a lower bound on the expected convergence time regardless of which of the two protocols is being used.
###### Theorem 20.
Suppose that $`m`$ is even. Let $`X(t)`$ be the process in Figure 2 with $`n=2`$. Let $`X(0)`$ be the assignment given by $`X(0)=(m,0)`$. Let $`T`$ be the first time at which $`X(t)`$ is a Nash equilibrium. Then $`𝔼[T]=\mathrm{\Omega }(\mathrm{log}\mathrm{log}m)`$. The same result holds for the process in Figure 1.
###### Proof.
Note that both protocols have the same behaviour since $`m`$ is even and, therefore, the situation $`x_1=x_2+1`$ cannot arise. For concreteness, focus on the protocol in Figure 2.
Let $`y(x)=\mathrm{max}_ix_im/2`$ and let $`y_t=y(X(t))`$ so $`y_0=m/2`$ and, for a Nash equilibrium $`x`$, $`y(x)=0`$. We will show that for any assignment $`x`$, $`\mathrm{Pr}(y_{t+1}>y(x)^{1/10}X(t)=x)1y_t^{1/4}`$. (There is nothing very special about the exact value “$`1/10`$” – this value is being used as part of an explicit “lack of concentration” inequality in the proof, noting that for a lower bound we essentially want to lower-bound the variances of the load distributions. This seems to require a somewhat ad-hoc approach, in contrast with the usage of concentration inequalities.)
Suppose $`X(t)=x`$ is an assignment with $`x_1x_2`$. As we have seen in Section 2, $`Y_{1,2}(x)`$ (the number of migrations from resource $`1`$ to resource $`2`$ in the round) is a binomial random variable
$$B(x_1,\frac{1}{2}\left(1\frac{x_2}{x_1}\right))=B(\frac{m}{2}+y_t,\frac{2y_t}{m+2y_t}).$$
In general, let $`T_t`$ be the number of migrations from the most-loaded resource in $`X(t)`$ to the least-loaded resource and note that the distribution of $`T_t`$ is $`B(\frac{m}{2}+y_t,\frac{2y_t}{m+2y_t})`$ with mean $`y_t`$. If $`T_t=y_t+\mathrm{}`$ or $`T_t=y_t\mathrm{}`$ then $`y_{t+1}=\mathrm{}`$. Thus $`\mathrm{Pr}(y_{t+1}>y_t^{1/10})=\mathrm{Pr}(|T_t𝔼[T_t]|>y_t^{1/10})`$. We continue by showing that this binomial distribution is sufficiently “spread out” in the region of its mode, that we can find an upper bound on $`\mathrm{Pr}(y_{t+1}y_t^{1/10})`$. This will lead to our lower bound on the expected time for $`(y_t)_t`$ to decrease below some constant (we use the constant 16).
$`\mathrm{Pr}(T_t=y_t)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\frac{1}{2}m+y_t}{y_t}}\right)\left({\displaystyle \frac{2y_t}{m+2y_t}}\right)^{y_t}\left({\displaystyle \frac{m}{m+2y_t}}\right)^{\frac{1}{2}m}`$
$`\mathrm{Pr}(T_t=y_t+j)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\frac{1}{2}m+y_t}{y_t+j}}\right)\left({\displaystyle \frac{2y_t}{m+2y_t}}\right)^{y_t+j}\left({\displaystyle \frac{m}{m+2y_t}}\right)^{\frac{1}{2}mj}`$
Suppose $`j>0`$.
$`{\displaystyle \frac{\mathrm{Pr}(T_t=y_t+j)}{\mathrm{Pr}(T_t=y_t)}}`$ $`=`$ $`\left({\displaystyle \frac{2y_t}{m+2y_t}}\right)^j\left({\displaystyle \frac{m}{m+2y_t}}\right)^j\left({\displaystyle \frac{y_t!(\frac{1}{2}m)!}{(y_t+j)!(\frac{1}{2}m+y_t(y_t+j))!}}\right)`$
$`=`$ $`\left({\displaystyle \frac{2y_t}{m}}\right)^j\left({\displaystyle \underset{\mathrm{}=1}{\overset{j}{}}}{\displaystyle \frac{\frac{1}{2}m+1\mathrm{}}{y_t+\mathrm{}}}\right)=\left({\displaystyle \frac{2y_t}{m}}\right)^j\left({\displaystyle \underset{\mathrm{}=1}{\overset{j}{}}}{\displaystyle \frac{m+22\mathrm{}}{2y_t+2\mathrm{}}}\right)`$
$`>`$ $`\left({\displaystyle \frac{2y_t}{m}}\right)^j\left({\displaystyle \underset{\mathrm{}=1}{\overset{j}{}}}{\displaystyle \frac{m2j}{2y_t+2j}}\right)=\left[\left({\displaystyle \frac{2y_t}{m}}\right)\left({\displaystyle \frac{m2j}{2y_t+2j}}\right)\right]^j.`$
Similarly, for $`j<0`$,
$`{\displaystyle \frac{\mathrm{Pr}(T_t=y_t+j)}{\mathrm{Pr}(T_t=y_t)}}`$ $`=`$ $`\left({\displaystyle \frac{2y_t}{m}}\right)^j\left({\displaystyle \underset{\mathrm{}=1}{\overset{|j|}{}}}{\displaystyle \frac{y_t+1\mathrm{}}{\frac{1}{2}m+\mathrm{}}}\right)=\left({\displaystyle \frac{m}{2y_t}}\right)^{|j|}\left({\displaystyle \underset{\mathrm{}=1}{\overset{|j|}{}}}{\displaystyle \frac{2y_t+22\mathrm{}}{m+2\mathrm{}}}\right)`$
$`>`$ $`\left({\displaystyle \frac{m}{2y_t}}\right)^{|j|}\left({\displaystyle \frac{2y_t2|j|}{m+2|j|}}\right)^{|j|}=\left[\left({\displaystyle \frac{m}{2y_t}}\right)\left({\displaystyle \frac{2y_t2|j|}{m+2|j|}}\right)\right]^{|j|}`$
$`=`$ $`\left[\left({\displaystyle \frac{2y_t}{m}}\right)\left({\displaystyle \frac{m2j}{2y_t+2j}}\right)\right]^j.`$
So for all $`j`$,
$$\frac{\mathrm{Pr}(T_t=y_t+j)}{\mathrm{Pr}(T_t=y_t)}>\left[\left(\frac{2y_t}{m}\right)\left(\frac{m2j}{2y_t+2j}\right)\right]^j=\left[\left(\frac{y_t}{y_t+j}\right)\left(\frac{m2j}{m}\right)\right]^j.$$
So, for all $`j`$ with $`|j|y_t^{1/4}`$, where $`y_t^{1/4}`$ is the positive fourth root of $`y_t`$, this is at least
$`\left({\displaystyle \frac{y_t}{y_t+y_t^{1/4}}}\right)^{y_t^{1/4}}\left({\displaystyle \frac{m2y_t^{1/4}}{m}}\right)^{y_t^{1/4}}`$
$``$ $`\left({\displaystyle \frac{y_t}{y_t+y_t^{1/4}}}\right)^{y_t^{1/4}}\left({\displaystyle \frac{2y_t2y_t^{1/4}}{2y_t}}\right)^{y_t^{1/4}}=\left({\displaystyle \frac{y_ty_t^{1/4}}{y_t+y_t^{1/4}}}\right)^{y_t^{1/4}}=\left({\displaystyle \frac{y_t+y_t^{1/4}2y_t^{1/4}}{y_t+y_t^{1/4}}}\right)^{y_t^{1/4}}`$
$`=`$ $`\left(1{\displaystyle \frac{2y_t^{1/4}}{y_t+y_t^{1/4}}}\right)^{y_t^{1/4}}\left(1{\displaystyle \frac{2y_t^{1/4}}{y_t}}\right)^{y_t^{1/4}}=\left(12y_t^{3/4}\right)^{y_t^{1/4}}`$
$``$ $`12y_t^{3/4}y_t^{1/4}=12y_t^{1/2}{\displaystyle \frac{1}{2}}`$
where the last inequality just requires $`y_t16`$.
Note that the mode of a binomial distribution is one or both of the integers closest to the expectation, and the distribution is monotonically decreasing as you move away from the mode. But, for $`|j|y_t^{1/4}`$, $`\mathrm{Pr}(T_t=y_t+j)\frac{1}{2}\mathrm{Pr}(T_t=y_t)`$, hence $`\mathrm{Pr}(T_t=y_t)2/(1+2y_t^{1/4})`$. Since $`\mathrm{Pr}(T_t=y_t+j)\mathrm{Pr}(T_t=y_t)`$, it follows that
$$\mathrm{Pr}(T_t[y_ty_t^{1/10},y_t+y_t^{1/10}])(2y_t^{1/10}+1)\mathrm{Pr}(T_t=y_t)<3y_t^{3/20}.$$
We say that the transition from $`y_t`$ to $`y_{t+1}`$ is a “fast round” if $`y_{t+1}y_t^{1/10}`$ (equivalently, it is a fast round if $`T_t[y_ty_t^{1/10},y_t+y_t^{1/10}]`$). Otherwise it is a slow round. Recall that $`y_0=m/2`$. Let
$$r=\mathrm{log}_{10}\left(\frac{\mathrm{log}(y_0)}{\mathrm{log}(12^{20/3})}\right).$$
If the first $`j`$ rounds are slow then $`y_jy_0^{10^j}`$. If $`jr`$ then $`y_0^{10^j}12^{20/3}`$ so the probability that the transition from $`y_j`$ to $`y_{j+1}`$ is the first fast round is at most $`3\left(y_0^{10^j}\right)^{3/20}1/4`$.
Also, if $`j<r`$ then these probabilities increase geometrically so that the ratio of the probability that the transition to $`y_{j+1}`$ is the first fast round and the probability that the transition to $`y_j`$ is the first fast round is
$$\frac{3\left(y_0^{10^{(j+1)}}\right)^{3/20}}{3\left(y_0^{10^j}\right)^{3/20}}=\left(y_0^{10^j10^{(j+1)}}\right)^{3/20}\left(y_0^{10^{(j+1)}}\right)^{3/20}122,$$
so $`_{j=0}^{r1}\mathrm{Pr}(\text{transition from }y_j\text{ to }y_{j+1}\text{ is the first fast round)}21/4=\frac{1}{2}`$. Therefore, with probability at least $`1/2`$, all of the first $`r`$ rounds are slow. In this case, $`\mathrm{arg}\mathrm{min}_t(y_t16)=\mathrm{\Omega }(\mathrm{log}\mathrm{log}(m))`$, which proves the theorem. ∎
We also have the following observation.
###### Observation 21.
Let $`X(t)`$ be the process in Figure 2 with $`m=n`$. Let $`X(0)`$ be the assignment given by $`X(0)=(2,0,1,\mathrm{},1)`$. Let $`T`$ be the first time at which $`X(t)`$ is a Nash equilibrium. Then $`𝔼[T]=\mathrm{\Omega }(n)`$.
The observation follows from the fact that the state does not change until one of the two tasks assigned to the first resource chooses the second resource.
## 5 Summary
We have analyzed a very simple, strongly distributed rerouting protocol for $`m`$ tasks on $`n`$ resources. We have proved an upper bound of $`(\mathrm{log}\mathrm{log}m+n^4)`$ on the expected convergence time (convergence to a Nash equilibrium), and for $`m>n^3`$ an upper bound of $`O(\mathrm{log}\mathrm{log}m)`$ on the time to reach an approximate Nash equilibrium. Our lower bound of $`\mathrm{\Omega }(\mathrm{log}\mathrm{log}m+n)`$ matches the upper bound as function of $`m`$. We have also shown an exponential lower bound on the convergence time for a related protocol that allows “neutral moves”.
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# Quaternions in molecular modeling
## I Introduction
Quaternions were introduced in the mid-nineteenth century by Hamilton \[1; 2\] as an extension of complex numbers and as a tool for manipulating 3-dimensional vectors. Indeed Maxwell used them to introduce vectors in his exposition of electromagnetic theory \[3, §§10–11\]. However, unlike complex numbers which occupy a central role in the development of algebra, quaternions found no similar place in mathematics and, with the introduction of modern vector notation by Gibbs , quaternions fell out of favor by the end of the nineteenth century. Nevertheless, quaternions excel as a way of representing rotations of objects in 3-dimensional space. They are economical to work with (both in terms of storage and computation); but more importantly they offer a clean conceptual framework which allow several problems involving rotations to be easily solved.
Basic quaternion algebra is well covered in Hamilton’s papers \[1; 2\], which are both accessible and readable. These papers may be supplemented with a wealth of on-line resources \[5; 6\]. Many authors over the past 20 years have “rediscovered” the application of quaternions to rotations and it is with some trepidation that this author inflicts another paper on the subject on the scientific community. However, within the molecular modeling community, quaternions are quite narrowly applied. This paper therefore briefly reviews quaternion algebra and then describes their applications to a broad range of rotational problems in molecular modeling. Much of this material has appeared before—but often scattered about in journals for fields unrelated to molecular modeling. I have, therefore, endeavored to organize the material, to generalize it, and to present it with a consistent notation, with the hope this affords a deeper appreciation of the power of quaternions in describing rotations and encourages their wider adoption in molecular modeling.
The outline of this paper is as follows. After introducing quaternions and their use in describing rotations, we tackle various applications. First we review the quaternion method for computing the least-squares fit of two conformations of the same molecule. We also see how to include molecular inversions and discuss why the least-squares fit is a poor choice to describe the orientation of a flexible molecule. We next show how to interpolate smoothly between two orientations and that this corresponds to rotating the molecule at constant angular velocity. In order to carry out statistics on orientational data, we give a robust definition of the mean orientation showing how to transform the deviations from the mean to 3-dimensional space so that familiar statistical tools may be employed. In Monte Carlo applications, we need to be able to select a random orientation uniformly; we show that this is trivially accomplished in quaternion space and we also consider the problem of making random incremental rotations. Finally, it is frequently useful to impose a grid on orientation space and we illustrate how this may be done with applications to quadrature and searching.
## II Quaternions
The original notation for quaternions paralleled the convention for complex numbers
$$𝗾=𝗾_\mathrm{𝟬}𝘂+𝗾_\mathrm{𝟭}𝗶+𝗾_\mathrm{𝟮}𝗷+𝗾_\mathrm{𝟯}𝗸,$$
which obey the conventional algebraic rule for addition and multiplication by scalars (real numbers) and which obey an associative non-commutative rule for multiplication where $`𝘂`$ is the identity element and
$`𝗶^\mathrm{𝟮}=𝗷^\mathrm{𝟮}=𝗸^\mathrm{𝟮}=𝘂,`$
$`\mathrm{𝗶𝗷}=\mathrm{𝗷𝗶}=𝗸,\mathrm{𝗷𝗸}=\mathrm{𝗸𝗷}=𝗶,\mathrm{𝗸𝗶}=\mathrm{𝗶𝗸}=𝗷.`$
It is frequently useful to regard quaternions as an ordered set of 4 real quantities which we write as
$$𝗾=[𝗾_\mathrm{𝟬},𝗾_\mathrm{𝟭},𝗾_\mathrm{𝟮},𝗾_\mathrm{𝟯}],$$
(1)
or as a combination of a scalar and a vector
$$𝗾=[𝗾_\mathrm{𝟬},𝐪],$$
(2)
where $`𝐪=[q_1,q_2,q_3]`$. A “scalar” quaternion has zero vector part and we shall write this as $`[q_0,\mathrm{𝟎}]=q_0𝘂=𝗾_\mathrm{𝟬}`$. A “pure” quaternion has zero scalar part $`[0,𝐪]`$.
In the scalar-vector representation, multiplication becomes
$$\mathrm{𝗽𝗾}=[𝗽_\mathrm{𝟬}𝗾_\mathrm{𝟬}𝐩𝐪,𝗽_\mathrm{𝟬}𝐪+𝗾_\mathrm{𝟬}𝐩+𝐩\times 𝐪],$$
where “$``$” and “$`\times `$” are the vector dot and cross products. The conjugate of a quaternion is given by
$$\overline{𝗾}=[q_0,𝐪];$$
the squared norm of a quaternion is
$$\left|𝗾\right|^2=𝗾\overline{𝗾}=𝗾_\mathrm{𝟬}^\mathrm{𝟮}+𝗾_\mathrm{𝟭}^\mathrm{𝟮}+𝗾_\mathrm{𝟮}^\mathrm{𝟮}+𝗾_\mathrm{𝟯}^\mathrm{𝟮},$$
and its inverse is
$$𝗾^\mathrm{𝟭}=\overline{𝗾}/\left|𝗾\right|^\mathrm{𝟮}.$$
Quaternions with $`\left|𝗾\right|=1`$ are called unit quaternions, for which we have $`𝗾^\mathrm{𝟭}=\overline{𝗾}`$.
The quaternion $`𝗾`$ can also be represented as a $`2\times 2`$ complex matrix,
$$\left(\begin{array}{cc}q_0+iq_1& q_2+iq_3\\ q_2+iq_3& q_0iq_1\end{array}\right),$$
or as a $`4\times 4`$ real matrix,
$$𝗤(𝗾)=\left(\begin{array}{cccc}𝗾_\mathrm{𝟬}& 𝗾_\mathrm{𝟭}& 𝗾_\mathrm{𝟮}& 𝗾_\mathrm{𝟯}\\ 𝗾_\mathrm{𝟭}& 𝗾_\mathrm{𝟬}& 𝗾_\mathrm{𝟯}& 𝗾_\mathrm{𝟮}\\ 𝗾_\mathrm{𝟮}& 𝗾_\mathrm{𝟯}& 𝗾_\mathrm{𝟬}& 𝗾_\mathrm{𝟭}\\ 𝗾_\mathrm{𝟯}& 𝗾_\mathrm{𝟮}& 𝗾_\mathrm{𝟭}& 𝗾_\mathrm{𝟬}\end{array}\right);$$
(3)
in these forms, quaternion multiplication becomes matrix multiplication.
The notation we adopt here is to use light-face italics for scalar quantities, bold roman for 3-dimensional vectors and $`3\times 3`$ matrices, bold sans serif for quaternions and $`4\times 4`$ matrices. Quaternion multiplication is indicated by $`\mathrm{𝗽𝗾}`$, while “$``$” is used to indicate matrix-vector and vector-vector (including quaternion-quaternion) contractions and in this context $`𝗾`$ and $`𝐯`$ are treated as column vectors. Thus, we may write $`\left|𝗾\right|^2=𝗾^𝖳𝗾`$. We also find that $`\mathrm{𝗽𝗾}=𝗽^𝖳𝗤(𝗾)`$, with $`𝗤`$ given by eq. (3). Consistent with eqs. (1) and (2), we shall number quaternion indices starting at 0 and vector indices from 1.
## III Rotations
The chief application of quaternions to molecular modeling lies in their use to represent rotations. Consider a unit quaternion
$$𝗾=[\mathrm{cos}(\theta /\mathrm{𝟮}),𝐯\mathrm{sin}(\theta /\mathrm{𝟮})],$$
(4)
where $`\left|𝐯\right|=1`$, and define an operator $`R_𝗾`$ on 3-dimensional vectors by
$$[0,R_𝗾(𝐱)]=𝗾[\mathrm{𝟬},𝐱]\overline{𝗾}.$$
(5)
Multiplying out the quaternion product, we find
$$R_𝗾(𝐱)=𝐑(𝗾)𝐱,$$
where $`𝐑(𝗾)`$ is the tensor
$`𝐑(𝗾)`$ $`=`$ $`(q_0^2\left|𝐪\right|^2)𝐈+2\mathrm{𝐪𝐪}+2q_0𝐈\times 𝐪`$ (6)
$`=`$ $`\mathrm{𝐯𝐯}+\mathrm{cos}\theta (𝐈\mathrm{𝐯𝐯})+\mathrm{sin}\theta 𝐈\times 𝐯,`$ (7)
where $`\mathrm{𝐚𝐚}`$ is the parallel projector \[$`(\mathrm{𝐚𝐚})𝐛=(𝐚𝐛)𝐚`$\] and $`𝐈\times 𝐚`$ is the cross operator \[$`(𝐈\times 𝐚)𝐛=𝐚\times 𝐛`$\] \[4, §113\]. Equation (7) is the conventional tensor representation for a right-handed rotation of $`\theta `$ about an axis $`𝐯`$ through the origin \[4, §126\]. Equation (6) may be written in component form as
$`𝐑(𝗾)`$ $`=`$ (11)
$`\left(\begin{array}{ccc}12q_2^22q_3^2& 2q_1q_22q_0q_3& 2q_1q_3+2q_0q_2\\ 2q_2q_1+2q_0q_3& 12q_3^22q_1^2& 2q_2q_32q_0q_1\\ 2q_3q_12q_0q_2& 2q_3q_2+2q_0q_1& 12q_1^22q_2^2\end{array}\right).`$
The definition, eq. (5), gives $`R_𝗽(R_𝗾(𝐱))=R_{\mathrm{𝗽𝗾}}(𝐱)`$, so that $`\mathrm{𝗽𝗾}`$ corresponds to composing rotations (with the rotation by $`𝗾`$ performed first). We also find that $`R_𝗾=R_𝗾`$; i.e., $`𝗾`$ and $`𝗾`$ give the same rotation—changing the sign of $`𝗾`$ is equivalent to increasing $`\theta `$ by $`2\pi `$ in eq. (4). Unit quaternions satisfy $`q_0^2+q_1^2+q_2^2+q_3^2=1`$ and the quaternion representation of rotations are as points on a hypersphere $`𝕊^3`$ with opposite points identified. For future reference, we note that the (three-dimensional) area of $`𝕊^3`$ is $`2\pi ^2`$.
Because $`𝗾`$ and $`𝗾`$ give the same rotation, some care needs to be taken when comparing two orientations represented by $`𝗾_𝗮`$ and $`𝗾_𝗯`$. The rotation, $`𝗾=𝗾_𝗯\overline{𝗾_𝗮}`$, moves from $`𝗾_𝗮`$ to $`𝗾_𝗯`$. When inverting eq. (4) to determine the rotation angle $`\theta `$ between the two orientations, we should, if necessary, change the sign of $`𝗾`$ to ensure that $`q_00`$, so that $`\theta [0,\pi ]`$. A simple metric for closeness is given by $`\mathrm{cos}(\theta /2)=\left|𝗾_𝗮^𝖳𝗾_𝗯\right|`$.
Describing rotations with quaternions has a number of benefits. They offer a compact representation of rotations. Compared to Euler angles, they are free of singularities. Rotations may be composed more efficiently using quaternions than by matrix multiplication. Also in contrast to rotation matrices, it is easy to maintain a quaternion’s unit normalization (merely divide it by $`\left|𝗾\right|`$). However the chief benefit is that the representation of a rotation as point on $`𝕊^3`$ allows us to derive many important results concerning rotations in a simple coordinate-free way.
There is one application where the matrix representation of rotations is more efficient that the quaternion representation. If we wish to apply the same rotation to many points, then we should form the rotation matrix using eq. (6) and transform the points by matrix multiplication.
The conventional representation for rotations that is most closely allied to quaternions is the axis-angle representation, where the rotation is given by a vector $`𝐬=\theta 𝐯`$ which denotes a rotation of $`\theta =\left|𝐬\right|`$ about an axis $`𝐯=𝐬/\left|𝐬\right|`$. It is useful to have an analytic relation between the quaternion and axis-angle representations and this is provided by the quaternion exponential ,
$$\mathrm{exp}([0,𝐬/2])=𝗾,$$
(12)
where $`𝗾`$ is given by eq. (4), This definition of the exponential follows from its series expansion. Similarly the inverse operation is given by the quaternion logarithm
$$\mathrm{ln}𝗾=[\mathrm{𝟬},𝐬/\mathrm{𝟮}+\mathrm{𝟮}\pi 𝗻𝐯],$$
(13)
where $`n`$ is an integer.
It is useful here to make a distinction between “orientation” and “rotation”. We imagine that our molecule has some arbitrary but definite reference state. We apply a rotation and a translation (jointly referred to as a “displacement”) to this reference state and so bring the molecule to a new orientation and position (jointly referred to as a “configuration”).
## IV Least-squares fit
Given two conformations of the same molecule, it is often useful to be able to determine how close the conformations are. In order to do this, we can rigidly move one conformation so that it nearly coincides up with the other and then determine the difference in the positions of the corresponding atoms. Thus, if we are given two sets of atomic positions $`\{𝐱_k\}`$ and $`\{𝐲_k\}`$ with $`k[1,N]`$ together with a set of atomic “weights” $`\{w_k\}`$, we wish to determine the (rigid) displacement $`T`$ which minimizes
$$E=\frac{1}{W}\underset{k}{}w_k\left|𝐲_kT(𝐱_k)\right|^2,$$
(14)
where $`W=_kw_k`$. Here $`w_k`$ is merely a statistical weight of an atom—it is not necessarily related to the atomic mass. The two sets of atomic positions are ordered which presumes that we can identify corresponding atoms. (This is not necessarily a simple matter, if, for example, we are dealing with a molecule with several identical branches.) The displacement $`T=(𝗾,𝐝)`$ may be expressed as a rotation about an axis through the origin followed by a translation, i.e., $`T(𝐱)=R_𝗾(𝐱)+𝐝`$.
This problem has been considered by many authors and a review of various approaches is given by Flower . Using quaternions to describe the rotation leads to an elegant and robust solution. An early use of quaternions in this context is to solve the problem formulated by Wahba \[8; 9\], the determination of the attitude of a spacecraft given the directions of several objects relative to the craft. The resulting “$`𝗾`$-method” is described by Keat \[10, §A.3\] and by Lerner \[11, §12.2.3\] who both credit the invention of the method to Paul B. Davenport (1968). The generalization to matching points (as opposed to directions) was considered by Faugeras and Hebert who independently found the same method for determining the orientation. Their method was subsequently rediscovered by Horn , by Diamond , and by Kearsley . The derivation of Faugeras and Hebert is one of the clearest, and we briefly summarize it here including the straightforward generalization of including arbitrary weights $`w_k`$.
If we demand that the variation of $`E`$ with respect to $`𝐝`$ vanish, we find that
$$𝐝=𝐲R_𝗾(𝐱),$$
(15)
where $`\mathrm{}`$ denotes the sample average,
$$X=\frac{1}{W}\underset{k}{}w_kX_k.$$
(16)
Equation (14) may now be written as
$$E=\frac{1}{W}\underset{k}{}w_k\left|𝐲_k^{}R_𝗾(𝐱_k^{})\right|^2,$$
(17)
where $`𝐱_k^{}=𝐱_k𝐱`$ and $`𝐲_k^{}=𝐲_k𝐲`$. Using eq. (5), eq. (17), becomes
$$E=\frac{1}{W}\underset{k}{}w_k\left|[0,𝐲_k^{}]𝗾[\mathrm{𝟬},𝐱_𝗸^{}]\overline{𝗾}\right|^\mathrm{𝟮}.$$
(18)
Because, the norm of a quaternion is unchanged on multiplying it by a unit quaternion, we may right-multiply the kernel of eq. (18) by $`𝗾`$ to give
$$E=\frac{1}{W}\underset{k}{}w_k\left|[0,𝐲_k^{}]𝗾𝗾[\mathrm{𝟬},𝐱_𝗸^{}]\right|^\mathrm{𝟮}.$$
(19)
We need to minimize eq. (19) subject to the constraint $`\left|𝗾\right|=1`$. Because the kernel is linear in $`𝗾`$, it can be written as
$$[0,𝐲_k^{}]𝗾𝗾[\mathrm{𝟬},𝐱_𝗸^{}]=𝗔_𝗸𝗾,$$
(20)
where $`𝗔_𝗸`$ is a $`4\times 4`$ skew matrix
$$𝗔_𝗸=𝗔(𝐲_𝗸^{}+𝐱_𝗸^{},𝐲_𝗸^{}𝐱_𝗸^{}),$$
with
$`𝗔(𝐚,𝐛)`$ $`=`$ $`\left(\begin{array}{cc}0& 𝐛^𝖳\\ 𝐛& 𝐈\times 𝐚\end{array}\right)`$
$`=`$ $`\left(\begin{array}{cccc}0& b_1& b_2& b_3\\ b_1& 0& a_3& a_2\\ b_2& a_3& 0& a_1\\ b_3& a_2& a_1& 0\end{array}\right).`$
Substituting this into eq. (19), we obtain
$$E=\frac{1}{W}\underset{k}{}w_k𝗾^𝖳𝗔_𝗸^𝖳𝗔_𝗸𝗾=𝗾^𝖳𝗕𝗾,$$
(23)
where $`𝗕=𝗔_𝗸^𝖳𝗔_𝗸`$ is a $`4\times 4`$ symmetric matrix which has real eigenvalues, $`0\lambda _0\lambda _1\lambda _2\lambda _3`$. Setting $`𝗾`$ to the eigenvector corresponding to $`\lambda _0`$ gives the minimum value for $`E=\lambda _0`$.
In summary, the best fit is achieved by subtracting the mean positions from the original sets of points to give $`\{𝐱_k^{}\}`$ and $`\{𝐲_k^{}\}`$, forming the matrices $`𝗔_𝗸`$ and $`𝗕`$, and determining the minimum eigenvalue $`\lambda _0`$ of $`𝗕`$. The optimal rotation is given by setting $`𝗾`$ to the corresponding eigenvector of $`𝗕`$ and the optimal translation is found from eq. (15). The mean squared error for this fit is $`\lambda _0`$.
This procedure has two attractive features. The rotation obtained is a proper rotation (without an inversion); this is usually the desired result. Secondly, degenerate molecules are treated satisfactorily. For example if one or both of the sets $`\{𝐱_k\}`$ and $`\{𝐲_k\}`$ is collinear, then the best fit is no longer unique. The result will be that there will be multiple minimum eigenvalues of $`𝗕`$ with distinct eigenvectors. The general solution is obtained by setting $`𝗾`$ to a linear combination of these eigenvectors. The method does require finding the eigenvalues and eigenvectors of a $`4\times 4`$ matrix. However there are many numerical libraries \[16; 17; 18\] which solve such problems and the results are accurate to round-off for small symmetric matrices such as $`𝗕`$. A fast method of determining just the required eigenvector and associated eigenvalue in order to determine the attitude of a spacecraft is given in \[19, §III\]. However, in applications to molecular modeling, it is probably preferable merely to invoke a library eigenvector routine.
Horn considered including a scaling in the transformation $`T`$ in eq. (14). This is quite easily accommodated. However there seems little need to include such an effect in molecular modeling.
Diamond considers the case where inversions are allowed. This is easily achieved by substituting $`𝐱_k^{}`$ for $`𝐱_k^{}`$ in eq. (17). Equation (23) then involves a matrix $`𝗕^{}`$ where
$$𝗕^{}=\mathrm{𝟮}\left|𝐱^{}\right|^\mathrm{𝟮}+\left|𝐲^{}\right|^\mathrm{𝟮}𝗜𝗕.$$
(24)
Consequently the rotation giving the best inverted fit is the eigenvector with the greatest eigenvalue of $`𝗕`$, $`\lambda _3`$. Because the sum of the eigenvalues of $`𝗕`$ is its trace, $`4\left|𝐱^{}\right|^2+\left|𝐲^{}\right|^2`$, we can express the mean squared error for the inverted fit as $`\frac{1}{2}(\lambda _0+\lambda _1+\lambda _2\lambda _3)`$. Thus, once the eigenvalues of $`𝗕`$ have been computed we immediately determine whether the inverted fit will be better than the proper fit.
Coutsias et al. provide an interesting extension of this method. Suppose the atomic positions $`\{𝐱_k\}`$ represent a model of a molecule which depends on a set of parameters $`\{\alpha _i\}`$, for example, the torsion angles of a protein backbone. By considering the gradient of $`E`$ in parameter space $`E/\alpha _i`$, they provide a method for determining the parameter values which result in the best fit to a given crystal structure.
One other interesting consequence of the result for the best fit is that the rotation is not a continuous function of the configurations of the molecules. Let us suppose that $`\{𝐱_k\}`$ gives the position of the atoms in a molecule in some predefined configuration and suppose that $`\{𝐲_k\}`$ gives the atom positions during the course of a dynamical simulation of the molecule. If the forces acting on the atoms are finite then $`𝐲_k`$ is a $`C^1`$ function (twice differentiable). During the course of the deformation of the molecule, $`𝗕`$ and its eigenvalues change. In the typical case, the two smallest eigenvalues exchange roles and $`𝗾`$ switches from one direction in $`^4`$ to an orthogonal direction. This results in the orientation of the best fit changing discontinuously by $`180^{}`$.
In modeling a flexible molecule, it is frequently useful to separate the external degrees of freedom, namely position and orientation, from the internal degrees of freedom. This allows, for example, translational and rotational symmetry to the system to be enforced and correlations between the motions of atoms within a molecule to be studied. This begs the question of how best to define the position and orientation of a molecule. Taking the position to be the center of mass is often the obvious choice. The position (so defined) evolves according to Newton’s second law driven by the total force on the molecule. It is not possible to keep track of the orientation in an analogous fashion by integrating the total angular momentum, because flexible bodies can change their orientation with zero angular momentum—witness the ability of a cat always to land on its feet. A possible definition of the orientation is the best fit orientation to a reference conformation; i.e., we define $`𝗼_𝗥(𝗔)`$ as the best fit orientation, expressed as a quaternion, of the molecule in conformation $`A`$ relative to a reference conformation $`R`$. Here again this choice has the attractive feature that the whole molecule is included in the definition. There are two problems with this prescription. Firstly, the difference in orientations between two conformations $`A`$ and $`B`$ depends, in general, on the choice of reference conformation, namely
$$𝗼_𝗥(𝗕)\overline{𝗼_𝗥(𝗔)}𝗼_𝗦(𝗕)\overline{𝗼_𝗦(𝗔)}.$$
(This is easily demonstrated for simple triatomic molecules.) Thus this definition of orientation entails a degree of “arbitrariness” absent in our definition of position. A second more serious defect arises from the discussion in the previous paragraph. Recovering the actual configuration of the molecule from the orientation defined in this way is numerically unstable (by a flip of $`180^{}`$!) whenever the lowest eigenvalues cross. This would also lead to large and discontinuous apparent internal motions of the molecule with small changes in the atoms’ true positions. A better choice would therefore be to make the fit to some rigid (or nearly rigid) subcomponent of the molecule . Although this still yields an arbitrary definition of orientation (depending on the choice of reference subcomponent), the resulting orientation varies continuously under continuous deformations of the molecule. An extensive discussion of how to separate the orientation from the internal motions of a flexible molecule is given by Littlejohn and Reinsch .
## V Interpolating rotations
The power of the quaternion representation of rotations is evident when we consider the problem of interpolating between two orientations of a molecule. (This application might arise in the animation of a molecular simulation.) Suppose we wish to interpolate between $`𝗾_𝗮`$ and $`𝗾_𝗯`$. Because these quaternions and their interpolants lie on the unit sphere $`𝕊^3`$, the shortest path will be a great circle whose parametric equation is given by
$$𝗾(\varphi )=\frac{𝗾_𝗮\mathrm{sin}(\theta \varphi )+𝗾_𝗯\mathrm{sin}(\varphi )}{\mathrm{sin}(\theta )},$$
(25)
where $`\mathrm{cos}\theta =𝗾_𝗮^𝖳𝗾_𝗯`$. In the computer animation community this “spherical linear interpolation” operation is denoted by $`Slerp(𝗾_𝗮,𝗾_𝗯;𝘂)=𝗾(𝘂\theta )`$ . As $`\varphi `$ is increased from $`0`$ to $`2\pi `$, $`𝗾(\varphi )`$ becomes successively $`𝗾_𝗮`$, $`𝗾_𝗯`$, $`𝗾_𝗮`$, $`𝗾_𝗯`$, and finally returns to $`𝗾_𝗮`$. During this operation the corresponding 3-dimensional rotation has increased by $`4\pi `$. If $`𝗾_𝗮^𝖳𝗾_𝗯\mathrm{𝟬}`$, then $`0\varphi \theta `$ takes $`𝗾(\varphi )`$ smoothly from $`𝗾_𝗮`$ to $`𝗾_𝗯`$. If, on the other hand, $`𝗾_𝗮^𝖳𝗾_𝗯<\mathrm{𝟬}`$, then a shorter path is found with $`0\varphi \theta \pi `$ which takes $`𝗾(\varphi )`$ smoothly from $`𝗾_𝗮`$ to $`𝗾_𝗯`$.
Equation (25) is derived using simple geometrical arguments applied to $`𝕊^3`$ and the same result is obtained for the great-circle interpolation for $`𝕊^n`$. For $`𝕊^3`$, the result can also be expressed as
$$𝗾(\varphi )=(𝗾_𝗯\overline{𝗾_𝗮})^{\varphi /\theta }𝗾_𝗮.$$
This relation has the interpretation: rotate to $`𝗾_𝗮`$ and then rotate a fraction $`\varphi /\theta `$ to the path from $`𝗾_𝗮`$ to $`𝗾_𝗯`$. The operation $`𝗾^𝘂`$ is defined by
$$𝗾^𝘂=\mathrm{exp}(𝘂\mathrm{ln}𝗾).$$
In fact this interpolation scheme results in the molecule undergoing rotation at constant angular velocity. In order to show this, consider a body rotating at $`\omega `$ about a unit axis $`𝐯`$. The evolution of the orientation $`𝗾`$ satisfies the differential equation
$$\dot{𝗾}=[0,(\omega /2)𝐯]𝗾.$$
(26)
This is easily solved (e.g., by using finite differences and passing to the limit $`\delta t0`$) to give
$`𝗾(𝘁)`$ $`=`$ $`\mathrm{exp}([0,(\omega t/2)𝐯])𝗾(\mathrm{𝟬})`$
$`=`$ $`[\mathrm{cos}(\omega t/2),𝐯\mathrm{sin}(\omega t/2)]𝗾(\mathrm{𝟬}),`$
which agrees with eq. (25) with the substitutions $`\varphi =\omega t/2`$, $`𝗾_𝗮=𝗾(\mathrm{𝟬})`$ and $`𝗾_𝗯=[\mathrm{𝟬},𝐯]𝗾(\mathrm{𝟬})`$.
If we wish to interpolate between two configurations of a rigid molecule, we are free to specify a point, $`𝐱_0`$, in the reference molecule which will move with constant velocity. If the initial and final configurations are given by $`T_a=(𝗾_𝗮,𝐝_𝗮)`$ and $`T_b=(𝗾_𝗯,𝐝_𝗯)`$, with $`𝗾_𝗮^𝖳𝗾_𝗯\mathrm{𝟬}`$, then the required interpolation is achieved by increasing $`u`$ from $`0`$ to $`1`$ with the orientation given by $`𝗾(𝘂\theta )`$ and the translation given by
$$(𝐝_a+R_{𝗾_𝗮}(𝐱_0))(1u)+(𝐝_b+R_{𝗾_𝗯}(𝐱_0))uR_{𝗾(𝘂\theta )}(𝐱_0).$$
## VI Mean Orientation
The mean of directional quantities has frequently presented difficulties . Let us assume we have $`N`$ samples of some directional quantity with weights $`w_k`$ for $`k[1,N]`$ and $`_kw_k=W`$. In the case where the samples are angles (e.g., the dihedral angles of a molecular bond) or directions (e.g., the orientations of a diatomic molecule), there is a well established procedure \[25, §2.2.1, §9.2.1\]: express the directions as unit vectors in $`^2`$ or $`^3`$, $`𝐧_k`$, and determine $`𝐧`$ where we take the sample average according to eq. (16). Now the mean direction is given by $`𝐧=𝐧/\left|𝐧\right|`$, while $`1\left|𝐧\right|`$, a quantity lying in $`[0,1]`$, is the “circular variance” \[25, §2.3.1\] or “spherical variance” \[25, §9.2.1\]. Here $`\mathrm{}`$ is defined as a simple weighted arithmetical average, eq. (16), while $`\mathrm{}`$ denotes the physically relevant mean of a quantity.
This procedure cannot be directly applied to unit quaternions used to represent rotations because of the indistinguishability of $`\pm 𝗾`$. Instead, we view $`\{𝗾_𝗸\}`$ as axes \[25, §1.1, §9.1\] in $`^4`$, and define $`𝗾`$ as the unit quaternion about which the weighted moment of inertia of $`\{𝗾_𝗸\}`$ is minimum \[26, §3\]. Thus we wish to minimize
$`L`$ $`=`$ $`{\displaystyle \frac{1}{W}}{\displaystyle \underset{k}{}}w_k𝗾_𝗸^𝖳(𝗜𝗾𝗾^𝖳)𝗾_𝗸`$
$`=`$ $`{\displaystyle \frac{1}{W}}{\displaystyle \underset{k}{}}w_k𝗾^𝖳(𝗜𝗾_𝗸𝗾_𝗸^𝖳)𝗾`$
$`=`$ $`𝗾^𝖳(𝗜\mathrm{𝗾𝗾}^𝖳)𝗾.`$
The minimum value of $`L`$ is given by the minimum eigenvalue of $`𝗜\mathrm{𝗾𝗾}^𝖳`$ and $`𝗾`$ is corresponding eigenvector. The resulting $`L`$, which is a quantity lying in $`[0,\frac{3}{4}]`$, then provides a measure of the variance of the rotations. This definition of the mean has a number of desirable properties: it is invariant when the signs of the $`𝗾_𝗸`$ are changed; it is independent of the order of the samples; and it transforms properly if the samples are transformed.
This prescription can also be applied to determine the mean direction of objects whose symmetry makes $`𝐧`$ and $`𝐧`$ indistinguishable (for example, the orientation of the diatomic molecule $`\mathrm{N}_2`$).
Suppose we wish to determine the mean configuration of a rigid molecule, i.e., the mean of $`\{T_k=(𝗾_𝗸,𝐝_𝗸)\}`$. We are free to choose a point $`𝐱_0`$ in the reference molecule whose position in the mean configuration coincides with its mean position. (Compare this with the discussion of interpolating configurations in the previous section.) A suitable definition for the mean configuration is then
$$T=(𝗾,𝐝+𝗥_𝗾(𝐱_\mathrm{𝟬})𝗥_𝗾(𝐱_\mathrm{𝟬})).$$
(27)
Frequently, we need more precise information about the distribution of configurations than its variance. We might need to know how much the rotation about different axes are correlated or whether rotational and translational motions are coupled. It is also desirable to be able to fit model distributions to a set of samples. For these purposes, it is useful to be able to map rotations onto $`^3`$ so that standard statistical tools can be employed. We require that the mapping be measure preserving (constant Jacobian) to simplify the use of the transformed rotations.
We have already introduced the axis-angle representation of rotations. We may make the restriction $`\left|𝐬\right|\pi `$ and so map the hemisphere $`q_00`$ of $`𝕊^3`$ onto a ball of radius $`\pi `$ in $`^3`$. Unfortunately, the mapping, eq. (12), does not have constant Jacobian. We can correct this by defining a new “turn” vector $`𝐮`$ with the properties
$`𝐮`$ $``$ $`𝐬,`$ (28a)
$`\left|𝐮\right|`$ $`=`$ $`\left({\displaystyle \frac{\left|𝐬\right|\mathrm{sin}\left|𝐬\right|}{\pi }}\right)^{1/3}.`$ (28b)
This is an extension of the Lambert azimuthal equal-area projection providing a measure-preserving mapping of the hemisphere $`q_00`$ of $`𝕊^3`$ onto the unit ball $`𝔹^3`$. Equation (28) is well behaved at the boundary, $`\left|𝐮\right|=1`$; however on this boundary antipodal points are identified. The inverse mapping has an infinite derivate at $`\left|𝐮\right|=\sqrt[3]{2n}`$ for integer $`n0`$ which corresponds to shells in $`𝐮`$ space which map to the origin. This inverse of eq. (28) is easily implemented via Newton’s method supplemented by a Taylor series at the origin and at $`\sqrt[3]{2n}`$.
This mapping was introduced to allow distributions of orientations to be fit using a mixture of Gaussians . Given a set of sample orientations $`\{𝗾_𝗸\}`$, we compute the mean orientation, $`𝗾`$. The deviations of the samples from the mean are then given by the rotations $`\{𝗾_𝗸\overline{𝗾}\}`$ and these are mapped to a set of turns $`\{𝐮_k\}`$. Because these are points in $`^3`$, we may fit them with a 3-dimensional Gaussian with zero mean and with covariance matrix $`𝐮^𝖳𝐮`$.
This procedure can be extended to fits of molecular configurations. In this case the deviations from the mean configuration, eq. (27), is mapped into a point in $`^6`$; the resulting Gaussian fit will capture the correlation between the translational and rotational degrees of freedom.
In closing this section, we mention an alternative way of fitting quaternion orientational data with analytic functions, namely in terms of spherical harmonics. The normal (3-dimensional) spherical harmonics can be generalized to 4 (and higher) dimensions \[29; 30\] and the orthogonality relation allows the coefficients of the harmonics to be computed simply. The $`\pm 𝗾`$ symmetry merely results in the odd harmonics dropping out. However in typical molecular interactions, the relative orientation of the molecules is tightly constrained which means that a large number of spherical harmonics will be needed to represent the orientational distribution. For such applications, a representation in terms of localized functions, such as Gaussians, is preferable.
## VII Random orientation
In Monte Carlo simulations it is sometimes necessary to select a molecule with a random and uniform position and orientation, for example, when attempting to insert a molecule into a simulation box during a grand canonical simulation . Choosing a random position is straightforward. However, we need to be careful to select the random orientation uniformly or else detailed balance will be violated (when balancing insertions and deletions). One possibility is to choose a random turn $`𝐮`$ in $`𝔹^3`$ and to convert this to a quaternion. However, it is much simpler to sample directly in quaternion space.
Let us first establish the requirement for “uniform” sampling of orientations. Composing 3-dimensional rotations is carried out by the multiplication of unit quaternions; but we know that $`\mathrm{𝗽𝗾}=𝗽^𝖳𝗤(𝗾)`$, where $`𝗤(𝗾)`$, given in eq. (3), is orthonormal if $`𝗾`$ is a unit quaternion. Thus 3-dimensional rotations map into a rigid rotation of $`𝕊^3`$; a uniform density on $`𝕊^3`$ is invariant to such rotations. It follows that the task of sampling a random orientation reduces to picking a random unit quaternion uniformly on $`𝕊^3`$.
Marsaglia provides one prescription: select $`x_1`$ and $`y_1`$ uniformly in $`(1,1)`$ until $`s_1=x_1^2+y_1^2<1`$; similarly, select $`x_2`$ and $`y_2`$ uniformly in $`(1,1)`$ until $`s_2=x_2^2+y_2^2<1`$; then
$$𝗾=[𝘅_\mathrm{𝟭},𝘆_\mathrm{𝟭},𝘅_\mathrm{𝟮}\sqrt{(\mathrm{𝟭}𝘀_\mathrm{𝟭})/𝘀_\mathrm{𝟮}},𝘆_\mathrm{𝟮}\sqrt{(\mathrm{𝟭}𝘀_\mathrm{𝟭})/𝘀_\mathrm{𝟮}}]$$
is uniformly distributed on $`𝕊^3`$.
A more transparent and symmetric method (which generalizes to sampling points on $`𝕊^n`$ \[34, §7.1\]) is to pick 4 normal deviates $`g_i`$ for $`i[0,4)`$ and to set
$$𝗽=[𝗴_\mathrm{𝟬},𝗴_\mathrm{𝟭},𝗴_\mathrm{𝟮},𝗴_\mathrm{𝟯}],𝗾=𝗽/\left|𝗽\right|.$$
Although this method is less efficient than Marsaglia’s, the overall impact in the context of a molecular simulation is probably tiny. Both of these methods return points uniformly over the whole of $`𝕊^3`$ rather that over just one hemisphere. In most applications, this is of no consequence.
Other representations of rotation yield more complex rules for obtaining random orientations. For example, with Euler angles, we would sample uniformly the first and third angles and the cosine of the second angle. If the orientation is given in axis-angle space, $`𝐬`$, then the axis, $`𝐬/\left|𝐬\right|`$, should be chosen uniformly on $`𝕊^2`$, and the rotation angle, $`\left|𝐬\right|`$, should be sampled from $`[0,\pi ]`$ with probability $`(2/\pi )\mathrm{sin}^2(\left|𝐬\right|/2)`$. Of course, this simplifies when $`𝐬`$ is transformed to $`𝐮`$ space, eq. (28), leading to a uniform distribution in $`𝔹^3`$.
A related problem is selection of random rotational moves for use in a Monte Carlo simulation . This method requires that detailed balance be satisfied, which, in the absence of torque bias, means that the probability of selecting the new orientation is symmetric under interchange of old and new orientations. Because we are typically interested in small changes in orientation, it is most convenient to select the rotation in axis-angle space as $`\mathrm{exp}([0,𝐬])`$ and to set the new orientation
$$𝗾^{}=\mathrm{exp}([\mathrm{𝟬},𝐬])𝗾,$$
where $`𝐬`$ is selected from an even distribution, $`p(𝐬)=p(𝐬)`$. (This result follows because the Jacobian factor is even in $`𝐬`$.) Usually, we wish the choice of rotation axis to be isotropic, and in that case we have $`p(𝐬)=p(\left|𝐬\right|)`$. Thus we might select $`𝐬`$ uniformly in a sphere of radius $`\mathrm{\Delta }`$. Rao et al. select $`\left|𝐬\right|`$ uniformly in $`[0,\mathrm{\Delta }]`$ (which results in a distribution which is singular at the origin in $`𝐬`$ space). An attractive choice of distribution is a 3-dimensional Gaussian
$$p(𝐬)=\frac{\mathrm{exp}(\frac{1}{2}\left|𝐬\right|^2/\mathrm{\Delta }^2)}{(2\pi )^{3/2}\mathrm{\Delta }^3}.$$
Not only is this simple to sample from, but it allows torque bias to be included in a simple manner. Torque bias is implemented by multiplying the a priori probability of selecting a move by $`\mathrm{exp}(\lambda \beta 𝐭^𝖳𝐬)`$, where $`\beta `$ is the inverse temperature, $`\lambda `$ is a constant (usually taken to be $`\frac{1}{2}`$), and $`𝐭`$ is the torque on the molecule. If the “starting” distribution is a Gaussian then the torque-bias factor merely shifts the Gaussian to give
$$p(𝐬)=\frac{\mathrm{exp}(\frac{1}{2}\left|𝐬\lambda \beta \mathrm{\Delta }^2𝐭\right|^2/\mathrm{\Delta }^2)}{(2\pi )^{3/2}\mathrm{\Delta }^3}.$$
This offers two simplifications over the original procedure : (a) it is trivial to sample from a shifted Gaussian; and (b) the acceptance probability, which involves the ratio of the forward and reverse a priori probabilities, is also easy to compute and, in particular, it does not require the evaluation of a normalization factor for the distribution. Similar considerations obviously apply to the application of force bias for translational moves, as has been discussed by Rossky et al. . Indeed, in the case of moving molecules, we would naturally perform a combined translational and orientational move applying both force- and torque-bias simultaneously. There are often strong gradients in the forces in molecular simulations and a direct application of force bias in this case can lead to poor sampling because certain transitions are effectively disallowed. In such cases, it is prudent to limit the effect of the bias by limiting the shift in the Gaussian, if necessary, to ensure that there a finite probability (at least 5–10%, say) of the sampled move being in the opposite direction to the force. This ensures that the molecule can effectively explore configuration space because small steps are always permitted and it provides a simpler “safety” mechanism than the distance scaling of $`\lambda `$ proposed by Mezei .
Finally, some care needs to be taken to treat the possibility of the orientation “wrapping” around. Suppose the sampled $`𝐬`$ has $`\left|𝐬\right|>\pi `$, then the resulting orientation is identical to the wrapped one, $`𝐬2\pi 𝐬/\left|𝐬\right|`$. To ensure that detailed balance is maintained, the acceptance probability should use the a priori probability for the reverse move $`𝐬`$ (rather than the negative of the wrapped move). A simple expedient for avoiding this problem is simply to reject any move with $`\left|𝐬\right|>\pi `$.
## VIII Grids for orientation
In many contexts, it is important to be able to represent the independent variables for a problem on a grid. It is therefore useful to be able to map orientations onto a grid. Possible applications are binning molecular data, implementing cavity bias in orientation , performing systematic searching of orientations (where the goal is to provide more regular coverage of orientation space than is achieved by random sampling), and performing integrals over orientation by numerical quadrature . Our goal is to provide a simple rule for covering orientation space with a grid while ensuring that the grid elements are approximately of equal volumes and are not unduly distorted. Here again, representing the orientation as a quaternion provides a reasonable solution.
Recall that unit quaternions lie on a hypersphere $`𝕊^3`$. Positions on $`𝕊^3`$ can be determined by 3 angle-like variables. However these are a poor basis for a grid because of singularities in the resulting coordinate system. Instead let imagine surrounding $`𝕊^3`$ by a tesseract (the 4-dimensional analogue of the cube) of edge length 2. This consists of 8 cells which are $`2\times 2\times 2`$ cubes tangent to $`𝕊^3`$. An exemplary cell is given by $`𝗽`$ with $`p_0=1`$, $`\left|p_{i0}\right|1`$. We need only consider half of the cells of the tesseract because of the identification of $`\pm 𝗽`$. Thus we choose to consider the four cells for which one of the components of $`𝗽`$ is $`+1`$.
This then forms the basis for a cubical grid for orientation space. This is attractive because cubical grids are simple to index into; they are easy to refine; they have an metric factor which is easy to compute; etc. The overall “wastefulness” of this grids relative to a cubic grid within a domain of $`^3`$ is given by the ratio of the volume of four cells of the tesseract ($`4\times 2^3`$) to the area (really a volume) of a hemisphere of $`𝕊^3`$, i.e., $`32/\pi ^23.24`$. This might seem rather profligate. However, if we managed to arrange the grid around the $`𝕊^3`$ without any wastage, the grid edge would be reduced by a factor of only $`\sqrt[3]{3.24}1.48`$.
Let us divide each of the cells of the tesseract into $`M^3`$ grid cubes (of side $`2/M`$). These cubes can then be projected to $`𝕊^3`$ by scaling $`𝗽`$ to a unit quaternion. This operation scales the volume of each of the grid cubes by $`\left|𝗽\right|^4`$—a factor of $`\left|𝗽\right|^3`$ is due to scaling a volume element linearly by $`\left|𝗽\right|^1`$ and the last factor of $`\left|𝗽\right|^1`$ arises from the distortion of the cube during this operation. The maximum scaling occurs at the corners of the tesseract, e.g., $`𝗽=[\mathrm{𝟭},\mathrm{𝟭},\mathrm{𝟭},\mathrm{𝟭}]`$, where $`\left|𝗽\right|=2`$, so that range of volumes for the grid elements is 16 with the maximum distortion being a factor of 2. Mapping between an arbitrary orientation $`𝗾`$ and a point in the grid is then achieved as follows. We identify the component $`q_l`$ of $`𝗾`$ which is largest in absolute value and set $`𝗽=𝗾/𝗾_𝗹`$, giving $`p_l=1`$ and $`p_{il}[1,1]`$. The grid then consists of $`4\times M\times M\times M`$ elements. The resolution of the grid, given by the maximum change in orientations between neighboring grid cells, is approximately $`4/M`$. (We need to multiply the grid cube edge by 2 to obtain the equivalent rotation angle, because, from eq. (4), we have $`𝗾=[\mathrm{𝟭},𝐯\theta /\mathrm{𝟮}]`$ for $`\theta `$ small.)
When the application is quadrature, it is natural to evaluate the function and to compute the metric factor $`\left|𝗽\right|^4`$ at the centers of the grid cubes. For binning, we assign the samples to the grid cube in the obvious way and again use the grid center to compute the metric factor to obtain a sample density.
The cubical grid defined above is suitable for quadrature and searching where the cost of function evaluations is small. Sometimes, however, the cost of function evaluations is so high that it is desirable to find an “optimal” set of grid points. For integrations over $`𝕊^2`$, this is a well-studied problem and various integration grid have been given that ensure accuracy to high order . For $`𝕊^3`$, various spherical $`t`$-designs are known \[42; 43\]. A $`t`$-design is a set of points on the sphere such that the average of a polynomial of degree $`t`$ over the sphere is given by averaging the value of polynomial at those points. Unfortunately $`t`$-designs for $`𝕊^3`$ are only known for $`t`$ up to 11 with the 11-design corresponding to 60 orientations.
In order to provide a denser coverage of the sphere we propose the following strategy: Consider $`N`$ sample orientations, corresponding to $`2N`$ points on $`𝕊^3`$. Define a “covering radius”, $`\alpha `$, as the maximum rotation needed to align an arbitrary orientation with one of the sample orientations. The “coverage”, $`c`$, is defined by the ratio of the area of $`2N`$ spherical caps of rotational extent $`\alpha `$ to the total area of $`𝕊^3`$, i.e.,
$$c=\frac{N(\alpha \mathrm{sin}\alpha )}{\pi }$$
(29)
\[compare with eq. (28b)\]. For a given $`N`$, the optimal configuration of sample orientations is obtained by minimizing $`\alpha `$—this gives the “thinnest” coverage, $`c`$. Finally, we weight each sample point according to the fraction of orientational space which is closest to it (i.e., in proportion to the volumes of the Voronoi cells); and we set a secondary goal of minimizing the variation in the weights. We expect the resulting sample points and weights to provide robust and accurate estimates of orientational integrals—particularly of experimentally or numerically determined functions which are bounded but which may not have bounded derivatives. The sample points are also suitable for searching orientation space optimally. Finding such optimal sets of points is difficult in practice. So, here, we propose some sets based on the regular 4-dimensional polytopes \[44; 45\], with the results summarized in table 1.
The 24-orientation set is obtained by placing two 24-cells (or icositetrachora) in their mutually dual configurations to give the set
8 permutations of $`[\pm 1,0,0,0],`$ (30a)
16 permutations of $`[\pm {\displaystyle \frac{1}{2}},\pm {\displaystyle \frac{1}{2}},\pm {\displaystyle \frac{1}{2}},\pm {\displaystyle \frac{1}{2}}],`$ (30b)
24 permutations of $`[\pm {\displaystyle \frac{1}{\sqrt{2}}},\pm {\displaystyle \frac{1}{\sqrt{2}}},0,0].`$ (30c)
(Each orientation is counted twice here because of the identification of $`\pm 𝗾`$.) The corresponding Voronoi tessellation is a truncated-cubic tetracontaoctachoron (or 48-cell) which consists of 48 regular truncated cubes . The set of orientations, eq. (30), is the direct symmetry group for the cube.
The vertices of the 600-cell (or hexacosichoron) are given by eqs. (30a) and (30b) together with
96 even permutations of $`[\pm {\displaystyle \frac{\sqrt{5}+1}{4}},\pm {\displaystyle \frac{\sqrt{5}1}{4}},\pm {\displaystyle \frac{1}{2}},0].`$
In this case, the Voronoi tessellation is the dual of the 600-cell, namely the 120-cell (or hecatonicosachoron). Because the Voronoi cells are dodecahedra which are nearly spherical, the resulting 60 orientations gives a particularly thin covering of orientation space. A good covering is also provided by adding the centers of the tetrahedral cells of the 600-cell.
For comparison, we list in table 1 the data for some of the ZCW3 orientation sets used by Edén and Levitt . These are obtained by the taking sets of points appropriate for integrating of a periodic unit cube \[49; 50; 51\] and mapping this set to the space of 3 Euler angles. There are two potential problems with this approach: (1) even though the metric of orientation space is treated properly, the mapping from Euler angles to orientation space is not distance-preserving and we expect this to degrade the properties of a mesh; and (2) because one of the Euler angles is not periodic, functions in orientation space do not obey the constraints assumed in constructing the sets of sample points. (More complete data for the ZCW3 sets is available in .)
Finally, table 1 provides various strategies for constructing an arbitrarily fine grid. We start with gridding the tesseract on which we easily impose a cubical grid (see above). However, the optimal sphere covering of $`^3`$ is body-centered cubic , and such a grid results in a thinner covering. Still better coverings can be found by starting with the 48-cell which has a typical cell (a truncated cube),
$$p_0=1,\left|p_{i0}\right|\sqrt{2}1,\left|p_1\right|+\left|p_2\right|+\left|p_3\right|1.$$
(31)
The other cells are obtained by multiplying $`𝗽`$ by the members of eq. (30). A cubic or body-centered cubic grid can easily be placed within each cell. For example, a body-centered lattice can be obtained with
$$p_0=1,𝐩=[k,l,m]\delta /2,$$
subject to the constraint eq. (31), where $`k`$, $`l`$, and $`m`$ are either all even integers or all odd integers. Table 1 gives three examples of such grids. The disadvantage of grids in 48-cells is that care must be taken to treat the faces of the cell correctly. The triangular faces of the truncated cubes slice cut through the grid cells at an angle and the octagonal faces fit together with a $`45^{}`$ twist. It is therefore necessary to resort to numerical methods to determine the volume of the Voronoi cells near the faces. The resulting data for the weights and examples of other body-centered cubic grids in the 48-cell with $`\alpha 0.65^{}`$ are given in .
One special searching problem is determining the volume of the smallest rectangular box (whose edges are parallel to the coordinate axes) into which a given molecule fits. This problem arises in the study of a single protein bathed in a solvent. In order to eliminate boundary effects, it is possible to construct a periodic system and, for efficiency, we wish the volume of the periodic cell to be minimum. We can solve this problem by systematically sampling over all orientations using our grid. However, because of the symmetries of a cube, eq. (30), there are 24 equivalent orientations which minimize the volume and we can restrict the search to $`1/24`$ of orientation space by searching only in eq. (31). We should point out that for the purposes of mimicking a single solute molecule in a solvent with a periodic system, the “best” computational box is not given by fitting a single image of the solute into a box but rather by the more challenging problem of optimally fitting the solute molecules into its neighboring images .
The emphasis in the section is on covering all orientation space with a grid. In many molecular modeling applications, the orientation may be quite restricted, e.g., when considering the orientation of a ligand in a protein binding pocket, and we may elect to restrict the integration (or search) to a set of orientations which differ from the mean rotation by at most $`\mathrm{\Theta }`$. If we express the deviation from the mean as a turn vector, eq. (28), integrations may be carried out in (three-dimensional) turn space with the range of integration restricted to the ball $`\left|𝐮\right|\sqrt[3]{(\mathrm{\Theta }\mathrm{sin}\mathrm{\Theta })/\pi }`$. Because the mapping to turn space is volume preserving, the integrals are exact. In addition, provided that $`\mathrm{\Theta }\pi /2`$, the mapping to turn space entails little distortion ($`12\%`$) and standard numerical methods for integrating in a ball $`𝔹^3`$ can be used.
## IX Discussion
Quaternions are an ideal “fiducial” representation of orientation in a molecular simulation. They provide an economical format for program input and output and as the internal representation of orientation. There is little redundancy in the representation—there is just the normalization constraint on its four elements and this is easily tested and corrected. At a given numerical precision, quaternions cover orientation space uniformly. Most operations involving orientation can be carried out directly and efficiently with quaternions and they can be converted to other representations as needed. The basic operation of composing rotations is most cheaply performed with quaternions. On the other hand, if we need to rotate a large molecule it is quicker to convert the quaternion to a rotation matrix, eq. (11), and to perform matrix-vector multiplication than to apply eq. (5) directly.
In comparison, other representation suffer serious drawbacks. Rotations cannot be easily composed when expressed as Euler angles. Picking a random orientation is more awkward when rotation matrices are used. In neither of these representations is it easy to interpolate between two orientations or to compute the mean orientation.
Although quaternions may be unfamiliar to some readers, we only needed to use quaternion algebra in the rule for composing rotations and in deriving the least squares fit. In carrying out the other tasks, we just used the fact that rotations are represented by opposite points on $`𝕊^3`$ and this provides a “natural” metric for rotations. In working with $`𝕊^3`$, we are able to carry over geometrical concepts from $`𝕊^2`$ or use straightforward extensions from Euclidean space, $`^3`$, to $`^4`$.
A curious and non-obvious property of rotations which is evident from their representation on $`𝕊^3`$, with $`\pm 𝗾`$ identified, is that rotations do not form a simply connected group. Thus, if we rotate an object by $`360^{}`$ it returns to its original orientation but with the sign of $`𝗾`$ changed. This means that we cannot continuously deform the path that the object took to reduce it to a point. However, we can do this if we rotate the object by $`720^{}`$. This property of rotations is an immediate consequence of their representation as a pair of points $`\pm 𝗾`$ on $`𝕊^3`$ and good visual illustrations of this property are provided by the Dirac belt trick and the Phillipine wine dance .
## Acknowledgment
This work was supported by the U.S. Army Medical Re-search and Materiel Command under Contract No. DAMD17-03-C-0082. The views, opinions, and findings contained in this report are those of the author and should not be construed as an official Department of the Army position, policy, or decision. No animal testing was conducted and no recombinant DNA was used.
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# THE CLEO RICH DETECTOR
## 1 INTRODUCTION
The CLEO II detector was revolutionary in that it was the first to couple a large magnetic tracking volume with a precision crystal electromagnetic calorimeter capable of measuring photons down to the tens of MeV level . CLEO II produced many ground-breaking physics results, but was hampered by its limited charged-hadron identification capabilities that were provided by a combination of dE/dx and time-of-flight measurements.
The CLEO III detector was designed to study decays of $`b`$ and $`c`$ quarks, $`\tau `$ leptons and $`\mathrm{{\rm Y}}`$ mesons produced in $`e^+e^{}`$ collisions near 10 GeV center-of-mass energy. CLEO III is an upgraded version of CLEO II, as the magnet, the calorimeter and the muon system were kept. It contained a new four-layer silicon strip vertex detector, a new wire drift chamber and a particle identification system based on Cherenkov ring imaging described herein. Information about CLEO is available elsewhere .
Design choices for particle identification were limited by radial space and the necessity of minimizing material in front of the CsI crystal calorimeter. The CsI imposed a hard radial outer limit, and the desire for maintaining excellent charged particle tracking imposed a radial lower limit, since at high momentum the error in momentum is inversely proportional to the square of the track length. The particle identification system was allocated only 20 cm of radial space, and this limited the technology choices. To retain the superior performance of the calorimeter, a material thickness of only 12% of a radiation length was allowed.
The CLEO III installation including the RICH detector occurred in the summer of 1999 and was used to study physics in the $`e^+e^{}`$ center-of-mass energy region around 9-12 GeV. The two-ring machines KEKB and PEP-II produced much more luminosity than CESR, and after some excellent initial results on $`B`$ decays and Upsilon spectroscopy, it became clear that there was much physics to be explored in charm decays and studies of charmonium. The transformation from CLEO III to CLEO-c was made in 2003, when the CESR accelerator began operating in the 3-5 GeV center-of-mass energy region. The inner double-sided silicon detector was replaced with a wire drift chamber that has significantly less material and is much better suited to the CLEO-c physics program; in any case, the silicon was showing effects of premature radiation damage. The magnetic field was lowered from 1.5 T to 1 T, mainly for accelerator related considerations.
The plan of this paper is as follows: The basic RICH detector design will be delineated, followed by descriptions of the individual components, namely the crystals, multi-wire chambers, superstructure, readout electronics, and support subsystems. Issues in the long-term operation of the RICH are then discussed. Finally, some details of the data analysis technique and physics performance are presented.
## 2 DETECTOR OVERVIEW
### 2.1 Design Choices
The severe radial spatial constraint forces the design to have a thin, few-cm detector for Cherenkov photons and a thin radiator. Otherwise the photons have too little distance to travel and it becomes very difficult to precisely measure the photon angles. In fact, the only thin, large-area photon detectors possible in our situation were wire chamber based, either using a reflective CsI photocathode or a gas mixture of methane (CH<sub>4</sub>) and triethylamine (TEA) vapor to convert ultraviolet photons.<sup>1</sup><sup>1</sup>1A TMAE-based photon detector would have been unacceptably thick in order to obtain the same detection efficiency, due to its long photon absorption length. Use of CsI would have allowed us to use a liquid freon radiator with quartz windows in the system and to work in the wavelength region from about 160–200 nm. However, at the time of decision, the use of CsI was far from proven and, in any case, would have imposed severe constraints on the construction process, which would have been both technically difficult and expensive. Thus we chose a multi-wire chamber filled with a mixture of CH<sub>4</sub> and TEA that uses Cherenkov photons in the vacuum ultraviolet (VUV) region, 135–165 nm, generated in a 1 cm thick LiF crystal.
The quantum efficiency of the CH<sub>4</sub> and TEA mixture peaks at $``$150 nm , as shown in Fig. 1, below the transmission cutoff for glasses and fused silica quartz. This short wavelength requires the use of alkali halide crystals as both the Cherenkov-emitting medium and the entrance window to the photosensitive volume . We chose LiF as the Cherenkov radiator because it has the lowest dispersion in the wavelength band of the CH<sub>4</sub>-TEA quantum efficiency . Transparent gases must be used between the radiator and the photon detector that are almost entirely free of O<sub>2</sub> and H<sub>2</sub>O, both of which have 1–10 Mbarn cross section for photon absorption below 175 nm .
Details of the design of the CLEO RICH have been discussed elsewhere . Here we briefly review the main elements. Cherenkov photons are produced in a LiF radiator. The photons then traverse a free space, an “expansion volume,” where the cone of Cherenkov light expands in size. The photons enter a detector consisting of spatially-segmented multi-wire chambers (MWC) filled with CH<sub>4</sub> gas mixed with TEA vapor, in which they are converted to electrons and multiplied. Finally, signals are picked up with sensitive low-noise electronics. No optical focusing elements are used; this is called “proximity-focusing” . The scheme is shown in the upper left of Fig. 2, while the placement in CLEO is shown in Fig. 3.
There are 30 individual photon detectors around the outer cylinder. They subtend the same azimuthal angle as the radiators, which are segmented into 30 rows and 14 rings on the inner cylinder. The expansion volume between the radiators and detectors is filled with pure N<sub>2</sub> gas. In fact, any transparent, pure gas could be used. Our choice was based on the availability of very pure gas from the super-conducting coil boil off.
The most distinctive features of the design of the CLEO RICH are: (1) ultra-thin photon detectors, with high spatial segmentation and almost full area coverage; (2) a large volume of high-purity, VUV-transparent crystals including a novel “sawtooth” geometry; and (3) exceptionally low-noise analog electronics readout.
### 2.2 Radiators
LiF was chosen over CaF<sub>2</sub> or MgF<sub>2</sub>, both of which are transparent in the VUV wavelength region, due to its lower dispersion leading to smaller chromatic error. Originally all the radiators were planned to be 1 cm thick planar pieces. However, since the refractive index of LiF at 150 nm is 1.5, all the Cherenkov light generated from tracks normal to the LiF surface would be totally internally reflected, as shown in Fig. 4 (top). We could have used these planar radiators, but we would have had to tilt them at about a 15 angle. Besides the obvious mechanical support problem, some tracks would produce Cherenkov photons that would cross through the entire thickness of another radiator tile causing a loss of efficiency and reconstruction problems. Instead, we developed novel radiators with a serrated top surface, called “sawtooth” radiators , as shown in Fig. 4 (bottom). Measured physical properties of the radiators have been described previously .
The overall radiator shape approximates a cylinder of radius 82 cm. Individual radiator crystals are placed in 14 coaxial rings of 30 crystals each, centered around the beam line and symmetrically positioned about the interaction point. The 30 crystals segments are parallel to the wire chambers. Inter-crystal gaps are typically 50–100 $`\mu `$m. The crystals are attached to the exterior surface of a 1.5 mm thick carbon fiber shell with a low outgassing epoxy.
The inner four rings are made of sawtooth radiators. Lengthy production time as well as cost limited our use of these novel objects.
### 2.3 Photon Detectors
The photon detectors have segmented cathode pads 7.5 mm (length) $`\times `$ 8.0 mm (width) etched onto G10 printed-circuit boards (PCBs). The pad array was formed from four individual boards, each with 24 $`\times `$ 80 pads, the latter split into two 24 $`\times `$ 40 pad sections with a 7 mm gap. Each board was individually flattened in an oven and then they were glued together longitudinally on a granite table where reinforcing G10 ribs were also glued on. The ribs have a box-like structure. There are 4 longitudinal ribs that traverse the entire length. Smaller cross ribs are placed every 12 cm for extra stiffening. The total length of the pad array is 2.46 m.
Wire planes were strung with 20 $`\mu `$m diameter gold plated tungsten with a 3% admixture of rhenium; the wire pitch was 2.66 mm, for a total of 72 wires per chamber. The wires were placed on and subsequently glued to precision ceramic spacers every 30 cm. The spacers extend 1 mm above the cathodes, and therefore are 3.5 mm from the CaF<sub>2</sub> windows. We achieved a tolerance of 50 $`\mu `$m on the wire to cathode distance. The spacers had slots in the center for the glue bead.
Eight 30 cm $`\times `$ 19 cm CaF<sub>2</sub> windows were glued together in precision jigs lengthwise to form a 2.4 m long window. Positive high voltage is applied to the anode wires, while negative high voltage is put on 100 $`\mu `$m wide silver traces deposited on the CaF<sub>2</sub>. The spacing between the traces is 2.5 mm, and they are connected together by a trace running across both edges. The pad plane is close to ground. The gain of the chamber depends on the amount of high voltage on wires and windows. We define “wire-gain,” in the normal manner, as the multiplication factor on a single electron. We run the system at typical wire-gains around 30,000. “Pad-gain” is our usual measure of multiplication and is calculated as being 75% of the wire-gain.
To maintain the ability of disconnecting any faulty part of a chamber, the high voltage is distributed independently to three groups of 24 wires and the windows are each powered separately.
### 2.4 Electronics
The position of Cherenkov photons is measured by sensing the induced avalanche charge on the cathode pad array. Since the pulse height distribution from single photons is expected to be exponential at low to moderate gas gains , the use of low noise electronics is required to ensure high detection efficiency. Pad clusters in the detector can be formed from single Cherenkov photons, overlaps of more than one Cherenkov photon, or charged tracks. In Fig. 5 we show the pulse height distribution for single photons and charged tracks. (Note, one ADC count corresponds to $``$200 electrons.) These can be separated by just examining the pulse height. The charged tracks give very large pulse heights because they are traversing $``$4.5 mm of the CH<sub>4</sub>-TEA mixture. We can distinguish somewhat between the charge due to single photons and two photons because of the pulse height shapes on adjacent pads.
To have as low noise electronics as possible, a dedicated VLSI chip, called VA\_RICH, based on a very successful chip developed for solid state applications, has been designed and produced for our application at IDE AS, Norway . We have fully characterized 3,600 64-channel chips, mounted on hybrid circuit boards. For moderate values of the input capacitance $`C_{in}`$, the equivalent noise charge measured $`ENC`$ is found to be about
$$ENC=130e^{}+(9e^{}/\text{pF})C_{in}.$$
(1)
The dynamic range of the chip is between 450,000 and 900,000 electrons, depending upon whether we choose a bias point for the output buffer suitable for signals of positive or negative polarity or we shift this bias point to have the maximum dynamic range for signals of a single polarity.
In our readout scheme we group 10 chips in a single readout cell communicating with data boards located in VME crates just outside the detector cylinder. Chips in the same readout cell share the same cable, which routes control signals and bias voltages from the data boards and output signals to the data boards. Two VA\_RICH chips are mounted using wire bonds on one hybrid circuit that is attached via two miniature connectors to the back of the cathode board of the photon detector.
The output of the VA\_RICH is transmitted to the data boards as an analog differential current, transformed into a voltage by transimpedance amplifiers and digitized by a 12-bit differential ADC. These receivers are part of very complex data boards which perform several important analog and digital functions. Each board contains 15 digitization circuits and three analog power supply sections providing the voltages and currents to bias the chips, and calibration circuitry. The digital component of these boards contains a sparsification circuit, an event buffer, memory to store the pedestal values, and the interface to the VME crate CPU.
Coherent noise is present. We eliminate this by measuring the pulse heights on all the channels and performing an average of the non-struck channels before the data sparsification step.<sup>2</sup><sup>2</sup>2 This algorithm is executed by a DSP located on the data boards before the data are sparsified. The pedestal width (r.m.s.) is reduced from 3.6 to 2.5 ADC counts after this coherent noise subtraction. The intrinsic noise of the system then is $``$500 electrons r.m.s.
## 3 CRYSTAL FABRICATION and TESTING
The design of the CLEO RICH demands a large number of high-quality, high-purity VUV-transparent crystals, used as the entrance window to the multi-wire photon detectors and as the solid Cherenkov radiator medium. In fact, nearly 420 kg of crystal was required, in the form of 240 full-sized windows and 420 radiator tiles. The full windows were 191.0 $`\times `$ 307.6 $`\times `$ 2.0 mm<sup>3</sup> in size, and the planar radiators were 174.3 $`\times `$ 169.8 $`\times `$ 10.0 mm<sup>3</sup>.
There were several significant issues in the production of these crystal pieces. First, the bulk absorption of VUV photons needed to be small, so high purity raw material was required. Second, the surface transmission needed to be appreciable, so high-quality polishing of the surfaces was required. These two elements are necessary to yield good transmission in the VUV region. Third, the finished plates needed to have large dimensions, so the active area of the detector remained high. Last, they needed to have good mechanical stability—free of cleaves, nascent cleaves, and geometric shape deformations—so the internal strain in the crystals needed to be minimal. This was particularly critical for the very thin 2 mm windows, which would cause the photon detector to fail if a crack occurred in situ. In order to meet these stringent requirements a R&D program in the areas of fluoride crystal growth and processing were undertaken.
### 3.1 Crystal Production
#### 3.1.1 Crystal Growth Process
Fluoride crystals were grown at the Optovac facilities using a modified Bridgman-Stockbarger method . In this technique, pure raw powder is packed in a graphite crucible which is heated in a vacuum furnace to a temperature far above its melting point. After a period of impurity removal, the crucible is reduced to a temperature just above the melting point, and then slowly lowered mechanically through a sharp thermal gradient into a second volume where the temperature is below the melting point. This effects a surface of solidification in the melt where the crystal lattice is actually grown. The resulting solid ingot is then held at a reduced temperature (roughly half that of the melting point ), at which it is annealed to reduce the bulk strain induced by the growth process as it is slowly brought to room temperature.
The growth process is the most important single step in crystal production, since it not only impacts the manufacturability directly, but the ultimate optical and mechanical properties of the finished crystal as well. Stockbarger recognized from the beginning the critical roles of both the high purity of the raw material and the sharpness of the thermal gradient between the two temperature regions. This latter impacts the design and operation of the furnace, as does the need for uniformity of temperature in the melt. Good control of the crystallization zone allows for low intrinsic strain growth, and for the impurities to segregate to the top of the ingot. Both are necessary for the large-area RICH crystals.
The actual temperature levels and durations in the thermal cycle were determined by repeated experimentation, and varied according to material used and ingot size. This optimization constituted an important line of development, in particular for the large 14.5 and 16-inch diameter ingots grown for CaF<sub>2</sub> windows.<sup>3</sup><sup>3</sup>3Large ingot growth also has important implications for other applications, such as 157 nm lithography for VLSI fabrication. The CaF<sub>2</sub> ingots were grown using a piece of single crystal as a seed, in order to reduce the number of domain boundaries between regions of different crystal orientation. This was never entirely successful, and all ingots grown were polycrystals.
The raw material used was either very pure grade LiF powder synthesized by Merck,<sup>4</sup><sup>4</sup>4 LiF Optipur (R) powder, Merck SA, Darmstadt D-64293, Germany. or crushed natural fluorite (CaF<sub>2</sub>), of sufficient purity for thin windows. In order to remove residual impurities from the fluorite, 2% PbO<sub>2</sub> was added as a getter that subsequently segregated from the melt. Any residual impurities reduce transmission significantly at 150 nm, since they will either create absorption (color) centers or scattering centers for light in the VUV region.
In order to monitor the ingot quality at growth, test bars were taken as a vertical slice from the periphery of each ingot. The transmission of these witness pieces from a representative early bad ingots containing relatively high impurity levels and good ingots produced after the growth process was optimized are shown in Fig. 6. This allowed the proofing of the ingots, as well as a diagnostic of growth problems.
#### 3.1.2 Crystal Machining
The machining process consisted of taking a single ingot and producing many “blanks” from it—each of which is a formed piece of crystal ready for polishing in order to have good optical transmission. To produce the blanks, the ingot is first mounted on a plaster base and sliced to approximate thickness on a band-saw using a diamond blade and a gravity-feed table. It is important to use a “well-cutting blade,”<sup>5</sup><sup>5</sup>5 A “well-cutting” blade is not only sharp (e.g., using a diamond abrasive), but is properly dressed and lubricated so as to prevent loading up with fine crystal silt as it cuts. Thus the blade will cut the ingot instead of rubbing and cracking it. This same criteria is applied to the other cutting tools and grinding wheels used. to thermalize the ingot and lubricant to room temperature, and to establish the proper feed rate of the blade through the ingot. Next, the slice is trimmed to approximate shape, and all edges are bevelled by hand. Then, the two faces of the slice are ground to within $``$150 $`\mu `$m of final thickness, using a Blanchard-type surface grinder with fine-grit superabrasive pellets. Finally, the four edges of the piece are ground to dimension using a SX CNC machine with a very fine diamond superabrasive wheel. The LiF radiators have two edges which are given a 6 angle.
For the case of windows and planar radiators, the blank is now finished and ready for polishing. For the case of sawtooth radiators, additional grinding steps (as well as a different polishing technique) are required, as described below.
The machining process is developed by determining the proper speeds and feed rates for each operation. Otherwise the crystal may be easily destroyed. Nevertheless, each of these machining steps does introduce some actual damage to the crystal. This consists of microscopic cleaves and irregularities at the surface and just below it, called “subsurface damage.” The goal of each step in the machining process is to cut away enough material so as to get under the subsurface damage created by the previous step, while itself creating at worst only a finer level of damage, which is to be removed by the subsequent operation. In this way, a piece of sufficient transmission is ultimately produced.
#### 3.1.3 Crystal Polishing
The polishing process consisted of placing a blank on a lapping table with a polyurethane polishing pad and a specific polishing media, such as fine-grit diamond powder suspended in glycol or a commercially-available polishing slurry. Fig. 7 shows a working polishing table. The blank is weighted but free to spin within a rotating holder on the rotating polishing pad, the net effect of which is randomized orbits of the fine abrasive with respect to the blank. This yields a uniform action as it cuts away stock (mechanically or chemically) to get under the subsurface damage. This procedure is repeated with a finer grit abrasive, in order to finish the surface.
This is a fine and somewhat delicately balanced procedure, taking many hours per surface. The effectiveness is governed principally by the total amount of stock removed, which is related to the amount of weight applied to the piece as it glides on the polishing pad. The rate of stock removal also depends on the material (CaF<sub>2</sub> or LiF) and on the orientation of the crystal. The effectiveness of different polishing times and weights was extensively studied.
#### 3.1.4 Production Yields
Nearly one thousand individual crystals were produced for the CLEO RICH. The early stages of the crystal production saw a very low yield of pieces. It was only after intense efforts developing all phases of production (growth, machining and polishing) that the yields grew to acceptable levels.
For CaF<sub>2</sub> production, the intrinsic strain in the thin large-area windows needed to be minimized. Optimization of the annealing cycle, major improvements to the machining process, and determination of proper weighting during polishing, gave rise to increased yields. The manufacturing yield of an ingot (i.e., the number of windows produced divided by the number possible in an ingot) grew from below 30% in the initial phases of the project to be routinely above 80%, and as high as 90%. In the wake of these improvements, the quality of the ingot (i.e., the fraction of windows produced that passed CLEO requirements) also rose to be close to 100%.
For LiF production, proper purity of the material needed to be maintained. Improvements parallel to those made for CaF<sub>2</sub>, as well as critical corrections to raw material handling, contributed to slowly increasing yields. For the first twenty LiF ingots grown, the manufacturing yield was under 40% with the quality about the same. Several episodes of material contamination then reduced the quality to effectively zero, as indicated in Fig. 6. This was due to insufficient segregation of impurities, or overabundance of impurities, both of which create color centers in the bulk material as well as enhanced domain boundaries. The latter creates “fault lines” along which the crystal may separate into multiple pieces, thereby becoming useless. As improvements were made during the second half of the project, the manufacturing yield grew to 70%, albeit with large fluctuations, due mainly to the continued presence of separating domain boundaries in some ingots. Also in later production, the quality of the ingots became quite high, far surpassing CLEO requirements even at 135 nm. This is also indicated in Fig. 6.
### 3.2 Crystals for Chamber Windows
There were 240 full-sized chamber window crystals required. A total of 300 were delivered, of which 272 were used in construction. To make the number of full-sized windows used, there were 393 individual crystals, of which 151 were full-sized CaF<sub>2</sub> crystals, 182 were half-sized CaF<sub>2</sub> crystals, and 60 were half-sized LiF window crystals. This small admixture of LiF windows was acceptable if the machining process succeeded, and sped up the chamber production schedule. An additional 16% were rejected for various reasons.
#### 3.2.1 Inspection and Testing
All window crystals were cleaned, inspected and tested individually. Fig. 1 shows typical transmission curves for a 2 mm thick CaF<sub>2</sub> window and for a 10 mm thick LiF planar radiator. The main features on these curves are the drop to zero transmission at the band edge for the crystals (120 nm for CaF<sub>2</sub> and 105 nm for LiF). There are also several possible absorption lines at 130, 142.5, and 175 nm from water impurities. The transmission above 180 nm rises to more than 90% where bulk impurities cause minimal loss and only surface reflections dominate. For speed of processing of the crystals, most parts were scanned over their surface at 3 wavelengths only, namely 135, 150, and 165 nm, for each piece on a grid of 30 points over its surface, using the VUV Spectrophotometer system at Syracuse University, described in Appendix A.
Fig. 8 shows the distribution of transmissions at each of these wavelengths for all window crystals. Variations in the transmission may result from bulk impurities which vary by ingot, from polishing non-uniformities, and from variations in the cleaning. The variation at 135 nm are clearly larger than at 165 nm, and reflect the greater influences of water and other impurities below 145 nm. On average, the transmission in the windows is 71.5, 80.1, and 85.3% at 135, 150, and 165 nm, respectively, which matches our design criteria of 72, 80, and 85% at these wavelengths. All windows were tested both before and after high voltage trace deposition.
The dominant mechanism for transmission loss was generally surface scattering or absorption, with bulk absorption accounting for $`<`$20% of the transmission loss. This was confirmed by transmission measurements of identically-prepared LiF and CaF<sub>2</sub> pieces of different thicknesses to separate the two effects.
On some crystals, cleaves were found to have propagated after production. In the case where the cleave was small and near an edge, the pieces were hand-worked to remove the cleave and get under the subsurface damage. If properly done, there was only a small probability that the cleave would reappear and propagate in a relevant amount of time. This was often successful and the best of these reworked pieces were used as chamber windows.
Fig. 9 shows the fraction of each window that was optically inactive. Factors which contribute to the inactive area include the window traces which typically contributed a 4% opacity, any mechanical defects such as missing corners, any reworked areas of cleave removal as mentioned above, or surface stains which were not easily cleaned with acetone. Additionally, due to production schedules and the detector installation deadline, the windows formed out of half-sized plates epoxied together in pairs had an extra glue joint resulting in an additional 2% loss of active surface area. The two peaks in Fig. 9 represent the whole and the half windows, with the tails in each distribution coming from defects, smudges, or coating problems.
#### 3.2.2 Deposition of Traces
The CaF<sub>2</sub> crystals serve a dual role as entrance windows into the wire chambers, as well as high voltage cathodes for these chambers. In order to apply a voltage to this plane, a 200 nm thick coating of nickel and silver was applied to the chamber side of the window. This coating was in the form of 100 $`\mu `$m wide strips spaced 2.54 mm apart, and connected at each of their ends by another metallized coating to act as a voltage distribution bus bar. During construction, a metal conductor was glued to this bus bar using conductive epoxy.
The metallized coating was performed by EMF;<sup>6</sup><sup>6</sup>6 Evaporated Metal Films Corp., Ithaca, NY 14850. the setup is shown in Fig. 10. The coating was a standard sputtering process, in which a stream of metal ions is created by placing loops of nickel or silver on tungsten filaments through which a 100 Amp current passes. The crystals and filaments are placed in a vacuum chamber for the sputtering process. The nickel was found to bond better to the CaF<sub>2</sub> surface, so a 50 nm layer of this was laid down first, followed by 150 nm layer of the silver. The strip features on the windows were defined by metal masks placed in contact with the windows which had 100 $`\mu `$m wide slots cut into them. These masks were held in contact with the crystal by placing the crystal on a machinist’s magnetic chuck and pulling on the steel mask with the field from the chuck. Careful preparation of the chucks were required to make them vacuum-compatible and prevent outgassing onto the crystals. The masks were heat-treated to make them flat enough to not rise more than $`25\mu `$m off the crystals. If the mask was not held in intimate contact with the crystals, the sputtered silver would tend to “feather” under the mask, making a broad coated area that was then optically opaque.
### 3.3 Crystals for Radiators
There were 420 full-sized radiator crystals required, of which 300 were of planar and 120 were of sawtooth geometry. A total of 436 were delivered, from which 420 were used. To make the full radiator array, there were 450 individual crystals, of which 270 were full-sized LiF planar crystals, 60 were half-sized LiF planar crystals, and 120 were full-sized LiF sawtooth crystals. Pairs of half crystals were alternately used as the end rings to make the radiator left-right symmetric. An additional 2% of the planar radiators were rejected for various reasons, whereas for sawtooth radiators there were 23% rejected, indicating the relative difficulty in production.
#### 3.3.1 Inspection and Testing
All radiator crystals were cleaned, inspected and tested individually. For planar radiators, the transmission was measured at 135 nm, 142 nm, 150 nm, and 165 nm, for each piece on a grid of $``$300 points over its surface, using the VUV Spectrophotometer system at SMU, described in Appendix A. On average, the transmission in the planar radiators is 65.5, 77.6, and 85.4% at 135, 150, and 165 nm, which matches our design criteria of 66, 77, and 85% at these wavelengths. However, the low value at 135 nm is misleading: it is an average which includes early ingots having low transmission due to an excess of impurities. Later ingots produced radiators with transmissions at 135 nm of up to 75%.
#### 3.3.2 Dielectric Coating of Radiator Crystals
Monte Carlo studies indicated that the Cherenkov angle reconstruction efficiency is enhanced if those radiator photons entering the RICH photon detectors which have first bounced off the bottom surface of a radiator crystal are suppressed. Correspondingly, we coated the bottom surfaces of all 120 sawtooth radiator crystals with a thin layer of polystyrene. This dielectric has an index of refraction well matched to that of LiF. Detailed tests confirmed that this material strongly absorbs at 150 nm.
#### 3.3.3 Techniques for Sawtooth Radiators
The unique geometry of the sawtooth radiator crystals cannot be accomplished using the fabrication methods described above. While the lower surface and edges of the sawtooth radiators are flat and may be polished using conventional orbital lapping techniques, special procedures were developed to produce high-quality VUV-polished faces in the “vee”-shaped grooves on the exit surface of the crystals. These new techniques could, in principle, be adapted to other new geometries.
Sawtooth radiator blanks are cut and ground to a thickness of 12.7 mm (i.e., 0.7 mm oversized) in a manner identical to plane radiators. The blanks are then mounted on a linear surface grinder to define the groove shapes. A set of ten 6-inch diameter grinding wheels are mounted on a single spindle. These wheels are edged with a 90 “vee” on their edges and rough-grit superabrasive. Grooves are ground into the upper surface of each blank to within 1.0 mm of the final depth. Since there are 19 grooves in a sawtooth piece, this is done in two operations. Following this step, the group of ten wheels is replaced with a single grinding wheel with a 96 included angle and a fine-grit bond. This single wheel is then used to finish-grind the last $``$0.9 mm of depth in each of the grooves, thereby minimizing groove-to-groove differences within a single radiator crystal. As with the planar pieces, the fine-grit operation was found to reduce subsurface damage from the grinding process. The wheels need to be maintained by dressing periodically to remove LiF build-up, and to re-define the profile of the edge which rounds after grinding. Care was taken to reduce any spurious vibrations in the spindle during this procedure.
Polishing is done using a conventional Bridgeport milling machine with automated travel. Unlike the grinding, which could address both left and right faces of each groove simultaneously, the polishing is applied to one face of all grooves first and then applied to the remaining face in turn. The milling head is set at a 42 angle from the vertical and a new head is mounted in place of a milling bit. The polishing head was a 6-inch diameter aluminum disk with a soft polishing pad attached on the bottom, as is used on the usual lapping tables. The crystals are polished by applying a small tool pressure to the groove faces with this rotating head, and wetting the part with a polishing slurry. The part passes back and forth several times under the rotating head which reaches down into the groove. Then the part is indexed to work the next groove. After all grooves are done, the tool pressure is re-adjusted and a second set of passes taken. The piece is rotated in order to polish the opposite faces on the crystal. This entire procedure is repeated with a smaller grit abrasive, in order to finish the piece.
In order to test the optical transparency of the sawtooth radiators, we compared them to a calibrated prism of $`42^{}`$ inclination angle which was conventionally polished and also met our transmission specifications.<sup>7</sup><sup>7</sup>7We fabricated two such prisms and measured the transmission of the two stacked together for a VUV beam incident at 15 angle to avoid total internal reflection. To compare a sawtooth radiator to the calibrated prism, we deflected the light incident on the flat face of the sawtooth or prism by $`15^{}`$ to avoid total internal reflection, and the light intensity through the sawtooth is compared to that of the prism. Using this method we scanned the left-faces of the teeth independent of the right-faces. Fig. 11 shows the results of a transmission scan from one side of a typical sawtooth crystal. It is possible to have a transmission that is better than our standard prism, hence a relative transmission greater than 100%. As can be seen from the figure, however, it was challenging to fully polish the sawtooth grooves all the way into the bottom of the valley and to the top of the peak. This manifests itself as a slight enhancement at the middle of the face, with a roll-off at the peak and valley. This roll-off is generally limited to 0.5–1.0 mm at the peak and valley, as indicated in the shape for each face shown in Fig. 11. A cumulative plot of these measurements, superimposing all faces for many sawtooth radiators, is shown in Fig. 12. The difference between the relative transmission and the average per face is plotted, where the average excludes a small region of roll-off at both ends of each face. This plot indicates the spatial uniformity of the polishing along the faces of the teeth. Fig. 13 gives the distribution of the average relative transmission per face. The width of this distribution indicates the repeatability of the sawtooth polishing process, on a per-tooth basis.
### 3.4 Radiation Damage in Fluoride Crystals
The crystals of the CLEO RICH are not expected to be exposed to significant levels of radiation. The RICH inner radius is 820 mm, and its outer radius is 1020 mm. At these radii, the CLEO-II detector saw approximately 0.05 Rad/day. This rate is expected to scale very nearly with instantaneous luminosity, so that in CLEO III the rate would be 0.1 Rad/day, or 37 Rad/year.<sup>8</sup><sup>8</sup>8The scaling with luminosity is only valid at the outer radii of the CLEO detector, where the penetrating radiation is shielded only by the inner detector material. The radiation doses to which the inner detectors will be exposed will actually not differ much in CLEO III from those in CLEO-II due to improved shielding near the beamline. In CLEO-c, the numbers are expected to be about a factor of four higher. We have investigated what effect could be expected on the transmission of the RICH crystals due to exposure to synchrotron radiation.
Twenty CaF<sub>2</sub> and 20 LiF samples of the same thicknesses, as used in the RICH, were exposed for various times to a <sup>60</sup>Co source at rates of 200 to 5000 rad/hr. Transmissions of the samples were measured before and after the exposure. The transmissions did not appear to depend significantly on dose rate, only on integrated dose. A control sample that travelled with the dosed samples but was not exposed to the source did not experience transmission changes greater than 1% at any wavelength. Fig. 14 shows the transmissions of the crystals after exposure to the Co<sup>60</sup> source. The observed absorption bands in the LiF at $``$250 nm and 170 nm correspond to known color center formation centers , as do the 190 nm and 250 nm centers observed in the CaF<sub>2</sub> . At the 100–500 Rad exposure expected in CLEO III and CLEO-c, however, the loss in transmission expected in our crystals would not exceed a few percent in the 135–165 nm range.
## 4 MULTI-WIRE CHAMBER CONSTRUCTION and TESTING
The main issues of concern guiding the construction of the multi-wire chamber photon detectors were: (1) field stability, requiring electrodes to be parallel over the full area of the detector, as well as have no local corona points, and high material cleanliness; (2) mechanical stress relief, to avoid cracking the thin CaF<sub>2</sub> windows or breaking an anode wire; (3) gas tightness, to prevent any VUV photon-absorbing gas from leaking into the expansion volume, or any impurities into the chamber gas; and (4) long-term stability, since in all practicality, the installed detector can never be accessed. Fig. 15 gives a schematic view of the multi-wire chamber.
To reduce dirt and dust, a clean room was set up, with air filtering, and the requisite cleanliness protocols (masks, nitrile gloves, booties, etc.). All chamber construction and testing of component parts occurred in this clean room environment. We purchased a set of three ten-foot long granite tables flat to better than 0.0005 inch over the entire table, in order to establish an accurate and stable reference. A mechanical prototype was constructed first, which was crucial in developing the detailed techniques and custom fixtures needed for the construction process.
The chamber construction procedure itself was a complex operation. The general technique was to assemble and test the two halves of each chamber (called the “wire plane” and the “window-plane”) separately, then mate them together and test the completed chamber as a single unit. This procedure allowed careful construction of the three electrode planes, helped in problem diagnosis, and facilitated the production schedule.
For most of the construction procedures, two structural epoxies were used, as well as two sealant epoxies. All were chosen for their strong adhesion properties, viscosity prior to curing, as well as ultra-low outgassing attributes, as indicated in Table 1. All materials were checked for chemical comparability using a special chamber held at elevated temperatures .
### 4.1 Wire Plane Construction and Testing
The wire plane must be constructed in such a way as to maintain the 1 mm gap between the anode wires and the cathode-pads as uniformly as possible over the $``$2.5 m chamber length along which the wires are strung. Gain variation at the wires would result if this spacing was not uniform, as well as variation in the pad-gain, as a consequence of the capacitive coupling to the cathode-pads . In addition, sparking could occur if the gain becomes excessively large at any place in the chamber. These requirement caused us to put a stringent limit of 25 $`\mu `$m on the gap variation.
#### 4.1.1 Cathode Board Design
The cathode board PCBs were designed to have cathode-pads on the front side and connectors for the electronics readout on the back side. Each board was about 205 $`\times `$ 615 mm<sup>2</sup> in area, and nominally 1.7 mm thick. Four boards made up a single chamber. Each contained two 24 $`\times `$ 40 arrays of pads. An 8 $`\times `$ 16 array of pads was routed to two linear connectors in a self-contained manner, so symmetry allows simple tiling of the plane.
As shown in Fig. 16, the pads on the front were connected to the traces on the back with small vias (0.5 mm diameter), which must all be sealed to prevent leaks. Rather than glue each via individually ($``$8000 per chamber), an additional layer of prepreg was added to the top of the lay-up stack, which covered all vias and had cut-outs for connectors and ground connections. This design provided an excellent solution for sealing the vias.
#### 4.1.2 Cathode Board Manufacturing and Flattening Procedure
The large-area cathode board PCBs were manufactured,<sup>9</sup><sup>9</sup>9 Speedy Circuits, Huntington Beach CA 92649. under pressure in a stack in an autoclave. After manufacture, the boards were tested for continuity, inspected for irregularities, and measured to determine mechanical size and deformations. The connectors were then soldered to the board.
For a large-size thin PCB of this type, there were three major deformations from a geometric plane: longitudinal bow, transverse bow, and twist. The magnitude of the longitudinal and transverse bows were about 1.5 mm and 1 mm on average, respectively. The twist was small on this scale. In our design, the transverse bow was the most problematic deformation.
A flattening procedure was developed, in order to remove this bow. The boards were baked under weight while supported along the long edge. The thermal cycle was 2 hr at 150C, followed by a slow 12 hr cool-down. The amount and placement of the $``$1 kg weights depended on the initial deformations of the board. The process was repeated as warranted. The effect of this procedure was to reduce the transverse bow to under half of its original value, on average. The longitudinal bow was not significantly changed.
#### 4.1.3 Cathode Board Assembly
Four cathode boards were assembled into a plane by gluing them together end to end on the granite table—weighted, edge-clamped, and pad face down to assure geometric planarity of the finished pad array.<sup>10</sup><sup>10</sup>10 There was a variation up to 25 $`\mu `$m in the thickness of the etched pads themselves, from board to board. The end-joint between PCBs was specially reinforced: a “vee” cut was milled on the back to allow a larger contact area for epoxy and to ensure that no epoxy came through to make the front (pad) surface irregular. (See Fig. 17.) This joint was covered on the back by a G10 strip in a separate gluing operation for additional strength. The appropriate ground straps between boards were added. Measurements were made of the flatness to monitor the gluing procedure.
Next, the four fiberglass box rib structures were screwed to “strongback” into inserts previously glued into the box ribs. The strongback was made of full-length 1 $`\times `$ 2 inch aluminum box channels. The box ribs were then epoxied longitudinally on the back of the cathode PCBs. Enough glue was applied that any surface non-uniformities in the cathode boards or ribs would be accommodated by the glue. In addition, G10 cross-pieces were epoxied transversely between the ribs. (See Fig. 15.) The main purpose was to provide requisite stiffness to the thin cathode board, when mounted in its final configuration. A second purpose was to remove the residual deformations in the cathode board (the longitudinal bow being removed more effectively than the transverse bow). Excessive twisting, for example, would increase the possibility of breaking wires.
The strongback remained connected throughout the whole construction procedure until the completed chambers were attached to the cylinder superstructure.
A precision ceramic spacer was epoxied each 30 cm along the cathode plane, near the window connections. This holds the wires at a precise distance from the cathode-pads, as well as allow for containment of the failure mode in which there was a broken wire. The ceramic spacer was cleaned before gluing, and care was taken in handling it, so as not to allow any grease or dirt to provide an eventual current path from the wires to the cathode-pads. The wires were eventually epoxied to the ceramic strip which had a slot running down the center to contain the glue bead. The area next to the spacer on the cathode board was passivated by a strip of Kapton, as indicated in Fig. 17. This extended over the edge of the nearest pad, and was done in order to remove the potential problem of any corona points that could lead to high voltage instability over time. Our experience with an early prototype showed that this was necessary.
The anode PCBs, on which the wires were to be soldered, were then epoxied into individually-milled grooves on top surface of cathode board, in order to have precision control of the wire to cathode-pad distance.
#### 4.1.4 Wire Stringing
The chamber wires were strung on a custom-made jig, consisting of a PCB and a precision comb aligned and fixed at each end of a rail structure, approximately 3 m long (longer than a chamber). The 70 central field wires were 20 $`\mu `$m diameter Au-plated W wires<sup>11</sup><sup>11</sup>11 LUMA Type 861-60, W with 3–5% Au and 3% Re, LUMA-METALL AB, 391 27 Kalmar, Sweden. held at 60 g tension. The two outside wires would produce higher fields so we used larger diameter 30 $`\mu `$m Au-plated W wire<sup>12</sup><sup>12</sup>12 Type F-77, 4% Au, Philips Elmet Corp., Lewiston, ME. held at 90 g tension to keep the gain approximately at the same level as the central wires. Each wire was held at the appropriate tension by means of a frictionless pulley with a weight. The wire was carefully positioned in the comb using transverse locator screws. It was then soldered at each end to the PCBs. When done for all wires, a “temporary” plane of wires was created. As much as possible, a single spool was used for the field wires of a given plane.
This temporary plane was tested for wire tension using the standard resonance frequency method. Any wire out of tolerance ($`\pm 2`$ g) was replaced. Fig. 18 shows the distribution of tensions for a representative sample of wires.
#### 4.1.5 Wire Transfer
After the cathode plane was prepared, as described above, the temporary wire plane was flipped over and lowered onto it by means of precision scissor jacks. This was a delicate operation, since the wires had to be aligned to the solder pads on the anode PCBs and just touching the ceramic spacers. When in position, each wire in turn was soldered to two consecutive pads on each end, approximately 5 mm in length, separated by 5 mm of non-pad surface. We chose to use solder with silver added.<sup>13</sup><sup>13</sup>13 Ersin Type Sn62, Eutectic, Tin/Lead/Silver, Multicore Solders, Westbury, NY 11590.
The quality of the solder joint was then assessed, and inspected for sharp points (important since the solder joints will sit in the chamber gas volume). The wire was cut behind the back solder joints, and checked for electrical continuity. Then anode PCBs were given a first cursory cleaning with isopropyl alcohol. The solder joints were covered completely in Delta Bond glue to prevent the solder from being attacked by the TEA in the gas. This process transferred the wires at tension from the temporary jig to the real cathode plane.
The anode PCBs were then populated with their requisite components, as indicated in Fig. 19. This was a mixture of discrete and surface-mount components. The novel feature here is that the large HV surface-mount capacitors (size 1812) were mounted on their side edge, in order to retain the tight-packing demanded by the wire spacing and the constraints on the overall chamber length. A technique was developed which used a special jig and both solder paste and solder cord, and this worked acceptably well for all chambers. Afterwards, the PCBs were brushed with isopropyl alcohol, and the electronic components tested.
The wires were then epoxied to the ceramic spacer. The glue bead was well-contained in a groove atop the ceramic spacer, such that the bead itself was smooth and there was no wicking along the wires. A bead of glue was also made between the two solder joints on the anode PCBs at each end, ensuring no loss of tension over time.
The last stage in the wire transfer procedure was a thorough final cleaning of the anode PCBs. This was accomplished by immersing the entire end of the wire-plane in a large ultrasonic bath filled with isopropyl alcohol, and covering all solder joints. The wire-plane (via its strongback) was mounted on a wall at a 23 angle in order to submerge the anode PCB in the bath. Each end was bathed for 30 mins at a time, and repeated three times with fresh isopropyl alcohol. Afterward, it was rinsed with isopropyl alcohol and dry nitrogen. This was not a completely efficient procedure, since in many cases, there remained a whitish residue around some components afterwards. This is a well-known effect due to solder resin, and simply needs to be cleaned by hand.
#### 4.1.6 High Voltage Testing
After ground connections were made, the wire-plane was placed into a polycarbonate testing box, filled with CH<sub>4</sub> gas and tested for high voltage stability. All wire-planes were tested to 1600 V, with $``$ 10 nA current draw, in conditions of $``$22C and $`<`$15% relative humidity.
In total, there were 35 wire-planes constructed. One was destroyed in subsequent testing, leaving 34 functional planes. They were stored in a low humidity tent until mating with the window-planes.
The wire connections have been very reliable (due to the conservative tension). One wire in the completed detector broke early in the operation due to the chamber exceeding the operating temperature, which went undetected due to insufficient slow control monitoring at the start of the experiment.
### 4.2 Window-plane Construction and Testing
After production of the windows as described in Section 3.2, the full window-plane was constructed. This consisted of a window frame into which the full-length ladder of eight window segments was epoxied.
The individual window segments, wrapped in teflon, were epoxied in pairs. They were butt-jointed with Torr-Seal, with a well-controlled glue bead. The pair was held together in a custom-made jig, under gentle compression, during the curing period. The interior surface of the butt-joint was “painted” using Hysol. When completed, this procedure was repeated for two pairs, and then once again, for a full-length ladder of eight window crystals.<sup>14</sup><sup>14</sup>14Half-plate windows were made into full-plate windows by the same technique.
The full window frame consisted of a long fiberglass side-rail, G10 end pieces, and an Ultem plastic hinge, all epoxied together, as shown in Fig. 15. The G10 end pieces had holes bored through for gas flow to the completed chamber.
The Ultem hinge was specially designed to allow some flexibility when the crystal window was glued to it, in order to take up any mismatch in thermal expansion between the crystals and the stiff frame. The full-length ladder of windows was lowered onto the frame, and glued to the Ultem hinge using Armstrong A-12. At the same time, a small metal foil was wrapped around the edge of each individual crystal window and glued to the interior trace bus-bar by conductive epoxy. On the exterior, this foil was soldered to a teflon-coated high voltage lead wire.<sup>15</sup><sup>15</sup>15 Gore Type F01A080 Wire, W.L. Gore & Associates, Inc., Newark, DE 19711.
After curing, another bead of epoxy (Hysol) was put on top of the existing one, in order to provide a secondary gas seal around the perimeter. Additionally, the butt-joints between window crystals were reinforced by gluing a narrow G10 strip over the joint, externally. This was again reinforced by a glue bead (Torr-Seal) on either side. The Torr-Seal was also used to tack the high-voltage wire along the frame. This operation may be seen in Fig. 20. The end-joints of the windows overlap the position of the ceramic spacers on the cathode plane to minimize the blockage of Cherenkov photons.
The critical mechanical item used in this construction procedure was a set of custom-made (multilayered) jigs, which established the proper referencing surfaces to hold the tolerance in window-to-cathode board spacing. In total, there were 34 window-planes constructed. After testing procedures, 33 usable planes were produced.
### 4.3 Full Chamber Construction and Testing
In the next step a wire-plane was matched with a mechanically suitable window-plane, and they were clamped together temporarily. The cathode-pad to window gap was measured at many points along the length of the chamber. The distribution about the nominal 4.5 mm gap size is shown in Fig. 21 for all chambers using approximately 32 measurements per chamber.
The mated chamber was tested with CH<sub>4</sub>-TEA gas, with nominal high voltage on both wires and windows for a full month. Four special test boxes were used in this testing phase, so that tests could be performed simultaneously. The chamber was scrutinized heavily during this period. Relative chamber gain was measured from the wires, upon excitation by a <sup>106</sup>Ru $`\beta `$-source. Current draws were monitored. Problems were diagnosed and fixed. Typically, if the chamber drew more than 10 nA, it was opened up and cleaned of stray dirt or residual construction materials acting as corona points. These were the most prevalent problems.
If a mated chamber passed this test period, the given window and wire-planes were epoxied together and the chamber was completed. The full chamber was then tested again for a month, under these same conditions.
One problem encountered later, when the chambers were mated with the radiators, was leaking of chamber gas. This was attributed to certain of the fiberglass side-rails, which were epoxy-starved and microscopically split under small torsion applied during handling. This situation was rectified, using a rather painful procedure, that required painting a coat of Hysol epoxy over the side-rail, and reinserting the chamber before the glue dried; this often had to be repeated.
A total of 33 chambers passed all tests. On the average, this phase of construction took $``$0.5 months per chamber. Parallel operations allowed a maximum of eight chambers to be constructed in a single month.
The last stage in individual chamber construction was mounting the electronics chip-carriers and cables on the back. The hybrids (described in Section 7) were screwed into small standoffs epoxied to the connectors, thus attaching them with a squeezing action instead of pressing action which could deflect the cathode board and possibly break wires.
## 5 RADIATOR CYLINDER CONSTRUCTION
### 5.1 Inner Carbon Fiber Cylinder
The RICH Inner Cylinder has $`12\mathrm{m}^2`$ of LiF crystals attached to a cylindrical $`1.64`$m outer diameter carbon fiber support shell of $`2.48`$m length and $`1.5`$mm skin thickness. The shell has mean density $`1.42\mathrm{g}/\mathrm{cm}^3`$, mean Young’s modulus $`E70`$GPa, and is built from wrapping multiple layers of pre-impregnated unidirectional tape<sup>16</sup><sup>16</sup>16 Type RS-3/AS-4 unidirectional tape, YLA, Inc., Benicia, CA 94510. around a drum-shaped steel mandrel. After serving as the form for the carbon fiber tape during autoclaving, the mandrel provides mechanical support for the shell during transportation, radiator assembly, insertion into the Outer Cylinder of photon detectors, and installation of the entire RICH detector into CLEO.
The mandrel rotates on a shaft that allows the shell to be rotated to an arbitrary azimuthal position. A large box-beam frame supports the shaft so that the total mass of the radiator system, including all its auxiliary mechanical fixturing, is $`3500`$kg.
### 5.2 Radiator Crystal Alignment and Mounting
The radiator cross-section is essentially a 30-sided polygon (“triacontagon”) with each side corresponding to a longitudinal row of 14 radiators. (Each row contains 13 full-sized crystals and two half-sized crystals at one end.) Construction was accomplished by first assembling and aligning rows of crystals and then attaching the rows sequentially to the support shell. All radiator assembly procedures are performed inside a class 100,000 clean room with typical 25% relative humidity to reduce the risk of contaminating crystal surfaces with particulate matter or excessive water vapor. Test crystals placed in the clean room were periodically measured for VUV transmission to verify that the ambient environment did no damage to radiator crystals.
To construct a radiator row, a set of 13+2 crystals were selected on the basis of VUV transmission and compatible mechanical dimensions. Each crystal was attached temporarily, by a $`50\mu `$m thick Kapton belt, to a picture frame jig that rests directly on top of it and that has vertical posts in each corner to provide for temporary attachment to a transport jig. Each crystal with its attached jig was then placed on an optical stage assembly with 6 degrees of freedom (3 orthogonal linear displacements and 3 Euler angle rotations). The 15 stage assemblies were all attached to a single rigid optical rail.
The open picture frame allows a pair of $`2.2`$cm diameter dowel pins to be placed directly on the crystal top surface. These pins support a precision level so that each surface is made horizontal to a typical angular precision of $`150\mu `$rad. A similar technique is used to limit vertical offsets between top surfaces of adjacent crystals to a precision of $`50\mu `$m.
A $`100\mu `$m clearance gap between adjacent crystals was reliably set by the thickness of the Kapton belts from adjacent crystal jigs. Pushing the same edge of each crystal flush against a set of tooling balls running parallel to the optical rail ensures that one edge of a crystal row was straight within $`75\mu `$m over its entire length.
After a crystal row has been aligned, a second “transport” optical rail is positioned over the crystal jigs. Quick-setting epoxy temporarily fixes the transport rail to the vertical posts of each crystal jig. The transport rail, along with the attached row of crystals, was then moved to the carbon fiber cylinder.
Linear translation optical stages positioned at opposite ends of the support shell near its surface at the “twelve o’clock” position receive the transport rail, as shown in Fig. 22. Reamed holes in these stages mate with dowel pins in the transport rail so that the axis of the new crystal row is aligned with the axis of the support shell. These stages eventually set the epoxy gap between the crystal bottom surfaces and the cylinder surface to $`25\mu `$m precision. Gaps between adjacent crystal rows on the cylinder were controlled by the azimuthal rotation of the mandrel, set by a micrometer actuator to a precision of $`50\mu `$m.
After a crystal row has been appropriately positioned about the shell surface, Armstrong A-12 was applied to the shell just below the center of each crystal. No special preparation of the shell surface was performed. The crystal row was then lowered by the stages to make a nominal $`150\mu `$m gap between shell and crystal, forcing the epoxy to make a contact patch with nominal surface area of $`100\mathrm{cm}^2`$. The jigging was left in place for a minimum of 8 hours before removal. The shell was then azimuthally rotated $`12^{}`$ to a new position and the entire alignment and epoxy process repeated until all 30 crystal rows were attached.
### 5.3 Radiator Transport
Care was taken to prevent the LiF crystals from experiencing excessive vibration or humidity during radiator transportation from its construction site (Dallas, TX) to the detector integration site (Syracuse, NY). A dedicated temperature-regulated shipping truck was used for the 2500 km trip. A local atmosphere close to 0% relative humidity (measured as $`<`$1%) was achieved by erecting a temporary plastic film cocoon around the radiator and flowing dry N<sub>2</sub> gas through the cocoon from gas bottles attached to the shipping frame.
Air bags attached to the four corners of the shipping/rotation fixture cushioned the radiator against mechanical shock and vibration. Acceleration of the radiator and temperature inside the truck as functions of time were measured, digitized and written to local non-volatile memory by special instrumentation<sup>17</sup><sup>17</sup>17 Model EDR-3C, Instrumented Sensor Technology, Inc., Okemos, MI 48864. installed on the mandrel support frame. No damage to the radiator was observed at the trip’s conclusion.
## 6 SUPERSTRUCTURE CONSTRUCTION
The guiding principles for the design and construction of the full cylinder mechanical superstructure were: (1) low mass, so as not to degrade the excellent performance of the electromagnetic calorimeter surrounding the RICH; (2) no internal obstructions, so as not to shadow the photon detectors; and (3) an excellent gas seal.
The overall design of the superstructure consisted of an “Outer Cylinder” for the multi-wire chambers, and an “Inner Cylinder” for the crystal radiators. The expansion gap separated these and was contained by two large end flanges .
### 6.1 End Flanges
For construction purposes, the two annular end flanges were each made of two concentric parts, as shown in Fig. 23. The outer half was a structural element for the Outer Cylinder, and the inner half was fixed to the carbon fiber shell of the Inner Cylinder. These two parts fit together with an double gas seal, as described below. The end flanges were made of structural aluminum alloy, and coated on the interior surface by low-outgassing black paint.
### 6.2 Outer Cylinder
The construction of the Outer Cylinder used a scaffold, on which the multi-wire chambers were mounted, as can be seen in Fig 24. It had long thin aluminum rails supported between the outer halves of the end flanges. The rails were machined to fit between the chambers, and contained an O-ring groove that was the primary gas seal for the expansion gap at the Outer Cylinder. As it was not very stiff but had to maintain the requisite compression on the O-ring, a clamp-strip was used to squeeze the chamber on the O-ring, with screws spaced on 1 inch centers. Hence the chambers themselves effectively formed the outer radius of the detector. This is another reason why the box rib reinforcements were needed.
During assembly, the Outer Cylinder was held on a large custom-designed A-Frame structure with three-point kinematic mounts, allowing free rotation of the cylinder. Initially, the inner area of the end flange was plugged, and the thirty chamber spots were taken up by “dummy panels,” stiff blanks that allowed for gas tightness testing of the cylinder. In this condition, the out-of-roundness deformation was measured, as shown in Fig. 25, to have a small eccentricity.
As chambers were completed, the dummy panels were removed and the chambers were mounted on the cylinder. A photograph of the interior during this procedure is given in Fig. 26. It was only at this time that the strong-backs could be removed. Leak-checking could then be done using He gas for high sensitivity. This was particularly important at the corners of the long single O-ring seal around the perimeter of the chamber. Often, the joints and corners had to be hand-worked to seal at the level of 10<sup>-4</sup> ml/s.
### 6.3 Cylinder Mating and Gas-Sealing
The Inner and Outer Cylinders were mated in a precision mechanical process. The inner half of one end flange was first epoxied to the carbon fiber shell while the other inner flange was free. The Outer Cylinder was rotated azimuthally to align the fixed inner-half flange with its corresponding outer half and O-rings were installed in both outer flanges. The Inner Cylinder remained on the mandrel and slipped into the Outer Cylinder, riding on two large box-beams. The halves of the end flanges were screwed together, providing the mechanical coupling between the Inner and Outer Cylinders. An O-ring provided the primary gas seal at this junction. Following our dictum to use redundant gas seals on all mated surfaces, a secondary gas seal was made by a glue bead around the joint. The previously free inner half of one end flange was epoxied to the carbon fiber shell as the mechanical connection was made.
Thereafter, the completed structure could be supported by the mandrel. It was then transported for installation in CLEO by a method similar to that used for the radiator cylinder. Chambers and gas seals were then all tested again, in situ. After additional iterations of hand-working and testing, the expansion volume was sealed. Most leaks occurred when the inner and outer cylinder were mated.
## 7 READOUT ELECTRONICS
### 7.1 System Description
The CLEO RICH electronics design is driven by two important considerations. First of all, we need to operate the chambers at moderate gain, to improve the stability of operation and the lifetime of the detector system. This requirement is very important as the system is designed to operate throughout the lifetime of CLEO without any access for repair. In this regime the single photon response of the MWC used as photon detectors has an exponential distribution, as shown previously in Fig. 5. Thus, low noise is a critical requirement to maintain good efficiency. On the other hand, the exponential distribution spans a wide dynamic range. Moreover we would like to be able to reconstruct the charge deposited by a minimum ionizing particle, shown as the bump on the right-hand side in the distribution of Fig. 5. This implies that a wide dynamic range is also very important.
The cathode-pad segmentation is such that the charge signal induced by the avalanche is spread around more than one pad, thus analog charge weighting produces a better spatial resolution and also allows for easier separation of the charge clusters produced by two nearby photons. Thus we chose to implement an analog readout.
Fig. 27 shows a schematic view of the readout architecture for a multi-wire chamber. Each chamber is divided into four sectors. Each sector contains three daisy-chained rows, connected with 50-conductor shielded cable to the back-end electronics, located in VME crates about 18 meters from the detector cylinder.
The front-end hybrids, shown in Fig. 28, are mounted on the back of the cathode board in a mother board-daughter board configuration. Five hybrids are daisy chained and share the same bias sources. Each hybrid has an independent differential output for faster data processing.
A photograph of a quarter-section of the cathode board with the hybrids attached is shown in Fig. 29.
### 7.2 The Front-End ASIC
The heart of the front end electronics is the VA\_RICH ASIC. It was designed for our application by the engineering team at IDE AS, Norway , and is an adaptation of the basic design of the very successful VA family, originally developed to process the signals of silicon microstrip detectors, tailored to our noise and dynamic range requirement.
Fig. 30 shows a conceptual diagram of this device. It features 64 individual channels including a semi-Gaussian preamplifier and shaper circuit. The peaking time can be adjusted by changing the biases of the shaper circuit around a typical value of 2 $`\mu `$s. We tuned the peaking time to match the time when a Level 1 decision is achieved. This section is followed by a sample and hold circuit, which is designed to hold the peak level out of the shaper until the output multiplexer is ready to transfer this level as a differential current. The individual inputs are connected through a 64:1 multiplexer to a calibration circuit, which allows injection of a test charge into each individual channel. The ASIC is implemented in 1.2 $`\mu `$m AMS CMOS technology.
The equivalent noise charge dependence upon the input capacitance was measured on prototype single chip hybrid carriers, using a set of calibrated capacitors. The measured performance matches the predictions from the ASIC simulation . In our application, we expect a total equivalent noise charge of 300 e<sup>-</sup>, without the cathode boards connected. Fig. 31 shows the analog output voltage on 500 $`\mathrm{\Omega }`$ resistors. Saturation occurs at an input charge level of 80 fC. In the linear region, the preamplifier and shaper gain is 40 mV/fC, for a load of 500 $`\mathrm{\Omega }`$.
### 7.3 Bench-Test Characterization
The total number of hybrids produced was 2200, over a period of two years. The hybrids produced were tested for functionality at IDEAS and then shipped to Syracuse for more complete characterization. Our tests involved a noise measurement by taking pedestal data, followed by a detailed mapping of the individual channel gain with calibration pulses of different amplitudes. Only 1800 were installed in the detector.
The noise measurement required great care in the grounding and shielding of the hybrids and careful routing of the analog power supplies. Our goal was to achieve a noise of the order of 400–600 e<sup>-</sup> with a simple and relatively quick set-up of the measurement.
After this initial characterization, we performed a burn-in test of the hybrids, maintaining them biased in their nominal working point at elevated temperature (70C) and performing electronic calibration cycles at regular intervals. In order to perform these tests at a rate compatible with our installation schedule, we produced a dedicated set-up where 32 hybrids could be biased and monitored in parallel. After one week in the burn-in set-up, the hybrids were tested for noise and gain with a quicker calibration procedure. The hybrids that were rejected upon this procedure were only a few per thousand.
### 7.4 CLEO RICH Data-Boards
Fig. 32 shows a picture of the CLEO RICH data boards. They are 9U mixed analog and digital environment boards that perform several very complex functions. Note that the board is physically composed of two different sections. The first is an analog section, providing the biases needed for the functionality of the VA\_RICH ASIC, the transimpedance receivers, 12-bit ADCs digitizing the serial analog information, and slow control monitoring ADCs. The second is a pure digital section, based on a common CLEO data acquisition framework, and containing two components specific to our application: a sequencer, based on the ALTERA MAX FPGA, that contains the firmware necessary to operate the VA\_RICH ASIC, and an Analog Devices ADSP-21061 DSP, used to perform the common mode suppression and data sparsification described below.
The analog section features 15 receiver channels, organized into three cells sharing the same ribbon cable interconnection to the front end hybrids. Each cell encompasses dedicated $`\pm `$2V analog power and several DACs that allow the adjustment of the bias currents and voltages that influence the working point of the analog front end, as well as a slow control section that monitors the values of these voltages and currents, and the temperature on different locations of the detector cylinder and on the data boards themselves. In addition, local pulse generators are used for the electronics calibration. All the cells share a common regulated $`\pm `$2V digital section.
When an event occurs, all the 15 channels in a board receive the synchronous serial information from the corresponding hybrid. In order to simplify the digital section the 5 ADC channels in a cell are multiplexed into a common FIFO buffer, with an offset of 23.4 ns, determined by the local 42 MHz clock.
To maintain the low thresholds that are required to optimize the efficiency, the hardware common mode noise suppression algorithm is crucial to suppress the adverse effects of coherent fluctuations of all the channels in an ASIC.
## 8 SUPPORT SYSTEMS
### 8.1 Gas System
The gas system supplies several distinct volumes. The system must: supply CH<sub>4</sub>-TEA to 30 separate chambers, supply ultra-pure N<sub>2</sub> to the expansion volume, supply ultra-pure N<sub>2</sub> to a sealed single volume surrounding all the chambers, called the “electronics volume,” since this is the region where the front-end hybrid boards are present. In addition we need to test the CH<sub>4</sub>-TEA mixture for the ability to detect photons and test the output N<sub>2</sub> for purity.
It is of primary importance that the gas system must not destroy any of the thin CaF<sub>2</sub> windows. We use computerized pressure and flow sensors with programmable logic controllers (PLC).
#### 8.1.1 System Design
The performance and mechanical integrity of the CLEO III RICH are critically dependent on the performance of its gas systems. In order to achieve its design resolution, the RICH must efficiently detect 14-21 photons emitted by charged tracks traversing the radiator. The efficiency of the RICH depends on the expansion volume transparency. The photosensitive detector gas and transparent nitrogen of the expansion volume are separated by fragile windows on the inner faces of the detector modules. These windows could be destroyed by a slight pressure difference between these volumes (of order 15” H<sub>2</sub>O). Such an overpressure would be catastrophic to the RICH.
The gas systems were carefully designed to protect against such damage. A highly automated design was chosen in order to minimize the possibility of operator error and to provide fast response to dangerous conditions. Most valves are controlled by the automated control system and most sensors are read out and monitored electronically.
The two major subsystems are the expansion volume gas system (EVGS), which purges the expansion volume with a large flow of ultra-pure nitrogen, and the detector gas system (DGS), which supplies the thirty detector modules with photo-sensitive methane/TEA gas. Pressures in the detector and expansion volumes are referenced to atmospheric pressure. Both the EVGS and DGS employ triply-redundant overpressure and underpressure protection. The first tier of protection is based on readings from pressure transmitters connected to the expansion volume and each of the thirty detector volumes, and has trip points at +0.75” H<sub>2</sub>O and -0.5” H<sub>2</sub>O. The second is based on readings from pressure switches which have trip points set at +1.5” H<sub>2</sub>O and -1.0” H<sub>2</sub>O. The final level of protection comes from mechanical relief valves with set points at +2.5” H<sub>2</sub>O and -1.5” H<sub>2</sub>O. Valves which perform the most critical operations, such as those which shut off the gas inputs, are equipped with position indicating sensors and, in the most critical, redundant valves are implemented.
#### 8.1.2 Control System
All critical operations of the RICH gas systems are controlled by a programmable logic controller (PLC). The PLC is an industrial control system designed to be reliable and modular. The control system utilizes 112 analog input, 8 analog output, 96 discrete output, and 128 discrete input channels. The gas system is controlled through a small number of buttons on the main control panel.
The basic operation of the EVGS and DGS is similar. Each has four states of operation: RUN, IDLE, STOP, and ALARM. RUN mode is the normal operating state of the system. In this state gas flows through the detector and the pressure is actively regulated by the PLC at 0.5” H<sub>2</sub>O to an accuracy of $`\pm 0.02`$” H<sub>2</sub>O. IDLE mode is similar to RUN mode except that the pressure is not actively regulated. IDLE mode is used in the system start up sequence. ALARM mode is the safe state to which the system defaults in case of an emergency. In this state, the input gas flow is shut off and the exhaust side of the detector is opened to track atmospheric pressure. After the system is stopped by an alarm it is important that it be restored to RUN or IDLE mode as soon as possible to prevent degradation of the CaF<sub>2</sub> and LiF crystal transparency when exposed to ambient humidity. After the offending condition is fixed, alarms may be reset using a button on the control panel. When the alarm is cleared, the system changes to STOP mode, which has the same physical configuration as ALARM mode. The system start up sequence is fully automated and is initiated by pressing a single button. The start up sequence consists of a purge of the input lines, followed by a ramp of the input flow, and finally a ramp of the pressure. All transitions in the system are gradual in order to protect the RICH from possible pressure spikes.
All critical components of the gas system are powered through an uninterruptable power supply which is backed up by a diesel-powered generator.
Interfaces to the operator and to the CLEO III detector control system are provided through the gas system monitor (GMON), consisting of LabView<sup>TM</sup> programs running on a PC. The GMON programs collect sensor readings, status, and alarm information from the PLC which it displays in several “active schematics” corresponding to each of the major components of the EVGS and DGS. The GMON relays this information to the CLEO III slow control system.
The GMON also handles the display of alarms. “Fatal” alarms are generated by the PLC and result in the system defaulting to ALARM mode in order to protect the system from damage. All other “non-fatal” alarms are generated by the GMON code, which monitors the status of all system parameters. The system is heavily instrumented with sensors in order to allow problems to be easily identified and repaired.
The GMON programs are also used to control some non-critical operations of the gas system, such as the gas quality monitoring, and to set some system parameters. The gas system does not depend on the PC or GMON programs in order to run, however.
#### 8.1.3 Expansion Volume Gas System (EVGS)
It is critical that the purity of the expansion volume be maintained in order for the RICH to achieve its design resolution. Most impurities will absorb the UV photons before they reach the detectors and must be kept below a few ppm concentration. This purity is obtained by flowing nitrogen at a rate of approximately 1500 $`\mathrm{}`$/h in a single-pass configuration. Boil-off from liquid nitrogen is used as the source for the expansion volume gas. Most of this system is constructed from electropolished stainless steel and other materials with low vapor pressure in order to minimize impurities.
Nitrogen gas from the dewar is processed by an automated purifier subsystem. This purifier injects a regulated flow of hydrogen corresponding to 40 ppm, which reacts with the oxygen inside a catalyst cartridge to form water. A slight excess of hydrogen is acceptable since it is known to be transparent over the wavelength range of interest. Water and other large molecules are removed by a large molecular sieve trap. The molecular sieve trap must be regenerated occasionally by baking at 300 C for a day. The purifier subsystem contains two parallel purifiers to allow for continuous operation while one of the sieves is being baked. The PLC controls the flow ramping and switching between purifiers as well as the temperature control and gas purge for the baking sieve.
The operating flow of nitrogen to the expansion volume is set with a manual metering valve near the input to the expansion volume. Gas exiting the expansion volume passes through ten 1/2” stainless steel lines to a 1” manifold outside of the CLEO III solenoid return yoke, where pressure is regulated by a PLC-controlled metering valve based on feedback from a pressure sensor.
#### 8.1.4 Detector Gas System (DGS)
The detector gas system is responsible for providing thirty separate photon detector volumes with a mixture of methane and TEA. Because TEA is corrosive to most materials, 316 stainless steel components were used wherever possible. Any other materials used in the system were first subjected to stringent chemical compatibility tests.
TEA vapor is introduced into the stream of methane gas by bubbling it through liquid TEA at 15 C. The TEA bubbler system is automated using programs on the PLC for reliable and continuous long-term operation. It utilizes a temperature-controlled bubbler chamber and elevated TEA reservoir with enough TEA for approximately two months of operation. The level of TEA in the bubbler chamber is regulated by the PLC which fills from the reservoir as needed. The reservoir may be depressurized and filled without interrupting the operation of the system. Sensors monitor the level and temperature of the TEA in the bubbler and the level of the TEA in the reservoir and notify the user when refilling is required.
In order to provide adequate pressure control, it is necessary to supply gas to the chambers in a parallel configuration rather than in series. The flow is split into thirty separate streams to provide even flow to all 30 detectors. The flow in each branch is set manually and monitored electronically. The flow is ramped up and down automatically during start up and purge operations. A total of ninety 1/4” stainless steel lines pass through the CLEO III solenoid return yoke steel to connect to the thirty photon detector modules: one input, one output, and one pressure sensing line.
Exhaust from the thirty modules is collected into a single exhaust manifold. Gas from this manifold passes through a PLC-controlled pressure regulating valve which maintains the average pressure of the thirty modules at 0.5” H<sub>2</sub>O.
#### 8.1.5 Gas Quality Monitoring
A nitrogen transparency monitoring system measures the UV transmission of the expansion volume gas as a function of photon wavelength. Monochromatic light is produced by passing light from a deuterium lamp through a computer-controlled monochromator. The monochromatic photon beam passes through a gas sample tube and is detected by a photomultiplier. The sample tube can be switched to sample the expansion volume exhaust gas or an ultra-pure argon reference gas. A comparison of these two gives the transparency. The GMON controls the transparency monitor and other gas monitoring devices. Automated transparency scans are taken once per day. Graphical panels allow the user to configure the scans and view the results.
Oxygen is the most likely and problematic contaminant, since it absorbs strongly at 150 nm. The oxygen concentration of the output gas from the expansion volume is also monitored using a precision oxygen sensor.
The CH<sub>4</sub>-TEA mixture is monitored by an electron capture detector. This device consists of a cylindrical proportional drift chamber which utilizes a beta source to produce ionization. Electrons produced near its outer diameter drift through a $``$5 cm long path to an anode wire at the center of the chamber where they are collected. A 0-5 VDC signal proportional to the time-integrated drift current is monitored by the control system. Impurities or changes in the gas composition can be detected as a change in the output of this device.
#### 8.1.6 Electronics Volume Purge System
The gas system also provides a purge flow of nitrogen gas around the electronics volume of the RICH in order to minimize the concentration of impurities adjacent to the gas volumes of the RICH. This helps to reduce the infiltration of oxygen into the expansion volume.
#### 8.1.7 Performance
The gas system has been in operation since the early commissioning of the RICH and has been running in its final form since the Summer of 1999. The most critical aspect of its performance is that it has protected the extremely fragile CaF<sub>2</sub> windows from damage. It has delivered nearly continuous service for several years, with only a few short down-times for repairs. Identification and diagnosis of problems has been straightforward due to the wealth of information provided by the system.
The expansion volume transparency exceeds the specifications for the RICH design. The EVGS typically delivers transparencies of $`>`$99.5% at 150 nm, as shown in Fig. 33. This wavelength is near the peak of the RICH sensitivity and is particularly vulnerable to oxygen contamination. The recovery time of the expansion volume gas system is quite short due to the high purge rate and careful choice of materials. After exposure to air, the expansion volume can typically be restored to acceptable transparency within a half day.
The RICH photon detectors show no evidence of aging due to reactions of TEA with materials in the gas system and RICH. The DGS has provided reliable and simple operation. In particular, maintaining and filling the liquid TEA bubbler has proved to be quite simple and trouble-free.
The highly automated design of the RICH gas system has proved to be very effective at protecting the RICH and simple to operate. It has exceeded expectations of performance and reliability.
### 8.2 High Voltage System
The High Voltage System for the CLEO RICH was required to supply 90 channels of anode field wire voltage (about $`+`$1500 V) and 240 channels of cathode window trace voltage (about $``$1200 V), with a current monitoring and trip circuit on each channel. The system was based on the LeCroy ViSyN System , and consists of a Model 1458 mainframe containing slots for 16 Model 1469 modules, which have either positive or negative polarity. In our implementation, we use 5 positive and 11 negative modules. The module outputs are ganged together such that there are eight HV channels per so-called “bulk” connector, a specially-designed LHV connector. There are three bulk connectors per module. The system allows voltage setting at the bulk connector level, so each channel in the bulk connector has the same voltage, but may have different current trip levels. The channels fan out to a patch-panel so the proper high voltage can be routed to each chamber channel. The value of the voltage setting on each of the 330 detector channels was determined by an optimization procedure, as discussed in Section 9.2.
Control of the high voltage system is accomplished by a stand-alone server program, on a dedicated linux station, that is integrated with the CLEO DAQ System in order to synchronize the state of the system with accelerator operation.
Ramping of voltages are made at rates of 40 V/s (increasing) and 50 V/s (decreasing) for both polarities. Care is taken to raise all voltages in a given chamber simultaneously, in order to minimize electro-mechanical deflections. All channels have current trip levels set to 10 $`\mu `$A, and are monitored at a rate of $``$2 kHz. Individual channel status, currents and voltages are read back and monitored as well. Any currents or voltages out of tolerance will cause an alarm, requiring immediate operator intervention. During detector operation, communication between the server and the mainframe is maintained via TCP/IP, and the status of all voltages and currents is displayed in a representational color-coded GUI.
### 8.3 Cooling System
The total power output of the VA\_RICH chips is approximately 360 W. In order to provide heat removal necessary for the mechanical stability of the chambers, four nylon cooling tubes run over the back of each chamber, and are each modestly thermally coupled to the VA\_RICH hybrids by two-component RTV.<sup>18</sup><sup>18</sup>18 RTV 577 / RTV 9811, GE Silicones, Wilton, CT 06897. Supplying these tubes with the hydrocarbon coolant PF-200IG is a specially-designed CLEO cooling farm . This coolant is used rather than water because of concern that a leak could reach the water sensitive CsI crystal, which would melt; an additional advantage is that PF-200IG is non-conductive. The RICH cooling circuit is one of several driven by the cooling farm through active-manifold platforms, with a coolant reservoir temperature of 18C and flow rate of $``$$`\mathrm{}`$/min. Each RICH chamber line is equipped with an embedded tip-sensitive $`100\mathrm{\Omega }`$ platinum resistance temperature sensor (RTD). For these RTDs, an array of “hockey-puck” style 4–20 mA transmitters route signals to a custom-built multiplexer circuit used to read out these transmitters every 10 s. The temperature signals are fed into a small logic controller for real-time control and monitoring. The typical range of exit temperatures from the 30 chambers when in operation is 21–23C, each held stable to $`\pm `$0.5C.
The status of the cooling system is continuously monitored by a LabVIEW program which provides web-based diagnostic tools for real-time viewing of the performance parameters on site and monitoring worldwide.
### 8.4 Slow Control System
A bit-serial data bus was implemented on the P3 connector on the RICH AVME data board using differential signaling technique. This bus is referred to as the SBUS. The maximum transfer rate of the SBUS is approximately 1 Mbit per second. Each RICH VME crate is equipped with one SBUS crate controller. Each SBUS crate controller can be addressed uniquely. All eight RICH VME crates are daisy chained to one SBUS controller module located in a separate SBUS crate at the pit level.
The SBUS system is responsible for monitoring the operational status of the VA\_RICH chips as well as that of the AVME data board. Important quantities such as the reference voltages, the reference currents, and the temperature on the chip carrier chain inside the RICH detector are being monitored periodically through the SBUS system. Analog signals from the RICH detector are converted into digital signals with a 10-bit ADC (TI TLC1542/3CDW) on the AVME data board for the SBUS read out. The SBUS can also be used to download some of the EEPROMs on the DVME data board.
A server program running on an on-line computer communicates with the SBUS controller module via CORBA calls. A monitor request is sent out every 15 minutes. This sampling interval is programmable and was chosen to ensure both adequate monitoring and maximum system wide stability. The data are shipped back to the server, archived locally in plain text files and then processed to provide information on the front-end electronics status in an easy-to-read HTML format on a Web server.
## 9 OPERATING EXPERIENCE
The CLEO RICH detector has been in operation since September 1999. All but $``$2% of the detector is functioning for data-taking. We lost 1.7% due to the breaking of one wire after about one year of operation due, most likely, to the un-monitored heating of chamber due to problems with slow control software. We have also lost 2% of the electronics chips and suffered one broken output cable, so the total number of lost channels is $``$5%.
### 9.1 Normal Operation
In order to ensure proper functioning of the RICH over the long period of its operation, a series of calibration and monitoring activities are performed periodically as a part of its normal operating procedure.
To ensure the highest detection efficiency of the VA\_RICH chips, the electronics pedestals need to be measured routinely. This is done by performing a special calibration procedure called “SmallCal.” During a SmallCal, the VA\_RICH chips are set in calibration mode and the electronics noise from all 230,400 electronic channels are read out. The pedestals are then calculated on-line and loaded to the VME crates to be used for sparsification during successive data-taking runs. This SmallCal procedure is carried out regularly every eight hours.
Another type of calibration, the WirePulse run, is performed once per week. The purpose of the WirePulse run is to measure the response of the RICH electronics to a given input signal. During a WirePulse run, the anode wires are pulsed with a sequence of several predefined waveforms of differing amplitudes. The output from all electronic channels are analyzed off-line (see below).
As detailed above in Sec. 8, the high voltage server program monitors the RICH HV status, by reading voltage and current values, and watching for trips. Also as detailed above, the RICH gas system and the cooling system adequately control critical operating parameters, but also provide for web-based monitors which are checked twice per hour during normal operation. For the gas system, the expansion volume transparency, the TEA bubbler temperature and the expansion volume oxygen content are closely monitored. For the cooling system, the chamber temperatures are watched. The operating parameters of the readout electronics is checked using the web-based utilities of the Slow Control system.
In addition, RICH performance is supervised offline, with a set of quality-monitoring plots produced from a prescaled sample of incoming events, from the so-called CLEO pass1 analysis. These are examined after each run, and include distributions of Cherenkov angles and photon yields from fast tracking, as well as the more pedestrian distributions of raw hits and cluster pulse heights (as described in Sec. 10).
This information has been sufficient to ensure that the CLEO RICH detector produced data with a high degree of stability over many years.
### 9.2 Chamber Gain Equalization
The high voltage operating point for all chambers was determined by a gain equalization procedure, which sought to make equal the pad-gain for single photons averaged over each window module individually. The gains were set to be below 25,000 in order to avoid discharges.
Gain changes were measured as a function of window and wire voltages. The gains were determined for each window-sized module by fitting an exponential curve to the pulse-height spectrum. We parameterized the gain as a function of the wire voltage, for each window voltage.
Gains of all of the chambers were varied by changing the voltages in an iterative manner in order to make them equal. Our goal was to have a pad-gain of $``$23,000. Fig. 34 shows the distribution of gains as measured after the gains had been equalized. We find a mean of 23,400 with a fitted r.m.s. spread of 10%. Fig. 34 shows the distribution of gains as measured after the gains had been equalized. We find a mean of 23,400 with a fitted r.m.s. spread of 10%. This has remained stable during detector operation.
### 9.3 Electronics Performance in CLEO III and CLEO-c
During data taking we routinely perform two sets of measurements to check the electronics performance. We measure the pedestal periodically, to verify that no baseline shifts occurred. These data also allow us to monitor the value of the total and incoherent noise of each channel. Typical noise distributions are shown in Fig. 35. The noise reduction produced by the coherent noise suppression can be clearly seen. The noise levels are quite low, the peak of the incoherent noise is at 425 electrons.
Fig. 36 shows a scatter plot of the intrinsic noise distribution across the detector. The horizontal axis corresponds to the length of the chambers and the vertical axis corresponds to the chamber number. It can be seen that the profile is rather uniform. The darker areas correspond to break points between different chamber sectors, and are associated with additional noise sources at the boundary between two adjacent sectors. Cross-talk or additional digital noise in the chip carrier at the end of a chain may cause this higher noise level.
The overall noise performance has been extremely stable throughout the years. Fig. 37 shows the time evolution of the mean value of the total noise and its very stable value.
An additional quantity that needs to be monitored is the gain of the front end electronics. In principle, we could undertake the same sort of electronics calibration that was performed to verify that the hybrids were compliant with all our specifications. However, the amount of data to be transferred is beyond the capabilities of the data acquisition system. Thus, we use a method that is much quicker and simpler, although not as accurate. We pulse the MWC wires and the capacitive coupling between wire and pad produces the current pulse amplified by the front end electronics. The advantage of this calibration procedure is that the current originates at the same location as the real signal. Thus it tests the integrity of the whole processing chain, including the wire bond between input channel and hybrid trace and all the connectors along the signal path. On the other hand, the pulse shape is not perfect because of the improper termination, thus the gain measurement is not very accurate, as illustrated in Fig. 38, that shows the signal distribution on all the pads for input pulses of +5 mV and $`5`$ mV respectively. The wire pulse distribution is also used to determine the number of dead channels. We have about 5% dead channels, and the number has been relatively stable throughout the duration of the experiment. These losses include one damaged multi-wire ribbon cable and ASIC damage as well as lost wirebonds.
The most notable source of ASIC loss has been occasional latch-up of these devices, due to the activation of parasitic PNPN paths in the device. The symptoms include large current draw and a different pedestal profile along the ASIC, much “flatter” than in a functioning ASIC. This phenomenon occurred more frequently in the early stages of the experiment, due to occasional faulty start-up procedures, that caused some device destruction. After we refined the start-up procedure, to prevent regenerative loops from occurring, the number of flat ASICs has remained stable, with some signs of recent recovery.
In general, we can say that the operation of this system has been very stable and reliable and no tuning of the bias voltages and currents has been necessary, since the initial adjustment performed upon installation.
## 10 DATA ANALYSIS AND PHYSICS PERFORMANCE
### 10.1 Introduction
The CLEO III detector was used for studies at the $`\mathrm{{\rm Y}}(1S)\mathrm{{\rm Y}}(5S)`$ resonances from August of 2000 to March of 2003. The CLEO detector was then modified by replacing the silicon strip vertex detector with a low mass wire chamber. The magnetic field was also lowered from 1.5 T to 1.0 T, to help increase the machine luminosity. Data was then taken from October of 2003 until April of 2005. The results in this section refer to the first period as CLEO III and the second period as CLEO-c. More CLEO-c data will be forthcoming.
Coherent noise suppression and data sparsification are performed on-line to eliminate the Gaussian part of the electric noise. A small non-Gaussian component of the coherent electric noise is eliminated off-line, by using an algorithm too complicated for use in the data board DSP. The incoherent part of non-Gaussian noise was eliminated by off-line pulse height thresholds adjusted to keep occupancy of each channel below 1%. Finally, we eliminate clusters of cathode pad hits that are extended along the anode wires, but are only 1–2 pads wide in the other direction.
We show in Fig. 39 the hit pattern in the detector for a Bhabha scattering event ($`e^+e^{}e^+e^{}`$) for track entering the plane (left image) and sawtooth (right image) radiators. The shapes of the Cherenkov “rings” are different in the two cases, resulting from refraction when leaving the LiF radiators. The hits in the centers of the images are produced by the electron passing the RICH multi-wire chambers.
### 10.2 Clustering of Hits
The entire detector contains 230,400 cathode pads, which are segmented into 240 modules of $`24\times 40`$ pads separated by the mounting rails and anode wire spacers. We cluster pad hits in each module separately. Pad hits touching each other either by a side or a corner form a “connected region.” Each charged track reconstructed in the CLEO tracking system is projected into the RICH and matched to the closest connected region. If the matching distance between the track projection and the center of the connected region is reasonably small and the total pulse height of the connected region sufficiently high, we associate this group of hits with the track. Local pulse height maxima in the remaining connected regions, so-called “bumps,” are taken as seeds for Cherenkov photons. We allow the pulse height maxima to touch each other by corners if the pulse height in the two neighboring pads is small relative to both bump heights. Hits adjacent to the bumps on the sides are assigned to them in order of decreasing bump pulse height.
To estimate the position of the photon conversion point we use the center-of-gravity method corrected for the bias towards the central pad. For many Cherenkov photons we are able to detect induced charge in only one pad. Since the pad dimensions are about $`8\times 8`$ mm<sup>2</sup>, the position resolution in this case is $`8\text{mm}/\sqrt{12}=2.3\text{mm}`$. For charged track intersections, which induce significant charge in many pads, the position resolution is $`0.76`$ mm. The position resolution for Cherenkov photons which generate multiple pad hits is somewhere in between these two values. In any case, the photon position error is not a significant contribution to the Cherenkov angle resolution (see below).
### 10.3 Corrections to the Track Direction
The resolution of the CLEO tracking system is very good in the bending view (the magnetic field is solenoidal in CLEO) . The track position and inclination angle along the beam axis is measured less precisely. The r.m.s. of the observed RICH hit residual is 1.7 mm. Since the RICH hit position resolution is 0.76 mm as measured by the residual in the perpendicular direction, the RICH can clearly help in pinning down the track trajectory. This, in turn, improves Cherenkov resolution, especially for the flat radiators for which we observe only half of the Cherenkov image and thus are quite sensitive to the tracking error. The improvement is as much as 50% in some parts of the detector.
### 10.4 Reconstruction of the Cherenkov Angle
Given the measured position of the Cherenkov photon conversion point in the RICH, the charged track direction and its intersection point with the LiF radiator, we calculate a Cherenkov angle for each photon-track combination ($`\theta _\gamma `$). We use the formalism outlined by Ypsilantis and Séguinot , except that we adopt a numerical method to find the solution to the equation for the photon direction, instead of simplifying it to a fourth-order polynomial. The latter would allow an analytical solution, but at the expense of introducing an additional source of error. Furthermore, using our numerical method, we calculate derivatives of the Cherenkov angle with respect to the measured quantities which allows us to propagate the detector errors and the chromatic dispersion to obtain an expected Cherenkov photon resolution for each photon independently ($`\sigma _\theta `$). This is useful since the Cherenkov angle resolution varies significantly even within one Cherenkov image. We use these estimated errors when calculating particle ID likelihoods and use them to weight each photon when measuring the average Cherenkov angle for a track.
### 10.5 Performance on Bhabha Events
We first view the physics performance on the simplest type of events, Bhabha events, and then subsequently in hadronic events. The distribution of Cherenkov angles measured for each photon in Bhabha events is shown in Fig. 40.
We note that Bhabha events have very low multiplicity compared with our normal hadronic events. They have two charged tracks present while the hadronic events have an average charged multiplicity of approximately 10 in CLEO III. In addition, the hadronic events have on the average 10 photons, mainly from $`\pi ^o`$ decays. All of these particles can interact in the calorimeter and the splash-back can hit the RICH photon detector.
The Cherenkov angle spectrum for single photons has an asymmetric tail and modest background. It is fit with a line-shape similar to that used when extracting photon signals from electromagnetic calorimeters . The functional form is
$$P(\theta |\theta _{exp},\sigma _\theta ,\alpha ,n)=$$
(2)
$`A\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\theta _{exp}\theta }{\sigma _\theta }}\right)^2\right]\mathrm{for}\theta <\theta _{exp}\alpha \sigma _\theta `$
$`A{\displaystyle \frac{\left(\frac{n}{\alpha }\right)^ne^{\frac{1}{2}\alpha ^2}}{\left(\frac{\theta _{exp}\theta }{\sigma _\theta }+\frac{n}{\alpha }\alpha \right)^n}}\mathrm{for}\theta >\theta _{exp}\alpha \sigma _\theta ,`$
$`A^1\sigma _\theta \left[{\displaystyle \frac{n}{\alpha }}{\displaystyle \frac{1}{n1}}e^{\frac{1}{2}\alpha ^2}+\sqrt{{\displaystyle \frac{\pi }{2}}}\left(1+\mathrm{erf}\left({\displaystyle \frac{\alpha }{\sqrt{2}}}\right)\right)\right].`$
Here $`\theta `$ is the measured angle, $`\theta _{exp}`$ is the “true” (or most likely) angle and $`\sigma _\theta `$ is the angular resolution. To use this formula, the parameter $`n`$ is fixed to value of about 5.
The data in Fig. 40 are fit using this signal shape plus a polynomial background function. We compare the results of these fits for the resolution parameter $`\sigma _\theta `$ as a function of radiator ring<sup>19</sup><sup>19</sup>19Effectively, this shows the dependence of the resolution on polar angle. for data and Monte Carlo simulation in Fig. 41. Here we use the symmetry of the detector about the center to map two full physical radiator rings into a single ring number, with ring 1 being closest to the middle. The single photon resolution averaged over the detector solid angles are 14.7 mr for the flat radiator and 12.2 mr for the sawtooth.
The number of photons per track within $`\pm 3\sigma `$ of the expected Cherenkov angle for each photon is shown in Fig. 42 and shown as a function of radiator row in Fig. 43. Averaged over the detector, and subtracting the background, we have a mean number of 10.6 photons with the flat radiators and 11.9 using the sawtooth radiators.
The resolution per track is obtained by taking a slice within $`\pm 3\sigma `$ of the expected Cherenkov angle for each photon and forming an average weighted by $`1/\sigma _\theta ^2`$. These track angles are shown in Fig. 44.
The r.m.s. spreads of these distributions are identified as the track resolutions. We obtain 4.7 mr for the flat radiators and 3.6 mr for the sawtooth. The resolutions as a function of radiator row are shown in Fig. 45.
The Cherenkov angular resolution is comprised of several different components. These include: error on the location of the photon emission point, the chromatic dispersion, the position error in the reconstruction of the detected photons, and finally the error on determining the charged track’s direction and position. These components are compared with the data in Fig. 46.
### 10.6 Performance on Hadronic Events in CLEO III
To resolve overlaps between Cherenkov images for different tracks we find the most likely mass hypotheses. Photons that match the most hypothesis within $`\pm 3\sigma `$ are then removed from consideration for the other tracks. To study the RICH performance in hadronic events in CLEO III<sup>20</sup><sup>20</sup>20Until now CLEO-c has been running below threshold for the production of $`D^{}`$’s. we use $`D^+\pi ^+D^0`$, $`D^0K^{}\pi ^+`$ events. The charge of the slow pion in the $`D^+`$ decay is opposite to the kaon charge in subsequent $`D^0`$ decay. Therefore, the kaon and pion in the $`D^0`$ decay can be identified without use of the RICH detector. The effect of the small combinatorial background is eliminated by fitting the $`D^0`$ mass peak in the $`K^{}\pi ^+`$ mass distribution to obtain the number of signal events for each momentum bin. The $`K^{}\pi ^+`$ invariant mass distribution selected by requiring that the $`K^{}\pi ^+\pi ^+`$ \- $`K^{}\pi ^+`$ mass difference be within 2.5 r.m.s. widths of the known mass difference is shown in Fig. 47. Here both the kaon and the pion are required to have momenta $`>0.6`$ GeV/c.
Single-photon Cherenkov angle distributions obtained on such identified kaons with the momentum above 0.7 GeV/c are plotted in Fig. 48. Averaged over all radiators, the single-photon resolution is 13.2 mr and 15.1 mr for sawtooth and flat radiators respectively. The background fraction within $`\pm 3\sigma `$ of the expected value is 12.8% and 8.4%. The background-subtracted mean photon yield is 11.8 and 9.6. Finally the per-track Cherenkov angle resolution is 3.7 mr and 4.9 mr.
### 10.7 Particle ID Likelihoods
For parts of the Cherenkov image for the sawtooth radiator, and for tracks intersecting more than one radiator there are some optical path ambiguities that impact the Cherenkov angle calculations. In the previous section we bypassed this problem by selecting the optical path that produces the closest Cherenkov angle to the expected one ($`\theta _{exp}^h`$) for the given particle hypothesis ($`h`$). There is some loss of information in this procedure, therefore, we use the likelihood method to perform particle identification instead of the per-track average angle. The likelihood method weights each possible optical path by the optical probability ($`P_{opt}`$), which includes length of the radiation path and the refraction probabilities obtained by the inverse ray tracing method:
$`L_h={\displaystyle \underset{j=1}{\overset{No.of\gamma s}{}}}\{P_{background}+\text{ }`$
$`{\displaystyle \underset{opt}{}}P_{opt}^jP_{signal}(\theta _\gamma ^{opt,j}|\theta _{exp}^h,\sigma _\theta ^{opt,j})\}`$
where, $`L_h`$ is the likelihood for the particle hypothesis $`h`$ ($`e`$, $`\mu `$, $`\pi `$, $`K`$ or $`p`$), $`P_{background}`$ is the background probability approximated by a constant and $`P_{signal}`$ is the signal probability given by the line-shape defined previously. In principle, the likelihood could include all hits in the detector. In practice, there is no point in inspecting hits which are far away from the regions where photons are expected for at least one of the considered hypotheses (we use $`\pm 5\sigma `$ cut-off).
An arbitrary scale factor in the likelihood definition cancels when we consider likelihood ratios for two different hypotheses. The likelihood conveniently folds in information about values of the Cherenkov angles and the photon yield for each hypothesis. For well separated hypotheses (typically at lower momenta) the photon yield that provides some discrimination. Since our likelihood definition does not know about the radiation momentum threshold, the likelihood ratio method can be only used when both hypotheses are sufficiently above the thresholds. When one hypothesis is below the radiation threshold we use the value of the likelihood for the hypothesis above the threshold to perform the discrimination.
The distribution of the $`2\mathrm{ln}\left(L_\pi /L_K\right)`$, is expected to behave as the difference $`\chi _K^2\chi _\pi ^2`$. This $`\chi ^2`$ difference obtained for 1.0-1.5 GeV/c kaons and pions identified with the $`D^{}`$ method is plotted in Fig. 49. Cuts at different values of this variable produce identification with different efficiency and fake rate. Pion fake rates for different values of kaon identification efficiency are plotted as a function of particle momentum in Fig. 50. Here when the fake rates get below a few percent there are other systematic effects that enter. For example, doubly Cabibbo suppressed decays where the $`D^o`$ decays into a $`K^+\pi ^{}`$ rather than a $`K^{}\pi ^+`$ have a relative branching fraction of 0.4% .
### 10.8 Efficiency and Fake Rates in CLEO-c
Here we use 180 pb<sup>-1</sup> integrated luminosity of CLEO-c data produced in $`e^+e^{}`$ collisions and recorded at the $`\psi ^{\prime \prime }`$ resonance (3.770 GeV). At this energy, the events consist of a mixture of pure $`D^+D^{}`$, $`D^o\overline{D}^o`$, three-flavor continuum event and $`\gamma \psi ^{}`$ events. There may also be small amounts of $`\tau ^+\tau ^{}`$ pairs and two-photon events.
In this study we select events containing at least one neutral $`D`$ candidate in the following decays $`D^oK^{}\pi ^{}\pi ^+\pi ^+`$, $`D^oK^{}\pi ^+`$ and $`D^oK^{}\pi ^+\pi ^o`$. (Charge conjugate modes are also used.) Event candidates in these modes are mostly signal with low background fractions. The $`D^o`$ candidate invariant mass plots are shown in Fig. 51. These mass plots are constructed by selecting decays where the sum of the measured energies is close to the electron beam energy and then using the measured beam energy to form the mass . We use this sample to look for two oppositely charged tracks present in the other side of the event not containing the tagged $`D`$. We then further select events where the $`\overline{D}^o`$ decays into $`K^\pm \pi ^{}`$. The momentum spectra of the kaon and the pion from this decay when the $`D`$ is produced on the $`\psi (3770)`$ is shown in Fig. 52.
We start by describing the analysis of the joint decays $`\overline{D}^oK^+\pi ^{}\pi ^+\pi ^{}`$ and $`D^oK^\pm \pi ^{}`$. Here we define the decay into $`K^{}\pi ^+`$ as “right” sign and the decay into $`K^+\pi ^{}`$ as “wrong” sign. In this case the wrong sign decays could result from one of three sources: background, doubly Cabibbo suppressed decays or $`D^o\overline{D}^o`$ mixing. We note that current measures of mixing limit it to $`<`$0.045% , while current measures of doubly Cabibbo suppressed decays are larger. For example, the modes $`K^{}\pi ^+`$, $`K^{}\pi ^{}\pi ^+\pi ^+`$ and $`K^{}\pi ^+\pi ^o`$ have rates of 0.35%, 0.42%, and 0.43%, respectively.
Fully reconstructed single tags for the $`K^+\pi ^{}\pi ^+\pi ^{}`$ mode (and its charge conjugate, which will not be explicitly mentioned in what follows) are reconstructed using the beam constrained mass. We use a 2$`\sigma `$ cut on $`\mathrm{\Delta }E`$ and require the mass to be between 1.86 and 1.870 GeV. We then form a double tag event using either right or wrong sign $`K\pi `$ decays. We use a tight cut on the $`K\pi `$ of $`|\mathrm{\Delta }E|`$, the difference between the measured energy and the beam energy within $`2\sigma `$, where $`\sigma `$ is the r.m.s. of the distribution. Since we need to use the RICH we impose a cut that both tracks be within $`|\mathrm{cos}(\theta )|<0.81`$, where $`\theta `$ is the angle of the track with respect to the beam line.
At first we do not use any RICH identification on the $`K\pi `$. Using a total of 180 pb<sup>-1</sup> we have 1158 such events where the kaon and pion are both in the RICH acceptance that give the right sign and the same sample yields 642 wrong sign events.
We now make three separate analyses: one where we identify only the kaon, one where we identify only the pion and one where we identify both the kaon and the pion. In the latter case we insist that there is significant discrimination in both cases or we do not accept the event.
Since the fake rates will be near 2%, the probability of getting both the kaon and the pion wrong is $`4\times 10^4`$, so that asking for a double identification is sufficient to ensure that we are getting the right answer for this level of tags.
The results are shown in Table 2. We find 1158 right sign events without using any particle identification, 970 right sign events with both particles identified, 15 wrong sign events with both particles identified, 25 events with only the kaon identified incorrectly and 36 events with only the pion identified incorrectly.
The 15 doubly identified wrong sign events are the combination of background, doubly Cabbibo suppressed decays and mixing. They correspond to a rate of these events of (1.5$`\pm `$0.4)%, consistent with them all being doubly Cabibbo suppressed decays, but somewhat larger . We subtract these events after correcting for the efficiency for the wrong sign candidates. This gives us a kaon fake rate of 1.1%, with an efficiency for events in the RICH of 88.5% and a pion fake rate of 3.7%, with an efficiency of 93.7%. We note that background is likely to be absent or small in these double tag events, but need to make a quantitative assessment.
To ascertain the background level we plot the beam constrained mass of the $`K^{}\pi ^+\pi ^+\pi ^{}`$ tag versus the $`K^\pm \pi ^{}`$ tag in Fig. 53. Events outside of the region where both masses are greater than 1.858 GeV are background. There is only one background event in the wrong sign plot indicating that the background is much less than one event.
We now consider the case where both neutral $`D`$’s decay into $`K^\pm \pi ^{}`$. One difference in this case with other cases is that doubly Cabbibo decays are forbidden due to Bose-Einstein statistics . The results are presented in Table 2. Here there is wrong sign doubly identified decay. This could be due to (a) background (b) $`D^o`$ mixing or (c) where both particles were incorrectly identified. We find zero events in the wrong sign doubly identified decay. We plot the beam constrained mass of the $`K^{}\pi ^+`$ tag versus the $`K^\pm \pi ^{}`$ tag in Fig. 54.
Our final mode uses $`K^{}\pi ^+\pi ^o`$ for the single tag. The results are also presented in Table 2. To ascertain the background level we plot the beam constrained mass of the $`K^{}\pi ^+\pi ^o`$ tag versus the $`K^\pm \pi ^{}`$ tag in Fig. 55. Although there appears to be some background in the right sign plot, the wrong sign shows no evidence of background.
The 4 doubly identified wrong sign events are the combination of background, doubly Cabbibo suppressed decays and mixing. They correspond to a rate of these events of (0.3$`\pm `$0.2)%, consistent with them all being doubly Cabibbo suppressed decays, but somewhat smaller.
Using all three of these mode combinations we find the rate of pions faking kaons of (1.10$`\pm `$0.37)%, with a pion efficiency for events in the RICH of (94.5$`\pm `$0.4)%. The rate of kaons faking pions is (2.47$`\pm `$0.38) %, with a kaon efficiency for events in the RICH of (88.4$`\pm `$0.6)%. The lower kaon efficiency arises because a significantly larger fraction of kaons than pions decay in this momentum region. It should be emphasized that these values are obtained for the entire running period between October, 2003 and February of 2005, and includes all possible system effects.
## 11 CONCLUSIONS
We have successfully constructed and operated a large, complex RICH detector in a particle physics experiment for over five years. The oxygen level has been kept below a few ppm in the “expansion volume” and the TEA photon conversion gas has been kept out, allowing for the Cherenkov photon yield to remain almost constant over the running period. We have lost some photon yield with a small $``$5% failure of electronics chips. One broken wire has caused an additional 1.7% loss and does somewhat effect the track efficiency. The total “cylindrical” detector thickness measured perpendicular to the axis is 13% of a radiation length.
The RICH is used during the normal course of most physics analyses using a standard set of criteria based on the minimum number of observed Cherenkov photons, usually 3, and the relative liklihood that a track is given type, either pion or kaon, for example.
The particle momenta for $`B`$ meson decay products seen by CLEO III are less than $`2.65`$ GeV/c. The detector provides excellent separation between pions and kaons at and below this cutoff. Separation between kaons and protons extends to even higher momentum, where it is used in charm studies. Thus, the physics performance has met design criteria. The RICH has provided crucially important particle separation in a number of important physics analyses including measurements of charmless hadronic two-body $`B`$ meson decays and the ratio $`\mathrm{\Gamma }(BDK)/\mathrm{\Gamma }(BD\pi )`$ , and measurement of the form-factors in $`D^o\pi ^{}\mathrm{}^+\nu `$ and $`D^oK^{}\mathrm{}^+\nu `$ decays .
CLEO is currently making an extensive study charm mesons and charmonium decays (called CLEO-c ). For these measurements the beam energy is lowered and the maximum particle momenta is about $`1.01.5`$ GeV/c. At these momentum the particle identification fake rates are at the 1% level.
## 12 ACKNOWLEDGEMENTS
The CLEO RICH project was funded primarily by the U. S. National Science Foundation which we deeply appreciate. We thank both the National Science Foundation and Department of Energy for supporting the University groups. We thank the late Tom Ypsilantis and Jacques Séguinot for early work on a similar system and for extensive discussions. Jeff Cherwinka helped with many engineering aspects of the system. Lee Greenler of PSL laboratories of Univ. of Wisconsin did much of the mechanical design. We thank Einar Nygard and Bjorn Sundal of IDEAS for their work on the front-end hybrid design. Paul Gelling contributed to the electronics infrastructure. We especially appreciate the efforts of Charles Brown, Lou Buda and Lester Schmutzer of the Syracuse Physics Dept. machine shop who made many of the components. We thank Peter Reed, Heather Lane, Dave Smith, Don Moulton at Optovac for their hard work during the four years of crystal production. Ken Powers helped with the plating of the windows. We thank the accelerator group at CESR for excellent efforts in supplying luminosity.
## Appendix A. VUV SPECTROPHOTOMETERS
The CH<sub>4</sub>-TEA photosensitive gas inside the RICH chambers has an appreciable photo-absorption cross-section in a narrow VUV band $`\lambda =150\pm 15`$nm. Satisfactory reconstruction of the Cherenkov cone geometry requires that the windows and the radiator crystals (especially their top surface), be sufficiently transparent in this wavelength band. A crystal surface which is poorly polished or contaminated would reduce significantly the number of photons emitted from the radiator crystals and subsequently degrade the Cherenkov angle measurement resolution.
Three VUV spectrophotometers used to measure the transmission of the crystals were built and placed at the three locations associated with crystal production and handling.
### A.1. LiF VUV Transmission Spectrophotometer at SMU
Prior to installation of individual crystals onto the RICH Inner Cylinder, each crystal’s transparency was measured using a specially constructed VUV spectrophotometer at four wavelengths: $`\lambda =135,\mathrm{\hspace{0.17em}142.5},\mathrm{\hspace{0.17em}150}\mathrm{and}165`$nm. The essential components of this transmission spectrophotometer are shown in Fig. 56. A vacuum monochromator with a deuterium lamp is the VUV light source. A test box 250 liter vacuum vessel contains the crystal to be measured, along with a two-dimensional stage driven by stepper motors. The crystal’s top surface is held perpendicular to the direction of the probe radiation and collimators produce a beam spot of radius $``$1 mm on the crystal surface. The stage can reliably position a crystal inside the vessel with a spatial precision of better than $`25\mu `$m along each of the stage’s mutually perpendicular axes. It moves the crystal perpendicular to the beam direction in a raster style and periodically stops at discrete points while the transparency measurement is made. The typical spacing between measurements points is 1 cm. For calibration purposes, during each motion along the width of the crystal, the stage both moves the crystal completely out of the beam and completely blocks the beam with its frame, so as to provide measurements along each row of the scan that correspond to 100% and 0% transparency, respectively.
Both, the monochromator and the vacuum vessel, are evacuated by dedicated turbomolecular vacuum pumps to minimize absorption of the probe radiation by any residual air or water molecules. A photomultiplier tube (PMT), whose front window is coated with an ultraviolet wave-shifting compound (sodium salicylate ), is attached to the vacuum vessel and measures the intensity of the transmitted VUV light from the monochromator after its passage through the crystal. PMT output signals are amplified, filtered, digitized and then processed by a LabVIEW-based program to determine the crystal transparency map in real time. Data is also written to disk for further off-line analysis.
The measurement process was necessarily highly automatic and the control software ran on a PC, which in turn communicated with the data acquisition system and the stage controllers. The software controlled the high voltage for the PMT, monitored pressures in the two vacuum systems (monochromator and test box), and through a standard GPIB interface adjusted the monochromator grating to the desired wavelength. A graphical user interface allowed the user to set efficiently a wide range of control parameters to measure radiator crystals in a variety of modes, to measure crystal ingot test bars, and to perform calibration runs of various types.
Extensive adjustment and calibration of the spectrophotometer system were performed to minimize systematic errors and to measure overall system time stability and reproducibility. For example, the linearity of the electronic readout chain was measured to be much better than 1% over the range of readout voltages and transmission measurements separated in time by more than 1 month of the same LiF sample show agreement within 1%. Overall, $`500,000`$ individual LiF transmission measurements were made.
### A.2. VUV Spectrophotometer at Syracuse
A second VUV spectrophotometer similar in design was constructed at Syracuse University for the purpose of testing the window crystals as well as the sawtooth radiators.
The window crystals were measured in a manner substantially similar to that described above for the planar radiators. Transmissions were measured at three wavelengths (135 nm, 150 nm, and 165 nm), and over a grid of 30 positions over the surface of the crystal window. For half-sized windows, there were 15 positions per crystal. The results are summarized in Section 3.2.
In this spectrophotometer, however, systematic uncertainties in the transmission measurements were reduced using different techniques. Before the light beam enters the vacuum tank, it is passed through a chopper which opens and closes at $``$40 Hz. The resulting phototube current appears to be a square wave, with an amplitude that provides a measure of the light output, while automatically subtracting stray light and dark current offsets. Also, in general, the transmission measurement is made by performing “crystal-in/crystal-out” measurements: we divide the light output when the beam passes through a crystal (as registered by the current in the phototube) by the light output when the crystal is removed from the beam. Measuring this ratio attempts to divide out drifts in the lamp output or gain variations in the PMT over time. Making this reference measurement frequently means that the measurement cannot drift appreciably from point-to-point within a scan. Individual transmission measurements have a “statistical” uncertainty of typically 0.5–1.0% due to averaging the photocurrent, and a 0.5–1.0% “systematic” uncertainty due to signal drift over a typical measurement interval.
The sawtooth radiators required a special technique in order to measure their transmission. For this, the spectrophotometer is configured as shown in Fig. 57. The monochromator light beam enters through the bottom of the vacuum tank, and passes through a 0.020 inch collimating slit. The pencil beam has a Gaussian shape with $`\sigma =25\mu \text{m}`$. The light reflects off of an adjustable mirror to produce a 15 deflection of the beam. It passes through a 42 VUV-polished prism, which is mounted on the two-dimensional stage. The light is wavelength-shifted by the sodium salicylate covered glass window, and detected in the photomultiplier tube.
The angular deflection of 15 is set so as to probe the transmission at the smallest possible angle that is not totally internally reflected and that accrues appreciable signal. Although the Fresnel coefficients are functions of incident angle at an interface, the important point is that the prism has the same shape as the sawtooth radiator grooves. Hence, by moving the stage such that the prism and sawtooth crystal passes through the beam, a relative transmission measurement using the prism as a standard may be made. Uncertainties were at the 1% level.
A third spectrophotometer, a duplicate to this system, was constructed and placed on site at the crystal manufacturing company in order to make quality assurance measurements of the window and planar radiator crystals before shipping.
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# Interface ordering and phase competition in a model Mott-insulator–band-insulator heterostructure
Electronic address: \]okapon@phys.columbia.edu
## Abstract
The phase diagram of model Mott-insulator–band-insulator heterostructures is studied using the semiclassical approximation to the dynamical-mean-field method as a function of thickness, coupling constant, and charge confinement. An interface-stabilized ferromagnetic phase is found, allow the study of its competition and possible coexistence with the antiferromagnetic order characteristic of the bulk Mott insulator.
Fabrication and investigation of heterostructures involving correlated-electron materials Imada98 ; Tokura00 are an important direction in material science. Understanding of the electronic properties near the interfaces and surfaces is not only of scientific interest but is also indispensable to realize electronic devices exploiting the unique properties of correlated-electron materials. A variety of heterostructures have been fabricated and studied including high-$`T_c`$ cuprates,Ahn99 ; Ahn02 Mott-insulator and band-insulator heterostructure,Ohtomo02 and superlattices of transition-metal oxides.Izumi01 ; Biswas00 ; Biswas01 Interestingly, the heterostructures comprising of a Mott-insulator and a band-insulator were reported to show metallic behavior.Ohtomo02
A fundamental question raised by those studies is “what electronic phases are realized at interfaces.” For the vacuum-bulk interface (i.e., the surface), Potthoff and Nolting,Potthoff99 Schwieger et al.,Schwieger03 and Liebsch Liebsch03 have argued that the reduced coordination may enhance correlation effects. The enhanced correlations could presumably induce surface magnetic ordering, although this possibility was not discussed in Refs. Potthoff99, ; Schwieger03, ; Liebsch03, . Matzdorf et al. proposed that ferromagnetic ordering is stabilized at the surface of two-dimensional ruthenates by a lattice distortion,Matzdorf00 but this is not yet observed. Surface ferromagnetism had been also discussed in a mean field treatment of the Hubbard model by Potthoff and Nolting.Potthoff95 Similarly, the effect of bulk strain on the magnetic ordering in perovskite manganites was discussed by Fang et al.Fang00
All these studies dealt with systems in which charge densities remain unchanged from the bulk values, and physics arising from the modulation of charge density was not addressed. A crucial aspect of the recently fabricated heterostructures is charge inhomogeneity, caused by the spreading of electrons from one region to another. Strongly correlated materials typically possess interesting density-dependent phase diagrams; raising the possibility of interesting phase behavior at interfaces. In this paper we use the semiclassical approximation (SCA)Okamoto05 to the dynamical-mean-field methodGeorges96 to explore the phase behavior of a model Mott-insulator–band-insulator heterostructure. The SCA is computationally inexpensive and a good representation of phase diagrams and transition temperatures in several models, allowing us to investigate the phase behavior for a wide range of parameters, and in particular to access the $`T>0`$ regime which Hartree-Fock and related approximations fail to represent adequately.Okamoto04a We observe antiferromagnetic ordering in regions with charge density $`1`$ characteristic of a bulk Mott insulator, while ferromagnetic ordering is found to be a surface effect supported by an intermediate charge density and a strong coupling. However, we have found that, despite the successes noted in previous work, the SCA overestimates ferromagnetism on the lattice we study here. Therefore, our results should be regarded as qualitative explanations of the type of phase behavior which may occur rather than as quantitative statements about the Hubbard-model phase diagram.
We study the model heterostructure introduced in Ref. Okamoto04b, . The Hamiltonian is a simplified representation of conduction bands of the systems studied in Ref. Ohtomo02, with the orbital degeneracy neglected. We consider heterostructures formed by varying the $`A`$-site of a $`AB`$O<sub>3</sub> perovskite lattice. The electrons of interest reside on the $`B`$-site ions, which form a simple cubic lattice with sites labeled by $`i`$ as $`\stackrel{}{r}_i=(x_i,y_i,z_i)=a(n_i,m_i,l_i)`$ with the lattice constant $`a`$ set to unity. We assume each $`B`$-site has a single orbital; electrons hop between nearest neighbor sites with the transfer $`t`$. The electrons interact via a on-site interaction $`U`$ and a long-ranged Coulomb repulsion. The heterostructure is defined by $`n`$ planes of charge +1 counterions placed on the $`A^{}`$ sublattice of $`A`$-site ions at positions $`\stackrel{}{r}_j^A^{}=a(n_j+1/2,m_j+1/2,l_j+1/2)`$, with $`\mathrm{}<n_j,m_j<\mathrm{}`$ and $`l_j=1,\mathrm{},n`$. (Coordinate $`z`$ will be shifted such that the center of the heterostructure, $`A^{}`$ sublattice, comes to $`z=0`$.) Charge neutrality requires that the areal density of electrons is $`n`$. The resulting Hamiltonian is $`H=H_{band}+H_{int}+H_{Coul}`$ with
$`H_{band}`$ $`=`$ $`t{\displaystyle \underset{ij,\sigma }{}}(d_{i\sigma }^{}d_{j\sigma }+H.c.),`$ (1)
$`H_{int}`$ $`=`$ $`U{\displaystyle \underset{i}{}}n_in_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{ij}{\sigma ,\sigma ^{}}}{}}{\displaystyle \frac{e^2n_{i\sigma }n_{j\sigma ^{}}}{\epsilon |\stackrel{}{r}_i\stackrel{}{r}_j|}},`$ (2)
$`H_{Coul}`$ $`=`$ $`{\displaystyle \underset{i,j,\sigma }{}}{\displaystyle \frac{e^2n_{i\sigma }}{\epsilon |\stackrel{}{r}_i\stackrel{}{r}_j^A^{}|}}.`$ (3)
Note that $`U0`$ on all sites. A dimensionless measure of the strength of the long-ranged Coulomb interaction is $`E_c=e^2/(\epsilon at)`$ with the dielectric constant $`\epsilon `$. In most of our analysis, we choose $`E_c=0.8`$. This corresponds to $`t0.3`$ eV, $`a4`$ Å, and $`\epsilon =15`$, which describe the system studied in Ref. Ohtomo02, . The charge profile is found not to depend in an important way on $`\epsilon `$, but the stability of magnetic orderings does because this is sensitive to the details of the charge density distribution as discussed later.
The basic object of our study is the electron Green’s function. In general, this is given by
$$G_\sigma (\stackrel{}{r},\stackrel{}{r^{}};\omega )=[\omega +\mu H_{band}H_{Coul}\mathrm{\Sigma }_\sigma (\stackrel{}{r},\stackrel{}{r^{}};\omega )]^1,$$
(4)
with the chemical potential $`\mu `$ and the electron self-energy $`\mathrm{\Sigma }`$. We consider heterostructures with either in-plane translational invariance or $`N_s`$-sublattice antiferromagnetism. The Green’s function and self-energy are therefore functions of the variables $`(z,\eta ,z^{},\eta ^{},\stackrel{}{k}_{})`$ where $`\eta `$ and $`\eta ^{}(=1,\mathrm{},N_s)`$ label the sublattice in layers $`z`$ and $`z^{}`$, respectively, and $`\stackrel{}{k}_{}`$ is a momentum in the (reduced) Brillouin zone. As in Ref. Okamoto04b, , we approximate the self-energy as the sum of a static Hartree term $`\mathrm{\Sigma }_\sigma ^H`$ arising from the long-ranged part of the Coulomb interaction and a dynamical part $`\mathrm{\Sigma }_\sigma ^D(\omega )`$ arising from local fluctuations. Generalizing the inhomogeneous case of the dynamical-mean-field theory (DMFT),Schwieger03 we assume that the self-energy is only dependent on layer $`z`$ and sublattice $`\eta `$. Thus, the dynamical part of the self-energy is written as
$$\mathrm{\Sigma }_\sigma ^D\mathrm{\Sigma }_\sigma ^D(z,\eta ;\omega ).$$
(5)
The $`z\eta `$-dependent self-energy is determined from the solution of a quantum impurity model Georges96 with the mean-field function fixed by the self-consistency condition
$$G_\sigma ^{imp}(z,\eta ;\omega )=N_s\frac{d^2k_{}}{(2\pi )^2}G_\sigma (z,\eta ,z,\eta ,\stackrel{}{k}_{};\omega ).$$
(6)
We shall use the DMFT, which is a mean-field approximation, to calculate magnetic phase diagrams of effectively two dimensional systems, so some mention of fluctuation effects is needed. In general, the calculated transition temperatures are to be understood as crossover scales below which the magnetic correlation length $`\xi `$ grows rapidly $`\xi \mathrm{exp}(2\pi \rho _s/T)`$ with spin stiffness $`\rho _s`$ Sachdev99 discussed in more detail below. True long-ranged order may be induced at $`T>0`$ by an Ising anisotropy or by coupling in a system made of repeated heterostructure units. In any event, the rapid growth of the correlations for $`T<\rho _s`$ means that the properties are effectively those of an ordered state. We will find that the low-$`T`$ $`\rho _s`$ is sufficiently large relative to the calculated transition temperature that fluctuation effects are not crucial.
In general for the heterostructure with $`L`$ layers with $`N_s`$ sublattices, one must solve $`L\times Ns`$ independent impurity models. Due to the self-consistency condition \[cf. Eq. (6)\] and to compute the charge density $`n_\sigma (z,\eta )=\frac{d\omega }{\pi }f(\omega )\mathrm{Im}G_\sigma ^{imp}(z,\eta ;\omega )`$ with $`f`$ the Fermi distribution function, it is required to invert the $`(L\times Ns)^2`$ Green’s function matrix at each momenta and frequency. This time consuming numerics restricts the size of the unit cell. In this study, we consider the commensurate magnetic states with up to two sublattices, $`N_s=1`$ and 2, on each layer and with the charge density independent of the sublattices, i.e., paramagnetic (PM), ferromagnetic (FM) states, and (layer-) antiferromagnetic (AF) state where antiferromagnetic (FM) planes with moment alternating from plane to plane. Note that the AF state extrapolates to the bulk AF state with the magnetic vector $`\stackrel{}{q}=(\pi ,\pi ,\pi )`$ at $`n\mathrm{}`$. By symmetry, the number of quantum impurity models one must solve is reduced to $`L`$. However, solution of the impurity models is a time consuming task, and an inexpensive solver is required. In Ref. Okamoto04b, , to study the evolution of the low-energy quasiparticle band and high-energy Hubbard bands as a function of position, we applied two-site DMFT Potthoff01 which is a simplified version of exact-diagonalization method. At $`T=0`$, this method is known to give reasonable result for Mott metal-insulator and magnetic transitions. However, small number of bath orbitals is known to be insufficient to describe the thermodynamics correctly.Okamoto05
For the investigations presented here, we use the semiclassical approximation, which is computationally inexpensive and has been found to be reasonably accurate for phase boundaries and excitation spectra of several models, including the half-filled Hubbard model and, at all fillings, for the $`d=3`$ and $`d=\mathrm{}`$ face-centered-cubic (FCC) lattices.Okamoto05 We note however that, perhaps because it does not properly include quasiparticle coherence, the method overemphasizes ferromagnetism at intermediate density, giving for the FCC lattices transition temperatures $`50\%`$ higher than those found by quantum Monte-Carlo,Okamoto05 and for cubic lattices, finding ferromagnetism at $`n_{tot}=0.5`$ and moderately large $`U`$ (of order the critical value for the Mott transition) when other methodsZitzler02 suggest ferromagnetism is confined to very large $`U`$ and $`n_{tot}`$ near 1. Thus, the application of the SCA to Hubbard heterostructure allows a convenient exploration of the interplay between different phases, but probably does not provide a quantitatively reliable picture of the phase diagram.
First, we investigate the magnetic behavior at finite temperature. The upper panel of Fig. 1 shows our calculated phase diagram in the interaction-temperature plane for heterostructures with various thicknesses for charge binding parameter $`E_c=0.8`$. The one-layer heterostructure is PM at weak to moderate interactions, and FM at strong interactions. The two- and three-layer heterostructures are AF at weak to intermediate interaction, and become FM at stronger interactions with almost the same $`T_C`$ for $`n=2`$ and 3. Antiferromagnetic Néel temperature $`T_N`$ is found to be strongly dependent on the layer thickness; it increases with the increase of layer thickness. These $`T_N`$’s are substantioally reduced from the the bulk values; maximum value $`T_N^{max}/t0.47`$ at $`U/t10`$, whereas Curie temperature $`T_C`$’s are almost the bulk 2$`d`$ values at $`n_{tot}0.5`$. The strong dependence of $`T_N`$ on thickness may be understood from the bulk phase diagram; antiferromagnetism is stabilized only very near to half-filling, and in the thinner heterostructures the charge-spreading effect reduces the density too much. This physics is seem from a different point of view in the lower panel of Fig. 1. Hartree-Fock studies of this and related models find a layer-AF phase. This phase is not found in our DMFT analysis.
The lower panel of Fig. 1 presents a detailed study of the $`n=2`$ heterostructure showing how changes in the charge confinement parameter $`E_c`$ affect the physics. The filled and open points (left-hand axis) show the variation of the Curie and Néel temperatures, respectively. The light solid and light broken lines (right-hand axis) show the variation of charge density on the central and next to central layers, respectively. It is seen that the AF ordering is rapidly destabilized with the decrease of $`E_c`$ (and concomitant decrease of charge density), and $`T_N`$ is seen to be correlated to that of the charge density at $`z=0`$. On the contrary, $`T_C`$ has a weak variation with $`E_c`$. This indicates that the FM ordering is favored by the intermediate charge density as discussed in the bulk single-band Hubbard model;Denteneer95 at large $`E_c`$, the magnetization is large on the outer layer and small in the inner layer, at small $`E_c`$, the situation is reversed.
These arguments are confirmed by our calculations of the spatial variation of the magnetization density, reported in Fig. 2, which shows numerical results for a 4-layer heterostructure with counterions at $`z=\pm 0.5`$ and $`\pm 1.5`$, and $`E_c=0.8`$. The upper panel of Fig. 2 shows the magnetization in the FM state. In DMFT (filled circles), only the layers near the interfaces ($`|z|1`$–2) have large polarization and inner layers in the heterostructure have small moments. This explains the weak $`n`$-dependence of $`T_C`$ of thick heterostructures (see the upper panel of Fig. 1). In HF (open circles), all layers in the heterostructure are highly polarized. In contrast in an AF state, result of the in-plane staggered magnetization by DMFT and HF agree well as shown in the lower panel of Fig. 2. For comparison, the total charge density is also plotted (filled squares). The in-plane staggered magnetization is large only at inner layers where the charge density is close to 1. Note that the staggered magnetization in the outer layers ($`|z|2`$) has the same sign as in the outermost layers ($`|z|=1`$) indicating that the outer layers are not intrinsically magnetic.
We now return to the issue of fluctuation effects. We study a model which is two dimensional and order parameters with a continuous spin rotation symmetry; thus, at $`T>0`$, rather the calculated $`T_C`$ means a crossover to a low-$`T`$ “almost ordered” phase characterized by a exponentially growing correlation length $`\xi \mathrm{exp}(2\pi \rho _s/T)`$Sachdev99 where the key stiffness $`\rho _s`$ is given by the second derivative of the free energy with respect to the order parameter orientation. We have computed $`\rho _s`$ for a single ferromagnetic plane with charge density 0.5 finding $`\rho _st/8`$ We observe that in all cases several planes exhibit the relevant order, so that the total stiffness is larger by a factor 2–3. Thus we see that the stiffness is large enough relative to the DMFT $`T_{C,N}`$’s that the purely two dimensional fluctuation effects are of minor importance, leading to $`\xi >\mathrm{exp}(4`$$`5)`$ for $`T<T_{C,N}/2`$.
Spin distributions presented in Fig. 2 can be understood from the single-particle spectral functions. In Fig. 3 are presented the DMFT results for the layer- and sublattice-resolved spectral functions $`A_\sigma (z,\eta ;\omega )=\frac{1}{\pi }\mathrm{Im}G_\sigma ^{imp}(z,\eta ;\omega +i0^+)`$ for the FM (upper panel) and the AF (lower panel) states of 4-layer heterostructure with the same parameters as in Fig. 2. These quantities can in principle be measured by spin-dependent photoemission or scanning tunneling microscopy. As noticed in Ref. Okamoto04b, , spectral function outside of the heterostructure ($`|z|2`$) is essentially identical to that of the free tight-binding model $`H_{band}`$, and electron density is negligibly small. With approaching the interfaces ($`|z|=2`$), the spectral function shifts downwards and begins to broaden. In the FM case, magnetic ordering is possible only near the interface ($`|z|1`$–2) carrying the intermediate charge density. Inside the heterostructure ($`|z|2`$), clear Hubbard gap exists due to the large $`U`$ and uniform polarization is hard to achieve. On the contrary, high charge density is necessary to keep the staggered magnetization in the AF case as seen as a difference between up and down spectra in the lower panel of Fig. 3.
So far, we have discussed the competition between ferromagnetic and antiferromagnetic phases and spatial distributions of charge and magnetic densities establishing that ferromagnetism is a surface effect. Finally, we discuss the interplay between phases. We note that the charge density or magnetization density typically varies over $`3`$ unit cell range, so for thin heterostructures and moderate charge confinement energies, the entire heterostructure exhibits a single phase, controlled by the instability exhibiting the highest transition temperature (cf. Fig. 1). However, for thicker heterostructure or stronger confinement we believe that one can observe an ordered state involving an antiferromagnetic center with a ferrimagnetic “skin.” A hint of this behavior can be observed in Fig. 2. Consider the central layer; this has an occupancy $`n_{tot}0.9`$, so from Fig. 1 and rescaling account for the different $`U`$, we would expect $`T_N0.1t`$. Now this layer is not isolated; hoppings to the layers of $`z=\pm 1`$ leads to a effective polarizing field of the order of $`2(1n_{tot})tm/n_{tot}`$ where the first factor is the number of layers, the second is the hopping amplitude renormalized by strong correlations, and the third factor is the relative spin polarization; putting these factors together gives a polarizing field of about $`0.15t`$, approximately equal to the AF coupling. Thus for the central layer of this heterostructure ferromagnetic and antiferromagnetic tendencies are very closely balanced, but for thicker systems (not at present computationally accessible) or perhaps for stronger charge confinement, an antiferromagnetic center will occur.
To summarize, we have presented a semiclassical DMFT study of magnetic phase behavior of a model Mott-insulator–band-insulator heterostructure in which the behavior is controlled by the spreading of the electronic charge out of the confinement region. Magnetic phase diagram is investigated as a function of layer thickness, temperature, and interaction strength. Ferromagnetic ordering is found to be a surface effect stabilized at an interface region with moderate charge density, while antiferromanetic ordering is found at a region with high density $`1`$ characteristic of the bulk Mott insulator, and Néel temperature is sensitive to the layer thickness and charge confinement energy. These magnetic orderings may coexist in very thick heterostructure exhibiting a ferromagnetic “skin” and an antiferromagnetic “core.”
We acknowledge fruitful discussions with G. Kotliar, P. Sun, J. Chakhalian, D. Vollhardt, and Th. Pruschke. This research was supported by JSPS (S.O.) and the DOE under Grant No. ER 46169 (A.J.M.).
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# Regular homotopy classes of locally generic mappings
## 1. Introduction
Our work was motivated by the paper of U. Pinkall , which classifies immersions of compact surfaces into $`^3`$ up to regular homotopy, allowing diffeomorphisms of the source manifold $`M^2`$. So two immersions $`f,g:M^2R^3`$ are considered equivalent if there is a diffeomorphism $`\phi `$ of $`M^2`$ such that $`f=g\phi `$. This notion is different from regular homotopy, yielding an interesting classification of immersed surfaces using the Arf invariant. That paper also gives generators for the abelian semigroup of immersed surfaces with the connected sum operation. Professor András Szűcs asked me what happens with Pinkall’s classification if we allow cross-cap (also called Whitney-umbrella) singularities. The notion of regular homotopy has to be revised and, unlike for immersions, for singular maps it turns out that all natural definitions are equivalent. In fact, we prove that for singular mappings (i.e. not immersions) of a closed connected surface the number of cross-caps totally determines the regular homotopy class, diffeomorphisms of the source manifold are not needed. Thus the approach of U. Pinkall and the classical regular homotopy classification give the same result for singular maps.
In the final section of our paper we present an application of the above result to the study of the path components of the space of those immersions of a closed connected surface into $`^4`$ whose projections into $`^3`$ are locally generic, i.e. may have cross-cap singularities.
It is a natural question if Theorem 2.6 generalizes to higher dimensions, for locally generic maps of a closed n-manifold $`M^n`$ into $`^{2n1}`$. Our methods of proof for Theorem 2.6 do not seem to work if $`n>2`$ since they rely heavily on the results of surface topology. However, I could prove the general result in the case when $`n>3`$ and $`M^n`$ is 2-connected. I will publish this in a separate paper. I want to emphasize that the results of Section 3 easily generalize for any closed manifold $`M^n`$ provided that the generalization of Theorem 2.6 holds true for $`M^n`$.
I would like to take this opportunity to express my gratitude to Professor András Szűcs for drawing my attention to this problem and for his constant support and encouragement. I would also like to thank Professor Balázs Csikós who read the first version of this paper and suggested several improvements.
## 2. The main result
Let us start by introducing the notion of a locally generic mapping of a closed $`n`$-manifold $`M^n`$ into a $`(2n1)`$-manifold $`N^{2n1}`$ for $`n2`$.
###### Definition 2.1.
$`f:M^nN^{2n1}`$ is called *locally generic* if it is an immersion except for cross-cap singularities. The set of singular points of $`f`$ in $`M^n`$ is denoted by $`S(f)`$. ($`\left|S(f)\right|<\mathrm{}`$ because $`M^n`$ is compact.) Let us denote by $`\text{Lgen}(M^n,N^{2n1})`$ the subspace of locally generic mappings in $`C^{\mathrm{}}(M^n,N^{2n1})`$ endowed with the $`C^{\mathrm{}}`$ topology.
Recall that a map $`f:M^nN^{2n1}`$ is called generic (or stable), if it is an immersion with normal crossings except in a finite set of points, moreover the singular points of $`f`$ are non-multiple cross-cap points. This explains our terminology. Whitney proved that the set of stable maps is dense open in $`C^{\mathrm{}}(M^n,^{2n1})`$ with respect to the $`C^{\mathrm{}}`$-topology. Studying the double point set of a generic mapping close to the locally generic mapping $`f`$ in the $`C^{\mathrm{}}`$-topology, one can easily verify that for $`M^n`$ closed $`\left|S(f)\right|`$ is an even integer. (For the cross-caps are precisely endpoints of double-point curves.)
###### Definition 2.2.
Two locally generic mappings $`f,g:M^nN^{2n1}`$ are called *regularly homotopic* ( denoted by $`fg`$ ) if there is a smooth mapping $`H:M^n\times [0,1]N^{2n1}`$ such that $`H_t`$ is locally generic for each $`t[0,1]`$, moreover $`H_0=f`$ and $`H_1=g`$. Here $`H_t(x)=H(x,t)`$ for $`xM^n`$ and $`t[0,1]`$.
The following definition will be especially useful in the case $`n=2`$.
###### Definition 2.3.
Two locally generic mappings $`f,g:M^nN^{2n1}`$ are called *image-homotopic* if there is a diffeomorphism $`\phi `$ of $`M^n`$ such that $`f\phi `$ is regularly homotopic to $`g`$. We denote this by $`\text{Im}(f)\text{Im}(g)`$.
###### Proposition 2.4.
If $`fg`$ or just $`\text{Im}(f)\text{Im}(g)`$, then $`\left|S(f)\right|=\left|S(g)\right|`$. In fact, if $`H:M^n\times [0,1]N^{2n1}`$ is a regular homotopy, then $`\left|S(H_t)\right|=\left|S(H_0)\right|`$ for every $`t[0,1]`$.
###### Proof.
If $`\phi `$ is a diffeomorphism of $`M^n`$, then $`S(f)=\phi (S(f\phi ))`$, so $`\left|S(f)\right|=\left|S(f\phi )\right|`$. Thus we can suppose that $`fg`$. From the definition of stability it is clear that every locally generic mapping $`h`$ has a neighborhood $`U_h`$ in the Whitney $`C^{\mathrm{}}`$ topology such that for every $`h^{}U_h`$ we have $`\left|S(h^{})\right|=\left|S(h)\right|`$ (because every $`pM^n`$ has a neighborhood $`U`$ such that $`h|U`$ is equivalent to $`h^{}|U`$ and since $`M^n`$ is closed, see \[1, p. 72\]). Thus the function $`\left|S(.)\right|:\text{Lgen}(M^n,N^{2n1})`$ is locally constant. So if $`H`$ is a regular homotopy connecting $`f`$ and $`g`$, this implies that $`H_t`$ is a continuous path in $`\text{Lgen}(M^n,N^{2n1})`$, along which $`\left|S(H_t)\right|`$ is constant. ∎
###### Remark 2.5.
We have defined the notions of regular homotopy and image-homotopy between two locally generic maps $`f`$ and $`g`$. Proposition 2.4 implies that if $`f`$ and $`g`$ are immersions then they are regularly homotopic, resp. image-homotopic as immersions iff they are those as locally generic maps.
If $`H:M^2\times [0,1]^3`$ is a regular homotopy between two locally generic mappings $`f=H_0`$ and $`g=H_1`$ and $`k=\left|S(f)\right|=\left|S(g)\right|`$, then there exists curves $`\gamma _1,\mathrm{}\gamma _k:[0,1]M^2`$ such that $`S(H_t)=\{\gamma _j(t):1jk\}`$ for every $`t[0,1]`$. We define the bijection $`i_H:S(f)S(g)`$ the following way : for $`1jk`$ let $`i_H(\gamma _j(0))=\gamma _j(1)`$.
Now we can state the main result of this paper yielding a converse of Proposition 2.4 for singular mappings.
###### Theorem 2.6.
Let $`M^2`$ be a closed connected surface and suppose that $`f,g:M^2^3`$ are locally generic mappings with $`\left|S(f)\right|=\left|S(g)\right|>0`$. Then $`fg`$. Moreover for any bijection $`i:S(f)S(g)`$ there exists a regular homotopy $`H`$ connecting $`f`$ and $`g`$ such that $`i_H=i`$.
Let us now list a few interesting corollaries of this. A surprising consequence of Theorem 2.6 is the following: If $`U`$ is the standard locally generic mapping of $`S^2`$ into $`^3`$ with two cross-cap points then for any two immersions $`f,g:M^2^3`$ the connected sums $`f\mathrm{\#}U`$ and $`g\mathrm{\#}U`$ become regularly homotopic as locally generic maps !
Another consequence of our theorem is that, given a closed connected surface $`M^2`$, we can easily produce a full list of representatives of all regular homotopy classes of locally generic maps $`M^2^3`$. (We only do this for singular maps, for immersions see .) First suppose that $`M^2`$ is orientable. If we denote by $`i_M`$ the standard embedding of $`M^2`$ into $`^3`$ then
$$i_M\mathrm{\#}\underset{n}{\underset{}{U\mathrm{\#}\mathrm{}\mathrm{\#}U}}$$
is a representative for the class of locally generic maps with $`2n>0`$ singular points. Now suppose that $`M^2`$ is a non-orientable surface of genus $`g`$. We denote the Boy surface by $`B`$. Now
$$\underset{g}{\underset{}{B\mathrm{\#}\mathrm{}\mathrm{\#}B}}\mathrm{\#}\underset{n}{\underset{}{U\mathrm{\#}\mathrm{}\mathrm{\#}U}}$$
is a locally generic mapping of $`M^2`$ into $`^3`$ with $`2n`$ cross-cap points.
Perhaps the following construction can be visualized more easily: Let us denote by $`V`$ the well-known locally generic mapping of $`P^2`$ into $`^3`$ having two singular points ($`VB\mathrm{\#}U`$). Then any singular locally generic mapping is regularly homotopic to one of the form $`i_N\mathrm{\#}U\mathrm{\#}\mathrm{}\mathrm{\#}U\mathrm{\#}V\mathrm{\#}\mathrm{}\mathrm{\#}V`$, where $`N^2`$ is orientable. (The left side of Figure 1 depicts $`A_2\mathrm{\#}U\mathrm{\#}U`$, where $`A_2`$ is the orientable surface of genus $`2`$. The right side of Figure 1 illustrates $`B\mathrm{\#}B\mathrm{\#}B\mathrm{\#}U\mathrm{\#}U\mathrm{\#}UV\mathrm{\#}V\mathrm{\#}V`$.)
Pinkall determined the abelian semigroup $`H`$ of immersed surfaces in $`^3`$ with the connected sum operation (see ). If we consider the extended semigroup $`\stackrel{~}{H}`$ of locally generic surfaces, then only $`U`$ is needed as a new generator with the following new relations (using Pinkall’s notation): $`S\mathrm{\#}U=T\mathrm{\#}U`$ and $`B\mathrm{\#}U=\overline{B}\mathrm{\#}U`$ .
Note that $`\stackrel{~}{H}H`$ is a sub-semigroup of $`\stackrel{~}{H}`$ and it easily follows from Theorem 2.6 that it is isomorphic to $`J_+`$, where $`J`$ denotes the semigroup of (closed connected) surfaces. As a corollary we may conclude that the Grothendieck group of $`\stackrel{~}{H}`$ is isomorphic to $``$.
## 3. Proof of the main result
The purpose of this section is to prove Theorem 2.6. First we recall the classification of immersions of an arbitrary two-dimensional manifold $`F^2`$ into $`^3`$ using Hirsch-Smale theory.
###### Theorem 3.1.
There is a 1-1 correspondence between the regular homotopy classes of immersions of a surface $`F^2`$ into $`^3`$ denoted by $`\text{Imm}(F^2,^3)`$ and $`H^1(F^2;_2)`$.
###### Proof.
By Hirsch (see ) there is a weak homotopy equivalence between the space $`\text{Imm}(F^2,^3)`$ and $`\mathrm{\Gamma }(\mu )`$. Here $`\mathrm{\Gamma }(\mu )`$ denotes the space of sections of the vector bundle $`\mu =\text{MONO}(TF^2,R^3)`$ over $`F^2`$ whose fiber over $`pF^2`$ consists of all linear injections from $`T_pF`$ to $`^3`$. Thus there is a bijection between the regular homotopy classes $`\pi _0(\text{Imm}(F^2,^3))`$ and $`\pi _0(\mathrm{\Gamma }(\mu ))`$. Fix an arbitrary Riemannian metric on $`F^2`$ and let $`\mu ^{}`$ be the bundle over $`F^2`$ whose general fiber over $`p`$ is the space of orthogonal injections of $`T_pF`$ into $`^3`$. Then the inclusion of $`\mu ^{}`$ into $`\mu `$ is a fiber homotopy equivalence (see \[4, p. 426\]), thus $`\pi _0(\mathrm{\Gamma }(\mu ))=\pi _0(\mathrm{\Gamma }(\mu ^{}))`$. Fixing a section $`s\mathrm{\Gamma }(\mu ^{})`$ every section $`t\mathrm{\Gamma }(\mu ^{})`$ can be obtained by the action of a unique element of $`C(F^2,SO(3))`$ on $`s`$. Thus $`\mathrm{\Gamma }(\mu ^{})`$ is homeomorphic to $`C(F^2,SO(3))`$, yielding $`\pi _0(\mathrm{\Gamma }(\mu ))=[F^2,SO(3)]`$. Since $`SO(3)`$ is homeomorphic to $`RP^3`$, it follows from obstruction theory that
$$[F^2,SO(3)]=[F^2,RP^3]=[F^2,RP^{\mathrm{}}]=[F^2,K(_2,1)]=H^1(F^2;_2),$$
where $`K(_2,1)=RP^{\mathrm{}}`$ is an Eilenberg-MacLane space. ∎
From now on $`M^2`$ denotes the closed connected surface mentioned in the statement of Theorem 2.6. If $`f`$ and $`g`$ are locally generic mappings of $`M^2`$ into $`^3`$ with $`\left|S(f)\right|=\left|S(g)\right|`$, then according to the lemma of homogeneity there exists a diffeotopy $`\{\phi _t:t[0,1]\}`$ of $`M^2`$ such that $`\phi _0=id_{M^2}`$ and $`\phi _1(S(g))=S(f)`$. Since $`f\phi _t`$ provides a regular homotopy between $`f`$ and $`f\phi _1`$, it is sufficient to prove Theorem 2.6 in the case $`S(f)=S(g)`$. Let $`S(f)=S(g)=\{p_1,\mathrm{},p_k\}`$, where $`k=\left|S(f)\right|>0`$ is an even integer. For each $`p_i`$ choose a sufficiently small open neighborhood $`D_i`$ diffeomorphic to an open 2-disc such that $`f|D_i`$ has the canonical form $`(x_1^2,x_2,x_1x_2)`$ in an appropriate pair of local coordinate-systems centered at $`p_i`$ and $`f(p_i)`$. Similarly $`g|D_i`$ should have the same canonical form in another pair of local coordinate-systems. Assume moreover that the discs $`D_1,\mathrm{},D_k`$ are pairwise disjoint. Denote by $`A`$ the disjoint union $`_{i=1}^kD_i`$, then $`F^2=M^2A`$ is a two-manifold with boundary.
###### Lemma 3.2.
Suppose that $`f`$ and $`g`$ are locally generic mappings of the surface $`M^2`$ into $`^3`$ such that $`S(f)=S(g)`$. Choose open discs $`D_1,\mathrm{},D_kM^2`$ centered at the points of $`S(f)`$ as above. Define $`F^2=M^2_{i=1}^kD_i`$. Then there exists a diffeomorphism $`d`$ of the pair $`(M^2,F^2)`$ such that the immersions $`(f|F^2)d`$ and $`g|F^2`$ are regularly homotopic and $`d`$ permutes $`S(f)`$. Moreover there is a diffeotopy $`d_t`$ of $`M^2`$ with $`d_0=id_{M^2}`$ and $`d_1=d`$.
###### Proof.
According to Theorem 3.1 the regular homotopy classes of $`f|F^2`$ and $`g|F^2`$ correspond to cohomology classes $`\alpha ,\beta H^1(F^2;_2)`$. (These will be shown to be non-zero later.) We construct a diffeomorphism $`d`$ of the pair $`(M^2,F^2)`$ such that for the induced automorphism $`d^{}`$ of $`H^1(F^2;_2)`$ it holds that $`d^{}(\alpha )=\beta `$. Using Theorem 3.1 again this gives the required result $`(f|F^2)dg|F^2`$.
We first note that
(3.1)
$$H_1(F^2;_2)=H_1(M^2;_2)\stackrel{k1}{\stackrel{}{_2\mathrm{}_2}},$$
as can be seen from the exact sequence of the pair $`(M^2,F^2)`$. (Recall that $`k=\left|S(f)\right|`$.) For each $`i`$, $`1ik`$ choose an embedded curve $`\gamma _i`$ in $`F^2`$ around $`D_i`$. Denote its homology class by $`[\gamma _i]=c_i`$. The classes $`c_1,\mathrm{},c_{k1}`$ can be chosen for the generators of the $`(k1)`$ $`_2`$ summands in 3.1. (Note that $`c_1+\mathrm{}+c_k=0`$ since $`F^2=\gamma _1\mathrm{}\gamma _k.`$) According to Pinkall we have that $`\alpha ,c_i=1`$ and $`\beta ,c_i=1`$ for every $`i\{\mathrm{\hspace{0.17em}1},\mathrm{},k\}`$, since $`\gamma _i`$ has a neighborhood homeomorphic to $`S^1\times [0,1]`$ which is mapped by $`f`$ and also by $`g`$ into a ”figure eight$`\times [0,1]`$”. Now we have two cases according to the orientability of $`M^2`$.
If $`M^2`$ is orientable of genus $`g`$ we denote the standard generators of $`H_1(M^2;_2)`$ by $`a_1,b_1,\mathrm{},a_g,b_g`$, and choose embedded curves $`\phi _1,\psi _1\mathrm{}\phi _g,\psi _g`$ in $`F^2`$ representing them. Define
$$H_a=\{\mathrm{\hspace{0.17em}1}ig:\alpha ,a_i\beta ,a_i\},$$
and
$$H_b=\{\mathrm{\hspace{0.17em}1}ig:\alpha ,b_i\beta ,b_i\}.$$
There is a simple (i.e. embedded) closed curve $`\delta `$ in $`F^2`$ that for each $`i`$, $`1ig`$ intersects transversally in one point the curve $`\phi _i`$ if $`iH_a`$ and is disjoint from $`\phi _i`$ if $`iH_a`$, moreover $`\delta `$ intersects transversally in one point $`\psi _i`$ if $`iH_b`$ and is disjoint from $`\psi _i`$ if $`iH_b`$. Note that the homology class of $`\delta `$ will be
(3.2)
$$\left(\underset{iH_a}{}b_i\right)+\left(\underset{iH_b}{}a_i\right).$$
Such a $`\delta `$ exists because in $`H_1(M^2;_2)`$ any class can be represented by a simple curve, and a simple curve representing the class 3.2 can be arranged to be transversal to all the curves $`\phi _j`$ and $`\psi _j`$ for $`j=1,\mathrm{},k`$ and to intersect each of them at most in one point. Now choose two points on $`\delta `$ very close to each other so that none of the curves $`\phi _i`$ and $`\psi _i`$ for $`1ig`$ intersects the shorter arc $`\delta ^{}`$ between them. Thus the following equalities hold :
(3.3)
$$\left|\delta \phi _i\right|=\{\begin{array}{cc}1\hfill & \text{if }iH_a\hfill \\ 0\hfill & \text{if }iH_a\hfill \end{array}$$
$$\left|\delta \psi _i\right|=\{\begin{array}{cc}1\hfill & \text{for }iH_b\hfill \\ 0\hfill & \text{for }iH_b\hfill \end{array}$$
$`\delta ^{}\phi _i=\delta ^{}\psi _i=\mathrm{}`$ for $`1ig`$.
Modify the arc $`\delta ^{}`$ ( by a homology ) so that it goes through the center of $`D_1`$ and $`D_2`$ avoiding the curves $`\phi _i`$ and $`\psi _i`$ for $`i=1,\mathrm{},g`$ as well as the discs $`D_i`$ for $`i=3,\mathrm{},k`$. This can be done since $`M^2_{i=1}^{2g}(\phi _i\psi _i)`$ is path-connected. From now on we will denote this new simple curve on $`M^2`$ by $`\delta `$ (see Figure 2). Note that the equalities 3.3 still hold. Now choose a tubular neighborhood $`T`$ of $`\delta `$ such that $`D_1,D_2T`$ and for $`iH_a`$ the curve $`\phi _i`$ and for $`jH_b`$ the curve $`\psi _j`$ intersects $`T`$ in a line segment. (See the left side of Figure 3.) We also select a slightly wider tubular neighborhood $`T^{}T`$. Define $`d`$ on $`T`$ to be a rotation of $`T=S^1\times [1,1]`$ by $`180^{}`$ interchanging $`D_1`$ and $`D_2`$ and also interchanging $`p_1`$ and $`p_2`$. The diffeomorphism $`d`$ acts identically on $`M^2T^{}`$. On $`T^{}T`$, which is homeomorphic to $`S^1\times ([2,2][1,1])`$, define $`d`$ as the rotation of $`S^1\times \{s\}`$ by $`(2\left|s\right|)\times 180^{}`$ for $`s[2,2][1,1]`$ (see the right side of Figure 3). The diffeomorphism $`d`$ is diffeotopic to $`id_{M^2}`$ : construct $`d_t:M^2M^2`$ similarly to $`d`$, just take a rotation by $`180^{}t`$ instead of rotating by $`180^{}`$ such that $`d_t(p_1)\delta ^{}`$ for every $`t[0,1]`$.
For $`iH_a`$ it is clear that $`d`$ is the identity on the image of $`\phi _i`$, thus $`d_{}(a_i)=a_i`$. On the other hand if $`iH_a`$, then $`d\phi _i`$ is homologous to the connected sum of $`\phi _i`$ and $`\gamma _1`$ (surrounding $`D_1`$), thus $`d_{}(a_i)=a_i+c_1`$. Similarly for $`iH_b`$ we have $`d_{}(b_i)=b_i`$, and if $`iH_b`$ then $`d_{}(b_i)=b_i+c_1`$. Finally it holds that $`d_{}(c_2)=c_1`$, $`d_{}(c_1)=c_2`$ and if $`i>2`$ then $`d_{}(c_i)=c_i`$. (If $`k=2`$ then $`d_{}(c_1)=c_2=c_1`$.) This can be verified by looking at the action of $`d`$ on the curves $`\gamma _1,\mathrm{},\gamma _k`$. Since $`d_{}`$ permutes the generators $`c_1,\mathrm{}c_{k1}`$, for $`1ik1`$ we have $`\alpha ,d_{}(c_i)=1`$. (Recall that $`\alpha ,c_i=1`$ and $`\beta ,c_i=1`$ for every $`i`$.) Thus $`\alpha ,d_{}(c_i)=\beta ,c_i`$. By the choice of $`H_a`$ we see that for $`iH_a`$ it holds that $`\alpha ,d_{}(a_i)=\alpha ,a_i+c_1=\alpha ,a_i+1=\beta ,a_i`$ and for $`iH_a`$ we have that $`\alpha ,d_{}(a_i)=\alpha ,a_i=\beta ,a_i`$. A similar argument holds for $`b_1,\mathrm{},b_g`$. Since $`a_1,b_1,\mathrm{},a_g,b_g`$ and $`c_1,\mathrm{},c_{k1}`$ form a basis of $`H_1(F^2;_2)`$, we have shown that $`d^{}(\alpha ),x=\alpha ,d_{}(x)=\beta ,x`$ for every $`xH_1(F^2;_2)`$. Thus $`d^{}(\alpha )=\beta `$. Hence $`(fd)|F^2`$ is regularly homotopic to $`g|F^2`$ as required. Also $`d`$ satisfies $`d(S(f))=S(f)`$.
Now suppose that $`M^2`$ is a non-orientable surface of genus $`g`$, i.e. a sphere with $`g`$ Moebius bands. Choose a curve $`\phi _i`$ in $`F^2`$ on the $`i`$-th Moebius band representing its homology generator for $`1ig`$. Then $`b_1=[\phi _1],\mathrm{},b_g=[\phi _g]`$ together with $`c_1,\mathrm{},c_{k1}`$ is the standard basis of $`H_1(F^2;_2)`$. Analogously to the orientable case let
$$H=\{\mathrm{\hspace{0.17em}1}ig:\alpha ,b_i\beta ,b_i\},$$
and similarly it is sufficient to construct a diffeomorphism $`d`$ of the pair $`(M^2,F^2)`$ with $`d_{}(c_i)=c_i`$ for $`1ik`$, $`d_{}(b_i)=b_i+c_1`$ for $`iH`$ and $`d_{}(b_i)=b_i`$ for $`iH`$. It is enough to show that for any fix $`1jg`$ we can find a diffeomorphism $`d_j`$ of the pair $`(M^2,F^2)`$ such that $`d_j(b_j)=b_j+c_1`$ and that $`d_j`$ is identical on every other homology-generator. (Then $`_{jH}d_j`$ is a good choice for $`d`$.)
For this end modify a small arc of $`\phi _j`$ using a homology such that it still lies in $`F^2`$ but gets close to $`D_1`$ and remains disjoint from all the other $`\phi _i`$ for $`ij`$. (We shall call this modified curve $`\phi _j`$ also.) This is possible since $`M^2_{i=1}^g\phi _i`$ is path-connected. Denote by $`T`$ a tubular neighborhood of $`\phi _j`$ in $`M^2`$ containing $`D_1`$ and disjoint from $`D_i`$ if $`i>1`$ and from $`\phi _i`$ if $`ij`$. (See Figure 4.) Then $`T`$ is homeomorphic to the Moebius band. Also choose a slightly larger tubular neighborhood $`T^{}`$ of $`T`$ with similar properties. Now think of $`T`$ as a rectangle with the vertical sides identified in the opposite direction and with $`D_1`$ in its center. Let $`d_j`$ be the reflection of the rectangle $`T`$ into its horizontal central axis going through the center of $`D_1`$ which is $`p_1`$. Then $`d_j`$ induces an orientation-preserving diffeomorphism ( a rotation ) of $`T=S^1`$, which can be extended to $`T^{}T=S^1\times [0,1]`$ being identical on $`T^{}=S^1\times \{1\}`$ as we have already seen. Finally $`d_j`$ is identical on $`M^2T^{}`$. This $`d_j`$ maps the curve $`\phi _j`$ (which is the horizontal central line in the rectangle except that it avoids $`D_1`$ (see Figure 4)) to a curve homologous to the connected sum of $`\phi _j`$ and $`\gamma _1`$, thus $`d_j(b_j)=b_j+c_1`$. Since $`\phi _i`$ for $`ij`$ and $`\gamma _i`$ for $`i>1`$ are fixed by $`d_j`$, it satisfies the required conditions. Concerning $`c_1`$ we have $`d_j(c_1)=c_1=c_1`$ since we are working with mod $`2`$ coefficients.
The diffeomorphism $`d_j`$ of $`M^2`$ is diffeotopic to $`id_{M^2}`$: Think of the Moebius band $`T`$ as the factor space $`S^1\times [0,1]/_{p\times \{1\}=(p)\times \{1\}}`$ (thus we identify the opposite points of one boundary component of an annulus). In this model define $`(d_j)_t`$ on $`T`$ as the diffeomorphism induced by the rotation of the annulus by $`180^{}t`$ degrees. On $`TT^{}`$ define $`(d_j)_t`$ as before. Finally on $`M^2T^{}`$ the diffeomorphism $`(d_j)_t`$ is the identity mapping. Then $`d_j=(d_j)_1`$ and $`(d_j)_0=id_{M^2}`$ as required. ∎
As a consequence of the above proof we obtain the following proposition (for the definition of the mapping class group see ):
###### Corollary 3.3.
Suppose that $`F^2`$ is a surface of genus $`g`$ with $`k>1`$ boundary components, where k is even. We denote, like as before, the homology classes represented by the boundary components of $`F^2`$ in $`H_1(F^2;_2)`$ by $`c_1,\mathrm{},c_k`$. The mapping class group $`M(F^2)`$ of $`F^2`$ acts on the set $`S=\{\alpha H^1(F^2,_2):\alpha ,c_i=1,\mathrm{\hspace{0.17em}1}ik1\}`$. (If $`\alpha S`$ then $`\alpha ,c_k=\alpha ,c_1+\mathrm{}+c_{k1}=1`$ since $`k`$ is even.) If $`M^2`$ is a closed surface of genus $`g`$ then there is a homomorphism $`m:M(F^2)M(M^2)`$ obtained by ”filling in the holes”. Then $`\mathrm{ker}(m)`$ acts transitively on $`S`$.
###### Lemma 3.4.
Let $`g`$ and $`h`$ be locally generic mappings of $`M^2`$ into $`^3`$ such that $`S(g)=S(h)`$ and $`g|F^2h|F^2`$, where $`F^2`$ is the complement of a small open neighborhood of the common singular set. Then $`gh`$.
###### Proof.
Recall that $`A=_{i=1}^kD_i`$. Since for each $`i`$ it holds that $`h|D_i`$ and $`g|D_i`$ have canonical forms in appropriate coordinate-systems, there is a regular homotopy $`H`$ between $`h`$ and a locally generic mapping $`\stackrel{~}{h}`$ such that for each $`t[0,1]`$ we have $`S(H_t)=S(h)`$ and that $`\stackrel{~}{h}|A=g|A`$. Thus $`H|(F^2\times [0,1])`$ is a regular homotopy between the immersions $`h|F^2`$ and $`\stackrel{~}{h}|F^2`$ showing that $`\stackrel{~}{h}|F^2g|F^2`$. So we can suppose that $`g|A=h|A`$.
Let $`H`$ be a regular homotopy between $`g|F^2`$ and $`h|F^2`$. For every $`i=1,\mathrm{},k`$ fix a smaller concentric closed disc $`B_i`$ in $`D_i`$ (hence $`p_iB_iD_i`$). Finally set $`F_l=F^2\left(_{i=1}^lD_i\right)`$. We will define recursively a sequence of regular homotopies $`H^l:F_l\times [0,1]^3`$ connecting $`g_l=g|F_l`$ and $`h_l=h|F_l`$ for $`l=0,\mathrm{},k`$ with the property $`H^0=H`$. Suppose that we have constructed $`H^l`$ for $`l<j`$. Let $`q`$ be a point in $`D_j`$. For $`t[0,1]`$ there is a one-parameter family of elements $`M_tGL(,3)`$ with $`M_td_qH_t^{j1}=d_qH_0^{j1}`$ and a vector $`v_t^3`$ with $`M_t(H_t^{j1}(q))+v_t=H_0^{j1}(q)`$. Here $`d_qH_t^{j1}`$ denotes the differential of the mapping $`H_t^{j1}`$ at the point $`q`$. Now define $`H_t^j|F_{j1}`$ to be equal to $`M_tH_t^{j1}+v_t`$. Since $`g(q)=h(q)`$ and $`d_qg=d_qh`$, we have that $`M_1=id_^3`$ and $`v_1=0`$, thus $`H_1^j|F_{j1}=h|F_{j1}`$. With this transformation of $`H^{j1}`$ we have achieved that $`H_t^j(q)=H_0^j(q)`$ and $`d_qH_t^j=d_qH_0^j`$ for every $`t[0,1]`$. Let $`U_q`$ be a small closed neighborhood of $`q`$ in $`F_{j1}`$ diffeomorphic to a closed 2-disc ($`qU_q`$). Using a standard argument of S. Smale we can suppose that the homotopy $`H^j`$ is kept fixed on a neighborhood of $`U_q`$ (see Hirsch , Lemma 2.5 on page 249). On $`B_j`$ define for every $`t[0,1]`$ the mapping $`H_t^j|B_j=H_0^j|B_j=g|B_j`$.
Denote the closure of the annulus $`D_jB_j`$ by $`T_j`$, and let $`I`$ be a radial line in $`T_j`$ containing $`q`$ (see Figure 5). A tubular neighborhood $`V_IT_j`$ of $`I`$ is obtained by taking $`V_I=(U_qD_j)\times I`$. Let $`H_t^j|V_I=g|V_I`$ for every $`t[0,1]`$. This is possible since $`H^j`$ is fix on $`U_q`$. We only have to define $`H^j`$ on the closed two-cell $`C_i=T_i\text{Int}(V_I)`$. For this purpose we will use Smale’s lemma (see Theorem 1.1 on page 245 of or Theorem 2.1 in ), which intuitively states the following: If we are given an immersed disk $`D^k`$ in $`^n`$ such that $`k<n`$ and we deform the boundary of the disk and the normal derivatives along the boundary, then we can deform the whole disk at the same time so as to induce the given deformation on the boundary and normal derivatives. Since $`H^j`$ is already defined on $`C_i`$ along with derivatives normal to $`C_i`$ we can use Smale’s lemma to get a regular homotopy $`G`$ on $`C_j`$ with $`G_0=g|C_j`$ and $`G`$ prescribed along the boundary. Finally $`G_1h|C_j`$ since they coincide on $`C_j`$ and the obstruction is an element of $`\pi _2(SO(3))=0`$. Putting this homotopy after $`G`$ we obtain the desired homotopy $`H^j|C_j`$. Thus we have constructed $`H^j`$ on the whole manifold $`F_j`$. This shows that $`H^k`$ is a regular homotopy connecting $`g`$ and $`h`$. ∎
###### Proof of Theorem 2.6.
We have seen using the lemma of homogeneity that it is sufficient to prove Theorem 2.6 under the assumption $`S(f)=S(g)`$. By Lemma 3.2 there exists a diffeomorphism $`d`$ such that $`(fd)|F^2g|F^2`$ and $`S(fd)=S(f)`$. Now applying Lemma 3.4 to $`h=fd`$ and $`g`$ we get $`fdg`$. But $`fd_t`$ is a regular homotopy between $`f`$ and $`fd`$ proving that $`fg`$.
It remains to show that the above regular homotopy $`H`$ joining $`f`$ and $`g`$ can be chosen in such a way that it defines a prescribed bijection $`i:S(f)S(g)`$, that is $`i_H=i`$. Recall that the bijection $`i_H:S(f)S(g)`$ depends only on the choice of the diffeotopy $`\phi _t`$ mentioned before Lemma 3.2, that clearly might induce any prescribed bijection between $`S(f)`$ and $`S(g)`$. If $`M^2`$ is orientable, then the diffeomorphism $`d:M^2M^2`$ of Lemma 3.2 swaps the singular points $`p_1`$ and $`p_2`$ (i.e. $`d(p_1)=p_2`$, $`d(p_2)=p_1`$ and $`d(p_i)=p_i`$ for $`i>2`$), and in the non-orientable case $`d(p_i)=p_i`$ for $`1ik`$. Finally the homotopy constructed in Lemma 3.4 between $`fd`$ and $`g`$ is a singularity fixing homotopy in the sense of Definition 3.7. This completes the proof of Theorem 2.6. ∎
The converse of Lemma 3.4 is true only in the following form:
###### Proposition 3.5.
Suppose that $`g`$ and $`h`$ are locally generic mappings of $`M^2`$ into $`^3`$ such that $`S(g)=S(h)`$. Denote by $`F^2`$ the complement of a small open neighborhood of the common singular set. Then $`gh`$ implies that $`\text{Im}(g|F^2)\text{Im}(h|F^2)`$.
###### Proof.
Let $`H`$ be a regular homotopy connecting $`g`$ and $`h`$. Then
$$S(H_t)=\{p_1(t),\mathrm{},p_k(t)\},$$
where $`p_i(t)`$ is a smooth curve in $`M^2`$. The lemma of homogeneity gives a diffeotopy $`\phi _t`$ of $`M^2`$ such that $`\phi _0=id_{M^2}`$ and $`\phi _t(S(H_0))=S(H_t)`$ for every $`t[0,1]`$. The homotopy $`G_t=H_t\phi _t`$ has the property that $`S(G_t)=S(G_0)`$ for $`t[0,1]`$ and connects $`g`$ with $`h\phi _1`$. Since $`\phi _1`$ permutes $`S(g)`$ (because $`\phi _1(S(g))=\phi _1(S(H_0))=S(H_1)=S(g)`$) we can choose $`\phi _1`$ to map $`F^2`$ onto itself. $`G_t|F^2`$ is a regular homotopy between the immersions $`g|F^2`$ and $`(h|F^2)(\phi _1|F^2)`$ which means by definition that $`\text{Im}(g|F^2)\text{Im}(h|F^2)`$. ∎
###### Remark 3.6.
Modify Definition 2.2 of regular homotopy the following way:
###### Definition 3.7.
Locally generic mappings $`f,g:M^2^3`$ are *regularly homotopic through a singularity fixing homotopy* – notation $`f_sg`$ – if $`S(f)=S(g)`$ and there exists a smooth mapping $`H:M^2\times [0,1]^3`$ such that $`H_0=f`$ and $`H_1=g`$ and for every $`t[0,1]`$ the mapping $`H_t`$ is locally generic with $`S(H_t)=S(f)`$. (That is the singular points are kept fixed.)
This gives a modification of the definition of image-homotopic maps:
###### Definition 3.8.
Locally generic mappings $`f,g:M^2^3`$ are *image homotopic through a singularity fixing homotopy* – notation $`\text{Im}(f)_s\text{Im}(g)`$ – if $`S(f)=S(g)`$ and there is a $`d:M^2M^2`$ diffeomorphism such that $`fd_sg`$. Note that $`d(S(f))=S(g)=S(f)`$ and $`d`$ can permute the points of $`S(f)`$.
Suppose that $`S(f)`$ or $`S(g)`$ is non-empty. The arguments above show that $`\text{Im}(f)_s\text{Im}(g)`$ if and only if $`\left|S(f)\right|=\left|S(g)\right|`$. To prove this we only have to use diffeomorphisms instead of diffeotopies since Lemma 3.4 remains true using the new definition.
On the other hand $`f_sg`$ implies that the immersions $`f|(M^2S(f))`$ and $`g|(M^2S(f))`$ are regularly homotopic. But there are locally generic mappings $`f,g:M^2^3`$ satisfying $`S(f)=S(g)`$ such that $`fg`$ but $`f|(M^2S(f))g|(M^2S(f))`$. Take for example $`M^2=RP^2`$ and choose two arbitrary points $`p,qRP^2`$. Denote $`RP^2\{p,q\}`$ by $`F^2`$. Using the notations of Lemma 3.2 we have that $`H_1(F^2;_2)=f_1,c_1`$. Define the cohomology classes $`\alpha ,\beta H^1(F^2;_2)`$ by the equalities $`\alpha (f_1)=0`$, $`\beta (f_1)=1`$ and $`\alpha (c_1)=\beta (c_1)=1`$. Then there exist locally generic mappings $`f,g:RP^2^3`$ satisfying $`S(f)=S(g)=\{p,q\}`$ such that $`f|F^2`$ and $`g|F^2`$ correspond to $`\alpha `$ and $`\beta `$ using the bijection of Theorem 3.1 : Denote by $`U`$ the locally generic mapping of $`S^2`$ to $`^3`$ with singular points $`p`$ and $`q`$ (this is unique up to singularity fixing homotopy). $`B`$ is the famous Boy surface, $`\overline{B}`$ is the mirror image of $`B`$ (see ). Then the connected sums $`f=B\mathrm{\#}U`$ and $`g=\overline{B}\mathrm{\#}U`$ satisfy the above conditions. Clearly $`f|F^2g|F^2`$, thus $`f_sg`$. This provides examples of locally generic mappings $`f`$ and $`g`$ such that $`\text{Im}(f)_s\text{Im}(g)`$, but $`f_sg`$.
## 4. Projections of regular homotopies
Suppose that $`M^2`$ is a closed connected surface. In the previous sections we examined the path-components of the space of locally generic mappings of $`M^2`$ into $`^3`$ endowed with the $`C^{\mathrm{}}`$ topology. This space, denoted by $`\text{Lgen}(M^2,^3)`$, is closely connected with the space $`\text{Imm}(M^2,^4)`$ of immersions of $`M^2`$ into $`^4`$ (also considered with the $`C^{\mathrm{}}`$ topology). This connection is realized by projections of $`^4`$ onto $`^3`$.
###### Definition 4.1.
A mapping $`\pi :^k^l`$ $`(k>l)`$ is a projection if it is linear and surjective. Denote the space of all projections from $`^k`$ to $`^l`$ by $`\text{Proj}(^k,^l)`$.
Our starting point is the following result of Mather :
###### Proposition 4.2.
Suppose that $`F\text{Imm}(M^2,^4)`$ is an immersion. Then for almost every $`\pi \text{Proj}(^4,^3)`$ (in the sense of Lebesgue measure) the mapping $`\pi F`$ is locally generic.
For every projection $`\pi \text{Proj}(^4,^3)`$ let
$$\text{Imm}_\pi (M^2,^4)=\{F\text{Imm}(M^2,^4):\pi F\text{Lgen}(M^2,^3)\}$$
be the subspace of $`\text{Imm}(M^2,^4)`$. There is a natural mapping
$$p_\pi :\text{Imm}_\pi (M^2,^4)\text{Lgen}(M^2,^3)$$
defined by the formula $`p_\pi (F)=\pi F`$ for every $`F\text{Imm}_\pi (M^2,^4)`$.
In this section our aim is to examine the path-components of $`\text{Imm}_\pi (M^2,^4)`$.
###### Definition 4.3.
Two immersions $`F,G\text{Imm}_\pi (M^2,^4)`$ are called *$`\pi `$-homotopic* (denoted by $`F_\pi G`$) if they are in the same path-component of $`\text{Imm}_\pi (M^2,^4)`$.
First let us recall that two immersions $`F,G\text{Imm}(M^2,^4)`$ are regularly homotopic if and only if $`e(F)=e(G)`$, where $`e(F)`$ denotes the (twisted) Euler-number of the normal bundle of the immersion $`F`$. It is clear that if $`F_\pi G`$ then $`e(F)=e(G)`$ and $`\left|S(\pi F)\right|=\left|S(\pi G)\right|`$. We are going to prove that if $`S(\pi F)`$ (and $`S(\pi G)`$) are non-empty then the converse also holds.
Suppose that $`F,G:M^2^4`$ are immersions. From Proposition 4.2 it is clear that for almost every projection $`\pi \text{Proj}(^4,^3)`$ both $`f=\pi F`$ and $`g=\pi G`$ are locally generic. Fix such a projection $`\pi `$.
We are going to define the sign of every point in $`S(f)`$ (and in $`S(g)`$). (An equivalent definition can be found in ). Take a cross-cap point $`pS(f)`$. Choose orientations of $`^4`$ and of $`^3`$. We will define the sign of $`p`$ as follows.
Fix local coordinates $`(x_1,x_2)`$ on a neighborhood $`U_p`$ of $`p`$ and $`(y_1,y_2,y_3)`$ centered at $`f(p)`$ such that $`f`$ has the following normal form:
$$y_1f=x_1^2,y_2f=x_2,y_3f=x_1x_2.$$
Suppose that $`U_p`$ is so small that $`F|U_p`$ is an embedding. The sign of $`p`$ will depend only on $`F|U_p`$ (thus the definition is local). Set $`D_\epsilon =\{y_1^2+y_2^2+y_3^2\epsilon \}^3`$ and choose $`\epsilon >0`$ sufficiently small such that $`\stackrel{~}{D}_\epsilon =f^1(D_\epsilon )U_p`$. The set $`\stackrel{~}{D}_\epsilon `$ is a closed disc neighborhood of $`p`$ in $`M^2`$ and $`\stackrel{~}{D}_\epsilon =f^1(D_\epsilon )`$ (see Lemma 2.2 in ). Let $`L`$ be the closure of the double point set of $`f|\stackrel{~}{D}_\epsilon ,i.e.`$ $`L=\{x\stackrel{~}{D}_\epsilon :f^1(f(x))\stackrel{~}{D}_\epsilon \{x\}\}\{p\}`$. Then $`L`$ is a one-dimensional smooth submanifold of $`\stackrel{~}{D}_\epsilon `$ and $`L\stackrel{~}{D}_\epsilon `$ consists of two points $`p_1`$ and $`p_2`$. We fix an orientation of $`\stackrel{~}{D}_\epsilon `$ and take an oriented base $`(u)`$ (resp. $`(v)`$) of the tangent space $`T_{p_1}(\stackrel{~}{D}_\epsilon )`$ (resp. $`T_{p_2}(\stackrel{~}{D}_\epsilon )`$). Then $`(df_{p_1}(u),df_{p_2}(v))`$ is a base of $`T_q(D_\epsilon )`$, where $`q=f(p_1)=f(p_2)`$. We may assume that $`(df_{p_1}(u),df_{p_2}(v),\xi )`$ is a positive basis of $`T_q^3`$, where $`\xi `$ is the outward normal vector of $`D_\epsilon `$, exchanging $`p_1`$ and $`p_2`$ if necessary. Now orient $`L`$ from $`p_2`$ to $`p_1`$.
Denote by $`\nu `$ a positive basis vector of $`T_pL`$. (If $`M^2`$ possess a Riemannian metric, then choose $`\nu `$ to be a unit-vector. This way $`\nu `$ is unique up to the orientation of $`^3`$.) Orient $`\mathrm{ker}\pi `$ in such a way that together with the orientation of $`\text{Im}\pi =^3`$ we obtain the fixed orientation of $`^4=\text{Im}\pi \mathrm{ker}\pi `$. Using this direct sum decomposition of $`^4`$ the mapping $`F:M^2^4`$ can be written in the form $`F=(\pi F,F^{})=(f,F^{})`$. Since $`F`$ is an immersion at the point $`p`$ and $`pS(f)`$ , we have $`dF^{}(\nu )0`$. After all this preparation we can now define the sign of $`p`$.
###### Definition 4.4.
The cross-cap point $`p`$ is *positive* if $`dF^{}(\nu )>0`$ and is *negative* if $`dF^{}(\nu )<0`$. We denote by $`p(F)`$ (resp. $`n(F)`$) the number of positive (resp. negative) cross-cap singularities of the locally generic mapping $`f=\pi F:M^2^3`$.
The following proposition is a special case of Proposition 2.5 in .
###### Proposition 4.5.
Suppose that $`^4`$ and $`^3`$ are oriented. Then we always have
$$e(F)=p(F)n(F),$$
where $`e(F)`$ is the (twisted) Euler-number of the normal bundle of the immersion $`F`$.
Thus the immersions $`F`$ and $`G`$ are regularly homotopic if and only if
$$p(F)n(F)=e(F)=e(G)=p(G)n(G).$$
On the other hand Theorem 2.6 states that if $`S(f)`$ and $`S(g)`$ are non-empty then $`fg`$ if and only if
$$p(F)+n(F)=\left|S(f)\right|=\left|S(g)\right|=p(G)+n(G).$$
Comparing the preceding two chains of equations we have that if both $`f`$ and $`g`$ are singular then $`[FG`$ and $`fg][p(F)=p(G)`$ and $`n(F)=n(G)]`$. The following theorem implies that in this case we can even find a regular homotopy between $`F`$ and $`G`$ whose projection is a regular homotopy between $`f`$ and $`g`$, i.e. $`F_\pi G`$.
###### Theorem 4.6.
Suppose that $`M^2`$ is a closed connected surface, $`F,G:M^2^4`$ are immersions and $`\pi :^4^3`$ is a projection such that $`f=\pi F`$ and $`g=\pi G`$ are both locally generic and singular. Then the following are equivalent:
$`(1)`$ There exists a regular homotopy $`H:M^2\times [0,1]^4`$ between $`F`$ and $`G`$ such that $`\stackrel{~}{H}=\pi H`$ is a regular homotopy between $`f`$ and $`g`$, i.e. $`F_\pi G`$.
$`(2)`$ The numbers of positive and negative cross-caps of $`f`$ and $`g`$ are the same, i.e. $`p(F)=p(G)`$ and $`n(F)=n(G)`$.
###### Proof.
First we prove the implication (1) $``$ (2). In this case $`fg`$, thus using Proposition 2.4 we have that $`\left|S(f)\right|=\left|S(g)\right|`$. From Definition 4.4 it is clear that the signs of the singular points do not change during a regular homotopy. (Here we did not use the assumption that $`\left|S(f)\right|>0,\left|S(g)\right|>0`$.)
Now we are going to prove the implication (2) $``$ (1). Since $`p(F)=p(G)`$ and $`n(F)=n(G)`$ there exists a bijection $`i:S(f)S(g)`$ that preserves the signs of the cross-cap points. By Theorem 2.6 there is a regular homotopy $`\stackrel{~}{H}`$ between $`f`$ and $`g`$ such that $`i_{\stackrel{~}{H}}=i`$. We shall construct a regular homotopy $`H`$ between the immersions $`F`$ and $`G`$ such that $`\pi H=\stackrel{~}{H}`$ as follows: Choose an arbitrary Riemannian metric on $`M^2`$. In the paragraph preceding Definition 4.4 we saw that in this case for every $`t[0,1]`$ and $`pS(\stackrel{~}{H}_t)`$ there is a unique positive unit-vector $`\nu _t(p)T_pM^2`$ tangent to the double-point curve $`L`$ crossing $`p`$. The singular sets $`S(\stackrel{~}{H}_t)`$ for $`t[0,1]`$ define curves $`\gamma _1,\mathrm{},\gamma _k`$ on $`M^2`$ such that $`S(H_t)=\{\gamma _j(t):1jk\}`$ for every $`t[0,1]`$. The points $`\gamma _i(0)`$ and $`\gamma _i(1)`$ have the same sign ($`1ik`$). Suppose for example that for a fix $`i`$ both $`\gamma _i(0)`$ and $`\gamma _i(1)`$ are positive cross-cap points. Introduce the notation $`\nu _i(t)=\nu _t(\gamma _i(t))`$, then $`dF^{}(\nu _i(0))>0`$ and $`dG^{}(\nu _i(1))>0`$. (Here $`F^{}`$ denotes the fourth coordinate function of $`F`$ in $`^4`$.) Using the Levi-Civita connection of the Riemannian manifold $`M^2`$ we may consider the exponential mapping on $`M^2`$. Since $`[0,1]`$ is compact, there exists $`\epsilon >0`$ such that for every $`t[0,1]`$ the mapping $`h:[\epsilon ,\epsilon ]\times [0,1]M^2`$ satisfying
$$h(s,t)=\mathrm{exp}_{\gamma _i(t)}(s\nu _i(t))$$
is defined and for every $`t[0,1]`$ the mapping $`h_t(s)=h(s,t)`$ is an embedding of $`[\epsilon ,\epsilon ]`$ into $`M^2`$ ($`h_t`$ is a geodetic curve). Define the function $`H_t^{}`$ on $`\text{Im}h_t`$ using the following formula:
$$H_t^{}(h_t(s))=(1t)F^{}(h_0(s))+tG^{}(h_1(s))$$
for $`s[\epsilon ,\epsilon ]`$. Note that $`H_0^{}|\text{Im}h_0=F^{}|\text{Im}h_0`$ and $`H_1^{}|\text{Im}h_1=G^{}|\text{Im}h_1`$. Thus we can extend $`H^{}`$ to an open neighborhood of $`_{t[0,1]}\left(\text{Im}h_t\times \{t\}\right)`$ in $`M^2\times [0,1]`$ as a smooth function. From the construction of $`H_t^{}`$ it is clear that
$$dH_t^{}(\nu _i(t))=(1t)dF^{}(\nu _i(0))+tdG^{}(\nu _i(1))>0,$$
which implies that the mapping $`H_t=(\stackrel{~}{H}_t,H_t^{})`$ is an immersion at the point $`\gamma _i(t)`$ for every $`t[0,1]`$. Repeat the preceding extension process for every $`1ik`$ and afterwards extend the obtained $`H^{}`$ to the whole cylinder $`M^2\times [0,1]`$ in such a way that $`H_0^{}=F^{}`$ and $`H_1^{}=G^{}`$. The mapping $`H=(\stackrel{~}{H},H^{})`$ is a regular homotopy connecting $`F`$ and $`G`$ whose projection is $`\stackrel{~}{H}`$. ∎
Putting together our previous results we obtain the following theorem:
###### Theorem 4.7.
If $`F,G\text{Imm}_\pi (M^2,^4)`$ then
$$F_\pi G[FG\text{and}\pi F\pi G].$$
###### Proof.
First we suppose that $`S(\pi F)`$ or $`S(\pi G)`$ is non-empty. Theorem 4.6 states that $`F_\pi G[p(F)=p(G)`$ and $`n(F)=n(G)]`$. We have seen in the paragraph preceding Theorem 4.6 that $`[p(F)=p(G)`$ and $`n(F)=n(G)][FG`$ and $`\pi F\pi G]`$.
Now we consider the case when both $`S(\pi F)`$ and $`S(\pi G)`$ are empty, i.e. $`f=\pi F`$ and $`g=\pi G`$ are immersions. If $`fg`$ then any regular homotopy connecting $`f`$ and $`g`$ can be lifted to a regular homotopy between $`F`$ and $`G`$, thus $`F_\pi G`$. This proves the implication $`F_\pi G[FG`$ and $`\pi F\pi G]`$. The other implication is trivial. ∎
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# Towards a characterization of Markov processes enjoying the time-inversion property
## 1 Introduction
Let $`\{(X_t,t0);(\mathrm{P}_x)_{x^n}\}`$ be a homogeneous Markov process with semigroup densities (assumed to exist and to be twice differentiable in the space and time variables):
$$P_t(x,dy)=p_t(x,y)dy.$$
(1.1)
For all $`x^n`$ and some $`\alpha >0`$, the process $`\{(t^\alpha X_{\frac{1}{t}},t>0);\mathrm{P}_x\}`$ is Markov and in general inhomogeneous.
###### Definition 1.1.
The process $`\{(X_t,t0);\mathrm{P}_x\}`$ is said to enjoy the time-inversion property of degree $`\alpha `$ if the Markov process $`\{(t^\alpha X_{\frac{1}{t}},t>0);\mathrm{P}_x\}`$ is homogeneous.
The problems associated to time-inversion of Markov processes are closely related to the so-called dual processes in probabilistic potential theory (see and ). A pair of dual processes is a pair of Markov processes whose resolvents are conjugate relative to some measure. The trajectories of these processes are connected by means of time-reversal. However, time-reversal in general violates homogeneity and changes the path properties of the process. In this paper, we are mainly concerned with the first issue, and in particular, whether one can characterize the semigroup densities of a Markov process such that homogeneity is preserved under time-inversion.
Celebrated examples of Markov processes, known to enjoy the time-inversion property for $`\alpha =1`$, are Brownian motions with drift in $`^n`$ and Bessel processes with drift (see and ). Gallardo and Yor recently worked out a sufficient condition on the semigroup densities for a Markov process to enjoy the time-inversion property. Their argument extended the class of processes to processes with jumps such as the Dunkl process and matrix-valued processes such as the Wishart process . The aim of the paper is to find a necessary and sufficient condition and to provide some new examples.
Section 2 contains the main theorem of the paper, which is proved in section 3 using straightforward analytical arguments. Section 4 considers an application of the theorem to Markov processes on $`_+`$. The result is shown to be strong enough to entirely characterize the class of diffusion processes on $`_+`$ that enjoy the time-inversion property and thus provide a different proof than that of Watanabe in . Section 5 gives new examples of processes that enjoy the time-inversion property. We review how the generalized Dunkl process fits the requirements of the theorem and then introduce a matrix-valued process with jumps. The relation of the latter process to the Wishart process mimics the relation between the one-dimensional Dunkl and Bessel processes.
## 2 Markov processes which enjoy time-inversion
Fix $`x^n`$. Recall that under $`\mathrm{P}_x`$, the process $`(t^\alpha X_{\frac{1}{t}},t>0)`$ is Markov, inhomogeneous, with transitional probability densities $`q_{s,t}^{(x)}(z,y),(s<t;z,y^n)`$, which satisfy the following relation:
$$E_x\left[f(t^\alpha X_{\frac{1}{t}})|s^\alpha X_{\frac{1}{s}}=z\right]=𝑑yf(y)q_{s,t}^{(x)}(z,y)$$
(2.1)
where
$$q_{s,t}^{(x)}(a,b)=t^{n\alpha }\frac{p_{\frac{1}{t}}(x,\frac{b}{t^\alpha })p_{\frac{1}{s}\frac{1}{t}}(\frac{b}{t^\alpha },\frac{a}{s^\alpha })}{p_{\frac{1}{s}}(x,\frac{a}{s^\alpha })}.$$
(2.2)
The process $`(t^\alpha X_{\frac{1}{t}},t>0)`$ is not uniquely defined from the knowledge of the semigroup densities $`p_t(x,y)`$. Actually, there exists at least one transformation that leaves the semigroup densities $`q_{s,t}^{(x)}(a,b)`$ unchanged: Doob’s $`h`$-transform.
###### Definition 2.1.
Doob’s $`h`$-transform is the transformation
$$T_h:\mathrm{P}_x|__t\frac{h(X_t)}{h(x)}e^{\nu t}\mathrm{P}_x|__t.$$
(2.3)
for some function $`h`$ and some constant $`\nu >0`$.
This remark leads to our first assertion.
###### Proposition 2.2.
Two processes related by $`h`$-transforms yield the same process by time-inversion.
It is hence legitimate to research for a criterium to classify all processes that enjoy the time-inversion property up to $`h`$-transforms. The following theorem gives a concise statement of our main result:
###### Theorem 2.3.
The Markov process $`(t^\alpha X_{\frac{1}{t}},t>0)`$ is homogeneous if and only if the semigroup densities of $`(X_t,t0)`$, assumed twice differentiable, are of the form:
$$p_t(x,y)=t^{\frac{n\alpha }{2}}\mathrm{\Phi }(\frac{x}{t^{\frac{\alpha }{2}}},\frac{y}{t^{\frac{\alpha }{2}}})\theta \left(\frac{y}{t^{\frac{\alpha }{2}}}\right)\mathrm{exp}\left\{\rho \left(\frac{x}{t^{\frac{\alpha }{2}}}\right)+\rho \left(\frac{y}{t^{\frac{\alpha }{2}}}\right)\right\},$$
(2.4)
or if they are in $`h`$-transform relationship with it. The functions $`\mathrm{\Phi },\theta ,\rho `$ have the following properties for $`\lambda >0`$:
1. $`\mathrm{\Phi }(\lambda x,y)=\mathrm{\Phi }(x,\lambda y)`$;
2. $`\rho (\lambda x)=\lambda ^{\frac{2}{\alpha }}\rho (x)`$;
3. $`\theta (\lambda y)=\lambda ^\beta \theta (y)`$ for some $`\beta `$;
Moreover, if the symmetry condition $`\mathrm{\Phi }(x,y)=\mathrm{\Phi }(y,x)`$ is satisfied, then the semigroup densities are related as follows:
$$q_t^{(x)}(a,b)=\frac{\mathrm{\Phi }(x,b)}{\mathrm{\Phi }(x,a)}\mathrm{exp}\left\{t\rho \left(x\right)\right\}p_t(a,b).$$
(2.5)
Such a decomposition of the semigroup densities shows furthermore that $`h`$-transforms are the only transformations that leave semigroup densities of the time-inverted process $`(t^\alpha X_{\frac{1}{t}},t>0)`$ unchanged.
###### Definition 2.4.
A Markov process $`(X_t,t0)`$ is called semi-stable with index $`\gamma `$ in the sense of Lamperti if:
$$\{(X_{ct},t0);\mathrm{P}_x\}\stackrel{(d)}{=}\{(c^\gamma X_t,t0);\mathrm{P}_{x/c^\gamma }\}.$$
(2.6)
As a consequence, the semigroup densities of a semi-stable process with index $`\gamma `$ have the following property:
$$p_t(x,y)=t^{n\gamma }p_1(\frac{x}{t^\gamma },\frac{y}{t^\gamma }).$$
(2.7)
Remark that the expression for the semigroup densities in Theorem 2.3 satisfies this property for $`\gamma =\alpha /2`$. This remark yields the following corollary:
###### Corollary 2.5.
A Markov process that enjoys the time-inversion property of degree $`\alpha `$ is a semi-stable process of index $`\alpha /2`$, or is in $`h`$-transform relationship with it. The converse is not true.
###### Remark 2.6.
Throughout the paper, we assume the semigroup densities $`p_t(x,y)`$ to be positive over the domain and regular enough to be at least twice differentiable in the space and time variables. This assumption is more a technical requirement for the proof than a necessity for the characterization of Markov processes enjoying the time-inversion property.
## 3 Proof of the Theorem
### 3.1 Sufficiency
If $`p_t(x,y)`$ satisfies the condition (2.4), then from formula (2.2) the semigroup densities $`q_{s,t}^{(x)}(a,b)`$ can be written as
$$q_{s,t}^{(x)}(a,b)=(ts)^{\frac{n\alpha }{2}}\underset{i=1}{\overset{3}{}}R_{s,t}^{(i)}$$
where
$`R_{s,t}^{(1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }(xt^{\frac{\alpha }{2}},\frac{b}{t^\alpha }t^{\frac{\alpha }{2}})}{\mathrm{\Phi }(xs^{\frac{\alpha }{2}},\frac{a}{s^\alpha }s^{\frac{\alpha }{2}})}}\mathrm{\Phi }({\displaystyle \frac{b}{t^\alpha }}\left({\displaystyle \frac{st}{ts}}\right)^{\frac{\alpha }{2}},{\displaystyle \frac{a}{s^\alpha }}\left({\displaystyle \frac{st}{ts}}\right)^{\frac{\alpha }{2}})`$
$`R_{s,t}^{(2)}`$ $`=`$ $`{\displaystyle \frac{\theta \left(\frac{b}{t^\alpha }t^{\frac{\alpha }{2}}\right)}{\theta \left(\frac{a}{s^\alpha }s^{\frac{\alpha }{2}}\right)}}\theta \left({\displaystyle \frac{a}{s^\alpha }}\left({\displaystyle \frac{st}{ts}}\right)^{\frac{\alpha }{2}}\right)`$
$`R_{s,t}^{(3)}`$ $`=`$ $`\mathrm{exp}\{\rho \left(xt^{\frac{\alpha }{2}}\right)+\rho \left({\displaystyle \frac{b}{t^\alpha }}t^{\frac{\alpha }{2}}\right)+\rho \left({\displaystyle \frac{b}{t^\alpha }}\left({\displaystyle \frac{st}{ts}}\right)^{\frac{\alpha }{2}}\right)`$
$`+\rho \left({\displaystyle \frac{a}{s^\alpha }}\left({\displaystyle \frac{st}{ts}}\right)^{\frac{\alpha }{2}}\right)\rho \left(xs^{\frac{\alpha }{2}}\right)\rho \left({\displaystyle \frac{a}{s^\alpha }}s^{\frac{\alpha }{2}}\right)\}`$
Using the properties of $`\mathrm{\Phi },\theta ,\rho `$ described in Theorem 2.3, we obtain:
$`R_{s,t}^{(1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }(x,b)}{\mathrm{\Phi }(x,a)}}\mathrm{\Phi }({\displaystyle \frac{b}{(ts)^{\frac{\alpha }{2}}}},{\displaystyle \frac{a}{(ts)^{\frac{\alpha }{2}}}})`$
$`R_{s,t}^{(2)}`$ $`=`$ $`(ts)^{\frac{\alpha }{2}\beta }\theta (b)`$
$`R_{s,t}^{(3)}`$ $`=`$ $`\mathrm{exp}\left\{(ts)\rho (x)+\rho \left({\displaystyle \frac{a}{(ts)^{\frac{\alpha }{2}}}}\right)+\rho \left({\displaystyle \frac{b}{(ts)^{\frac{\alpha }{2}}}}\right)\right\}.`$
Hence there is no separate dependence on $`s`$ and $`t`$, but only on the difference $`ts`$, which allows to conclude that
$$q_{s,t}^{(x)}(a,b)=q_{ts}^{(x)}(a,b)$$
and proves homogeneity for the process $`(t^\alpha X_{\frac{1}{t}},t>0)`$. If in addition $`\mathrm{\Phi }(x,y)=\mathrm{\Phi }(y,x)`$, then
$`q_{s,t}^{(x)}(a,b)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }(x,b)}{\mathrm{\Phi }(x,a)}}{\displaystyle \frac{e^{(ts)\rho (x)}}{(ts)^{\frac{n\alpha }{2}}}}\mathrm{\Phi }({\displaystyle \frac{a}{(ts)^{\frac{\alpha }{2}}}},{\displaystyle \frac{b}{(ts)^{\frac{\alpha }{2}}}})\theta \left({\displaystyle \frac{b}{(ts)^{\frac{\alpha }{2}}}}\right)`$
$`\mathrm{exp}\left\{\rho \left({\displaystyle \frac{a}{(ts)^{\frac{\alpha }{2}}}}\right)+\rho \left({\displaystyle \frac{b}{(ts)^{\frac{\alpha }{2}}}}\right)\right\},`$
in which we recognize (2.5).
### 3.2 Necessity
For simplicity of notation, we prove the necessity of condition (2.4) in the case $`\alpha =1`$. The extension to $`\alpha >0`$ is immediate by the change of variables $`xx^\alpha \frac{x}{x}`$. Recall first the following definition of homogeneous functions:
###### Definition 3.1.
A function $`f:^n`$ that satisfies for $`x=(x_1,\mathrm{},x_n)`$, $`\lambda _+`$,
$$f(\lambda x)=\lambda ^\beta f(x)$$
is called homogeneous of degree $`\beta `$. If $`fC^1(^n)`$, Euler’s Homogeneous Function Theorem gives a necessary and sufficient condition for the function $`f(x)`$ to be homogeneous:
$$\underset{i=1}{\overset{n}{}}x_i\frac{}{x_i}f(x)=\beta f(x).$$
Consider the function of $`2n+1`$ variables: $`l(x,y,t)=\mathrm{ln}(p_t(x,y)),x,y^n`$. From (2.2), $`l`$ must satisfy for $`s=th`$:
$$\frac{}{t}\left[n\mathrm{ln}\frac{1}{t}+l(x,\frac{b}{t},\frac{1}{t})+l(\frac{b}{t},\frac{a}{th},\frac{1}{th}\frac{1}{t})l(x,\frac{a}{th},\frac{1}{th})\right]=0.$$
(3.1)
#### 3.2.1 The kernel $`\mathrm{\Phi }(x,y)`$
Taking derivatives with respect to $`b_i`$ and $`a_j`$ for some $`0i,jn`$ yields:
$$\frac{}{b_i}\frac{}{a_j}\frac{}{t}l(\frac{b}{t},\frac{a}{th},\frac{1}{th}\frac{1}{t})=0.$$
If we set $`\varphi (x,y,t)=\frac{}{x_i}\frac{}{y_j}l(x,y,t)`$, the latter becomes
$$\frac{1}{t^2(th)}\left(\varphi +\frac{b}{t}_1\varphi \right)\frac{1}{t(th)^2}\left(\varphi +\frac{a}{th}_2\varphi \right)$$
$$+\frac{1}{t(th)}\left(\frac{1}{t^2}\frac{1}{(th)^2}\right)\varphi =0,$$
with the notation $`_1=(\frac{}{x_1},\mathrm{},\frac{}{x_n})^T,_2=(\frac{}{y_1},\mathrm{},\frac{}{y_n})^T`$, $``$ the derivative with respect to the time variable and $`\varphi =\varphi (\frac{b}{t},\frac{a}{th},\frac{1}{th}\frac{1}{t})`$. For clarity, we change variables to
$$z_1=\frac{b}{t},z_2=\frac{a}{th},t_1=\frac{1}{t},t_2=\frac{1}{th}.$$
Then $`\varphi =\varphi (z_1,z_2,t_1+t_2)`$ and
$$t_1\left(\varphi +z_1_1\varphi +t_1\varphi \right)=t_2\left(\varphi +z_2_2\varphi +t_2\varphi \right),$$
or equivalently
$$t_1\left(\varphi +z_1_1\varphi +(t_1+t_2)\varphi \right)=t_2\left(\varphi +z_2_2\varphi +(t_1+t_2)\varphi \right).$$
We change variables once more: $`u=\frac{t_1}{t_2},v=t_1+t_2`$ to get:
$$\varphi +z_1_1\varphi +v\varphi =u\left(\varphi +z_2_2\varphi +v\varphi \right).$$
(3.2)
This gives the following proposition.
###### Proposition 3.2.
Equation (3.2) is satisfied if and only if
$$\varphi (\lambda x,\mu y,\lambda \mu t)=\frac{1}{\lambda \mu }\varphi (x,y,t).$$
(3.3)
###### Proof.
As the LHS of equation (3.2) is independent of $`u`$, one can readily take the equivalent condition:
$`\varphi +z_1_1\varphi +v\varphi `$ $`=`$ $`0,`$ (3.4)
$`\varphi +z_2_2\varphi +v\varphi `$ $`=`$ $`0.`$ (3.5)
Let $`g(\lambda )=\varphi (\lambda z_1,z_2,\lambda v)`$ and $`h(\mu )=\varphi (z_1,\mu z_2,\mu v)`$. Equation (3.4) implies $`g(\lambda )+\lambda g^{}(\lambda )=0,`$ which solves to $`g(\lambda )=\lambda ^1g(1)`$ and hence
$$\varphi (\lambda z_1,z_2,\lambda v)=\frac{1}{\lambda }\varphi (z_1,z_2,v).$$
Equation (3.5) implies $`h(\mu )+\mu h^{}(\mu )=0,`$ which solves to $`h(\mu )=\mu ^1h(1)`$ and hence
$$\varphi (z_1,\mu z_2,\mu v)=\frac{1}{\mu }\varphi (z_1,z_2,v).$$
Combining the latter two equations yields (3.3).
Conversely, if $`g`$ is homogeneous of degree $`1`$ and $`h`$ is homogeneous of degree $`1`$, we get:
$$g(\lambda )+\lambda g^{}(\lambda )=h(\mu )+\mu h^{}(\mu )=0,$$
which is equivalent to (3.4) and (3.5) and concludes the proof. ∎
The scaling property (3.3) implies moreover the equivalent formulation
$$\varphi (x,y,t)=\frac{1}{t}\varphi _1(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}}),$$
(3.6)
where $`\varphi _1(x,y)=\varphi (x,y,1)`$ must satisfy
$$\varphi _1(\lambda x,y)=\varphi _1(x,\lambda y).$$
(3.7)
Under the change of variables $`\overline{x}={\displaystyle \frac{x}{\sqrt{t}}}`$, $`\overline{y}={\displaystyle \frac{y}{\sqrt{t}}}`$, the kernel $`\overline{l}(\overline{x},\overline{y},t)=l(x,y,t)`$ satisfies
$$\varphi _1(\overline{x},\overline{y})=\frac{}{\overline{x}_i}\frac{}{\overline{y}_j}\overline{l}(\overline{x},\overline{y},t).$$
Note that the previous results remain valid for all second derivatives of $`\overline{l}(\overline{x},\overline{y},t)`$ with respect to $`\overline{x}_i`$ and $`\overline{y}_j`$, $`i,j\{1,\mathrm{},n\}`$. By integration over $`\overline{x}_i`$ and $`\overline{y}_j`$, one can thus already make an assumption on the general shape of the kernel $`\overline{l}(\overline{x},\overline{y},t)`$:
$$\overline{l}(\overline{x},\overline{y},t)=k(\overline{x},\overline{y})+\overline{\phi }_1(\overline{x},t)+\overline{\phi }_2(\overline{y},t),$$
where $`k:^{2n}`$ does not depend explicitly on time and $`\overline{\phi }_1,\overline{\phi }_2:^{n+1}`$. An equivalent representation for $`l(x,y,t)`$ gives
$$l(x,y,t)=k(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})+\phi _1(x,t)+\phi _2(y,t),$$
(3.8)
with $`\overline{\phi }_1(\overline{x},t)=\phi _1(x,t)`$ and $`\overline{\phi }_2(\overline{y},t)=\phi _2(y,t)`$.
The scaling property (3.7) of $`\varphi _1(x,y)`$ now translates to $`k(x,y)`$ as follows:
$$\frac{}{x_i}\frac{}{y_j}k(\lambda x,y)=\frac{}{x_i}\frac{}{y_j}k(x,\lambda y),i,j\{1,\mathrm{},n\}.$$
(3.9)
By integration, $`k(x,y)`$ must satisfy the following condition,
$$k(\lambda x,y)=k(x,\lambda y)+\zeta _1(x,\lambda )+\zeta _2(y,\lambda ),$$
(3.10)
along with the properties $`\zeta _1(x,1)+\zeta _2(y,1)=0`$ and
$$\zeta _1(x,\alpha \lambda )+\zeta _2(y,\alpha \lambda )=\zeta _1(x,\alpha )+\zeta _1(\alpha x,\lambda )+\zeta _2(y,\lambda )+\zeta _2(\lambda y,\alpha ).$$
###### Proposition 3.3.
The functions $`\zeta _1(x,\lambda )`$ and $`\zeta _2(y,\lambda )`$ are equivalently defined by
$`\zeta _1(x,\alpha )`$ $`=`$ $`\zeta _1(x)\zeta _1(\alpha x),`$
$`\zeta _2(x,\alpha )`$ $`=`$ $`\zeta _2(x)\zeta _2(\alpha x),`$ (3.11)
with a slight abuse of notation.
###### Proof.
The first property of the functions $`\zeta _1(x,\lambda ),\zeta _2(y,\lambda )`$ implies
$$\zeta _1(x,1)=\zeta _2(y,1)=\zeta $$
which is set to 0 without loss of generality. The second property leads to
$`\zeta _1(x,\alpha \lambda )`$ $`=`$ $`\zeta _1(x,\alpha )+\zeta _1(\alpha x,\lambda )+F(\alpha ,\lambda ),`$
$`\zeta _2(y,\alpha \lambda )`$ $`=`$ $`\zeta _2(y,\lambda )+\zeta _2(\lambda y,\alpha )F(\alpha ,\lambda ),`$ (3.12)
for some function $`F:^2`$. Moreover, by deriving the latter expression involving $`\zeta _1(x,\lambda )`$ successively by $`x_i`$ and $`\lambda `$, we get
$$\frac{}{x_i}\frac{}{\lambda }\zeta _1(\alpha x,\lambda )=\frac{}{x_i}\frac{}{\lambda }\zeta _1(x,\alpha \lambda ).$$
(3.13)
On the other hand, taking the derivatives with respect to $`\lambda `$ and $`\alpha `$ leads to
$$\zeta _1^{}(x,\alpha \lambda )+\alpha \lambda \zeta _1^{\prime \prime }(x,\alpha \lambda )=x_{\alpha x}\zeta _1^{}(\alpha x,\lambda )+\frac{}{\alpha }\frac{}{\lambda }F(\alpha ,\lambda ),$$
with denoting the derivative with respect to the last coordinate in $`\zeta _1(x,\lambda )`$. Using (3.13), we obtain the following:
$$\frac{}{\alpha }\frac{}{\lambda }F(\alpha ,\lambda )=\zeta _1^{}(x,\alpha \lambda )+\alpha \lambda \zeta _1^{\prime \prime }(x,\alpha \lambda )x_x\zeta _1^{}(x,\alpha \lambda ),$$
which shows that the second derivative $`\frac{}{\alpha }\frac{}{\lambda }F(\alpha ,\lambda )`$ is a function of the product $`\alpha \lambda `$ only. Hence with a slight abuse of notation,
$$F(\alpha ,\lambda )=F(\alpha \lambda )F_1(\alpha )F_2(\lambda ),$$
where we set $`F_1(1)=F_2(1)=0`$ without loss of generality. Setting respectively $`\alpha =1`$ and $`\lambda =1`$ yields from (3.2.1), $`F_1(\lambda )=F(\lambda )`$ and $`F_2(\lambda )=F(\lambda )`$, hence $`F(1)=0`$. Equation (3.2.1) then becomes
$`\zeta _1(x,\alpha \lambda )F(\alpha \lambda )`$ $`=`$ $`\zeta _1(x,\alpha )F(\alpha )+\zeta _1(\alpha x,\lambda )F(\lambda ),`$
$`\zeta _2(y,\alpha \lambda )+F(\alpha \lambda )`$ $`=`$ $`\zeta _2(y,\lambda )+F(\lambda )+\zeta _2(\lambda y,\alpha )+F(\alpha ).`$
$`F(\lambda )`$ turns out to simply shift the functions $`\zeta _1(x,\lambda )`$ and $`\zeta _2(y,\lambda )`$. It can thus be included as part of those functions, so without loss of generality $`F(\lambda )=0`$ and the latter equations become
$`\zeta _1(x,\alpha \lambda )`$ $`=`$ $`\zeta _1(x,\alpha )+\zeta _1(\alpha x,\lambda ),`$
$`\zeta _2(y,\alpha \lambda )`$ $`=`$ $`\zeta _2(y,\lambda )+\zeta _2(\lambda y,\alpha ).`$ (3.14)
The latter functions can thus be further reduced to functions of one variable. With a slight abuse of notation, setting $`\zeta _i(x)=\zeta _i(x,\frac{1}{x})`$ and choosing respectively $`\lambda =\frac{1}{\alpha x}`$ and $`\alpha =\frac{1}{\lambda y}`$ yields the resulting equations (3.3). ∎
This last proposition combined with (3.10) gives
$$k(\lambda x,y)+\zeta _1(\lambda x)\zeta _2(y)=k(x,\lambda y)+\zeta _1(x)\zeta _2(\lambda y).$$
Now setting $`\mathrm{\Phi }(x,y)=\mathrm{exp}\left(k(x,y)+\zeta _1(x)\zeta _2(y)\right)`$ completes the proof of the first condition of the theorem, i.e.
$$\mathrm{\Phi }(\lambda x,y)=\mathrm{\Phi }(x,\lambda y).$$
#### 3.2.2 The function $`\rho (x)`$
Going back to (3.1), we replace the explicit form of $`l(x,y,t)`$ and use the scaling property of $`\zeta _1(x,\lambda ),\zeta _2(y,\lambda )`$ and $`k(x,y)`$ to obtain
$`{\displaystyle \frac{}{t}}`$ $`[n\mathrm{ln}{\displaystyle \frac{1}{t}}+\zeta _1(x,\sqrt{t})\zeta _1(x,\sqrt{th})\zeta _2(b,{\displaystyle \frac{1}{\sqrt{t}}})+\zeta _2(a,{\displaystyle \frac{1}{\sqrt{th}}})`$ (3.15)
$`+\zeta _1({\displaystyle \frac{b}{\sqrt{h}}},\sqrt{{\displaystyle \frac{th}{t}}})\zeta _2({\displaystyle \frac{a}{\sqrt{h}}},\sqrt{{\displaystyle \frac{t}{th}}})`$
$`+\phi _1(x,{\displaystyle \frac{1}{t}})\phi _1(x,{\displaystyle \frac{1}{th}})+\phi _1({\displaystyle \frac{b}{t}},{\displaystyle \frac{h}{t(th)}})`$
$`+\phi _2({\displaystyle \frac{b}{t}},{\displaystyle \frac{1}{t}})\phi _2({\displaystyle \frac{a}{th}},{\displaystyle \frac{1}{th}})+\phi _2({\displaystyle \frac{a}{th}},{\displaystyle \frac{h}{t(th)}})]=0.`$
Note that the terms composed of the function $`k(x,y)`$ do not depend on $`t`$ and thus cancel out with the time derivative.
Recall that the variables $`a,b,x`$ are independent from each other. So taking the derivative of (3.15) with respect to $`x_i`$ gives
$$t_1^2\frac{}{x_i}\phi _1^{}(x,t_1)\frac{\sqrt{t_1}}{2}\frac{}{x_i}\zeta _1^{}(x,\frac{1}{\sqrt{t_1}})=t_2^2\frac{}{x_i}\phi _1^{}(x,t_2)\frac{\sqrt{t_2}}{2}\frac{}{x_i}\zeta _1^{}(x,\frac{1}{\sqrt{t_2}})$$
with $`t_1=\frac{1}{t}`$, $`t_2=\frac{1}{th}`$ and still denotes the derivative with respect to the last coordinate. Since the latter is valid for all $`t_1,t_2>0`$, each side of the equation must be independent of $`t`$. Hence, for some differentiable function $`\phi _{11}:^n`$,
$$\frac{}{t}\frac{}{x_i}\phi _1(x,t)=\frac{1}{t^2}\frac{}{x_i}\phi _{11}(x)\frac{}{t}\frac{}{x_i}\zeta _1(x,\frac{1}{\sqrt{t}}),$$
which integrates to
$$\phi _1(x,t)=\frac{1}{t}\phi _{11}(x)\zeta _1(x,\frac{1}{\sqrt{t}})+h_1(x)+\tau _1(t).$$
There are so far no further conditions to add on the function $`\phi _1(x,t)`$.
With the new result for the shape of $`\phi _1(x,t)`$ and using the scaling property of $`\zeta _1(x,\lambda )`$ and $`\zeta _2(y,\lambda )`$, we derive (3.15) with respect to $`b_i`$ to obtain
$`{\displaystyle \frac{}{b_i}}{\displaystyle \frac{}{t}}`$ $`[\zeta _1(b,{\displaystyle \frac{1}{t}})\zeta _2(b,{\displaystyle \frac{1}{\sqrt{t}}})+{\displaystyle \frac{t(th)}{h}}\phi _{11}\left({\displaystyle \frac{b}{t}}\right)`$
$`+h_1\left({\displaystyle \frac{b}{t}}\right)+\phi _2({\displaystyle \frac{b}{t}},{\displaystyle \frac{1}{t}})]=0.`$
For convenience, let $`\phi _{22}:^{n+1}`$ be such that
$$\phi _2(x,t)=\zeta _2(\frac{x}{t},\sqrt{t})\zeta _1(\frac{x}{t},t)+\frac{1}{t}\phi _{11}(x)h_1(x)+\phi _{22}(x,t).$$
(3.16)
The former partial derivative equation then becomes
$$\frac{}{b_i}\frac{}{t}\left[\frac{t^2}{h}\phi _{11}\left(\frac{b}{t}\right)+\phi _{22}(\frac{b}{t},\frac{1}{t})\right]=0,$$
(3.17)
which develops to
$$t_1^2\left\{1+t_1\frac{}{t_1}+z\right\}_i\phi _{22}(z,t_1)=\frac{1}{h}\left\{1z\right\}_i\phi _{11}(z),$$
with the notation $`z=\frac{b}{t},t_1=\frac{1}{t}`$. Since the LHS is independent of $`h`$, both sides of the equation must cancel out. This is the case if $`_i\phi _{22}(z,t)`$ is homogeneous of degree $`1`$ and $`_i\phi _{11}(z)`$ is homogeneous of degree 1, which means by integration over $`z_i`$,
$$\phi _{11}(\lambda z)=\lambda ^2\phi _{11}(z)+\overline{\phi }\mathrm{ln}\lambda ,\overline{\phi }.$$
(3.18)
Let $`\rho (z)=\phi _{11}(z)`$, then, conditioned on showing $`\overline{\phi }=0`$, we recover the second condition of the theorem, that is
$$\rho (\lambda z)=\lambda ^2\rho (z).$$
#### 3.2.3 The function $`\theta (y)`$
In order to further investigate the properties of $`\phi _{22}(z,t)`$, we derive (3.15) by $`a_i`$:
$`{\displaystyle \frac{}{a_i}}{\displaystyle \frac{}{t}}`$ $`[\phi _{22}({\displaystyle \frac{a}{th}},{\displaystyle \frac{h}{t(th)}})\phi _{22}({\displaystyle \frac{a}{th}},{\displaystyle \frac{1}{th}})`$
$`+\zeta _1(a,{\displaystyle \frac{t}{h}})\zeta _2(a,{\displaystyle \frac{t}{h}})]=0,`$
using the scaling properties of $`\zeta _1(x,\lambda ),\zeta _2(y,\lambda )`$ and $`\phi _{11}(x)`$. Since $`_i\phi _{22}(z,t)`$ is homogeneous of degree $`1`$, this leads to
$$\frac{}{t}\left[_i\phi _{22}(a,\frac{h}{t})+_i\zeta _1(a,\frac{t}{h})_i\zeta _2(a,\frac{t}{h})\right]=0.$$
Hence, by integration over time, we obtain
$$_i\phi _{22}(a,\frac{h}{t})+_i\zeta _1(a,\frac{t}{h})_i\zeta _2(a,\frac{t}{h})=_ih_2(a),$$
for some differentiable function $`h_2:^n`$. Integration with respect to $`a_i`$ yields
$$\phi _{22}(a,t)=h_2(a)+\tau _2(t)\zeta _1(a,\frac{1}{t})+\zeta _2(a,\frac{1}{t}),$$
for some real function $`\tau _2(t)`$. Furthermore the homogeneity condition on $`_i\phi _{22}(z,t)`$ translates to $`_ih_2(z)`$ as follows:
$$\lambda _ih_2(\lambda a)=_ih_2(a)+_i\zeta _2(a,\lambda )_i\zeta _1(a,\lambda ),$$
by the scaling properties of $`_i\zeta _1(x,\lambda )`$ and $`_i\zeta _2(y,\lambda )`$. Integration over $`a_i`$ then yields
$$h_2(\lambda a)=h_2(a)+H(\lambda )+\zeta _2(a,\lambda )\zeta _1(a,\lambda ),$$
for $`H:`$ with the properties $`H(1)=0`$ and $`H(\lambda _1\lambda _2)=H(\lambda _1)+H(\lambda _2)`$. The function $`H(\lambda )`$ is thus of the form $`H(\lambda )=\beta \mathrm{ln}\lambda `$ for some $`\beta `$. Using Proposition 3.3, the latter becomes
$$h_2(\lambda a)+\zeta _2(\lambda a)\zeta _1(\lambda a)=\beta \mathrm{ln}\lambda +h_2(a)+\zeta _2(a)\zeta _1(a).$$
Let $`\theta (z)=\mathrm{exp}\left(h_2(z)+\zeta _2(a)\zeta _1(a)\right)`$. We then recover condition 3 of the theorem, i.e.
$$\theta (\lambda z)=\lambda ^\beta \theta (z).$$
#### 3.2.4 The function $`\tau (t)`$
The results for $`\phi _1(x,t)`$ and $`\phi _2(y,t)`$ summarize so far to:
$`\phi _1(x,t)`$ $`=`$ $`{\displaystyle \frac{\rho (x)}{t}}\zeta _1(x)+\zeta _1\left({\displaystyle \frac{x}{\sqrt{t}}}\right)+h_1(x)+\tau _1(t)`$
$`\phi _2(y,t)`$ $`=`$ $`{\displaystyle \frac{\rho (y)}{t}}+\zeta _2(y)\zeta _2\left({\displaystyle \frac{y}{\sqrt{t}}}\right)h_1(y)+\tau _2(t)+h_2(y).`$
This implies that $`l(x,y,t)=\mathrm{ln}(p_t(x,y))`$ has the form
$$l(x,y,t)=\mathrm{ln}\mathrm{\Phi }(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})+\mathrm{ln}\frac{\eta (y)}{\eta (x)}+\frac{\rho (x)+\rho (y)}{t}+\tau _1(t)+\tau _2(t),$$
(3.19)
where $`\eta (x)=\mathrm{exp}\left(\zeta _1(x)h_1(x)\right)`$ defines an $`h`$-transform and can thus be neglected. It remains to explicitly formulate the form of the function $`\tau :_+`$ that we define by
$$\tau (t)=\tau _1(t)+\tau _2(t)+\frac{1}{2}(n+\beta )\mathrm{ln}t$$
to satisfy the resulting partial differential equation (3.15):
$`{\displaystyle \frac{}{t}}`$ $`[\overline{\phi }{\displaystyle \frac{t(th)}{h}}({\displaystyle \frac{t}{th}}\mathrm{ln}{\displaystyle \frac{1}{t}}+{\displaystyle \frac{th}{t}}\mathrm{ln}{\displaystyle \frac{1}{th}})`$ (3.20)
$`+\tau \left({\displaystyle \frac{1}{t}}\right)+\tau \left({\displaystyle \frac{h}{t(th)}}\right)\tau \left({\displaystyle \frac{1}{th}}\right)]=0.`$
The former equation becomes after derivation,
$$\frac{\overline{\phi }}{t_2t_1}\left(t_1(2\mathrm{ln}t_2+1)+t_2(2\mathrm{ln}t_1+1)\right)=t_1^2\tau ^{}(t_1)+(t_2^2t_1^2)\tau ^{}(t_2t_1)t_2^2\tau ^{}(t_2),$$
for $`t_1=\frac{1}{t}`$, $`t_2=\frac{1}{th}`$, which hints that $`\overline{\phi }=0`$. Assuming $`\overline{\phi }`$ to be non-zero and setting $`g(t)={\displaystyle \frac{t^2}{2\overline{\phi }}}\tau ^{}(t)+{\displaystyle \frac{1}{2}}`$, we get indeed
$$t_1\mathrm{ln}t_2+t_2\mathrm{ln}t_1=t_1\left(g(t_2t_1)g(t_1)+g(t_2)\right)+t_2\left(g(t_2t_1)+g(t_1)g(t_2)\right),$$
which has no solution. So $`\overline{\phi }=0`$. Now dividing the former equation with $`\overline{\phi }=0`$ by $`t_1`$ and taking the limit $`t_10`$ leads to the differential equation
$$\tau ^{\prime \prime }(t)=0,$$
which solves for $`\tau (t)=\nu t,\nu `$. However, $`e^{\nu t}`$ can always be included in an $`h`$-transform, so $`\nu `$ is set to 0 and $`\tau (t)`$ becomes trivial.
#### 3.2.5 The general case $`\alpha >0`$
Combining the different factors, we obtain for the semigroup densities,
$$p_t(x,y)=t^{\frac{1}{2}\left(n+\beta \right)}\mathrm{\Phi }(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})\theta (y)\mathrm{exp}\left\{\frac{\rho (x)}{t}+\frac{\rho (y)}{t}\right\},$$
which simplifies to, using the scaling property of $`\theta (y)`$ and $`\rho (x)`$,
$$p_t(x,y)=t^{\frac{n}{2}}\mathrm{\Phi }(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})\theta \left(\frac{y}{\sqrt{t}}\right)\mathrm{exp}\left\{\rho \left(\frac{x}{\sqrt{t}}\right)+\rho \left(\frac{y}{\sqrt{t}}\right)\right\}.$$
For $`\alpha >0`$, we apply the change of variables $`xx^\alpha \frac{x}{x}`$. The semigroup densities become
$$p_t^\alpha (x,y)=p_t(x^{\frac{1}{\alpha }},y^{\frac{1}{\alpha }})J(y)$$
where $`J(y)`$ is the Jacobian of the inverse transformation, that is $`J(y)=\alpha ^ny_1^{\frac{1}{\alpha }1}\mathrm{}y_n^{\frac{1}{\alpha }1}`$, and we use the slight abuse of notation $`x^{\frac{1}{\alpha }}=(x_1^{\frac{1}{\alpha }},\mathrm{},x_n^{\frac{1}{\alpha }})`$. The semigroup densities can be recast into
$$p_t^\alpha (x,y)=t^{\frac{n\alpha }{2}}\stackrel{~}{\mathrm{\Phi }}(\frac{x}{t^{\frac{\alpha }{2}}},\frac{y}{t^{\frac{\alpha }{2}}})\stackrel{~}{\theta }\left(\frac{y}{t^{\frac{\alpha }{2}}}\right)\mathrm{exp}\left\{\stackrel{~}{\rho }\left(\frac{x}{t^{\frac{\alpha }{2}}}\right)+\stackrel{~}{\rho }\left(\frac{y}{t^{\frac{\alpha }{2}}}\right)\right\}$$
where $`\stackrel{~}{\mathrm{\Phi }}(x,y)=\mathrm{\Phi }(x^{\frac{1}{\alpha }},y^{\frac{1}{\alpha }})`$ satisfies condition 1, $`\stackrel{~}{\theta }(y)=\theta \left(y^{\frac{1}{\alpha }}\right)J(y)`$ satisfies condition 2 for $`\stackrel{~}{\beta }=\frac{\beta +n}{\alpha }n`$ and $`\stackrel{~}{\rho }(x)=\rho \left(y^{\frac{1}{\alpha }}\right)`$ satisfies condition 3 of equation (2.4).
## 4 Application to diffusions on $`_+`$
The case of the diffusions on $`_+`$ was entirely characterized by Watanabe in . It was shown that only Bessel processes in the wide sense (which we recall the definition below) enjoy the time-inversion property of degree 1.
###### Definition 4.1.
For some $`\nu >1`$ and $`c0`$, the diffusion process generated by
$$=\frac{1}{2}\frac{^2}{x^2}+\left(\frac{2\nu +1}{2x}+\frac{h_c^{}(x)}{h_c(x)}\right)\frac{}{x}$$
(4.1)
is called Bessel process in the wide sense. The function $`h_c(x)`$ is given by
$$h_c(x)=2^\nu \mathrm{\Gamma }(\nu +1)(\sqrt{2c}x)^\nu I_\nu (\sqrt{2c}x),$$
(4.2)
where $`I_\nu `$ is the modified Bessel function.
###### Remark 4.2.
The Bessel process in the wide sense is in $`h`$-transform relationship, for $`hh_c`$, with the Bessel process.
We show that the result in is a consequence of Theorem 2.3, which has the following one-dimensional formulation:
###### Theorem 4.3.
The Markov process $`(t^\alpha X_{\frac{1}{t}},t>0)`$ on $`_+`$ is homogeneous if and only if the semigroup densities of $`(X_t,t0)`$ are of the form:
$$p_t(x,y)=t^{\frac{\alpha }{2}(1+\beta )}\varphi \left(\frac{xy}{t^\alpha }\right)y^\beta \mathrm{exp}\left\{\frac{k^2}{2}\left(\frac{x^{\frac{2}{\alpha }}}{t}+\frac{y^{\frac{2}{\alpha }}}{t}\right)\right\}$$
(4.3)
for $`k>0`$, or if it is in $`h`$-transform relationship with it. Moreover, the semigroup densities are related as follows:
$$q_t^{(x)}(a,b)=\frac{\varphi \left(xb\right)}{\varphi \left(xa\right)}\mathrm{exp}\left(t\frac{k^2}{2}x^{\frac{2}{\alpha }}\right)p_t(a,b).$$
(4.4)
###### Proof.
Theorem 2.3 formulated on $`_+`$ gives $`\mathrm{\Phi }(x,y)=\varphi (xy)`$, $`\rho (x)=\frac{k^2}{2}x^{\frac{2}{\alpha }},k>0`$, and $`\theta (y)=y^\beta `$ for conditions 1, 2 and 3 to be satisfied. Equation (4.4) is then a consequence of (2.2). ∎
We identify further the class of diffusion processes and provide a different proof for the result in .
###### Proposition 4.4.
If $`(X_t,t0)`$ is a diffusion process and $`(tX_{\frac{1}{t}},t>0)`$ is homogeneous and conservative, then both are necessarily (possibly time-scaled) Bessel processes in the wide sense.
###### Proof.
If $`(X_t,t0)`$ is a diffusion process, then its infinitesimal generator has the following general structure:
$$=\frac{d}{𝔪(dy)}\frac{d}{d𝔰(y)}$$
where $`𝔪(dy)`$ is called the speed measure and $`𝔰(y)`$ the scale measure density (see ). From the assumptions of the proposition, $`(X_t,t0)`$ enjoys the time-inversion property of degree 1. As a consequence of (4.3) in Theorem 4.3, the speed measure is thus absolutely continuous with respect to the Lebesgue measure. Let $`\mu (x)=\frac{1}{𝔪(x)𝔰(x)}`$ and $`s(x)=\frac{𝔰^{}(x)}{𝔪(x)𝔰(x)^2}`$. The infinitesimal generator can then be expressed as
$$=s(y)\frac{^2}{y^2}+\mu (y)\frac{}{y}$$
where it remains to identify the functions $`s(y)`$ and $`\mu (y)`$. For a fixed $`x>0`$, let $`^{(x)}`$ be the infinitesimal generator of $`(tX_{\frac{1}{t}},t>0)`$. From equation (4.4), $`^{(x)}`$ has the following relationship with $``$:
$$^{(x)}:f(b)\frac{1}{\varphi (xb)}\left(\varphi (xb)f(b)\right)\frac{k^2}{2}x^2f(b),$$
which develops to
$$^{(x)}f(b)=s(b)f^{\prime \prime }(b)+\left\{\mu (b)+2xs(b)\frac{\varphi ^{}(xb)}{\varphi (xb)}\right\}f^{}(b)+U(x,b)f(b).$$
For the process to be conservative, we require $`U(x,b)=0`$, which implies no killing in the interior of the domain, that is
$$s(b)x^2\frac{\varphi ^{\prime \prime }(xb)}{\varphi (xb)}+\mu (b)x\frac{\varphi ^{}(xb)}{\varphi (xb)}\frac{k^2}{2}x^2=0.$$
We change variables to $`z=xb`$ to obtain
$$s(z/x)\varphi ^{\prime \prime }(z)+\frac{\mu (z/x)}{x}\varphi ^{}(z)\frac{k^2}{2}\varphi (z)=0.$$
Since the latter must be valid for all $`x>0`$, we are led to set: $`s(b)=\frac{\sigma ^2}{2}`$ for $`\sigma >0`$ and $`\mu (b)=\frac{\sigma ^2}{2}\frac{2\nu +1}{b}`$ for $`\nu >1`$. This yields the following equation
$$\frac{1}{2}\varphi ^{\prime \prime }(z)+\frac{2\nu +1}{2z}\varphi ^{}(z)\frac{k^2}{2\sigma ^2}\varphi (z)=0.$$
The general solution (non-singular at 0 and up to a constant factor) is expressed through the modified Bessel function of the first kind as follows:
$$\varphi (z)=\left(\frac{kz}{\sigma }\right)^\nu I_\nu \left(\frac{kz}{\sigma }\right).$$
Gathering the different factors in (4.3) leads to
$$p_t(x,y)=Nt^{\frac{1+\beta }{2}}\left(\frac{xy}{t}\right)^\nu I_\nu \left(\frac{k}{\sigma }\frac{xy}{t}\right)y^\beta \mathrm{exp}\left\{\frac{k^2}{2}\left(\frac{x^2}{t}+\frac{y^2}{t}\right)\right\},$$
where $`N`$ is a normalization factor. The additional condition that $`\underset{t0}{lim}p_t(x,y)=\delta (xy)`$ implies $`\beta =2\nu +1`$ and $`k=\frac{1}{\sigma }`$, which leads to the semigroup densities of a time-scaled Bessel process of dimension $`\nu `$:
$$p_t(x,y)=\frac{y}{\sigma ^2t}\left(\frac{y}{x}\right)^\nu I_\nu \left(\frac{xy}{\sigma ^2t}\right)\mathrm{exp}\left\{\frac{x^2+y^2}{2\sigma ^2t}\right\}.$$
The infinitesimal generator for $`(tX_{\frac{1}{t}},t>0)`$ is given by
$$^{(x)}=\frac{\sigma ^2}{2}\frac{^2}{b^2}+\sigma ^2\left\{\frac{2\nu +1}{2b}+x\frac{\varphi ^{}(xb)}{\varphi (xb)}\right\}\frac{}{b},$$
where one recognizes expression (4.1) for $`h_c(b)=\varphi (xb)`$ and a time-scale $`t\sigma ^2t`$. ∎
Proposition 4.4 smoothly extends to any power $`\alpha >0`$ of the Bessel process through the mapping $`xx^\alpha \frac{x}{x}`$. In particular, it is worth remarking that the case $`\alpha =2`$ gives rise to squares of Bessel processes, which leads to the following result:
###### Proposition 4.5.
If $`(X_t,t0)`$ is a diffusion process and $`(t^2X_{\frac{1}{t}},t>0)`$ is homogeneous and conservative, then both are necessarily (possibly time-scaled) squares of Bessel processes in the wide sense.
## 5 Examples
### 5.1 Generalized Dunkl processes and Jacobi-Dunkl processes
#### 5.1.1 Multidimensional Dunkl processes
We briefly review the construction of the Dunkl process in $`^n`$ (see ).
###### Definition 5.1.
The Dunkl process in $`^n`$ is the Markov càdlàg process with infinitesimal generator
$$\frac{1}{2}^{(k)}=\frac{1}{2}\underset{i=1}{\overset{n}{}}T_i^2$$
(5.1)
where $`T_i,1in,`$ is a one-dimensional differential-difference operator defined for $`uC^1(^n)`$ by
$$T_iu(x)=\frac{u(x)}{x_i}+\underset{\alpha R_+}{}k(\alpha )\alpha _i\frac{u(x)u(\sigma _\alpha x)}{\alpha ,x}.$$
(5.2)
$`,`$ is the usual scalar product. $`R`$ is a root system in $`^n`$ and $`R_+`$ a positive subsystem. $`k(\alpha )`$ is a non-negative multiplicity function defined on $`R`$ and invariant by the finite reflection group $`W`$ associated with $`R`$. $`\sigma _\alpha `$ is the reflection operator with respect to the hyperplane $`H_\alpha `$ orthogonal to $`\alpha `$ such that $`\sigma _\alpha x=x\alpha ,x\alpha `$ and for convenience $`\alpha ,\alpha =2`$ (see ).
A result obtained by M. Rösler yields the semigroup densities as follows:
$$p_t^{(k)}(x,y)=\frac{1}{c_kt^{\gamma +n/2}}\mathrm{exp}\left(\frac{|x|^2+|y|^2}{2t}\right)D_k(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})\omega _k(y)$$
(5.3)
where $`D_k(x,y)>0`$ is the Dunkl kernel, $`\omega _k(y)={\displaystyle \underset{\alpha R_+}{}}|\alpha ,y|^{2k(\alpha )}`$ the weight function which is homogeneous of degree $`2\gamma =2{\displaystyle \underset{\alpha R_+}{}}k(\alpha )`$ and $`c_k={\displaystyle _^n}e^{\frac{|x|^2}{2}}\omega _k(x)𝑑x`$.
Following a thorough study of the properties of the one-dimensional Dunkl process in , it was remarked in that the Dunkl process in $`^n`$ enjoys the time-inversion property of degree 1. Considering that the Dunkl kernel satisfies
$$D_k(x,y)=D_k(y,x)\mathrm{and}D_k(\mu x,y)=D_k(x,\mu y),$$
(5.4)
the proof is straightforward with
$$\mathrm{\Phi }(x,y)D_k(x,y),\theta (y)\omega _k(y),\rho (x)\frac{|x|^2}{2}.$$
(5.5)
By equation (2.5), the semigroup densities of the time-inverted process is even in $`h`$-transform relationship with the semigroup densities of the original Dunkl process:
$$q_t^{(x)}(a,b)=\frac{D_k(x,b)}{D_k(x,a)}\mathrm{exp}\left(\frac{|x|^2}{2}t\right)p_t^{(k)}(a,b).$$
(5.6)
#### 5.1.2 Generalized Dunkl processes
In an attempt to generalize the Dunkl process, we extend the definition of the infinitesimal generator to
$$^{(k,\lambda )}f(x)=\frac{1}{2}\mathrm{\Delta }f(x)+\underset{\alpha R_+}{}k(\alpha )\frac{f(x),\alpha }{x,\alpha }+\underset{\alpha R_+}{}\lambda (\alpha )\frac{f(\sigma _\alpha x)f(x)}{x,\alpha ^2}$$
(5.7)
where $`\mathrm{\Delta }`$ is the usual Laplacian, $`fC^2(^n)`$ and $`\lambda (\alpha )`$ is a non-negative multiplicity function defined on $`R`$ and invariant by the finite reflection group $`W`$, similarly as $`k(\alpha )`$. We retrieve the Dunkl process for $`\lambda (\alpha )=k(\alpha )`$. Note that these processes are no longer martingales for $`\lambda (\alpha )k(\alpha )`$.
The one-dimensional case was introduced in , where the semigroup densities were explicitly derived,
$$p_t^{(k,\lambda )}(x,y)=\frac{1}{t^{k\frac{1}{2}}}y^{k\frac{1}{2}}\mathrm{exp}\left(\frac{x^2+y^2}{2t}\right)D_{k,\lambda }\left(\frac{xy}{t}\right)$$
(5.8)
with the generalized Dunkl kernel, $`(\nu =k\frac{1}{2},\mu =\sqrt{\nu ^2+4\lambda })`$,
$$D_{k,\lambda }\left(z\right)=1_{\{y_{}\}}\frac{1}{2z^\nu }(I_\nu I_\mu )\left(z\right)+1_{\{y_+\}}\frac{1}{2z^\nu }(I_\nu +I_\mu )\left(z\right),$$
(5.9)
for $`z=\frac{xy}{t}`$. From (5.4), the generalized Dunkl kernel satisfies
$$D_{k,\lambda }(x,y)=D_{k,\lambda }(y,x)\mathrm{and}D_{k,\lambda }(\mu x,y)=D_{k,\lambda }(x,\mu y),$$
(5.10)
which readily implies that the generalized Dunkl process also enjoys the time-inversion property of degree 1. The semigroup densities were derived as an application of the skew-product representation of the generalized Dunkl process $`(X_t,t0)`$ in terms of its absolute value (a Bessel process) and an independent Poisson process $`N_t^{(\lambda )}`$:
$$X_t\stackrel{(d)}{=}|X_t|(1)^{N_{A_t}^{(\lambda )}}$$
(5.11)
where $`A_t=_0^t\frac{ds}{X_s^2}`$.
In the $`n`$-dimensional case, the application of the skew-product representation derived by Chybiryakov shows that the generalized Dunkl process enjoys time-inversion for some specific root systems. We first recall one of the main results of .
###### Proposition 5.2.
Let $`(X_t,t0)`$ be the generalized Dunkl process generated by (5.7), with $`(X_t^W,t0)`$ its radial part, i.e. the process confined to a Weyl chamber. Let $`R_+\{\alpha _1,\mathrm{},\alpha _l\}`$ for some $`l`$ be the corresponding positive root system and let $`(N_t^i,t0),i=1,\mathrm{},l`$ be independent Poisson processes of respective intensities $`\lambda (\alpha _i)`$. Then $`X_t`$ may be represented as $`Y_t^l`$, which is defined by induction as follows:
$$Y_t^0=X_t^W\mathrm{and}Y_t^i=\sigma _{\alpha _i}^{N_{A_t^i}^i}Y_t^{i1},i=1,\mathrm{},l,$$
where $`A_t^i={\displaystyle _0^t}{\displaystyle \frac{ds}{Y_s^{i1},\alpha _i^2}}`$.
The proof follows the argument in , while replacing $`k(\alpha _i)`$ by $`\lambda (\alpha _i)`$ appropriately.
From now on, let $`R_+\{\alpha _1,\mathrm{},\alpha _l\}`$ for some $`ln`$ be an orthogonal positive root system, that is $`\alpha _i,\alpha _j=2\delta _{ij}`$. For this particular root system, one can show that the generalized Dunkl process enjoys time-inversion of degree 1. We first prove the following absolute continuity relation:
###### Lemma 5.3.
Let $`(X_t^W,t0)`$ be the radial Dunkl process with infinitesimal generator
$$_k^Wf(x)=\frac{1}{2}\mathrm{\Delta }f(x)+\underset{i=1}{\overset{l}{}}k(\alpha _i)\frac{\alpha _i,f(x)}{\alpha _i,x}.$$
(5.12)
Fixed $`\nu \{1,\mathrm{},l\}`$. Let $`k^{}(\alpha )`$ be another coefficient function on the root system $`R_+`$ such that $`k^{}(\alpha _\nu )>k(\alpha _\nu )`$ and $`k^{}(\alpha _i)=k(\alpha _i)`$ for $`i\nu `$. Then, denoting $`\mathrm{P}_x^{(k)}`$ the law of the radial Dunkl process $`X_t^W`$ starting from $`x`$, we have
$`\mathrm{P}_x^{(k^{})}|__t`$ $`=`$ $`\left({\displaystyle \frac{\alpha _\nu ,X_t^W}{\alpha _\nu ,x}}\right)^{k^{}(\alpha _\nu )k(\alpha _\nu )}`$
$``$ $`\mathrm{exp}\left[{\displaystyle \frac{\left(k^{}(\alpha _\nu )\frac{1}{2}\right)^2\left(k(\alpha _\nu )\frac{1}{2}\right)^2}{2}}{\displaystyle _0^t}{\displaystyle \frac{ds}{\alpha _\nu ,X_s^W^2}}\right]\mathrm{P}_x^{(k)}|__t.`$
###### Proof.
Let $`k_0(\alpha )`$ be a coefficient such that $`k_0(\alpha _\nu )=\frac{1}{2}`$ for some fixed $`i\{1,\mathrm{},l\}`$. $`X_t^W`$ has the following martingale decomposition (see ):
$$X_t^W=x+B_t^{(k_0)}+\underset{i=1}{\overset{l}{}}k_0(\alpha _i)_0^t\frac{ds}{\alpha _i,X_s^W}\alpha _i$$
where $`B_t^{(k_0)}`$ is a $`(\mathrm{P}_x^{(k_0)},_t)`$-Brownian motion. Consider the local martingale
$$L_t^{(k^{})}=\mathrm{exp}\left(\left(k^{}(\alpha _\nu )\frac{1}{2}\right)_0^t\frac{\alpha _\nu ,dB_s^{(k_0)}}{\alpha _\nu ,X_s^W}\frac{\left(k^{}(\alpha _\nu )\frac{1}{2}\right)^2}{2}_0^t\frac{ds}{\alpha _\nu ,X_s^W^2}\right)$$
for some coefficient function $`k^{}(\alpha )`$ such that $`k^{}(\alpha _\nu )>\frac{1}{2}`$ and $`k^{}(\alpha _i)=k_0(\alpha _i)`$ for $`i\nu `$. The Itô formula for $`\mathrm{ln}\left(\alpha _\nu ,X_t^W\right)`$ combined with the orthogonality of the roots yields
$$L_t^{(k^{})}=\left(\frac{\alpha _\nu ,X_t^W}{\alpha _\nu ,x}\right)^{k^{}(\alpha _\nu )\frac{1}{2}}\mathrm{exp}\left(\frac{\left(k^{}(\alpha _\nu )\frac{1}{2}\right)^2}{2}_0^t\frac{ds}{\alpha _\nu ,X_s^W^2}\right).$$
Define the new law $`\mathrm{P}_x^{(k^{})}|__t=L_t^{(k^{})}\mathrm{P}_x^{(k_0)}|__t`$. By the Girsanov theorem,
$$B_t^{(k^{})}=B_t^{(k_0)}\left(k^{}(\alpha _\nu )\frac{1}{2}\right)_0^t\frac{ds}{\alpha _\nu ,X_s^W}$$
is a $`(\mathrm{P}_x^{(k^{})},_t)`$-Brownian motion and hence,
$$X_t^W=x+B_t^{(k^{})}+\underset{i=1}{\overset{l}{}}k^{}(\alpha _i)_0^t\frac{ds}{\alpha _i,X_s^W}\alpha _i$$
is a radial Dunkl process under $`(\mathrm{P}_x^{(k^{})},_t)`$.
Define $`k`$ as another coefficient on the root system that satisfies the conditions enunciated in the lemma. The absolute continuity relation is then a consequence of the successive application of the latter result to the indices $`k^{}`$ and $`k`$. ∎
Now as an application of Proposition 5.2, we prove the following:
###### Proposition 5.4.
Let $`R_+\{\alpha _1,\mathrm{},\alpha _l\}`$ be an orthogonal positive root system for $`ln`$. Then the generalized Dunkl process $`(X_t,t0)`$ generated by (5.7) enjoys the time-inversion property of degree 1.
###### Proof.
Using orthogonality of the roots, remark that
$$\alpha _i,\sigma _jx^2=\alpha _i,x\alpha _j,x\alpha _j^2=\alpha _i,x^2,$$
which implies in particular
$$\alpha _i,Y_t^i^2=\alpha _i,X_t^W^2,$$
so that the inductive representation of $`X_t`$ in Proposition 5.2 becomes
$$X_t=\underset{i=1}{\overset{l}{}}\sigma _{\alpha _i}^{N_{A_t^i}^i}X_t^W\mathrm{for}A_t^i=_0^t\frac{ds}{X_s^W,\alpha _i^2}.$$
The radial part of a Dunkl process enjoys the time-inversion property of degree 1. We need to show that the semigroup densities of the generalized Dunkl process are related to the semigroup densities of its radial parts. For $`fC^2(^n)`$,
$$\mathrm{E}_x\left[f(Y_t^i)\right]=\mathrm{E}_x\left[f(Y_t^{i1})1_{\{N_{A_t^i}^i\mathrm{is}\mathrm{even}\}}\right]+\mathrm{E}_x\left[f(\sigma _{\alpha _i}Y_t^{i1})1_{\{N_{A_t^i}^i\mathrm{is}\mathrm{odd}\}}\right].$$
Since $`\mathrm{P}(N_u^i\mathrm{is}\mathrm{even})=\frac{1}{2}\left(1+\mathrm{exp}(2\lambda (\alpha _i)u)\right)`$, we obtain
$`\mathrm{E}_x\left[f(Y_t^i)\right]`$ $`=`$ $`\mathrm{E}_x\left[f(Y_t^{i1}){\displaystyle \frac{1}{2}}\left(1+\mathrm{exp}(2\lambda (\alpha _i)A_t^i)\right)\right]`$
$`+`$ $`\mathrm{E}_x\left[f(\sigma _{\alpha _i}Y_t^{i1}){\displaystyle \frac{1}{2}}\left(1\mathrm{exp}(2\lambda (\alpha _i)A_t^i)\right)\right].`$
The expectation $`\mathrm{E}_x\left[f(X_t)\right]`$ can thus be evaluated by induction on $`i\{1,\mathrm{},l\}`$. It follows that the semigroup densities of $`X_t`$ can be expressed as the product of the semigroup densities of its radial parts times a function involving expectations of the form
$$\mathrm{E}_x^{(k)}\left[\mathrm{exp}\left(2\lambda (\alpha _\nu )A_t^\nu \right)|X_t^W=y\right],$$
for $`\nu \{1,\mathrm{},l\}`$. From Lemma 5.3,
$$\mathrm{E}_x^{(k)}\left[\mathrm{exp}\left(2\lambda (\alpha _\nu )A_t^\nu \right)|X_t^W=y\right]=\frac{p_t^{(k^{})}(x,y)}{p_t^{(k)}(x,y)}\left(\frac{\alpha _\nu ,y}{\alpha _\nu ,x}\right)^{k(\alpha _\nu )k^{}(\alpha _\nu )},$$
where $`k^{}(\alpha _i)=k(\alpha _i)`$ for $`i\nu `$ and $`k^{}(\alpha _\nu )=\frac{1}{2}+\sqrt{\left(k(\alpha _\nu )\frac{1}{2}\right)^2+4\lambda (\alpha _\nu )}`$. The form of the semigroup densities in (5.3) implies that the expectation is a ratio of Dunkl kernels,
$$\mathrm{E}_x^{(k)}\left[e^{2\lambda (\alpha _\nu )A_t^\nu }|X_t^W=y\right]=\frac{c_k}{c_k^{}}\frac{D_k^{}(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})}{D_k(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})}\frac{w_k^{}\left(\frac{y}{\sqrt{t}}\right)}{w_k\left(\frac{y}{\sqrt{t}}\right)}\left(\frac{\alpha _\nu ,\frac{y}{\sqrt{t}}}{\alpha _\nu ,\frac{x}{\sqrt{t}}}\right)^{(kk^{})(\alpha _\nu )},$$
which reduces to
$$\mathrm{E}_x^{(k)}\left[e^{2\lambda (\alpha _\nu )A_t^\nu }|X_t^W=y\right]=\frac{c_k}{c_k^{}}\frac{D_k^{}(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})}{D_k(\frac{x}{\sqrt{t}},\frac{y}{\sqrt{t}})}\left(\alpha _\nu ,\frac{y}{\sqrt{t}}\alpha _\nu ,\frac{x}{\sqrt{t}}\right)^{(k^{}k)(\alpha _\nu )}$$
by definition of $`k^{}(\alpha )`$. The conditional expectation thus satisfies condition 1 of Theorem 2.3. As a consequence, the conditions of Theorem 2.3 are satisfied for
$$\mathrm{\Phi }(x,y)D_{k,\lambda }(x,y),\theta (y)\omega _k(y),\rho (x)\frac{|x|^2}{2}$$
where $`D_{k,\lambda }(x,y)`$ is a generalized Dunkl kernel given explicitly in terms of radial Dunkl kernels and satisfying equivalent conditions (see (5.4)). ∎
#### 5.1.3 Jacobi-Dunkl processes
Gallardo et al. derived the Jacobi-Dunkl process as the hyperbolic analog of the one-dimensional Dunkl process. It is defined as the process generated by
$$^{(\alpha ,\beta )}f(x)=\frac{^2f(x)}{x^2}+\frac{A^{}(x)}{A(x)}\frac{f(x)}{x}+\frac{}{x}\left(\frac{A^{}(x)}{A(x)}\right)\left(\frac{f(x)f(x)}{2}\right),$$
(5.14)
where $`A(x)=\left(\mathrm{sinh}^2(x)\right)^{\alpha +\frac{1}{2}}\left(\mathrm{cosh}^2(x)\right)^{\beta +\frac{1}{2}}`$. From the expression of the semigroup densities developed in , this process does not enjoy the time-inversion property. Its radial part however corresponds to the Jacobi process of index $`(\alpha ,\beta )`$ on $`_+`$ (see ). The infinitesimal generator of the Jacobi process, expressed by
$$^{(\alpha ,\beta )}f(x)=\frac{1}{2}\frac{^2f(x)}{x^2}+\frac{A^{}(x)}{A(x)}\frac{f(x)}{x},$$
(5.15)
is in $`h`$-transform relationship with the Laplacian operator for $`h(x)=\sqrt{A(x)}`$. Since the Brownian motion enjoys the time-inversion property of degree 1, so does the Jacobi process by Theorem 2.3.
### 5.2 Matrix-valued processes
#### 5.2.1 Eigenvalue processes
Dyson in described the eigenvalues of a Hermitian Brownian motion as the joint evolution of independent Brownian motions conditioned never to collide (see also and ). It was further remarked that the process version of the Gaussian orthogonal ensemble does not admit such a representation for its eigenvalues. This work was extended by König and O’Connell in to the process version of the Laguerre ensemble, denominated the Laguerre process and defined as follows:
###### Definition 5.5.
Let $`B_t`$ be an $`n\times m`$ matrix with independent standard complex Brownian entries. The Laguerre process is the matrix-valued process defined by $`\{X_t=B_t^{}B_t,t0\}`$, where $`B_t^{}`$ is the transpose of $`B_t`$.
From , the eigenvalues of the Laguerre process evolve like $`m`$ independent squared Bessel processes conditioned never to collide. No such representation however exists for the case where the entries of $`B_t`$ are real Brownian motions, i.e. the Wishart process considered by Bru in .
The main result of is that both of the above mentioned eigenvalues processes (in the complex Brownian case) can be obtained as the $`h`$-transform of processes with $`m`$ independent components. The joint eigenvalues process is thus in $`h`$-transform relationship with a process that enjoys the time inversion property of degree 1 in the case of the $`m`$-dimensional Brownian motion and degree 2 in the case of the $`m`$-dimensional squared Bessel process, as made explicit in the following proposition:
###### Proposition 5.6.
Let $`p_t(x_i,y_i)(i=1,\mathrm{},m)`$ be the semigroup densities of squared Bessel processes (respectively Brownian motions), and let
$$h(x)=\underset{i<j}{\overset{m}{}}(x_jx_i)$$
(5.16)
for $`x=(x_1,\mathrm{},x_m)`$. Then, the semigroup densities of the joint eigenvalues process of the Laguerre process (respectively the Hermitian Brownian motion) are given by
$$\stackrel{~}{p}_t(x,y)=\frac{h(y)}{h(x)}det\left(p_t(x_i,y_j)\right)_{i,j=1}^m$$
(5.17)
with respect to the Lebesgue measure $`dy={\displaystyle \underset{j=1}{\overset{m}{}}}dy_j`$.
It follows immediately by Theorem 2.3 that the eigenvalues processes enjoy the time-inversion property. Moreover by Proposition 2.2, they yield the same process under time-inversion as the $`m`$-dimensional Brownian motion or the $`m`$-dimensional squared Bessel process respectively.
#### 5.2.2 Wishart processes
The Wishart process WIS($`\delta ,tI_m,\frac{1}{t}x`$), introduced by Bru in , is a continuous Markov process taking values in the space of real symmetric positive definite $`m\times m`$ matrices $`S_m^+`$. It is solution to the following stochastic differential equation:
$$dX_t=\sqrt{X_t}dB_t+dB_t^{}\sqrt{X_t}+\delta I_mdt,X_0=x,$$
(5.18)
where $`B_t`$ is an $`m\times m`$ matrix with Brownian entries and $`I_m`$ the identity matrix. Further results have been obtained in and . In , among other major findings about the Wishart process, the transition probability densities expressed with respect to the Lebesgue measure $`dy={\displaystyle \underset{ij}{}}(dy_{ij})`$ were derived in terms of generalized Bessel functions (we refer to the appendix for the definition):
$$p_t(x,y)=\frac{1}{(2t)^{\frac{m(m+1)}{2}}}\mathrm{exp}\left(\frac{1}{2t}Tr(x+y)\right)\left(\frac{det(y)}{det(x)}\right)^{\frac{\delta m1}{4}}\stackrel{~}{𝐈}_{\frac{\delta m1}{2}}\left(\frac{xy}{4t^2}\right),$$
(5.19)
for $`x,yS_m^+`$ and $`\delta >m1`$. From the shape of its densities, the Wishart process was stated in as an example of Markov processes enjoying the time-inversion property of degree 2. The hypothesis of Theorem 2.3 is indeed satisfied for $`n=\frac{1}{2}m(m+1)`$ and
$$\mathrm{\Phi }(x,y)\left(det(x)det(y)\right)^{\frac{\delta m1}{4}}\stackrel{~}{𝐈}_{\frac{\delta m1}{2}}\left(\frac{xy}{4}\right),$$
$$\theta (y)\frac{1}{2^n}(det(y))^{\frac{\delta m1}{2}},\rho (x)\frac{1}{2}Tr(x).$$
(5.20)
Next we use a skew-product representation, as for the Dunkl process, to elaborate on the Wishart process and derive a matrix-valued process with jumps. The skew-product representation allows the expression of the semigroup densities in terms of the Wishart transition probability densities.
###### Definition 5.7.
Let $`(N_t^{(\lambda )},t0)`$ be a Poisson process with intensity $`\lambda `$. Let $`(X_t,t0)`$ be a Wishart process WIS($`\delta ,tI_m,\frac{1}{t}x`$) independent of the Poisson process. The skew-Wishart process $`(X_t^{(\lambda )},t0)`$ is defined through the skew-product
$$X_t^{(\lambda )}=X_t(1)^{N_{A_t}^{(\lambda )}}$$
(5.21)
where $`A_t={\displaystyle _0^t}Tr(X_s^1)𝑑s`$.
###### Proposition 5.8.
The transition probability densities of the skew-Wishart process are related to the semigroup densities $`p_t(x,y)`$ of the Wishart process $`X_t`$ as follows
$`p_t^{(\lambda )}(x,y)`$ $`=`$ $`p_t(x,|y|)\{1_{\{yS_m^+\}}{\displaystyle \frac{1}{2}}(1+\left({\displaystyle \frac{\stackrel{~}{𝐈}_\nu ^{}}{\stackrel{~}{𝐈}_\nu }}\right)\left({\displaystyle \frac{xy}{4t^2}}\right))`$ (5.22)
$`+1_{\{yS_m^{}\}}{\displaystyle \frac{1}{2}}(1\left({\displaystyle \frac{\stackrel{~}{𝐈}_\nu ^{}}{\stackrel{~}{𝐈}_\nu }}\right)\left({\displaystyle \frac{xy}{4t^2}}\right))\}.`$
for $`\nu =\frac{\delta m1}{2}`$, $`\nu ^{}=\sqrt{\nu ^2+4\lambda }`$ and $`|y|=y(1_{\{yS_m^+\}}1_{\{yS_m^{}\}})`$.
###### Proof.
Let $`(𝐏_t)_{t>0}`$ be the semigroup of the skew-Wishart process. For $`x>0`$ and $`fC_c(M_m())`$,
$`𝐏_tf(x)`$ $`=`$ $`𝐄_x\left[f(X_t^{(\lambda )})\right]`$
$`=`$ $`𝐄_x\left[f(X_t)1_{\{N_{A_t}^{(\lambda )}\mathrm{is}\mathrm{even}\}}\right]+𝐄_x\left[f(X_t)1_{\{N_{A_t}^{(\lambda )}\mathrm{is}\mathrm{odd}\}}\right].`$
With $`\mathrm{P}(N_u^{(\lambda )}\mathrm{is}\mathrm{even})=\frac{1}{2}(1+\mathrm{exp}(2\lambda u))`$, we have
$$𝐏_tf(x)=𝐄_x\left[f(X_t)\frac{1}{2}(1+\mathrm{exp}(2\lambda A_t))\right]+𝐄_x\left[f(X_t)\frac{1}{2}(1\mathrm{exp}(2\lambda A_t))\right].$$
(5.23)
Let $`𝐐_x^{(\nu ^{})}`$ with $`\nu ^{}=\frac{\delta ^{}m1}{2}`$ denote the probability law of a Wishart process WIS$`(\delta ^{},tI,\frac{1}{t}x)`$, and $`𝐐_x^{(\nu )}`$ with $`\nu =\frac{\delta m1}{2}`$ the probability law of $`X_t`$. According to Theorem 1.2 (Remark 2.3) in , the probability laws are related as follows:
$$𝐐_x^{(\nu ^{})}|__t=\left(\frac{detX_t}{detx}\right)^{\frac{\nu ^{}\nu }{2}}\mathrm{exp}\left(\frac{\nu ^2\nu ^2}{2}_0^tTr(X_s^1)𝑑s\right)𝐐_x^{(\nu )}|__t,$$
from which we deduce
$$\frac{𝐩_t^{(\nu ^{})}(x,y)}{𝐩_t^{(\nu )}(x,y)}=\left(\frac{dety}{detx}\right)^{\frac{\nu ^{}\nu }{2}}𝐐_x^{(\nu )}\left[\mathrm{exp}\left(\frac{\nu ^2\nu ^2}{2}_0^tTr(X_s^1)𝑑s\right)|X_t=y\right].$$
Thus, from the expression of the semigroup densities in (5.19), we have
$$𝐄_x^{(\nu )}\left[\mathrm{exp}\left(2\lambda _0^tTr(X_s^1)𝑑s\right)|X_t=y\right]=\left(\frac{\stackrel{~}{𝐈}_{\sqrt{\nu ^\mathrm{𝟐}+\mathrm{𝟒}\lambda }}}{\stackrel{~}{𝐈}_\nu }\right)\left(\frac{xy}{4t^2}\right).$$
Combining the latter with (5.23) yields the semigroup densities for the skew-Wishart process. ∎
The skew-Wishart is an example of matrix-valued process with jumps that enjoys the time-inversion property of degree 2. Indeed, by setting
$`\mathrm{\Phi }(x,y)`$ $``$ $`(det(x)det(|y|))^{\frac{\nu }{2}}\{1_{\{yS_m^+\}}{\displaystyle \frac{1}{2}}(\stackrel{~}{𝐈}_\nu +\stackrel{~}{𝐈}_\nu ^{})\left({\displaystyle \frac{xy}{4}}\right)`$
$`+1_{\{yS_m^{}\}}{\displaystyle \frac{1}{2}}(\stackrel{~}{𝐈}_\nu \stackrel{~}{𝐈}_\nu ^{})\left({\displaystyle \frac{xy}{4}}\right)\},`$
$$\theta (y)\frac{1}{2^n}(det(|y|))^\nu ,\rho (x)\frac{1}{2}Tr(|x|),$$
(5.24)
the conditions of Theorem 2.3 are satisfied for $`\alpha =2`$.
## Appendix A Generalized hypergeometric functions
Using the notation in Muirhead , hypergeometric functions of matrix arguments are defined for a real symmetric $`m\times m`$ matrix $`X`$, $`a_i`$ and $`b_j\backslash \{0,\frac{1}{2},1,\mathrm{},\frac{m1}{2}\}`$ by
$${}_{p}{}^{}𝐅_{q}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_q;X)=\underset{k=0}{\overset{\mathrm{}}{}}\underset{\kappa }{}\frac{(a_1)_\kappa \mathrm{}(a_p)_\kappa }{(b_1)_\kappa \mathrm{}(b_q)_\kappa }\frac{C_\kappa (X)}{k!}$$
(A.1)
where the second summation is over all partitions $`\kappa =(k_1,\mathrm{},k_m),k_1\mathrm{}k_m0`$, of $`k=_{i=1}^mk_i,k!=k_1!\mathrm{}k_m!`$ and the generalized Pochhammer symbols are given by
$$(a)_\kappa =\underset{i=1}{\overset{m}{}}\left(a\frac{i1}{2}\right)_{k_i},(a)_k=a(a+1)\mathrm{}(a+k1),(a)_0=1.$$
$`C_\kappa (X)`$ is the zonal polynomial corresponding to $`\kappa `$, which is a symmetric, homogeneous polynomial of degree $`k`$ in the eigenvalues of $`X`$ that satisfies
$$C_\kappa (YX)=C_\kappa (\sqrt{Y}X\sqrt{Y})$$
(A.2)
for some $`YS_m^+`$. The function $`{}_{p}{}^{}𝐅_{q}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_q;YX)`$ thus makes sense. Finally, we define the generalized modified Bessel function by
$$\stackrel{~}{𝐈}_\nu (X)=\frac{\left(det(X)\right)^{\frac{\nu }{2}}}{𝚪_m\left(\nu +\frac{m+1}{2}\right)}{}_{0}{}^{}𝐅_{1}^{}(\nu +\frac{m+1}{2};X)$$
(A.3)
where the generalized gamma function is given as a product of the usual gamma functions,
$$𝚪_m\left(\alpha \right)=\pi ^{\frac{m(m1)}{4}}\underset{i=1}{\overset{m}{}}\mathrm{\Gamma }\left(\alpha \frac{i1}{2}\right)$$
(A.4)
for Re$`(\alpha )>\frac{m1}{2}`$. Note that the generalized modified Bessel for $`m=1`$ relates to the usual modified Bessel function $`I_\nu (x)`$ by $`\stackrel{~}{𝐈}_\nu (x)=I_\nu (2\sqrt{x})`$ (see ).
## Acknowledgements
The author wishes to thank Marc Yor for helpful discussions and references.
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# Bound states in Yukawa theory
## 1 Introduction
During the more than 50 years that have passed since Bethe and Salpeter formulated their famous equation , the calculation of relativistic bound states has proved to be one of the truely hard problems in quantum field theory. For most of the numerous proposals for its solution, a model theory of two scalar bosons interacting via the exchange of a third scalar has served as a first testing ground. In the case of the Bethe-Salpeter equation itself, it was shown by Wick and Cutkosky that this model theory has an analytical solution in the popular “ladder approximation” to the equation, in case that the exchanged boson is massless . It is hence quite surprising that for one of the most natural generalizations of this model where two spin 1/2 fermions interact through the exchange of a scalar boson, not a single consistent formalism has been devised for its solution to date.
In this introduction, we will give a brief review of the most important intents to calculate relativistic two-fermion bound states in this model, i.e., in Yukawa theory. It was clear even before Wick and Cutkosky’s solution of the purely scalar model that the case of fermionic constituents means a lot more to the Bethe-Salpeter equation than just a technical complication due to the inclusion of spin degrees of freedom : the fact that the kernel of the equation in the ladder approximation (after a Wick rotation) is not of Fredholm type in this case represents a fundamental difficulty for analytic as much as for numerical investigations. In the case of equal constituent masses and zero boson mass, it was shown by Goldstein (later corrected and improved upon by Green ), originally for the subspace of pseudoscalar bound states only, that the ladder approximation gives a *continuum* of coupling constants corresponding to a “tightly bound” state with zero mass, i.e., such that the binding energy completely compensates the masses of the constituents, while one expects a series of discrete values if a discrete spectrum of energy eigenvalues is to exist. The definite version of the argument is reviewed in Ref. (see the references there). In the case of Yukawa theory, this is essentially all that is known about the solutions of the Bethe-Salpeter equation.
The exchange of a scalar boson between spin 1/2 fermions has come to play an important role in the description of the nucleon-nucleon interaction in the context of one-boson exchange (OBE) models. Although the scalar boson exchange is used as an effective and approximate description of an actual two-pion exchange in this case, it is believed to give the dominant attractive contribution to the intermediate-range potential. In the numerous numerical calculations for the deuteron system within OBE models, the scalar boson exchange is always accompanied by the exchange of other mesons, most importantly pions and rho and omega mesons. For this type of calculations, the Bethe-Salpeter equation as well as several of its three-dimensional reductions have been employed, notably the Blankenbecler-Sugar-Logunov-Tavkhelidze (BSLT) equation and the Gross or spectator equation . The latter equations avoid problems with abnormal solutions, the one-body limit, and others that arise in the Bethe-Salpeter equation (in the ladder approximation) , see, however, Ref. . Among the vast literature on OBE models, we mention the work of Tjon and collaborators who use the Bethe-Salpeter equation, the work of Gross and collaborators who base their analysis on the Gross equation, and the quite complete “Bonn model” by Holinde, Machleidt, et al. reviewed in Ref. , which employs the Bloch-Horowitz scheme for the bound state calculations, a formalism that is *not* based on the Bethe-Salpeter equation.
In all these approaches, a momentum cutoff is introduced in order to make the equations well-defined. It is interpreted phenomenologically as representing the spatial extension of the nucleons and mesons. Although we have not found a corresponding calculation in the published literature, it appears plausible in view of the singular nature of the Bethe-Salpeter kernel in the ladder approximation that the introduction of a momentum cutoff would also be necessary in the case of a pure Yukawa interaction, and that the corresponding numerical results depend on the value of the cutoff. Note, however, the work of Gari and collaborators who dispense with a momentum cutoff and rather determine the nucleon and meson structures self-consistently through loop corrections within the same nucleon-meson model. It is perhaps not a coincidence that their approach employs the Okubo transformation technique , which is unrelated to the Bethe-Salpeter equation but very similar to the formalism proposed in the present paper.
More recently, calculations of bound states in Yukawa theory have been attempted in light-front quantization, both in a Tamm-Dancoff approximation and in the so-called covariant light-front dynamics . Just as in the equal-time quantized theory, an additional momentum cutoff has to be introduced for the one-boson exchange, and the solutions for the bound states turn out to depend on this cutoff. It was pointed out by van Iersel and Bakker that the light front ladder approximation is incomplete as long as instantaneous terms are not taken into account. One possibility to include such terms was presented in Ref. and may help to solve the renormalization problem in the future.
In this paper, we propose to use a generalization of the Gell-Mann–Low theorem for the calculation of relativistic bound states. After the obligatory application to the purely scalar model in a previous publication , we now turn to Yukawa theory. Surprisingly, the application of the generalized Gell-Mann–Low theorem turns out to be straightforward and, except for the necessary technical complications due to the spin degrees of freedom, completely analogous to the scalar case. No inconsistencies whatsoever have been found, and no necessity for the introduction of a cutoff (for the interaction between the fermions) arises. For the numerical calculation of the bound state spectrum, we will focus on the case of a massless exchanged boson, for the following reasons: first, the massless case is the most singular one, and if the formalism works in this case, it is expected to be applicable to the case of a massive exchanged boson as well. Second, the solutions in the non-relativistic limit are known analytically in this case (Coulomb potential) and our numerical solutions can be compared against them. In particular, one would like to find the degeneracies characteristic of the non-relativistic limit. Finally, we plan to extract the lowest-order fine and hyperfine structure from our effective Hamiltonian in the near future, and this can be done analytically only in the massless (Coulomb) case.
The organization of the paper is as follows: in the next section, we introduce the generalization of the Gell-Mann–Low theorem on which our approach is based, and comment on aspects relevant to the present work. We also present, in the same section, the effective Hamiltonians in the zero-, one-, and two-fermion sectors which result from the application of the generalized Gell-Mann–Low theorem to lowest non-trivial order. They describe the renormalization of the vacuum energy, the fermion mass renormalization, and the effective potential for the two-fermion dynamics, respectively. In Section 3, we perform a series of non-trivial analytical checks on the results obtained in Section 2. In particular, we investigate the non-relativistic and one-body limits of the effective Schrödinger equation, replace fermions by antifermions, and consider the case of identical constituents. All the properties of our effective Hamiltonian description turn out to be in accord with physical expectations. We present numerical bound state solutions in Section 4, for fine structure constants between zero and one and arbitrary ratios of the constituent masses. The eight lowest-lying states are calculated in each case, including states with non-zero relative orbital angular momentum and mixtures of spin singlet and triplet states. The non-relativistic, one-body, and equal-mass limits are discussed in detail. Section 5 contains our conclusions. There are also four appendices. Appendices A and B are concerned with the appearing loop corrections and their regularization in different schemes from a Hamiltonian, not manifestly covariant perspective. They prepare the ground for future calculations at higher orders in the expansion of the effective Hamiltonian, where regularization and renormalization are expected to become central issues. Appendices C and D present several formulas which are used in the analytical separation of angular and spin variables in the effective Schrödinger equation.
## 2 The Bloch-Wilson Hamiltonian
The physical idea behind the Bloch-Wilson or effective Hamiltonian is very similar to the Born-Oppenheimer approximation: the integration over the “light” degrees of freedom generates the dynamics for the “heavy” degrees of freedom (the constituents in our case), where the “light” degrees of freedom are given by the so-called interacting “virtual clouds” around the constituents. The virtual clouds are, in turn, created by the constituents as described here through an adiabatic process. Technically, a generalization of the Gell-Mann–Low theorem is used which we will briefly describe in the following (for further details, see Ref. ).
The adiabatic evolution operator $`U_ϵ`$ provides a map from an eigenstate of the free Hamiltonian to an exact eigenstate of the full interaction Hamiltonian, *if its application to the free eigenstate, after a suitable normalization, is well defined*. This is the content of the original theorem by Gell-Mann and Low . To this end, the full Hamiltonian $`H`$ is split into a free part $`H_0`$ and an interacting part $`H_1`$. The adiabatic damping is implemented by replacing $`H=H_0+H_1`$ through
$`H(t)=H_0+e^{ϵ|t|}H_1.`$ (1)
An eigenstate $`|\varphi `$ of $`H_0=lim_t\mathrm{}H(t)`$ is then evolved according to the time-dependent Schrödinger equation with Hamiltonian $`H(t)`$ from $`t\mathrm{}`$ to the state $`U_ϵ|\varphi `$ at $`t=0`$. According to the adiabatic theorem, in the limit $`ϵ0`$ the state $`U_ϵ|\varphi `$ is an eigenstate of $`H(t=0)=H`$. A possible infinite phase (see, for example, ) is taken care of by imposing the normalization condition
$`\varphi |\psi =1,`$ (2)
so that
$`|\psi ={\displaystyle \frac{U_ϵ|\varphi }{\varphi |U_ϵ|\varphi }}.`$ (3)
As a generalization of this theorem, it can be shown that the adiabatic evolution operator also maps $`H_0`$-invariant subspaces $`\mathrm{\Omega }_0`$, $`H_0\mathrm{\Omega }_0\mathrm{\Omega }_0`$, to $`H`$-invariant subspaces $`\mathrm{\Omega }`$, $`H\mathrm{\Omega }\mathrm{\Omega }`$, *if its application to the $`H_0`$-invariant subspace, after a suitable normalization, is well defined*. When a subspace is invariant under a hermitian operator, this operator can be diagonalized in the subspace, in other words, the subspace is a direct sum of eigenspaces of the operator, in this case the Hamiltonian. Since the only important property of the adiabatic evolution operator in this context is its $`H`$-invariant image, there is still an infinity of possible “normalizations” of the operator that map between the same subspaces $`\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }`$. Here we choose the normalization condition
$`P_0U_{\text{BW}}=\mathrm{𝟏}_{\mathrm{\Omega }_0},`$ (4)
where $`P_0`$ is the orthogonal projector to $`\mathrm{\Omega }_0`$ and $`\mathrm{𝟏}_{\mathrm{\Omega }_0}`$ the identity operator in $`\mathrm{\Omega }_0`$. $`U_{\text{BW}}`$ is the “Bloch-Wilson operator”, the normalized version of the adiabatic evolution operator given explicitly by
$`U_{\text{BW}}=U_ϵ(P_0U_ϵ)^1.`$ (5)
Eq. (4) naturally generalizes the normalization condition (2) in the original Gell-Mann–Low theorem. The subspace $`\mathrm{\Omega }_0`$ has to be small enough in order that it is possible in practice to diagonalize the Hamiltonian in $`\mathrm{\Omega }=U_{\text{BW}}\mathrm{\Omega }_0`$, at least numerically. On the other hand, it has to be large enough for the Bloch-Wilson operator to be well defined. This is in principle a subtle issue and may depend on the normalization condition chosen. We will limit ourselves here to show that everything is well defined for a natural choice of the subspace (for the field theory and the perturbative order considered here, see the discussion below). In the case of a two-constituent bound state, a natural choice for the $`H_0`$-invariant subspace $`\mathrm{\Omega }_0`$ is the space of all states of the two constituents as free particles. This will be mapped by the normalized adiabatic evolution operator to the $`H`$-invariant subspace $`\mathrm{\Omega }`$, a direct sum of eigenstates of the full Hamiltonian, which is expected to coincide with the space of all physical two-particle states, scattering as well as bound states.
A further simplification is obtained by realizing that the Bloch-Wilson operator effects a similarity transformation between the subspaces $`\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }`$. Since the free subspace $`\mathrm{\Omega }_0`$ is usually much easier to work in, it is convenient to similarity transform the full Hamiltonian back to this subspace, thereby defining an effective or Bloch-Wilson Hamiltonian $`H_{\text{BW}}`$ in $`\mathrm{\Omega }_0`$,
$`H_{\text{BW}}=P_0HU_{\text{BW}},`$ (6)
where we have taken into account that $`P_0|_\mathrm{\Omega }`$ is the inverse map to $`U_{\text{BW}}`$ as a consequence of the normalization (4). In the above example, the effective Hamiltonian acts on the subspace of all states of the free constituents. Its eigenvalues are exactly equal to the eigenvalues of the full Hamiltonian (in the corresponding subspace $`\mathrm{\Omega }`$), and there is a one-to-one relation between the corresponding eigenstates. Mathematically, the time-independent Schrödinger equation for the effective Hamiltonian takes the form of a non-relativistic Schrödinger equation for the two constituents, with the relativistic expression for the kinetic energies and a complicated (generally non-local and non-hermitian) interaction term. With the normalization chosen, the eigenstates $`|\varphi `$ of the effective Hamiltonian have the immediate meaning of wavefunctions for the constituents:
$`|\varphi =P_0(U_{\text{BW}}|\varphi ),`$ (7)
hence $`|\varphi `$ is the two-particle component of the complicated exact eigenstate of $`H`$, $`U_{\text{BW}}|\varphi `$.
Obviously, the Bloch-Wilson operator $`U_{\text{BW}}`$ cannot be determined exactly in any interesting practical application. However, the adiabatic evolution operator $`U_ϵ`$, and hence $`U_{\text{BW}}=U_ϵ(P_0U_ϵ)^1`$, have a well-known perturbative expansion, the Dyson series:
$`U_ϵ={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i)^n}{n!}}{\displaystyle _{\mathrm{}}^0}𝑑t_1\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑t_ne^{ϵ(|t_1|+\mathrm{}+|t_n|)}T[H_1(t_1)\mathrm{}H_1(t_n)],`$ (8)
where
$`H_1(t)=e^{iH_0t}H_1e^{iH_0t}.`$ (9)
Every order in this expansion determines an approximation to the effective Hamiltonian and consequently of its eigenvalues and eigenstates. To compare with, in the Bethe-Salpeter approach an infinite part of the Dyson series has to be summed up to find an approximation to the bound states. The sum is performed in practice implicitly by solving an integral equation (the homogeneous Bethe-Salpeter equation) with a perturbative approximation to its kernel, and the counterpart in our Hamiltonian approach is the solution of the (approximated) effective Schrödinger equation. All of the above is exemplified in the application to a simple scalar model in .
Finally, let us comment on the sense in which the adiabatic evolution operator, or rather its normalized version $`U_{\text{BW}}`$, is well-defined in $`\mathrm{\Omega }_0`$. In the proof of the generalized Gell-Mann–Low theorem, as well as in the proof of the original theorem, the adiabatic evolution operator is treated as a formal power series (in $`H_1`$) given by the Dyson series. Consequently, $`U_{\text{BW}}`$ is considered well-defined if every term of the corresponding series is well-defined (finite) disregarding convergence properties of the series as a whole. However, the UV divergencies of the usual covariant perturbation theory cannot be escaped: they appear in the context of the Bloch-Wilson Hamiltonian (in the most straightforward approach) as divergencies of the three-dimensional momentum integrals for large momenta. In the scalar model considered before it was shown that these divergencies can be dealt with (to the lowest non-trivial order considered there) rather easily, and, in fact, in the same way as in covariant perturbation theory. On the other hand, divergencies of the IR type are considered to potentially invalidate the existence of the Bloch-Wilson operator (for the specific subspace $`\mathrm{\Omega }_0`$ chosen). In the time-independent version presented in , which is obtained by performing the time integrals in the Dyson series, the IR divergencies arise from vanishing energy denominators. They can usually be avoided by appropriately enlarging the subspace $`\mathrm{\Omega }_0`$. However, the usefulness of the entire approach rests on the possibility of choosing a subspace $`\mathrm{\Omega }_0`$ which is manageably small, ideally the space of all states of the two constituents as free particles in the example mentioned above. Here, as for the scalar model before, we take a pragmatic approach: we will be content with the fact that no IR divergence arises to the perturbative order and for the choice of $`\mathrm{\Omega }_0`$ considered.
We will now apply these ideas to Yukawa theory. To simplify matters, we consider two different fermion species, $`A`$ and $`B`$, which interact with a scalar boson. Hence the Hamiltonian of the theory consists of the free Hamiltonians for the fermions $`A`$ and $`B`$ and the scalar boson, and the (normal-ordered) interaction term
$`H_1=g{\displaystyle d^3x}\mathbf{:}\left[\overline{\psi }_A(𝐱)\psi _A(𝐱)+\overline{\psi }_B(𝐱)\psi _B(𝐱)\right]\phi (𝐱)\mathbf{:}.`$ (10)
We are interested in bound states of one fermion $`A`$ and one fermion $`B`$. Other bound states, of identical fermions or fermion and antifermion (e.g., in a theory with only one fermion species), are simply related to bound states of type $`AB`$ and can be treated by exactly the same method. They will be discussed in some detail in the next section. The application of the generalized Gell-Mann–Low theorem to lowest non-trivial order follows in very close analogy the purely scalar case . In particular, for a subspace $`\mathrm{\Omega }_0`$ with a fixed number of (free) particles $`P_0H_1P_0=0`$, and the effective Hamiltonian becomes to lowest non-trivial order
$`H_{BW}=H_0P_0i{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}P_0H_1(0)H_1(t)P_0+𝒪(H_1^4)`$ (11)
(the limit $`ϵ0`$ being understood). For the rest of this chapter, we will discuss in detail the cases where $`\mathrm{\Omega }_0`$ is the subspace of zero, one (fermionic), and two ($`AB`$) particles. We will go through the essential steps briefly, present the corresponding results, and take the opportunity to comment on several issues that have not received due attention in . For illustrative purposes, it may be helpful to have a look at the diagrams presented there which are the same in the present case, only that spin labels have to be added to the external fermion lines.
### 2.1 Vacuum state and zero-point energy renormalization
Beginning with the vacuum state, we consider the Bloch-Wilson Hamiltonian for the subspace $`\mathrm{\Omega }_0=|0`$, or, equivalently, for the orthogonal projector on $`\mathrm{\Omega }_0`$, $`P_0=|00|`$. In this case, the application of the generalized version of the Gell-Mann–Low theorem is equivalent to the original Gell-Mann–Low theorem. The result for the Bloch-Wilson Hamiltonian is $`H_{\text{BW}}=E_VP_0`$, where (by Wick’s theorem)
$$\begin{array}{c}E_VE_0\hfill \\ \hfill =ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\mathrm{\Delta }_F(0t,𝐱𝐱^{})\text{tr}\left[S_F^A(0t,𝐱𝐱^{})S_F^A(t0,𝐱^{}𝐱)\right]\\ \hfill +ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\mathrm{\Delta }_F(0t,𝐱𝐱^{})\text{tr}\left[S_F^B(0t,𝐱𝐱^{})S_F^B(t0,𝐱^{}𝐱)\right].\end{array}$$
(12)
Here $`E_0`$ is the vacuum energy of the free theory and $`S_F^{A,B}`$ and $`\mathrm{\Delta }_F`$ are the fermionic (for particles $`A,B`$) and the bosonic Feynman propagators. In our conventions,
$`S_F^A(t,𝐱)`$ $`=i{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{p\gamma +m_A}{p^2m_A^2+iϵ}e^{i(p_0t𝐩𝐱)}}`$ (13)
$`={\displaystyle \frac{d^3p}{(2\pi )^3\mathrm{\hspace{0.17em}2}E_𝐩^A}\left[\theta (t)\left(p\gamma +m_A\right)e^{i(E_𝐩^At𝐩𝐱)}\theta (t)\left(p\gamma m_A\right)e^{i(E_𝐩^At𝐩𝐱)}\right]}`$ (14)
in the covariant and non-covariant representation, respectively. We are using the following shorthands for the relativistic kinetic energies:
$`E_𝐩^{A,B}=\sqrt{m_{A,B}^2+𝐩^2},\omega _𝐩=\sqrt{\mu ^2+𝐩^2}`$ (15)
($`\mu `$ is the boson mass). The bosonic propagator $`\mathrm{\Delta }_F(t,𝐱)`$ is simply given by the $`m_A`$-coefficient of $`S_F^A(t,𝐱)`$ (and the substitution $`m_A\mu `$). In Eq. (12), note the opposite sign as compared to the bosonic case .
We can use time translation invariance of the integrands on the r.h.s. of Eq. (12) to move the vertices from times $`t<0`$ and $`t=0`$ to $`t=0`$ and $`t>0`$. Mathematically, this corresponds to replacing $`t`$ by $`t`$ and exchanging $`𝐱`$ and $`𝐱^{}`$, just as we did for the one-particle states in the scalar case . If we, furthermore, use the invariance of the bosonic propagator under the change of the direction of propagation and apply these manipulations to half the r.h.s. of Eq. (12), we arrive (in the limit $`ϵ0`$) at the usual expression for the one-loop correction to the vacuum energy in covariant perturbation theory ,
$`E_VE_0`$ $`=i{\displaystyle \frac{g^2}{2}}{\displaystyle d^4xd^4x^{}\delta (x^0)\mathrm{\Delta }_F(xx^{})\text{tr}\left[S_F^A(xx^{})S_F^A(x^{}x)\right]}`$
$`+i{\displaystyle \frac{g^2}{2}}{\displaystyle d^4xd^4x^{}\delta (x^0)\mathrm{\Delta }_F(xx^{})\text{tr}\left[S_F^B(xx^{})S_F^B(x^{}x)\right]}.`$ (16)
The appearance of $`\delta (x^0)`$ eliminates a factor
$`2\pi \delta (0)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t`$ (17)
(due to time translation invariance) in the result, and leaves us with
$`(2\pi )^3\delta ^{(3)}(0)={\displaystyle d^3x}`$ (18)
(due to spatial translation invariance), to be interpreted as the volume of space.
We can immediately express Eqs. (12) and (16) in momentum space, by use of the non-covariant and covariant expressions for the Feynman propagators in momentum space, Eqs. (14) and (13), respectively. The results are
$`E_VE_0`$
$`=[g^2{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^3}}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^A\mathrm{\hspace{0.17em}2}E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩+𝐩^{}}}}{\displaystyle \frac{4\left(E_𝐩^AE_𝐩^{}^A+𝐩𝐩^{}+m_A^2\right)}{E_𝐩^A+E_𝐩^{}^A+\omega _{𝐩+𝐩^{}}}}`$
$`+g^2{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^3}}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^B\mathrm{\hspace{0.17em}2}E_𝐩^{}^B\mathrm{\hspace{0.17em}2}\omega _{𝐩+𝐩^{}}}}{\displaystyle \frac{4\left(E_𝐩^BE_𝐩^{}^B+𝐩𝐩^{}+m_B^2\right)}{E_𝐩^B+E_𝐩^{}^B+\omega _{𝐩+𝐩^{}}}}](2\pi )^3\delta ^{(3)}(0)`$
$`=\{{\displaystyle \frac{g^2}{2}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}{\displaystyle \frac{d^4p^{}}{(2\pi )^4}}{\displaystyle \frac{4(pp^{}+m_A^2)}{[(pp^{})^2\mu ^2+iϵ][p^2m_A^2+iϵ][p_{}^{}{}_{}{}^{2}m_A^2+iϵ]}}`$
$`+{\displaystyle \frac{g^2}{2}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}{\displaystyle \frac{d^4p^{}}{(2\pi )^4}}{\displaystyle \frac{4(pp^{}+m_B^2)}{[(pp^{})^2\mu ^2+iϵ][p^2m_B^2+iϵ][p_{}^{}{}_{}{}^{2}m_B^2+iϵ]}}\}(2\pi )^3\delta ^{(3)}(0).`$ (19)
The equivalence between the two expressions can, of course, be established directly by integrating over $`p_0`$ and $`p_0^{}`$ in the manifestly covariant expression. However, there is a subtlety involved in the integration which may be of interest for future calculations at higher orders and will be described in detail in Appendix A. The results (19) are highly UV divergent and will be treated as formal expressions only. It is, however, important to remark that they are IR finite in the limit $`\mu 0`$ (we are here only interested in massive constituents $`m_A,m_B0`$).
### 2.2 One-fermion states and mass renormalization
We will now turn to the one-fermion states, considering for concreteness fermions of type $`A`$. The $`H_0`$-invariant subspace considered is, correspondingly,
$`\mathrm{\Omega }_0=\text{span}\left\{|𝐩_A,r|𝐩_A^3,r=1,2\right\},`$ (20)
where $`|𝐩_A,r`$ stands for a state of one fermion $`A`$ with 3-momentum $`𝐩_A`$ and spin orientation $`r`$ (we do not fix the basis to be used in spin space yet). The one-fermion states are normalized in a non-covariant fashion,
$`𝐩_A,r|𝐩_A^{},s=(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rs}.`$ (21)
The application of the generalized Gell-Mann–Low theorem to lowest non-trivial order yields the following matrix elements of the Bloch-Wilson Hamiltonian in $`\mathrm{\Omega }_0`$,
$$\begin{array}{c}𝐩_A,r|H_{\text{BW}}|𝐩_A^{},s=\left(E_V+E_{𝐩_A}^A\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rs}\hfill \\ \hfill ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)S_F^A(0t,𝐱𝐱^{})\psi _{𝐩_A^{},s}^A(t,𝐱^{})\right]\mathrm{\Delta }_F(0t,𝐱𝐱^{})\\ \hfill ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱^{})S_F^A(t0,𝐱^{}𝐱)\psi _{𝐩_A^{},s}^A(0,𝐱)\right]\mathrm{\Delta }_F(t0,𝐱^{}𝐱),\end{array}$$
(22)
where the fermion wave functions $`\psi _{𝐩,r}^A(t,𝐱)`$ are given by
$`\psi _{𝐩,r}^A(t,𝐱)={\displaystyle \frac{u_A(𝐩,r)}{\sqrt{2E_𝐩^A}}}e^{iE_𝐩^At+i𝐩𝐱}`$ (23)
with the Dirac spinors normalized to
$`\overline{u}_A(𝐩,r)u_A(𝐩,s)=2m_A\delta _{rs}.`$ (24)
Using 3-momentum conservation at the vertices and, consequently, $`E_{𝐩_A}^A=E_{𝐩_A^{}}^A`$, the vertices in the last integral in Eq. (22) can be translated from times $`t<0`$ and $`t=0`$ to $`t=0`$ and $`t>0`$, just as we did in the scalar case and also for the corrections to the vacuum energy above. Equation (22) can then be written in the covariant form
$$\begin{array}{c}𝐩_A,r|H_{\text{BW}}|𝐩_A^{},s=\left(E_V+E_{𝐩_A}^A\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rs}\hfill \\ \hfill ig^2d^4xd^4x^{}\delta (x^0)\left[\overline{\psi }_{𝐩_A,r}^A(x)S_F^A(xx^{})\psi _{𝐩_A^{},s}^A(x^{})\right]\mathrm{\Delta }_F(xx^{}),\end{array}$$
(25)
where the integral over $`d^3x`$ serves to implement 3-momentum conservation.
Finally, then, the matrix elements of $`H_{\text{BW}}`$ take the form
$$\begin{array}{c}𝐩_A,r|H_{\text{BW}}|𝐩_A^{},s=\left(E_V+E_{𝐩_A}^A\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rs}\hfill \\ \hfill +\frac{1}{2E_{𝐩_A}^A}\left[\overline{u}_A(𝐩_A,r)G(𝐩_A)u_A(𝐩_A,s)\right](2\pi )^3\delta (𝐩_A𝐩_A^{}),\end{array}$$
(26)
where $`G(𝐩_A)`$ is written in momentum space by use of the non-covariant and covariant expressions (14) and (13) for the Feynman propagators in momentum space in Eqs. (22) and (25), respectively, to give the following equivalent expressions
$`G(𝐩)`$ $`=g^2{\displaystyle \frac{d^3p^{}}{(2\pi )^3}\frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}\frac{E_𝐩^{}^A\gamma _0𝐩^{}𝜸+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}E_𝐩^A}}`$
$`g^2{\displaystyle \frac{d^3p^{}}{(2\pi )^3}\frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}\frac{E_𝐩^{}^A\gamma _0𝐩^{}𝜸+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}+E_𝐩^A}}`$
$`=ig^2{\displaystyle \frac{d^4p^{}}{(2\pi )^4}\frac{p^{}\gamma +m_A}{[(pp^{})^2\mu ^2+iϵ][p_{}^{}{}_{}{}^{2}m_A^2+iϵ]}}|_{p_0=E_𝐩^A}.`$ (27)
The equivalence between the two expression is shown directly in Appendix A by performing the integration over $`p_0^{}`$ in the manifestly covariant form.
We will now show that the contribution to $`H_{\text{BW}}`$ in the second line of Eq. (26) can be absorbed into a renormalization of the mass $`m_A`$. To this end, consider the covariant expression for $`G(𝐩)`$ in Eq. (27): from covariance arguments, one has that $`G(𝐩)`$ is of the form
$`G(𝐩)=\left[G_1(p^2)p\gamma +G_0(p^2)m_A\right]_{p_0=E_𝐩^A},`$ (28)
hence by use of the Dirac equation
$`\overline{u}_A(𝐩,r)G(𝐩)u_A(𝐩,s)`$ $`=\overline{u}_A(𝐩,r)\left[G_1(m_A^2)m_A+G_0(m_A^2)m_A\right]u_A(𝐩,s)`$
$`=2m_A^2\left[G_1(m_A^2)+G_0(m_A^2)\right]\delta _{rs}.`$ (29)
Consequently, we define
$`\mathrm{\Delta }m_A^2=2m_A^2\left[G_1(m_A^2)+G_0(m_A^2)\right]`$ (30)
in order to write
$`𝐩_A,r|H_{\text{BW}}|𝐩_A^{},s=\left(E_V+E_{𝐩_A}^A+{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2E_{𝐩_A}^A}}\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rs}.`$ (31)
We are now in a position to perform the mass renormalization, *completely within the Hamiltonian framework*. First, define the renormalized or physical mass through
$`\left[E_V+E_{𝐩_A}^A+{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2E_{𝐩_A}^A}}\right]_{𝐩_A=0}=E_V+M_A,`$ (32)
then
$`M_A=m_A+{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2m_A}}+𝒪(g^4).`$ (33)
We can use Eq. (33) to express $`m_A`$ in terms of $`M_A`$ in Eq. (31), in particular in $`E_{𝐩_A}^A`$, taking into account that $`\mathrm{\Delta }m_A^2`$ is of order $`g^2`$. Working consequently to order $`g^2`$, we arrive at
$`E_V+E_{𝐩_A}^A+{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2E_{𝐩_A}^A}}=E_V+\sqrt{M_A^2+𝐩_A^2}+𝒪(g^4).`$ (34)
Remarkably, through the mass renormalization (to the order presently considered) we have obtained an expression for the energy which is exactly covariant, in contradistinction to Eq. (31).
Now, the arguments to arrive at the manifestly covariant expression in Eq. (27) (see also Appendix A) and from there to the crucial equation (28), are formal in the sense that the (unregularized) expression for $`G(𝐩)`$ is UV divergent. We show in Appendix B that a careful derivation using different regularization schemes leads to the same results. Observe, again, the absence of IR divergencies in the limit $`\mu 0`$ in the expressions (27) (or in the corresponding explicit expressions presented in Appendix B).
### 2.3 Two-fermion states and effective potential
Finally, in order to obtain the effective Schrödinger equation for $`AB`$ bound states, we consider the Bloch-Wilson Hamiltonian for the subspace
$`\mathrm{\Omega }_0=\text{span}\left\{|𝐩_A,r;𝐩_B,s|𝐩_A,𝐩_B^3,r,s=1,2\right\}`$ (35)
of all (free) states of one fermion $`A`$ and one fermion $`B`$, non-covariantly normalized as in Eq. (21). Then the matrix elements of the Bloch-Wilson Hamiltonian to lowest non-trivial order turn out to be
$`𝐩_A,r;𝐩_B,s|H_{\text{BW}}|𝐩_A^{},r^{};𝐩_B^{},s^{}`$
$`=\left(E_V+\sqrt{M_A^2+𝐩_A^2}+\sqrt{M_B^2+𝐩_B^2}\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rr^{}}(2\pi )^3\delta (𝐩_B𝐩_B^{})\delta _{ss^{}}`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_B,s}^B(0,𝐱)\psi _{𝐩_B^{},s^{}}^B(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱^{})\psi _{𝐩_A^{},r^{}}^A(t,𝐱^{})\right]`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)\psi _{𝐩_A^{},r^{}}^A(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_B,s}^B(t,𝐱^{})\psi _{𝐩_B^{},s^{}}^B(t,𝐱^{})\right],`$ (36)
where $`M_A`$ and $`M_B`$ are the renormalized masses defined as in Eq. (32), and the wave functions $`\psi _{𝐩,r}^{A,B}(t,𝐱)`$ have been introduced in Eq. (23).
Due to the non-conservation of the perturbative energies, generally $`E_{𝐩_A}^A+E_{𝐩_B}^BE_{𝐩_A^{}}^A+E_{𝐩_B^{}}^B`$ for $`𝐩_A+𝐩_B=𝐩_A^{}+𝐩_B^{}`$, and the corresponding lack of time translation invariance, the expressions with Feynman propagators in Eq. (36) can*not* be converted to on-shell Feynman diagrams. It has not been much emphasized in that the mixing of states with different perturbative energies is essential for the formation of bound states, because it is imperative for the localization of the constituents in relative position space: consider, e.g., states with total momentum $`𝐩_A+𝐩_B=0`$, then a continuous superposition of states with different relative momenta $`𝐩=𝐩_A`$ is necessary to obtain a wavefunction of finite extension in relative position space.
Equation (36) leads to the following effective Schrödinger equation,
$$\begin{array}{c}\left(\sqrt{M_A^2+𝐩_A^2}+\sqrt{M_B^2+𝐩_B^2}\right)\varphi (𝐩_A,r;𝐩_B,s)\hfill \\ \hfill +\underset{r^{},s^{}=1}{\overset{2}{}}\frac{d^3p_A^{}}{(2\pi )^3}\frac{d^3p_B^{}}{(2\pi )^3}𝐩_A,r;𝐩_B,s|V|𝐩_A^{},r^{};𝐩_B^{},s^{}\varphi (𝐩_A^{},r^{};𝐩_B^{},s^{})\\ \hfill =\left(EE_V\right)\varphi (𝐩_A,r;𝐩_B,s)\end{array}$$
(37)
for the two-particle wave function in momentum space
$`\varphi (𝐩_A,r;𝐩_B,s)=𝐩_A,r;𝐩_B,s|\varphi ,|\varphi \mathrm{\Omega }_0.`$ (38)
Through the use of the non-covariant representations (14) of the Feynman propagators in momentum space, the effective potential can be written as
$$\begin{array}{c}𝐩_A,r;𝐩_B,s|V|𝐩_A^{},r^{};𝐩_B^{},s^{}\hfill \\ \hfill =\frac{g^2}{\sqrt{2E_{𝐩_A}^A\mathrm{\hspace{0.17em}2}E_{𝐩_B}^B\mathrm{\hspace{0.17em}2}E_{𝐩_A^{}}^A\mathrm{\hspace{0.17em}2}E_{𝐩_B^{}}^B}}\frac{1}{2\omega _{𝐩_A𝐩_A^{}}}\left(\frac{1}{E_{𝐩_A}^A+\omega _{𝐩_A𝐩_A^{}}E_{𝐩_A^{}}^A}+\frac{1}{E_{𝐩_B}^B+\omega _{𝐩_B𝐩_B^{}}E_{𝐩_B^{}}^B}\right)\\ \hfill \times \left[\overline{u}_A(𝐩_A,r)u_A(𝐩_A^{},r^{})\right]\left[\overline{u}_B(𝐩_B,s)u_B(𝐩_B^{},s^{})\right](2\pi )^3\delta (𝐩_A+𝐩_B𝐩_A^{}𝐩_B^{}).\end{array}$$
(39)
The masses $`m_{A,B}`$ appearing in the kinetic energies $`E_𝐩^{A,B}`$ in the potential term can be replaced by their renormalized counterparts $`M_{A,B}`$ to the present order in the perturbative expansion. Eq. (39) differs from its scalar analogue only through the products of Dirac spinors (cf. Ref. ).
The effective Hamiltonian commutes with the total 3-momentum operator $`𝐏=𝐩_A+𝐩_B`$, hence we consider total momentum eigenstates from now on. In particular, we specialize to the center-of-mass system $`𝐏=0`$ where the effective Schrödinger equation becomes
$$\begin{array}{c}\left(\sqrt{M_A^2+𝐩^2}+\sqrt{M_B^2+𝐩^2}\right)\varphi (𝐩;r,s)g^2\underset{r^{},s^{}=1}{\overset{2}{}}\frac{d^3p^{}}{(2\pi )^3}\frac{1}{\sqrt{2E_𝐩^A\mathrm{\hspace{0.17em}2}E_𝐩^B\mathrm{\hspace{0.17em}2}E_𝐩^{}^A\mathrm{\hspace{0.17em}2}E_𝐩^{}^B}}\hfill \\ \hfill \times \frac{1}{2\omega _{𝐩𝐩^{}}}\left(\frac{1}{E_𝐩^A+\omega _{𝐩𝐩^{}}E_𝐩^{}^A}+\frac{1}{E_𝐩^B+\omega _{𝐩𝐩^{}}E_𝐩^{}^B}\right)\\ \hfill \times \left[\overline{u}_A(𝐩,r)u_A(𝐩^{},r^{})\right]\left[\overline{u}_B(𝐩,s)u_B(𝐩^{},s^{})\right]\varphi (𝐩^{};r^{},s^{})=(EE_V)\varphi (𝐩;r,s),\end{array}$$
(40)
with the relative momentum $`𝐩=𝐩_A=𝐩_B`$ and
$`\varphi (𝐩_A,r;𝐩_B,s)=\varphi (𝐩_A;r,s)(2\pi )^3\delta (𝐩_A+𝐩_B).`$ (41)
Observe in the potential term in Eqs. (39) and (40) the square roots of the kinetic energies which are characteristic of the non-locality of the interaction, and the differences of energies in the denominators which are due to the retardation of the interaction. The latter lead to a further non-locality in the effective potential, and also introduce non-hermiticity in the effective Hamiltonian.
For the discussion of the non-relativistic and the one-body limits in the next section, and also for the numerical solution of the effective Schrödinger equation, it is convenient to cast Eq. (40) into 2-spinorial form. To this end, the Dirac spinors are expressed in terms of Pauli spinors, most conveniently in the Dirac-Pauli representation,
$`u_A(𝐩,r)=\sqrt{E_𝐩^A+M_A}\left(\begin{array}{c}\chi _A(𝐩,r)\\ {\displaystyle \frac{𝐩𝝈}{E_𝐩^A+M_A}}\chi _A(𝐩,r)\end{array}\right).`$ (44)
Here the Pauli spinors are normalized in the usual way,
$`\chi _A^{}(𝐩,r)\chi _A(𝐩,s)=\delta _{rs}.`$ (45)
The Pauli spinors may or may not depend on the momentum $`𝐩`$. The possibility of a momentum dependence is important if one wishes to employ helicity eigenspinors. In the present work, however, we will not make use of a momentum-dependent basis.
Eq. (40) can now be rewritten in spinorial form as
$$\begin{array}{c}\left(\sqrt{M_A^2+𝐩^2}+\sqrt{M_B^2+𝐩^2}\right)\varphi (𝐩)\hfill \\ \hfill g^2\frac{d^3p^{}}{(2\pi )^3}\sqrt{\frac{E_𝐩^A+M_A}{2E_𝐩^A}\frac{E_𝐩^B+M_B}{2E_𝐩^B}\frac{E_𝐩^{}^A+M_A}{2E_𝐩^{}^A}\frac{E_𝐩^{}^B+M_B}{2E_𝐩^{}^B}}\\ \hfill \times \frac{1}{2\omega _{𝐩𝐩^{}}}\left(\frac{1}{E_𝐩^A+\omega _{𝐩𝐩^{}}E_𝐩^{}^A}+\frac{1}{E_𝐩^B+\omega _{𝐩𝐩^{}}E_𝐩^{}^B}\right)\\ \hfill \times [1\frac{𝐩𝝈_A}{E_𝐩^A+M_A}\frac{𝐩^{}𝝈_A}{E_𝐩^{}^A+M_A}][1\frac{𝐩𝝈_B}{E_𝐩^B+M_B}\frac{𝐩^{}𝝈_B}{E_𝐩^{}^B+M_B}]\varphi (𝐩^{})=(EE_V)\varphi (𝐩),\end{array}$$
(46)
where the spinorial wave function $`\varphi (𝐩)`$ is defined as
$`\varphi (𝐩)={\displaystyle \underset{r,s}{}}\varphi (𝐩;r,s)\left[\chi _A(𝐩,r)\chi _B(𝐩,s)\right].`$ (47)
As usual, $`𝝈_A`$ is understood to act on $`\chi _A(𝐩,r)`$ only.
## 3 Limiting cases, identical fermions and antifermions
In the present section, we will perform a series of non-trivial analytical checks on the effective Schrödinger equation. We will begin with the non-relativistic limit: if the wave function $`\varphi (𝐩;r,s)`$ is strongly suppressed for $`𝐩^2M_r^2`$, where
$`M_r={\displaystyle \frac{M_AM_B}{M_A+M_B}}`$ (48)
is the (renormalized) reduced mass, we can approximate the effective Schrödinger equation in the center-of-mass frame, Eq. (46), by
$`{\displaystyle \frac{𝐩^2}{2M_r}}\varphi (𝐩){\displaystyle \frac{d^3p^{}}{(2\pi )^3}\frac{4\pi \alpha }{\mu ^2+(𝐩𝐩^{})^2}\varphi (𝐩^{})}=\left(EE_VM_AM_B\right)\varphi (𝐩),`$ (49)
where we have introduced the “fine structure constant”
$`\alpha ={\displaystyle \frac{g^2}{4\pi }}`$ (50)
(for details on the approximation of the energy denominators, see Ref. ). Observe in particular that the effective Hamiltonian acts trivially on the spin degrees of freedom in this limit, as expected from non-relativistic scattering processes where spin orientations remain unchanged. As a consequence, Eq. (49) has the same form as in the case of scalar constituents .
The non-relativistic limit of Eq. (46) is hence precisely what we expect on physical grounds, a non-relativistic Schrödinger equation with the usual Yukawa potential (after Fourier transforming to position space). As discussed in detail in Ref. , the limit is attained for $`\alpha 1`$ and $`\mu M_r`$, both conditions being necessary.
We will now consider the so-called one-body limit “$`M_B\mathrm{}`$” where $`M_B^2M_A^2`$ and the wave function $`\varphi (𝐩;r,s)`$ is negligibly small for $`𝐩^2M_B^2`$. In this case, we can approximate the effective Schrödinger equation (46) by
$$\begin{array}{c}\sqrt{M_A^2+𝐩^2}\varphi (𝐩)g^2\frac{d^3p^{}}{(2\pi )^3}\sqrt{\frac{E_𝐩^A+M_A}{2E_𝐩^A}\frac{E_𝐩^{}^A+M_A}{2E_𝐩^{}^A}}\hfill \\ \hfill \times \frac{1}{2\omega _{𝐩𝐩^{}}}\left(\frac{1}{\omega _{𝐩𝐩^{}}}+\frac{1}{E_𝐩^A+\omega _{𝐩𝐩^{}}E_𝐩^{}^A}\right)\left[1\frac{𝐩𝝈_A}{E_𝐩^A+M_A}\frac{𝐩^{}𝝈_A}{E_𝐩^{}^A+M_A}\right]\varphi (𝐩^{})\\ \hfill =\left(EE_VM_B\right)\varphi (𝐩).\end{array}$$
(51)
The effective Hamiltonian acts trivially on the spin of particle $`B`$ in this limit. Eq. (51) has the form of a relativistic equation for fermion $`A`$ in an external potential, independent of the mass (except for a constant shift in the energy) and the spin orientation of particle $`B`$, in accord with physical expectations .
We can, however, go one step further and compare Eq. (51) with the relativistic equation for particle $`A`$ in the external potential due to a fixed source which exchanges spinless bosons of mass $`\mu `$ with particle $`A`$. The latter physical situation can also be described within the same general formalism, providing an internal consistency check for the application of the generalized Gell-Mann–Low theorem.
To this end, we begin by defining the Hamiltonian $`H^{}`$ of the fixed source system, which consists of the Hamiltonians corresponding to free fermions $`A`$ and scalar bosons of mass $`\mu `$, and the interaction term
$`H_1^{}=g{\displaystyle d^3x}\mathbf{:}\overline{\psi }_A(𝐱)\psi _A(𝐱)\phi (𝐱)\mathbf{:}+g\phi (\mathrm{𝟎}).`$ (52)
The form of the second contribution to the interaction Hamiltonian results from replacing the dynamical fermion field $`B`$ in Eq. (10) with a fixed source at $`𝐱=\mathrm{𝟎}`$. To see that, we consider $`\psi _B(𝐱)`$ as a classical Dirac field with probability density given by
$`\rho (𝐱)=\psi _B^{}(𝐱)\psi _B(𝐱)\overline{\psi }_B(𝐱)\psi _B(𝐱),`$ (53)
the latter approximate equality holding when $`\psi _B(𝐱)`$ describes a particle (not an anti-particle) and the relevant momenta satisfy $`𝐩^2M_B^2`$. For a fermion $`B`$ localized at $`𝐱=\mathrm{𝟎}`$ we then have
$`\overline{\psi }_B(𝐱)\psi _B(𝐱)=\delta (𝐱),`$ (54)
from which Eq. (52) follows.
The Bloch-Wilson Hamiltonian in the one-fermion sector is represented diagrammatically to lowest non-trivial order \[see Eq. (11)\] in Ref. (adding spin labels to the external lines for the present case). The corresponding algebraic expressions lead to a mass renormalization for fermion $`A`$ which is identical to the one discussed in the previous section, and the following matrix elements of the effective Hamiltonian:
$`𝐩_A,r|H_{\text{BW}}^{}|𝐩_A^{},r^{}`$
$`=\left(E_V^{}+\sqrt{M_A^2+𝐩_A^2}\right)(2\pi )^3\delta (𝐩_A𝐩_A^{})\delta _{rr^{}}`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3x\mathrm{\Delta }_F(0t,\mathrm{𝟎}𝐱)\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱)\psi _{𝐩_A^{},r^{}}^A(t,𝐱)\right]}`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3x\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)\psi _{𝐩_A^{},r^{}}^A(0,𝐱)\right]\mathrm{\Delta }_F(0t,𝐱\mathrm{𝟎})},`$ (55)
where the (renormalized) vacuum energy $`E_V^{}`$ differs from $`E_V`$ as determined from Eq. (12) for the case of two dynamical fermions, precisely because fermion $`B`$ has now been replaced by a fixed source which is formally considered to be part of the vacuum.
Using the non-covariant expression (14) for the propagators and Eq. (23) for the fermionic wave functions, Eq. (55) leads to the effective Schrödinger equation
$$\begin{array}{c}\sqrt{M_A^2+𝐩^2}\varphi (𝐩)g^2\frac{d^3p^{}}{(2\pi )^3}\sqrt{\frac{E_𝐩^A+M_A}{2E_𝐩^A}\frac{E_𝐩^{}^A+M_A}{2E_𝐩^{}^A}}\hfill \\ \hfill \times \frac{1}{2\omega _{𝐩𝐩^{}}}\left(\frac{1}{\omega _{𝐩𝐩^{}}}+\frac{1}{E_𝐩^A+\omega _{𝐩𝐩^{}}E_𝐩^{}^A}\right)\left[1\frac{𝐩𝝈}{E_𝐩^A+M_A}\frac{𝐩^{}𝝈}{E_𝐩^{}^A+M_A}\right]\varphi (𝐩^{})\\ \hfill =\left(EE_V^{}\right)\varphi (𝐩),\end{array}$$
(56)
where the wave function is now defined as
$`\varphi (𝐩)={\displaystyle \underset{r}{}}\varphi (𝐩,r)\chi (𝐩,r){\displaystyle \underset{r}{}}𝐩,r|\varphi \chi (𝐩,r).`$ (57)
Equation (56) is identical to Eq. (51) except for an irrelevant shift in the vacuum energy and the fact that the wave function in Eq. (51) includes the orientation of the spin of fermion $`B`$ which, however, has no influence on the dynamics of particle $`A`$. Self-consistency of the method in the one-body limit is hence established.
We will now investigate how the effective Schrödinger equation changes when we replace one of the constituents, say fermion $`A`$, by the corresponding antiparticle. First of all, consider the one-$`\overline{A}`$-antifermion sector where the matrix elements of the Bloch-Wilson Hamiltonian to lowest non-trivial order are given by
$$\begin{array}{c}𝐩_{\overline{A}},r|H_{\text{BW}}|𝐩_{\overline{A}}^{},s=\left(E_V+E_{𝐩_{\overline{A}}}^A\right)(2\pi )^3\delta (𝐩_{\overline{A}}𝐩_{\overline{A}}^{})\delta _{rs}\hfill \\ \hfill +ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s}^{\overline{A}}(t,𝐱^{})S_F^A(t0,𝐱^{}𝐱)\psi _{𝐩_{\overline{A}},r}^{\overline{A}}(0,𝐱)\right]\mathrm{\Delta }_F(0t,𝐱𝐱^{})\\ \hfill +ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s}^{\overline{A}}(0,𝐱)S_F^A(0t,𝐱𝐱^{})\psi _{𝐩_{\overline{A}},r}^{\overline{A}}(t,𝐱^{})\right]\mathrm{\Delta }_F(t0,𝐱^{}𝐱)\end{array}$$
(58)
\[compare with Eq. (22) and note the change in sign\]. Here we need the antifermion wave functions defined by
$`\overline{\psi }_{𝐩,r}^{\overline{A}}(t,𝐱)={\displaystyle \frac{\overline{v}_A(𝐩,r)}{\sqrt{2E_𝐩^A}}}e^{iE_𝐩^At+i𝐩𝐱}.`$ (59)
We can convert Eq. (58) to the form corresponding to a particle $`A`$ by introducing the charge conjugate wave functions
$`\psi _{𝐩,r}^{\overline{A},C}(t,𝐱)=C\left[\overline{\psi }_{𝐩,r}^{\overline{A}}(t,𝐱)\right]^T`$ (60)
(the superindex $`T`$ stands for transposition), with the charge conjugation matrix $`C=i\gamma ^0\gamma ^2`$ in the Dirac-Pauli representation. Eq. (58) is then transformed into expression (22) for the subspace of one fermion $`A`$ (including the sign) with the wave function $`\psi _{𝐩,r}^A(t,𝐱)`$ being replaced by $`\psi _{𝐩,r}^{\overline{A},C}(t,𝐱)`$, which corresponds to the replacement of the particle spinor $`u_A(𝐩,r)`$ with the charge conjugate antiparticle spinor
$`v_A^C(𝐩,r)=C\left[\overline{v}_A(𝐩,r)\right]^T`$ (61)
for the description of the spin orientation of the antifermion. Since $`v_A^C(𝐩,r)`$ is a positive-energy solution of the Dirac equation, the mass renormalization from $`m_A`$ to $`M_A`$ for antifermions is identical to the one for fermions $`A`$.
We now turn to the two-particle sector with an antifermion $`\overline{A}`$ and a fermion $`B`$. The matrix elements of the corresponding Bloch-Wilson Hamiltonian to lowest non-trivial order read
$`𝐩_{\overline{A}},r;𝐩_B,s|H_{\text{BW}}|𝐩_{\overline{A}}^{},r^{};𝐩_B^{},s^{}`$
$`=\left(E_V+\sqrt{M_A^2+𝐩_A^2}+\sqrt{M_B^2+𝐩_B^2}\right)(2\pi )^3\delta (𝐩_{\overline{A}}𝐩_{\overline{A}}^{})\delta _{rr^{}}(2\pi )^3\delta (𝐩_B𝐩_B^{})\delta _{ss^{}}`$
$`+ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_B,s}^B(0,𝐱)\psi _{𝐩_B^{},s^{}}^B(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_{\overline{A}}^{},r^{}}^{\overline{A}}(t,𝐱^{})\psi _{𝐩_{\overline{A}},r}^{\overline{A}}(t,𝐱^{})\right]`$
$`+ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_{\overline{A}}^{},r^{}}^{\overline{A}}(0,𝐱)\psi _{𝐩_{\overline{A}},r}^{\overline{A}}(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_B,s}^B(t,𝐱^{})\psi _{𝐩_B^{},s^{}}^B(t,𝐱^{})\right],`$ (62)
to be compared with Eq. (36). Again, the introduction of the charge conjugate wave functions converts Eq. (62) to the form of Eq. (36), with the wave function $`\psi _{𝐩,r}^A(t,𝐱)`$ replaced by $`\psi _{𝐩,r}^{\overline{A},C}(t,𝐱)`$, or the spinor $`u_A(𝐩,r)`$ by $`v_A^C(𝐩,r)`$. The effective Schrödinger equation for constituents $`\overline{A}B`$ can then be written in the spinorial form of Eq. (46), where in the definition (47) of the wave function the Pauli spinor $`\chi _A(𝐩,r)`$ has to be replaced by the charge conjugate spinor
$`\xi _A^C(𝐩,r)=i\sigma ^2\xi _A^{}(𝐩,r),`$ (63)
$`\xi _A(𝐩,r)`$ being the Pauli spinor that describes the spin orientation of the corresponding negative-energy solution of the Dirac equation. The most important result is that, just as for the mass renormalization, the interaction via a scalar boson is “charge conjugation blind”, i.e., does not distinguish between a fermion $`A`$ and its antifermion $`\overline{A}`$. This is, in fact, the expected behaviour.
The same arguments are used to describe an antifermion $`\overline{A}`$ interacting with a static source which leads to a Bloch-Wilson Hamiltonian and an effective Schrödinger equation analogous to Eqs. (55) and (56). Consequently, the one-body limit is self-consistent also in the case of a (light) antifermion $`\overline{A}`$. More interesting is the case of an antisource corresponding to the one-body limit $`M_B\mathrm{}`$ for an antifermion $`\overline{B}`$: a simple-minded argument replaces $`\psi _B(𝐱)`$ in the interaction Hamiltonian with a classical negative-energy Dirac field, which leads to a probability density
$`\rho (𝐱)=\psi _B^{}(𝐱)\psi _B(𝐱)\overline{\psi }_B(𝐱)\psi _B(𝐱),`$ (64)
if we suppose that the relevant momenta satisfy $`𝐩^2M_B^2`$. Equation (64) with $`\rho (𝐱)=\delta (𝐱)`$ would be in conflict with the one-body limit $`M_B\mathrm{}`$ for constituents $`A\overline{B}`$ or $`\overline{A}\overline{B}`$. The reason is, of course, that $`\rho (𝐱)`$ is to be interpreted physically as a charge density (when multiplied with the charge of fermion $`B`$), and is *negative* for antifermions, hence it is $`\overline{\psi }_B(𝐱)\psi _B(𝐱)`$ which turns out to be positive and is to be replaced by $`\delta (𝐱)`$. The latter results are related to the use of anticommutators in the quantization of the Dirac field.
In order to have a formally satisfactory description, we define an approximately localized antifermion state with spin orientation $`s`$ as
$`|𝐱=\mathrm{𝟎},\mathrm{\Delta }x;s={\displaystyle \frac{d^3p_{\overline{B}}}{(2\pi )^3}\left(\frac{8\pi }{3}\mathrm{\Delta }x^2\right)^{3/4}e^{\mathrm{\Delta }x^2𝐩_{\overline{B}}^2/3}|𝐩_{\overline{B}},s},`$ (65)
where $`|𝐩_{\overline{B}},s`$ denotes a $`\overline{B}`$-antifermion 3-momentum eigenstate. If we choose $`M_B`$ large enough for
$`\mathrm{\Delta }x^2M_B^21`$ (66)
to hold, we obtain
$`𝐱=\mathrm{𝟎},\mathrm{\Delta }x;s|\mathbf{:}\overline{\psi }_B(𝐱)\psi _B(𝐱)\mathbf{:}|𝐱=\mathrm{𝟎},\mathrm{\Delta }x;s=\left({\displaystyle \frac{3}{2\pi \mathrm{\Delta }x^2}}\right)^{3/2}e^{3𝐱^2/(2\mathrm{\Delta }x^2)}.`$ (67)
In the limit $`\mathrm{\Delta }x0`$ the right-hand side of Eq. (67) tends towards $`\delta (𝐱)`$. Equation (66) implies that we need $`M_B\mathrm{}`$ (even faster) in this limit. We take Eq. (67) as justification to replace $`[\mathbf{:}\overline{\psi }_B(𝐱)\psi _B(𝐱)\mathbf{:}]`$ with $`\delta (𝐱)`$ for a fixed antisource, leading to the interaction Hamiltonian (52). The one-body limit $`M_B\mathrm{}`$ is hence fully consistent also for the case of an antifermion $`\overline{B}`$.
We will now consider bound states of identical fermions, $`A`$-fermions to be concrete. To this end, we calculate the Bloch-Wilson Hamiltonian to lowest non-trivial order for the subspace
$`\mathrm{\Omega }_0=\text{span}\left\{|𝐩_A,r;𝐩_A^{},s|𝐩_A,𝐩_A^{}^3,r,s=1,2\right\}`$ (68)
of two-$`A`$-fermion states. As a consequence of the identity of the constituents, “crossed” diagrams appear which carry a relative minus sign due to the antisymmetry
$`|𝐩_A,r;𝐩_A^{},s=|𝐩_A^{},s;𝐩_A,r.`$ (69)
However, mass renormalization works exactly as before, and also the effective Schrödinger equation is the same as in the case of $`AB`$ bound states (with $`M_B=M_A`$) when we take into account the antisymmetry of the wave function
$`\varphi (𝐩_A,r;𝐩_A^{},s)=\varphi (𝐩_A^{},s;𝐩_A,r).`$ (70)
In the center-of-mass system, we have consequently
$`\varphi (𝐩)=\varphi (𝐩)^t`$ (71)
with the definitions (41) and (47), where the transposition $`t`$ refers to the tensor product,
$`\left[\chi _A(𝐩,r)\chi _A(𝐩,s)\right]^t=\chi _A(𝐩,s)\chi _A(𝐩,r).`$ (72)
The consequences of antisymmetry for the solutions of the effective Schrödinger equation will be discussed in the next section. In the non-relativistic limit, however, the situation is particularly simple because the spin degrees of freedom do not participate in the dynamics: the solutions of Eq. (49) with even orbital angular momentum have symmetric spatial wave functions and hence antisymmetric spin states (total spin zero), while solutions with odd orbital angular momentum have antisymmetric orbital wave functions and hence necessarily symmetric spin states (total spin one).
In a theory which only contains $`A`$-fermions and scalar bosons, the results for the mass renormalization and the $`AA`$ bound states are the same as the ones presented above for a theory with $`A`$\- and $`B`$-fermions, only the (irrelevant) corrections to the free vacuum energy, $`E_VE_0`$, are different (the $`B`$-fermion vacuum loops are missing).
The last case we will consider in this section is the one of $`A\overline{A}`$ bound states, of one fermion and the corresponding antifermion. In this case, there are additional contributions from the virtual annihilation diagrams to the effective Hamiltonian, so that the effective potential reads
$`𝐩_A,r;𝐩_{\overline{A}},s|V|𝐩_A^{},r^{};𝐩_{\overline{A}}^{},s^{}`$
$`=ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s^{}}^{\overline{A}}(0,𝐱)\psi _{𝐩_{\overline{A}},s}^{\overline{A}}(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱^{})\psi _{𝐩_A^{},r^{}}^A(t,𝐱^{})\right]`$
$`+ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)\psi _{𝐩_A^{},r^{}}^A(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s^{}}^{\overline{A}}(t,𝐱^{})\psi _{𝐩_{\overline{A}},s}^{\overline{A}}(t,𝐱^{})\right]`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)\psi _{𝐩_{\overline{A}},s}^{\overline{A}}(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s^{}}^{\overline{A}}(t,𝐱^{})\psi _{𝐩_A^{},r^{}}^A(t,𝐱^{})\right]`$
$`ig^2{\displaystyle _{\mathrm{}}^0}𝑑te^{ϵ|t|}{\displaystyle d^3xd^3x^{}\left[\overline{\psi }_{𝐩_{\overline{A}}^{},s^{}}^{\overline{A}}(0,𝐱)\psi _{𝐩_A^{},r^{}}^A(0,𝐱)\right]}`$
$`\times \mathrm{\Delta }_F(0t,𝐱𝐱^{})\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱^{})\psi _{𝐩_{\overline{A}},s}^{\overline{A}}(t,𝐱^{})\right]`$
$`={\displaystyle \frac{g^2}{\sqrt{2E_{𝐩_A}^A\mathrm{\hspace{0.17em}2}E_{𝐩_{\overline{A}}}^A\mathrm{\hspace{0.17em}2}E_{𝐩_A^{}}^A\mathrm{\hspace{0.17em}2}E_{𝐩_{\overline{A}}^{}}^A}}}\{{\displaystyle \frac{1}{2\omega _{𝐩_A𝐩_A^{}}}}({\displaystyle \frac{1}{E_{𝐩_A}^A+\omega _{𝐩_A𝐩_A^{}}E_{𝐩_A^{}}^A}}+{\displaystyle \frac{1}{E_{𝐩_{\overline{A}}}^A+\omega _{𝐩_{\overline{A}}𝐩_{\overline{A}}^{}}E_{𝐩_{\overline{A}}^{}}^A}})`$
$`\times \left[\overline{u}_A(𝐩_A,r)u_A(𝐩_A^{},r^{})\right]\left[\overline{v}_A(𝐩_{\overline{A}}^{},s^{})v_A(𝐩_{\overline{A}},s)\right]`$
$`{\displaystyle \frac{1}{2\omega _{𝐩_A+𝐩_{\overline{A}}}}}\left({\displaystyle \frac{1}{\omega _{𝐩_A+𝐩_{\overline{A}}}+E_{𝐩_A}^A+E_{𝐩_{\overline{A}}}^A}}+{\displaystyle \frac{1}{\omega _{𝐩_A^{}+𝐩_{\overline{A}}^{}}E_{𝐩_A^{}}^AE_{𝐩_{\overline{A}}^{}}^A}}\right)`$
$`\times \left[\overline{u}_A(𝐩_A,r)v_A(𝐩_{\overline{A}},s)\right]\left[\overline{v}_A(𝐩_{\overline{A}}^{},s^{})u_A(𝐩_A^{},r^{})\right]\}(2\pi )^3\delta (𝐩_A+𝐩_{\overline{A}}𝐩_A^{}𝐩_{\overline{A}}^{}).`$ (73)
We will leave the discussion of the effect of virtual annihilation on the bound state energies in different theories and the related instability of these states for a future investigation.
To close this section, we will briefly comment on the corresponding results in a purely bosonic theory, where the constituents are chosen to be charged scalar bosons. The consistency of the non-relativistic and one-body limits in this case has been shown in detail in Ref. . If we substitute one or both of the constituents by antibosons, there are, compared to the fermionic case, no additional minus signs from anticommutation relations to take into account and, of course, no spinor structures, consequently the whole argument is much simpler than for fermionic constituents. The results are, however, finally the same: mass renormalization is identical for antiparticles and for particles, and the interaction due to scalar boson exchange is universally attractive and does not distinguish particles from antiparticles.
As for a static antisource, consider the charge density (properly to be multiplied by the charge of boson $`B`$)
$`\rho (x)=\varphi _B^{}(x)i{\displaystyle \frac{}{t}}\varphi _B(x)\varphi _B(x)i{\displaystyle \frac{}{t}}\varphi _B^{}(x)2M_B\varphi _B^{}(x)\varphi _B(x)`$ (74)
for a classical negative-energy solution of the Klein-Gordon equation, in case that the relevant momenta fulfill $`𝐩^2M_B^2`$. The probability density is the negative of $`\rho (x)`$, hence we would replace
$`\mathbf{:}\varphi _B^{}(𝐱)\varphi _B(𝐱)\mathbf{:}={\displaystyle \frac{1}{2M_B}}\delta (𝐱)`$ (75)
in the interaction Hamiltonian for a localized antisource. A more formal argument proceeds in analogy with Eqs. (65)–(67) for the fermionic case, with the result (75). The one-boson limit $`M_B\mathrm{}`$ is then consistent also in the case of an antiboson $`\overline{B}`$.
Finally, in the case of identical bosonic constituents, the effective Schrödinger equation is the same as for bosonic $`AB`$ bound states, only that the wave function has to be symmetric under particle exchange in this case. Since there are no spin degrees of freedom in the scalar bosonic case, the spatial wave function has to be symmetric, hence only even angular momenta are allowed.
## 4 Numerical solution
In order to actually solve the effective Schrödinger equation in the form (46), it is convenient first to separate off the angular and spin degrees of freedom. A direct numerical solution of Eq. (46) would lead to numerical instabilities for equal and opposite momenta, due to the presence of a singularity in the integrand. The effective Hamiltonian is rotationally invariant, hence it is natural to consider total angular momentum eigenstates. To make contact to the usual spectroscopy, we choose to couple first the individual spins to a total spin $`𝐒`$ and then couple this spin with the relative orbital angular momentum $`𝐋`$ to the total angular momentum $`𝐉`$. The usual construction with Clebsch-Gordan coefficients yields simultaneous eigenstates of $`𝐉^2`$, $`J_z`$, $`𝐒^2`$, and $`𝐋^2`$ which we will denote as $`{}_{}{}^{\mathrm{\hspace{0.17em}2}S+1}𝒴_{lM}^{J}(\widehat{𝐩})`$, $`\widehat{𝐩}𝐩/|𝐩|`$. Explicit expressions are given in Appendix C.
The Hamiltonian (46) contains the helicity operators $`\widehat{𝐩}𝝈_A`$ and $`\widehat{𝐩}𝝈_B`$. These operators are hermitian and unitary, and in particular
$`(\widehat{𝐩}𝝈_A)^2=(\widehat{𝐩}𝝈_B)^2=1.`$ (76)
The helicity operators are invariant under spatial rotations, however, they are odd under spatial parity transformations which maintain the spin directions unchanged. Since the Schrödinger equation (46) contains only even powers of helicity operators, the effective Hamiltonian is parity even (the intrinsic parities of the constituent fermions have no use in the present context, and we will not consider them in the following).
We hence have the conservation of total angular momentum $`J`$ and spatial parity $`(1)^l`$, but a priori not of relative orbital angular momentum $`l`$ or total spin $`S`$. For given $`J`$, $`l`$ can take the values $`J`$ (for $`S=0`$) and $`J,J\pm 1`$ (for $`S=1`$), $`l=J,J1`$ being excluded for $`J=0`$, $`S=1`$. Taking into account the conservation of $`(1)^l`$, we can hence conclude without any explicit calculation that the effective Hamiltonian may mix states with $`l=J`$, $`S=0`$ and with $`l=J`$, $`S=1`$ on the one hand (we will call this “S-coupling”), and states $`l=J1`$, $`S=1`$ and $`l=J+1`$, $`S=1`$ (“L-coupling”) on the other. The effective Schrödinger equation will then decay into pairs of coupled one-dimensional equations. In the special case $`J=0`$, neither of the two mixings is possible. For future use we remark that the use of helicity eigenstates is expected to diagonalize the effective Hamiltonian in the S-coupled sector, thus slightly simplifying the calculations, although there is no reason why L-coupling should not occur in this case.
For the actual solution of the effective Schrödinger equation (46), we need explicit expressions for the application of the helicity operators. Again, it is clear from the fact that the helicity operators preserve total angular momentum and change parity, that the application of a helicity operator maps S-coupled states to L-coupled states and vice versa. The explicit expressions, as well as a rather pedestrian way to derive them, are presented in Appendix C. Finally, the integration over the angles, i.e., over $`\widehat{𝐩}^{}`$, in Eq. (46) can be performed with the help of a partial wave decomposition combined with the spherical harmonics addition theorem,
$`V(p,p^{},\mathrm{cos}\theta )={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi }{2l+1}}a_l(p,p^{}){\displaystyle \underset{m=l}{\overset{l}{}}}Y_{lm}(\widehat{𝐩})Y_{lm}^{}(\widehat{𝐩}^{}),`$ (77)
where $`p|𝐩|`$, $`\theta `$ denotes the angle between $`𝐩`$ and $`𝐩^{}`$, and
$`a_l(p,p^{})={\displaystyle \frac{2l+1}{2}}{\displaystyle _1^1}d\mathrm{cos}\theta P_l(\mathrm{cos}\theta )V(p,p^{},\mathrm{cos}\theta ).`$ (78)
(Note that we have changed conventions relative to the ones employed in our earlier work .)
The whole procedure outlined above can be carried through independently of the boson mass $`\mu `$. In what follows, I will focus on the case of massless bosons $`\mu =0`$, as discussed in the introduction. The explicit form of the effective Schrödinger equation in this case is, again, given in Appendix C. No additional difficulties are expected in the massive case in principle, even though the existence of a critical coupling constant changes the qualitative features of the spectrum.
The (pairs of) one-dimensional integral equations can now be solved numerically. To this end, the equations were converted to (continuous) matrix form and the corresponding two-dimensional integrals approximated by finite sums over a discrete two-dimensional grid. The distribution of abscissas took the logarithmic singularity of the integrand and the long range of the wave functions in configuration space into account. We have approximated the solutions by a finite linear combination of an appropriately chosen set of basis functions, the same we had used before in the scalar case . Two parameters that determine the shape of the basis functions were optimized variationally. The orthogonality of the basis functions could be retained numerically to 11 to 14 decimal places. Both energies and wave functions converged with increasing number of integration points and basis functions. However, in our experience convergence does not guarantee the correctness of the solutions if the choice of basis is inappropriate. For this reason, we have also checked the residual $`r_i(p)`$ of the solutions $`\varphi _i(p)`$ defined as $`(EH_{\text{BW}})\varphi _i(p)=r_i(p)`$. The determination of the residual $`r_i(p)`$ and the “point-wise” convergence of the wave functions were limited essentially by the redundancy of the grid points (up to 400) with respect to the number of basis functions (up to 40). The analogue of the Gibbs phenomenon in Fourier series was, in the worst case, of the order of two percent.
The numerical solutions (for $`\mu =0`$) are shown in Figs. 1 and 2 for two extreme mass ratios, $`M_A=M_B`$ and $`M_B\mathrm{}`$ (with $`M_A`$ fixed), as functions of the (Yukawa theory) fine structure constant $`\alpha =g^2/4\pi `$.
Between these two extremes, the eigenvalues for fixed $`\alpha `$ can be seen to vary smoothly with the mass ratio in Fig. 3 for $`\alpha =1`$.
In all figures, the energy eigenvalues are represented normalized to twice the non-relativistic ionization energy in a Coulomb potential, $`M_r\alpha ^2`$, where $`M_r`$ denotes the reduced mass.
In Figs. 1 and 2, it is seen that in the small-coupling limit $`\alpha 0`$, the energies tend to the non-relativistic Coulomb values
$`{\displaystyle \frac{M_r\alpha ^2}{2}}{\displaystyle \frac{1}{n^2}},n=1,2,3,\mathrm{},`$ (79)
as expected for $`\mu =0`$. In particular, we observe the characteristic degeneracies in this limit. For instance, we expect and find that the following states tend to the same energy eigenvalue $`M_r\alpha ^2/8`$ (principal quantum number $`n=2`$) for $`\alpha 0`$: the ground states in the L-coupled $`J=0`$ and $`J=2`$ sectors (states $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{0}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{2}^{}`$ in the usual spectroscopic notation $`n{}_{}{}^{\mathrm{\hspace{0.17em}2}S+1}L_{J}^{}`$) and in the S-coupled $`J=1`$ sector ($`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$), as well as the first excited states in the S-coupled $`J=0`$ and $`J=1`$ sectors ($`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$) and in the L-coupled $`J=1`$ sector ($`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}S_{1}^{}`$). To be precise, the lowest-energy eigenstates in the S-coupled $`J=1`$ sector are linear combinations of the $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$ states, the coefficients depending on the mass ratio. In Fig. 2, only three of these degenerate states are visible because there are always two curves lying on top of each other (see the discussion of the one-body limit below).
Let us now discuss the equal-mass case of Fig. 1 in detail. The binding is *weaker* in this case than predicted by the non-relativistic formula. Compared to an electromagnetic interaction (exchange of photons) as in positronium, where we eliminate the contribution of the virtual annihilation diagram, the sign of the relativistic corrections is just the opposite (with the exception of the $`1{}_{}{}^{\mathrm{\hspace{0.17em}3}}S_{1}^{}`$ state). Also, the ordering of the levels is different.
In the case of equal masses, the effective Hamiltonian possesses an additional symmetry under the exchange of the fermions $`A`$ and $`B`$,
$`\varphi (𝐩)\varphi (𝐩)^t`$ (80)
in terms of the spinorial wave function in the center-of-mass system (see Eq. (72)). This symmetry is implicit in our discussion of bound states of identical fermions in Section 3, where we saw that in the case of identical constituents the effective Schrödinger equation is the same as for $`AB`$ bound states, hence it must possess exchange symmetry. The exchange parity for the angular momentum eigenstates is $`(1)^l(1)^{S+1}`$. Together with the symmetry under spatial parity $`(1)^l`$ (remember that we do not take the intrinsic parities of the fermions into account), this new symmetry forbids the $`S`$-coupling, hence $`S`$ becomes a good quantum number in this case. This can also be seen explicitly in the expressions for the matrix elements of the effective potential in the S-coupled sector, Eq. (173) in Appendix C, for $`M_A=M_B`$. In the case of identical fermionic constituents, the wave function must be antisymmetric with respect to particle exchange, hence, for $`J`$ even, in the S-coupled sector only $`S=0`$ states are possible, while we have both $`l=J1`$ and $`l=J+1`$ states in the L-coupled sector, in particular the coupling between these states remains. On the other hand, for $`J`$ odd, we have only $`S=1`$ states in the S-coupled sector, and *no* possible state in the L-coupled sector. Among the eight states shown in Fig. 1, five are antisymmetric under particle exchange, namely the states $`1{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$, $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$, $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{0}^{}`$, $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$, and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{2}^{}`$. We can see explicitly in these examples that the absence of S-coupling is necessary for the antisymmetry of the states.
In the one-body limit $`M_B\mathrm{}`$ depicted in Fig. 2, the sign of the relativistic corrections changes for several states with respect to the equal-mass case. Within numerical accuracy, there are always two exactly degenerate states. The reason for the degeneracy can be seen from the effective Schrödinger equation in the one-body limit, Eq. (51): the effective Hamiltonian is invariant under rotations which involve the spatial coordinates and the spin $`𝐬_A`$ of particle $`A`$ only, and *independently* under rotations of spin $`𝐬_B`$ (which does not affect the dynamics). The total angular momentum of fermion $`A`$, $`𝐣_A=𝐋+𝐬_A`$, is then a conserved quantity, and it is natural to consider simultaneous eigenstates of $`𝐣_A^2`$, $`j_{A,z}`$, $`𝐋^2`$, and $`s_{B,z}`$. Since $`j_A=l\pm 1/2`$ and spatial parity $`(1)^l`$ is conserved as before, it follows that $`l`$ is a good quantum number in this limit.
As a further consequence, states that only differ in the value of $`j_{A,z}`$ or $`s_{B,z}`$ are degenerate, and one may as well consider simultaneous eigenstates of $`𝐉^2`$, $`J_z`$, $`𝐣_A^2`$, and $`𝐋^2`$, where $`𝐉=𝐣_A+𝐬_B`$ with eigenvalues $`J=j_A\pm 1/2`$. Now states which only differ in the eigenvalue of $`𝐉^2`$ or $`J_z`$ are degenerate. We can express the latter states in terms of the simultaneous eigenstates of $`𝐉^2`$, $`J_z`$, $`𝐋^2`$, and $`𝐒^2`$ that we are using in the numerical calculations. From the beforegoing discussion, we expect the absence of L-coupling and the degeneracy of one of the S-coupled states which is eigenstate of $`𝐣_A^2`$ with eigenvalue $`j_A=J+1/2`$ and $`l=J`$, with an L-coupled state with $`𝐉^2`$-eigenvalue $`J+1`$ and $`l=J`$ (hence necessarily $`j_A=J+1/2`$). The other S-coupled state with $`l=J`$ is an $`𝐣_A^2`$-eigenstate with $`j_A=J1/2`$, and is degenerate with an L-coupled state with $`𝐉^2`$-eigenvalue $`J1`$, $`l=J`$, and $`j_A=J1/2`$. Explicit expressions for the eigenstates in the different coupling schemes and their relations are given in Appendix C, where we follow a different line of reasoning starting from the explicit expressions for the matrix elements of the effective potential in the S- and L-coupled sectors, Eqs. (173) and (176).
These expectations are fully borne out in the results of the numerical calculations. Among the eight states calculated in Fig. 2, we expect and find the following to be degenerate: $`1{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$ and $`1{}_{}{}^{\mathrm{\hspace{0.17em}3}}S_{1}^{}`$, $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}S_{1}^{}`$, $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{0}^{}`$ with one linear combination of $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$, and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{2}^{}`$ with the orthogonal linear combination of $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$. The coefficients of these linear combinations are predicted and found to be $`\sqrt{1/3}`$ and $`\sqrt{2/3}`$, and $`\sqrt{2/3}`$ and $`\sqrt{1/3}`$, respectively (see Eq. (192) in Appendix C), and correspond to eigenstates of $`𝐣_A^2`$ with eigenvalues $`j_A=1/2`$ and $`3/2`$. Unlike for the Dirac equation with an electromagnetic Coulomb potential, states with the same $`j_A`$ but different $`l`$, here $`j_A=1/2`$ and $`l=0`$ (states $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}S_{0}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}S_{1}^{}`$) or $`l=1`$ (states $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{0}^{}`$ and the first linear combination of $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$), are not degenerate. In addition, the ordering of the $`l=1`$ states is opposite to the electromagnetic case.
Figure 3 shows the smooth transition between the two extreme mass ratios $`M_A=M_B`$ and $`M_A/M_B=0`$ for fixed $`\alpha =1`$, and the appearance of the characteristic degeneracies in the limit $`M_A/M_B0`$. One would like to compare the coefficients of the mixing of the S-coupled states $`2{}_{}{}^{\mathrm{\hspace{0.17em}1}}P_{1}^{}`$ and $`2{}_{}{}^{\mathrm{\hspace{0.17em}3}}P_{1}^{}`$ against theoretical predictions. However, the diagonalization of the effective potential matrix (173) is generally not possible without solving the entire equation, i.e., it cannot be isolated from the $`p`$-dependence of the wave function. To the order $`\alpha ^4`$ of the first relativistic corrections, the diagonalization can still be performed analytically and turns out to factorize from the “radial” $`p`$-dependence. The results are presented in Appendix D and provide very satisfactory approximations to the corresponding results of the numerical calculations. In particular, the analytical results indicate that there is no hyperfine splitting to the order $`\alpha ^4`$ and $`M_A/M_B`$, for the energy levels that are degenerate in the one-body limit. In fact, even for $`\alpha =1`$ no hyperfine splitting of the order $`M_A/M_B`$ is visible in Fig. 3.
## 5 Conclusions
We have presented what appears to be the first consistent treatment of bound states in Yukawa theory. It is the result of a straightforward application of the generalized Gell-Mann–Low theorem. The consistency of the method has been checked thoroughly. In particular, we have shown that mass renormalization can be performed exactly as in a manifestly covariant formulation, even though the renormalization conditions were imposed entirely within our Hamiltonian framework. We have checked the non-relativistic and one-body limits, replaced the fermionic constituents by antifermions, and considered the case of identical constituents. In all these cases, the formalism generates the correct results in a very natural way. In the numerical calculations, no abnormal solutions have been found (nor were there expected to be any, due to the absence of relative time or energy as a dynamical variable), and all the characteristic degeneracies in the non-relativistic and one-body limits show up in the numerical results with very good accuracy.
In general terms, the framework presented here has several advantages over other formulations of quantum field theoretic bound state equations. As we have shown, the derivation of the effective Schrödinger equation is straightforward and presents no essential complications in the case of fermionic constituents as compared to scalar bosons. In principle, the complete bound state spectrum can be obtained as we have demonstrated by numerically determining the eight lowest-lying states (corresponding to the non-relativistic principal quantum numbers $`n=1`$ and $`n=2`$). The wave functions for these states are also obtained in the course of the (approximate) diagonalization of the effective Hamiltonian.
Several rather technical issues, which are nonetheless expected to be important for related work in the near future, have been treated in detail in the appendices. In particular, the relation between manifestly covariant and non-covariant representations of the relevant loop integrals has been established (Appendix A). We have discussed dimensional, Pauli-Villars, Schwinger proper time and covariant and non-covariant cutoff regularizations for the appearing one-loop integrals from a non-covariant Hamiltonian perspective (Appendix B). Finally, we have presented explicit expressions for the angular momentum eigenstates in different coupling schemes, for the application of helicity operators and the coefficient functions in a partial wave expansion, all necessary ingredients for the separation of angular and spin variables in the effective Schrödinger equation (Appendix C).
The results of this work, if only as a point of departure for the application to more realistic physical situations in the future, bear on fundamental issues in nuclear and high energy physics, as for example the nucleon-nucleon interaction. In this respect, one interesting particular result of the numerical computations is the qualitative difference between the relativistic bound state spectra for scalar (boson exchange) and electromagnetic interactions.
### Acknowledgments
One of us (A.W.) gratefully acknowledges support by CIC-UMSNH and Conacyt grant 32729-E. Part of the research (by N.L.) was done at the Department of Physics and Astronomy of the University of Pittsburgh in the group of Eric Swanson and Steve Dytman.
## Appendix A Covariant and non-covariant representations of loop integrals
In this appendix, we will relate different expressions for loop integrals in momentum space, where the manifestly covariant representations arise directly from the momentum space Feynman rules, while the non-covariant representations result naturally from the application of the Gell-Mann–Low theorem.
Let us begin with the lowest-order correction to the vacuum energy density Eq. (19). The equivalence will be established by performing the integrations over $`p_0`$ and $`p_0^{}`$ through the use of complex integration theory. For greater clarity, we will first discuss the analogous problem in a purely scalar theory . The corresponding formula differs from Eq. (19) by the global sign (for both expressions) and the numerators of the integrands which are simply equal to one. Considering the first term (for particle $`A`$) for concreteness, we begin with the integral over $`p_0^{}`$,
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dp_0^{}}{2\pi }}{\displaystyle \frac{1}{[(pp^{})^2\mu ^2+iϵ][p_{}^{}{}_{}{}^{2}m_A^2+iϵ]}}.`$ (81)
Standard application of the residue theorem, closing the integration contour through the usual large semicircle in the upper half plane, gives
$`i\left\{{\displaystyle \frac{1}{2E_𝐩^{}^A\left[(p_0+E_𝐩^{}^A)^2\omega _{𝐩𝐩^{}}^2+iϵ\right]}}+{\displaystyle \frac{1}{2\omega _{𝐩𝐩^{}}\left[(p_0\omega _{𝐩𝐩^{}})^2(E_𝐩^{}^A)^2+iϵ\right]}}\right\}.`$ (82)
This expression becomes much more transparent after a decomposition in partial fractions with respect to $`p_0`$,
$$\begin{array}{c}\frac{i}{2\omega _{𝐩𝐩^{}}\mathrm{\hspace{0.17em}2}E_𝐩^{}^A}[\frac{1}{p_0+\omega _{𝐩𝐩^{}}+E_𝐩^{}^Ai\eta }\frac{1}{p_0\omega _{𝐩𝐩^{}}E_𝐩^{}^A+i\eta }\hfill \\ \hfill +2\pi i\delta (p_0\omega _{𝐩𝐩^{}}+E_𝐩^{}^A)],\end{array}$$
(83)
where we have used the formula
$`{\displaystyle \frac{1}{\omega i\eta }}{\displaystyle \frac{1}{\omega +i\eta }}=2\pi i\delta (\omega )`$ (84)
(for $`\eta 0`$). Equation (83) *cannot be correct* as it stands: the original integral (81) is even under $`p_0p_0`$ (by the substitution $`p_0^{}p_0^{}`$ of the integration variable), and this symmetry is manifestly broken by the delta function in Eq. (83).
However, the term with the delta function only contributes for $`p_0=\omega _{𝐩𝐩^{}}E_𝐩^{}^A`$, and this is precisely the value of $`p_0`$ where the two poles in the upper half plane coincide. The correct evaluation of the residue at the double pole gives for the integral in this case
$`{\displaystyle \frac{i}{2\omega _{𝐩𝐩^{}}\mathrm{\hspace{0.17em}2}E_𝐩^{}^A}}\left[{\displaystyle \frac{1}{p_0+\omega _{𝐩𝐩^{}}+E_𝐩^{}^Ai\eta }}{\displaystyle \frac{1}{p_0\omega _{𝐩𝐩^{}}E_𝐩^{}^A+i\eta }}\right]_{p_0=\omega _{𝐩𝐩^{}}E_𝐩^{}^A},`$ (85)
so the result (83) is wrong for this value of $`p_0`$, and the term with the delta function has to be omitted in (83). The integration over $`p_0`$ can then be performed easily, using the residue theorem again. The result, after substituting $`𝐩^{}𝐩^{}`$, is the one expected from Eq. (19) or rather its scalar analogue. Note that the use of Eq. (83) including the delta function would lead to additional (incorrect) terms in Eq. (19).
In the Yukawa case, Eq. (19) proper, the argument is nearly identical, only the expressions are slightly more complicated. The result for the naive $`p_0^{}`$-integration, after the decomposition in partial fractions, reads
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dp_0^{}}{2\pi }}{\displaystyle \frac{4(pp^{}+m_A^2)}{[(pp^{})^2\mu ^2+iϵ][p_{}^{}{}_{}{}^{2}m_A^2+iϵ]}}`$
$`={\displaystyle \frac{4i}{2\omega _{𝐩𝐩^{}}\mathrm{\hspace{0.17em}2}E_𝐩^{}^A}}\{2E_𝐩^{}^A+[E_𝐩^{}^A(E_𝐩^{}^A+\omega _{𝐩𝐩^{}})𝐩𝐩^{}+m_A^2][{\displaystyle \frac{1}{p_0+\omega _{𝐩𝐩^{}}+E_𝐩^{}^Ai\eta }}`$
$`{\displaystyle \frac{1}{p_0\omega _{𝐩𝐩^{}}E_𝐩^{}^A+i\eta }}]+[E_𝐩^{}^A(E_𝐩^{}^A\omega _{𝐩𝐩^{}})𝐩𝐩^{}+m_A^2]2\pi i\delta (p_0\omega _{𝐩𝐩^{}}+E_𝐩^{}^A)\}.`$ (86)
By a calculation of the residue at the double pole for the special case $`p_0=\omega _{𝐩𝐩^{}}E_𝐩^{}^A`$, one can again show that the term with the delta function in Eq. (86) is spurious. Integration over $`p_0`$ of the rest gives the desired result.
We now turn to the lowest-order corrections to the mass, Eq. (27). As before, we begin with the scalar case where all the numerators in Eq. (27) are replaced by one . Then the result of the $`p_0^{}`$-integration in the manifestly covariant expression is given precisely by Eq. (83) above, with $`p_0`$ to be replaced by $`E_𝐩^A`$. The delta function in Eq. (83) is, of course, again spurious, although it cannot give any contribution anyway as long as $`\mu ^2<4m_A^2`$.
In Yukawa theory, where Eq. (27) properly applies, the $`p_0^{}`$-integration gives the following result , after a decomposition in partial fractions (with respect to $`p_0`$),
$$\begin{array}{c}_{\mathrm{}}^{\mathrm{}}\frac{dp_0^{}}{2\pi }\frac{p^{}\gamma +m_A}{[(pp^{})^2\mu ^2+iϵ][p_{}^{}{}_{}{}^{2}m_A^2+iϵ]}\hfill \\ \hfill =\frac{i}{2\omega _{𝐩𝐩^{}}\mathrm{\hspace{0.17em}2}E_𝐩^{}^A}[\frac{E_𝐩^{}^A\gamma _0𝐩^{}𝜸+m_A}{p_0+\omega _{𝐩𝐩^{}}+E_𝐩^{}^Ai\eta }\frac{E_𝐩^{}^A\gamma _0𝐩^{}𝜸+m_A}{p_0\omega _{𝐩𝐩^{}}E_𝐩^{}^A+i\eta }\\ \hfill +(E_𝐩^{}^A\gamma _0𝐩^{}𝜸+m_A)2\pi i\delta (p_0\omega _{𝐩𝐩^{}}+E_𝐩^{}^A)],\end{array}$$
(87)
putting $`p_0=E_𝐩^A`$. Again, the delta function turns out to be spurious (by explicitly considering the case of a double pole in the upper half plane), which leads to the non-covariant expression in Eq. (27).
## Appendix B Regularization of one-loop integrals
The aim of this appendix is to derive Eq. (28) starting from the non-covariant expression for $`G(𝐩)`$ in Eq. (27) in a suitably regularized form, so that all integrals appearing in the derivation are well-defined. By far the simplest way is to go through the manifestly covariant form also presented in Eq. (27), regularized correspondingly for the present purpose. For the complications that arise in a direct derivation in the non-covariant formulation for the simpler case of a purely scalar theory where it has to be shown that $`G(𝐩)`$ is actually independent of $`𝐩`$ (and only depends on $`p^2|_{p_0=E_𝐩^A}=m_A^2`$), see Ref. .
### B.1 Dimensional regularization
The technically simplest regularization scheme (although not the most natural one in the present context, see below) is dimensional regularization. The idea is to continuously change the dimension of space to smaller values where all the integrals are well-defined, to be save to spatial dimensions smaller than two, i.e., space-time dimensions $`D<3`$. We can then establish the relation between the non-covariant and covariant expression in Eq. (27) for these dimensions and show that the dependence on $`D`$ is analytical, with a simple pole appearing at $`D=4`$. As a consequence, Eq. (28) can be shown to hold true for space-time dimensions arbitrarily close to (but smaller than) four by analytical continuation. However, the definition of the integrals in arbitrary (non-integer) dimensions is made precise only in the (Euclidean) covariant formulation, which makes this form of regularization somewhat unnatural for the non-covariant expressions.
In detail, we begin by establishing the relation between the non–covariant and covariant expressions for $`G(𝐩)`$ in Eq. (27), but for $`(D1)`$ spatial dimensions instead of three where, to begin with, $`D<3`$ in order that all integrals are well-defined. The analogue of Eq. (27) in $`D`$ dimensions can then be shown by integrating over $`p_0^{}`$ in the covariant form exactly as detailed in Appendix A \[see Eq. (87) and the following remarks\]. The $`d^{D1}p^{}`$-integration is not touched in this process. With the covariant form at hand, we can introduce Feynman parameters in the usual way:
$`G_\epsilon ^{\text{DR}}(𝐩)`$ $`=ig^2{\displaystyle \frac{d^Dp^{}}{(2\pi )^D}_0^1𝑑x\frac{p^{}\gamma +m_A}{\left[(pp^{})^2x\mu ^2x+p_{}^{}{}_{}{}^{2}(1x)m_A^2(1x)+iϵ\right]^2}}|_{p_0=E_𝐩^A}`$
$`=ig^2{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^Dp^{}}{(2\pi )^D}\frac{p^{}\gamma +m_A}{\left[(p^{}xp)^2+x(1x)p^2x\mu ^2(1x)m_A^2+iϵ\right]^2}}|_{p_0=E_𝐩^A}`$
$`=ig^2{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{xp\gamma +m_A}{\left[q^2+x(1x)p^2x\mu ^2(1x)m_A^2+iϵ\right]^2}}|_{p_0=E_𝐩^A},`$ (88)
denoting the dimensionally regularized form of $`G(𝐩)`$ as $`G_\epsilon ^{\text{DR}}(𝐩)`$, where $`\epsilon =4D`$. In the last step, we have shifted the integration variable to $`q=p^{}xp`$ and used that for the term $`q\gamma `$ emerging in the numerator, the integrand is odd. These manipulations are unproblematic as long as we stick to dimensions $`D<3`$ where the integrals are well-defined.
It is now convenient to decompose $`G_\epsilon ^{\text{DR}}(𝐩)`$ in two parts in analogy with Eq. (28),
$`G_\epsilon ^{\text{DR}}(𝐩)=\left[G_{1,\epsilon }^{\text{DR}}(p^2)p\gamma +G_{0,\epsilon }^{\text{DR}}(p^2)m_A\right]_{p_0=E_𝐩^A}.`$ (89)
However, all space-time vectors in this equation are still $`D`$-dimensional. Replacing $`p^2|_{p_0=E_𝐩^A}`$ by $`m_A^2`$ and Wick rotating to Euclidean space gives
$`G_{0,\epsilon }^{\text{DR}}(m_A^2)=g^2{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^Dq_E}{(2\pi )^D}\frac{1}{\left[q_E^2+(1x)^2m_A^2+x\mu ^2\right]^2}}`$ (90)
The corresponding expression for $`G_{1,\epsilon }^{\text{DR}}(m_A^2)`$ is equal to Eq. (90) except for an additional factor of $`x`$ in the integrand. It is these latter integrals which are rigorously defined for arbitrary, continuous values of $`D`$. They are actually defined in such a way as to make them analytical functions of $`D`$, with a simple pole at $`D=4`$. We can use analytic continuation to define all the beforegoing integrals for $`3D<4`$, leaving the established relations intact, in particular Eq. (89). In the limit $`\epsilon 0`$ or $`D4`$, via the standard formulae of dimensional regularization,
$`G_{0,\epsilon }^{\text{DR}}(m_A^2)={\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle _0^1}𝑑x\left[{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )\mathrm{ln}{\displaystyle \frac{(1x)^2m_A^2+x\mu ^2}{\kappa ^2}}\right],`$ (91)
plus terms which tend to zero in this limit. In Eq. (91), $`\kappa `$ is the renormalization scale, and $`g`$ is left dimensionless also for $`D4`$. The corresponding expression for $`G_{1,\epsilon }^{\text{DR}}(m_A^2)`$ is, again, equal to Eq. (91) except for an additional factor of $`x`$ in the integrand.
After an integration by parts, the evaluation of the integrals over the Feynman parameter $`x`$ is straightforward and yields
$`G_{0,\epsilon }^{\text{DR}}(m_A^2)={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )\mathrm{ln}{\displaystyle \frac{m_A^2}{\kappa ^2}}+2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`2\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\},`$ (92)
$`G_{1,\epsilon }^{\text{DR}}(m_A^2)={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )]{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{m_A^2}{\kappa ^2}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`({\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}(2{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}`$ (93)
(for $`\epsilon 0`$). Observe the simple pole at $`\epsilon =0`$ or $`D=4`$ and the absence of IR divergences for $`\mu 0`$.
The result for $`G_{0,\epsilon }^{\text{DR}}(m_A^2)`$ gives immediately the explicit dimensionally regularized expression for $`\mathrm{\Delta }m_A^2`$ in the purely scalar case considered in Ref. (where the coupling constant $`g`$ has the dimension of mass). For Yukawa theory, we define $`\mathrm{\Delta }m_A^2`$ by
$`\mathrm{\Delta }m_A^2=2m_A^2\left[G_{0,\epsilon }^{\text{DR}}(m_A^2)+G_{1,\epsilon }^{\text{DR}}(m_A^2)\right],`$ (94)
so that
$`\overline{u}_A^\epsilon (𝐩,r)G_\epsilon ^{\text{DR}}(𝐩)u_A^\epsilon (𝐩,s)=\mathrm{\Delta }m_A^2\delta _{rs}`$ (95)
from the analogue of Eq. (29), where the Dirac equation in $`D`$ dimensions is used. Explicitly, from Eqs. (92) and (93),
$`{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2m_A^2}}`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{3}{2}}[{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )]{\displaystyle \frac{3}{2}}\mathrm{ln}{\displaystyle \frac{m_A^2}{\kappa ^2}}+{\displaystyle \frac{7}{2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{3}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`(4{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\},`$ (96)
which can now be used to redefine $`m_A`$ in terms of $`M_A`$ as in Eq. (33) and the limit $`\epsilon 0`$ be taken with the (finite) physical mass $`M_A`$ held fixed.
### B.2 Pauli-Villars regularization
We now consider Pauli-Villars regularization which turns out to be the most natural regularization scheme in the context of Hamiltonian (non-covariant) perturbation theory. It is effected by subtracting from the expression in Eq. (27), denoted for the time being as $`G(𝐩,\mu )`$ to make its dependence on the boson mass $`\mu `$ explicit, a similar contribution for a ficticious “heavy” boson of mass $`\mathrm{\Lambda }`$ to define
$`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)=G(𝐩,\mu )G(𝐩,\mathrm{\Lambda }).`$ (97)
It is then easy to verify by power counting that $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ is UV finite, both from the non-covariant and the manifestly covariant expression.
The $`p_0^{}`$-integration in the covariant expression can then be performed again as in Appendix A, only that the following integration over $`d^3p^{}`$ is now well-defined. This establishes the equivalence of the non-covariant and manifestly covariant expressions for $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ \[i.e., the regularized form of Eq. (27)\] rigorously. Starting from the covariant form, we now introduce Feynman parameters and shift the integration variable as in Eq. (88) to obtain
$`G(𝐩,\mu )=ig^2{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{xp\gamma +m_A}{\left[q^2+x(1x)p^2x\mu ^2(1x)m_A^2+iϵ\right]^2}}|_{p_0=E_𝐩^A}`$ (98)
and the analogous expression for $`G(𝐩,\mathrm{\Lambda })`$ \[remembering that its difference $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ is well-defined in 4 dimensions\]. From Eq. (98) and the corresponding expression for $`G(𝐩,\mathrm{\Lambda })`$, it is clear that $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ is of the form described in Eq. (28). The mass renormalization can then already be effected at this stage as detailed following Eq. (28), reading $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ instead of $`G(𝐩)`$.
To evaluate $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ analytically, it is easiest to apply dimensional regularization to $`G(𝐩,\mu )`$ and $`G(𝐩,\mathrm{\Lambda })`$ separately and analytically continue the result for the difference $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ back to four space-time dimensions. We start with the form obtained in Eq. (98) above, replace $`p^2|_{p_0=E_𝐩^A}`$ by $`m_A^2`$, and separate $`G_\mathrm{\Lambda }^{\text{PV}}(𝐩)`$ in two parts $`G_{1,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ and $`G_{0,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$, analogous to Eq. (28). The continuation to $`D`$ space-time dimensions and a Wick rotation to Euclidean space leads to the expression (90) for the contribution $`G_0(m_A^2,\mu )`$. The integrations over $`d^Dq_E`$ and the Feynman parameter then yield the results given in Eqs. (92) and (93) for the contributions $`G_0(m_A^2,\mu )`$ and $`G_1(m_A^2,\mu )`$, respectively. The expressions for $`G_0(m_A^2,\mathrm{\Lambda })`$ and $`G_1(m_A^2,\mathrm{\Lambda })`$ are obtained by replacing $`\mu \mathrm{\Lambda }`$. In the (relevant) limit of large $`\mathrm{\Lambda }`$, one analytically continues the last terms in Eqs. (92) and (93) to $`\mathrm{\Lambda }^2>4m_A^2`$ and expands the complete expression in powers of $`m_A^2/\mathrm{\Lambda }^2`$ (up to second order). This gives, after several cancellations,
$`G_0(m_A^2,\mathrm{\Lambda })`$ $`{\displaystyle \frac{g^2}{(4\pi )^2}}\left\{{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\kappa ^2}}+1\right\},`$
$`G_1(m_A^2,\mathrm{\Lambda })`$ $`{\displaystyle \frac{g^2}{(4\pi )^2}}\left\{{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{2}{\epsilon }}\gamma _E+\mathrm{ln}(4\pi )\right]{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\kappa ^2}}+{\displaystyle \frac{1}{4}}\right\}`$ (99)
for $`\mathrm{\Lambda }^2m_A^2`$. The same results can be obtained somewhat easier by replacing $`(1x)^2m_A^2+x\mathrm{\Lambda }^2x\mathrm{\Lambda }^2`$ in the argument of the logarithm in Eq. (91) \[and analogously for $`G_1(m_A^2,\mathrm{\Lambda })`$\] from the start.
We hence obtain, finally, for $`G_0(m_A^2)`$ and $`G_1(m_A^2)`$ in Pauli-Villars regularization
$`G_{0,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`2\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\},`$ (100)
$`G_{1,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+{\displaystyle \frac{5}{4}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`(2{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}`$ (101)
($`\mathrm{\Lambda }^2m_A^2`$). We obtain a logarithmic UV divergence with $`\mathrm{\Lambda }\mathrm{}`$ and, again, the absence of IR divergences for $`\mu 0`$. We emphasize that we have used dimensional regularization only as a convenient calculational tool here, and that $`G_{0,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ and $`G_{1,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ are well-defined in four space-time dimensions from the start.
The result for $`G_{0,\mathrm{\Lambda }}^{\text{PV}}(m_A^2)`$ gives directly the expression for $`\mathrm{\Delta }m_A^2`$ in Pauli-Villars regularization for the purely scalar case of Ref. . For Yukawa theory, we have from Eq. (30)
$`{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2m_A^2}}`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{3}{2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+{\displaystyle \frac{9}{4}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{3}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$ (102)
$`(4{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}.`$ (103)
### B.3 Schwinger proper time regularization
The other two regularization schemes that we will discuss, Schwinger proper time and momentum cutoff regularization, are only effective in Euclidean space. Hence we have to write the non-covariant expression in Eq. (27) in Euclidean space in order to regularize. After evaluation, it can then be analytically continued to physical values of the external variables.
Let us begin with the diagrammatic expressions in Eq. (22). We rewrite the first contribution to the mass renormalization by use of the covariant representations (13) of the propagators and Eq. (23) for the fermionic wave functions as
$$\begin{array}{c}ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(0,𝐱)S_F^A(0t,𝐱𝐱^{})\psi _{𝐩_A^{},s}^A(t,𝐱^{})\right]\mathrm{\Delta }_F(0t,𝐱𝐱^{})\hfill \\ \hfill =ig^2\frac{d^3p^{}}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}d^3xd^3x^{}\frac{e^{i𝐩_A𝐱}}{\sqrt{2E_{𝐩_A}^A}}\frac{dp_0^{}}{2\pi }\frac{\left[\overline{u}_A(𝐩_A,r)\left(p^{}\gamma +m_A\right)u_A(𝐩_A^{},s)\right]e^{i𝐩^{}(𝐱𝐱^{})}}{p_0^{}{}_{}{}^{2}𝐩^{}{}_{}{}^{2}m_A^2+iϵ}\\ \hfill \times \frac{dk_0}{2\pi }\frac{e^{i𝐤(𝐱𝐱^{})}}{k_0^2𝐤^2\mu ^2+iϵ}\frac{e^{i𝐩_A^{}𝐱^{}}}{\sqrt{2E_{𝐩_A^{}}^A}}_{\mathrm{}}^0dte^{ϵt}e^{i(p_0^{}+k_0p_0)t}|_{p_0=E_{𝐩_A^{}}^A}\end{array}$$
(104)
In the latter expression we Wick rotate in the mathematically positive sense
$`p_0^{}p_0^{}=ip_0^E,k_0k_0=ik_0^E,`$ (105)
where $`p_0^{}^E`$ and $`k_0^E`$ are real after the rotation. The sense of the rotation is determined by the position of the poles in the integrand (by the $`iϵ`$-prescription). Since $`p_0^{}^E`$ and $`k_0^E`$ take both positive and negative values, it is imperative to Wick rotate $`t`$, too, keeping the exponents imaginary in order that the integrand do not blow up. This implies the rotation
$`tt=it^E`$ (106)
in the negative sense, where $`\mathrm{}<t^E0`$. Keeping $`p_0`$ fixed at $`p_0=E_{𝐩_A^{}}^A`$ would then lead to a divergence in the $`t^E`$-integration, so that we have to rotate, in addition,
$`p_0p_0=ip_0^E.`$ (107)
After performing the integration, we then analytically continue the result to $`p_0^E=iE_{𝐩_A^{}}^A`$.
As a result of this Wick rotation, the expression (104) is written as
$$\begin{array}{c}\frac{g^2}{\sqrt{2E_{𝐩_A}^A\mathrm{\hspace{0.17em}2}E_{𝐩_A^{}}^A}}\frac{d^3p^{}}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}d^3xd^3x^{}e^{i𝐩_A𝐱}\hfill \\ \hfill \times \frac{dp_0^E}{2\pi }\frac{\left[\overline{u}_A(𝐩_A,r)\left(p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A\right)u_A(𝐩_A^{},s)\right]e^{i𝐩^{}(𝐱𝐱^{})}}{(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2}\\ \hfill \times \frac{dk_0^E}{2\pi }\frac{e^{i𝐤(𝐱𝐱^{})}}{(k_0^E)^2+𝐤^2+\mu ^2}e^{i𝐩_A^{}𝐱^{}}_{\mathrm{}}^0dt^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}|_{p_0^EiE_{𝐩_A^{}}^A},\end{array}$$
(108)
where we have defined $`ϵ^E=iϵ`$ to assure the convergence of the $`t^E`$-integration. Alternatively, we can assume $`ϵ`$ to have a positive imaginary part from the beginning. As for the $`\gamma `$ matrices, we have chosen the common convention where
$`p_0^E\gamma _0^E+𝐩^{}𝜸^E=p^{}\gamma ,`$ (109)
i.e., $`\gamma _0^E=i\gamma _0`$ and $`\gamma _i^E=\gamma ^{iE}=\gamma ^i`$. Finally, for the momentum-space representation, we integrate over $`d^3x`$, $`d^3x^{}`$ and $`d^3k`$ (taking advantage of the three-dimensional $`\delta `$-function resulting from the $`x`$\- and $`x^{}`$-integrations). The resulting contribution to $`G(𝐩)`$ \[compare with Eq. (26)\] is
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}\frac{dp_0^E}{2\pi }\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2}\hfill \\ \hfill \times \frac{dk_0^E}{2\pi }\frac{1}{(k_0^E)^2+(𝐩𝐩^{})^2+\mu ^2}_{\mathrm{}}^0dt^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}|_{p_0^EiE_𝐩^A}.\end{array}$$
(110)
Turning now to the second contribution to the mass renormalization in Eq. (22),
$$\begin{array}{c}ig^2_{\mathrm{}}^0𝑑te^{ϵ|t|}d^3xd^3x^{}\left[\overline{\psi }_{𝐩_A,r}^A(t,𝐱^{})S_F^A(t0,𝐱^{}𝐱)\psi _{𝐩_A^{},s}^A(0,𝐱)\right]\mathrm{\Delta }_F(t0,𝐱^{}𝐱)\hfill \\ \hfill =ig^2\frac{d^3p^{}}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}d^3xd^3x^{}\frac{e^{i𝐩_A𝐱^{}}}{\sqrt{2E_{𝐩_A}^A}}\frac{dp_0^{}}{2\pi }\frac{\left[\overline{u}_A(𝐩_A,r)\left(p^{}\gamma +m_A\right)u_A(𝐩_A^{},s)\right]e^{i𝐩^{}(𝐱𝐱^{})}}{p_0^{}{}_{}{}^{2}𝐩^{}{}_{}{}^{2}m_A^2+iϵ}\\ \hfill \times \frac{dk_0}{2\pi }\frac{e^{i𝐤(𝐱𝐱^{})}}{k_0^2𝐤^2\mu ^2+iϵ}\frac{e^{i𝐩_A^{}𝐱}}{\sqrt{2E_{𝐩_A^{}}^A}}_{\mathrm{}}^0dte^{ϵt}e^{i(p_0^{}+k_0p_0)t}|_{p_0=E_{𝐩_A}^A},\end{array}$$
(111)
and going through the steps following Eq. (104), we arrive at the following contribution to $`G(𝐩)`$:
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}\frac{dp_0^E}{2\pi }\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2}\hfill \\ \hfill \times \frac{dk_0^E}{2\pi }\frac{1}{(k_0^E)^2+(𝐩𝐩^{})^2+\mu ^2}_{\mathrm{}}^0dt^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}|_{p_0^EiE_𝐩^A}.\end{array}$$
(112)
It is easy to check the correctness of the results (110) and (112) by performing the integrations over $`p_0^E`$ and $`k_0^E`$ (using the residue theorem) and finally over $`t^E`$, leading to the non-covariant expression in Eq. (27).
After this preparation, we are now in a position to introduce the Schwinger proper time regularization by replacing
$`{\displaystyle \frac{1}{(p_0^E)^2+(E_𝐩^{}^A)^2}}={\displaystyle _0^{\mathrm{}}}𝑑\alpha e^{\alpha [(p_0^E)^2+(E_𝐩^{}^A)^2]}{\displaystyle _{1/\mathrm{\Lambda }^2}^{\mathrm{}}}𝑑\alpha e^{\alpha [(p_0^E)^2+(E_𝐩^{}^A)^2]},`$ (113)
and analogously for the other (scalar) propagator. We emphasize that only in Euclidean space the modification of the lower limit for the parameter integration corresponds to the exponential suppression of the propagator for large momenta. Consequently, the proper time regularization of the first contribution Eq. (110) reads
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}_{1/\mathrm{\Lambda }^2}^{\mathrm{}}𝑑\alpha 𝑑\beta e^{\alpha (E_𝐩^{}^A)^2\beta \omega _{𝐩𝐩^{}}^2}_{\mathrm{}}^0𝑑t^Ee^{ϵ^Et^E}e^{ip_0^Et^E}\hfill \\ \hfill \times \frac{dp_0^E}{2\pi }[p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A]e^{\alpha (p_0^E)^2+ip_0^Et^E}\frac{dk_0^E}{2\pi }e^{\beta (k_0^E)^2+ik_0^Et^E}|_{p_0^EiE_𝐩^A}.\end{array}$$
(114)
The integrations over $`p_0^E`$ and $`k_0^E`$ in Eq. (114) are readily performed. After shifting the Euclidean time variable $`t^Et^E=t^E+2i\alpha \beta p_0^E/(\alpha +\beta )`$, one may analytically continue $`p_0^E`$ to $`iE_{𝐩_A}^A`$ directly in the integrand to obtain
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}_{1/\mathrm{\Lambda }^2}^{\mathrm{}}\frac{d\alpha d\beta }{4\pi \sqrt{\alpha \beta }}e^{\alpha (E_𝐩^{}^A)^2\beta \omega _{𝐩𝐩^{}}^2+\alpha \beta (E_𝐩^A)^2/(\alpha +\beta )}\hfill \\ \hfill \times _{\mathrm{}}^{2\alpha \beta E_𝐩^A/(\alpha +\beta )}dt^E[(\frac{\beta E_𝐩^A}{\alpha +\beta }\frac{t^E}{2\alpha })\gamma _0𝐩^{}𝜸+m_A]e^{(\alpha +\beta )(t^E)^2/(4\alpha \beta )}.\end{array}$$
(115)
A little algebra shows the first exponent in this expression to be negative. Integration over $`t^E`$ gives
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}_{1/\mathrm{\Lambda }^2}^{\mathrm{}}\frac{d\alpha d\beta }{4\pi \sqrt{\alpha \beta }}e^{\alpha (E_𝐩^{}^A)^2\beta \omega _{𝐩𝐩^{}}^2}\{\frac{\beta }{\alpha +\beta }\gamma _0\hfill \\ \hfill +\sqrt{\frac{\pi \alpha \beta }{\alpha +\beta }}[\frac{\beta E_𝐩^A}{\alpha +\beta }\gamma _0𝐩^{}𝜸+m_A]e^{\alpha \beta (E_𝐩^A)^2/(\alpha +\beta )}\text{erfc}(\sqrt{\frac{\alpha \beta }{\alpha +\beta }}E_𝐩^A)\},\end{array}$$
(116)
where the complementary error function is defined as
$`\text{erfc}(x)={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _x^{\mathrm{}}}𝑑te^{t^2}={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^x}𝑑te^{t^2}.`$ (117)
For the second contribution to the mass renormalization we start from Eq. (112), implement the Schwinger proper time regularization as in Eq. (113) and integrate over $`p_0^E`$, $`k_0^E`$, and $`t^E`$ (after a shift and the continuation $`p_0^EiE_𝐩^A`$) to arrive at
$$\begin{array}{c}g^2\frac{d^3p^{}}{(2\pi )^3}_{1/\mathrm{\Lambda }^2}^{\mathrm{}}\frac{d\alpha d\beta }{4\pi \sqrt{\alpha \beta }}e^{\alpha (E_𝐩^{}^A)^2\beta \omega _{𝐩𝐩^{}}^2}\{\frac{\beta }{\alpha +\beta }\gamma _0\hfill \\ \hfill +\sqrt{\frac{\pi \alpha \beta }{\alpha +\beta }}[\frac{\beta E_𝐩^A}{\alpha +\beta }\gamma _0𝐩^{}𝜸+m_A]e^{\alpha \beta (E_𝐩^A)^2/(\alpha +\beta )}\text{erfc}\left(\sqrt{\frac{\alpha \beta }{\alpha +\beta }}E_𝐩^A\right)\}.\end{array}$$
(118)
The results (116) and (118) represent the proper time regularization of the non-covariant expressions in Eq. (27). When we add them up, we arrive at the proper time regularized equivalent to Eq. (27),
$$\begin{array}{c}G_\mathrm{\Lambda }^{\text{PT}}(𝐩)=g^2\frac{d^3p^{}}{(2\pi )^3}_{1/\mathrm{\Lambda }^2}^{\mathrm{}}\frac{d\alpha d\beta }{2\pi }\sqrt{\frac{\pi }{\alpha +\beta }}\left[\frac{\beta E_𝐩^A}{\alpha +\beta }\gamma _0𝐩^{}𝜸+m_A\right]\hfill \\ \hfill \times e^{\alpha (E_𝐩^{}^A)^2\beta \omega _{𝐩𝐩^{}}^2+\alpha \beta (E_𝐩^A)^2/(\alpha +\beta )},\end{array}$$
(119)
where we have used that \[see Eq. (117)\]
$`\text{erfc}(x)+\text{erfc}(x)={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑te^{t^2}=2.`$ (120)
Next, we show that the sum (119) of the non-covariant expressions coincides with the Schwinger proper time regularized equivalent of the covariant form in Eq. (27). To this end, first perform a Wick rotation to rewrite the covariant expression in Eq. (27) in Euclidean space,
$`G(𝐩)=g^2{\displaystyle \frac{d^4p^E}{(2\pi )^4}\frac{p^E\gamma ^E+m_A}{\left[(p^Ep^E)^2+\mu ^2\right]\left[(p^E)^2+m_A^2\right]}}|_{p_0^EiE_𝐩^A},`$ (121)
where all the products of Euclidean 4-vectors (including squares) are understood to be Euclidean scalar products, i.e., taken with the positive Euclidean metric. The proper time regularization is introduced as in Eq. (113) and yields
$$\begin{array}{c}G_\mathrm{\Lambda }^{\text{PT}}(𝐩)=g^2_{1/\mathrm{\Lambda }^2}^{\mathrm{}}𝑑\alpha 𝑑\beta \frac{d^4p^E}{(2\pi )^4}\left[p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A\right]\hfill \\ \hfill \times e^{\alpha [(p_0^E)^2+𝐩^2+m_A^2]\beta [(p_0^Ep_0^E)^2+(𝐩𝐩^{})^2+\mu ^2]}|_{p_0^EiE_𝐩^A}.\end{array}$$
(122)
Integrating over $`p_0^E`$ and continuing $`p_0^E`$ to $`iE_𝐩^A`$ results in Eq. (119), thus establishing the equivalence of the non-covariant and covariant proper time regularized expressions.
Note that there is a much faster way to establish this equivalence by adding up Eq. (114) and its counterpart from Eq. (112), substituting $`t^Et^E`$ in the latter expression, and using the formula
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^Ee^{ϵ^E|t^E|}e^{i\omega t^E}=2\pi \delta (\omega ).`$ (123)
This procedure directly leads to the form (122), so it indeed represents a considerable shortcut. However, we were interested in writing down the proper time regularized equivalent of Eq. (27), which is why we followed the calculation through to expressions (116) and (118).
To establish the form (28), we start again from the covariant expression (122), but integrate over all four components of $`p^E`$ this time with the result
$`G_\mathrm{\Lambda }^{\text{PT}}(𝐩)={\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle _{1/\mathrm{\Lambda }^2}^{\mathrm{}}}{\displaystyle \frac{d\alpha d\beta }{(\alpha +\beta )^2}}\left[{\displaystyle \frac{\beta }{\alpha +\beta }}p\gamma +m_A\right]e^{\alpha m_A^2\beta \mu ^2+\alpha \beta p^2/(\alpha +\beta )}|_{p_0=E_𝐩^A},`$ (124)
from where relation (28) can be established. Note that the exponent in Eq. (124) is manifestly negative for $`p^2|_{p_0=E_𝐩^A}=m_A^2`$.
In order to evaluate the functions $`G_{0,\mathrm{\Lambda }}^{\text{PT}}(m_A^2)`$ and $`G_{1,\mathrm{\Lambda }}^{\text{PT}}(m_A^2)`$ in proper time regularization explicitly, we change variables $`(\alpha ,\beta )`$ to $`(\rho ,\sigma )`$, where
$$\rho =\alpha +\beta ,\sigma =\frac{\beta }{\alpha +\beta },$$
$$0<\sigma <1,\rho P(\sigma ,\mathrm{\Lambda }^2)\mathrm{max}(\frac{1}{\sigma \mathrm{\Lambda }^2},\frac{1}{(1\sigma )\mathrm{\Lambda }^2}).$$
(125)
The integration over $`\rho `$ can now be performed and leads to
$`G_\mathrm{\Lambda }^{\text{PT}}(𝐩)={\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle _0^1}𝑑\sigma \left[\sigma p\gamma +m_A\right]E_1\left([(1\sigma )^2m_A^2+\sigma \mu ^2]P(\sigma ,\mathrm{\Lambda }^2)\right),`$ (126)
where we have introduced the exponential integral function
$`E_1(x)={\displaystyle _x^{\mathrm{}}}𝑑t{\displaystyle \frac{e^t}{t}}={\displaystyle _1^{\mathrm{}}}𝑑t{\displaystyle \frac{e^{xt}}{t}},x>0.`$ (127)
Using the expansion of $`E_1(x)`$ for small values of the argument , we can approximate $`G_\mathrm{\Lambda }^{\text{PT}}(𝐩)`$ by
$`{\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle _0^1}𝑑\sigma \left[\sigma p\gamma +m_A\right]\left[\gamma _E+\mathrm{ln}\left([(1\sigma )^2m_A^2+\sigma \mu ^2]P(\sigma ,\mathrm{\Lambda }^2)\right)\right],`$ (128)
all higher terms in the expansion being suppressed by powers of $`1/\mathrm{\Lambda }`$. For the sake of readability, the demonstration of this latter assertion is relegated to the end of this section. The results (92) and (93) and the integral of $`\mathrm{ln}(\mathrm{\Lambda }^2P(\sigma ,\mathrm{\Lambda }^2))`$ which is elementary, can then be used to establish that
$`G_{0,\mathrm{\Lambda }}^{\text{PT}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}\gamma _E\mathrm{ln}2+1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`2\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\},`$ (129)
$`G_{1,\mathrm{\Lambda }}^{\text{PT}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{1}{2}}[\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}\gamma _E\mathrm{ln}2]+1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`(2{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}.`$ (130)
As before, $`G_{0,\mathrm{\Lambda }}^{\text{PT}}(m_A^2)`$ directly gives $`\mathrm{\Delta }m_A^2`$ in the purely scalar theory, while in the Yukawa case we have
$`{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2m_A^2}}`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{3}{2}}[\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}\gamma _E\mathrm{ln}2]+2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{3}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`(4{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}.`$ (131)
Finally, we will give the demonstration of Eq. (128) which was relegated to this point. To begin with, let us note that the approximation (128) is not trivial because the function $`P(\sigma ,\mathrm{\Lambda }^2)`$ takes large values for $`\sigma `$ close to the limits of integration, 0 and 1. We further remark that the extension of the integration domain to $`0<\sigma <1`$, $`\rho 2/\mathrm{\Lambda }^2`$ \[compare with Eq. (125)\], in which case the approximation corresponding to Eq. (128) is straightforward, leads to a different result.
This being said, we start from Eq. (126) and integrate by parts in order to get rid of the exponential integral function. The result can be written as
$$\begin{array}{c}G_\mathrm{\Lambda }^{\text{PT}}(𝐩)=\frac{g^2}{(4\pi )^2}_0^1𝑑\sigma \left[\frac{\sigma ^2}{2}p\gamma +\sigma m_A\right]e^{[(1\sigma )^2m_A^2+\sigma \mu ^2]P(\sigma ,\mathrm{\Lambda }^2)}\hfill \\ \hfill \times \frac{d}{d\sigma }\mathrm{ln}\left([(1\sigma )^2m_A^2+\sigma \mu ^2]P(\sigma ,\mathrm{\Lambda }^2)\right).\end{array}$$
(132)
In the exponent in this expression, we keep the divergent terms (for $`\sigma 0`$ and $`\sigma 1`$) and expand the rest in powers of $`\sigma `$ and $`(1\sigma )`$, respectively, leading to
$`{\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle _0^{1/2}}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]e^{m_A^2/(\mathrm{\Lambda }^2\sigma )}{\displaystyle \frac{d}{d\sigma }}[\mathrm{ln}\left({\displaystyle \frac{(1\sigma )^2m_A^2+\sigma \mu ^2}{\mathrm{\Lambda }^2}}\right)\mathrm{ln}\sigma ]`$
$`+{\displaystyle _{1/2}^1}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]e^{\mu ^2/(\mathrm{\Lambda }^2(1\sigma ))}{\displaystyle \frac{d}{d\sigma }}[\mathrm{ln}\left({\displaystyle \frac{(1\sigma )^2m_A^2+\sigma \mu ^2}{\mathrm{\Lambda }^2}}\right)\mathrm{ln}(1\sigma )]\},`$ (133)
the higher terms in the expansion of the exponential of the finite terms being suppressed by powers of $`\mathrm{\Lambda }`$, as we shall see shortly. First, note that we can, by way of the substitutions
$`y={\displaystyle \frac{1}{2\sigma }}\text{and}y={\displaystyle \frac{1}{2(1\sigma )}}`$ (134)
in the first and second integral in Eq. (133), respectively, express the contributions that contain the derivatives of $`\mathrm{ln}\sigma `$ and $`\mathrm{ln}(1\sigma )`$ in terms of the exponential integral functions
$`E_n(x)={\displaystyle _1^{\mathrm{}}}𝑑t{\displaystyle \frac{e^{xt}}{t^n}},x>0,n=1,2,3.`$ (135)
For small $`x`$, we have the expansions
$`E_1(x)`$ $`=\gamma _E\mathrm{ln}x+𝒪(x),`$
$`E_2(x)`$ $`=1+𝒪(x\mathrm{ln}x),`$
$`E_3(x)`$ $`=1/2+𝒪(x),`$ (136)
which lead to
$`{\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle _0^{1/2}}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]{\displaystyle \frac{e^{m_A^2/(\mathrm{\Lambda }^2\sigma )}}{\sigma }}`$
$`+{\displaystyle _{1/2}^1}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]{\displaystyle \frac{e^{\mu ^2/(\mathrm{\Lambda }^2(1\sigma ))}}{1\sigma }}\}`$
$`={\displaystyle \frac{g^2}{(4\pi )^2}}\left[{\displaystyle \frac{1}{2}}p\gamma +m_A\right]\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}\right)\gamma _E\mathrm{ln}21\right]`$ (137)
in the limit of large $`\mathrm{\Lambda }`$. In an analogous way, we can see that the higher orders in the expansion of the exponential of the finite terms (for $`\sigma 0`$ and $`\sigma 1`$) in Eq. (132) are suppressed by powers of $`\mathrm{\Lambda }`$, taking into account that the expressions resulting from this expansion and the first terms in the integrals (133) are continuous bounded functions over the intervals in question, and that furthermore the higher orders in the expansion carry inverse powers of $`\mathrm{\Lambda }`$.
As far as the first terms in the integrals (133) are concerned, the limit $`\mathrm{\Lambda }\mathrm{}`$ can be taken naively there (note that the $`\mathrm{\Lambda }`$-dependence in the argument of the logarithm is spurious since its derivative with respect to $`\sigma `$ gives zero), as we will show in a moment. If we suppose, for the time being, that this is indeed correct, the exponentials can be replaced by one, and an integration by parts yields
$`{\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle _0^{1/2}}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]e^{m_A^2/(\mathrm{\Lambda }^2\sigma )}{\displaystyle \frac{d}{d\sigma }}\mathrm{ln}\left({\displaystyle \frac{(1\sigma )^2m_A^2+\sigma \mu ^2}{\mathrm{\Lambda }^2}}\right)`$
$`+{\displaystyle _{1/2}^1}d\sigma [{\displaystyle \frac{\sigma ^2}{2}}p\gamma +\sigma m_A]e^{\mu ^2/(\mathrm{\Lambda }^2(1\sigma ))}{\displaystyle \frac{d}{d\sigma }}\mathrm{ln}\left({\displaystyle \frac{(1\sigma )^2m_A^2+\sigma \mu ^2}{\mathrm{\Lambda }^2}}\right)\}`$
$`={\displaystyle \frac{g^2}{(4\pi )^2}}\left\{\left[{\displaystyle \frac{1}{2}}p\gamma +m_A\right]\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}\right)+{\displaystyle _0^1}𝑑\sigma \left[\sigma p\gamma +m_A\right]\mathrm{ln}\left({\displaystyle \frac{(1\sigma )^2m_A^2+\sigma \mu ^2}{\mathrm{\Lambda }^2}}\right)\right\}`$ (138)
for $`\mathrm{\Lambda }\mathrm{}`$ (the $`\mathrm{\Lambda }`$-dependence cancels between the two terms). Using Eq. (137) and the results (92) and (93) for the remaining integration, we finally arrive at the expressions (129) and (130), thus confirming the approximation (128).
Let us now show in detail that it is, indeed, correct to replace the exponential on the left-hand side of Eq. (138) by one. To this end, take the integrals on the left-hand side of Eq. (138) as they stand, perform the substitutions Eq. (134) and decompose the result in partial fractions to find, after a somewhat lengthy calculation,
$$\begin{array}{c}\frac{g^2}{(4\pi )^2}_1^{\mathrm{}}dy\{\frac{1}{2}p\gamma [\frac{1}{2y^3}\frac{M_++M_{}}{2y^2}+\frac{M_+^2+M_{}^2}{y}\frac{M_{}^2}{y+M_+/2}\frac{M_+^2}{y+M_{}/2}]\hfill \\ \hfill +m_A[\frac{1}{y^2}\frac{M_++M_{}}{y}+\frac{M_{}}{y+M_+/2}+\frac{M_+}{y+M_{}/2}]\}e^{2(m_A^2/\mathrm{\Lambda }^2)y}\\ \hfill \frac{g^2}{(4\pi )^2}_1^{\mathrm{}}dy\{\frac{1}{2}p\gamma [\frac{1}{2y^3}\frac{M_++M_{}2}{2y^2}\\ \hfill \frac{M_+^2+M_{}^2}{y}+\frac{M_+^2}{y1/(2+2M_+)}+\frac{M_{}^2}{y1/(2+2M_{})}]\\ \hfill +m_A[\frac{1}{y^2}+\frac{M_++M_{}}{y}\frac{M_+}{y1/(2+2M_+)}\frac{M_{}}{y1/(2+2M_{})}]\}e^{2(\mu ^2/\mathrm{\Lambda }^2)y},\end{array}$$
(139)
where we have introduced the notations
$`M_\pm ={\displaystyle \frac{\mu ^2}{2m_A^2}}1\pm \sqrt{{\displaystyle \frac{\mu ^2}{2m_A^2}}\left({\displaystyle \frac{\mu ^2}{2m_A^2}}2\right)}.`$ (140)
For simplicity, we have written this expression for the case $`\mu ^24m_A^2`$ where the square roots are real. To compare with Eqs. (129) and (130) in the end, one has to analytically continue the result to smaller values of $`\mu ^2`$.
The integrals in Eq. (139) are well-defined with and without the exponentials. However, this is not true (in all cases) for the integrals over the partial fractions individually. It is clear then, that the exponentials can be considered as regulating factors for the integrals over the partial fractions. After summing up the individual results, the limit $`\mathrm{\Lambda }\mathrm{}`$ can safely be taken, thus removing the regulator. In the end, this is equivalent to replacing the exponentials by one in Eq. (138). The integrals in (139) can also be explicitly calculated with the help of Eqs. (135) and (136) in the limit $`\mathrm{\Lambda }\mathrm{}`$. As a result, Eq. (138) is recovered \[using Eqs. (92) and (93) to evaluate the integral in Eq. (138)\].
These comments conclude the demonstration of the correctness of the approximation (128) and hence of the results (129), (130) and (131) for the Schwinger proper time regularization. It is clear from the above that the proper time regularization will not be the method of choice in a Hamiltonian (not explicitly covariant) approach, considering the difficulties in establishing the regularized form of the non-covariant expressions and the evaluation of the corresponding (one-loop) integral as compared to the other regularizations discussed before.
### B.4 Momentum cutoff regularization
At last, we discuss the probably simplest regularization scheme available, the use of a momentum cutoff. As it turns out, it does not have very simple properties when used in the present context. We start with the non-covariant expressions written in Euclidean space after a Wick rotation in Eqs. (110) and (112). The momentum cutoff is implemented by restricting the integration domain (in Euclidean space) to $`(p^E)^2\mathrm{\Lambda }^2`$, hence the complete regularized expression is
$$\begin{array}{c}G_\mathrm{\Lambda }^{\text{MC}}(𝐩)=g^2^\mathrm{\Lambda }\frac{d^4p^E}{(2\pi )^4}\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2}\frac{dk_0^E}{2\pi }\frac{1}{(k_0^E)^2+(𝐩𝐩^{})^2+\mu ^2}\hfill \\ \hfill \times \left[_{\mathrm{}}^0𝑑t^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}+_{\mathrm{}}^0𝑑t^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}\right]_{p_0^EiE_𝐩^A}.\end{array}$$
(141)
We can perform the integrations over $`k_0^E`$ (with the help of the residue theorem) and $`t^E`$ in Eq. (141), separately for the two contributions, to obtain the result
$`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ $`=g^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p^E}{(2\pi )^4}}{\displaystyle \frac{i}{2\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{[p_0^Ep_0^Ei\omega _{𝐩𝐩^{}}][(p_0^E)^2+(E_𝐩^{}^A)^2]}}|_{p_0^EiE_𝐩^A}`$
$`g^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p^E}{(2\pi )^4}}{\displaystyle \frac{i}{2\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{[p_0^Ep_0^E+i\omega _{𝐩𝐩^{}}][(p_0^E)^2+(E_𝐩^{}^A)^2]}}|_{p_0^EiE_𝐩^A}.`$ (142)
Adding up the integrands in Eq. (142) gives
$`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)=g^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p^E}{(2\pi )^4}}{\displaystyle \frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{[(p_0^Ep_0^E)^2+(𝐩𝐩^{})^2+\mu ^2][(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2]}}|_{p_0^EiE_𝐩^A},`$ (143)
which obviously coincides with the cutoff regularization of the Euclidean version (121) of the covariant form in Eq. (27), thus establishing the equivalence of the non-covariant and covariant expressions in Eq. (27) in the momentum cutoff regularized form. One can use rotations in four-dimensional (Euclidean) space, taking into account the form of the integrand in Eq. (143) as well as the invariance of the integration measure and the integration domain, to show that $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ is in fact of the form of Eq. (28).
In order to obtain the three-dimensional form of the non-covariant expression in cutoff regularization, we start with Eq. (142) and perform the $`p_0^E`$-integration in
$`{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p^E}{(2\pi )^4}}={\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle _{\sqrt{\mathrm{\Lambda }^2𝐩^2}}^{\sqrt{\mathrm{\Lambda }^2𝐩^2}}}{\displaystyle \frac{dp_0^E}{2\pi }}`$ (144)
(with the $`d^3p^{}`$-integration restricted to $`𝐩^2\mathrm{\Lambda }^2`$). To this end, it is easiest to decompose the integrand in partial fractions with respect to $`p_0^E`$, for the first non-covariant contribution
$$\begin{array}{c}\frac{i}{2\omega _{𝐩𝐩^{}}}\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{[p_0^Ep_0^Ei\omega _{𝐩𝐩^{}}][(p_0^E)^2+(E_𝐩^{}^A)^2]}=\frac{i}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}\hfill \\ \hfill \times [\frac{1}{E_𝐩^{}^A\omega _{𝐩𝐩^{}}+ip_0^E}(\frac{(p_0^E+i\omega _{𝐩𝐩^{}})\gamma _0^E+𝐩^{}𝜸^E+m_A}{p_0^Ep_0^Ei\omega _{𝐩𝐩^{}}}\frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{p_0^EiE_𝐩^{}^A})\\ \hfill +\frac{1}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}ip_0^E}(\frac{(p_0^E+i\omega _{𝐩𝐩^{}})\gamma _0^E+𝐩^{}𝜸^E+m_A}{p_0^Ep_0^Ei\omega _{𝐩𝐩^{}}}\frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{p_0^E+iE_𝐩^{}^A})].\end{array}$$
(145)
The decomposition in partial fractions for the second contribution in Eq. (142) can be obtained immediately from the one above by realizing that the two contributions are complex conjugate to each other (if we consider $`p_0^E`$ and the Euclidean $`\gamma `$-matrices as real, for example by using the Majorana representation).
The $`p_0^E`$-integration is then straightforward, although the result is rather lengthy, indicating that the momentum cutoff would not usually be the regularization of choice in the present not manifestly covariant context, either. However, if we extended the $`p_0^E`$-integration to the whole real axis while sticking somewhat arbitrarily to $`𝐩^2\mathrm{\Lambda }^2`$, we would obtain the following simple result for the sum of the two contributions:
$`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$ $`=g^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}ip_0^E}}|_{p_0^EiE_𝐩^A}`$
$`g^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}+ip_0^E}}|_{p_0^EiE_𝐩^A}.`$ (146)
Eq. (146) results directly from the non-covariant expression in Eq. (27) by restricting the 3-momentum integration to $`𝐩^2\mathrm{\Lambda }^2`$. It is hence, in a not manifestly covariant approach, a very natural regularization scheme. Although it is apparent that four-dimensional Euclidean rotational invariance is broken in Eq. (146) and that hence $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$ will *not* be of the form (28) for any finite value of $`\mathrm{\Lambda }`$, it is not clear a priori whether the form (28) could not be recovered in the limit $`\mathrm{\Lambda }\mathrm{}`$. In the special case of a purely scalar theory, it was shown in Ref. that covariance is indeed reestablished for $`\mathrm{\Lambda }\mathrm{}`$ in the sense that $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$ becomes independent of $`𝐩`$ in this limit (and only depends on the square of the four-vector, $`p^2|_{p_0=E_𝐩^A}=m_A^2`$).
In the following, we will consider the integration (144) of the integrand (145) and its complex conjugate counterpart arising from the second non-covariant contribution in the limit of large $`\mathrm{\Lambda }`$ and compare it with $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$. To make this comparison mathematically precise, we divide the integration over the spatial momentum $`𝐩^{}`$ into the two regions $`𝐩^2<K^2`$ and $`K^2<𝐩^2<\mathrm{\Lambda }^2`$ with an intermediate scale $`K`$ which fulfills $`𝐩^2,m_A^2,\mu ^2K^2\mathrm{\Lambda }^2`$. Then, in the limit $`\mathrm{\Lambda }\mathrm{}`$, we can approximate the integral over the first region,
$`{\displaystyle ^K}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle _{\sqrt{\mathrm{\Lambda }^2𝐩^2}}^{\sqrt{\mathrm{\Lambda }^2𝐩^2}}}{\displaystyle \frac{dp_0^E}{2\pi }}{\displaystyle ^K}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dp_0^E}{2\pi }},`$ (147)
hence the integral over $`𝐩^2<K^2`$ of Eq. (145) and its complex conjugate, taking special care of the analytic continuation of the logarithms resulting from the $`p_0^E`$-integration, tend to
$`g^2{\displaystyle ^K}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}ip_0^E}}|_{p_0^EiE_𝐩^A}`$
$`g^2{\displaystyle ^K}{\displaystyle \frac{d^3p^{}}{(2\pi )^3}}{\displaystyle \frac{1}{2E_𝐩^{}^A\mathrm{\hspace{0.17em}2}\omega _{𝐩𝐩^{}}}}{\displaystyle \frac{iE_𝐩^{}^A\gamma _0^E+𝐩^{}𝜸^E+m_A}{E_𝐩^{}^A+\omega _{𝐩𝐩^{}}+ip_0^E}}|_{p_0^EiE_𝐩^A}.`$ (148)
The first corrections to this result, of relative order $`(K/m_A)(K/\mathrm{\Lambda })`$ \[which can be suppressed in the limit of large $`\mathrm{\Lambda }`$ through a suitable choice of $`K`$, e.g., $`K=(m_A^3\mathrm{\Lambda })^{1/4}`$\], cancel among the two contributions. Incidentally, the second corrections, of order $`(K/m_A)(K/\mathrm{\Lambda })^2`$, also vanish.
Now, for the other integration region $`K^2<𝐩^2<\mathrm{\Lambda }^2`$, we have $`𝐩^2,m_A^2,\mu ^2𝐩^2,\mathrm{\Lambda }^2`$. A lengthy calculation leads to the following result: the leading terms in a systematic expansion vanish for both contributions, the subleading terms cancel among the two contributions for the $`\gamma _0^E`$-coefficients or as a result of the $`(𝐩^{}𝐩^{})`$-symmetry of integrand, integration measure and domain for the $`𝜸^E`$-coefficient, and it is hence the subsubleading terms that give the dominant contributions which turn out to be
$$\begin{array}{c}g^2_K^\mathrm{\Lambda }\frac{d^3p^{}}{(2\pi )^3}\frac{p_0^E\gamma _0^E+𝐩𝜸^E+2m_A}{4\pi 𝐩^2}\left[\frac{\pi }{2|𝐩^{}|}\frac{\mathrm{arcsin}(|𝐩^{}|/\mathrm{\Lambda })}{|𝐩^{}|}+\frac{\sqrt{1𝐩^2/\mathrm{\Lambda }^2}}{\mathrm{\Lambda }}\right]\hfill \\ \hfill +g^2_K^\mathrm{\Lambda }\frac{d^3p^{}}{(2\pi )^3}\frac{3p_0^E\gamma _0^E𝐩𝜸^E}{6\pi \mathrm{\Lambda }^2}\frac{\sqrt{1𝐩^2/\mathrm{\Lambda }^2}}{\mathrm{\Lambda }}.\end{array}$$
(149)
The first term in square brackets in Eq. (149) coincides with the dominant contribution from the part of $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$ that stems from the second integration region $`K^2<𝐩^2<\mathrm{\Lambda }^2`$ \[cf. Eq. (146)\], and hence, when added to Eq. (148), combines to $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$. The other terms in Eq. (149) are readily integrated in the limit $`K/\mathrm{\Lambda }0`$ through the substitution $`y=|𝐩^{}|/\mathrm{\Lambda }`$. The final result is
$$\begin{array}{c}G_\mathrm{\Lambda }^{\text{NC}}(𝐩)=g^2^\mathrm{\Lambda }\frac{d^4p^E}{(2\pi )^4}\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{[(p_0^Ep_0^E)^2+(𝐩𝐩^{})^2+\mu ^2][(p_0^E)^2+𝐩^{}{}_{}{}^{2}+m_A^2]}|_{p_0^EiE_𝐩^A}\hfill \\ \hfill +\frac{g^2}{(4\pi )^2}\left(\frac{1}{2}\mathrm{ln}2\right)\left(E_𝐩^A\gamma _0𝐩𝜸+2m_A\right)\frac{g^2}{(4\pi )^2}\frac{1}{12}\left(3E_𝐩^A\gamma _0+𝐩𝜸\right)\end{array}$$
(150)
in the limit $`\mathrm{\Lambda }\mathrm{}`$.
The result (150) has two important consequences: first, together with the expressions (159) and (160) for $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ below, it provides an explicit expression for $`G_\mathrm{\Lambda }^{\text{NC}}(𝐩)`$ in the limit of large $`\mathrm{\Lambda }`$. Second, and more importantly, it shows that Lorentz invariance *is* broken, even in the limit $`\mathrm{\Lambda }\mathrm{}`$, through the last term in Eq. (150). For example, it can explicitly be shown that this term leads to a non-covariant contribution to the energies of the one-particle states by following the procedure detailed in Eqs. (31)–(34). It is then clear that the non-covariant cutoff regularization is not suitable for a fermionic theory. On the other hand, in a purely scalar theory, given by the $`m_A`$-coefficient in Eq. (150), Lorentz invariance is recovered in the limit $`\mathrm{\Lambda }\mathrm{}`$, in agreement with the result in Ref. which was obtained by a different method.
Let us now return to the expression for $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ given in Eq. (143). In order to calculate $`G_0(m_A^2)`$ and $`G_1(m_A^2)`$ in cutoff regularization, one could use four-dimensional spherical coordinates for the $`d^4p^E`$-integration (see also below), which would also explicitly confirm that $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ is of the form of Eq. (28) for any finite value of $`\mathrm{\Lambda }`$. However, as long as one is only interested in the result for large enough $`\mathrm{\Lambda }`$, it is much quicker in the present situation to introduce Feynman parameters and proceed in analogy with Eq. (88) \[in $`(D=4)`$-dimensional Euclidean space and for the integration domain $`(p^E)^2\mathrm{\Lambda }^2`$\]. We thus arrive at
$`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)=g^2{\displaystyle _0^1}𝑑x{\displaystyle ^\mathrm{\Lambda }^{}}{\displaystyle \frac{d^4q^E}{(2\pi )^4}}{\displaystyle \frac{(q^E+xp^E)\gamma ^E+m_A}{\left[(q^E)^2+x(1x)(p^E)^2+x\mu ^2+(1x)m_A^2\right]^2}}|_{p_0^EiE_𝐩^A}`$ (151)
with $`q^E`$ to be integrated over the four-dimensional Euclidean domain $`(q^E+xp^E)^2\mathrm{\Lambda }^2`$ (as a result of the shift of the integration variable to $`q^E=p^Exp^E`$). To evaluate this expression it is easiest to calculate the difference
$`\mathrm{\Delta }G_{\mathrm{\Lambda },\mathrm{\Lambda }^{}}^{\text{MC}}(𝐩)=G_\mathrm{\Lambda }^{\text{MC}}(𝐩)G_\mathrm{\Lambda }^{}^{\text{MC}}(𝐩)`$ (152)
between $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$ and
$`G_\mathrm{\Lambda }^{}^{\text{MC}}(𝐩)=g^2{\displaystyle _0^1}𝑑x{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4q^E}{(2\pi )^4}}{\displaystyle \frac{xp^E\gamma ^E+m_A}{\left[(q^E)^2+x(1x)(p^E)^2+x\mu ^2+(1x)m_A^2\right]^2}}|_{p_0^EiE_𝐩^A}`$ (153)
with the integration domain restricted to $`(q^E)^2\mathrm{\Lambda }^2`$ \[observe that in this case the first term in the numerator in Eq. (151) does not contribute because of the symmetry of the integration domain\]. The $`d^4q^E`$-integral in $`G_\mathrm{\Lambda }^{}^{\text{MC}}(𝐩)`$ can be evaluated by standard methods.
Incidentally, it would in principle be possible to define cutoff regularization via Eq. (153) \[instead of Eq. (151)\], a choice that would obviously fulfill Eq. (28), too. The non-covariant expression corresponding to Eq. (141) in this alternative regularization is
$$\begin{array}{c}G_\mathrm{\Lambda }^{}^{\text{MC}}(𝐩)=\hfill \\ \hfill g^2_0^1𝑑x^\mathrm{\Lambda }^{}\frac{d^4p^E}{(2\pi )^4}\frac{dk_0^E}{2\pi }\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{\left[\left(k_0^E\right)^2x+\omega _{𝐩𝐩^{}}^2x+\left(p_0^E\right)^2(1x)+\left(E_𝐩^{}^A\right)^2(1x)\right]^2}\\ \hfill \times \left[_{\mathrm{}}^0𝑑t^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}+_{\mathrm{}}^0𝑑t^Ee^{ϵ^Et^E}e^{i(p_0^E+k_0^Ep_0^E)t^E}\right]_{p_0^EiE_𝐩^A},\end{array}$$
(154)
where the $`d^4p^E`$-integration is now over $`(p^Exp^E)^2\mathrm{\Lambda }^2`$. When comparing this expression with Eq. (141), it becomes clear that this regularization is not too natural in the present non-covariant context. Still, we can integrate over $`k_0^E`$ (again with the help of the residue theorem) and $`t_E`$ in Eq. (154). The result for the first contribution is
$$\begin{array}{c}g^2_0^1𝑑x^\mathrm{\Lambda }^{}\frac{d^4p^E}{(2\pi )^4}\frac{p_0^E\gamma _0^E+𝐩^{}𝜸^E+m_A}{\left[\left(p_0^Ep_0^E\right)^2x+\omega _{𝐩𝐩^{}}^2x+\left(p_0^E\right)^2(1x)+\left(E_𝐩^{}^A\right)^2(1x)\right]^2}\hfill \\ \hfill \times \left\{\frac{1}{2}+i\frac{(p_0^Ep_0^E)\sqrt{x}\left[\left(p_0^Ep_0^E\right)^2x+3\omega _{𝐩𝐩^{}}^2x+3\left(p_0^E\right)^2(1x)+3\left(E_𝐩^{}^A\right)^2(1x)\right]}{4\left[\omega _{𝐩𝐩^{}}^2x+\left(p_0^E\right)^2(1x)+\left(E_𝐩^{}^A\right)^2(1x)\right]^{3/2}}\right\},\end{array}$$
(155)
to be analytically continued to $`p_0^EiE_𝐩^A`$. The result for the second contribution can again be obtained by complex conjugation from the above (considering $`p_0^E`$ and the Euclidean $`\gamma `$-matrices as real). The integrations over both $`x`$ and $`p_0^E`$ in Eq. (155) look forbidding. However, the sum of the two non-covariant contributions is readily seen to give $`G_\mathrm{\Lambda }^{}^{\text{MC}}(𝐩)`$ as defined in Eq. (153), after shifting the four-momentum integration variable to $`q^E=p^Exp^E`$.
Returning to the explicit calculation of $`G_\mathrm{\Lambda }^{\text{MC}}(𝐩)`$, we note that the difference $`\mathrm{\Delta }G_{\mathrm{\Lambda },\mathrm{\Lambda }^{}}^{\text{MC}}(𝐩)`$ is an integral over the difference of the four-balls $`(q^E+xp^E)^2\mathrm{\Lambda }^2`$ and $`(q^E)^2\mathrm{\Lambda }^2`$. We introduce four-dimensional (Euclidean) spherical coordinates with the fourth axis oriented in direction of $`p^E`$ and the corresponding polar angle denoted as $`\chi `$. Then for the four-ball $`(q^E+xp^E)^2\mathrm{\Lambda }^2`$ the integration along the radius runs up to $`R(\chi )`$,
$`R(\chi )=\mathrm{\Lambda }x|p^E|\mathrm{cos}\chi +𝒪(1/\mathrm{\Lambda })`$ (156)
in the limit of large $`\mathrm{\Lambda }`$. Furthermore, the denominator of the integrand in the region between the two four-balls can be approximated by $`\mathrm{\Lambda }^4`$ in this limit. We then have
$`\mathrm{\Delta }G_{\mathrm{\Lambda },\mathrm{\Lambda }^{}}^{\text{MC}}(𝐩)=g^2{\displaystyle _0^1}𝑑x\mathrm{\hspace{0.17em}4}\pi {\displaystyle _0^\pi }𝑑\chi \mathrm{sin}^2\chi {\displaystyle _\mathrm{\Lambda }^{R(\chi )}}𝑑rr^3{\displaystyle \frac{r\mathrm{cos}\chi \gamma _4^E+x(p^E\gamma ^E)+m_A}{(2\pi )^4\mathrm{\Lambda }^4}},`$ (157)
while the other components of $`q^E`$ in the numerator do not contribute for symmetry reasons (integration over the other angular coordinates). The integrations in Eq. (157) yield, in the limit of large $`\mathrm{\Lambda }`$,
$`\mathrm{\Delta }G_{\mathrm{\Lambda },\mathrm{\Lambda }^{}}^{\text{MC}}(𝐩)={\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle \frac{|p^E|\gamma _4^E}{4}}={\displaystyle \frac{g^2}{(4\pi )^2}}{\displaystyle \frac{p\gamma }{4}},`$ (158)
where we have used $`|p^E|\gamma _4^E=p^E\gamma ^E`$ and replaced $`p_0^EiE_𝐩^Aip^0`$. In particular, for a purely scalar theory, the difference $`\mathrm{\Delta }G_{\mathrm{\Lambda },\mathrm{\Lambda }^{}}^{\text{MC}}(𝐩)`$ tends to zero for $`\mathrm{\Lambda }\mathrm{}`$.
Now from Eqs. (152) and (153), the above result (158) and the result of the integration over the Feynman parameter $`x`$ in Eqs. (92) and (93), we finally get the explicit expressions
$`G_{0,\mathrm{\Lambda }}^{\text{MC}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`2\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\},`$ (159)
$`G_{1,\mathrm{\Lambda }}^{\text{MC}}(m_A^2)`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+{\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$
$`(2{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}`$ (160)
($`\mathrm{\Lambda }^2m_A^2`$). Again, the result for $`G_{0,\mathrm{\Lambda }}^{\text{MC}}(m_A^2)`$ gives directly the expression for $`\mathrm{\Delta }m_A^2`$ in momentum cutoff regularization for the purely scalar case, while for Yukawa theory we have
$`{\displaystyle \frac{\mathrm{\Delta }m_A^2}{2m_A^2}}`$ $`={\displaystyle \frac{g^2}{(4\pi )^2}}\{{\displaystyle \frac{3}{2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_A^2}}+{\displaystyle \frac{7}{4}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}({\displaystyle \frac{3}{2}}{\displaystyle \frac{\mu ^2}{m_A^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^4}{m_A^4}})\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_A^2}}`$ (161)
$`(4{\displaystyle \frac{\mu ^2}{m_A^2}})\sqrt{{\displaystyle \frac{\mu ^2}{m_A^2}}\left(1{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu ^2}{m_A^2}}\right)}\text{arcctg}\sqrt{{\displaystyle \frac{\mu ^2/(4m_A^2)}{1\mu ^2/(4m_A^2)}}}\}.`$ (162)
## Appendix C Separation of angular variables and spin
In terms of the well–known eigenstates $`\chi _{S,m_S}`$ of total spin $`𝐒=𝐬_A+𝐬_B`$ ($`S=0`$ or $`S=1`$), one has the following explicit expressions for the eigenstates of $`𝐉^2`$ and $`J_z`$ with eigenvalues $`J(J+1)`$ and $`M`$:
$`{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=Y_{JM}(\widehat{𝐩})\chi _{00},`$
$`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})`$ $`={\displaystyle \frac{1}{\sqrt{2J(2J1)}}}[\sqrt{(JM1)(JM)}Y_{J1,M+1}(\widehat{𝐩})\chi _{1,1}`$
$`+\sqrt{2(JM)(J+M)}Y_{J1,M}(\widehat{𝐩})\chi _{10}`$
$`+\sqrt{(J+M1)(J+M)}Y_{J1,M1}(\widehat{𝐩})\chi _{11}](J1),`$
$`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`={\displaystyle \frac{1}{\sqrt{2J(J+1)}}}[\sqrt{(JM)(J+M+1)}Y_{J,M+1}(\widehat{𝐩})\chi _{1,1}`$
$`+\sqrt{2}MY_{JM}(\widehat{𝐩})\chi _{10}`$
$`\sqrt{(JM+1)(J+M)}Y_{J,M1}(\widehat{𝐩})\chi _{11}](J1),`$
$`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩})`$ $`={\displaystyle \frac{1}{\sqrt{2(J+1)(2J+3)}}}[\sqrt{(J+M+1)(J+M+2)}Y_{J+1,M+1}(\widehat{𝐩})\chi _{1,1}`$
$`\sqrt{2(JM+1)(J+M+1)}Y_{J+1,M}(\widehat{𝐩})\chi _{10}`$
$`+\sqrt{(JM+1)(JM+2)}Y_{J+1,M1}(\widehat{𝐩})\chi _{11}].`$ (163)
To determine the action of the helicity operators on these eigenfunctions, the explicit spinor representation
$`\widehat{𝐩}𝝈=\left(\begin{array}{cc}\mathrm{cos}\vartheta & \mathrm{sin}\vartheta e^{i\phi }\\ \mathrm{sin}\vartheta e^{i\phi }& \mathrm{cos}\vartheta \end{array}\right)`$ (166)
can be used, together with the following special instances of the spherical harmonics addition relation :
$`\mathrm{cos}\vartheta Y_{lm}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{(lm)(l+m)}{(2l1)(2l+1)}}}Y_{l1,m}(\widehat{𝐩})`$
$`+\sqrt{{\displaystyle \frac{(lm+1)(l+m+1)}{(2l+1)(2l+3)}}}Y_{l+1,m}(\widehat{𝐩}),`$
$`\mathrm{sin}\vartheta e^{i\phi }Y_{lm}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{(lm1)(lm)}{(2l1)(2l+1)}}}Y_{l1,m+1}(\widehat{𝐩})`$
$`\sqrt{{\displaystyle \frac{(l+m+1)(l+m+2)}{(2l+1)(2l+3)}}}Y_{l+1,m+1}(\widehat{𝐩}),`$
$`\mathrm{sin}\vartheta e^{i\phi }Y_{lm}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{(l+m1)(l+m)}{(2l1)(2l+1)}}}Y_{l1,m1}(\widehat{𝐩})`$
$`+\sqrt{{\displaystyle \frac{(lm+1)(lm+2)}{(2l+1)(2l+3)}}}Y_{l+1,m1}(\widehat{𝐩}).`$ (167)
It is then a straightforward, though somewhat lengthy exercise to derive the following formulas for the application of the helicity operators to the total angular momentum eigenstates:
$`(\widehat{𝐩}𝝈_A){}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_A){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_A){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_A){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}),`$ (168)
and
$`(\widehat{𝐩}𝝈_B){}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})+\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_B){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_B){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈_B){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}).`$ (169)
In the special case $`J=0`$, the states $`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})`$ and $`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩})`$ do not exist, and on the right–hand sides for the application of one of the helicity operators to $`{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})`$ and $`{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩})`$, only one term remains.
By use of Eq. (77), where now
$$\begin{array}{c}a_l(p,p^{})=\frac{2l+1}{2}_1^1d\mathrm{cos}\theta P_l(\mathrm{cos}\theta )\frac{1}{2|𝐩𝐩^{}|}\hfill \\ \hfill \times \left(\frac{1}{E_p^A+|𝐩𝐩^{}|E_p^{}^A}+\frac{1}{E_p^B+|𝐩𝐩^{}|E_p^{}^B}\right)\end{array}$$
(170)
(for $`\mu =0`$), the effective Schrödinger equation Eq. (46) decouples into pairs of coupled one-dimensional integral equations. For the S-coupled states, we introduce the wave function
$`\varphi (𝐩)={}_{}{}^{S}\varphi _{0}^{J}(p){}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})+{}_{}{}^{S}\varphi _{1}^{J}(p){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}).`$ (171)
The effective Schrödinger equation for the coefficient functions becomes
$$\begin{array}{c}\left(\sqrt{M_A^2+p^2}+\sqrt{M_B^2+p^2}\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{0}^{J}(p)\\ {}_{}{}^{S}\varphi _{1}^{J}(p)\end{array}\right)\hfill \\ \hfill \frac{g^2}{2\pi ^2}_0^{\mathrm{}}𝑑p^{}p_{}^{}{}_{}{}^{2}\sqrt{\frac{E_p^A+M_A}{2E_p^A}\frac{E_p^B+M_B}{2E_p^B}\frac{E_p^{}^A+M_A}{2E_p^{}^A}\frac{E_p^{}^B+M_B}{2E_p^{}^B}}\\ \hfill \times \frac{1}{2J+1}\left(\begin{array}{cc}{}_{}{}^{S}V_{00}^{J}(p,p^{})& {}_{}{}^{S}V_{01}^{J}(p,p^{})\\ {}_{}{}^{S}V_{10}^{J}(p,p^{})& {}_{}{}^{S}V_{11}^{J}(p,p^{})\end{array}\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{0}^{J}(p^{})\\ {}_{}{}^{S}\varphi _{1}^{J}(p^{})\end{array}\right)\\ \hfill =\left(EE_V\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{0}^{J}(p)\\ {}_{}{}^{S}\varphi _{1}^{}(p)\end{array}\right),\end{array}$$
(172)
with
$`{}_{}{}^{S}V_{00}^{J}(p,p^{})`$ $`=\left[1+{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]a_J(p,p^{})`$
$`\left[{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}+{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]`$
$`\times \left[J{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+(J+1){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{S}V_{11}^{J}(p,p^{})`$ $`=\left[1+{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]a_J(p,p^{})`$
$`\left[{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}+{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]`$
$`\times \left[(J+1){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+J{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{S}V_{01}^{J}(p,p^{})`$ $`={}_{}{}^{S}V_{10}^{J}(p,p^{})`$
$`=\left[{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]`$
$`\times \sqrt{J(J+1)}\left[{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right].`$ (173)
On the other hand, with the wave function
$`\varphi (𝐩)={}_{}{}^{L}\varphi _{J1}^{J}(p){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J1,M}^{J}(\widehat{𝐩})+{}_{}{}^{L}\varphi _{J+1}^{J}(p){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{J+1,M}^{J}(\widehat{𝐩}),`$ (174)
the effective Schrödinger equation for the L-coupled states becomes
$$\begin{array}{c}\left(\sqrt{M_A^2+p^2}+\sqrt{M_B^2+p^2}\right)\left(\begin{array}{c}{}_{}{}^{L}\varphi _{J1}^{J}(p)\\ {}_{}{}^{L}\varphi _{J+1}^{J}(p)\end{array}\right)\hfill \\ \hfill \frac{g^2}{2\pi ^2}_0^{\mathrm{}}𝑑p^{}p_{}^{}{}_{}{}^{2}\sqrt{\frac{E_p^A+M_A}{2E_p^A}\frac{E_p^B+M_B}{2E_p^B}\frac{E_p^{}^A+M_A}{2E_p^{}^A}\frac{E_p^{}^B+M_B}{2E_p^{}^B}}\\ \hfill \times \frac{1}{2J+1}\left(\begin{array}{cc}{}_{}{}^{L}V_{J1,J1}^{J}(p,p^{})& {}_{}{}^{L}V_{J1,J+1}^{J}(p,p^{})\\ {}_{}{}^{L}V_{J+1,J1}^{J}(p,p^{})& {}_{}{}^{L}V_{J+1,J+1}^{J}(p,p^{})\end{array}\right)\left(\begin{array}{c}{}_{}{}^{L}\varphi _{J1}^{J}(p^{})\\ {}_{}{}^{L}\varphi _{J+1}^{J}(p^{})\end{array}\right)\\ \hfill =\left(EE_V\right)\left(\begin{array}{c}{}_{}{}^{L}\varphi _{J1}^{J}(p)\\ {}_{}{}^{L}\varphi _{J+1}^{J}(p)\end{array}\right),\end{array}$$
(175)
where
$`{}_{}{}^{L}V_{J1,J1}^{J}(p,p^{})`$ $`=(2J+1){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}`$
$`\left[{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}+{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]a_J(p,p^{})`$
$`+{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}`$
$`\times {\displaystyle \frac{1}{2J+1}}\left[{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+4J(J+1){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{L}V_{J+1,J+1}^{J}(p,p^{})`$ $`=(2J+1){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}`$
$`\left[{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}+{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}\right]a_J(p,p^{})`$
$`+{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}`$
$`\times {\displaystyle \frac{1}{2J+1}}\left[4J(J+1){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{L}V_{J1,J+1}^{}(p,p^{})`$ $`={}_{}{}^{L}V_{J+1,J1}^{}(p,p^{})`$
$`={\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}`$
$`\times {\displaystyle \frac{2\sqrt{J(J+1)}}{2J+1}}\left[{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right].`$ (176)
In the special case $`J=0`$, $`{}_{}{}^{S}\varphi _{1}^{J}(p){}_{}{}^{L}\varphi _{J1}^{J}(p)0`$, and the equations for $`{}_{}{}^{S}\varphi _{0}^{J}(p)`$ and $`{}_{}{}^{L}\varphi _{J+1}^{J}(p)`$ decouple.
The functions $`a_l(p,p^{})`$ are quite straightforwardly calculated from Eq. (170), for any given value of $`l`$. The results for the functions which are relevant for the numerical calculations presented here ($`l3`$) are:
$`a_0(p,p^{})`$ $`={\displaystyle \frac{1}{4pp^{}}}\left\{\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)+\mathrm{ln}\left({\displaystyle \frac{E_p^B+(p+p^{})E_p^{}^B}{E_p^B+|pp^{}|E_p^{}^B}}\right)\right\},`$
$`a_1(p,p^{})`$ $`={\displaystyle \frac{3}{4pp^{}}}\{{\displaystyle \frac{(E_p^AE_p^{}^A+E_p^BE_p^{}^B)(p+p^{}|pp^{}|)}{2pp^{}}}2`$
$`+{\displaystyle \frac{p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2}{2pp^{}}}\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)`$
$`+{\displaystyle \frac{p^2+p_{}^{}{}_{}{}^{2}(E_p^BE_p^{}^B)^2}{2pp^{}}}\mathrm{ln}\left({\displaystyle \frac{E_p^B+(p+p^{})E_p^{}^B}{E_p^B+|pp^{}|E_p^{}^B}}\right)\},`$
$`a_2(p,p^{})`$ $`={\displaystyle \frac{5}{8pp^{}}}\{{\displaystyle \frac{3[2(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2](E_p^AE_p^{}^A)(p+p^{}|pp^{}|)}{4p^2p_{}^{}{}_{}{}^{2}}}`$
$`+{\displaystyle \frac{3[2(p^2+p_{}^{}{}_{}{}^{2})(E_p^BE_p^{}^B)^2](E_p^BE_p^{}^B)(p+p^{}|pp^{}|)}{4p^2p_{}^{}{}_{}{}^{2}}}`$
$`{\displaystyle \frac{(E_p^AE_p^{}^A+E_p^BE_p^{}^B)[(p+p^{})^3|pp^{}|^3]}{4p^2p_{}^{}{}_{}{}^{2}}}`$
$`{\displaystyle \frac{3[2(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2(E_p^BE_p^{}^B)^2]}{2pp^{}}}`$
$`+\left[{\displaystyle \frac{3[p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2]^2}{4p^2p_{}^{}{}_{}{}^{2}}}1\right]\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)`$
$`+[{\displaystyle \frac{3[p^2+p_{}^{}{}_{}{}^{2}(E_p^BE_p^{}^B)^2]^2}{4p^2p_{}^{}{}_{}{}^{2}}}1]\mathrm{ln}\left({\displaystyle \frac{E_p^B+(p+p^{})E_p^{}^B}{E_p^B+|pp^{}|E_p^{}^B}}\right)\},`$
$`a_3(p,p^{})`$ $`={\displaystyle \frac{7}{8pp^{}}}\{[{\displaystyle \frac{5[3(p^2+p_{}^{}{}_{}{}^{2})^23(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2+(E_p^AE_p^{}^A)^4](E_p^AE_p^{}^A)}{8p^3p_{}^{}{}_{}{}^{3}}}`$
$`+{\displaystyle \frac{5[3(p^2+p_{}^{}{}_{}{}^{2})^23(p^2+p_{}^{}{}_{}{}^{2})(E_p^BE_p^{}^B)^2+(E_p^BE_p^{}^B)^4](E_p^BE_p^{}^B)}{8p^3p_{}^{}{}_{}{}^{3}}}`$
$`{\displaystyle \frac{3(E_p^AE_p^{}^A+E_p^BE_p^{}^B)}{2pp^{}}}](p+p^{}|pp^{}|)`$
$`[{\displaystyle \frac{5(E_p^AE_p^{}^A)[3(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2]}{24p^3p_{}^{}{}_{}{}^{3}}}`$
$`+{\displaystyle \frac{5(E_p^BE_p^{}^B)[3(p^2+p_{}^{}{}_{}{}^{2})(E_p^BE_p^{}^B)^2]}{24p^3p_{}^{}{}_{}{}^{3}}}][(p+p^{})^3|pp^{}|^3]`$
$`+{\displaystyle \frac{(E_p^AE_p^{}^A+E_p^BE_p^{}^B)[(p+p^{})^5|pp^{}|^5]}{8p^3p_{}^{}{}_{}{}^{3}}}`$
$`{\displaystyle \frac{5[p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2]^2+5[p^2+p_{}^{}{}_{}{}^{2}(E_p^BE_p^{}^B)^2]^2}{4p^2p_{}^{}{}_{}{}^{2}}}+{\displaystyle \frac{8}{3}}`$
$`+\left[{\displaystyle \frac{5[p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2]^3}{8p^3p_{}^{}{}_{}{}^{3}}}{\displaystyle \frac{3[p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2]}{2pp^{}}}\right]`$
$`\times \mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)`$
$`+\left[{\displaystyle \frac{5[p^2+p_{}^{}{}_{}{}^{2}(E_p^BE_p^{}^B)^2]^3}{8p^3p_{}^{}{}_{}{}^{3}}}{\displaystyle \frac{3[p^2+p_{}^{}{}_{}{}^{2}(E_p^BE_p^{}^B)^2]}{2pp^{}}}\right]`$
$`\times \mathrm{ln}\left({\displaystyle \frac{E_p^B+(p+p^{})E_p^{}^B}{E_p^B+|pp^{}|E_p^{}^B}}\right)\}.`$ (177)
In the actual numerical calculations we did not rely on these explicit expressions, but rather used a quasi-algebraic method which increases speed and accuracy. To this end, the integrand in Eq. (170) is written, after changing variables from $`\mathrm{cos}\theta `$ to $`x=|𝐩𝐩^{}|`$, as a (finite) Laurent series in $`x`$ around $`E_p^{}^AE_p^A`$ (or $`E_p^{}^BE_p^B`$ for the second term). This requires expressing $`P_l(\mathrm{cos}\theta )`$ as a polynomial in $`x`$, a task left to the computer. Each term in the Laurent series has a known integral which is simply inserted. Any value of $`l`$ can be handled easily by this method.
In the one-body limit $`M_B\mathrm{}`$, the matrix for the effective potential in the S-coupled sector, Eq. (173), tends to
$$\begin{array}{c}\left(\begin{array}{cc}{}_{}{}^{S}V_{00}^{J}(p,p^{})& {}_{}{}^{S}V_{01}^{J}(p,p^{})\\ v{}_{}{}^{S}V_{10}^{J}(p,p^{})& {}_{}{}^{S}V_{11}^{J}(p,p^{})\end{array}\right)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)a_J(p,p^{})\hfill \\ \hfill \frac{p}{E_p^A+M_A}\frac{p^{}}{E_p^{}^A+M_A}[\left(\begin{array}{cc}J& \sqrt{J(J+1)}\\ \sqrt{J(J+1)}& J+1\end{array}\right)\frac{a_{J1}(p,p^{})}{2J1}\\ \hfill +\left(\begin{array}{cc}J+1& \sqrt{J(J+1)}\\ \sqrt{J(J+1)}& J\end{array}\right)\frac{a_{J+1}(p,p^{})}{2J+3}].\end{array}$$
(178)
This matrix is diagonalized by the orthogonal linear combinations
$`\left(\begin{array}{c}{}_{}{}^{S}\varphi _{0}^{J}(p)\\ {}_{}{}^{S}\varphi _{1}^{J}(p)\end{array}\right)=\left(\begin{array}{c}\sqrt{J}\\ \sqrt{J+1}\end{array}\right){\displaystyle \frac{{}_{}{}^{S}\varphi _{J1/2}^{J}(p)}{\sqrt{2J+1}}}`$ (183)
and
$`\left(\begin{array}{c}{}_{}{}^{S}\varphi _{0}^{J}(p)\\ {}_{}{}^{S}\varphi _{1}^{J}(p)\end{array}\right)=\left(\begin{array}{c}\sqrt{J+1}\\ \sqrt{J}\end{array}\right){\displaystyle \frac{{}_{}{}^{S}\varphi _{J+1/2}^{J}(p)}{\sqrt{2J+1}}},`$ (188)
with the respective eigenvalues
$`(2J+1)\left({\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}\right),`$
$`(2J+1)\left({\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right).`$ (189)
Alternatively, we can write the wave function as
$`\varphi (𝐩)={}_{}{}^{S}\varphi _{J1/2}^{J}(p){}_{J1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ (190)
or
$`\varphi (𝐩)={}_{}{}^{S}\varphi _{J+1/2}^{J}(p){}_{J+1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ (191)
with
$`{}_{J1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})+\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}),`$
$`{}_{J+1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+1}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}1}}𝒴_{JM}^{J}(\widehat{𝐩})+\sqrt{{\displaystyle \frac{J}{2J+1}}}{}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J}(\widehat{𝐩}),`$ (192)
anticipating the notation $`{}_{j_A}{}^{}𝒴_{lM}^{J}(\widehat{𝐩})`$. In the special case $`J=0`$, of course, $`{}_{}{}^{S}\varphi _{1}^{J}(p)0`$, $`{}_{}{}^{S}\varphi _{J1/2}^{J}(p)0`$, and $`{}_{J1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ does not exist, since there is only one state in the S-coupled $`J=0`$ sector.
The matrix (176) of the effective potential for L-coupled states, on the other hand, becomes diagonal in the limit $`M_B\mathrm{}`$, hence there is no L-coupling in this limit. The diagonal matrix elements tend to
$`{}_{}{}^{L}V_{J1,J1}^{J}(p,p^{})`$ $`=(2J+1)\left({\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{a_J(p,p^{})}{2J+1}}\right),`$
$`{}_{}{}^{L}V_{J+1,J+1}^{J}(p,p^{})`$ $`=(2J+1)\left({\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}{\displaystyle \frac{p}{E_p^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{a_J(p,p^{})}{2J+1}}\right).`$ (193)
Then, from Eqs. (172) and (175), $`{}_{}{}^{S}\varphi _{J1/2}^{J}(p)`$ in Eq. (183) and $`{}_{}{}^{L}\varphi _{J1/2}^{J1}(p){}_{}{}^{L}\varphi _{J}^{J1}(p)`$ in Eq. (174),
$`\varphi (𝐩)={}_{}{}^{L}\varphi _{J1/2}^{J1}(p){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J1}(\widehat{𝐩}),`$ (194)
fulfill exactly the same one-dimensional equation in the limit $`M_B\mathrm{}`$ (for $`J1`$),
$$\begin{array}{c}\sqrt{M_A^2+p^2}\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J1/2}^{J}(p)\\ {}_{}{}^{L}\varphi _{J1/2}^{J1}(p)\end{array}\right)\frac{g^2}{2\pi ^2}_0^{\mathrm{}}𝑑p^{}p_{}^{}{}_{}{}^{2}\sqrt{\frac{E_p^A+M_A}{2E_p^A}\frac{E_p^{}^A+M_A}{2E_p^{}^A}}\hfill \\ \hfill \times \left(\frac{a_J(p,p^{})}{2J+1}\frac{p}{E_p^A+M_A}\frac{p^{}}{E_p^{}^A+M_A}\frac{a_{J1}(p,p^{})}{2J1}\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J1/2}^{J}(p^{})\\ {}_{}{}^{L}\varphi _{J1/2}^{J1}(p^{})\end{array}\right)\\ \hfill =\left(EE_VM_B\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J1/2}^{J}(p)\\ {}_{}{}^{L}\varphi _{J1/2}^{J1}(p)\end{array}\right).\end{array}$$
(195)
Likewise, $`{}_{}{}^{S}\varphi _{J+1/2}^{J}(p)`$ in Eq. (188) and $`{}_{}{}^{L}\varphi _{J+1/2}^{J+1}(p){}_{}{}^{L}\varphi _{J}^{J+1}(p)`$ in Eq. (174),
$`\varphi (𝐩)={}_{}{}^{L}\varphi _{J+1/2}^{J+1}(p){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J+1}(\widehat{𝐩}),`$ (196)
both fulfill the equation
$$\begin{array}{c}\sqrt{M_A^2+p^2}\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J+1/2}^{J}(p)\\ {}_{}{}^{L}\varphi _{J+1/2}^{J+1}(p)\end{array}\right)\frac{g^2}{2\pi ^2}_0^{\mathrm{}}𝑑p^{}p_{}^{}{}_{}{}^{2}\sqrt{\frac{E_p^A+M_A}{2E_p^A}\frac{E_p^{}^A+M_A}{2E_p^{}^A}}\hfill \\ \hfill \times \left(\frac{a_J(p,p^{})}{2J+1}\frac{p}{E_p^A+M_A}\frac{p^{}}{E_p^{}^A+M_A}\frac{a_{J+1}(p,p^{})}{2J+3}\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J+1/2}^{J}(p^{})\\ {}_{}{}^{L}\varphi _{J+1/2}^{J+1}(p^{})\end{array}\right)\\ \hfill =\left(EE_VM_B\right)\left(\begin{array}{c}{}_{}{}^{S}\varphi _{J+1/2}^{J}(p)\\ {}_{}{}^{L}\varphi _{J+1/2}^{J+1}(p)\end{array}\right).\end{array}$$
(197)
As a result, every state is (at least) twofold degenerate in the one-body limit.
The states $`{}_{J1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ and $`{}_{J+1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ defined in Eq. (192), as well as $`{}_{J1/2}{}^{}𝒴_{JM}^{J1}(\widehat{𝐩}){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J1}(\widehat{𝐩})`$ and $`{}_{J+1/2}{}^{}𝒴_{JM}^{J+1}(\widehat{𝐩}){}_{}{}^{\mathrm{\hspace{0.17em}3}}𝒴_{JM}^{J+1}(\widehat{𝐩})`$, can be rewritten in terms of the eigenstates $`Y_{lm}^{j_A}(\widehat{𝐩})`$ of $`𝐣_A^2`$, $`j_{A,z}`$, and $`𝐋^2`$ (where $`𝐣_A=𝐋+𝐬_A`$),
$`Y_{lm}^{l1/2}(\widehat{𝐩})`$ $`={\displaystyle \frac{1}{\sqrt{2l+1}}}\left(\begin{array}{c}\sqrt{lm+1/2}Y_{l,m1/2}(\widehat{𝐩})\\ \sqrt{l+m+1/2}Y_{l,m+1/2}(\widehat{𝐩})\end{array}\right)(l1),`$ (200)
$`Y_{lm}^{l+1/2}(\widehat{𝐩})`$ $`={\displaystyle \frac{1}{\sqrt{2l+1}}}\left(\begin{array}{c}\sqrt{l+m+1/2}Y_{l,m1/2}(\widehat{𝐩})\\ \sqrt{lm+1/2}Y_{l,m+1/2}(\widehat{𝐩})\end{array}\right),`$ (203)
as
$`{}_{J1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{JM}{2J}}}Y_{J,M+1/2}^{J1/2}(\widehat{𝐩})\left(\begin{array}{c}0\\ 1\end{array}\right)+\sqrt{{\displaystyle \frac{J+M}{2J}}}Y_{J,M1/2}^{J1/2}(\widehat{𝐩})\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (208)
$`{}_{J1/2}{}^{}𝒴_{JM}^{J1}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+M}{2J}}}Y_{J,M+1/2}^{J1/2}(\widehat{𝐩})\left(\begin{array}{c}0\\ 1\end{array}\right)\sqrt{{\displaystyle \frac{JM}{2J}}}Y_{J,M1/2}^{J1/2}(\widehat{𝐩})\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (213)
$`{}_{J+1/2}{}^{}𝒴_{JM}^{J}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{J+M+1}{2(J+1)}}}Y_{J,M+1/2}^{J+1/2}(\widehat{𝐩})\left(\begin{array}{c}0\\ 1\end{array}\right)\sqrt{{\displaystyle \frac{JM+1}{2(J+1)}}}Y_{J,M1/2}^{J+1/2}(\widehat{𝐩})\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (218)
$`{}_{J+1/2}{}^{}𝒴_{JM}^{J+1}(\widehat{𝐩})`$ $`=\sqrt{{\displaystyle \frac{JM+1}{2(J+1)}}}Y_{J,M+1/2}^{J+1/2}(\widehat{𝐩})\left(\begin{array}{c}0\\ 1\end{array}\right)+\sqrt{{\displaystyle \frac{J+M+1}{2(J+1)}}}Y_{J,M1/2}^{J+1/2}(\widehat{𝐩})\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (223)
hence they are simultaneous eigenstates of $`𝐉^2`$, $`J_z`$, $`𝐣_A^2`$, and $`𝐋^2`$.
Using the formulas
$`(\widehat{𝐩}𝝈)Y_{lm}^{l1/2}(\widehat{𝐩})`$ $`=Y_{l1,m}^{l1/2}(\widehat{𝐩}),`$
$`(\widehat{𝐩}𝝈)Y_{lm}^{l+1/2}(\widehat{𝐩})`$ $`=Y_{l+1,m}^{l+1/2}(\widehat{𝐩}),`$ (224)
proven, e.g., with the help of Eqs. (166) and (167), and consistent with the fact that the application of helicity operators conserves angular momentum $`𝐣_A=𝐋+𝐬_A`$ and changes spatial parity $`(1)^l`$, it is easily seen that for the wave functions (190) and (194) the effective Schrödinger equation in the one-body limit, Eq. (51), is equivalent to Eq. (195), while for the wave functions (191) and (196), Eq. (51) is equivalent to Eq. (197).
The coefficient functions $`a_l(p,p^{})`$ of the partial wave expansion in the one-body limit $`M_B\mathrm{}`$ are defined by
$$\begin{array}{c}a_l(p,p^{})=\frac{2l+1}{2}_1^1d\mathrm{cos}\theta P_l(\mathrm{cos}\theta )\frac{1}{2|𝐩𝐩^{}|}\left(\frac{1}{|𝐩𝐩^{}|}+\frac{1}{E_p^A+|𝐩𝐩^{}|E_p^{}^A}\right),\hfill \end{array}$$
(225)
see Eq. (51). The explicit expressions in the cases $`l=0,1,2`$, which are the ones relevant for the numerical results presented here, are
$`a_0(p,p^{})`$ $`={\displaystyle \frac{1}{4pp^{}}}\left\{\mathrm{ln}{\displaystyle \frac{p+p^{}}{|pp^{}|}}+\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)\right\},`$
$`a_1(p,p^{})`$ $`={\displaystyle \frac{3}{4pp^{}}}\{{\displaystyle \frac{(E_p^AE_p^{}^A)(p+p^{}|pp^{}|)}{2pp^{}}}2`$
$`+{\displaystyle \frac{p^2+p_{}^{}{}_{}{}^{2}}{2pp^{}}}\mathrm{ln}{\displaystyle \frac{p+p^{}}{|pp^{}|}}+{\displaystyle \frac{p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2}{2pp^{}}}\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)\},`$
$`a_2(p,p^{})`$ $`={\displaystyle \frac{5}{8pp^{}}}\{{\displaystyle \frac{3[2(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2](E_p^AE_p^{}^A)(p+p^{}|pp^{}|)}{4p^2p_{}^{}{}_{}{}^{2}}}`$
$`{\displaystyle \frac{(E_p^AE_p^{}^A)[(p+p^{})^3|pp^{}|^3]}{4p^2p_{}^{}{}_{}{}^{2}}}{\displaystyle \frac{3[2(p^2+p_{}^{}{}_{}{}^{2})(E_p^AE_p^{}^A)^2]}{2pp^{}}}`$
$`+\left[{\displaystyle \frac{3(p^2+p_{}^{}{}_{}{}^{2})^2}{4p^2p_{}^{}{}_{}{}^{2}}}1\right]\mathrm{ln}{\displaystyle \frac{p+p^{}}{|pp^{}|}}`$
$`+[{\displaystyle \frac{3[p^2+p_{}^{}{}_{}{}^{2}(E_p^AE_p^{}^A)^2]^2}{4p^2p_{}^{}{}_{}{}^{2}}}1]\mathrm{ln}\left({\displaystyle \frac{E_p^A+(p+p^{})E_p^{}^A}{E_p^A+|pp^{}|E_p^{}^A}}\right)\}.`$ (226)
They coincide with the limit $`M_B\mathrm{}`$ of the corresponding expressions in Eq. (177).
The consistency of the one-body limit is hence fully established at the level of the effective Schrödinger equation after the separation of angular and spin variables.
## Appendix D Approximate diagonalization of the effective potential matrices
To order $`\alpha ^4`$, i.e., for the lowest-order relativistic corrections, we can approximate
$`{\displaystyle \frac{p}{E_p^A+M_A}}={\displaystyle \frac{p}{2M_A}},{\displaystyle \frac{p}{E_p^B+M_B}}={\displaystyle \frac{p}{2M_B}},`$ (227)
and analogously for $`p^{}`$ instead of $`p`$, in the potential terms of the effective Schrödinger equations (172) and (175). Furthermore, terms containing
$`{\displaystyle \frac{p}{E_p^A+mM_A}}{\displaystyle \frac{p}{E_p^B+M_B}}{\displaystyle \frac{p^{}}{E_p^{}^A+M_A}}{\displaystyle \frac{p^{}}{E_p^{}^B+M_B}}`$ (228)
can be neglected. The matrix (176) in the L-coupled sector can then be approximated by the following diagonal matrix
$`{\displaystyle \frac{1}{2J+1}}\left(\begin{array}{cc}{}_{}{}^{L}V_{J1,J1}^{J}(p,p^{})& {}_{}{}^{L}V_{J1,J+1}^{J}(p,p^{})\\ {}_{}{}^{L}V_{J+1,J1}^{J}(p,p^{})& {}_{}{}^{L}V_{J+1,J+1}^{J}(p,p^{})\end{array}\right)`$ (231)
$`=\left(\begin{array}{cc}{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}& 0\\ 0& {\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\end{array}\right)\left({\displaystyle \frac{pp^{}}{4M_A^2}}+{\displaystyle \frac{pp^{}}{4M_B^2}}\right)\left(\begin{array}{cc}{\displaystyle \frac{a_J(p,p^{})}{2J+1}}& 0\\ 0& {\displaystyle \frac{a_J(p,p^{})}{2J+1}}\end{array}\right).`$ (236)
In particular, to order $`\alpha ^4`$ there is *no* L-coupling. We hence expect, at least for moderate values of $`\alpha `$, the mixing for the solutions in the L-coupled sector (with $`J1`$) to be rather small.
The matrix elements (173) in the S-coupled sector can be approximated by
$`{}_{}{}^{S}V_{00}^{J}(p,p^{})`$ $`=a_J(p,p^{})\left[{\displaystyle \frac{pp^{}}{4M_A^2}}+{\displaystyle \frac{pp^{}}{4M_B^2}}\right]\left[J{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+(J+1){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{S}V_{11}^{J}(p,p^{})`$ $`=a_J(p,p^{})\left[{\displaystyle \frac{pp^{}}{4M_A^2}}+{\displaystyle \frac{pp^{}}{4M_B^2}}\right]\left[(J+1){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}+J{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right],`$
$`{}_{}{}^{S}V_{01}^{J}(p,p^{})`$ $`={}_{}{}^{S}V_{10}^{J}(p,p^{})=\left[{\displaystyle \frac{pp^{}}{4M_A^2}}{\displaystyle \frac{pp^{}}{4M_B^2}}\right]\sqrt{J(J+1)}\left[{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}\right].`$ (237)
This approximate matrix, somewhat fortunately, can be diagonalized through a $`p`$\- and $`p^{}`$-independent transformation. After some algebraic labor, one obtains the following eigenvectors:
$`v_{J,x}^{(+)}`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\left(1{\displaystyle \frac{1}{\sqrt{1+4J(J+1)x^2}}}\right)^{1/2}\\ \left(1+{\displaystyle \frac{1}{\sqrt{1+4J(J+1)x^2}}}\right)^{1/2}\end{array}\right),`$ (240)
$`v_{J,x}^{()}`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\left(1+{\displaystyle \frac{1}{\sqrt{1+4J(J+1)x^2}}}\right)^{1/2}\\ \left(1{\displaystyle \frac{1}{\sqrt{1+4J(J+1)x^2}}}\right)^{1/2}\end{array}\right)`$ (243)
($`J1`$), where we have expressed the mass dependence through the parameter
$`x={\displaystyle \frac{M_B^2M_A^2}{M_B^2+M_A^2}}.`$ (244)
The corresponding eigenvalues are
$`{\displaystyle \frac{1}{2J+1}}{}_{}{}^{S}V_{J,x}^{(+)}(p,p^{})`$ $`={\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{pp^{}}{4M_A^2}}{\displaystyle \frac{1}{1+x}}[(1+{\displaystyle \frac{\sqrt{1+4J(J+1)x^2}}{2J+1}}){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}`$
$`+(1{\displaystyle \frac{\sqrt{1+4J(J+1)x^2}}{2J+1}}){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}],`$
$`{\displaystyle \frac{1}{2J+1}}{}_{}{}^{S}V_{J,x}^{()}(p,p^{})`$ $`={\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{pp^{}}{4M_A^2}}{\displaystyle \frac{1}{1+x}}[(1{\displaystyle \frac{\sqrt{1+4J(J+1)x^2}}{2J+1}}){\displaystyle \frac{a_{J1}(p,p^{})}{2J1}}`$
$`+(1+{\displaystyle \frac{\sqrt{1+4J(J+1)x^2}}{2J+1}}){\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}}].`$ (245)
These formulas are valid for $`x0`$ or $`M_BM_A`$. Comparing these results with the approximate diagonal matrix elements for L-coupling,
$`{\displaystyle \frac{1}{2J1}}{}_{}{}^{L}V_{JJ}^{J1}(p,p^{})`$ $`={\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{pp^{}}{4M_A^2}}{\displaystyle \frac{2}{1+x}}{\displaystyle \frac{a_{J1}(p,p^{})}{2J1}},`$
$`{\displaystyle \frac{1}{2J+3}}{}_{}{}^{L}V_{JJ}^{J+1}(p,p^{})`$ $`={\displaystyle \frac{a_J(p,p^{})}{2J+1}}{\displaystyle \frac{pp^{}}{4M_A^2}}{\displaystyle \frac{2}{1+x}}{\displaystyle \frac{a_{J+1}(p,p^{})}{2J+3}},`$ (246)
one notes the coincidence in the one-body limit $`x1`$. For $`x`$ close to one, Eq. (245) can be expanded in $`x`$ around one. Consequently, the energy levels that are degenerate at $`x=1`$ split for $`M_AM_B`$ by terms of the order of (at least) $`(M_A/M_B)^2`$, which points to the fact that there is *no hyperfine splitting* of the levels in the strict sense, i.e., to orders $`\alpha ^4`$ and $`M_A/M_B`$.
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# 1 Introduction
## 1 Introduction
One of the most important open questions of neutrino physics is whether neutrinos are Dirac or Majorana particles. From a theoretical perspective, large Majorana mass terms appear quite naturally for the right-handed neutrinos, since they are complete gauge singlets. This leads directly to the best investigated (and therefore most popular) explanations for the huge ratio of observed mass scales in the see-saw mechanism . In its simplest form, neutrino masses get suppressed by a factor $`v_{\mathrm{EW}}/M_{}`$ with $`v_{\mathrm{EW}}`$ denoting the vacuum expectation value (VEV) of the Higgs boson and $`M_{}`$ being the scale at which $`BL`$ symmetry (baryon minus lepton number) is assumed to be broken. However, it is important to keep in mind that the suppression factor $`v/M_{}`$ (with $`M_{}`$ now being the GUT scale or a related scale) arises rather generally whenever neutrino masses arise from integrating out heavy degrees of freedom with mass $`M_{}`$. This statement holds independently of the nature of neutrino masses, in particular both for Dirac and Majorana masses. Indeed, there exist a couple of appealing models where small Dirac masses are explained in this way by using extra, heavy degrees of freedom, or by relating the Yukawa coupling $`Y_\nu `$ to the ratio of gravitino mass (or other soft masses) and GUT (or compactification) scale . Another possibility is using localization in extra dimensions, and explaining the suppression by a small overlap of the corresponding zero-mode profiles along extra dimensions (see, e.g., ). Further support for Dirac neutrinos was found in certain orbifold compactifications of the heterotic string, where it is difficult to obtain the standard see-saw . For recent overviews and further references of various possibilities of explaining small Dirac masses see .
Cosmological arguments do not give a preference for Dirac or Majorana masses either. For instance, even if one requires the explanation of the observed baryon asymmetry to be related to neutrino properties, one finds that successful baryogenesis can work both for Majorana and Dirac neutrinos. Dirac neutrinos evade also constraints from primordial nucleosynthesis, since the right-handed degrees of freedom decouple with a low temperature so that their energy density is relatively suppressed . The question whether neutrinos are Dirac or Majorana particles is therefore one of the main motivations for improved neutrino-less double beta decay experiments. Both options should therefore be studied seriously until this question is clarified by experiments.
We investigate in this Letter RG effects under the assumption that neutrinos are Dirac particles, and that the small Yukawa couplings are explained by some mechanism operating at a high, e.g. GUT or compactification, scale, denoted by $`M_{\mathrm{GUT}}`$ in the following. The radiative corrections to the leptonic CP violation has been studied in . We extend this analysis to include all leptonic mixing parameters, and derive analytic formulae describing the renormalization group (RG) running of the leptonic mixing parameters. The radiative corrections are compared to analogous corrections in the quark sector, and we will show that generically RG running in the lepton sector is stronger than in the quark sector since the coefficients of the renormalization group equations (RGEs) are enhanced due to the fact that mass hierarchy is milder and the mixing angles are larger in the lepton sector. We compare the size of the radiative corrections to the accuracy of present and future neutrino experiments, and find that they are particularly relevant if neutrino masses are degenerate and $`\mathrm{tan}\beta `$ is large in the MSSM.
## 2 Analytic formulae
Using the standard parameterization (see, e.g., ) for leptonic (and quark) mixing the RGEs for the leptonic mixing angles read
$`\dot{\theta }_{12}`$ $`=`$ $`{\displaystyle \frac{Cy_\tau ^2}{32\pi ^2}}{\displaystyle \frac{m_1^2+m_2^2}{m_2^2m_1^2}}\mathrm{sin}(2\theta _{12})\mathrm{sin}^2\theta _{23}+𝒪(\theta _{13}),`$ (1)
$`\dot{\theta }_{13}`$ $`=`$ $`{\displaystyle \frac{Cy_\tau ^2}{32\pi ^2}}{\displaystyle \frac{1}{\left(m_3^2m_1^2\right)\left(m_3^2m_2^2\right)}}\{(m_2^2m_1^2)m_3^2\mathrm{cos}\delta \mathrm{cos}\theta _{13}\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23})`$ (2)
$`+[m_3^4(m_2^2m_1^2)m_3^2\mathrm{cos}(2\theta _{12})m_1^2m_2^2]\mathrm{cos}^2\theta _{23}\mathrm{sin}(2\theta _{13})\},`$
$`\dot{\theta }_{23}`$ $`=`$ $`{\displaystyle \frac{Cy_\tau ^2}{32\pi ^2}}{\displaystyle \frac{\left[m_3^4m_1^2m_2^2+(m_2^2m_1^2)m_3^2\mathrm{cos}(2\theta _{12})\right]}{(m_3^2m_1^2)(m_3^2m_2^2)}}\mathrm{sin}(2\theta _{23})+𝒪(\theta _{13}),`$ (3)
where the dot indicates the logarithmic derivative w.r.t. the renormalization scale $`\mu `$, e.g. $`\dot{\theta }_{12}=\mathrm{d}\theta _{12}/\mathrm{d}t=\mu \mathrm{d}\theta _{12}/\mathrm{d}\mu `$, and
$$C=\{\begin{array}{cc}1,\hfill & \text{(MSSM)},\hfill \\ 3/2,\hfill & \text{(SM)}.\hfill \end{array}$$
(4)
Here, we have neglected the tiny electron and muon Yukawa couplings, as well as the neutrino Yukawa couplings, against the $`\tau `$ Yukawa coupling $`y_\tau `$. Furthermore, in $`\dot{\theta }_{12}`$ and $`\dot{\theta }_{23}`$ we only kept the leading order term of an expansion in the reactor mixing angle $`\theta _{13}`$. We have checked numerically that this is a good approximation for realistic values of $`\theta _{13}`$.
It is instructive to compare Eqs. (1)–(3) to the corresponding expressions for Majorana neutrinos. Technically one obtains Eqs. (1)–(3) by discarding all terms which depend on the Majorana phases in Eqs. (8)–(10) of Ref. . One could thus say that the running of the Dirac mixing parameters resembles the running of Majorana mixing parameters averaged over the Majorana phases $`\phi _1`$ and $`\phi _2`$<sup>1</sup><sup>1</sup>1Stated differently, the running in the Dirac case is roughly half of the maximal running in the Majorana case. The factor 1/2 can be understood from the structure of the RGE: in the Dirac case, the mass matrix gets rotated by only one term (cf. Eq. (A.3d)),
$$\mathrm{\Delta }m_\nu =Cm_\nu \left(Y_e^{}Y_e\right)+\text{flavor-trivial terms},$$
while in the Majorana case there are two terms
$$\mathrm{\Delta }m_\nu =C\left[m_\nu \left(Y_e^{}Y_e\right)+\left(Y_e^{}Y_e\right)^Tm_\nu \right]+\text{flavor-trivial terms},$$
with $`C=3/2`$ in SM and two-Higgs models , and $`C=1`$ in the MSSM .. This means in particular that strong damping effects for the evolution of $`\theta _{12}`$, as observed for Majorana neutrinos, cannot occur in the Dirac case.
From Eqs. (1)–(3), we can immediately recognize several features of the RG evolution. First, for a strong mass hierarchy, the running of the angles is negligible. For $`m_1=0`$, the angles always run less than $`1^{}`$ (except for $`\theta _{23}`$ which runs more if $`\mathrm{tan}\beta 40`$). Second, for $`m_10.02\mathrm{eV}`$, $`\theta _{12}`$ has the strongest RG evolution. Third, as is obvious from Eqs. (1) and (3), in the MSSM $`\theta _{12}`$ always increases when running downwards while $`\theta _{23}`$ increases for a normal and decreases for an inverted mass hierarchy. This means that these two angles are radiatively enhanced (for normal mass ordering) which may, at least partially, be the reason for their large size. Whether $`\theta _{13}`$ increases or decreases depends on $`\delta `$.
The evolution of the Dirac phase $`\delta `$ is described by<sup>2</sup><sup>2</sup>2The evolution of the weak basis CP invariant has already been studied in .
$$\dot{\delta }=\dot{\delta }^{(1)}\theta _{13}^1+\dot{\delta }^{(0)}+\dot{\delta }^{(1)}+𝒪\left(\theta _{13}^2\right),$$
(5)
with the first two coefficients $`\dot{\delta }^{(k)}`$ given by
$`\dot{\delta }^{(1)}`$ $`=`$ $`{\displaystyle \frac{Cy_\tau ^2}{32\pi ^2}}{\displaystyle \frac{(m_2^2m_1^2)m_3^2}{\left(m_3^2m_1^2\right)\left(m_3^2m_2^2\right)}}\mathrm{sin}(\delta )\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23}),`$ (6a)
$`\dot{\delta }^{(0)}`$ $`=`$ $`0.`$ (6b)
Moreover, the term linear in $`\theta _{13}`$ contains
$`\dot{\delta }^{(1)}`$ $`=`$ $`{\displaystyle \frac{Cy_\tau ^2}{16\pi ^2}}{\displaystyle \frac{m_2^2\left(m_3^2m_1^2\right)^2}{\left(m_2^2m_1^2\right)\left(m_3^2m_1^2\right)\left(m_3^2m_2^2\right)}}\mathrm{cot}(\theta _{12})\mathrm{sin}(2\theta _{23})\mathrm{sin}\delta +\mathrm{},`$ (6c)
which becomes relevant if $`\theta _{13}`$ is not too small.
As usual, $`\delta `$ is undefined for $`\theta _{13}=0`$. Clearly, if $`\delta `$ vanishes for some scale, the (lepton sector of the) theory is CP invariant at this scale and thus remains CP invariant for all scales. Hence, the statement $`\delta =0`$ cannot depend on the renormalization scale, which can also be seen in our formulae. Likewise, if $`\theta _{13}`$ is zero at some given scale, the theory must again be CP invariant for all scales<sup>3</sup><sup>3</sup>3This is in contrast to the case of Majorana neutrinos, where the memory to CP violation can be stored in the Majorana phases, and $`\theta _{13}`$ can cross zero even in the presence of leptonic CP violation .. From this we conclude that $`\theta _{13}`$ can never cross zero if we have at some scale $`\theta _{13}0`$ and $`\mathrm{sin}\delta 0`$. If $`\theta _{13}`$ approaches zero, then we can see from (5) that $`\delta `$ runs quickly to a value such that the coefficient in (2) changes its sign and $`\theta _{13}`$ increases again. Thus, the only case where $`\theta _{13}`$ can cross zero is the CP conserving case. This is interesting for future precision measurements of $`\theta _{13}`$, since the assumption of leptonic CP violation at any scale leads to the conclusion that the weak-scale value of $`\theta _{13}`$ does not vanish. We illustrate the corresponding large effects in the evolution of $`\delta `$ in Fig. 1. For all our plots, we use the software packages REAP and MPT associated with .
We can understand this feature also differently. In the above approximation, we can write $`U_{e3}=\theta _{13}e^{\mathrm{i}\delta }`$ and by inserting the RGEs for $`\theta _{13}`$ and $`\delta `$, we find in the limit $`\theta _{13}0`$
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}U_{e3}`$ $`=`$ $`\dot{\theta }_{13}e^{\mathrm{i}\delta }\mathrm{i}\theta _{13}e^{\mathrm{i}\delta }\dot{\delta }`$ (7)
$``$ $`{\displaystyle \frac{Cy_\tau ^2}{32\pi ^2}}{\displaystyle \frac{(m_2^2m_1^2)m_3^2}{\left(m_3^2m_1^2\right)\left(m_3^2m_2^2\right)}}\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23}).`$
For $`\theta _{13}0`$ we find thus that the RG change of $`U_{e3}`$ is along the real axis and $`U_{e3}`$ can therefore only become zero if it is real. The imaginary part of Eq. (7) allows to determine a minimal value of $`\theta _{13}`$ as $`(\theta _{13})_{\mathrm{min}}\theta _{13}(\mu )\mathrm{sin}\delta (\mu )`$ where any scale $`\mu `$ can be used.
Furthermore, let $`t_0`$ denote the turning point of $`\theta _{13}`$ characterized by $`\delta =\pi /2`$. The ‘asymptotic’ behavior $`\delta (tt_0)\delta (t_0t)=\pi \delta (t_0t)`$ is a consequence of the fact that $`\dot{\delta }`$ is an odd function in $`\theta _{13}`$<sup>4</sup><sup>4</sup>4We have introduced $`\pi `$ in order to keep $`\theta _{13}`$ positive as we use the convention that $`\theta _{13}`$ is always positive, and a possible sign flip is absorbed in a jump of $`\delta `$.. This allows to understand why $`\delta `$ approaches $`\pi \delta (m_Z)`$ for large $`\mu `$ in Fig. 1.
The evolution of the mass eigenvalues is given by
$`16\pi ^2\dot{m}_1`$ $`=`$ $`\{Cy_\tau ^2[\mathrm{cos}^2\theta _{12}\mathrm{cos}^2\theta _{23}\mathrm{sin}^2\theta _{13}+\mathrm{sin}^2\theta _{12}\mathrm{sin}^2\theta _{23}`$ (8a)
$`{\displaystyle \frac{1}{2}}\mathrm{cos}\delta \mathrm{sin}\theta _{13}\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23})]+\alpha _\nu \}m_1,`$
$`16\pi ^2\dot{m}_2`$ $`=`$ $`\{Cy_\tau ^2[\mathrm{sin}^2\theta _{12}\mathrm{cos}^2\theta _{23}\mathrm{sin}^2\theta _{13}+\mathrm{cos}^2\theta _{12}\mathrm{sin}^2\theta _{23}`$ (8b)
$`+{\displaystyle \frac{1}{2}}\mathrm{cos}\delta \mathrm{sin}\theta _{13}\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23})]+\alpha _\nu \}m_2,`$
$`16\pi ^2\dot{m}_3`$ $`=`$ $`\left\{Cy_\tau ^2\mathrm{cos}^2\theta _{13}\mathrm{cos}^2\theta _{23}+\alpha _\nu \right\}m_3.`$ (8c)
$`\alpha _\nu `$ represents the flavor-independent part of the RGE, and is given in (A.5l). Clearly, the dominant RG effect of the masses is a common rescaling governed by $`\alpha _\nu `$. In addition, for large $`\mathrm{tan}\beta `$ in the MSSM, there are corrections specific to the individual $`m_i`$. In leading order in $`\theta _{13}`$, these effects tend to decrease $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ for a normal hierarchy and to increase $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ for an inverted hierarchy when running down.
Fig. 2 shows an extreme example of the evolution of the mass eigenvalues and the corresponding $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$ and $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$. For large $`\mathrm{tan}\beta `$, there are substantial deviations from the flavor-independent scaling of the mass eigenvalues. The latter can be approximately inferred from the curve of $`m_1`$ in Fig. 2 (a).
## 3 Radiative corrections and future precision experiments
An important question is if future experiments will reach a precision which allows to draw interesting conclusions from quantum corrections. There exist, for example, models where $`\theta _{13}`$ vanishes at the GUT scale, but RG corrections still lead to a finite value of $`\theta _{13}`$ at low energies. A certain finite value of $`\theta _{13}`$ is therefore guaranteed unless the initial values at the GUT scale and the independent RG effects are adjusted to cancel each other. From the discussion of the previous section, we would also know that the CP phase $`\delta `$ would vanish for all scales for Dirac neutrinos, while it could become finite for Majorana neutrinos. A finite value of $`\delta `$ and $`\theta _{13}`$ would thus exclude either Dirac neutrinos or $`\theta _{13}(M_{GUT})=0`$. Similar arguments can be made for other quantities of the leptonic mixing matrix. $`\theta _{23}`$, for example, is within current experimental errors compatible with $`45^{}`$. However, as in the case of $`\theta _{13}`$, certain deviations are expected from RG effects even if $`45^{}`$ is exactly predicted at the GUT scale. Future precision measurements of neutrino oscillations may therefore allow interesting tests of flavor models and related renormalization group effects.
Atmospheric neutrino data and results from the K2K long-baseline accelerator experiment currently determine $`\mathrm{\Delta }m_{31}^2=(2.2_{0.4}^{+0.6})\times 10^3\mathrm{eV}^2`$ and $`\theta _{23}45^{}`$ , whereas solar neutrino data , combined with the results from the KamLAND reactor experiment lead to $`\mathrm{\Delta }m_{21}^2=(8.2_{0.3}^{+0.3})\times 10^5\mathrm{eV}^2`$ and $`\mathrm{tan}^2\theta _{12}=0.39_{0.04}^{+0.05}`$ . These results are to a good approximation still described by two independent two flavor oscillations. The key parameter for genuine three flavor effects is the mixing angle $`\theta _{13}`$ which is so far only known to be small from the CHOOZ and Palo Verde experiments. The current bound for $`\theta _{13}`$ depends on the value of the atmospheric mass squared difference and it gets weaker for $`\mathrm{\Delta }m_{31}^22\times 10^3\mathrm{eV}^2`$. However, in that region an additional constraint on $`\theta _{13}`$ from global solar neutrino data becomes important . At the current best fit value of $`\mathrm{\Delta }m_{31}^2=2.2\times 10^3\mathrm{eV}^2`$ the $`3\sigma `$ bound is $`\mathrm{sin}^2\theta _{13}0.041`$ .
Genuine three flavor oscillation effects occur only for a finite value of $`\theta _{13}`$ and establishing a finite value of $`\theta _{13}`$ is therefore one of the next milestones in neutrino physics. Leptonic CP violation is another three flavor effect which can only be tested if $`\theta _{13}`$ is finite. There exists therefore a very strong motivation to establish a finite value of $`\theta _{13}`$ and then leptonic CP violation . Different experimental projects are therefore under construction or are being planned in order to achieve these goals. It is useful to divide the future into what can be achieved with specific current or next generation projects and what may be achieved with long term projects. The MINOS project, which started already data taking, and the CNGS projects ICARUS and OPERA , which are completing construction can be considered as ‘‘current projects’’. Beyond that exist other, more ambitious ‘‘next generation’’ projects like the superbeam experiments J-PARC to SuperKamiokande (T2K) and the NuMI off-axis experiment NO$`\nu `$. In addition there are ‘‘next generation’’ plans for new reactor neutrino experiments with a near and far detector. A first interesting question concerns improvements of $`\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}^2\theta _{23}`$. In Tab. 1 we show the relative precision which can be obtained in the future in comparison to the current precision, as obtained from a global fit to SuperKamiokande (SK) atmospheric and K2K long-baseline data . We observe from these numbers, that the accuracy on $`\mathrm{\Delta }m_{31}^2`$ can be improved by one order of magnitude, whereas the accuracy on $`\mathrm{sin}^2\theta _{23}`$ will be improved only by a factor two.
Tab. 1 depends on the value of $`\mathrm{\Delta }m_{31}^2`$ and the sensitivity suffers for all experiments for low values of $`\mathrm{\Delta }m_{31}^2`$. T2K will provide a relatively precise determination of $`\mathrm{\Delta }m_{31}^2`$ for $`\mathrm{\Delta }m_{31}^22\times 10^3\mathrm{eV}^2`$. Although NO$`\nu `$A can put a comparable lower bound on $`\mathrm{\Delta }m_{31}^2`$, the upper bound is significantly weaker, and similar to the bound from MINOS . The reason for this is a strong correlation between $`\mathrm{\Delta }m_{31}^2`$ and $`\theta _{23}`$, which disappears only for $`\mathrm{\Delta }m_{31}^23\times 10^3\mathrm{eV}^2`$. From Tab. 1 one can also see that only T2K is able to improve the current bound on $`\mathrm{sin}^2\theta _{23}`$. The main reason for the rather poor performance on $`\mathrm{sin}^2\theta _{23}`$ is the fact that these experiments are mostly sensitive to $`\mathrm{sin}^22\theta _{23}`$. This implies that for $`\theta _{23}\pi /4`$ it is very hard to achieve a good accuracy on $`\mathrm{sin}^2\theta _{23}`$, although $`\mathrm{sin}^22\theta _{23}`$ can be measured with relatively high precision .
An important task of the next generation long baseline and reactor experiments of the coming years will be to establish the three flavored-ness of the oscillations. The sensitivity to a finite value of the key parameter $`\theta _{13}`$ is only modest for MINOS, OPERA and ICARUS. Double Chooz, T2K and NO$`\nu `$A can do much better. The $`\mathrm{sin}^22\theta _{13}`$-limits of the beam experiments are, however, strongly affected by parameter correlations and degeneracies, whereas new reactor experiments provide a ‘‘clean’’ measurement of $`\mathrm{sin}^22\theta _{13}`$ . Altogether these experiments will provide an improvement by about a factor ten for $`\mathrm{sin}^22\theta _{13}`$ over the existing limit. In addition, the KamLAND (and solar neutrino) data will also further increase the accuracy for $`\mathrm{\Delta }m_{21}^2`$ and $`\theta _{12}`$. An accuracy of $`5\%`$ for $`\mathrm{\Delta }m_{12}^2`$ and $`20\%`$ for $`\mathrm{sin}^2\theta _{12}`$ is expected. Further improvements are possible, e.g. by loading the SuperKamiokande detector with Gadolinium, which might lead to an error of $`3\%`$ for $`\mathrm{\Delta }m_{21}^2`$ and $`15\%`$ for $`\mathrm{sin}^2\theta _{12}`$, both at 99%CL .
Beyond the next generation accelerator and reactor based oscillation experiments exist much more ambitious projects like the JHF-HyperKamiokande project, beta beams and neutrino factories<sup>5</sup><sup>5</sup>5See for a comparison and for references.. Such experiments clearly require further R&D before they can be built. However, assuming current knowledge, such facilities appear to be possible and they will lead to a precision at the level of percent or even below. With a neutrino factory, for example, a sensitivity to a finite value of $`\mathrm{sin}^22\theta _{13}`$ might be achievable down to $`10^5`$.
It is interesting to compare these perspectives with RG effects. To illustrate the RG effects, we start with initial values for the mixing parameters at the GUT scale, $`M_{\mathrm{GUT}}=3\times 10^{16}\mathrm{GeV}`$, assuming that these values find an explanation in a more fundamental theory.<sup>6</sup><sup>6</sup>6One could, for instance, enjoy the possibility of fixing the initial values by continuous (e.g. ) or discrete (e.g. ) symmetries. In this case, RG effects add to the corrections arising from the breakdown of those symmetries. These initial values are then compared with the corresponding mixing parameters at $`m_Z`$. In all our illustrations, we assume $`m_{\mathrm{SUSY}}=1\mathrm{TeV}`$, and a normal mass hierarchy. The simple expressions (Eqs. (1)--(3) and (5)) allow a quick estimate of the RG effects. A more precise evaluation requires a numerical analysis for which we use the Mathematica package REAP , which is publicly available<sup>7</sup><sup>7</sup>7See http://www.ph.tum.de/~rge/.
The mixing angles $`\theta _{12}`$ and $`\theta _{23}`$ turn out to be rather unstable for a degenerate spectrum (cf. Fig. 3). As a consequence, a Dirac version of quark-lepton complementarity can only work for certain mass eigenvalues and ratios of the Higgs VEVs and $`\mathrm{tan}\beta `$ (for the discussion of the RG effects in the see-saw Majorana case see ). This means stability of $`\theta _{12}`$ is only given in models with hierarchical masses and/or small $`\mathrm{tan}\beta `$. Note also that for an inverted hierarchy $`\theta _{12}`$ is unstable. This means that concerning $`\theta _{12}`$ RG effects are in general an issue. RG corrections to the special value $`\theta _{23}=45^{}`$ can be comparable to the precision of upcoming experiments. Again, this happens for a quite degenerate spectrum and/or large $`\mathrm{tan}\beta `$.
The running of $`\theta _{13}`$ depends crucially on its initial value. We illustrate this by plotting the radiative correction to $`\mathrm{sin}^22\theta _{13}`$ in Fig. 4. Most important is the second term in Eq. (2) which is dominant for not too small $`\theta _{13}`$. As a consequence we find that, for $`\theta _{13}=0`$ at the high scale, running in general does not generate a measurable value at the low scale. Only for the most optimistic sensitivities, a quite degenerate spectrum and large $`\mathrm{tan}\beta `$ this conclusion can be avoided. On the other hand, if $`\theta _{13}`$ is not tiny, RG effects can be comparable to the precision of upcoming experiments (except for small $`\mathrm{tan}\beta `$).
Finally, let us discuss corrections to $`\delta `$. From the previous discussion in Sec. 2 it is clear that small $`\theta _{13}`$ corresponds to an unstable configuration with large RG effects, even for a hierarchical spectrum (cf. Fig. 5 (b)). In particular, RG effects are generically comparable with the precision of future experiments such as the combination T2K+NO$`\nu `$A+Reactor-II, T2HK and NuFact-II (see and references therein).
Let us finally mention that RG effects for Dirac neutrinos will always result in a rescaling of the mass eigenvalues. Beyond that, in the framework of the SM, mixing parameters are quite stable. The only exceptions are $`\theta _{12}`$ for very degenerate masses, and $`\delta `$ for tiny $`\theta _{13}`$. On the other hand, in the MSSM, RG effects are increased by $`\mathrm{tan}^2\beta `$, i.e. by up to three orders of magnitude.
## 4 Summary
Assuming Dirac neutrinos, we have derived renormalization group equations for leptonic mixing parameters. The results share several features with the corresponding equations for Majorana neutrinos. However, Dirac running is more predictive, as the Majorana phases are unphysical in this case. This makes it possible to specify the amount of renormalization group evolution unambiguously as soon as $`m_1`$ and $`\delta `$ (and $`\mathrm{tan}\beta `$) are known. The renormalization group evolution alone does not yield an explanation of the largeness of the leptonic mixing angles (for an analogous and very clear discussion for Majorana neutrinos see ). Yet it may account for radiative enhancement of $`\theta _{12}`$, and possibly also of $`\theta _{23}`$, since both can increase significantly in the MSSM when running down.
Most importantly, we find that in phenomenological studies renormalization group effects for leptonic mixing angles can in general not be neglected. This can be understood from the fact that $`\dot{\theta }_{ij}=f(m_i,\theta _{ij},\delta )/(m_i^2m_j^2)`$ which becomes singular if $`m_im_j`$ and vanishes if the mixing angles are zero. We have thus traced back the relative enhancement of the quantum corrections of leptonic mixing parameters as compared to quark mixings to two reasons. First, the mass hierarchy which suppresses the renormalization group running, can be much weaker. Second, the mixing angles are larger so that the parameters are further apart from the renormalization group stable situation where all mixings are zero.
As there is no suppression of the running by phases, the renormalization group corrections should in general be taken into account even for a strong hierarchy to accommodate the precision of future experiments. Renormalization group corrections are especially relevant if the mass spectrum is non-hierarchical, and $`\mathrm{tan}\beta `$ is large in the MSSM. Hence, similarly to the case of Majorana neutrinos , also in the Dirac case the non-observation of deviations of the angles from special points, e.g. of $`\theta _{12}`$ from $`\pi /4\vartheta _{12}`$ (with $`\vartheta _{12}`$ being the Cabibbo angle), of $`\theta _{13}`$ from 0 and $`\theta _{23}`$ from $`\pi /4`$, may restrict the parameters such as the absolute neutrino mass scale. The current experimental data already has the necessary precision to indicate disfavored parameter regions. It may also point to exactly realized symmetries and our formulae can hence be used to identify possible symmetries. Whenever a symmetry is exact and fixes some mixing parameters, those mixing parameters have to be stable under quantum corrections. For instance, it has recently been pointed out that for Majorana neutrinos and an inverted hierarchy the configuration $`m_3=\theta _{13}=0`$ is stable. From the analytic expressions it is obvious that this statement also applies to the Dirac case. Likewise, a quick inspection of the RGEs excludes most of the proposed symmetries from being exact. Our formulae are basis-independent, and thus allow to understand certain features of the underlying theory, such as symmetries, in a basis-independent way. We have discussed this only for the case of CP symmetry, but it is obvious how the analysis can be carried over to other symmetries. In this context, it would be interesting to see if infrared fixed points with large mixings, as discussed in , can be obtained for (non-standard) Dirac neutrinos as well . In this case, one may hope for a scenario where the large mixings are a consequence of running, and the mechanism of generation of neutrino masses is still related to the scale where gauge couplings meet.
We conclude that in the light of future precision experiments, flavor physics might enter into an era of ‘‘precision model building’’. It seems possible to determine the mixing parameters to a remarkable accuracy, precise enough such that flavor models and the corresponding renormalization group effects become to a certain degree distinguishable. For a specific parameter and a desired accuracy, our formulae allow to estimate the renormalization group effects, and to judge to which extent a numerical analysis is required.
#### Acknowledgements
We would like to thank S. Antusch, A. Dedes and J. Kersten for valuable discussions. One of us (M.R.) would like to thank the Aspen Center for Physics for support, and the CERN theory group for hospitality. This work was partially supported by the EU 6th Framework Program MRTN-CT-2004-503369 ‘‘Quest for Unification’’ and MRTN-CT-2004-005104 ‘‘ForcesUniverse’’. This work was also supported by the Deutsche Forschungsgemeinschaft in the ‘‘Sonderforschungsbereich 375 für Astro-Teilchenphysik’’ and under project number RO-2516/3-1.
## Appendix A Mixing parameters RGEs for Dirac masses
### A.1 Lagrangian
The Yukawa couplings are given by
$$_{\mathrm{Yuk}}=\left(Y_\nu \right)_{gf}\overline{N_R^g}\stackrel{~}{\varphi }^{}\mathrm{}_L^f+\left(Y_e\right)_{gf}\overline{e_R^g}\varphi ^{}\mathrm{}_L^f+\left(Y_u\right)_{gf}\overline{u_R^g}\stackrel{~}{\varphi }^{}Q_L^f+\left(Y_d\right)_{gf}\overline{d_R^g}\varphi ^{}Q_L^f$$
(A.1)
in the SM extended by right-handed neutrinos where $`\stackrel{~}{\varphi }=i\sigma _2\varphi ^{}`$. In the extended MSSM, the Yukawa couplings are analogously defined in the superpotential by
$$𝒲_{\mathrm{Yuk}}=\left(Y_\nu \right)_{gf}N_R^{Cg}\varphi ^{(2)}ϵ^T\mathrm{}_L^f+\left(Y_e\right)_{gf}e_R^{Cg}\varphi ^{(1)}ϵ\mathrm{}_L^f+\left(Y_u\right)_{gf}u_R^{Cg}\varphi ^{(2)}ϵ^TQ_L^f+\left(Y_d\right)_{gf}d_R^{Cg}\varphi ^{(1)}ϵQ_L^f.$$
(A.2)
The left-handed lepton and quark doublets are denoted by $`\mathrm{}_L`$ and $`Q_L`$, respectively. We assume that there is no Majorana mass term for the right-handed neutrinos.
### A.2 $`𝜷`$-functions
The relevant $`\beta `$-functions for the down-type quark, up-type quark, charged lepton and neutrino Yukawa coupling matrices $`Y_d`$, $`Y_u`$, $`Y_e`$ and $`Y_\nu `$ read at one-loop
$`(4\pi )^2\dot{Y}_d`$ $`=`$ $`Y_d\left\{C_d^dY_d^{}Y_d+C_d^uY_u^{}Y_u+\alpha _d\right\},`$ (A.3a)
$`(4\pi )^2\dot{Y}_u`$ $`=`$ $`Y_u\left\{C_u^dY_d^{}Y_d+C_u^uY_u^{}Y_u+\alpha _u\right\},`$ (A.3b)
$`(4\pi )^2\dot{Y}_e`$ $`=`$ $`Y_e\left\{C_e^eY_e^{}Y_e+C_e^\nu Y_\nu ^{}Y_\nu +\alpha _{\mathrm{}}\right\},`$ (A.3c)
$`(4\pi )^2\dot{Y}_\nu `$ $`=`$ $`Y_\nu \left\{C_\nu ^eY_e^{}Y_e+C_\nu ^\nu Y_\nu ^{}Y_\nu +\alpha _\nu \right\},`$ (A.3d)
where
$`C_d^d`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 3,\hfill & \text{(MSSM)}\hfill \end{array}`$ $`C_d^u`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 1,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.4e)
$`C_u^d`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 1,\hfill & \text{(MSSM)}\hfill \end{array}`$ $`C_u^u`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 3,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.4j)
$`C_e^e`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 3,\hfill & \text{(MSSM)}\hfill \end{array}`$ $`C_e^\nu `$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 1,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.4o)
$`C_\nu ^e`$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 1,\hfill & \text{(MSSM)}\hfill \end{array}`$ $`C_\nu ^\nu `$ $`=\{\begin{array}{cc}3/2,\hfill & \text{(SM)}\hfill \\ 3,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.4t)
and
$`\alpha _d`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{4}g_1^2\frac{9}{4}g_2^28g_3^2+T_{\mathrm{SM}},\hfill & \text{(SM)}\hfill \\ 3\mathrm{Tr}(Y_d^{}Y_d)+\mathrm{Tr}(Y_e^{}Y_e)\frac{7}{15}g_1^23g_2^2\frac{16}{3}g_3^2,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.5c)
$`\alpha _u`$ $`=`$ $`\{\begin{array}{cc}\frac{17}{20}g_1^2\frac{9}{4}g_2^28g_3^2+T_{\mathrm{SM}},\hfill & \text{(SM)}\hfill \\ \mathrm{Tr}(Y_\nu ^{}Y_\nu )+3\mathrm{Tr}(Y_u^{}Y_u)\frac{13}{15}g_1^23g_2^2\frac{16}{3}g_3^2,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.5f)
$`\alpha _{\mathrm{}}`$ $`=`$ $`\{\begin{array}{cc}\frac{9}{4}g_1^2\frac{9}{4}g_2^2+T_{\mathrm{SM}},\hfill & \text{(SM)}\hfill \\ 3\mathrm{Tr}(Y_e^{}Y_e)+\mathrm{Tr}(Y_\nu ^{}Y_\nu )\frac{9}{5}g_1^23g_2^2,\hfill & \text{(MSSM)}\hfill \end{array}`$ (A.5i)
$`\alpha _\nu `$ $`=`$ $`\{\begin{array}{cc}\frac{9}{20}g_1^2\frac{9}{4}g_2^2+T_{\mathrm{SM}},\hfill & \text{(SM)}\hfill \\ \mathrm{Tr}(Y_\nu ^{}Y_\nu )+3\mathrm{Tr}(Y_u^{}Y_u)\frac{3}{5}g_1^23g_2^2,\hfill & \text{(MSSM)}.\hfill \end{array}`$ (A.5l)
Here, we define $`T_{\mathrm{SM}}\mathrm{Tr}\left[Y_e^{}Y_e+Y_\nu ^{}Y_\nu +3Y_d^{}Y_d+3Y_u^{}Y_u\right]`$, and use GUT normalization for $`g_1`$.
### A.3 General derivation
In this subsection, we will perform a general analysis applicable for any Dirac masses, and treat the evolution of lepton and quark masses and mixings only as a special case.
We derive the running of mixing parameters for a RGE of the form
$$16\pi ^2\frac{\mathrm{d}}{\mathrm{d}t}H=F^{}H+HF+fH,$$
(A.6)
where $`f`$ is real and $`H`$ is Hermitean, so that we can diagonalize it in a ‘reference basis’,
$$U^{}HU=D.$$
(A.7)
In the application in the main part, $`F`$ corresponds either to $`C_d^uY_u^{}Y_u+C_d^dY_d^{}Y_d`$ (or $`C_\nu ^eY_e^{}Y_e+C_\nu ^\nu Y_\nu ^{}Y_\nu `$ for the lepton sector), and $`H`$ to $`Y_d^{}Y_d`$ (or $`Y_\nu ^{}Y_\nu `$). The ‘reference basis’ is the basis where $`Y_u^{}Y_u`$ (or $`Y_e^{}Y_e`$) is diagonal at $`t=t_0`$. $`U`$ denotes then to the CKM matrix $`U_{\mathrm{CKM}}`$ (or the MNS matrix $`U_{\mathrm{MNS}}`$). $`f`$ denotes the diagonal parts of the $`\beta `$-function, $`f=2\alpha _d`$ (or $`f=2\alpha _\nu `$).
Now we perform an analysis very similar to what is done in which is based on . We can differentiate the relation $`H=UDU^{}`$,
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}(UDU^{})`$ $`=`$ $`\dot{U}DU^{}+UD\dot{U}^{}+U\dot{D}U^{}`$ (A.8)
$`\stackrel{!}{=}`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left(F^{}UDU^{}+UDU^{}F+fUDU^{}\right).`$
Multiplying by $`U^{}`$ from the left and by $`U`$ from the right yields
$$U^{}\dot{U}D+D\dot{U}^{}U+\dot{D}=\frac{1}{16\pi ^2}\left(F^{}D+DF^{}+fD\right),$$
(A.9)
where $`F^{}=U^{}FU`$. For the quark case, $`F^{}=C_d^dD+C_d^uU^{}Y_u^{}Y_uU`$. We will see below that only the off-diagonal components are relevant for the RGEs of the mixing parameters.
The evolution of $`U`$ can be written as
$$\frac{\mathrm{d}}{\mathrm{d}t}U=UX,$$
(A.10)
where $`X`$ is anti-Hermitean. Inserting this relation yields
$$\dot{D}+XD+DX^{}=\frac{1}{16\pi ^2}\left(F^{}D+DF^{}+fD\right),$$
(A.11)
or, by using the anti-Hermitecity of $`X`$,
$$\dot{D}=\frac{1}{16\pi ^2}\left(fD+F^{}D+DF^{}\right)XD+DX.$$
(A.12)
Denoting the entries of $`D`$ by $`y_i^2`$, i.e. $`D=\mathrm{diag}(y_1^2,y_2^2,y_3^2)`$, we find
$$\frac{\mathrm{d}}{\mathrm{d}t}y_i^2=\frac{1}{16\pi ^2}\left[fy_i^2+(F_{ii}^{}+F_{ii}^{})y_i^2\right],$$
(A.13)
i.e. the terms proportional to $`X`$ have dropped out. This equation corresponds to a RGE for the running mass eigenvalues, defined by $`m_i(t)=|y_i(t)|v`$ with $`v`$ fixed, of the form
$$(4\pi )^2\dot{m}_i=(\mathrm{Re}F_{ii}^{}+\alpha )m_i.$$
(A.14)
By analyzing the off-diagonal parts we obtain
$$y_i^2X_{ij}X_{ij}y_j^2=\frac{1}{16\pi ^2}\left[(F^{})_{ij}y_j^2+y_i^2F_{ij}^{}\right].$$
(A.15)
This can be converted into equations for real and imaginary part of $`X`$, which, since $`F`$ is Hermitean, can be combined to
$$16\pi ^2X_{ij}=\frac{y_j^2+y_i^2}{y_j^2y_i^2}F_{ij}^{}.$$
(A.16)
The diagonal parts of $`X`$ remain undetermined. However, this is not a problem, since they only influence the RG evolution of the unphysical phases.<sup>8</sup><sup>8</sup>8Note that the Majorana phases are unphysical in the the Dirac case as well.
So far, we have derived the differential change of the CKM matrix due to the RG corrections for $`Y_d^{}Y_d`$ (cf. Eq. (A.10)). In the Majorana neutrino case, the analogous differential equation already describes the evolution of the MNS matrix since $`Y_e^{}Y_e`$ doesn’t get rotated by the RGE.<sup>9</sup><sup>9</sup>9This is only true at leading order in $`M^1`$ where $`M`$ denotes the scale of the effective neutrino mass operator (e.g. the see-saw scale) . For Dirac neutrinos, $`Y_e^{}Y_e`$ gets rotated only by terms proportional to the squares of Dirac Yukawa couplings, hence those rotations can safely be neglected. In the quark sector, the radiative rotation of $`Y_u^{}Y_u`$ represents an important effect, as we will argue in the following.
### A.4 Contribution from the change of $`𝒀_𝒖`$
Here, we specialize to the quark sector as the analogous effect is irrelevant for Dirac neutrinos.
The RGE for $`Y_u`$ contains non-diagonal terms so that continuous re-diagonalization is required. Since the mixing matrix $`U_{\mathrm{CKM}}`$ is defined as the matrix which diagonalizes $`Y_d^{}Y_d`$ in the basis in which $`Y_u`$ is diagonal, $`U_{\mathrm{CKM}}`$ receives an additional contribution from the running of $`Y_u`$,
$$\frac{\mathrm{d}}{\mathrm{d}t}U_{\mathrm{CKM}}=U_{\mathrm{CKM}}X+\text{term stemming from the change of}Y_u.$$
(A.17)
To evaluate this change, we can essentially repeat the steps of the previous subsection. Introducing a matrix $`\stackrel{~}{U}`$ which diagonalizes $`Y_u^{}Y_u`$ in the reference basis (implying $`\stackrel{~}{U}(t=t_0)=\mathrm{𝟙}`$), i.e.
$$\stackrel{~}{U}^{}Y_u^{}Y_u\stackrel{~}{U}=\mathrm{diag}(\stackrel{~}{y}_1^2,\stackrel{~}{y}_2^2,\stackrel{~}{y}_3^2),$$
(A.18)
we arrive at
$$\frac{\mathrm{d}}{\mathrm{d}t}\stackrel{~}{U}=\stackrel{~}{U}\stackrel{~}{X},$$
(A.19)
where the off-diagonal entries of $`\stackrel{~}{X}`$ are given by
$$16\pi ^2\stackrel{~}{X}_{ij}=\frac{\stackrel{~}{y}_i^2+\stackrel{~}{y}_j^2}{\stackrel{~}{y}_j^2\stackrel{~}{y}_i^2}\stackrel{~}{F}_{ij}.$$
(A.20)
Completely analogous to A.3,
$$\stackrel{~}{F}^{}=\stackrel{~}{U}^{}\stackrel{~}{F}\stackrel{~}{U},$$
(A.21)
and at $`t=t_0`$
$$\stackrel{~}{F}^{}=C_u^dUDU^{}+C_u^uY_u^{}Y_u.$$
(A.22)
Again, only the off-diagonal terms influence the RGEs of the mixing angles.
### A.5 Mixing parameter RGEs in the quark sector
As $`U_{\mathrm{CKM}}=\stackrel{~}{U}^1U=\stackrel{~}{U}^{}U`$, the RGE for the CKM matrix reads
$$\frac{\mathrm{d}}{\mathrm{d}t}U_{\mathrm{CKM}}=\stackrel{~}{X}^{}U_{\mathrm{CKM}}+U_{\mathrm{CKM}}X.$$
(A.23)
To proceed, we label the mixing parameters by
$$\{\xi _k\}=\{\theta _{12},\theta _{13},\theta _{23},\delta ,\delta _e,\delta _\mu ,\delta _\tau ,\phi _1,\phi _2\},$$
(A.24)
and evaluate the derivative of $`U_{\mathrm{CKM}}`$,
$$\dot{U}_{\mathrm{CKM}}=\dot{U}_{\mathrm{CKM}}(\{\dot{\xi }_k\},\left\{\xi _k\right\}).$$
(A.25)
Observe that the resulting expression is linear in $`\dot{\xi }_k`$. By solving a system of linear equations of the form
$$\underset{k}{}A_{TX}^{(k)}\dot{\xi }_k+\mathrm{i}S_{TX}^{(k)}\dot{\xi }_k=R_X,$$
(A.26)
where
$$R_X=U_{\mathrm{MNS}}T+X^{}U_{\mathrm{MNS}},$$
(A.27)
we thus obtain a set of linear equations for the $`\dot{\xi }_k`$. RGEs for the matrix elements have been derived in refs. .
From these, we obtain the RGEs for the mixing angles in the quark sector. Neglecting all Yukawa coupling except for $`y_t`$ and $`y_b`$, they read
$`\dot{\vartheta }_{12}`$ $`=`$ $`{\displaystyle \frac{C_d^uy_{t}^{}{}_{}{}^{2}}{64\pi ^2}}\mathrm{cos}(\vartheta _{12})\{[(3\mathrm{cos}2\vartheta _{13})\mathrm{cos}2\vartheta _{23}2\mathrm{cos}^2\vartheta _{13}]\mathrm{sin}\vartheta _{12}`$ (A.28a)
$`+4\mathrm{cos}\delta _{\mathrm{CP}}\mathrm{cos}\vartheta _{12}\mathrm{sin}\vartheta _{13}\mathrm{sin}2\vartheta _{23}\},`$
$`\dot{\vartheta }_{13}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\vartheta _{13}}{64\pi ^2}}\left[2C_u^dy_b^2+C_d^uy_t^2\left(1+\mathrm{cos}2\vartheta _{23}\right)\right],`$ (A.28b)
$`\dot{\vartheta }_{23}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\vartheta _{23}}{64\pi ^2}}\left[C_u^dy_b^2\left(1+\mathrm{cos}2\vartheta _{13}\right)+2C_d^uy_t^2\right].`$ (A.28c)
It turns out that finite $`y_s`$ and $`y_c`$ corrections yield an important but sub-dominant effect for $`\dot{\vartheta }_{12}`$. The dominant term in the RGE of $`\delta _{\mathrm{CP}}`$ is
$$\dot{\delta }_{\mathrm{CP}}=\frac{C_d^uy_s^2y_t^2}{8\pi ^2\left(y_b^2y_s^2\right)}\mathrm{cos}\vartheta _{12}\mathrm{cos}\vartheta _{23}\mathrm{sin}\delta \mathrm{sin}\vartheta _{12}\mathrm{sin}\vartheta _{23}\times \vartheta _{13}^1.$$
(A.29)
### A.6 Mixing parameter RGEs in the (Dirac) neutrino sector
In order to derive analogous RGEs for the leptonic mixing parameters, observe that the RG change of $`Y_e^{}Y_e`$ is quadratic in neutrino Yukawa couplings, i.e. strongly suppressed. Thus, we can safely neglect the $`\stackrel{~}{X}`$ contribution,
$$\frac{\mathrm{d}}{\mathrm{d}t}U_{\mathrm{MNS}}=\stackrel{~}{X}^{}U_{\mathrm{MNS}}+U_{\mathrm{MNS}}XU_{\mathrm{MNS}}X,$$
(A.30)
where $`X`$ is now related to $`F^{}`$ by Eq. (A.16), and $`F^{}=C_\nu ^\nu D+C_\nu ^eU_{\mathrm{MNS}}^{}Y_e^{}Y_eU_{\mathrm{MNS}}`$ at $`t=t_0`$.
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# Evolution systems for paraxial wave equations of Schrödinger-type with non-smooth coefficients11footnote 1Work supported by FWF grants P16820-N04 and Y237-N13
## 1 Introduction
The paraxial equations in models of wave propagation are based on parabolic symbol approximations in theories of wave operators. They have been extensively applied in integrated optics, underwater acoustic tomography as well as reflection seismic imaging (cf. ). They have also entered the analysis of time-reversal mirror experiments with waves taking into account stochastic variations in the wave speed (cf. ). The paraxial equations can be derived from the reduced wave or Helmholtz equation and, since they split the wave fields according to a prescribed principal direction of propagation, are also called one-way wave equations. In particular, the leading-order parabolic symbol approximation, also referred to as the narrow-angle or beam-propagation approximation, leads to model equations of Schrödinger-type. The well-posedness of the one-way wave Cauchy problems has been discussed by . The methodologies developed to date, however, have assumed smoothness of the wave speed function (i.e., the coefficients in the wave operators).
In the analysis presented here, we depart from this smoothness assumption by allowing the coefficients to be of any Hölder regularity between zero and one, but typically non-Lipschitz. The existence of distributional solutions to second order strictly hyperbolic equations in general may fail below Log-Lipschitz regularity of the coefficients (cf. ). In case of Hölder regularity $`2`$ or higher a constructive approach for hyperbolic evolution equations has been developed in , thereby extending results on propagation of singularities. The particular equation considered here is derived from such a second order equation, but the existence of its solution is not restricted by the same conditions. Indeed, we exploit the framework of Sobolev space techniques, in particular, multiplication of distributions in scales of Sobolev spaces, to construct a strongly continuous evolution system. The novelty of the paper lies in the method of construction which not only provides a solution concept for the paraxial wave equation with low coefficient regularity, but also allows us to investigate how the coefficient regularity influences the solution.
The class of coefficients of Hölder regularity between zero and one arises in a variety of geophysical applications. Perhaps the most fundamental one concerns the study of thermo-chemical boundary layers and phase transitions in Earth’s lowermost mantle — the so-called $`D^{\prime \prime }`$ layer overlaying the core-mantle boundary (see Figure 1; cf. for recent images of the phase transformation in $`D^{\prime \prime }`$). Such phase transitions can only be probed by earthquake generated seismic waves through scattering off these. The relevant scattered wave constituents appear as precursors to, for example, the core-reflected compressional PcP phase and the horizontally polarized shear ScS phase. The scattering is most prominent at large opening angles (towards grazing incidence).
It is to this situation that the paraxial approximation and coefficient dependence considered here applies. In this context the principal direction, $`z`$, of propagation is perpendicular to the direction of backscattering, $`x`$. The central question is whether an imprint of the coefficient regularity on the (regularity of the) scattered wave occurs. A positive answer to this question, as provided in this paper, allows to infer a lack of lateral regularity in the medium from a lack of regularity of the measured data.
Seismic inverse scattering has been formulated mathematically in terms of evolution equations with respect to the depth variable in cf. (with smooth symbols in the single scattering approximation). At the basis of these models are one-way wave equations, which are typically of the form
$$\left(_z\pm iB(z,x,D_t,D_x)\right)u=f,$$
where $`B`$ is (microlocally) a pseudodifferential operator with principal part
$$b(z,x,\tau ,\xi )=\tau c(z,x)^1\sqrt{1\tau ^2c(z,x)^2|\xi |^2}.$$
Here, $`t`$ is time, $`z0`$ denotes depth, $`x^d`$ are lateral spatial directions, and $`D=i^1`$. Approximation of the square root to leading order results in the standard Schrödinger-type paraxial equation
$$\left(_z+\frac{i}{c(z,x)}D_t\right)w+\frac{1}{2i}D_t^1\underset{j=1}{\overset{d}{}}D_{x_j}\left(c(z,x)D_{x_j}w\right)=0,$$
where $`c(z,x)`$ is the local speed of propagation and $`z`$ plays the role of the evolution parameter. In the frequency ($`\tau `$-)domain the above equation is transformed with a so-called co-moving frame of reference according to $`\widehat{w}(z,x,\tau )=\widehat{u}(z,x,\tau )`$ $`\mathrm{exp}(i\tau T(z,x))`$, where $`T(z,x)=_0^z𝑑z^{}/c(z^{},x)`$. Then the paraxial equation attains the form
$$\left(_zi\overline{A}\right)\widehat{u}=0$$
in which $`\overline{A}`$ is given by
$$\overline{A}\widehat{u}=\frac{1}{2}\underset{j=1}{\overset{d}{}}e^{i\tau T(z,x)}D_{x_j}\left(c(z,x)\tau ^1D_{x_j}\left(e^{i\tau T(z,x)}\widehat{u}\right)\right).$$
The second-order differential operator $`\overline{A}`$ can be written as the sum of the self-adjoint operator
$$A(\tau ;z,x,D_x):=\frac{1}{2}\underset{j=1}{\overset{d}{}}D_{x_j}c(z,x)\tau ^1D_{x_j}$$
and a symmetric perturbation. We observe that if $`D_xT(z,.)_{L^{\mathrm{}}(^d)}<K`$, where $`K`$ is an appropriate constant, then it is guaranteed that this perturbation is $`A`$-bounded with relative bound less than $`1`$ (cf. \[14, Section X.2\]). From the viewpoint of generators of strongly continuous contraction semigroups on $`L^2(^d)`$ the simplification of $`\overline{A}`$ by $`A`$ is of no consequence.
###### Remark 1.1 (Directional decomposition and one-way wave equations).
Directional decomposition leads to the introduction of one-way wave equations . One-way wave equations approximate solutions to the wave equation microlocally, relative to a principal direction of propagation. (This principal direction does not need to be defined globally; one can introduce curvilinear coordinates and an associated Riemannian metric to generate such directions locally.) The validity of one-way wave propagation breaks down when singularities tend to propagate in a direction perpendicular to the principal direction (that is, in a transverse direction). Hence, to make the statement concerning approximation above, precise, one needs to introduce a microlocal attenuation . The mentioned procedures and results require smooth coefficients and symbols, and can be proven by making use of the calculus of pseudodifferential operators and Fourier integral operators with complex phase.
It has been demonstrated that, in special cases, one can weaken the condition of smooth coefficients. For example, if the coefficients are independent of the coordinate along the principal direction, one can allow a step function (in a transverse coordinate) and still solve the associated scattering problem by methods of one-way wave equations. The approach to carry out such an evaluation can be found in . Indeed, scattering in the transverse directions can, at least in special cases, be incorporated in the one-way wave equation. Moreover, in the case of wave propagation in random media, the (stochastic) paraxial equation naturally appears as well . However, a general result concerning directional decomposition for – or recomposition to solutions of – the wave equation with non-smooth coefficients has not been obtained.
In this paper, we address the general problem of “transverse scattering” by a one-way wave equation, which constitutes one component in the development of a general theory referred to above.
###### Remark 1.2 (Regularity and the second order wave equation).
Both, mode decoupling of a second order wave equation into one-way wave equations as well as the derivation of the narrow beam approximation outlined above, require higher order differentiablity of the coefficient $`c(z,x)`$ with respect to all variables to make sense (due to truncation of symbol expansions). However, the resulting paraxial wave equation displays precisely the same coefficient regularity as in the original second order wave equation. Hence as a model equation it still reflects the correct medium properties on all scales. The (exact) solutions of the paraxial equation then serve as a narrow beam approximation to solutions of the original wave equation. In particular, the regularity properties are comparable on the same scales.
The fine tuned well-posedness theorem for wave equations by Colombini and Lerner in assumes Log-Lipschitz regularity of the coefficients in the principal part. Moreover, their results are sharp in the sense that counterexamples to solvability exist when the coefficient regularity is below Log-Lipschitz (but still of any continuity type arbitrarily close to such). In order to indicate how their results relate to ours, we repeat the key energy estimate from :
A function $`aL^{\mathrm{}}(^d)`$ is said to be a *Log-Lipschitz* function if
$$a_{LL}:=\underset{x^d}{sup}|a(x)|+\underset{\genfrac{}{}{0pt}{}{xy^d,}{|xy|1/2}}{sup}\frac{|a(x)a(y)|}{|(xy)\mathrm{log}(|xy|)|}<\mathrm{}.$$
Colombini and Lerner consider second order wave operators of the form
$$Pu:=_t^2u\underset{1j,kn}{}_{x_j}\left(a_{jk}(x,t)_{x_k}u\right)+M(x,t,_t,_x)u,$$
where $`(a_{jk})_{1j,kn}`$ is real symmetric with Log-Lipschitz components and satisfies with some $`1\delta _0>0`$ the strong ellipticity condition
$$\underset{1j,kn}{}a_{jk}(x,t)\xi _j\xi _k\delta _0|\xi |^2(\xi ^n,(x,t)^{n+1}),$$
and $`M(x,t,_t,_x)`$ is a first-order differential operator with Hölder-continuous coefficients.
Let $`P`$ be as above and $`\theta ]0,1/4]`$. There exist $`\beta >0`$ and $`C>0`$ such that for $`u𝒞^{\mathrm{}}(^{n+1})`$ and $`t[0,1/\beta ]`$ the *energy estimate*
$$\begin{array}{c}\underset{0st}{sup}_tu(.,s)_{H^{\theta \beta s}}+\underset{0st}{sup}u(.,s)_{H^{1\theta \beta s}}\hfill \\ \hfill C(_0^tPu(.,s)_{H^{\theta \beta s}}ds+_tu(.,0)_{H^\theta }+u(.,0)_{H^{1\theta }}),\end{array}$$
holds (cf. \[4, Equation (2.6)\]), where $`\beta `$ depends only on $`\delta _0`$, on the Log-Lipschitz norm of the $`a_{jk}`$, and on the Hölder norms of the coefficients in $`M(x,t,_t,_x)`$.
We show in the sequel that for the paraxial wave equation the condition on the coefficient regularity can be relaxed. For example, if $`\epsilon >0`$ any function in $`H^{1+\epsilon }(^2)`$ of local behavior like $`x|x|^{1/2+\epsilon }`$ is not Log-Lipschitz continuous but satsifies the assumptions of our main results below.
The low coefficient regularity in our model conditions has its price in terms of a few technical aspects of the current paper: Additional care is needed in identifying the appropriate distribution and function spaces that allow for the description of mixed regularity properties and for the rigorous formulation of a solution concept. The existence proof then consists in showing a series of functional analytic properties to establish an evolution system of operators; among these the basic self-adjointness property — in the disguise of an elliptic regularity lemma — is derived by employing rather delicate regularity properties of multiplication in certain subspaces of the space of distributions.
The plan of the paper is as follows. In Section 2 we present the precise form of the operator and specify our (low) regularity assumptions on the coefficients. The solution will be sought as a continuous map of depth into the space of temperate $`L^2`$-valued distributions. Section 3 is devoted to the construction of the evolution system in the frequency domain. First, we prove that $`A(\tau ;z,x,D_x)`$ generates a unitary group at fixed $`z`$ and $`\tau `$. The determination of its domain requires delicate use of the duality product of distributions as well as a bootstrap argument involving multiplication in scales of Sobolev spaces. This leads to the construction of an evolution system at fixed $`\tau `$. Finally, strongly continuous dependence on the frequency parameter $`\tau `$ is established, based on a difference approximation. Again a subtle interplay between regularity arguments and distributional products is at the heart of our arguments. The strong continuity enables us, in Section 4, to construct a solution of the evolution system in frequency domain with distributional data. Finally, existence, uniqueness, and regularity of solutions to the original Cauchy problem is obtained. As an application to inverse regularity analysis, we obtain that a lack of $`H^2`$-regularity in the observed solution implies the existence of a region in which the lateral regularity of the medium is at most $`H^1`$ on the Sobolev scale.
## 2 The Cauchy problem: function spaces and <br>coefficient regularity
We recall the definition of temperate distributions on $``$ with values in a Banach space $`E`$ (cf. \[23, Section 39.3\]; but note that we use a different topology here): let $`𝒮^{}(;E)`$ be the space of continuous linear maps $`𝒮()E`$, equipped with the topology of pointwise convergence; the Fourier transform $``$ on $`𝒮()`$ is extended to $`𝒮^{}(;E)`$ by setting $`(G)(\varphi )=G(\varphi )`$ ($`G𝒮^{}(;E)`$, $`\varphi 𝒮()`$), which is easily seen to be an isomorphism (for the locally convex structure).
We denote the time variable by $`t`$ and introduce coordinates $`z[0,\mathrm{})`$ for depth (the evolution direction in our context) and $`x^d`$ for the lateral variation, where $`d2`$. As basic space of the *wave components* we consider
(2.1)
$$𝒱:=𝒞([0,\mathrm{}),𝒮^{}(;L^2(^d)).$$
Its elements are continuous maps $`zu(z)`$ of the depth variable $`z`$ into temperate distributions of time $`t`$ valued in $`L^2`$-functions of the lateral variables $`x`$. When we need to keep track of precise regularity information in the lateral variation of the waves, we may employ the Sobolev-scale $`H^s(^d)`$ ($`s`$) and define
(2.2)
$$𝒱^s:=𝒞([0,\mathrm{}),𝒮^{}(;H^s(^d)).$$
Let $`_t:𝒮^{}(;H^s(^d))𝒮^{}(;H^s(^d))`$ denote (partial) Fourier transform with respect to the time variable. We extend $`_t`$ to an isomorphism $`\stackrel{~}{_t}`$ of (the locally convex structure of) $`𝒱^s`$ by
$$(\stackrel{~}{_t}u)(z):=_t(u(z,.))z0.$$
We consider the following Cauchy problem for a prospective solution $`u𝒱`$ (div and grad with respect to $`x^d`$):
(2.3) $`Pu:=_zui\text{div}(C(z,x,D_t)\text{grad}u)`$ $`=f𝒱`$
(2.4) $`u_{z=0}`$ $`=u_0𝒮^{}(;L^2(^d)).`$
Here, $`C`$ is a pseudodifferential operator in $`t`$ with parameters $`z`$ and $`x`$. While in the classical paraxial wave equation it is of order $`1`$, here we may assume it is of some order $`m`$. The precise conditions are collected in the following
###### Assumption 1.
The symbol of $`C(z,x,D_t)`$ is of the form
(2.5)
$$C(z,x,\tau )=c(z,x,\tau )I_d=\left(c_0+\underset{l=1}{\overset{N}{}}c_l(z,x)h_l(\tau )\right)I_d,$$
where $`N`$, $`I_d`$ is the $`d\times d`$ identity matrix and the following hold:
1. For $`l=1,\mathrm{},N`$: $`h_l`$ is a real-valued smooth symbol (of order $`m`$) on $``$, i.e., for all $`k_0`$ an estimate $`|_\tau ^kh_l(\tau )|=O(|\tau |^{mk})`$ holds when $`|\tau |`$ is large; in addition, we assume that $`|h_l(\tau )|\eta _0`$ near $`\tau =0`$ ($`l=1,\mathrm{},N`$) with some constant $`\eta _0>0`$ (this can be achieved by adding a cut-off function without changing the relevant frequency range).
2. $`c_0`$ is a positive constant.
3. there is an $`r(0,1)`$ such that for $`1lN`$: $`c_l`$ is in $`𝒞^1([0,\mathrm{}),H^{r+1}(^d))`$ and real-valued.
4. for all $`(z,x,\tau )[0,\mathrm{})\times ^d\times `$: $`c(z,x,\tau )c_0`$.
###### Remark 2.1.
(i) The operator action corresponding to a typical term $`c_l(z,x)h_l(D_t)`$ in the sum decomposition (2.5) on any element $`w`$ of $`𝒱`$ is given as follows: for all $`\varphi 𝒮()`$
$$\left(h_l(D_t)w(z)\right)(\varphi )=w(z)((h_l^1\varphi ))L^2(^d),$$
which is then multiplied by the function $`c_l(z,.)`$.
(ii) Note that parts (i-iii) of the Assumption imply, for any $`z`$ and $`\tau `$ fixed, the following (Zygmund-)Hölder-continuity:
$$c(z,.,\tau )c_0H^{r+1}(^d)C_{}^{r+1\frac{d}{2}}(^d)$$
(cf. \[9, Proposition 8.6.10\]). Thus the coefficients have lateral Hölder-regularity of order $`r+1\frac{d}{2}`$. In the most relevant case from geophysics, $`d=2`$, this yields coefficients in $`C_{}^r(^2)`$, but not necessarily Lipschitz continuous. (This includes the situation of a boundary layer, as discussed in the introduction, where the coefficient is smooth in one of the two variables.) Part (iv) implies uniform ellipticity of the lateral differential operator.
(iii) Note that in dimension $`d3`$ Sobolev regularity of order $`r+1`$ would not imply continuity of the coefficients (in case $`1<r+1<d/2`$).
###### Example 2.2.
We consider a model with two-dimensional lateral variation ($`d=2`$) in the medium properties of low Hölder regularity depending on depth. In (2.5) we put $`c_l=0`$ when $`l2`$ and let $`c_1`$ be of the form
$$c_1(z,x)=\chi (z,x)|x|^{\alpha (z)}$$
with the following specifications: $`\chi 𝒞^1([0,\mathrm{})\times ^d)`$ such that $`\chi (z,x)=\chi _0(z)`$ when $`|x|R_1`$ and $`\chi (z,x)=0`$ when $`|x|R_2`$ for certain radii $`0<R_1<R_2`$ and some $`\chi _0𝒞^1([0,\mathrm{}))`$; $`\alpha 𝒞^1([0,\mathrm{}))`$ with some uniform positive lower bound $`\alpha _0`$, i.e., $`\alpha (z)\alpha _0>0`$ for all $`z`$. Then we may choose any $`r`$ such that $`0<r<\alpha _0`$ and obtain the following regularity properties at arbitrary fixed values of $`z`$ and $`\tau `$:
$$c(z,.,\tau )c_0H^{r+1}(^2)C_{}^{\alpha (z)}(^2)H^{r+1}(^2)C_{}^{\alpha _0}(^2).$$
(Observe that locally in two dimensions, for any $`0<s<1`$, the function $`|x|^s`$ belongs to $`C_{}^s`$ and to $`H^{s+1\epsilon }`$ for every $`\epsilon >0`$ but not to $`H^{s+1}`$.)
Applying $`\stackrel{~}{_t}`$ to (2.3-2.4) we obtain an equivalent formulation of the Cauchy problem in the frequency domain:
(2.6) $`\stackrel{~}{P}v=_zvi\text{div}(c(z,x,\tau )\text{grad}v)`$ $`=g𝒱`$
(2.7) $`v_{z=0}`$ $`=v_0𝒮^{}(;L^2(^d)).`$
Equation (2.6) is an evolution equation for depth $`z`$ with the second-order operator
(2.8)
$$A(\tau ;z,x,D_x)v:=\text{div}(c(z,x,\tau )\text{grad}v)$$
acting in the lateral $`x`$-domain and smoothly depending on the “external” parameter $`\tau `$. Note that $`A(\tau ;z,x,D_x)`$ is uniformly elliptic by Assumption 1 (iv).
###### Remark 2.3.
Note that $`P`$ in (2.3) commutes with convolution in the time variable. Therefore, damping (or cut-off) of high frequencies in the data of (2.6-2.7) corresponds to time-smoothing the data in the original problem (2.3-2.4): more precisely, if a frequency filter $`\widehat{\chi }𝒮()`$ is applied by $`g:=\stackrel{~}{_t}(f)\widehat{\chi }(\tau )`$, $`v_0:=_t(u_0)\widehat{\chi }(\tau )`$ and $`v`$ is a solution to (2.6-2.7) then $`u:=\stackrel{~}{_t}^1(v)`$ solves (2.3-2.4) with the data changed to $`f\underset{(t)}{}\chi `$ and $`u_0\underset{(t)}{}\chi `$.
## 3 Evolution system
In this section, we will show that, up to any finite depth $`Z>0`$, the family of unbounded operators $`iA(\tau ;z,x,D_x)`$ ($`\tau `$, $`z0`$) generates a strongly continuous evolution system (or fundamental solution) $`\{U(\tau ;z_1,z_2):0z_1z_2Z\}`$ on $`L^2(^2)`$ in the sense of \[20, Chapter 4\] which, in addition, is strongly continuous in the frequency variable $`\tau `$. In a first step we freeze both parameters, $`\tau `$ as well as $`z`$, and construct a strongly continuous (semi-)group of operators on $`L^2(^d)`$.
### 3.1 Unitary group at frozen values of $`\tau `$ and $`z`$
#### Notational simplifications:
By abuse of notation we will employ the short-hand symbols $`A:=A(\tau ;z,x,D_x)`$, $`c(x):=c(z,x,\tau )`$, and $`c_1(x)`$ now denoting $`_{l1}c_l(z,x)h_l(\tau )`$. To summarize, using the above conventions and Assumption 1, we have
(3.9)
$$Av=\text{div}(c(x)\text{grad}v)$$
as unbounded, formally self-adjoint operator on $`L^2(^d)`$ with coefficient
(3.10)
$$c(x)=c_0+c_1(x)\text{ with }c_0>0\text{ and }0c_1H^{r+1}(^d).$$
We will show that $`A`$ is a self-adjoint operator with domain $`D(A)=H^2(^d)`$.
###### Remark 3.1.
Note that self-adjointness of $`A`$ with domain $`H^2`$ would be immediate from uniform ellipticity in case the coefficient were smooth. On the other hand, self-adjointness on some domain could be obtained in an abstract fashion via quadratic forms under mere $`L^{\mathrm{}}`$-assumptions (\[15, Section VIII.6\]). However, in accordance with our focus on the interplay of the coefficient regularity class with qualitative solution properties, we will give an explicit domain description in terms of Sobolev spaces, which in addition is uniform with respect to $`\tau `$ and $`z`$.
Observe that, due to the low coefficient regularity, we also have to establish that $`A`$ is well-defined on all of $`H^2(^d)`$ as an operator into $`L^2(^d)`$. This is included in the following lemma as the special case $`s=0`$.
###### Lemma 3.2.
Let $`0s<r<1`$ and $`vH^{s+2}(^d)`$. Then $`AvH^s(^d)`$.
###### Proof.
Each component of $`\text{grad}v`$ is in $`H^{s+1}(^d)`$. Hence multiplying with the $`H^{r+1}`$-coefficient $`c_1`$, as well as with the constant $`c_0`$, is well-defined within $`H^{s+1}(^d)`$ since this is an algebra. ∎
We set $`D(A):=H^2(^d)`$ and note that an integration by parts immediately yields that $`A`$ is symmetric, i.e., $`D(A)D(A^{})`$ and $`A^{}_{D(A)}=A`$. We proceed to show that also $`D(A^{})D(A)`$ by which self-adjointness will be established.
By definition, the adjoint operator has domain
$$D(A^{})=\{vL^2\text{for some }wL^2:\psi |w=A\psi |v\text{ for all }\psi H^2\},$$
where $`|`$ denotes the inner product in $`L^2(^d)`$. Let $`vD(A^{})`$ and choose a sequence $`(v_k)_k`$ in $`H^2(^d)`$ which converges to $`v`$ in $`L^2(^d)`$. On the one hand, there exists $`wL^2(^d)`$ such that we have for all test functions $`\psi 𝒟(^d)`$
$$A^{}v_k|\psi =v_k|A\psi v|A\psi =w|\psi \text{ when }k\mathrm{}.$$
Thus, $`(A^{}v_k)_k`$ has the distributional limit $`wL^2(^d)`$. On the other hand, since $`c_1H^1(^d)`$ and $`\text{grad}v_k\text{grad}v`$ in $`H^1(^d)`$ (as $`k\mathrm{}`$) we may employ the continuous duality product of distributions (cf. \[11, Proposition 5.2\]) $`H^1\times H^1W_{\text{loc}}^{1,1}`$ and deduce that $`A^{}v_k=Av_k\text{div}(c(x)\text{grad}v)=Av`$ in $`W_{\text{loc}}^{2,1}`$, hence in the sense of distributions. By uniqueness of distributional limits, we deduce that $`Av=wL^2(^d)`$. We obtain
$$D(A^{})=\{vL^2AvL^2\}.$$
The assertion $`D(A^{})D(A)=H^2(^d)`$ follows now from the following result. For later reference, it is stated in slightly more general terms. (We first consider the important case $`d=2`$ and leave the case of one-dimensional lateral variation for a remark below.)
###### Lemma 3.3 (Elliptic regularity).
Let $`0s<r<1`$ and $`vH^s(^2)`$ such that $`AvH^s(^2)`$. Then $`v`$ belongs to $`H^{s+2}(^2)`$.
The proof will be based on repeated use of the following three facts, which we collect in a preparatory list of “ingredients”:
Let $`w_jH^{s_j}(^2)`$ ($`j=1,2`$) such that $`s_1+s_20`$. Then
$$w_1w_2H^{s_0}(^2),$$
where
$$s_0=\{\begin{array}{cc}\mathrm{min}(s_1,s_2,s_1+s_21)\hfill & \text{ if }s_1\pm 1,s_2\pm 1,\text{ and }s_1+s_20\hfill \\ \mathrm{min}(s_1,s_2,s_1+s_21\epsilon )\hfill & \text{ with }\epsilon >0\text{ arbitrary, otherwise}.\hfill \end{array}$$
This is included in the statement of \[9, Theorem 8.3.1\]. As can be seen from the proof therein, one also obtains continuity of the multiplication $`H^{s_1}\times H^{s_2}H^{s_0}`$ (with respect to the corresponding Sobolev-norms).
We can find a function $`F𝒞^{\mathrm{}}()`$, $`F(0)=0`$, such that
$$\frac{1}{c(x)}=\frac{1}{c_0}+F(c_1(x)).$$
In particular, we obtain $`1/c1/c_0=F(c_1)H^{r+1}(^2)`$. Since $`r>0`$, this follows from \[9, Theorem 8.5.1\] once $`F`$ is given. To find $`F`$, we simply set $`F(y):=y/(c_0(c_0+y))`$ when $`yc_0/2`$ and extend it in a smooth way to $``$ such that $`F(0)=0`$.
If $`vL^2(^2)`$ the equation
$$Av=\text{div}(c\text{grad}v)=\mathrm{\Delta }(cv)\text{div}(v\text{grad}c)$$
holds in $`𝒟^{}(^2)`$, where the occurring products are defined as follows: using $`\text{grad}vH^1`$ we get $`c\text{grad}vW_{\text{loc}}^{1,1}`$ by the duality product \[11, Proposition 5.2\]; $`cvL^2`$ since $`cL^{\mathrm{}}`$; and $`v\text{grad}cL^1`$ because $`_jcH^rL^2`$.
Under the stronger assumption $`vH^{r+1}(^2)`$ we have, in addition, that
$$Av=\text{grad}c\text{grad}v+c\mathrm{\Delta }v$$
in $`𝒟^{}(^2)`$, with the meaning of the products on the right-hand side as follows: since $`\mathrm{\Delta }vH^{r1}`$, Fact A applies and yields $`c\mathrm{\Delta }vH^{r1}`$; furthermore, $`\text{grad}c`$ and $`\text{grad}v`$ both lie in $`H^r`$, so that another application of Fact A shows that their (Euclidean inner) product belongs to $`H^{\mathrm{min}(r,2r1)}`$.
###### Remark 3.4.
Note that formula (C2) represents $`A`$ as an operator with a (Hölder-) continuous coefficient in its principal part and Sobolev regularity in the lower orders. We observe that in such situation, \[8, Theorem 17.1.1\] gives local solvability in $`H^2`$ for right-hand sides in $`L^2`$. However, the latter does not imply $`H^2`$-regularity of any $`L^2`$-solution. For the pure regularity question, it also seems that methods based on perturbations of constant coefficient operators do not apply either, since $`A`$ is not necessarily of constant strength (cf. \[7, Chapter XIII\]).
###### Proof of Lemma 3.3.
To begin with, we only know that $`v`$ as well as $`Av`$ belong to $`H^s(^2)`$. The proof proceeds in three steps, successively revealing higher regularity.
*Claim 1:* $`vH^{s+r}(^2)`$
We have $`\text{grad}cH^r`$, so that application of Fact A, noting that $`r1<0`$, gives $`v\text{grad}cH^{s+r1}`$, hence $`\text{div}(v\text{grad}c)H^{s+r2}`$. Since $`AvH^s`$ we deduce from equation (C1) that $`\mathrm{\Delta }(cv)H^{s+r2}`$, which in turn implies that $`cvH^{s+r}`$. Now invoke the decomposition from Fact B and write
$$v=\frac{1}{c_0}cv+F(c_1)cv.$$
The first part clearly is in $`H^{s+r}`$, and for the second summand the same is true by Fact A. (Note that $`\mathrm{min}(s+r,s+2r\epsilon )=s+r`$ if $`0<\epsilon <r`$.)
*Claim 2:* $`vH^{r+1}(^2)`$
We may start from $`vH^{r+s}`$ by claim 1 and proceed inductively to show that
$$vH^{r+\mathrm{min}(1,s+jr/2)}j0.$$
Claim 2 then follows upon choosing $`j`$ sufficiently large (i.e., $`j2(1s)/r`$ steps will be required).
The case $`j=0`$ is just claim 1, so we assume that the assertion holds for some $`j0`$. If $`s+jr/21`$ it is trivially satisfied for larger values of $`j`$, therefore we assume $`t_j:=s+jr/2<1`$ and that
$$vH^{r+t_j}.$$
Fact A gives $`v\text{grad}cH^{r+t_j}H^rH^{\mathrm{min}(r,t_j+2r1\epsilon )}H^{\mathrm{min}(r,t_j+3r/21)}`$ upon choosing $`\epsilon <r/2`$. Thus, using the short-hand notation $`r_j:=\mathrm{min}(r+1,t_j+3r/2)r+1`$ we may infer that $`\text{div}(v\text{grad}c)H^{r_j2}`$. Since $`r_j2r1<0<s`$ we also have $`AvH^sH^{r_j2}`$. Equation (C1) now implies that $`\mathrm{\Delta }(cv)H^{r_j2}`$, hence $`cvH^{r_j}`$. Again by Fact B, combined with Fact A, we obtain
$$v=\frac{1}{c_0}cv+F(c_1)cvH^{r+1}H^{r_j}H^{r_j}=H^{r+\mathrm{min}(1,t_j+r/2)},$$
which means ($``$) for $`j+1`$ in place of $`j`$.
*Claim 3:* $`vH^{s+2}(^2)`$
We use a similar strategy as in the proof of claim 2 and prove inductively that
$$vH^{\mathrm{min}(s+2,1+(j+1)r/2)}j1.$$
Claim 3 then follows when $`j`$ is chosen sufficiently large (i.e., $`j2(1+s)/r1`$ steps are required).
The basic case $`j=1`$ corresponds to claim 2, so we proceed with some $`j1`$, under the additional assumption $`rs_j:=(j+1)r/2<s+1`$ to exclude trivial cases. In other words, we assume that
$$vH^{s_j+1}.$$
Therefore, again by Fact A and choosing a possibly occurring $`\epsilon <r/2`$, we deduce
$$\text{grad}c\text{grad}vH^rH^{s_j}H^{\mathrm{min}(r,s_j+r/21)}=H^{\mathrm{min}(r,q_j)},$$
where we have introduced $`q_j:=s_j+r/21`$. Exploiting equation (C2) and noting that $`AvH^s`$, $`s<r`$, we extract the information that $`c\mathrm{\Delta }vH^{\mathrm{min}(s,q_j)}`$. Once again we use the decomposition corresponding to Fact B and the statement of Fact A to deduce
$$\mathrm{\Delta }vH^{\mathrm{min}(s,q_j)}+H^{r+1}H^{\mathrm{min}(s,q_j)}H^{\mathrm{min}(s,q_j)}$$
and a fortiori that $`\mathrm{\Delta }vH^{\mathrm{min}(s+2,q_j+2)}`$. But $`q_j=(j+1)r/2+r/21`$ hence $`q_j+2=1+(j+2)r/2`$ and ($``$) is proven with $`j+1`$ in place of $`j`$. ∎
###### Remark 3.5.
The one-dimensional analogue of Lemma 3.3 is more elementary: $`Av=(cv^{})^{}H^s()`$ implies $`cv^{}H^{s+1}()`$ and, since Fact B is valid for $`d=1`$ as well, we obtain $`v^{}H^{s+1}()`$; thus, $`vH^{s+2}()`$.
We summarize the intermediate conclusions from the discussion so far in a separate statement, where we appeal to Stone’s theorem providing us with the exponential unitary group $`T(z)=\mathrm{exp}(izA)`$.
###### Theorem 3.6.
Let $`c𝒞(^d)`$ satisfy (3.10) and define $`Av=\text{div}(c(x)\text{grad}v)`$ with domain $`D(A)=H^2(^d)`$ in $`L^2(^d)`$. Then $`A`$ is self-adjoint and $`iA`$ generates a strongly continuous unitary group $`(T(z))_z`$ on $`L^2(^d)`$. Moreover, we have the following resolvent estimate, valid for $`\lambda \{0\}`$:
(3.10)
$$(iA\lambda )^1v_{L^2}\frac{v_{L^2}}{|\lambda |}\text{ for all }vL^2(^d).$$
We briefly recall how $`(T(z))_z`$ can be used to construct solutions to the Cauchy problem on $`^d\times `$
$$_zviAv=g,v(0)=v_0.$$
Let $`v_0L^2(^d)`$ and $`gL^1(,L^2(^d))`$ then the *mild solution*
(3.11)
$$v(z):=T(z)v_0+_0^zT(z\rho )g(\rho )𝑑\rho $$
is in $`𝒞(,L^2(^d))`$.
If $`v_0H^2(^d)`$ and $`g𝒞(,H^2(^d))`$ or $`g𝒞^1(,L^2(^d))`$ then $`v`$ belongs to $`𝒞^1(,L^2(^d))`$ and is the unique classical solution with pointwise values in $`H^2(^d)`$ (cf. \[13, Section 4.2, Corollaries 2.5 and 2.6\]).
###### Remark 3.7.
The mild solution (3.11) defines a weak solution in the following sense: $`v(0)=v_0`$ and for all $`\varphi 𝒟(^{d+1})`$
$$_{}\left(v(z)|_z\varphi (z,.)iA\varphi (z,.)\right)dz=_{}g(z)|\varphi (z,.)dz,$$
where $`|`$ denotes the inner product in $`L^2(^d)`$. To see this, one approximates the mild solution by classical solutions $`(v_k)_k`$ to equations with regularized right-hand side and initial data (\[13, Section 4.2, Theorem 2.7\]): $`L^2`$-convergence of $`v_k(z)v(z)`$ (as $`k\mathrm{}`$), uniformly when $`z`$ varies in compact intervals, together with the convergence $`g_kg`$ in $`L^1(,L^2(^d))`$ implies convergence in the integral formula above.
### 3.2 Evolution system at fixed frequency $`\tau `$
Let $`\tau `$ be fixed, but arbitrary. We consider the $`z`$-parameterized family of unbounded self-adjoint operators in (2.8) and put
$$A(\tau ;z):=A(\tau ;z,x,D_x)(z0).$$
Let $`Z>0`$ be arbitrary. We will check that, for every $`\tau `$, $`(iA(\tau ;z))_{z0}`$ defines an evolution system (or fundamental solution) $`(U(\tau ;z_1,z_2))_{Zz_1z_20}`$ on $`L^2(^d)`$ by applying \[20, Section 4.4, Corollary to Theorem 4.4.2, p. 102\] (cf. also \[13, Sections 5.3-5.5\]). We have to check that the corresponding hypotheses are satisfied.
First, observe that $`D(A(\tau ;z))=H^2(^d)`$ is independent of the evolution parameter $`z`$ (and of $`\tau `$), and every $`iA(\tau ;z)`$ is the skew-adjoint generator of a strongly continuous (unitary) semigroup $`(T(\tau ,z;\zeta )_{\zeta 0})`$ on $`L^2(^d)`$. Furthermore, the resolvent estimates (3.10), valid for all $`z`$ (and $`\tau `$), immediately imply that $`(A(\tau ;z))_{z0}`$ is a stable family of generators *with stability constants $`1`$ and $`0`$ for all $`\tau `$* (cf. \[20, Definition 4.3.1\]).
Finally, we have to check that for all $`vH^2(^d)`$ the map
$$[0,\mathrm{})zA(\tau ;z)vL^2(^d)$$
is continuously differentiable. We may use equation (C2) from Fact C to write (with grad taken with respect to $`x`$ only)
$$A(\tau ;z)v=\text{grad}c(z,x,\tau )\text{grad}v+c(z,x,\tau )\mathrm{\Delta }v.$$
By Assumption 1,(ii-iii), we have
$$\text{grad}c(.,.,\tau )𝒞^1([0,\mathrm{}),H^r(^d)),c(.,.,\tau )c_0h_0(\tau )𝒞^1([0,\mathrm{}),H^{r+1}(^d)).$$
Since $`\text{grad}vH^1(^d)`$ and $`\mathrm{\Delta }vL^2(^d)`$ the multiplication rules plus continuity properties in Fact A apply (where in case $`H^rH^1`$ we choose $`\epsilon <r`$, if $`d=2`$) and yield that $`A(\tau ;.)v𝒞^1([0,\mathrm{}),L^2(^d))`$.
Thus, all hypotheses of \[20, Section 4.4, Corollary to Theorem 4.4.2, p. 102\] are fulfilled. Note that the evolution system is constructed as the strong operator limit of discretizations based on the unitary semigroups of each generator, hence is contractive. This implies the following intermediate result.
###### Proposition 3.8.
Let $`Z>0`$. Then for all $`\tau `$ the family $`(iA(\tau ;z))_{z0}`$ defines a unique evolution system $`(U(\tau ;z_1,z_2))_{Zz_1z_20}`$ on $`L^2(^2)`$ with the following properties: The map $`(z_1,z_2)U(\tau ;z_1,z_2)`$ is strongly continuous, $`U(\tau ;z,z)=I`$, $`U(\tau ;z_1,z_2)`$ is contractive, and
(3.12)
$$U(\tau ;z_1,z_2)U(\tau ;z_2,z_3)=U(\tau ;z_1,z_3)0z_3z_2z_1Z;$$
moreover, $`H^2(^d)`$ is invariant under $`U(\tau ;z_1,z_2)`$, for all $`vH^2(^d)`$ the map $`(z_1,z_2)U(\tau ;z_1,z_2)v`$ is continuously differentiable, separately in both variables, and the following equations hold:
(3.13) $`{\displaystyle \frac{}{z_1}}U(\tau ;z_1,z_2)v`$ $`=A(\tau ;z_1)U(\tau ;z_1,z_2)v`$
(3.14) $`{\displaystyle \frac{}{z_2}}U(\tau ;z_1,z_2)v`$ $`=U(\tau ;z_1,z_2)A(\tau ;z_2)v.`$
At this stage, we obtain solutions to a version of the Cauchy problem (2.6-2.7) at fixed frequency $`\tau `$, i.e.,
(3.15) $`_zviA(\tau ;z)v`$ $`=gL^1([0,Z],L^2(^d))`$
(3.16) $`v_{z=0}`$ $`=v_0L^2(^d).`$
The *mild solution* is defined by
(3.17)
$$v(z):=U(\tau ;z,0)v_0+_0^zU(\tau ;z,\rho )g(\rho )𝑑\rho $$
and belongs to $`𝒞([0,Z],L^2(^d))`$ (\[13, Section 5.5, Definition 5.1\]).
###### Remark 3.9.
(i) In the case of classical solutions, we have the following regularity property: If $`v_0H^2(^d)`$ and $`g𝒞([0,Z],H^2(^d))`$ or $`g𝒞^1([0,Z],L^2(^d))`$ then $`v𝒞^1([0,Z],L^2(^d))`$ is the unique *$`H^2`$-valued solution* and satisfies the equation in the strong sense (cf. \[13, Section 5.5, Theorems 5.2 and 5.3\]).
(ii) Observe that, at frozen value of $`\tau `$, one may apply \[10, Chapter 3, Theorem 10.1 and Remark 10.2\] directly by putting $`H=L^2`$, $`V=H^1`$, and
$$a(z;u,v)=\underset{j=1}{\overset{d}{}}c(z,.)_ju|_jvu,vV.$$
It suffices to assume $`c𝒞^1([0,Z],L^{\mathrm{}}(^d))`$, then for any initial value $`v_0H^1(^d)`$ and right-hand side $`gL^2([0,Z]\times ^d)`$ such that $`_zgL^2([0,Z],H^1)`$ there is a unique solution $`v𝒞([0,Z],H^1)𝒞^1([0,Z],H^1)`$ to the Cauchy problem (3.15-3.16). However, our approach allows for a precise investigation of the $`\tau `$-dependence, which is needed to solve the full Cauchy problem (2.6-2.7) with distributional data as well as to transform back to the original problem (2.3-2.4) in Section 4. Furthermore, our results show that lateral $`H^2`$-regularity of the data is preserved in the solution.
We thus have established an evolution system in the $`L^2`$-setting. Note that by Lemma 3.3 we have, in fact, that $`A(\tau ;z,x,D_x)`$ is an unbounded operator on $`H^s`$ with domain $`H^{s+2}`$ for any $`0s<r`$. If we were able to establish an evolution system on $`H^s`$ then the regularity information encoded into $`A`$ would be more directly preserved.
### 3.3 Frequency dependence of the evolution system
Throughout this subsection, let $`Z>0`$ be arbitrary but fixed. So far, the frequency parameter $`\tau `$ was arbitrary, but fixed, throughout the construction of the evolution system $`(U(\tau ;z_1,z_2))_{Zz_1z_20}`$. We will prove that the dependence on all parameters $`(\tau ,z_1,z_2)`$ jointly is strongly continuous. In the sequel, let $`L(E,F)`$ (resp. $`L(E)`$) denote the set of bounded linear operators between the Banach spaces $`E`$ and $`F`$ (resp. on $`E`$).
We begin with an observation on the general level of semigroups and evolution systems.
###### Lemma 3.10.
Assume that
(3.18)
$$(\tau ,z)A(\tau ;z)\text{ is continuous }\times [0,\mathrm{})L(H^2(^d),L^2(^d))$$
(with respect to the operator norm) and
(3.19)
$$(\tau ,z,\zeta )T(\tau ,z;\zeta )\text{ is strongly continuous }\times [0,\mathrm{})\times [0,\mathrm{})L(L^2(^d)),$$
where $`(T(\tau ,z;\zeta ))_{\zeta 0}`$ denotes the semi-group generated by $`A(\tau ;z)`$. Then the map $`(\tau ,z_1,z_2)U(\tau ;z_1,z_2)`$ is strongly continuous from $`B:=\times \{(z_1,z_2):Zz_1z_20\}`$ into $`L(L^2(^d))`$.
###### Proof.
We inspect the basic construction of the evolution system from the family of semigroups in the proof of \[13, Section 5.3, Theorem 3.1\] and keep track of the additional parameter $`\tau `$ in our case. For all $`(\tau ,z_1,z_2)B`$ we obtain $`U(\tau ;z_1,z_2)`$ as the strong limit of $`U_n(\tau ;z_1,z_2)`$ (as $`n\mathrm{}`$), where $`U_n(\tau ;.,.)`$ is the evolution system defined as follows: put $`z_n^j=jZ/n`$ ($`j=0,\mathrm{},n`$) then for $`\tau `$, $`0yzZ`$ let
$$U_n(\tau ;z,y):=T(\tau ,z_n^l;zy)\text{if }z_n^lyzz_n^{l+1},$$
and
$$\begin{array}{c}U_n(\tau ;z,y):=T(\tau ,z_n^k;zz_n^k)\underset{l+1jk1}{}T(\tau ,z_n^j;Z/n)T(\tau ,z_n^l;z_n^{l+1}y)\hfill \\ \hfill \text{if }z_n^lyz_n^{l+1}z_n^kzz_n^{k+1},k>l.\end{array}$$
By (3.19) the right-hand side of each formula is strongly continuous with respect to $`(\tau ,z,y)`$, and the boundary values, when $`k=l+1`$ and $`y=z_n^{l+1}`$ or $`z=z_n^{l+1}`$, match. Hence $`U_n`$ is strongly continuous on $`B`$ and $`U_n(\tau ;z,y)=1`$.
As in \[13, (3.13) on p. 136\] we have the following integral representation for the action on any $`vH^2`$
$$U_n(\tau ;z,y)vU_m(\tau ;z,y)v=_y^zU_n(\tau ;z,\rho )\left(A_n(\tau ;\rho )A_m(\tau ;\rho )\right)U_m(\tau ;\rho ,y)v𝑑\rho ,$$
where $`A_n(\tau ;\rho )`$ is the piecewise constant approximation of $`A(\tau ;\rho )`$ with $`A_n(\tau ;\rho ):=A(\tau ;z_n^k)`$, when $`z_n^k\rho <z_n^{k+1}`$, and $`A_n(\tau ;Z)=A(\tau ;Z)`$. By (3.18) we have $`A_n(\tau ;\rho )A(\tau ;\rho )_{L(H^2,L^2)}0`$ uniformly for $`(\tau ,\rho )`$ in compact sets. Passing to the limit $`m\mathrm{}`$ in the integral representation above yields the estimate
$$U_n(\tau ;z,y)vU(\tau ;z,y)v_{L^2}v_{H^2}_y^zA_n(\tau ;\rho )A(\tau ;\rho )_{L(H^2,L^2)}𝑑\rho .$$
By the uniform convergence of $`A_n(\tau ;\rho )`$ (as $`n\mathrm{}`$) we thus obtain (local) uniform convergence of $`U_n(\tau ;z,y)v`$, which proves the asserted continuity of $`(\tau ,z,y)U(\tau ;z,y)v`$. ∎
We have to establish conditions (3.18-3.19) in the specific context of the assumptions described in Section 2. In due course, we will make repeated use of $`(\tau ,z)`$-parameterized variants of Facts A-C, stated in Subsection 3.1. Note that, in particular, the function $`F`$ used in Fact B does not depend on $`(\tau ,z)`$.
###### Lemma 3.11.
If $`A(\tau ;z)`$ ($`\tau `$, $`z0`$) is defined by (2.8) then Assumption 1 implies Lipschitz-continuity of the map in condition (3.18).
###### Proof.
Let $`M(\tau ,z)`$ denote the operator of multiplication of pairs $`(v_1,v_2)H^1\times H^1`$ by the scalar function $`c(\tau ,z,x)c(\tau _0,z_0,x)`$. We write $`A(\tau ;z)A(\tau _0;z_0)=\text{div}M(\tau ,z)\text{grad}`$ as a composition of operators and get the following norm inequality
$$\begin{array}{c}A(\tau ;z)A(\tau _0;z_0)_{L(H^2,L^2)}\hfill \\ \hfill \text{div}_{L(H^1\times H^1,L^2)}M(\tau ,z)_{L(H^1\times H^1)}\text{grad}_{L(H^2,H^1\times H^1)}\\ \hfill \sqrt{2}M(\tau ,z)_{L(H^1\times H^1)}.\end{array}$$
To estimate $`M(\tau ,z)(v_1,v_2)_{H^1\times H^1}`$ it suffices to find an upper bound of
$`(c(z,.,\tau )c(z_0,.,\tau _0))v_{H^1}`$ for $`vH^1`$. We have
$$\begin{array}{c}(c(z,x,\tau )c(z_0,x,\tau _0))v(x)=\hfill \\ \hfill \left(\begin{array}{c}zz_0\\ \tau \tau _0\end{array}\right)_0^1\text{grad}_{(z,\tau )}c(z_0+\sigma (zz_0),x,\tau _0+\sigma (\tau \tau _0))𝑑\sigma v(x),\end{array}$$
which, upon taking the $`H^1`$-norm with respect to $`x`$ and assuming $`\mathrm{max}(|zz_0|,|\tau \tau _0|)1`$, yields
$$\begin{array}{c}(c(z,.,\tau )c(z_0,.,\tau _0))v_{H^1}\mathrm{max}(|zz_0|,|\tau \tau _0|)\hfill \\ \hfill sup(_zc(z^{},.,\tau ^{})v_{H^1}+_\tau c(z^{},.,\tau ^{})_{H^1}),\end{array}$$
where the supremum is taken over $`(z^{},\tau ^{})[z_01,z_0+1]\times [\tau _01,\tau _0+1]`$. Assumption 1 implies that $`_zc`$, $`_\tau c`$ both are continuous functions of $`(z,\tau )`$ valued in $`H^{r+1}(^d)`$, which combined with Fact A gives
$$\begin{array}{c}_zc(z^{},.,\tau ^{})v_{H^1}+_\tau c(z^{},.,\tau ^{})_{H^1}\hfill \\ \hfill C_1v_{H^1}(_zc(z^{},.,\tau ^{})_{H^{r+1}}+_\tau c(z^{},.,\tau ^{})_{H^{r+1}})C_2v_{H^1}\end{array}$$
with positive constants $`C_1`$, $`C_2`$ and for all $`(z^{},\tau ^{})[z_01,z_0+1]\times [\tau _01,\tau _0+1]`$. Combining all estimates we deduce that there is $`C_3>0`$ such that $`|zz_0|+|\tau \tau _0|1`$ implies
$$M(\tau ,z)(v_1,v_2)_{H^1\times H^1}C_3\mathrm{max}(|zz_0|,|\tau \tau _0|)(v_1,v_2)_{H^1\times H^1},$$
which proves the asserted Lipschitz-continuity. ∎
###### Lemma 3.12.
If $`A(\tau ;z)`$ ($`\tau `$, $`z0`$) is defined by (2.8) then Assumption 1 implies the continuity condition (3.19).
###### Proof.
We apply the Kato-Trotter theorem on convergence of semi-groups (cf. \[26, Chapter IX, Section 12, Theorem 1\]). According to this theorem, we obtain
$`T(\tau ,z;\zeta )T(\tau _0,z_0;\zeta )`$ strongly as $`(\tau ,z)(\tau _0,z_0)`$,
uniformly on any compact interval containing $`\zeta `$,
thus (3.19) by uniformity, provided that we show strong continuity of the resolvent map $`(\tau ,z)(\lambda iA(\tau ;z))^1=:R(\lambda ,iA(\tau ;z))`$ for some $`\lambda >0`$.
Fix $`\lambda >0`$ and let $`fL^2(^d)`$ be arbitrary. Define $`u(\tau ,z):=R(\lambda ,iA(\tau ;z))fH^2(^d)`$, so that $`u`$ solves
(3.20)
$$\lambda u(\tau ,z)iA(\tau ;z)u(\tau ,z)=f.$$
Adding the difference $`iA(\tau ;z)u(\tau ,z)iA(\tau _0;z_0)u(\tau ,z)`$ yields
$$(\lambda iA(\tau _0;z_0))u(\tau ,z)=f+i(A(\tau ;z)A(\tau _0;z_0))u(\tau ,z)=:f+iw(\tau ,z).$$
Hence, $`u(\tau ,z)=R(\lambda ,iA(\tau _0;z_0))(f+iw(\tau ,z))`$ and it suffices to prove that $`w(\tau ,z)0`$ in $`L^2(^d)`$ as $`(\tau ,z)(\tau _0,z_0)`$. Applying (C2) from Fact C in Subsection 3.1, we may write
$$\begin{array}{c}w(\tau ,z)=\text{grad}\left(c(z,x,\tau )c(z_0,x,\tau _0)\right)\text{grad}u(\tau ,z)\hfill \\ \hfill +\left(c(z,x,\tau )c(z_0,x,\tau _0)\right)\mathrm{\Delta }u(\tau ,z).\end{array}$$
By Assumption 1, the difference $`c(z,.,\tau )c(z_0,.,\tau _0)`$ tends to $`0`$ in $`H^{r+1}(^d)`$. In view of Fact A this implies $`w(\tau ,z)0`$, if $`\text{grad}u(\tau ,z)`$ as well as $`\mathrm{\Delta }u(\tau ,z)`$ stays bounded. To prove the latter, we take the $`L^2`$-inner product with $`u`$ on both sides of equation (3.20) and obtain
$$\lambda u_{L^2}^2i\underset{j}{}c_{x_j}u|_{x_j}u=f|u.$$
Note that taking real parts here yields the estimate (3.10), which is $`u_{L^2}f_{L^2}/\lambda `$ . If we take absolute values of the imaginary parts, we may use the lower bound $`c(z,x,\tau )c_0`$ and the resolvent estimate to deduce $`_j_{x_j}u(\tau ,z)_{L^2}^2f_{L^2}^2/(\lambda c_0)`$, uniformly in $`(\tau ,z)`$.
Finally, the boundedness of $`\mathrm{\Delta }u(\tau ,z)_{L^2}`$ is revealed in several steps. First, note that (C2) from Fact C applied to (3.20) yields
$$c(z,.,\tau )\mathrm{\Delta }u(\tau ,z)=i(f\lambda u(\tau ,z))\text{grad}c(z,.,\tau )\text{grad}u(\tau ,z).$$
The first term on the right-hand side is bounded in $`L^2`$, uniformly for all $`(\tau ,z)`$, whereas the second term is uniformly bounded in $`H^{r1}`$ by Fact A. Hence $`c\mathrm{\Delta }u(\tau ,z)`$ is a bounded family in $`H^{r1}`$ and, combining Facts A and C, we find that $`\mathrm{\Delta }u(\tau ,z)`$ is uniformly bounded in $`H^{r1}`$ as well. Therefore, $`u(\tau ,z)`$ has a uniform bound in $`H^{r+1}`$-norms. From here we may proceed as in Claim 3 from the proof of Proposition 3.3 (with $`s=0`$). Indeed, the arguments used there preserve uniform boundedness properties throughout, since we have such in $`H^{r+1}`$ already. Thus, $`u(\tau ,z)`$ is uniformly bounded in $`H^2`$, in particular, $`\mathrm{\Delta }u(\tau ,z)`$ is bounded uniformly for all $`(\tau ,z)`$, which completes the proof. ∎
We summarize the preceding results in the announced continuity statement for the evolution system.
###### Theorem 3.13.
Let $`(U(\tau ;z_1,z_2))_{Zz_1z_20}`$ be the evolution system generated by the family of operators $`A(\tau ;z)`$ ($`\tau `$, $`z0`$), defined in (2.8) and satisfying Assumption 1. Then $`(\tau ,z_1,z_2)U(\tau ;z_1,z_2)`$ is strongly continuous $`\times \{(z_1,z_2):Zz_1z_20\}L(L^2(^d))`$.
## 4 Solution of the Cauchy problem
In this section we present our main results: existence and uniqueness of solutions to the Cauchy problem (2.6-2.7) in the frequency domain and to (2.3-2.4) in the time domain.
If $`E`$ is a Banach space, let $`𝒞_b(,E)`$ denote the space of $`E`$-valued continuous bounded functions. Observe that for any $`G𝒞_b(,L^2)𝒮^{}(;L^2)`$ the expression $`U(\tau ;z_1,z_2)G(\tau )`$ is well-defined pointwise for all $`(\tau ,z_1,z_2)`$. Therefore, collecting the results obtained so far we arrive at the following assertion.
###### Proposition 4.1.
If $`v_0𝒞_b(,L^2(^d))`$ and $`g𝒞_b([0,Z]\times ,L^2(^d))`$ then the formula
(4.21)
$$v(z,\tau ):=U(\tau ;z,0)v_0(\tau )+_0^zU(\tau ;z,\rho )g(\rho ,\tau )𝑑\rho $$
defines a mild solution $`v𝒞^1([0,Z],𝒞_b(,L^2(^d)))𝒱`$ to (2.6-2.7). Moreover, when $`v`$ is a strong solution then $`u:=_t^1v`$ is a strong solution of (2.3-2.4) with initial data $`u_0=_t^1v_0`$ and right-hand side $`f=_t^1g`$.
For example, the hypotheses leading to strong solvability are satisfied if $`v_0𝒞_b(,H^2)`$ and $`g𝒞_b([0,Z]\times ,H^2)`$ or $`g𝒞^1([0,Z],𝒞_b(,L^2))`$. Of course, using functions that are bounded and continuous with respect to the frequency variable $`\tau `$ here is just one simple way to ensure that all constructions described above work with all involved objects staying temperate. More generally, it would suffice to consider elements in $`𝒱`$ whose distributional action (with respect to the frequency variable) is given by (weak) integration over a continuous function (times the test function).
To apply formula (4.21) to the original Cauchy problem (2.3-2.4) we only need to state conditions on the data $`u_0`$ and $`f`$ that imply $`_tu_0𝒞_b(,H^2)`$ and $`_tf𝒞_b([0,Z]\times ,H^2)`$ or $`_tf𝒞^1([0,Z],𝒞_b(,L^2))`$. Note that, for example, in our physical application such conditions would be met if the source or force terms are active only for some finite time interval and vanish otherwise. Then by the uniqueness of $`H^2`$-valued solutions, as stated in Remark 3.9, we obtain the following result.
###### Theorem 4.2.
Assume that the right-hand side $`f`$ in equation (2.3) satisfies either $`f𝒞([0,Z],L^1(,H^2(^2)))`$ or $`f𝒞^1([0,Z],L^1(,L^2(^2)))`$.
Then for every $`u_0L^1(,H^2(^2))`$ the Cauchy problem (2.3-2.4) has a unique strong solution $`u𝒞^1([0,Z],𝒮^{}(,L^2(^2)))`$ which is $`H^2`$-valued in the following sense: for all $`z[0,Z]`$ and $`\varphi 𝒮^{}()`$ we have $`u(z),\varphi H^2(^2)`$.
Moreover, $`u`$ is obtained by inverse partial Fourier transform (with respect to $`\tau `$) of $`v𝒞^1([0,Z],𝒞_b(,L^2(^2)))`$ as defined in Equation (4.21), where $`v_0:=_tu_0`$, $`g:=_tf`$, and $`U`$ is the evolution system from Proposition 3.8. In addition, $`v(z,\tau )`$ belongs to $`H^2(^2)`$ for every $`(z,\tau )`$.
#### Inverse analysis of medium regularity:
We conclude with a brief indication of a potential application of Theorem 4.2 to an inverse analysis of medium regularity in wave propagation. Suppose that we are in a model situation where parts (i-ii) and (iv) of Assumption 1 in Section 2 are satisfied and the regularity property (iii) of the medium is in question; assume that we have the a priori knowledge that $`c_l𝒞^1([0,\mathrm{}),L^{\mathrm{}}(^2))`$ for all $`l`$. Let the sources of a seismic experiment be calibrated to produce data in accordance with the hypotheses in Theorem 4.2. If then the measured wave solution $`u`$ (or $`v`$) fails to display the asserted $`H^2`$-regularity then we may conclude that for some $`c_l`$ (and near some depth $`z`$) there is no $`r(0,1)`$ such that the $`H^{r+1}`$-regularity holds; in other words, the lateral regularity of the medium there cannot be better than $`H^1`$ (on the Sobolev scale). If we were in the possession of analogous $`H^s`$-results ($`s>r`$) for the Cauchy problem it would enable us to draw sharper conclusions in such a “inverse regularity analysis”. Note that the exact location of the most singular region need not be known. For the application hence, precise imaging of the singularities is not required prior to the regularity analysis.
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# Crowell’s Derived Group and Twisted Polynomials
## 1 Survey of Crowell’s derived group.
The derived group of a permutation representation unifies in a single theory the concepts of knot group, Alexander matrix and covering space, three concepts that Crowell and Fox believed were central to knot theory (see preface to ). Crowell presented his idea in as a culmination of many years of thought.
The operator notation that we use is more compact than the notation of . It is convenient for our applications, and we hope that it will suggest further applications.
Throughout $`G`$ will denote a multiplicative group acting on the right of a nonempty set $`\mathrm{\Gamma }`$. The action $`\mathrm{\Gamma }\times G\mathrm{\Gamma }`$ is denoted by $`(\gamma ,g)\gamma g`$. It determines in the usual way a representation $`\rho :GS_\mathrm{\Gamma }`$, where $`S_\mathrm{\Gamma }`$ is the group of permutations of $`\mathrm{\Gamma }`$.
###### Definition 1.1.
The *derived group (of the permutation representation $`\rho `$)* is the free group with basis $`\mathrm{\Gamma }\times G=\{g^\gamma gG,\gamma \mathrm{\Gamma }\}`$ modulo the relations $`(gh)^\gamma =g^\gamma h^{\gamma g}`$, for all $`g,hG`$ and $`\gamma \mathrm{\Gamma }`$. We denote it by $`G_\rho `$.
It is helpful to regard the exponent $`\gamma `$ in $`g^\gamma `$ as a coordinate. Before giving examples, we mention basic combinatorial facts about derived groups that are proved in .
###### Lemma 1.2.
(i) $`g^\gamma =1`$ if and only if $`g`$ is trivial in $`G`$.
(ii) $`(g^\gamma )^1=(g^1)^{\gamma g}`$, for any $`gG`$.
(iii) Every nontrivial element $`uG_\rho `$ has a one and only one *normal form*
$$u=g_1^{\gamma _1}g_2^{\gamma _2}\mathrm{}g_n^{\gamma _n},$$
where $`g_k1`$ and $`\gamma _{k+1}\gamma _kg_k`$, for $`1k<n`$.
As Crowell warned, $`G_\rho `$ should not be confused with the commutator subgroup $`[G,G]`$ (also known as the derived group of $`G`$). The reason for the name is a connection between $`G_\rho `$ and the derived module of a homomorphism.
###### Definition 1.3.
Let $`P:G\mathrm{\Gamma }`$ be a group homomorphism. Let $`\mathrm{\Gamma }\times G\mathrm{\Gamma }`$ be the $`G`$-action on $`\mathrm{\Gamma }`$ given by $`(\gamma ,g)\gamma P(g)`$, and let $`\rho :GS_\mathrm{\Gamma }`$ be the associated permutation representation. The *derived module* of $`P`$ is the $`[\mathrm{\Gamma }]`$-module $`G_\rho /[G_\rho ,G_\rho ]`$.
Derived groups can be regarded in two different ways as universal objects in appropriate categories. We briefly outline each approach.
Let $`A`$ be a multiplicative group. We recall that a *crossed product* is a function $`f:\mathrm{\Gamma }\times GA`$ with the property that $`f(\gamma ,g_1g_2)=f(\gamma ,g_1)f(\gamma g_1,g_2)`$, for each $`\gamma \mathrm{\Gamma },g_1,g_2G`$. The function $`:\mathrm{\Gamma }\times GG_\rho `$ determined by $`(\gamma ,g)g^\gamma `$ is an example of a crossed product. Moreover, given any crossed product $`f:\mathrm{\Gamma }\times GA`$ there is a unique group homomorphism $`\overline{f}:G_\rho A`$ such that $`\overline{f}=f.`$ This universal property characterizes $`G_\rho `$ in the usual way.
A second point of view, one that is more topological, involves covering spaces. Let $`B,b_0`$ be a connected space and base point such that $`\pi _1(B,b_0)G`$. Let $`p:EB`$ be the covering space corresponding to the permutation representation $`\rho :GS_\mathrm{\Gamma }`$. The preimage $`p^1(b_0)`$ can be identified with $`\mathrm{\Gamma }`$. Let $`(\gamma ,g)\mathrm{\Gamma }\times G`$ denote the relative homotopy class of paths lying above $`g`$ beginning at $`\gamma `$ and ending at $`\gamma g`$. Then $`\mathrm{\Gamma }\times G`$ is a groupoid under concatenation, and
$$(\gamma _1,g_1)(\gamma _1g_1,g_2)=(\gamma _1,g_1g_2),$$
for $`g_1,g_2G,\gamma \mathrm{\Gamma }`$. The function $`:\mathrm{\Gamma }\times GG_\rho `$ sending $`(\gamma ,g)`$ to $`g^\gamma `$ is a morphism of groupoids. It has the universal property that for any group $`A`$ and groupoid morphism $`f:\mathrm{\Gamma }\times GA`$, there is a unique homomorphism $`\overline{f}:G_\rho A`$ such that $`\overline{f}=f`$.
In a sense, $`G_\rho `$ is the smallest group generated by the monoid of path liftings. It is shown in that $`G_\rho `$ is isomorphic to the free product of $`\pi _1(E)`$ and a free group $`F`$. When $`B`$ is a cell complex, free generators of $`F`$ can be identified with the edges of a maximal tree in the $`1`$-skeleton of the induced cell complex of $`E`$.
Under certain conditions the derived group has a natural structure as a group with operators. The notion, reviewed here, was developed by Noether, who attributed the idea to Krull. A detailed treatment can be found in . An *$`\mathrm{\Omega }`$-group* is a group $`K`$ together with a set $`\mathrm{\Omega }`$ (*operator set*) and a function $`\varphi :K\times \mathrm{\Omega }K,(g,\omega )g^\omega `$, such that for each fixed $`\omega `$, the restricted map $`\varphi _\omega :KK`$ is an endomorphism. If $`K`$ is an $`\mathrm{\Omega }`$-group and $`\mathrm{\Omega }_0`$ is a subset of $`\mathrm{\Omega }`$, then we can regard $`K`$ as an $`\mathrm{\Omega }_0`$-group by restricting the action.
An $`\mathrm{\Omega }`$-group is *finitely generated* (respectively, *finitely presented*) if it is generated (resp. presented) by finitely many $`\mathrm{\Omega }`$-orbits of generators (resp. generators and relators). The commutator subgroup of a knot group is an example of a finitely presented $``$-group. Free $`\mathrm{\Omega }`$-groups are defined in the standard way.
Assume that a group $`G`$ acts on a multiplicative semigroup $`\mathrm{\Gamma }`$ so that $`(\gamma \gamma ^{})g=\gamma (\gamma ^{}g)`$ for all $`\gamma ,\gamma ^{}\mathrm{\Gamma }`$ and $`gG`$. Let $`G_\rho `$ be the associated derived group. We can define a map $`\varphi :G_\rho \times \mathrm{\Gamma }G_\rho `$ on generators of $`G_\rho `$ by $`\varphi _\gamma (g^\gamma ^{})=g^{\gamma \gamma ^{}}`$, extending multiplicatively. One checks that $`G_\rho `$ is then a $`\mathrm{\Gamma }`$-group. Its abelianization is a left $`[\mathrm{\Gamma }]`$-module with $`(n_i\gamma _i)g^\gamma =n_ig^{\gamma _i\gamma }`$ for $`n_i\gamma _i[\mathrm{\Gamma }]`$. If $`G`$ has presentation $`g_1,\mathrm{},g_mr_1,\mathrm{},r_n`$, then $`G_\rho `$ is finitely presented as a $`\mathrm{\Gamma }`$-group by the orbits of $`g_1^1,\mathrm{},g_m^1`$ and $`r_1^1,\mathrm{},r_n^1`$. We will ususally abbreviate $`g_i^1`$ by $`g_i`$. The relators $`r_j^1`$ are easily written in terms of generators $`g_i^\gamma `$ of $`G_\rho `$, as we see in the following example.
###### Example 1.4.
The group $`\pi _k`$ of the figure-eight knot $`k=4_1`$ has presentation
$$\pi _k=a,b\overline{a}ba\overline{b}ab\overline{a}\overline{b}a\overline{b},$$
(1.1)
where $`\overline{}`$ denotes inverse. Let $`P:\pi _kt`$ be the abelianization homomorphism mapping $`a`$ and $`b`$ to $`t`$. The derived group of the associated permutation representation $`\rho :\pi _kS_{}`$ has $``$-group presentation
$$\pi _{k,\rho }=a,b\overline{a}^{t^1}b^{t^1}a\overline{b}ab^t\overline{a}^t\overline{b}a\overline{b}.$$
Here the $``$-group generators $`a`$ and $`b`$ represent orbits $`\{a^{t^i}\}_i`$ and $`\{b^{t^i}\}_i`$ of group generators, while the single $``$-group relator $`\delta (r)`$ represents an orbit $`\{r^{t^i}\}_i`$ of group relators.
The relator $`\delta (r)`$ is a rewrite of $`r^1`$, where $`r`$ is the relator in 1.1. It is easily computed by a variant of Fox calculus, using the axioms:
$$\delta (g)=g,\delta (\overline{g})=\overline{g}^{P(\overline{g})}$$
$$\delta (wg)=\delta (w)g^{P(w)},\delta (w\overline{g})=\delta (w)\overline{g}^{P(w\overline{g})}.$$
Here $`g`$ is any generator while $`w`$ is a word in generators and their inverses. It is easy to check that $`r^{t^i}`$ can be rewritten as $`\delta (r)^{t^i}`$, that is, $`\delta (r)`$ with each generator $`g^{t^j}`$ replaced by $`g^{t^{i+j}}`$.
For this particular representation $`\rho `$, we refer to the group $`\pi _{k,\rho }`$ as the *Alexander group* of the knot, and denote it by $`𝒜_k`$. Abelianizing $`𝒜_k`$ gives the derived module of the abelianization homomorphism. It is the Alexander module of the knot, a $`[t^{\pm 1}]`$-module with matrix presentation described by the usual Fox partial derivatives (see ):
$$\left(\right(\frac{r}{a})^P(\frac{r}{b})^P)=(t^1+3tt^13+t).$$
The superscript indicates that each term of the formal sum is to be replaced by its $`P`$-image.
The $`0`$th characteristic polynomial of the matrix, $`t^23t+1`$ (well defined up to multiplication by $`\pm t^i`$), is the ($`1`$st) Alexander polynomial $`\mathrm{\Delta }_k(t)`$ of the knot. The index shift is a consequence of the fact that we removed a redundant relator from the Wirtinger presentation for $`\pi _k`$ when obtaining 1.1.
Replacing $`\mathrm{\Gamma }=`$ by $`\mathrm{\Gamma }=/rtt^r`$ results in another derived group that is identical to the one above except that exponents $`t^i`$ are taken modulo $`r`$. In this case, $`G_\rho `$ modulo the normal subgroup generated by $`x,x^t,\mathrm{},x^{t^{r2}}`$ produces the fundamental group of the $`r`$-fold cyclic cover of $`𝕊^3k`$. Killing the entire family of generators, $`x,x^t,\mathrm{},x^{t^{r1}}`$ results in the fundamental group of the $`r`$-fold cover $`M_r`$ of $`𝕊^3`$ branched over $`k`$.
## 2 Core group as a derived group.
A presentation for the *core group* $`C_k`$ of a knot can be obtained from any diagram. Generators correspond to arcs, while at each crossing the generators must satisfy a relation indicated in Figure 1. The relation does not depend on the orientation of the arcs.
The core group was discovered independently by A.J. Kelly and M. Wada . It appears in a paper by R.A. Fenn and C.P. Rourke . Wada proved in that $`C_k`$ is isomorphic to the free product of $`\pi _1M_2`$ with an infinite cyclic group (see for an elementary proof).
The *$`\pi `$-orbifold group* $`O_k`$ is another group associated to any knot $`k`$. It is the quotient group
$$\pi _k/x^2,$$
where $`x`$ is any meridian element and $``$ denotes normal closure. Letting $`\tau `$ denote the nontrivial covering transformation of $`M_2`$, the $`\pi `$-orbifold group has a topological interpretation as the group of lifts of $`\tau `$ and the identity transformation to the universal cover of $`M_2`$. Consequently, it fits into a short exact sequence
$$1\pi _1M_2O_k\stackrel{P}{}/21$$
(see for example , page 133). The homomorphism $`P`$ induces a permutation representation $`\rho :O_kS_{/2}`$. Hence the derived group $`O_{k,\rho }`$ is defined.
###### Proposition 2.1.
The derived group $`O_{k,\rho }`$ is isomorphic to the core group $`C_k`$.
###### Proof.
Consider a Wirtinger presentation for the group of $`k`$,
$$\pi =a,b,\mathrm{}r,s,\mathrm{},$$
where as usual generators correspond to arcs of a diagram, and relations correspond to crossings (see Figure 2). Since any two (meridian) generators are conjugate, the $`\pi `$-orbifold group has presentation
$$O_k=a,b,\mathrm{}r,s,\mathrm{},a^2,b^2,\mathrm{}.$$
The homomorphism $`P:O_ktt^2`$ takes each generator to $`t`$. Hence the derived group of the corresponding permutation representation $`\rho `$ has group presentation
$$O_{k,\rho }=a^{t^i},b^{t^i},\mathrm{}\delta (r)^{t^i},\delta (s)^{t^i},\mathrm{},\delta (a^2)^{t^i},\delta (b^2)^{t^i},\mathrm{},$$
where $`i`$ is taken modulo $`2`$, and as before $`\delta (r),\delta (s),\mathrm{}`$ indicate rewritten relators.
Consider a typical Wirtinger relation $`bd=ab`$ corresponding to a positve crossing, as in Figure 2. (The argument or a negative crossing is similar.) The relation determines a pair of relations in the presentation of the derived group. Abbreviating $`a^1,b^1,\mathrm{}`$ by $`a,b,\mathrm{},`$ the relations appear as:
$$bd^t=ab^t,b^td=a^tb.$$
(2.1)
The relators $`a^2,b^2,\mathrm{}`$ also determine pairs of relators:
$$aa^t,a^ta,bb^t=b^tb,\mathrm{}$$
or equivalently,
$$a^t=\overline{a},b^t=b,\mathrm{}.$$
(2.2)
Using the relations 2.2, we can eliminate $`a^t,b^t,\mathrm{},`$ from the presentation for $`O_{k,\rho }`$. The relations 2 become $`b\overline{d}=a\overline{b}`$ (the second relation being redundant). Since this relation is equivalent to $`d=b\overline{a}b`$, we have arrived at a presentation of the core group. Hence $`O_{k,\rho }C_k`$.
## 3 Twisted invariants.
Suppose the permutation representation $`\rho `$ of $`G`$ corresponds to an action of $`G`$ on a left $`R`$-module $`V`$. Then, regarding $`R`$ as a multiplicative semigroup, the derived group $`G_\rho `$ has the structure of an $`R`$-group. The requisite function $`G_\rho \times RG_\rho `$ is defined on generators by $`(g^v,r)g^{rv}`$ and extended multiplicatively. This is well defined on $`G_\rho `$ since $`r`$ takes $`(gh)^v=g^vh^{v\rho (g)}`$ to $`(gh)^{rv}=g^{rv}h^{rv\rho (g)}`$.
Note that in $`G_\rho `$ we do *not* have the relation $`g^vg^w=g^{v+w}`$. If we impose this relation for all $`gG`$ and $`v,wV`$ we obtain a quotient group with an induced $`R`$-group structure. Its abelianization is an $`R`$-module, since we recover the additivity of the $`R`$ action from
$$\begin{array}{cc}\hfill (r_1+r_2)(vg)& =((r_1+r_2)v)g=(r_1v+r_2v)g\hfill \\ & =(r_1v)g+(r_2v)g=r_1(vg)+r_2(vg),\hfill \end{array}$$
where the third equality follows from the imposed relations.
###### Example 3.1.
Twisted Alexander polynomials, introduced for knots in and defined for finitely presented groups in , can be obtained by replacing $`\rho `$ with more general representations. We return to the example above of the figure-eight knot (example 1.4). Let $`P_1:\pi _kt`$ be the abelianization homomorphism considered there, and let $`P_2:\pi _k\mathrm{SL}_2()`$ be the discrete faithful representation
$$a\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),b\left(\begin{array}{cc}1& 0\\ w& 1\end{array}\right),$$
where $`\omega =\frac{1\sqrt{3}}{2}`$ (see Chapter 6 of ). Together, $`P_1`$ and $`P_2`$ determine a homomorphism $`P:\pi _kGL_2([t^{\pm 1}])`$ described by:
$$a\left(\begin{array}{cc}t& t\\ 0& t\end{array}\right)b\left(\begin{array}{cc}t& 0\\ t\omega & t\end{array}\right).$$
Denote the image of $`P`$ by $`\mathrm{\Gamma }`$.
The derived module of $`P:\pi _k\mathrm{\Gamma }`$ is a $`[\mathrm{\Gamma }]`$-module with presentation matrix given by a $`1\times 2`$ Jacobian:
$$\left(\right(\frac{r}{a})^P(\frac{r}{b})^P).$$
Each entry is a formal sum of matrices. Add its terms to form a single $`2\times 2`$-matrix. The result is:
$$\left(\left(\begin{array}{cc}\overline{\omega }t+1+\overline{\omega }\frac{1}{t}& t2\omega +\frac{1}{t}\\ i\sqrt{3}\overline{\omega }(t1)& \overline{\omega }t+3\frac{1}{t}\end{array}\right)\left(\begin{array}{cc}2+\frac{1}{t}& \omega t+1+\omega \frac{1}{t}\\ \overline{\omega }t1& 2t3+\frac{1}{t}\end{array}\right)\right).$$
(3.1)
Wada showed that the following procedure produces a polynomial invariant of knots. Remove the inner parentheses to form a $`2\times 4`$ *twisted Alexander matrix* $`A_\rho `$. Delete the first two columns (corresponding to the generator $`a`$), take the determinant, and divide by $`\mathrm{Det}(tIP(a))`$. The resulting polynomial is Wada’s invariant,
$$W_{k,\rho }=\frac{(t1)^2(t^24t+1)}{(t1)^2}=t^24t+1.$$
The result is the same if we reverse the roles of $`a`$ by $`b`$. (See .)
Adding terms and removing inner parentheses to obtain $`B`$ appears ad hoc. The following alternative approach, based on the derived group of an action on an $`R`$-module, gives $`B`$ more naturally.
The group $`\pi _k`$ acts on the free $`[t^{\pm 1}]`$-module $`V=[t^{\pm 1}][t^{\pm 1}]`$ via $`P:\pi _k\mathrm{GL}_2([t^{\pm 1}])`$: for any $`vV`$ and $`g\pi _k`$, we define $`vg`$ to be $`vP(g)`$. This gives a permutation representation $`\rho _V`$ of $`\pi _k`$ on $`V`$. By remarks above, the derived group $`\pi _{\pi _k,\rho _V}`$ is a $`[t^{\pm 1}]`$-group, and the abelianization becomes a left $`[t^{\pm 1}]`$-module if we impose the extra relations $`(v_1+v_2)g=v_1g+v_2g`$ for every $`v_1,v_2V`$ and $`g\pi _k`$. It is generated by $`e_1a,e_2a,e_1b,e_2b`$, where $`e_1=(1,0),e_2=(0,1)`$ is the standard basis for $`V`$, and one easily checks that $`A_\rho `$ is a presentation matrix for it.
###### Remark 3.2.
(i) Twisted Alexander polynomials of any $`d`$-component link $`\mathrm{}`$ can be found in a similar way, using the abelianization $`P_1:\pi _{\mathrm{}}^d`$.
(ii) The greatest common divisor of the six $`2\times 2`$ minors of $`A`$ is an invariant of $`k`$ and the representation $`P`$. It is the $`1`$st twisted Alexander polynomial $`\mathrm{\Delta }_1`$, defined by P. Kirk and C. Livingston . The greatest common divisor of the $`1\times 1`$ minors (that is, the entries of $`A`$) is the $`0`$th twisted polynomial $`\mathrm{\Delta }_0`$. The relationship $`W_{k,\rho }=\mathrm{\Delta }_1/\mathrm{\Delta }_0`$ holds generally.
The reader is cautioned that an invariant does not result by simply taking the greatest common divisor of the two $`2\times 2`$ minors corresponding to the entries of 3.1, contrary to Proposition 1.3 of . Such a quantity is not preserved when the presentation of $`\pi _k`$ is altered by a Tietze move, in contradiction to the assertion on page 213 of . Indeed, in the above example, the greatest common divisor of the two minors is $`(t1)^2(t^24t1)`$. However, if in the presentation of $`\pi _k`$ we replace $`a`$ by the new generator $`c=\overline{a}b`$, accomplished by a sequence of Tietze moves, then the polynomial $`t^24t+1`$ is obtained.
The Alexander group $`𝒜_k`$ is the free product of the commutator subgroup $`[\pi _k,\pi _k]`$ with a countable-rank free group $`F`$. One can choose $`F`$ to be the free group generated by any $``$-group family $`\{x^{t^i}\}`$ of group generators, where $`x`$ is a Wirtinger generator (cf. Section 1; also see ). Notice, in particular, that $`𝒜_k/x`$ is isomorphic to $`[\pi _k,\pi _k]`$. Here $`x`$ denotes the normal $``$-subgroup generated by $`x`$, that is, the smallest normal subgroup containing the orbit $`\{x^{t^i}\}`$.
For any knot $`k`$, the commutator subgroup $`[\pi _k,\pi _k]`$ is a finitely generated $``$-group. As an ordinary group, it need not be finitely generated. A theorem of J. Stallings states that if it is, then $`k`$ is fibered, and $`[\pi _k,\pi _k]`$ is free of rank $`2g_k`$, where $`g_k`$ denotes the genus of $`k`$. The converse is also true: if $`k`$ is fibered, then $`[\pi _k,\pi _k]`$ is a free group of rank $`2g_k`$ .
A consequence is that if $`k`$ is fibered, then its Alexander polynomial is monic with degree equal to twice the genus of $`k`$. H. Goda, T. Kitano and T. Morifuji extended this classical result by showing that if $`k`$ is fibered and $`\rho `$ is obtained as in Example 3.1 from the abelianization homomorphism $`P_1`$ and homomorphism $`P_2:\pi _kSL_{2n}(𝔽),𝔽`$ a field, then Wada’s invariant $`W_{k,\rho }(t)`$ is a rational function of monic polynomials. A similar result was obtained by J.C. Cha . We prove a companion result for Alexander groups, generalizing the theorem of Neuwirth and Stallings.
Consider a knot $`k𝕊^3`$ with abelianization homomorphism $`P_1:\pi _kt`$ and $`P_2:\pi _k\mathrm{SL}_n(R)`$ any homomorphism. We assume that $`R`$ is a unique factorization domain. Let $`P:\pi _k\mathrm{GL}_n(R[t^{\pm 1}])`$ be the product homomorphism, as above, with image $`\mathrm{\Gamma }`$. Let $`\rho :\pi _kS_\mathrm{\Gamma }`$ be the associated permutation representation. We denote by $`(\pi )`$ the collection of permutation representations that arise this way.
The derived group $`\pi _{k,\rho }`$ is the *$`\rho `$-twisted Alexander group* of $`k`$, and we denote it by $`𝒜_{k,\rho }`$. When $`P_2`$ is trivial, $`\pi _{k,\rho }`$ reduces to the Alexander group of $`k`$. The twisted Alexander group is a $`\mathrm{\Gamma }`$-group. The image $`\mathrm{\Gamma }_2`$ of $`P_2`$ is a subgroup of $`\mathrm{\Gamma }`$. We can also regard $`𝒜_{k,\rho }`$ as a $`\mathrm{\Gamma }_2`$-group, letting only elements of $`\mathrm{\Gamma }_2`$ act on it. Throughout the remainder of the section, $`x`$ will denote a Wirtinger generator of $`\pi _k`$ corresponding to a meridian of $`k`$.
###### Theorem 3.3.
If a knot $`k`$ is fibered, then, for any representation $`\rho (\pi )`$, the twisted Alexander group $`𝒜_{k,\rho }`$ modulo $`x`$ is a free $`\mathrm{\Gamma }_2`$-group of rank $`2g_k`$. Conversely, if $`𝒜_{k,\rho }`$ modulo $`x`$ is a finitely generated $`\mathrm{\Gamma }_2`$-group for some $`\rho (\pi )`$, then $`k`$ is fibered.
###### Proof.
Assume that $`k`$ is fibered. Then $`\pi _k`$ is an HNN extension of a free group $`F=F(a_1,\mathrm{},a_{2g})`$, where $`g=g_k`$ is the genus of $`k`$, and it has a presentation of the form
$$x,a_1,\mathrm{},a_{2g}xa_1x^1=\varphi (a_1),\mathrm{},xa_{2g}x^1=\varphi (a_{2g})$$
$`(2.3)`$
for some automorphism $`\varphi `$ of $`F`$ (see , for example). The twisted Alexander group $`𝒜_{k,\rho }`$ has a $`\mathrm{\Gamma }`$-group presentation
$$x,a_1,\mathrm{},a_{2g}xa_1^{tP_2(x)}\overline{x}^{P_2(xa_1\overline{x})}=\stackrel{~}{\varphi (a_1)},\mathrm{},xa_1^{tP_2(x)}\overline{x}^{P_2(xa_{2g}\overline{x})}=\stackrel{~}{\varphi (a_{2g})},$$
where $`\stackrel{~}{}`$ denotes the $`\delta `$-rewrite, explained above. The quotient $`\mathrm{\Gamma }`$-group $`𝒜_{k,\rho }/x`$ has presentation
$$a_1,\mathrm{},a_{2g}a_1^{tP_2(x)}=\stackrel{~}{\varphi (a_1)},\mathrm{},a_{2g}^{tP_2(x)}=\stackrel{~}{\varphi (a_{2g})}.$$
Since $`P_2(x)`$ is invertible, each $`a_i^t`$ can be expressed as $`\stackrel{~}{\varphi (a_i)}^{P_2(x)^1}`$. Generally, the relators can be used one at a time to express each $`a_i^{t^j},j0`$, as a word in the $`\mathrm{\Gamma }_2`$-orbits of $`a_1,\mathrm{},a_{2g}`$. Hence $`𝒜_{k,\rho }/x`$ is isomorphic to the free $`\mathrm{\Gamma }_2`$-group generated by $`a_1,\mathrm{},a_{2g}`$.
Conversely, assume that $`𝒜_{k,\rho }/x`$ is a finitely generated $`\mathrm{\Gamma }_2`$-group for some $`\rho (\pi )`$. The knot group $`\pi _k`$ is generated by $`x,a_1,\mathrm{},a_m`$, where $`a_1,\mathrm{},a_m[\pi _k,\pi _k]`$. Then $`𝒜_{k,\rho }/x`$ is generated by elements of the form $`a_i^{t^jM}`$, where $`j`$ and $`M`$ ranges over matrices in $`\mathrm{\Gamma }_2`$, while $`[\pi _k,\pi _k]`$ is generated by the elements $`a_i^{t^j}`$. The mapping $`a_i^{t^jM}a_i^{t^j}`$ determines a surjection $`f:𝒜_{k,\rho }/x[\pi _k,\pi _k]`$. By assumption, we can find generators $`a_i^{t^jM}`$ of $`𝒜_{k,\rho }/x`$ such that the values of $`j`$ are bounded. Consequently, the image of $`f`$ is finitely generated. Hence $`k`$ is fibered. ∎
## 4 Virtual knots and links.
Like their classical counterparts, virtual knots and links are equivalence classes of diagrams, the equivalence relation generated by Reidemeister moves. However, in the case of virtual knots and links, we also allow virtual crossings (indicated by a small circle about the crossing), and the set of Reidemeister moves is suitably generalized. The notion is due to Kauffman, and the reader is referred to or for details. A result of M. Goussarov, M. Polyak and O. Viro assures us that if two classical diagrams—that is, diagrams without virtual crossings—are equivalent by generalized Reidemeister moves, then they are equivalent by the usual classical Reidemeister moves. In this sense, virtual knot theory extends the classical theory.
Many invariants of classical knots and links are also defined in the virtual category. The knot group $`\pi _{\mathrm{}}`$ is such an invariant. Given any diagram of a virtual knot or link $`\mathrm{}`$, a Wirtinger presentation is obtained by assigning a generator to each arc. (An arc is a maximal connected component of the diagram containing no classical under-crossing.) A relation is associated to each classical crossing in the usual way.
Let $`\mathrm{}=\mathrm{}_1\mathrm{}\mathrm{}_d`$ be an oriented virtual link of $`d`$ components with group $`\pi _{\mathrm{}}`$. Regard $`^d`$ as a multipicative group freely generated by $`u_1,\mathrm{},u_d`$. Let $`P:\pi _{\mathrm{}}^d`$ be the abelianization homomorphism mapping the class of the $`i`$th oriented meridian to $`u_i`$. Let $`\rho `$ be the associated permutation representation $`\rho :\pi _{\mathrm{}}S_^d`$. The derived group is the Alexander group $`𝒜_{\mathrm{}}`$ introduced in .
Let $`D`$ be a diagram of $`\mathrm{}`$. An *edge* of $`D`$ is a maximal segment of an arc going from one classical crossing to the next. (Thus at each crossing, the arc passing over is broken into distinct edges.) The *extended group* $`\stackrel{~}{\pi }_{\mathrm{}}`$ has generators $`a,b,c,\mathrm{}`$ corresponding to edges together with an additional generator $`x`$ not associated with any edge. Relations come in pairs, corresponding to classical crossings: $`ab=cd,\overline{x}bx=c`$, if the crossing is positive, and $`ab=cd,\overline{x}dx=a`$, if the crossing is negative (see Figure 3). It is a straightforward matter to check that the group so defined is unchanged if a generalized Reidemeister move is applied to the diagram, and we leave this to the reader.
The extended Alexander group $`\stackrel{~}{𝒜}_{\mathrm{}}`$ of the link $`\mathrm{}`$ was defined in to be the $`^{d+1}`$-group with generators $`a,b,c,\mathrm{}`$, corresponding to edges as above, and relations at each classical crossing: $`ab^{u_i}=cd^{u_j},c^v=b`$, if the crossing is positive; $`ab^{u_i}=cd^{u_j},a^v=d`$, if the crossing is negative. We assume that generators $`a,d`$ correspond to edges on the $`i`$th component of the link while $`b,c`$ correspond to edges on the $`j`$th component. One regards each of $`a,b,c,\mathrm{}`$ as families of generators indexed by monomials in $`u_1,\mathrm{},u_d,v`$, generators for the group $`^{d+1}`$ written multiplicatively. Similarly, each relation $`r`$ of the $`^{d+1}`$-group presentation is a family of group relations indexed by monomials.
Let $`P:\stackrel{~}{\pi }_{\mathrm{}}^{d+1}`$ be the abelianization homomorphism mapping generators corresponding to edges on the $`i`$th component to $`u_i`$, and mapping $`x`$ to $`v`$, and let $`\rho :\stackrel{~}{\pi }_{\mathrm{}}S_{^{d+1}}`$ be the associated permutation representation. Consider the derived group $`\stackrel{~}{\pi }_{\mathrm{},\rho }`$. It is generated as a $`^{d+1}`$-group by $`a,b,c,\mathrm{}`$ and $`x`$. Modulo $`x`$, the $`\delta `$-rewrite of the relations for $`\stackrel{~}{\pi }_{\mathrm{},\rho }`$ are easily seen to be those of the extended Alexander group. Indeed, the $`\delta `$-rewrite of $`ab=cd`$ is $`ab^{u_i}=cd^{u_j}`$, while the rewrite of $`xc\overline{x}=b`$ is $`xc^v\overline{x}^{u_j}=b`$. Killing $`x`$, that is, killing its $`^{d+1}`$-orbit, yields the relations for $`\stackrel{~}{𝒜}_{\mathrm{}}`$ at a positive classical crossing. A similar argument applies to negative crossings. Summarizing:
###### Proposition 4.1.
Let $`\mathrm{}=\mathrm{}_1\mathrm{}\mathrm{}_d`$ be an oriented virtual link of $`d`$ components, and let $`\stackrel{~}{\pi }_{\mathrm{}}`$ be its extended group. Let $`\stackrel{~}{\pi }_{\mathrm{},\rho }`$ be the derived group of the permutation representation associated to the abelianization homomorphism $`P:\stackrel{~}{\pi }_{\mathrm{}}^{d+1}`$. Then as $`^{d+1}`$-groups, the extended Alexander group $`\stackrel{~}{𝒜}_{\mathrm{}}`$ is isomorphic to $`\stackrel{~}{\pi }_{\mathrm{},\rho }/x`$, the derived group modulo the normal $`^{d+1}`$-subgroup generated by $`x`$.
Proposition 4.2 yields new invariants for virtual links.
###### Proposition 4.2.
(i) Let $`\mathrm{}`$ be an oriented virtual link. Then
$$1x\stackrel{~}{\pi }_{\mathrm{}}\stackrel{p}{}\pi _{\mathrm{}}1$$
is an exact sequence of groups, where $``$ denotes normal closure and $`p`$ is the quotient map sending $`x1`$.
(ii) If $`\mathrm{}`$ is a classical link, then
$$1𝒜_{\mathrm{}}\stackrel{~}{\pi }_{\mathrm{}}\stackrel{\chi }{}t1$$
is an exact sequence of groups, where $`\chi `$ is the homomorphism mapping $`xt`$ and $`a,b,c,\mathrm{}1`$.
###### Proof.
Killing $`x`$ converts the relations for $`\stackrel{~}{\pi }_{\mathrm{}}`$ into the Wirtinger relations for $`\pi _{\mathrm{}}`$. Hence (i) is proved.
In order to prove (ii), assume that $`D`$ is a classical diagram for $`\mathrm{}`$; that is, a diagram without virtual crossings. By the Seifert smoothing algorithm, we can label the edges of $`D`$ with integers $`\nu (a),\nu (b),\nu (c),\mathrm{}`$ such that at any crossing such as in Figure 1, whether positive or negative, edges corresponding to $`a,c`$ receive the same label, say $`\nu `$, while edges corresponding to $`b,d`$ receive $`\nu +1`$. Such an assignment of integers will be called an *Alexander numbering*. (See for details.)
After replacing each generator $`a,b,c,\mathrm{}`$ with $`x^{\nu (a)}ax^{\nu (a)},`$ $`x^{\nu (b)}bx^{\nu (b)},`$ $`x^{\nu (c)}cx^{\nu (c)}\mathrm{}`$, the relations at a positive crossing become $`axb\overline{x}=cxd\overline{x}`$ and $`c=b`$. At a negative crossing, the relations become $`axb\overline{x}=cxd\overline{x}`$ and $`a=d`$. Notice in particular that $`\stackrel{~}{\pi }_{\mathrm{}}`$ is generated by $`x`$ and symbols corresponding to the arcs rather than the edges of $`D`$. Moreover, the form of the relations allows us to describe the kernel of $`\chi `$ using the Reidemeister-Schreier method: it is generated as a $`t`$-group by symbols corresponding to the arcs of $`D`$ together with relations $`ab^t=bd^t`$ at a positive crossing and $`ab^t=ca^t`$ at a negative crossing. This is also a description of the Alexander group $`𝒜_{\mathrm{}}`$ of the link. ∎
Proposition 4.2 motivates the following.
###### Definition 4.3.
A virtual link is *almost classical* if some diagram for the link admits an Alexander numbering.
###### Remark 4.4.
(i) Reassuringly, classical implies almost classical. However, we see in Example 5.1 that the converse does not hold.
(ii) The conclusion of Proposition 4.2(ii) holds for almost classical links. The proof is similar.
###### Example 4.5.
The virtual knot $`k`$ in Figure 4 appears in . Its extended group has generators $`a,\mathrm{},h`$. Using the relations corresponding to the classical crossings, one finds after some simplification that $`\stackrel{~}{\pi }_k`$ has group presentation
$$x,a,daxd\overline{a}\overline{x}\overline{a}xa^2xa=dxaxd,dxaxd=xadxa.$$
The kernel of the homomorphism $`\chi :\stackrel{~}{\pi }_kt`$ mapping $`a,d1`$ and $`xt`$ is a $`t`$-group. The following $``$-group presentation can be found using the Reidemeister-Schreier method:
$$\mathrm{ker}(\chi )=a,dad^t\overline{a}^t\overline{a}(a^2)^ta^{t^2}=da^td^{t^2},da^td^{t^2}=a^td^ta^{t^2}$$
(4.1)
If $`k`$ were almost classical, then the group presented by 4.1 would be isomorphic to the Alexander group $`𝒜_k`$ of $`k`$. We compute a presentation of $`𝒜_k`$ directly from Figure 4, and find
$$𝒜_k=c,dcd^tc^{t^2}=dc^td^{t^2},cd^tc^{t^2}=dc^td^{t^2}$$
(4.2)
Each presentation (4.1), (4.2) determines a $`[t^{\pm 1}]`$-module with $`2\times 2`$-relation matrix. Both matrices have trivial determinant, which is the $`0`$th characteristic polynomial. However, the $`1`$st characteristic polynomial corresponding to (4.1) is $`1`$ while that corresponding to (4.2) is $`t^2t+1`$. Since these are unequal, the two groups are not isomorphic, and hence $`k`$ is not almost classical.
###### Remark 4.6.
It follows that the knot in Example 4.5 is not classical. Heather Dye has shown us an alternative proof using minimal surface representations, as in .
The proof of Theorem 4.1 of shows that if $`\mathrm{}`$ is an almost classical oriented virtual link of $`d`$ components, then a presentation for the Alexander group $`𝒜_{\mathrm{}}`$ obtained from any diagram of the link can be converted into a presentation for the extended Alexander group $`\stackrel{~}{𝒜}_{\mathrm{}}`$ by replacing each occurrence of $`u_i`$ by $`u_iv,1id`$. As a consequence, each virtual Alexander polynomial $`\mathrm{\Delta }_i(u_1,\mathrm{},u_d,v)`$ is a polynomial in the $`d`$ variables $`u_1v,\mathrm{},u_dv`$.
The arguments are similar for twisted groups and virtual Alexander polynomials. Let $`P_1:\stackrel{~}{\pi }_{\mathrm{}}^{d+1}`$ be the abelianization representation, and let $`P_2:\stackrel{~}{\pi }_{\mathrm{}}\mathrm{GL}_n()`$ be a linear representation. Together, $`P_1`$ and $`P_2`$ determine a representation $`P:\stackrel{~}{\pi }_{\mathrm{}}\mathrm{GL}_n([u_1^{\pm 1},\mathrm{},u_d^{\pm 1},v^{\pm 1}])`$. As above, we denote the image of $`P`$ by $`\mathrm{\Gamma }`$ and the induced permutation representation of $`\pi _{\mathrm{}}`$ on $`\mathrm{\Gamma }`$ by $`\rho `$. A twisted Alexander matrix $`A_\rho `$ and associated polynomials $`\mathrm{\Delta }_{\mathrm{},\rho ,i}=\mathrm{\Delta }_{\mathrm{},\rho ,i}(u_1,\mathrm{},u_d,v)`$ can now be defined as in Example 3.1.
If the diagram for $`\mathrm{}`$ has $`N`$ classical crossings, $`A_\rho `$ is an $`Nn\times Nn`$-matrix. The greatest common divisor of the determinants of all $`(Nni)\times (Nni)`$ minors produces the *$`i`$th twisted virtual Alexander polynomial* $`\mathrm{\Delta }_{\mathrm{},\rho ,i}`$. As is the case with $`\mathrm{\Delta }_{\mathrm{},0}`$, the twisted polynomial $`\mathrm{\Delta }_{\mathrm{},\rho ,0}`$ must vanish if $`\mathrm{}`$ is classical, and the argument is similar (see ).
###### Theorem 4.7.
If $`\mathrm{}`$ is almost classical, then a presentation for the twisted Alexander group $`𝒜_{\mathrm{},\rho }`$ obtained from any diagram of $`\mathrm{}`$ can be converted into a presentation for the extended Alexander group $`\stackrel{~}{𝒜}_{\mathrm{},\rho }`$ by replacing each occurrence of $`u_i`$ by $`u_iv,1id`$.
###### Corollary 4.8.
Assume that $`\mathrm{}`$ is almost classical. Then each twisted virtual Alexander polynomial $`\mathrm{\Delta }_{\mathrm{},\rho ,i}(u_1,\mathrm{},u_d,v)`$ is a polynomial in the $`d`$ variables $`u_1v,\mathrm{},u_dv`$.
## 5 Remarks about almost classicality.
There are different ways to define what is means for a virtual knot to be almost classical. Our definition is motivated by group-theoretical concerns. The following two examples are intended to convince the reader that the definition is subtle and worthy of further study.
###### Example 5.1.
*Almost classical does not imply classical.* Consider the diagram $`D`$ for a virtual knot $`k`$ in Figure 5. An Alexander numbering for $`D`$ is shown, and hence $`k`$ is almost classical. The $`1`$st virtual Alexander polynomial $`\mathrm{\Delta }_1(u,v)`$ is easily computed and seen to be $`2uv`$. If $`k`$ were classical, then its Alexander polynomial would be $`2t`$. However, $`2t`$ is not reciprocal, and hence it cannot be the Alexander polynomial of any classical knot. From this we conclude that $`k`$ is not classical.
One can weaken the definition of almost classical by asking only for a mod-$`2`$ Alexander numbering, that is, replacing $``$ by $`/2`$ in the definition of Alexander numbering.
###### Example 5.2.
*Existence of a mod-2 Alexander numbering does not imply almost classical.* Figure 6 is a diagram of a virtual knot $`k`$ with a mod-$`2`$ Alexander numbering. The 0th virtual Alexander polynomial of $`k`$ is
$$\mathrm{\Delta }_{k,0}(u,v)=1u^2uvv^2+u^3v+uv^3+u^2v^2u^3v^3.$$
Since it is not a polynomial in $`uv`$, the knot $`k`$ admits no diagram with an Alexander numbering.
The knot in Example 5.2 was introduced by Kauffman in . It is obtained from a standard diagram of a trefoil by “virtualizing a single crossing,” switching the crossing and flanking it by virtual crossings. As observed in , it is a nontrivial virtual knot with unit Jones polynomial. Such knots were further studied in .
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# Coupled-channel model for charmonium levels and an option for 𝑋(3872)
## I Introduction
The charmonium spectroscopy has again become a very interesting field. On one hand, the $`2^1S_0`$ $`\eta _c^{}`$ state was discovered by Belle Belle and confirmed by BABAR Babar and CLEO Cleo1 , and the $`1^1P_1`$ $`h_c`$ was observed in Fermilab FNAL and by CLEO Cleo2 . The masses of these long missing states are in perfect agreement with the predictions of quark model. At the same time, a new state $`X(3872)`$ was found by Belle BelleX and CDF CDF .
This discovery has attracted much attention. As the state is just at the $`D\overline{D}^{}`$ threshold, it was immediately suggested TornqvistX that it might be a $`D\overline{D}^{}`$ molecule bound by pion exchange (”deuson”), considered long ago in TornqvistD and, much earlier, in VO and RGG . This requires $`1^{++}`$ quantum numbers, and this assignment seems to be favoured by the data abe ; bauer ; olsen . The discovery channel is $`\pi ^+\pi ^{}J/\psi `$ with dipion most probably originating from the $`\rho `$. Together with the observation of $`X`$ in the $`\omega J/\psi `$ channel Belleomega this opens fascinating possibilities for strong isospin violation, which is also along the lines of the deuson model.
Other options for $`X(3872)`$ are under discussion in the literature, see e.g. CP ; BG ; PS ; Bugg . The most obvious possibility of $`X`$ being a $`c\overline{c}`$ state seems to be ruled out by its mass: the state is too high to be a $`1D`$ charmonium, and too low to be a $`2P`$ one BG . This assumes that we do know the spectrum of higher charmonia, namely, the fine splittings and the role of coupling to $`D`$-meson pairs. It is the latter issue which is addressed in the present paper.
The mechanism of open-flavour strong decay is not well-understood. The simplest model for light-quark pair creation is the so-called $`{}_{}{}^{3}P_{0}^{}`$ model, suggested many years ago Micu . It assumes that the pair is created with vacuum ($`{}_{}{}^{3}P_{0}^{}`$) quantum numbers uniformly in space. The application of this model has a long history Orsay ; BO ; KI . Systematic studies Barnes1 ; Barnes2 of the decays of light and strange quarkonia show that with a $`{}_{}{}^{3}P_{0}^{}`$ -type amplitude calculated widths agree with data to within $`2540\%`$. Recently the charmonia decays Barnes3 and decays of $`D`$\- and $`D_s`$-mesons CS were considered in the framework of the $`{}_{}{}^{3}P_{0}^{}`$ model.
There exist also microscopic models of strong decays, which relate the pair-creation interaction to the interaction responsible for the formation of the spectrum, by constructing the current-current interaction due to confining force and one-gluon exchange. Among these is the Cornell model Eichten which assumes that confinement has Lorentz vector nature. The model ABarnes assumes that the confining interaction is the scalar one, while one-gluon-exchange is, of course, Lorentz vector. Possible mechanisms of strong decays were studied in the framework of Field Correlator Method (FCM) Simonov , and an effective $`{}_{}{}^{3}P_{0}^{}`$ operator for open-flavour decay has emerged from this study, with the strength computed in terms of FCM parameters (string tension and gluonic correlation length).
Most of the above-mentioned papers are devoted to computing the widths, and only a few consider the effects of virtual hadronic loops on the spectra. The Cornell model Eichten has presented a detailed analysis of charmonia with coupling to $`D`$-mesons taken into account. The recent update nEichten of the Cornell model has presented splittings caused by coupling to mesonic channels for $`1D`$ and $`2P`$ $`c\overline{c}`$ levels, confirming the previous result: $`X(3872)`$ is well above the range of $`1D`$ levels and well below the range of $`2P`$ ones. The paper Barnes0 has reported first results for hadronic shifts of lower charmonia due to mixing with $`D`$-meson pairs, calculated within the $`{}_{}{}^{3}P_{0}^{}`$ model. The shifts appear to be alarmingly large.
Meanwhile, phenomenological coupled-channel models like port ; Tornqvist ; Markushin ; Pennington accumulate experience on the possibilities to generate nontrivial effects due to the coupling to hadronic channels. As a recent example one should mention the analyses Ds of new $`D_{sJ}`$ states with masses considerably lower than quark model predictions, and coupling to mesonic channels being responsible for these anomalously low masses. It is interesting to note that the coupled-channel calculations performed in the framework of chiral Lagrangian approach ST has arrived at the same conclusions.
In this paper the coupled-channel model for charmonia levels is presented, based on the nonrelativistic quark model for $`c\overline{c}`$ spectrum and $`{}_{}{}^{3}P_{0}^{}`$-type model for pair-creation. In Section II the dynamics of coupled channels is briefly outlined. Section III introduces the quark model. Sections IV and V contain the results which are discussed in Section VI. The paper ends with a short summary.
## II Dynamics of coupled channels
The details of coupled-channel model can be found e.g. in Eichten , Ker . Here I review the essentials.
In what follows the simplest version of coupled-channel model is employed. Namely, it is assumed that the hadronic state is represented as
$$|\mathrm{\Psi }=\left(\genfrac{}{}{0pt}{}{\underset{\alpha }{}c_\alpha |\psi _\alpha }{_i\chi _i|M_1(i)M_2(i)}\right),$$
(1)
where the index $`\alpha `$ labels bare confined states $`|\psi _\alpha `$ with the probability amplitude $`c_\alpha `$, and $`\chi _i`$ is the wave function in the $`i`$-th two-meson channel $`|M_1(i)M_2(i)`$. The wave function $`|\mathrm{\Psi }`$ obeys the equation
$$\widehat{}|\mathrm{\Psi }=M|\mathrm{\Psi },\widehat{}=\left(\begin{array}{cc}\widehat{H}_c& \widehat{V}\\ \widehat{V}& \widehat{H}_{M_1M_2}\end{array}\right),$$
(2)
where $`\widehat{H}_c`$ defines the discrete spectrum of bare states, with $`\widehat{H}_c|\psi _\alpha =M_\alpha |\psi _\alpha `$. The part $`\widehat{H}_{M_1M_2}`$ includes only the free-meson Hamiltonian, so that the direct meson-meson interaction (e.g., due to $`t`$\- or $`u`$-channel exchange forces) is neglected. The term $`\widehat{V}`$ is responsible for dressing of the bare states.
Consider one bare state $`|\psi _0`$ (the generalization to multi-level case is straightforward). The interaction part is given by the transition form factor $`f_i(𝐩)`$,
$$\psi _0|\widehat{V}|M_{i1}M_{i2}=f_i(𝐩_i),$$
(3)
where $`𝐩_i`$ is the relative momentum in $`i`$-th mesonic channel. Then (2) leads to the system of coupled equations for $`c_0(M)`$ and $`\chi _{i,M}(𝐩_i)`$:
$$\{\begin{array}{c}c_0(M)M_0+\underset{i}{}f_i(𝐩)\chi _{i,M}(𝐩)d^3p=Mc_0(M),\\ \left(m_{i1}+m_{i2}+\frac{p^2}{2\mu _i}\right)\chi _{i,M}(𝐩_i)+c_0(M)f_i(𝐩_i)=M\chi _{i,M}(𝐩_i).\end{array}$$
(4)
Here $`\mu _i=\frac{m_{i1}m_{i2}}{m_{i1}+m_{i2}}`$ is the reduced mass in the system of mesons with the masses $`m_{i1}`$ and $`m_{i2}`$, and $`M_0`$ is the mass of the bare state. In what follows the formalism will be applied to the system of charmed mesons, so the nonrelativistic kinematics is employed in (4). The fully relativistic version of coupled-channel model is presented in Tornqvist .
With the help of (4) one easily calculates the $`t`$-matrix in the mesonic system:
$$t_{ik}(𝐩_i,𝐩_k^{},M)=\frac{f_i(𝐩_i)f_k(𝐩_k^{})}{MM_0+g(M)},g(M)=\underset{i}{}g_i(M),$$
$$g_i(M)=\frac{f_i(𝐩)f_i(𝐩)}{\frac{p^2}{2\mu _i}E_ii0}d^3p,$$
(5)
where $`E_i=Mm_{i1}m_{i2}`$. The quantity $`g(M)`$ is often called the hadronic shift of the bare state $`|\psi _0`$, as the masses of physical states are defined, in accordance with eq.(5), from the equation
$$MM_0+g(M)=0.$$
(6)
Let the eq.(6) have the solution $`M_B`$ with $`M_B`$ smaller than the lowest mesonic threshold, so there is a bound state with the wave function
$$|\mathrm{\Psi }_B=\left(\genfrac{}{}{0pt}{}{\mathrm{cos}\theta |\psi _0}{\mathrm{sin}\theta _i\chi _{iB}(𝐩_i)}\right),\mathrm{\Psi }_B|\mathrm{\Psi }_B=1,\mathrm{cos}\theta =\psi _0|\mathrm{\Psi }_B.$$
(7)
Here $`_i\chi _{iB}(𝐩_i)`$ is normalized to unity, and $`\mathrm{cos}\theta `$ defines the admixture of the bare state $`|\psi _0`$ in the physical state $`|\mathrm{\Psi }_B`$. The explicit expression for this admixture reads
$$Z\mathrm{cos}^2\theta =\left(1+\underset{i}{}\frac{f_i(𝐩)f_i^{}(𝐩^{})d^3p}{(\frac{p^2}{2\mu _i}+ϵ_i)^2}\right)^1=\left(1+\frac{g(M)}{M}|_{M=M_B}\right)^1,$$
(8)
$`ϵ_i=m_{i1}+m_{i2}M_B,ϵ_i>0`$. As far as I know, this $`Z`$-factor was first introduced by S.Weinberg in SWein many years ago as the field renormalization factor which defines the probability to find the physical deuteron $`|d`$ in a bare elementary-particle state $`|d_0`$, $`Z=|d_0|d|^2`$.
Even more detailed information is contained in the continuum counterpart of the factor $`Z`$, the spectral density $`w(M)`$ of the bare state, given by
$$w(M)=\underset{i}{}w_i(M),w_i(M)=4\pi \mu _ip_i|c(M)|^2\mathrm{\Theta }(Mm_{i1}m_{i2}),$$
(9)
where $`c(M)`$ is the probability amplitude to find the bare state in the continuum wave function $`|\mathrm{\Psi }_M`$. With $`c(M)`$ found from the system of equations (4), one can calculate $`w(M)`$:
$$w(M)=\frac{1}{2\pi i}\left(\frac{1}{MM_0+g^{}(M)}\frac{1}{MM_0+g(M)}\right).$$
(10)
As shown in bhm , the normalization condition for the distribution $`w(M)`$ follows from the completeness relation for the total wave function (1) projected onto bare state channel, and reads:
$$_{m_{01}+m_{02}}^{\mathrm{}}w(M)𝑑M=1Z,$$
(11)
if the system possesses a bound state, and
$$_{m_{01}+m_{02}}^{\mathrm{}}w(M)𝑑M=1,$$
(12)
if there is no bound state ($`m_{01}+m_{02}`$ is the lowest threshold). In the case of bound state present, all the information on the factor $`Z`$ is encoded, due to eq. (11), in the $`w(M)`$ too. On the other hand, the analysis in terms of $`w(M)`$ can be performed in the case of resonance as well, as exemplified in evidence .
In the latter case the $`t`$-matrix poles are situated in the complex plane. While the positions of the poles are the fundamental quantities, another quantities are useful for practical purposes. Namely, one defines the visible resonance mass $`M_R`$ from the equation
$$M_RM_0+\mathrm{Re}g(M_R)=0,$$
(13)
and calculates the visible width as
$$\mathrm{\Gamma }=2\mathrm{}\mathrm{Im}g(M_R),\mathrm{}=\left(1+\frac{\mathrm{Re}g(M)}{M}|_{M=M_R}\right)^1.$$
(14)
Clearly this brings the $`t`$-matrix into Breit-Wigner form, i.e. in the form in which experimental data are usually delivered. In what follows the factor $`\mathrm{}`$ will be called the renormalization factor.
There are some limitations of course, as not the every peak has the Breit-Wigner shape. In the case of overlapping resonances the formulae (13) and (14) do not work. The special case of near-threshold $`S`$-wave resonance is not described by Breit-Wigner or Flattè formula, and the scattering length parametrization is more appropriate Flatte .
The quantities $`Z`$ and $`w(M)`$ are the ones of immediate relevance. Indeed, there is no hope that, say, the elastic $`D\overline{D}`$ scattering will be measured some time. Our knowledge on mesonic resonances comes from external reactions, like $`e^+e^{}`$ annihilation, $`\gamma \gamma `$ collisions, $`B`$-meson decays etc. etc. Assuming that such reactions proceed via intermediate $`q\overline{q}`$ states, one obtains that the cross-section is proportional to $`w(M)`$:
$$\sigma (\mathrm{mesons})\mathrm{\Gamma }_{0r}w(M),$$
(15)
where $`\mathrm{\Gamma }_{0r}`$ is the width of the bare state corresponding to the external reaction. Such formulae were used in Eichten to describe the $`e^+e^{}`$ annihilation into charmed mesons. In the limit of narrow resonance eq.(15) is reduced to the standard Breit-Wigner formula
$$\sigma (\mathrm{mesons})\frac{1}{2\pi }\frac{\mathrm{\Gamma }_r^0\mathrm{\Gamma }}{(MM_0)^2+\frac{1}{4}\mathrm{\Gamma }^2},$$
(16)
where $`\mathrm{\Gamma }=2\mathrm{I}\mathrm{m}g(M_0)`$ is the (small) width of the resonance. Similarly, for the bound state case the width $`\mathrm{\Gamma }_r`$ for a given reaction is renormalized as
$$\mathrm{\Gamma }_r=Z\mathrm{\Gamma }_r^0.$$
(17)
## III The quark model
This section specifies the form factors $`f_i(𝐩)`$. The pair-creation model employed is the $`{}_{}{}^{3}P_{0}^{}`$ one, that is the pair-creation Hamiltonian is the nonrelativistic reduction of
$$H_q=g_qd^3x\overline{\psi }_q\psi _q,$$
(18)
for a given flavour $`q`$, but two important points make it different from the model used in Barnes1 ; Barnes2 ; Barnes3 .
The approach Barnes1 ; Barnes2 ; Barnes3 assumes that the pair creation is flavour-independent, which yields for the constant $`g_q`$ the form
$$g_q=\gamma 2m_q,$$
(19)
where $`\gamma `$ is the effective strength of pair-creation. The factor $`2m_q`$ implies enhancement of strange quarks creation comparing to light quarks one. There are no fundamental reasons to have such enhancement. Moreover, such factor is absent in microscopical models of pair creation, like Eichten and ABarnes . So, throughout the present study, I use the effective strength $`\gamma `$ for the creation of light ($`u`$\- and $`d`$-) flavours, while for strange quarks the effective strength $`\gamma _s=\frac{m_q}{m_s}\gamma `$ is used, where $`m_q`$ and $`m_s`$ are the constituent masses of light and strange quarks correspondingly.
The authors of Barnes1 ; Barnes2 ; Barnes3 argue that the assumption of flavour-independence gives a reasonably accurate description of known decays. One should have in mind, however, that the calculations Barnes1 ; Barnes2 ; Barnes3 are performed with the so-called SHO wavefunctions, i.e. with the wave functions of harmonic oscillator, and with the same oscillator parameter $`\beta `$ for all states. This assumption looks implausible, as the behaviour of form factors is defined by scales of wavefunctions, which, in turn, are defined by quark model.
In what follows the standard nonrelativistic potential model is introduced, with the Hamiltonian
$$H_0=\frac{p^2}{m_c}+V(r)+C,V(r)=\sigma r\frac{4}{3}\frac{\alpha _s}{r},$$
(20)
$`m_c`$ is the mass of charmed quark. This Hamiltonian should be supplied by Fermi-Breit-type relativistic corrections, including spin-spin, spin-orbit and tensor force, which cause splittings in the $`{}_{}{}^{2S+1}L_{J}^{}`$ multiplets. In the first approximation these splitting should be calculated as perturbations, using the eigenfunctions of the Hamiltonian (20). The same interaction $`V(r)`$ should be used in spectra and wavefunction calculations of $`D`$ ($`D_s`$) mesons.
In the first approximation the pair-creation amplitude is to be calculated with the eigenfunctions of the zero-order Hamiltonian. Use of the SHO wavefunctions simplify these calculations drastically. So the procedure adopted is to find the SHO wavefunctions (of the form $`\mathrm{exp}(\frac{1}{2}\beta ^2r^2)`$ multiplied by appropriate polynomials) for each orbital momentum $`L`$ and radial quantum number $`n`$, with the effective value of oscillator parameter $`\beta `$ for each $`L`$ and $`n`$, which reproduce the r.m.s. of the states.
I use the following set of potential model parameters:
$$\alpha _s=0.55,\sigma =0.175GeV^2,m_c=1.7GeV,C=0.271GeV,$$
$$m_q=0.33GeV,m_s=0.5GeV.$$
(21)
The spin-dependent force is taken in the form
$$V_{sd}=V_{HF}+\frac{2\alpha _s}{m_c^2r^3}𝐋𝐒\frac{\sigma }{2m_c^2r}𝐋𝐒+\frac{4\alpha _s}{m_c^2r^3}T,$$
(22)
where $`𝐋𝐒`$ and $`T`$ are spin-orbit and tensor operators correspondingly, and $`V_{HF}`$ is the contact hyperfine interaction,
$$V_{HF}(r)=\frac{32\pi \alpha _s}{9m_c^2}\stackrel{~}{\delta }(r)𝐒_q𝐒_{\overline{q}},$$
(23)
where, following the lines of Barnes3 , Gaussian-smearing of the hyperfine interaction is introduced,
$$\stackrel{~}{\delta }(r)=\left(\frac{\kappa }{\sqrt{\pi }}\right)^3e^{\kappa ^2r^2},$$
(24)
with $`\kappa =1.45`$ GeV.
The masses and effective values of oscillator parameters $`\beta `$ for the model (21) are listed in the Table I. The effective values of oscillator parameter for $`D`$-mesons are $`\beta _D=0.385`$ GeV, and $`\beta _{D_s}=0.448`$ GeV.
One should not take the numbers given in last column too seriously, especially for higher states, as the fine splittings are not well-known, and the expression (22) is surely too naive. Moreover, various much more sophisticated approaches, which reproduce the splittings in $`1P`$ multiplet, give different predictions for higher multiplets, as discussed in detail in Alla .
The $`D`$-meson masses taken are $`M_D=1.867`$ GeV, $`M_D^{}=2.008`$ GeV, $`M_{D_s}=1.969`$ GeV, $`M_{D_S^{}}=2.112`$ GeV, so that the mass difference between neutral and charged $`D`$-mesons is not taken into account. The pair-creation strength for light quarks $`\gamma =0.322`$ is used. The $`{}_{}{}^{3}P_{0}^{}`$ amplitudes are listed in the Appendix A.
## IV Lower charmonia and $`D`$-levels
The single-level version of the coupled-channel model is used in what follows, with the exception of $`2^3S_11^3D_1`$ levels.
All the physical charmonium masses below threshold are known. The hadronic shifts and bare masses were calculated from the equation 6,
$$M_0=M_{phys}+\delta ,\delta =g(M_{phys})=\underset{i}{}g_i(M_{phys}),$$
(25)
where the sum is over mesonic channels $`D\overline{D}`$, $`D\overline{D}^{}`$, $`D^{}\overline{D}^{}`$, $`D_s\overline{D}_s`$, $`D_s\overline{D}_s^{}`$ and $`D_s^{}\overline{D}_s^{}`$ (an obvious shorthand notation is used here and in what follows: $`D\overline{D}^{}D\overline{D}^{}+\overline{D}D^{}`$, and $`D_s\overline{D}_s^{}D_s\overline{D}_s^{}+\overline{D}_sD_s^{}`$).
Besides, as the position of the $`1^{}`$ $`\psi (3770)`$ state is well-established, the bare mass of $`1^3D_1`$ state was reconstructed by means of eq.13, and the visible width was calculated as (14).
The parameters of the underlying quark model (21) and the pair-creation strength $`\gamma =0.322`$ were chosen to reproduce, with reasonable accuracy, the model masses of $`1S`$, $`1P`$ and $`2S`$ states, and the width of $`\psi (3770)`$.
The results for bound states are given in Table II together with corresponding values of $`Z`$-factors. The shifts are much smaller than in Barnes3 , but still substantial.
Let me now discuss the $`1^3D_1`$ level, lying above $`D\overline{D}`$ threshold. The bare mass is calculated to be $`4.018`$ GeV. The calculated width of the $`\psi (3770)`$ is $`25.5`$ MeV, which compares well with the PDG value of $`23.6\pm 2.7`$ MeV PDG . Note that it is the visible width, while naive calculations would give
$$\mathrm{\Gamma }_0=2\mathrm{Im}g(M_R)=34.3MeV.$$
(26)
So the effect of coupling to mesonic channels on the width of $`\psi (3770)`$ is not small, $`\mathrm{}=0.743`$.
The mass of $`\psi (3770)`$ is less than $`100`$ MeV higher than the mass of $`\psi ^{}(3686)`$, with $`D\overline{D}`$ threshold opening in between, so one could in principle expect that these states are mixed due to coupling to mesonic channels. This mixing is not large in the given model, as the mixing disappears if the mass difference between various $`D`$-mesons is neglected (see the corresponding spin-orbit recoupling coefficients listed in Appendix A). The relevant formulae for two-level mixing scheme are given in the Appendix B, and here I quote the results. With physical masses of $`\psi ^{}(3686)`$ and $`\psi (3770)`$ the bare masses are reconstructed as $`M(2^3S_1)=3.899`$ GeV, and $`M(1^3D_1)=4.016`$ GeV, so the masses are shifted only by a few MeV due to the mixing. The width of the $`\psi (3770)`$ becomes only $`18.4`$ MeV, but it is the visible width, and the deviation from the value $`25.5`$ MeV obtained without mixing is mainly due to the illegitimate attempt to fit the system of two overlapping states with a single-Breit-Wigner lineshape. The states indeed do overlap, as shown at Fig.1, where the spectral densities of bare $`2^3S_1`$ and $`1^3D_1`$ bare states are plotted. The lineshape of $`1^3D_1`$ is distorted due to the mixing, and the lineshape of the $`2^3S_1`$ is drastically changed, displaying, instead of small smooth background, a peak at the mass of about $`3760`$ MeV. The $`10`$ MeV difference between the peak position and the mass of $`\psi (3770)`$ is again due to the prescription of visible width: single-Breit-Wigner approximation is not appropriate both for $`{}_{}{}^{3}S_{1}^{}`$ and $`{}_{}{}^{3}D_{1}^{}`$ lineshapes. As to the lower state $`\psi ^{}(3686)`$, the admixture $`Z_1`$ of the $`2^3S_1`$ bare state is $`0.742`$, while the admixture $`Z_2`$ of the bare $`1^3D_1`$ state is $`0.00063`$.
The calculated mass of bare $`1^3D_1`$ state is $`4.0164.018`$ GeV. So the value of $`4.026`$ GeV given in the Table I looks quite acceptable. No other $`1D`$-states are known, so in what follows I take, with reservations mentioned in the previous section, the values of bare masses from the Table I. With the bare $`1^1D_2`$ state having mass of $`4.043`$ GeV, the physical state is a bound state,
$$M(1^1D_2)=3.800GeV,Z(1^1D_2)=0.712.$$
(27)
Similarly,
$$M(1^3D_2)=3.806GeV,Z(1^3D_2)=0.689$$
(28)
for the mass $`4.043`$ GeV of the bare $`1^3D_2`$ state. The $`1^3D_3`$ state is allowed to decay into $`D\overline{D}`$, but the width is extremely small, as the $`D\overline{D}`$ system is in the $`F`$-wave:
$$M(1^3D_3)=3.812GeV,\mathrm{\Gamma }(1^3D_3)0.7MeV,\mathrm{}(1^3D_3)=0.717.$$
(29)
## V $`2P`$-levels
Coupled-channel effects do not cause dramatic changes for the charmonia states discussed in the previous section. The situation with $`2P`$-levels promises more, as $`2P`$-charmonia are expected to populate the mass range of $`3.904.00`$ GeV, where more charmed meson channels start to open, and some of these channels are the $`S`$-wave ones.
The importance of $`S`$-wave channels follows from the eq.(5). The form factor $`f(p)`$ of the $`S`$-wave mesonic channel behaves as some constant at small $`p`$, so the derivative of the hadronic shift $`g(M)`$ with respect to the mass is large for the masses close to the $`S`$-wave threshold. As the result, hadronic shift due to the coupling to $`S`$-wave channel displays rather vivid cusp-like near-threshold behaviour.
The physical masses and widths of $`2P`$ states were calculated with bare masses given by Table I, and the results are listed in Table III. Two different values are given for the mass and width of the $`2^3P_0`$ state, for two different choices of the bare mass, see below. Looking at the numbers one would say that nothing dramatic has happened due to $`S`$-wave thresholds. Indeed, all the shifts are about $`200`$ MeV, and the renormalization factors are about $`0.50.6`$. It is the behaviour of spectral density which reveals the role of $`S`$-wave thresholds.
For the $`2^1P_1`$ case the $`S`$-wave thresholds are $`D\overline{D}^{}`$ with the spin-orbit recoupling coefficient $`C=1/\sqrt{2}`$ and multiplicity 4, the $`D_s\overline{D}_s^{}`$ with $`C=1/\sqrt{2}`$ and multiplicity 2, the $`D^{}\overline{D}^{}`$ with $`C=1`$ and multiplicity 2, and $`D_s^{}\overline{D}_s^{}`$ with $`C=1`$ and multiplicity 1. If the resonance is in the mass range $`3.904.00`$ GeV (and it appears to be so), then the relevant thresholds are $`D\overline{D}^{}`$ and $`D^{}\overline{D}^{}`$. The spectral density of the bare $`2^1P_1`$ state is shown at Fig.2. It displays the relatively steep rise near $`D\overline{D}^{}`$ threshold, and a beautiful well-pronounced cusp due to the opening of $`D^{}\overline{D}^{}`$ channel.
In the $`2^1P_1`$ case the $`S`$-wave strength is shared equally between $`D\overline{D}^{}`$ and $`D^{}\overline{D}^{}`$ channels, while for the $`2^3P_2`$ case all $`S`$-wave strength is concentrated in the $`D^{}\overline{D}^{}`$, channel. As the result, the cusp due to the opening of $`D^{}\overline{D}^{}`$ channel is more spectacular in the $`2^3P_2`$ case, as shown at Fig.3.
The case of $`2^3P_1`$ is even more interesting. Here, similarly to the $`2^3P_2`$ case, all the $`S`$-wave strength is concentrated in two channels, $`D\overline{D}^{}`$, $`D_s\overline{D}_s^{}`$, and the multiplicity of the former is 4. So the strongest $`S`$-wave threshold is the $`D\overline{D}^{}`$ one, well below the resonance. The behaviour of the $`2^3P_1`$ bare state spectral density is shown at Fig.4, and is very peculiar: together with a clean and relatively narrow resonance, there is a near-threshold peak, rising at the flat background. The $`D\overline{D}^{}`$ scattering length appears to be negative and large,
$$a_{D\overline{D}^{}}=8fm,$$
(30)
signalling the presence of virtual state very close to the $`D\overline{D}^{}`$ threshold, with the energy $`ϵ=0.32`$ MeV. So the coupling to mesonic channels has generated not only the resonance, but, in addition, a virtual state very close to physical region.
The near-threshold peak should fade with the increase of bare state mass, and strengthen otherwise. There are uncertainties in the fine splitting estimates, so the mass of the bare $`2^3P_1`$ state could easily be about $`10`$ MeV larger or smaller. The dependence of spectral density behaviour on the bare mass is shown at Fig.5. The peak becomes less pronounced for the bare mass of $`4.190`$ GeV, but the scattering length remains rather large, $`a5.2`$ fm, which corresponds to the energy of virtual state of about $`0.76`$ MeV (compare this with the scattering length in the $`1^+`$ channel, $`|a|1`$ fm). $`10`$ MeV decrease of the bare state mass leads to incredibly large scattering length, $`a17.8`$ fm, and virtual state with the energy $`0.065`$ MeV. Further decrease of bare state mass leads to moving the state to the physical sheet, i.e. to appearance of the bound state. This happens at the bare mass of about $`4.160`$ GeV, which seems, in the present model, to be beyond acceptable range for fine splitting. Similarly, the bound state appears if the pair-creation strength is increased by several per cent.
The $`2^3P_0`$ level is a disaster, as it always happens with scalars. The bare mass is considerably lower that the c.o.g., and the uncertainty in the fine splitting estimate is large. The $`S`$-wave $`D^{}\overline{D}^{}`$ and $`D_s^{}\overline{D}_s^{}`$ channels are too high. The $`S`$-wave $`D\overline{D}`$ channel is too low. The only relevant $`S`$-wave channel is $`D_s\overline{D}_s`$ ($`C=\sqrt{\frac{3}{2}}`$, multiplicity 1), and corresponding threshold is at $`3.938`$ GeV, i.e. around the region where the resonance is expected. So, depending on the position of the bare state, variety of spectral density behaviour can be achieved, as shown at Fig.6, where spectral density is plotted for $`M_0(2^3P_0)=4.108`$ GeV and $`M_0(2^3P_0)=4.140`$ GeV. Note that two curves of Fig.6 correspond to two completely different situations. The curve for $`M_0(2^3P_0)=4.108`$ GeV displays the normal resonance behaviour with a tiny cusp due to the opening of the $`D_s\overline{D}_s`$ channel. The curve for $`M_0(2^3P_0)=4.140`$ GeV is not resonance-like at all. The formal exercise of calculating the $`\mathrm{}`$-factor and visible width does not make much sense, and, as suggested in Flatte , such excitation curve should be analysed in terms of scattering length approximation, and not in terms of Breit-Wigner or Flattè distributions.
## VI Discussion
The quark model (21) is not the result of the sophisticated fit, it is rather a representative example. A serious fit should include proper treatment of fine and hyperfine splittings, as well as the mixing of bare $`2^3S_1`$ and $`1^3D_1`$ states. The calculations should be performed with more realistic wavefunctions, and not with the SHO ones. Relativistic corrections should be taken into account in calculations of bare spectra, and more realistic model should be used for $`D`$-meson wavefunctions. More QCD-motivated model should be employed for the pair-creation Hamiltonian, and loop integrals $`g(M)`$ are to be treated relativistically. Nevertheless, there are gross features which are model-independent.
In various pair-creation models, the shifts within the each orbital multiplet are approximately the same, and differ only due to the mass difference of bare states in the multiplet and different masses of charmed mesons. This is model-independent, as the charmed quark is heavy, and the wavefunctions within each $`nL`$ multiplet do not differ much. The same is true for the wavefunctions of $`D`$ and $`D^{}`$ (and $`D_s`$ and $`D_s^{}`$) mesons due to heavy quark spin symmetry.
In addition, the shifts for all states are more or less the same. There are no a priori reasons for this, but the numbers given in Table II are stable up to overall multiplier $`\gamma `$ as soon as the scales $`\beta `$ behave as expected from quark model, and the value of $`\gamma `$ is constrained by experimental value of the $`\psi (3770)`$ width.
Indeed, the present results appear to be very similar to the ones given in Eichten ; nEichten . One might suggest that, from phenomenological point of view, the effect of coupling to mesonic channels can be approximated by adding a negative constant to the potential. But it is not the whole story, as the coupling to charmed mesons generates the admixture of $`D`$-mesons in the wavefunctions of bound states, which affects such quantities as $`e^+e^{}`$ widths (in accordance with eq17) and the rates of radiative transitions, as discussed in detail in Eichten .
The $`e^+e^{}`$ widths of $`J/\psi `$ and $`\psi ^{}(3686)`$ are more or less accurately described by Van Royen-Weisskopf formula with QCD correction, so the renormalization (17) is not harmless, even if it is as mild as 10 $`\%`$ reduction required for $`J/\psi `$ by the results of Table II. For $`\psi ^{}(3686)`$ $`Z`$-factor is 0.743, so renormalization is larger.
In the nonrelativistic quark model the $`e^+e^{}`$ width of the $`{}_{}{}^{3}D_{1}^{}`$ state is zero. The mixing communicates the $`e^+e^{}`$ width of the $`2^3S_1`$ state to the $`\psi (3770)`$ region, with the $`e^+e^{}D\overline{D}`$ crossection given by eq.15 where the spectral density of $`2^3S_1`$ state $`w_S(M)`$ is substituted. It is reasonable to estimate the $`e^+e^{}`$ width of the $`\psi (3770)`$ using the peak value of $`w_S(M)`$, which yields less than $`1/3`$ of the measured value $`0.26\pm 0.04`$ keV PDG . Similar result was obtained in Eichten . The bare $`1^3D_1`$ state has the small $`e^+e^{}`$ width of its own as relativistic correction, and the bare $`2^3S_1`$ and $`1^3D_1`$ states are to be mixed by tensor force. The scale of the admixture required to reproduce the relatively large $`e^+e^{}`$ width of the $`\psi (3770)`$ is not small, as shown in Rosner . The coupling to charmed mesons is able to explain only about $`1/3`$ of the observed $`e^+e^{}`$ width of $`\psi (3770)`$, so direct mixing between bare $`2^3S_1`$ and $`1^3D_1`$ bare levels is still needed, which would reduce the $`e^+e^{}`$ width of the $`\psi ^{}`$ further. While the problem of leptonic widths is an open problem for coupled-channel model, it is clear that the values of $`Z`$-factors considerably smaller than given in Table II would destroy fragile agreement with the data on $`e^+e^{}`$ widths achieved by quark model practitioners.
The $`2^3S_1`$$`1^3D_1`$ mixing due to coupling to meson channels is small, and this is also model-independent, as it vanishes when the masses and wavefunctions of $`D`$-mesons are taken to be the same. Note, however, that both $`\psi (3770)`$ and $`\psi ^{}`$ states contain considerable admixture of four-quark component of their own, in the form of various $`D`$-meson pairs. Thus, if the mass difference between charged and neutral $`D`$-mesons is taken into account, some isospin violation could be generated. It is argued in Voloshincc that a small admixture of $`I=1`$ component is needed to explain some discrepancies in the observed properties of $`\psi (3770)`$ and $`\psi ^{}`$. The question of whether such admixture can be generated by coupled-channel effects certainly deserves attention.
The shifts of $`D`$-levels are more or less the same, the relevant thresholds are $`P`$-wave ones, and the $`2^3S_1`$$`1^3D_1`$ mixing does not affect the shift of the bare $`{}_{}{}^{3}D_{1}^{}`$ level drastically. So the combined effect of fine splitting and splitting due to coupled-channel effects on other $`D`$-levels is not large. In particular, it means that, with the mass of $`\psi (3770)`$ as experimental input, the physical $`{}_{}{}^{3}D_{2}^{}`$ or $`{}_{}{}^{3}D_{3}^{}`$ state cannot be placed as high as $`3.872`$ GeV and, therefore, cannot be identified as $`X(3872)`$, unless something drastic happens with fine splittings in the $`1D`$-multiplet.
Due to the presence of $`S`$-wave thresholds, the situation with $`2P`$-levels is more interesting. The coupling to $`D^{}\overline{D}^{}`$ channels generates pronounced cusps in the spectral densities of bare $`2^1P_1`$ and $`2^3P_2`$ levels, and the coupling to $`D\overline{D}^{}`$ channel generates the strong threshold effect in the $`1^{++}`$ wave. Within the given model, it is a virtual state with the energy less than $`1`$ MeV. The mechanism of generating such a state is quite peculiar: it is one and a same bare $`2^3P_1`$ state, which gives rise both to the $`1^{++}`$ resonance with the mass of about $`3990`$ Mev, and a virtual state at the $`D\overline{D}^{}`$ threshold, i.e. where the $`X(3872)`$ is observed. Due to the presence of strong $`S`$-wave threshold, the hadronic shift appears to be large enough to destroy the one-to-one correspondence between bare and physical states, as widely discussed in connection with light scalar mesons Tornqvist ; Pennington .
To what extent this prediction is robust? Changes of the underlying quark model parameters or of the value of pair-creation strength can shift this extra state either to the physical sheet or away from the physical region. In the latter case, however, the $`D\overline{D}^{}`$ scattering length remains large. One should have in mind that if such dynamical generation of extra state at $`D\overline{D}^{}`$ threshold is possible in the charmonia, the $`1^{++}`$ channel is the most appropriate place for this phenomenon. The latter statement is model independent.
First, note that the scalar charmonium does not decay into $`D\overline{D}^{}`$ at all, and the tensor one decays into $`D\overline{D}^{}`$ in the $`D`$-wave. As to $`1^{++}`$ and $`1^+`$ levels, they both have the desired $`S`$-wave decay mode. Apply now the heavy quark spin selection rule Voloshinspin , which suggests that the spin of a heavy quark pair is conserved in the decay. The $`S`$-wave decay mode comes from the four-quark state $`c\overline{c}q\overline{q}`$ with all relative angular momenta equal to zero. Then the total angular momentum $`J=1`$ can result only from the quark spins. The combination of $`S_{c\overline{c}}=1`$ and $`S_{q\overline{q}}=1`$ is $`C`$-even, while the combination of $`S_{c\overline{c}}=0`$ and $`S_{q\overline{q}}=1`$ is $`C`$-odd. It is a simple algebra exercise to show that, symbolically,
$$(c\overline{c})_{S=1}(q\overline{q})_{S=1}\frac{1}{\sqrt{2}}(D\overline{D}^{}+\overline{D}D^{}),$$
(31)
and
$$(c\overline{c})_{S=0}(q\overline{q})_{S=1}\frac{1}{2}(D^{}\overline{D}D\overline{D}^{})+\frac{1}{\sqrt{2}}D^{}\overline{D}^{},$$
(32)
and, independently of the pair-creation model, all the $`S`$-wave strength of the $`{}_{}{}^{3}P_{1}^{}`$ decay is concentrated in the $`D\overline{D}^{}`$ channel, while in the $`{}_{}{}^{1}P_{1}^{}`$ decay it is shared equally between $`D\overline{D}^{}`$ and $`D^{}\overline{D}^{}`$. Thus the threshold attraction in the $`1^{++}`$ $`D\overline{D}^{}`$ channel is always much stronger than in the $`1^+`$ one. It would be interesting to see if the pair-creation model of Eichten is able to generate large $`D\overline{D}^{}`$ scattering length.
The model SwansonX contains the detailed analysis of the $`X(3872)`$ as a state bound both by pion exchange and quark exchange in the form of transitions $`D\overline{D}^{}J/\psi \rho ,J/\psi \omega `$, with the latter contributions being important for the binding. In fact, this model has predicted the $`J/\psi \omega `$ decay mode of the $`X(3872)`$, and, after observation Belleomega of the decay $`X(3872)\pi ^+\pi ^{}\pi ^0J/\psi `$, it has become almost official model of $`X(3872)`$. This is challenged by preliminary data from Belle olsen on large $`D^0D^0\pi ^0`$ rate, more than ten times larger than $`\pi ^+\pi ^{}J/\psi `$ one, while the model SwansonX claims the opposite.
One could question validity of the naive quark-exchange model. Besides, the pion exchange is definitely attractive in the $`1^{++}`$ channel, but there are uncertainties in the actual calculations; the details of binding depend on the cutoff scale $`\mathrm{\Lambda }`$, as recognized in SwansonX . The attraction found in the coupled-channel model is large, and could help binding without large quark-exchange kernels, and, correspondingly, without large $`\pi ^+\pi ^{}J/\psi `$ rate.
From practical point of view, the wavefunctions of both models are not very distinguishable. Indeed, the near-threshold virtual state of the coupled-channel model owes its existence to the bare $`1^{++}`$ state, but the near-threshold admixture of the bare state in the wavefunction is extremely small, as seen from Fig.4 (recall that the spectral density is normalized to unity). So the decays like $`D\overline{D}\pi `$ and $`D\overline{D}\gamma `$ would proceed via $`D^{}`$ decays, as described in Voloshindecay . As to short-distance decays and exclusive production, the rates of these are governed by large scattering length. This phenomenon was called low-energy universality in Braaten1 . Consider, for example, the near-threshold production of $`D\overline{D}^{}`$ pairs in the reaction $`BD\overline{D}^{}K`$. As explained in Section II, the $`D\overline{D}^{}`$ invariant mass distribution is proportional to $`w(M)`$ in the coupled channel model, with the lineshape plotted at Fig.5. Now compare these curves with the ones presented in Braaten2 ; Braaten3 with the scattering length approximation for the $`D\overline{D}^{}`$ amplitude, and observe that the low-energy universality indeed takes place.
The lineshape for $`D\overline{D}^{}`$ production depends only on the modulus of scattering length, and the cases of bound state and virtual state are not distinguishable. If some inelastic channel like $`\pi \pi J/\psi `$ is present, then, as shown in Braaten3 , the lineshape for this channel depends on whether there is a bound state or virtual state (cusp). The area under the cusp is much smaller than the area under the resonance, so, in principle, the data could distinguish between $`X(3872)`$ as a bound state or a virtual state.
To conclude the discussion of coupled-channel virtual state as $`X(3872)`$ I note that, with $`7`$ MeV difference between $`D^0\overline{D}^0`$ and $`D^+\overline{D}^{}`$ thresholds, the coupled-channel model is able to generate considerable isospin violation. The coupled-channel calculations which resolve the mass difference between charged and neutral $`D`$-mesons are in progress now.
Last year yet another new state was reported BelleY as an enhancement in $`\omega J/\psi `$ mode, with the mass of $`3941`$ MeV and the width of about $`90`$ MeV ($`Y(3940`$). As there are two $`1^{++}`$ physical states per one bare $`2^3P_1`$, it is tempting to identify the $`1^{++}`$ resonance with the $`Y(3940)`$ state, following the suggestion of Bugg . Nevertheless, with the given set of quark model parameters the resonance is $`4050`$ MeV higher than $`3940`$ MeV, and is narrow.
One might suggest the $`0^{++}`$ assignment for $`Y(3940)`$, as the coupled-channel model places the scalar just at the right place. However, as the $`D_s\overline{D}_s`$ channel is opening at $`3938`$ MeV, the width of the state cannot be as large as observed $`90`$ MeV independently of what channel saturates it, unless, for some strange reasons, the state couples weakly to $`D_s\overline{D}_s`$. It is similar to $`a_0(980)/f_0(980)`$ case, where, in spite of large coupling to light pseudoscalars ($`\pi \pi `$ or $`\pi \eta `$), the visible width is small due to strong affinity to $`K\overline{K}`$ threshold.
## VII Conclusions
The spectrum of charmonium states below $`4`$ GeV is calculated taking into account coupling to the pairs of lowest $`D`$ and $`D_s`$ mesons. The analysis is performed within the $`{}_{}{}^{3}P_{0}^{}`$ model for light-quark pair-creation. It appears that quite moderate modification of quark model parameters is needed to describe charmonia below $`D\overline{D}`$ threshold, while coupling to open charm does not cause drastic effects on $`1D`$-levels. This is in contrast to $`2P`$-levels, where opening of strong $`S`$-wave channels leads to pronounced threshold effects.
In particular, coupling of the bare $`2^3P_1`$ state to $`D\overline{D}^{}`$ channel generates, together with the $`1^{++}`$ resonance with the mass of $`3990`$ MeV, a near-threshold virtual state with the energy of about $`0.3`$ MeV, which corresponds to the extremely large $`D\overline{D}^{}`$ scattering length $`a8`$ fm, with the possibility to identify this state with $`X(3872)`$. The admixture of the bare $`c\overline{c}`$ state in the near-threshold wavefunction is very small, so it is essentially the $`D\overline{D}^{}`$ state. The $`1^{++}`$ channel appears to be the only one where such state can be formed.
###### Acknowledgements.
I am grateful to Yu.A.Simonov for useful discussions. This research is supported by the grant NS-1774.2003.2, as well as of the Federal Programme of the Russian Ministry of Industry, Science, and Technology No 40.052.1.1.1112.
Appendix A
The $`{}_{}{}^{3}P_{0}^{}`$ form factors are defined following the lines of ABarnes ; Barnes1 , adapted for charmonium transitions. Let the $`c\overline{c}`$ meson $`A`$ decay to $`q\overline{c}`$ meson $`B`$ and $`c\overline{q}`$ meson $`C`$. Then the spin-space part of the amplitude in the c.m. frame of the initial meson $`A`$ is given by
$$I(𝐩)=d^3k\varphi _A(𝐤𝐩)<s_qs_{\overline{q}}|\widehat{O}(𝐤)|0>\varphi _B^{}(𝐤r_q𝐩)\varphi _C^{}(𝐤r_q𝐩),$$
(A.1)
$`\varphi _A`$ is the wavefunction of the initial meson in the momentum space, $`\varphi _B`$ and $`\varphi _C`$ are the wavefunctions of the final mesons $`B`$ and $`C`$, $`𝐩=𝐩_B`$, $`r_q=\frac{m_q}{m_q+m_c}`$, $`m_c`$ is the mass of charmed quark, $`m_q`$ is the mass of light quark. $`\widehat{O}(𝐤)`$ is the $`{}_{}{}^{3}P_{0}^{}`$ operator:
$$\widehat{O}(𝐤)=2\gamma \sigma 𝐤.$$
(A.2)
The form factors for the transition $`ABC`$ are given below in the $`lS`$ basis, where $`l`$ is the orbital momentum in the final meson system, and $`\stackrel{}{S}=\stackrel{}{J}_B+\stackrel{}{J}_C`$ is the total spin of the mesons $`B`$ and $`C`$. In the narrow width approximation these form factors $`f_{lS}`$ define the partial widths $`\mathrm{\Gamma }_{ABC}`$ as
$$\mathrm{\Gamma }_{ABC}=2\mathrm{Im}g_{BC}(M_A)=2\pi p_{BC}\mu _{BC}\underset{lS}{}|f_{lS}|^2.$$
(A.3)
The form factors calculated with SHO wavefunctions take the form
$$f_{lS}=\frac{\gamma }{\pi ^{1/4}\beta _A^{1/2}}\mathrm{exp}\left(\frac{p^2(r_q1)^2}{\beta _B^2+2\beta _A^2}\right)𝒫_{lS},$$
(A.4)
where $`\beta _A`$ and $`\beta _B=\beta _C`$ are the oscillator parameters of initial meson $`A`$ and final mesons $`B`$ and $`C`$, $`\gamma `$ is the pair-creation strength, and $`r_q=m_q/(m_c+m_q)`$, $`m_c`$ is the mass of charmed quark, $`m_q`$ is the mass of light quark. The polynomial $`𝒫_{LS}`$ is a channel-dependent one:
$$𝒫_{lS}=f_lC_{lS},$$
(A.5)
where $`C_{lS}`$ are spin-orbit recoupling coefficients for specific mesonic channels, and $`f_l`$ are:
$$f_P(1S1S+1S)=\frac{2^3}{3^2}\frac{\lambda \beta ^3}{\beta _B^3\beta _A}p$$
(A.6)
$$f_P(2S1S+1S)=\frac{2^{5/2}}{3^{3/2}}\frac{\beta ^3}{\beta _A\beta _B^3}p\left(\lambda +\left(\frac{10}{9}\lambda \frac{4}{9}\right)\frac{\beta ^2}{\beta _A^2}+\frac{2p^2}{3\beta _A^2}\lambda (\lambda 1)^2\right)$$
(A.7)
$$f_S(1P1S+1S)=\frac{2^4}{3^{5/2}}\frac{\beta ^5}{\beta _B^3\beta _A^2}\left(1+\lambda (\lambda 1)\frac{p^2}{\beta ^2}\right)$$
(A.8)
$$f_D(1P1S+1S)=\frac{2^4}{3^25^{1/2}}\frac{\lambda (\lambda 1)\beta ^3}{\beta _A^2\beta _B^3}p^2$$
(A.9)
$$f_P(1D1S+1S)=\frac{2^{11/2}}{3^3}\frac{(\lambda 1)\beta ^5}{\beta _A^3\beta _B^3}p\left(1+\frac{3}{5}\lambda (\lambda 1)\frac{p^2}{\beta ^2}\right)$$
(A.10)
$$f_F(1D1S+1S)=\frac{2^4}{3^{3/2}5^{1/2}7^{1/2}}\frac{\lambda (\lambda 1)^2\beta ^3}{\beta _A^3\beta _B^3}p^3$$
(A.11)
$$f_S(2P1S+1S)=$$
$$\frac{2^{7/2}5^{1/2}}{3^{5/2}}\frac{\beta ^5}{\beta _B^3\beta _A^2}\left(1\frac{2\beta ^2}{3\beta _A^2}+\lambda (\lambda 1)\frac{p^2}{\beta ^2}\frac{2p^2}{3\beta _A^2}(\lambda 1)(2\lambda 1)\frac{2p^4}{5\beta ^2\beta _A^2}\lambda (\lambda 1)^3\right)$$
(A.12)
$$f_D(2P1S+1S)=$$
$$\frac{2^{7/2}}{3^2}\frac{\beta ^3}{\beta _A^2\beta _B^3}p^2\left(\lambda (\lambda 1)\frac{2\beta ^2}{15\beta _A^2}(\lambda 1)(7\lambda 2)\frac{2p^2}{5\beta _A^2}\lambda (\lambda 1)^3\right)$$
(A.13)
Here
$$\lambda =\frac{\beta _B^2+2r_q\beta _A^2}{\beta _B^2+2\beta _A^2},$$
(A.14)
and
$$\beta ^2=\frac{3\beta _A^2\beta _B^2}{\beta _B^2+2\beta _A^2}.$$
(A.15)
For the transitions to strange mesons, one should replace $`r_q`$ by $`r_s=m_s/(m_c+m_s)`$, and insert the multiplier $`\frac{m_q}{m_s}`$, as explained in the main text.
The decays of light mesons were considered in Barnes1 , which corresponds to $`r_q=1/2`$. In the case of $`\beta _A=\beta _B=\beta `$, and $`r_q=1/2`$ the expressions for the amplitudes listed above are equal to those of Barnes1 up to the factor of $`1/2`$. In general, there are two graphs with different topologies which contribute to the $`{}_{}{}^{3}P_{0}^{}`$ amplitude, and the sum of both is quoted in Barnes1 , while in actual calculations each graph contributes with the individual flavour factor. In the case of charmonia transitions only one graph contributes, so that the amplitude for the transition into given charge channel is equal to (A.4) with the flavour factor of unity.
The mass difference between neutral and charged mesons is not taken into account, so the sum over charge states is equivalent to introducing the multiplicity factor 2 for $`D\overline{D}`$ and $`D^{}\overline{D}^{}`$ channels, and 4 for $`D\overline{D}^{}`$ channel. The multiplicity factor for $`D_s\overline{D}_s`$ and $`D_s^{}\overline{D}_s^{}`$ is 1, and it is 2 for $`D_s\overline{D}_s^{}`$ channel.
Spin-orbit recoupling coefficients are tabulated in the Appendix A of Barnes1 . In the single-level coupled-channel calculations the squares of spin-orbit recoupling coefficients are needed, given in Table IV. Note that the multiplicity factors for $`D\overline{D}^{}`$ channels are twice as large as of $`D\overline{D}`$ and $`D^{}\overline{D}^{}`$ ones.
The case of $`2^3S_1`$$`1^3D_1`$ mixing requires explicit expressions for spin-orbit recoupling coefficients, as relative signs are important in the two-level mixing scheme. These are given below:
$$C_{10}(^3S_1^1S_0+^1S_0)=1$$
$$C_{11}(^3S_1^3S_1+^1S_0)=\sqrt{2}$$
$$C_{10}(^3S_1^3S_1+^3S_1)=\sqrt{\frac{1}{3}}$$
$$C_{12}((^3S_1^3S_1+^3S_1)=\sqrt{\frac{20}{3}}$$
(A.16)
$$C_{10}(^3D_1^1S_0+^1S_0)=\sqrt{\frac{5}{12}}$$
$$C_{11}(^3D_1^3S_1+^1S_0)=\sqrt{\frac{5}{24}}$$
$$C_{10}(^3D_1^3S_1+^3S_1)=\frac{\sqrt{5}}{6}$$
$$C_{12}(^3D_1^3S_1+^3S_1)=\frac{1}{6}$$
$$C_{32}(^3D_1^3S_1+^3S_1)=\sqrt{\frac{28}{5}}$$
(A.17)
Appendix B
The formulae necessary to describe the $`2^3S_1`$$`1^3D_1`$ mixing are collected here.
Two sets of form factors, $`f_S`$ and $`f_D`$, are introduced, which describe the transitions between mesons and $`2^3S_1`$, $`1^3D_1`$ levels. The system of coupled channel equation similar to (4) leads to the $`D\overline{D}`$ $`t`$-matrix:
$$t(𝐩,𝐩^{},M)=\underset{\mu ,\nu }{}f_{\mu ,D\overline{D}}(𝐩)\tau _{\mu \nu }(M)f_{\nu ,D\overline{D}}(𝐩^{}),$$
(B.1)
where sum is over $`S`$ and $`D`$ states, and
$$\tau _{SS}(M)=\frac{MM_D+g_{DD}(M)}{\mathrm{\Delta }(M)},$$
$$\tau _{DD}(M)=\frac{MM_S+g_{SS}(M)}{\mathrm{\Delta }(M)},$$
$$\tau _{DS}(M)=\tau _{SD}(M)=\frac{g_{SD}(M)}{\mathrm{\Delta }(M)},$$
(B.2)
$$\mathrm{\Delta }(M)=[MM_S+g_{SS}(M)][MM_D+g_{DD}(M)]g_{SD}^2,$$
(B.3)
$`M_S`$ and $`M_D`$ are the masses of bare states, and
$$g_{\mu \nu }(M)=\underset{i}{}g_{i,\mu \nu }(M),g_{i,\mu \nu }(M)=\frac{f_{\mu ,i}(𝐩)f_{\nu ,i}(𝐩)}{\frac{p^2}{2\mu _i}E_ii0}d^3p,$$
(B.4)
the index $`i`$ labels mesonic channels.
The visible physical masses are defined, similarly to (13), from the equation
$$\mathrm{Re}\mathrm{\Delta }(M_{phys})=0.$$
(B.5)
There are two solutions of this equation, $`M_b`$ below $`D\overline{D}`$ threshold, corresponding to the bound state, and $`M_a`$ above threshold corresponding to the resonance. The visible width is defined for the latter as
$$\mathrm{\Gamma }=2\mathrm{}\mathrm{Im}\mathrm{\Delta }(M_a),$$
$$\mathrm{}^1=[1+d_{SS}(M_a)][M_aM_D+\mathrm{Re}g_{DD}(M_a)]$$
$$+[1+d_{DD}(M_a)][M_aM_S+\mathrm{Re}g_{SS}(M_a)]2\mathrm{R}\mathrm{e}g_{SD}(M_A)d_{SD}(M_A),$$
$$d_{\mu \nu }(M)=\frac{\mathrm{Re}g_{\mu \nu }(M)}{M}.$$
(B.6)
The probabilities to find bare states in the wavefunction of bound state are
$$Z_S=[M_bM_D+g_{DD}(M_b)]𝒵^1,Z_D=[M_bM_S+g_{SS}(M_b)]𝒵^1,$$
(B.7)
and $`𝒵^1`$ is obtained from $`\mathrm{}^1`$ given in (B.6) with the replacement $`M_aM_b`$. The spectral densities of bare states are:
$$w_S(M)=\frac{1}{2\pi i}\left(\frac{MM_D+g_{DD}^{}}{\mathrm{\Delta }^{}(M)}\frac{MM_D+g_{DD}}{\mathrm{\Delta }(M)}\right),$$
(B.8)
and
$$w_D(M)=\frac{1}{2\pi i}\left(\frac{MM_S+g_{SS}^{}}{\mathrm{\Delta }^{}(M)}\frac{MM_S+g_{SS}}{\mathrm{\Delta }(M)}\right),$$
(B.9)
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# Squeezed Fock state by inconclusive photon subtraction
## 1 Introduction
Beam splitters (BS) and avalanche photodetectors (APDs) play a fundamental role in quantum information processing. These key elements, among the other applications, can be used in order to generate non-Gaussian states from Gaussian ones and to distill continuous-variable entanglement .
In this paper we focus our attention on the output state recently obtained experimentally by J. Wenger et al. by means of photon subtraction on a squeezed vacuum state $`S(r)|0`$, $`S(r)`$ being the squeezing operator. More precisely, when a Gaussian state, such as $`S(r)|0`$, is mixed with the vacuum at a beam splitter and, then, on/off photodetection is performed on the reflected beam, an unknown number of photons is subtracted from the input state and the output state is no longer Gaussian, i.e., the input state is de-Gaussified: this is due to the fact that the positive operator valued measure (POVM) describing the APD is non-Gaussian. Since the actual number of detected photons cannot be resolved by the APD, in we referred to this process as to inconclusive photon subtraction (IPS). In general the IPS process can be characterized by two parameters: the beam splitter transmissivity $`\tau `$ and the quantum efficiency $`\eta `$ of the APD. As we will see, the conditional output state obtained by IPS on a squeezed vacuum is close to the squeezed Fock state $`S(r)|1`$, which is otherwise difficult to produce by Hamiltonian processes. For this reason, we address IPS as an effective resource to generate those squeezed Fock states. We find that the IPS conditional state reduces to $`S(r)|1`$ in the limit $`\tau ,\eta 1`$, whereas for different values of the transmissivity and of the quantum efficiency it remains close to this target state, showing a high fidelity for a wide range of the parameters. Finally, since the IPS state obtained from the squeezed vacuum is, in general, non classical and mixed, we study how the purity and the nonclassical depth of the IPS state depend on $`\tau `$, $`\eta `$, and on the input squeezing parameter $`r`$.
The paper is structured as follows: in Section 2 we review the main elements of the IPS process on a single mode of radiation. The fidelity between the IPS conditional state and the squeezed Fock state $`S(z)|1`$, as well as its purity are then investigated in Section 3, whereas section 4 is devoted to the analysis of the nonclassicality of the IPS state. Finally, Section 5 closes the paper with some concluding remarks.
## 2 The inconclusive photon subtraction process
The scheme of the inconclusive photon subtraction (IPS) process is sketched in figure 1. An input state $`\varrho ^{(\mathrm{in})}`$ is mixed with the vacuum state $`\varrho _0=|00|`$ at a beam splitter (BS) with transmissivity $`\tau `$ and, then, on/off avalanche photodetection (APD) with quantum efficiency $`\eta `$ is performed on the reflected beam. Since the APD con only distinguish the presence from the absence of light, this measurement is inconclusive, namely does not resolve the number of the detected photons. In this way, when the detector clicks, an unknown number of photon is subtracted from the initial state and we obtain the IPS state $`\varrho ^{(\mathrm{out})}`$. Since the whole process is characterized by $`\tau `$ and $`\eta `$, we will refer to them also as IPS transmissivity and IPS quantum efficiency.
If the input state of the mode $`a`$ is the squeezed vacuum state $`\varrho _r^{(\mathrm{in})}=|0,r0,r|`$, where $`|0,r=S(r)|0`$, $`S(r)=\mathrm{exp}\{\frac{1}{2}r(a_{}^{}{}_{}{}^{2}a^2)\}`$ being the squeezing operator (for the sake of the simplicity, without lack of generality, we can assume $`r`$ as real), its (Gaussian) characteristic function $`\chi _r^{(\mathrm{in})}(𝚲_a)\chi [\varrho _r^{(\mathrm{in})}](𝚲_a)`$ reads
$$\chi _r^{(\mathrm{in})}(𝚲_a)=\mathrm{exp}\left\{\frac{1}{2}𝚲_a^T𝝈_r𝚲_a\right\}$$
(1)
where $`𝚲=(𝗑_a,𝗒_a)^T`$, $`(\mathrm{})^T`$ being the transposition operation, and
$$𝝈_r=\frac{1}{2}\left(\begin{array}{cc}\mathrm{cosh}r+\mathrm{sinh}r& 0\\ 0& \mathrm{cosh}r\mathrm{sinh}r\end{array}\right),$$
(2)
is the covariance matrix. Analogously, the vacuum state $`\varrho _0=|00|`$ of the mode $`b`$ is described by the (Gaussian) characteristic function
$$\chi _0(𝚲_b)\chi [\varrho _0](𝚲_b)=\mathrm{exp}\left\{\frac{1}{2}𝚲_b^T𝝈_0𝚲_b\right\},$$
(3)
where $`𝝈_0=\frac{1}{2}\mathrm{𝟙}_2`$, $`\mathrm{𝟙}_2`$ being the $`2\times 2`$ identity matrix. Since the initial two-mode state $`\varrho _r^{(\mathrm{in})}\varrho _0`$ is Gaussian, under the action of the BS its $`4\times 4`$ covariance matrix
$$𝝈_{\mathrm{in}}=\left(\begin{array}{cc}𝝈_r& \mathrm{𝟎}\\ & \\ \mathrm{𝟎}& 𝝈_0\end{array}\right)$$
(4)
transforms as follows
$$𝝈_{\mathrm{in}}𝝈^{}𝑺_{\mathrm{BS}}^T𝝈_{\mathrm{in}}𝑺_{\mathrm{BS}}\left(\begin{array}{cc}𝑨& 𝑪\\ & \\ 𝑪^T& 𝑩\end{array}\right),$$
(5)
where $`𝑨`$, $`𝑩`$, and $`𝑪`$ are $`2\times 2`$ matrices and
$$𝑺_{\mathrm{BS}}=\left(\begin{array}{cc}\sqrt{\tau }\mathbb{\hspace{0.17em}1}_2& \sqrt{1\tau }\mathbb{\hspace{0.17em}1}_2\\ & \\ \sqrt{1\tau }\mathbb{\hspace{0.17em}1}_2& \sqrt{\tau }\mathbb{\hspace{0.17em}1}_2\end{array}\right),$$
(6)
is the symplectic transformation associated to the evolution operator of the BS. Now, the on/off photodetector with quantum efficiency $`\eta `$ can be described by the POVM $`\{\mathrm{\Pi }_{\mathrm{off}}(\eta ),\mathrm{\Pi }_{\mathrm{on}}(\eta )\}`$, with
$$\mathrm{\Pi }_{\mathrm{off}}(\eta )=\underset{k=0}{\overset{\mathrm{}}{}}(1\eta )^k|kk|,\mathrm{\Pi }_{\mathrm{on}}(\eta )=𝕀\mathrm{\Pi }_{\mathrm{off}}(\eta ),$$
(7)
which corresponds to the characteristic functions
$`\chi [\mathrm{\Pi }_{\mathrm{off}}(\eta )](𝚲)\chi _\eta ^{(\mathrm{off})}(𝚲)={\displaystyle \frac{1}{\eta }}\mathrm{exp}\left\{\frac{1}{2}𝚲^T𝝈_\mathrm{M}𝚲\right\},`$ (8)
$`\chi [\mathrm{\Pi }_{\mathrm{on}}(\eta )](𝚲)\chi _\eta ^{(\mathrm{on})}(𝚲)=2\pi \delta ^{(2)}(𝚲)\chi _\eta ^{(\mathrm{off})}(𝚲),`$ (9)
respectively, $`\delta ^{(2)}(𝚲)`$ being the 2-dim Dirac’s delta function, and
$$𝝈_\mathrm{M}=\frac{2\eta }{2\eta }\mathbb{\hspace{0.17em}1}_2.$$
(10)
The probability of a click in the detector is then given by
$`p_{\mathrm{on}}(r,\tau ,\eta )`$ $`=\mathrm{Tr}_{ab}[\varrho _{r,\tau }^{}𝕀\mathrm{\Pi }_{\mathrm{on}}(\eta )]`$ (11)
$`={\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle _^4}d^2𝚲_ad^2𝚲_b\chi [\varrho _{r,\tau }^{}](𝚲_a,𝚲_b)\chi [𝕀](𝚲_a)\chi _\eta ^{(\mathrm{on})}(𝚲_b)`$ (12)
$`=1\left(\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}\right)^1=1\left(\sqrt{1+(1\tau _{\mathrm{eff}}^2)\mathrm{sinh}^2r}\right)^1,`$ (13)
where $`\chi [\varrho _{r,\tau }^{}](𝚲_a,𝚲_b)`$ is the two-mode characteristic function associated to the state $`\varrho _{r,\tau }^{}U_{\mathrm{BS}}\varrho _r^{(\mathrm{in})}\varrho _0U_{\mathrm{BS}}^{}`$, $`\chi [𝕀](𝚲)=2\pi \delta ^{(2)}(𝚲)`$, and $`\tau _{\mathrm{eff}}\tau _{\mathrm{eff}}(\tau ,\eta )=1\eta (1\tau )`$. Note that when $`\tau _{\mathrm{eff}}1`$, the probability (13) can be approximated at the first order in $`\tau _{\mathrm{eff}}`$ as follows
$$p_{\mathrm{on}}(r,\tau ,\eta )=(1\tau _{\mathrm{eff}})\mathrm{sinh}^2r+o\left[(1\tau _{\mathrm{eff}})^2\right].$$
(14)
Finally, the output state
$$\varrho _{r,\tau ,\eta }^{(\mathrm{out})}=\frac{\mathrm{Tr}_b[\varrho _{r,\tau }^{}𝕀\mathrm{\Pi }_{\mathrm{on}}(\eta )]}{p_{\mathrm{on}}(r,\tau ,\eta )},$$
(15)
conditioned to a click of the on/off photodetector, has the following characteristic function $`\chi _{r,\tau ,\eta }^{(\mathrm{out})}(𝚲_a)\chi [\varrho _{r,\tau ,\eta }^{(\mathrm{out})}](𝚲_a)`$:
$`\chi _{r,\tau ,\eta }^{(\mathrm{out})}(𝚲_a)`$ $`={\displaystyle \frac{1}{2\pi p_{\mathrm{on}}(r,\tau ,\eta )}}{\displaystyle _^2}d^2𝚲_b\chi [\varrho _{r,\tau }^{}](𝚲_a,𝚲_b)\chi _\eta ^{(\mathrm{on})}(𝚲_b)`$ (16)
$`={\displaystyle \frac{1}{p_{\mathrm{on}}(r,\tau ,\eta )}}\left\{\mathrm{exp}\left\{\frac{1}{2}𝚲_a^T𝚺_1𝚲_a\right\}{\displaystyle \frac{\mathrm{exp}\left\{\frac{1}{2}𝚲_a^T𝚺_2𝚲_a\right\}}{\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}}}\right\},`$ (17)
with $`𝚺_1=𝑨`$ and $`𝚺_2=𝑨𝑪(𝑩+𝝈_\mathrm{M})^1𝑪^T`$. Note that the output state is no longer a Gaussian state, namely its characteristic function is no longer Gaussian: for this reason the IPS process is also referred to as de-Gaussification process .
In general, a Gaussian state described by the characteristic function \[in Cartesian notation, namely $`𝚲=(𝗑,𝗒)^T`$\]
$$\chi (𝚲)=\mathrm{exp}\left\{\frac{1}{2}𝚲^T𝝈𝚲\right\}$$
(18)
with covariance matrix
$$𝝈=\left(\begin{array}{cc}𝖺& 𝖼\\ 𝖼& 𝖻\end{array}\right),$$
(19)
can be also written in the complex notation as follows:
$$\chi (\lambda )=\mathrm{exp}\left\{𝒜|\lambda |^2\lambda ^2^{}\lambda _{}^{}{}_{}{}^{2}\right\},$$
(20)
with
$$𝒜=\frac{1}{2}(𝖺+𝖻),=\frac{1}{4}(𝖻𝖺+2i𝖼),$$
(21)
where we introduced the complex number $`\lambda =\frac{1}{\sqrt{2}}(𝗑+i𝗒)`$. In this way, the characteristic function (17) can be written as follows:
$$\chi _{r,\tau ,\eta }^{(\mathrm{out})}(\lambda )=\frac{\mathrm{exp}\left\{𝒜_1|\lambda |^2_1\lambda ^2_1^{}\lambda _{}^{}{}_{}{}^{2}\right\}}{p_{\mathrm{on}}(r,\tau ,\eta )}\frac{\mathrm{exp}\left\{𝒜_2|\lambda |^2_2\lambda ^2_2^{}\lambda _{}^{}{}_{}{}^{2}\right\}}{p_{\mathrm{on}}(r,\tau ,\eta )\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}},$$
(22)
where $`𝒜_k`$ and $`_k`$ are refers to the covariance matrix $`𝚺_k`$, $`k=1,2`$ respectively. Finally, using the definition
$$W[\varrho ](\alpha )=\frac{1}{\pi ^2}_{}d^2\lambda \chi [\varrho ](\lambda )\mathrm{exp}\left\{\lambda ^{}\alpha \alpha ^{}\lambda \right\},$$
(23)
which relates the Wigner function $`W[\varrho ](\alpha )`$ of a state $`\varrho `$ to its characteristic function $`\chi [\varrho ](\lambda )`$, one can obtain the Wigner function $`W_{r,\tau ,\eta }^{(\mathrm{out})}(\alpha )W[\varrho _{r,\tau ,\eta }^{(\mathrm{out})}](\alpha )`$. As for the characteristic function, to pass from the complex, $`W[\varrho ](\alpha )`$, to the Cartesian notation, $`W[\varrho ](x,y)`$, one should put $`\alpha =\frac{1}{\sqrt{2}}(x+iy)`$ . In figure 2 (a) we report $`W_{r,\tau ,\eta }^{(\mathrm{out})}(x,y)`$ for fixed $`r`$, $`\tau `$, and $`\eta `$: as it is apparent from the plot the Wigner function is not Gaussian, and may assume negative values . In Section 4 we will investigate this effect by analyzing the nonclassicality of the conditioned state. In figure 2 (b) we show the Wigner function $`\chi _z^{(\mathrm{SqF})}(x,y)`$ associated to the squeezed Fock state $`\varrho _z^{(\mathrm{SqF})}=S(z)|11|S^{}(z)`$, whose characteristic function $`\chi _z^{(\mathrm{SqF})}(\lambda )\chi [\varrho _z^{(\mathrm{SqF})}](\lambda )`$ reads (we assume $`z`$ as real)
$$\chi _z^{(\mathrm{SqF})}(\lambda )=\left[12\left(𝒜_0|\lambda |^2+_0\lambda ^2+_0^{}\lambda _{}^{}{}_{}{}^{2}\right)\right]\mathrm{exp}\left\{𝒜_0|\lambda |^2_0\lambda ^2_0^{}\lambda _{}^{}{}_{}{}^{2}\right\},$$
(24)
with $`𝒜_0=2(\mathrm{cosh}^2z+\mathrm{sinh}^2z)`$ and $`_0=2\mathrm{cosh}z\mathrm{sinh}z`$. Since the Wigner functions of the IPS squeezed vacuum and of the squeezed number state are quite similar, one can think of using the IPS process to produce the state $`\varrho _r^{(\mathrm{SqF})}`$; motivated by this consideration, in the next section we will analyses the fidelity between this states.
Figure 3 shows $`W_{r,\tau ,\eta }^{(\mathrm{out})}(x,y)`$ with fixed $`r`$ and $`\eta `$ and different values of the IPS transmissivity $`\tau `$; the plots on the right of the same figure compare the $`W_{r,\tau \eta }^{(\mathrm{out})}(0,y)`$ (solid lines) with $`W_r^{(\mathrm{SqF})}(0,y)`$ (dashed line).
Finally, the effect of the quantum efficiency $`\eta `$ on the output state is shown in figure 4, where we plot as reference the value of the Wigner function $`W_{r,\tau ,\eta }^{(\mathrm{out})}`$ at the center of the complex plane as a function of the transmissivity $`\tau `$ and different values of $`\eta `$: we can see that the main effect on the output state is due to $`\tau `$.
## 3 Fidelity and purity
The fidelity between the pure state $`\varrho _z^{(\mathrm{SqF})}`$ and the IPS state $`\varrho _{r,\tau ,\eta }^{(\mathrm{out})}`$ is defined as follows:
$`F_{\tau ,\eta }(z,r)`$ $`=\mathrm{Tr}[\varrho _z^{(\mathrm{SqF})}\varrho _{r,\tau ,\eta }^{(\mathrm{out})}]`$ (25)
$`={\displaystyle \frac{1}{2\pi }}{\displaystyle _^2}d^2𝚲\chi _z^{(\mathrm{SqF})}(𝚲)\chi _{r,\tau ,\eta }^{(\mathrm{out})}(𝚲),`$ (26)
$`={\displaystyle \frac{1}{p_{\mathrm{on}}(r,\tau ,\eta )}}\{_1{\displaystyle \frac{_2}{\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}}},\}`$ (27)
where
$$_k=\frac{𝒜_k^2𝒜_0^24(_k^2_0^2)}{[(𝒜_0+𝒜_k)^24(_0+_k)^2]^{3/2}}$$
(28)
and $`𝒜_h`$ and $`_h`$, $`h=0,1,2`$, have been introduced in equations (24) and (22), respectively. The analytic expression of $`F_{\tau ,\eta }(z,r)`$ is quite cumbersome, but, on the other hand, we can draw some interesting consideration by addressing its expansion at the first order in the transmissivity $`\tau `$ when $`\tau 1`$ and $`\eta =1`$, namely
$`F_{\tau ,1}(z,r)={\displaystyle \frac{1}{\mathrm{cosh}^3(rz)}}`$
$`\left[{\displaystyle \frac{9\mathrm{cosh}(r+z)3\mathrm{cosh}(3rz)}{8\mathrm{cosh}(rz)}}{\displaystyle \frac{1}{4}}\right](1\tau )+o\left[(1\tau )^2\right].`$ (29)
In fact, from the expansion (3) we conclude that the maximum of the fidelity is achieved when $`z=r`$.
In figure 5 we plot $`F_{\tau ,\eta }(r)F_{\tau ,\eta }(r,r)`$ as a function of the IPS transmissivity and for different values of $`r`$. We can see that $`F_{\tau ,\eta }`$ reaches it maximum when the IPS transmissivity approaches $`1`$, namely in the single-photon subtraction limit . Moreover, when the squeezing parameter $`r`$ increases the fidelity decreases: this is due to the increasing (unknown) number of subtracted photons which reduces the purity of the IPS state itself. In figure 6 we plot the purity $`\mu _{\tau ,\eta }(r)`$ of the IPS squeezed vacuum $`\varrho _{r,\tau ,\eta }^{(\mathrm{out})}`$, defined as follows :
$`\mu _{\tau ,\eta }(r)=\mathrm{Tr}\left[(\varrho _{r,\tau ,\eta }^{(\mathrm{out})})^2\right]=\pi {\displaystyle _{}}d^2\alpha \left[W_{r,\tau ,\eta }^{(\mathrm{out})}(\alpha )\right]^2`$ (30)
$`={\displaystyle \frac{1}{2p_{\mathrm{on}}(r,\tau ,\eta )}}\{{\displaystyle \frac{1}{\sqrt{𝒜_1^24_1^2}}}+{\displaystyle \frac{1}{\eta ^2\mathrm{Det}[𝑩+𝝈_\mathrm{M}]\sqrt{𝒜_1^24_1^2}}}`$
$`{\displaystyle \frac{4}{\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}\sqrt{(𝒜_1+𝒜_2)^24(_1+_2)^2}}}\}.`$ (31)
## 4 Nonclassicality of the IPS squeezed vacuum state
As a measure of nonclassicality of the IPS state $`\varrho _{r,\tau ,\eta }^{(\mathrm{out})}`$ we consider the nonclassical depth
$$𝒯=\frac{1\overline{s}}{2},$$
(32)
$`\overline{s}`$ being the maximum $`s`$ for which the generalized quasi-probability function
$$W_s(\alpha )=\frac{1}{\pi }_{}d^2\lambda \chi (\lambda )\mathrm{exp}\left\{\frac{1}{2}s+\lambda ^{}\alpha \alpha ^{}\lambda \right\}$$
(33)
is a probability distribution, i.e. positive semidefinite and non singular. As a matter of fact, one has $`𝒯=1`$ for number states and $`𝒯=0`$ for coherent states. Moreover, the nonclassical depth can be interpreted as the minimum number of thermal photons which has to be added to a quantum state in order to erase all the quantum features of the state . In the case of $`\varrho _{r,\tau ,\eta }^{(\mathrm{out})}`$, we have \[for the sake of simplicity we do not write explicitly the dependence on $`r`$,$`\tau `$ and $`\eta `$ in the symbol $`W_s^{(\mathrm{out})}(\alpha )`$\]
$$W_s^{(\mathrm{out})}(\alpha )=\frac{1}{p_{\mathrm{on}}(r,\tau ,\eta )}\left\{𝒢_1(\alpha )\frac{𝒢_2(\alpha )}{\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}}\right\}$$
(34)
where we defined
$$𝒢_k(\alpha )=\frac{2\mathrm{exp}\left\{{\displaystyle \frac{2(2𝒜_ks)|\alpha |^2+4_k^{}\alpha ^2+4_k\alpha _{}^{}{}_{}{}^{2}}{(2𝒜_ks)^216|_k|^2}}\right\}}{\pi \sqrt{(2𝒜_ks)^216|_k|^2}}.$$
(35)
At first we note that in order to have $`W_s^{(\mathrm{out})}(\alpha )`$ normalizable, equation the following condition should be satisfied
$$s2𝒜_k(k=1,2).$$
(36)
Furthermore, since $`W_s^{(\mathrm{out})}(\alpha )`$ is a difference between two Gaussian functions with the center in the origin of the complex plane, one can easily see that, in general, this function has a minimum in $`\alpha =0`$ and that this minimum can be negative. For this reason and thanks to other simple considerations about the symmetries of $`W_s^{(\mathrm{out})}(\alpha )`$ with respect to the point $`\alpha =0`$, we can focus our attention in the origin of the complex plane, obtaining this further condition for the positivity:
$$𝒢_1(0)\frac{𝒢_2(0)}{\eta \sqrt{\mathrm{Det}[𝑩+𝝈_\mathrm{M}]}}0,$$
(37)
which, together with the conditions (36), brings to
$$\overline{s}(\tau ,\eta )=\frac{2\eta (4\eta )\tau }{2(1\tau )\eta },$$
(38)
and, then, to the following expression for the nonclassical depth:
$$𝒯(\tau ,\eta )=\frac{2\tau }{2(1\tau )\eta }.$$
(39)
Since $`𝒯(\tau ,\eta )0`$, the conditional state is nonclassical for any non-zero value of the IPS transmissivity and efficiency. Note that equation (39) depends only on $`\tau `$ and $`\eta `$, whereas it is independent on the squeezing parameter $`r`$. Notice, however, the nonclassical depth does not measure the extension of the negativity region, but only the presence of negative values. Therefore it is not surprising that equation (39) does not depend on $`r`$. We plot $`𝒯(\tau ,\eta )`$ in figure 7. Since the usual Wigner function is obtained when $`s=0`$ in (33), from equation (38) we can see that $`W_{r,\tau ,\eta }^{(\mathrm{out})}(\alpha )`$ becomes semi-positive definite when $`\tau =(2\eta )/(4\eta )`$.
## 5 Concluding remarks
We have analyzed in details the state obtained subtracting photons from the squeezed vacuum by means of linear optics, namely using beam splitters and avalanche photodetectors. We referred to the whole photon-subtraction process as to inconclusive photon subtraction (IPS), since avalanche photodetectors are not able to resolve the number of detected photons. We found that the IPS conditional state obtained from a squeezed vacuum state is close to the squeezed Fock state $`S(r)|1`$ and approaches this target state when only one photon is subtracted, namely, using a high transmissivity beam splitter for the IPS. Moreover, when the transmissivity and the quantum efficiency are not unitary, the output state remains close to the target state, showing a high fidelity for a wide range of the parameters. The purity and the nonclassicality of the IPS squeezed vacuum state have been also considered: we found that the relevant parameter is the transmissivity $`\tau `$, while the IPS efficiency $`\eta `$ only slightly affects the output state. We conclude that IPS, which was recently experimentally implemented , can be effectively used to produce a nonclassical state such as the squeezed Fock state $`S(r)|1`$, whose generation would be, otherwise, quite challenging.
## References
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# A new approach to the Darboux-Bäcklund transformation versus the standard dressing method
## 1 Introduction
There are several methods to construct the Darboux matrix (which generates soliton solutions) ). However, these methods are technically difficult when applied to the matrix versions of the spectral problems which are naturally represented in Clifford algebras . Some of these problems are avoided in our recent paper . In the present paper we develop the ideas of in the matrix case. We extend our approach on the multisoliton case and consider the reduction problem and the discrete case. We also show that our approach, although different, is to some extent equivalent to the standard dressing method. We compare our method with the Zakharov-Shabat approach and the Neugebauer-Meinel approach .
We consider the spectral problem
$$\mathrm{\Psi }_{,\mu }=U_\mu \mathrm{\Psi },(\mu =1,\mathrm{},m)$$
(1)
(with no assumptions on $`U_\mu `$ except rational dependence on $`\lambda `$) and the Darboux transformation
$$\stackrel{~}{\mathrm{\Psi }}=D\mathrm{\Psi },$$
(2)
which means that
$$\stackrel{~}{\mathrm{\Psi }},_\mu =\stackrel{~}{U}_\mu \stackrel{~}{\mathrm{\Psi }},$$
(3)
where $`\stackrel{~}{U}_\mu `$ and $`U_\mu `$ have the same rational dependence on $`\lambda `$ ($`U_\mu `$ and $`\mathrm{\Psi }`$ are $`n\times n`$ matrices but our approach works well also in the Clifford numbers case ).
The construction of the Darboux transformation is well known (especially in the matrix case) . The first step is the equation for $`D`$ resulting from (1),(2) and (3):
$$D,_\mu +DU_\mu =\stackrel{~}{U}_\mu D.$$
(4)
In our erlier paper we proposed the following procedure. We assume that there exist two different values of $`\lambda `$, say $`\lambda _+`$ and $`\lambda _{}`$, satisfying
$$D^2(\lambda _\pm )=0.$$
(5)
Denoting $`\mathrm{\Psi }(\lambda _\pm )=\mathrm{\Psi }_\pm `$, $`D(\lambda _\pm )=D_\pm `$, evaluating (4) at $`\lambda =\lambda _\pm `$ and multiplying (4) by $`D_\pm `$ from the right, we get:
$$D_\pm ,_\mu D_\pm +D_\pm U_\mu (\lambda _\pm )D_\pm =0.$$
(6)
We assume that $`\mathrm{\Psi }(\lambda _\pm )`$ are invertible (which is obviously true in the generic case). It is not difficult to check that $`D_\pm `$ given by
$$D_\pm =\phi _\pm \mathrm{\Psi }_\pm d_\pm \mathrm{\Psi }_\pm ^1,d_\pm ^2=0,$$
(7)
(where $`d_\pm =\mathrm{const}`$ and $`\phi _\pm `$ are scalar functions) satisfy equations (5), (6). Assuming that $`D`$ is linear in $`\lambda `$, i.e.,
$$D(\lambda )=A_0+A_1\lambda ,$$
(8)
we can easily express $`A_0,A_1`$ by $`D_\pm `$ to get
$$D(\lambda )=\frac{\lambda \lambda _{}}{\lambda _+\lambda _{}}\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1+\frac{\lambda \lambda _+}{\lambda _{}\lambda _+}\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1.$$
(9)
## 2 One-soliton case and the Zakharov-Shabat approach
We confine ourselves to the case linear in $`\lambda `$ (see (8)). The condition (5) can be easily realized if
$$D^2(\lambda )=\sigma (\lambda \lambda _+)(\lambda \lambda _{})I$$
(10)
where $`\sigma 0`$ is a constant, $`\lambda _+\lambda _{}`$ and $`I`$ is the identity matrix. The identity matrix will be sometimes omitted (i.e., for $`a𝐂`$ we write $`aI=a`$). In the case (10) from (5) and (9) it follows that
$$D_+D_{}+D_{}D_+=\sigma (\lambda _+\lambda _{})^2.$$
(11)
###### Lemma 1
$`D`$ of the form (8) satisfies (10) if and only if $`n`$ is even and
$$D=𝒩\left(\lambda \lambda _++(\lambda _+\lambda _{})P\right)$$
(12)
where the matrices $`𝒩`$ and $`P`$ satisfy
$$P^2=P,𝒩^2=\sigma ,𝒩P𝒩^1=IP.$$
(13)
In this case the Darboux matrices (9) and (12) are equivalent.
Proof: We denote $`𝒩:=A_1`$. From (8) we get
$$D^2\left(\lambda \right)=A_0^2+\left(A_0𝒩+𝒩A_0\right)\lambda +𝒩^2\lambda ^2,$$
i.e., $`D^2(\lambda )`$ is a quadratic polynomial. It is proportional to the identity matrix $`I`$ (compare (10)) iff
$$𝒩^2=\sigma ,A_0𝒩+𝒩A_0=\sigma \left(\lambda _++\lambda _{}\right),A_0^2=\sigma \lambda _+\lambda _{}.$$
(14)
Multiplying the second equation by $`𝒩A_0`$ we get
$$\sigma ^2\lambda _+\lambda _{}+\left(𝒩A_0\right)^2+\sigma \left(\lambda _++\lambda _{}\right)𝒩A_0=0.$$
Hence $`\left(𝒩A_0+\sigma \lambda _+\right)\left(𝒩A_0+\sigma \lambda _{}\right)=0`$, and, denoting $`Q:=𝒩A_0+\sigma \lambda _+`$, we have
$$Q^2=\left(\lambda _+\lambda _{}\right)\sigma Q$$
which means that $`Q=\left(\lambda _+\lambda _{}\right)\sigma P`$, where $`P^2=P`$. Therefore, taking into account $`𝒩^2=\sigma `$, we get (12). Now, we take into account the third equation of (14). First, $`A_0^2P=\sigma \lambda _+\lambda _{}P`$ yields $`\lambda _{}\left(\lambda _+\lambda _{}\right)𝒩P𝒩P=0`$. Then the equation $`A_0^2=\sigma \lambda _+\lambda _{}`$ is equivalent to $`\lambda _+\left(\lambda _+\lambda _{}\right)\left(\sigma \left(IP\right)𝒩P𝒩\right)=0`$. Therefore $`𝒩P𝒩^1=IP`$. This equality means that $`\mathrm{ker}P=𝒩^1\mathrm{Im}P`$ which implies $`dim\mathrm{ker}P=dim\mathrm{Im}P`$. Thus $`n`$ is even which complets the proof. $`\mathrm{}`$
The case $`\lambda _+=\lambda _{}`$ can be treated in a similar way and it leads to the nilpotent case :
$$D=𝒩(\lambda \lambda _++M),M^2=0,𝒩^2=\sigma ,M=𝒩M𝒩^1.$$
Our method is closely related to the standard dressing transformation . The Darboux matrix (12) can be rewritten as
$$D=\left(\lambda \lambda _+\right)𝒩\left(I+\frac{\lambda _+\lambda _{}}{\lambda \lambda _+}P\right).$$
(15)
We recognize the standard one-soliton Darboux matrix in the Zakharov-Shabat form . We point out that usually one considers the Darboux matrix $`𝒟=(\lambda \lambda _+)^1D`$ which is equivalent to $`D`$ given by (12) because the multiplication of $`D`$ by a constant factor leaves the equation (4) invariant . $`𝒩`$ is known as the normalization matrix and $`P`$ is a projector expressed by the background wave function:
$$\mathrm{ker}P=\mathrm{\Psi }(\lambda _+)V_{ker},\mathrm{im}P=\mathrm{\Psi }(\lambda _{})V_{im},$$
(16)
$`V_{ker}`$ and $`V_{im}`$ are some constant vector spaces, $`\lambda _+`$ and $`\lambda _{}`$ are constant complex parameters. The last constraint of (13) has the following interpretation. Let $`𝒩P𝒩^1=IP`$. Then
$$\begin{array}{c}v\mathrm{im}P(IP)v=0P𝒩^1v=0𝒩^1v\mathrm{ker}P\hfill \\ v\mathrm{ker}PPv=0P𝒩^1v=𝒩^1v𝒩^1v\mathrm{im}P\hfill \end{array}$$
Hence, $`dim\mathrm{im}P=dim\mathrm{ker}P=dn/2`$, which implies $`dimV_{im}=dimV_{ker}`$. In this case, given a projector $`P`$, one can always find a corresponding $`𝒩`$. Indeed, let $`v_1,\mathrm{},v_d`$ be a basis in $`\mathrm{im}P`$ and $`w_k:=𝒩^1v_k`$ ($`k=1,\mathrm{},d`$) an associated basis in $`\mathrm{ker}P`$. By virue of $`𝒩^2=\sigma `$ we have $`𝒩^1w_k=\sigma ^1v_k`$. Therefore
$$𝒩^1(v_1,\mathrm{},v_d,w_1,\mathrm{},w_d)=(w_1,\mathrm{},w_d,v_1/\sigma ,\mathrm{},v_d/\sigma )$$
(where $`(v_1,v_2,\mathrm{})`$ denotes the matrix with columns $`v_1,v_2,\mathrm{}`$) and, finally,
$$𝒩=(v_1,\mathrm{},v_d,w_1,\mathrm{},w_d)(w_1,\mathrm{},w_d,v_1/\sigma ,\mathrm{},v_d/\sigma )^1.$$
(17)
The $`𝒩`$ obtained in this way depends on the choice of the bases $`v_1,\mathrm{},v_d`$ and $`w_1,\mathrm{},w_d`$ (we can put $`Av_k`$, $`detA0`$, in the place of $`v_k`$ and $`Bw_j`$, $`detB0`$, in the place of $`w_j`$). In other words, $`𝒩`$ is given up to nondegenerate $`d\times d`$ matrices $`A`$ and $`B`$.
The formulas (9) and (12) coincide after appropriate identification of the parameters. Indeed, comparing coefficients by powers of $`\lambda `$ we have:
$$\begin{array}{c}𝒩=\frac{\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1}{\lambda _+\lambda _{}},\hfill \\ 𝒩(\lambda _++(\lambda _+\lambda _{})P)=\frac{\lambda _+\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1\lambda _{}\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1}{\lambda _+\lambda _{}},\hfill \end{array}$$
(18)
and after straightforward computation we get
$$\begin{array}{c}P=(\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1)^1\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1,\hfill \\ IP=(\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1\phi _+\mathrm{\Psi }_+d_+\mathrm{\Psi }_+^1)^1\phi _{}\mathrm{\Psi }_{}d_{}\mathrm{\Psi }_{}^1.\hfill \end{array}$$
(19)
Taking into account the assumption (11) we have:
$$P=\frac{D_{}D_+}{D_+D_{}+D_{}D_+}=\frac{D_{}D_+}{\sigma (\lambda _+\lambda _{})^2}.$$
(20)
The above results are valid for $`n\times n`$ matrix linear problems. Now, we focus on the $`2\times 2`$ case. Because the elements $`d_+`$, $`d_{}`$ are nilpotent ($`d_\pm =0`$), then there exist vectors $`v_+,v_{}`$ such that
$$d_+v_+=0,d_{}v_{}=0.$$
(21)
Then from (19) it follows immediately $`P\mathrm{\Psi }_+v_+=0`$ and $`(IP)\mathrm{\Psi }_{}v_{}=0`$, i.e., $`\mathrm{\Psi }_+v_+`$ span $`\mathrm{ker}P`$ and $`\mathrm{\Psi }_{}v_{}`$ span $`\mathrm{im}P`$. Hence, $`v_+V_{ker}`$ and $`v_{}V_{im}`$.
It is not difficult to check that the general form of $`2\times 2`$ matrices $`d_\pm `$ such that $`d_\pm ^2=0`$ is given by
$$d_\pm =\left(\begin{array}{cc}a_\pm b_\pm & b_\pm ^2\\ a_\pm ^2& a_\pm b_\pm \end{array}\right)=\left(\begin{array}{c}b_\pm \\ a_\pm \end{array}\right)\left(\begin{array}{cc}a_\pm & b_\pm \end{array}\right),$$
(22)
where $`a_\pm ,b_\pm `$ are complex numbers. Therefore, to satisfy (21), we can take
$$v_+=\left(\begin{array}{c}\hfill b_+\\ \hfill a_+\end{array}\right),v_{}=\left(\begin{array}{c}\hfill b_{}\\ \hfill a_{}\end{array}\right).$$
(23)
We have almost unique correspondence (i.e., up to a scalar factor) between $`v_+`$ and $`d_+`$ and between $`v_{}`$ and $`d_{}`$.
Denoting
$$\mathrm{\Psi }_+v_+\left(\begin{array}{c}B_+\\ A_+\end{array}\right),\mathrm{\Psi }_{}v_{}\left(\begin{array}{c}B_{}\\ A_{}\end{array}\right),$$
we get the explicit formula for $`P`$
$$P=\left(\begin{array}{cc}0& B_{}\\ 0& A_{}\end{array}\right)\left(\begin{array}{cc}B_+& B_{}\\ A_+& A_{}\end{array}\right)^1=\frac{\left(\begin{array}{cc}A_+B_{}& B_+B_{}\\ A_+A_{}& B_+A_{}\end{array}\right)}{A_{}B_+A_+B_{}}$$
(24)
The corresponding $`𝒩`$ reads (compare (17)):
$$𝒩=\frac{1}{A_{}B_+A_+B_{}}\left(\begin{array}{cc}\sigma A_{}B_{}A_+B_+& B_+^2\sigma B_{}^2\\ \sigma A_{}^2A_+^2& A_+B_+\sigma A_{}B_{}\end{array}\right)$$
(25)
Although we can reduce our approach to the explicit formulas (24) and (25) the main advantage of our method consists in expressing the Darboux transformation in terms of $`\mathrm{\Psi }_\pm d_\pm \mathrm{\Psi }_\pm ^1`$ and avoiding difficulties with parameterizing kernel and image of the projector $`P`$ which is especially troublesome in the Clifford algebras case.
## 3 Reductions
Let us consider the unitary reduction
$$U_\mu ^{}(\overline{\lambda })=U_\mu (\lambda ).$$
(26)
If $`U_\mu `$ is a polynom in $`\lambda `$, then the condition (26) means that the coefficients of this polynom by powers of $`\lambda `$ are $`u(n)`$-valued.
One can easily prove that (26) implies $`\mathrm{\Psi }^{}(\overline{\lambda })\mathrm{\Psi }(\lambda )=C(\lambda )`$, where $`C(\lambda )`$ is a constant matrix ($`C,_\nu =0`$). The matrix $`C`$ can be fixed by a choice of the initial conditions. Usually we confine ourselves to the case
$$\mathrm{\Psi }^{}(\overline{\lambda })\mathrm{\Psi }(\lambda )=k(\lambda )I,$$
(27)
where $`k(\lambda )`$ is analytic in $`\lambda `$. From (27) we can derive $`\overline{k(\overline{\lambda })}=k(\lambda )`$. By virtue of (2), the Darboux matrix have to satisfy the analogical constraint:
$$D^{}(\overline{\lambda })D(\lambda )=p(\lambda )I.$$
(28)
Assuming that $`D`$ is a polynom with respect to $`\lambda `$, compare (8), we get that $`p(\lambda )`$ is a polynom with constant real coefficients, i.e., $`\overline{p(\overline{\lambda })}=p(\lambda )`$ and $`p,_\nu =0`$.
###### Lemma 2
If $`D`$ is linear in $`\lambda `$ and (28) holds, then roots of the equation $`detD(\lambda )=0`$ satisfy the quadratic equation $`p(\lambda )=0`$.
Proof: Let $`p(\lambda )=\alpha \lambda ^2+\beta \lambda +\gamma `$. From (8), (28) it follows
$$A_0^{}A_0=\gamma ,A_1^{}A_1=\alpha ,A_0^{}A_1+A_1^{}A_0=\beta $$
(29)
which can be easily reduced to a single equation for $`S:=A_0A_1^1`$. Namely,
$$\alpha S^2+\beta S+\gamma =0.$$
(30)
Therefore, the eigenvalues of $`S`$ have to satisfy the equation $`p\left(\lambda \right)=0`$. Indeed, if $`S\stackrel{}{v}=\mu \stackrel{}{v}`$, then $`\left(\alpha \mu ^2+\beta \mu +\gamma \right)\stackrel{}{v}=0`$. On the other hand, the equation $`detD\left(\lambda \right)=0`$ can be rewritten as
$$0=det\left(\lambda IS\right)detA_1,$$
(31)
which means that the roots of $`detD\left(\lambda \right)=0`$ coincide with eigenvalues of $`S`$. $`\mathrm{}`$
###### Lemma 3
We assume (10). Then the reduction (27) imposes the following constraints on the Darboux matrix (9):
$$\lambda _{}=\lambda _+^{},d_{}^{}d_+=0,$$
(32)
and (for $`n=2`$) $`v_+v_{}=0`$.
In particular, by virtue of (5), we can take $`d_{}=fd_+^{}`$, where $`f`$ is a scalar function.
Proof: Let us denote zeros of the polynom $`p(\lambda )`$ by $`\lambda _1`$, $`\lambda _2`$. Because $`\overline{p(\overline{\lambda })}=p(\lambda )`$ there are two possibilities: either $`\lambda _2=\overline{\lambda }_1`$ or $`\lambda _1`$, $`\lambda _2`$ are real. From (10) we have
$$\left(detD\left(\lambda \right)\right)^2=\sigma ^n\left(\lambda \lambda _+\right)^n\left(\lambda \lambda _{}\right)^n.$$
(33)
Therefore, in the case (10), Lemma 2 means that $`\lambda _+`$, $`\lambda _{}`$ coincide with $`\lambda _1`$, $`\lambda _2`$.
Suppose that $`\lambda _+𝐑`$. Then from (28) we have $`(D(\lambda _+))^{}D(\lambda _+)=0`$ which implies $`D_+D(\lambda _+)=0`$ (because for any vector $`v𝐂^n`$ the scalar product $`vD_+^{}D_+v=0`$, hence $`D_+vD_+v=0`$, and, finally $`D_+v=0`$). Therefore $`\lambda _+`$ (and, similarly, $`\lambda _{}`$) cannot be real. Thus $`\lambda _{}=\lambda _+^{}`$. In this case (28) reads
$$\left(D\left(\lambda _{}\right)\right)^{}D\left(\lambda _+\right)=0.$$
(34)
Using (7) and (27) (assuming $`k\left(\lambda _\pm \right)0`$) we get
$$\left(D\left(\lambda _{}\right)\right)^{}=\overline{\phi }_{}\left(\mathrm{\Psi }_{}^{}\right)^1d_{}^{}\mathrm{\Psi }_{}^{}=\overline{\phi }_{}\mathrm{\Psi }_+d_{}^{}\mathrm{\Psi }_+^1$$
and (34) assumes the form $`\phi _+\overline{\phi }_{}\mathrm{\Psi }_+d_{}^{}d_+\mathrm{\Psi }_+^1=0`$. Hence $`d_{}^{}d_+=0`$.
Finally, in the case $`n=2`$, we use (22). Then the condition $`d_{}^{}d_+=0`$ is equivalent to $`a_+\overline{a}_{}+b_+\overline{b}_{}=0`$, i.e., $`v_+v_{}=0`$. $`\mathrm{}`$
Another very popular reduction is given by
$$U_\mu (\lambda )=JU_\mu (\lambda )J^1,J^2=c_0I,$$
(35)
then one can prove that $`\mathrm{\Psi }(\lambda )=J\mathrm{\Psi }(\lambda )C(\lambda )`$, and we choose such initial conditions that $`C(\lambda )=J^1`$, i.e.,
$$\mathrm{\Psi }(\lambda )=J\mathrm{\Psi }(\lambda )J^1,D(\lambda )=JD(\lambda )J^1.$$
(36)
Such choice of $`C(\lambda )`$ is motivated by a natural requirement that $`\mathrm{\Psi },\stackrel{~}{\mathrm{\Psi }},D`$ are elements of the same loop group (by the way, the formula (27) has the same motivation).
###### Lemma 4
We assume (10). Then the reduction (36) imposes the following constraints on the Darboux matrix (9):
$$\lambda _{}=\lambda _+,\phi _+=\phi _{},d_+=J^1d_{}J,$$
(37)
and (for $`n=2`$) $`v_{}=Jv_+`$.
Proof: From (36) it follows that $`detD(\lambda )=detD(\lambda )`$ which means that the set of roots of the equation $`detD(\lambda )=0`$ is invariant under the transformation $`\lambda \lambda `$. Therefore $`\lambda _{}=\lambda _+`$. Then, using once more (36) we get $`D_{}=JD_+J^1`$ and $`\mathrm{\Psi }_{}=J\mathrm{\Psi }_+J^1`$. Hence $`\phi _+d_+=\phi _{}J^1d_{}J`$. Thus $`\phi _+=c_0\phi _{}`$, where $`c_0`$ is a constant. Without loss of the generality we can take $`c_0=1`$ (redefining $`d_\pm `$ if necessary). In the case $`n=2`$ the kernels of $`d_\pm `$ are 1-dimensional. Therefore $`0=d_+v_+=J^1d_{}Jv_+`$ implies $`v_{}=c_1Jv_+`$, where $`c_1=\mathrm{const}`$. We can take $`v_+=Jv_{}`$. $`\mathrm{}`$
Other types of reductions (compare ) can be treated in a similar way.
## 4 The multi-soliton Darboux matrix
In this section we generalize the approach of . First, we relax the assumption (5). Second, we consider the $`N`$-soliton case (the Darboux matrix is a polynom of order $`N`$):
$$D(\lambda )=A_0+A_1\lambda +\mathrm{}A_N\lambda ^N.$$
(38)
The condition (5) will be replaced by:
$$D(\lambda _k)T(\lambda _k)=0$$
(39)
We denote $`D_kD(\lambda _k)`$, $`T_kT(\lambda _k)`$, $`\mathrm{\Psi }_k\mathrm{\Psi }(\lambda _k)`$ and $`U_{k\mu }U_\mu (\lambda _k)`$. Evaluating (4) at $`\lambda =\lambda _k`$ and multiplying the resulting equation by $`T_k`$ from the right we get:
$$D_k,_\mu T_k+D_kU_{k\mu }T_k=0$$
(40)
To solve the equation (40) we define $`d_k`$ and $`h_k`$ by
$$D_k=\mathrm{\Psi }_kd_k\mathrm{\Psi }_k^1,T_k=\mathrm{\Psi }_kh_k\mathrm{\Psi }_k^1$$
(41)
$$D_k,_\mu =\mathrm{\Psi }_k,_\mu d_k\mathrm{\Psi }_k^1+\mathrm{\Psi }_kd_k,_\mu \mathrm{\Psi }_k^1\mathrm{\Psi }_kd_k\mathrm{\Psi }_k^1\mathrm{\Psi }_k,_\mu \mathrm{\Psi }_k^1.$$
Therefore
$$D_k,_\mu =U_{k\mu }D_k+\mathrm{\Psi }_kd_k,_\mu \mathrm{\Psi }_k^1D_kU_{k\mu },$$
and, taking into account (39) and (41), we rewrite (40) as follows
$$\mathrm{\Psi }_kd_k,_\mu h_k\mathrm{\Psi }_k^1=0.$$
(42)
Finally, as a straightforward consequence of (39) and (42) we get the following constraints on $`d_k`$ and $`h_k`$:
$$d_kh_k=0,d_kh_k,_\mu =0.$$
(43)
In we confined ourselves to the case $`T(\lambda )=D(\lambda )`$, i.e, $`d_k=\phi _kd_{0k}`$, ($`\phi _k`$ scalar functions, $`d_{0k}`$ constant elements satisfying $`d_{0k}^2=0`$), $`h_k=d_k`$. Now we are going to obtain the general solution of (43) in the case of $`2\times 2`$ matrices.
###### Lemma 5
Let $`d`$ and $`h`$ are $`2\times 2`$ matrices depending on $`x^1,\mathrm{},x^n`$ such that $`dh=0`$, $`dh,_\mu =0`$ and $`d0`$, $`h0`$. Then there exist constants $`c^1,c^2`$ and scalar functions $`q^1,q^2,p^1,p^2`$ (depending on $`x^1,\mathrm{},x^n`$) such that
$$\begin{array}{c}d=\left(\begin{array}{cc}q^1c^2& q^1c^1\\ q^2c^2& q^2c^1\end{array}\right)=\left(\begin{array}{c}q^1\\ q^2\end{array}\right)\left(\begin{array}{cc}c^2& c^1\end{array}\right)qc^{},\hfill \\ h=\left(\begin{array}{cc}c^1p^1& c^1p^2\\ c^2p^1& c^2p^2\end{array}\right)=\left(\begin{array}{c}c^1\\ c^2\end{array}\right)\left(\begin{array}{cc}p^1& p^2\end{array}\right)cp^T\hfill \end{array}$$
(44)
Proof: The columns of $`h`$ are orthogonal to the rows of $`d`$. If $`det(d)0`$, then, obviously, $`h=0`$ in contrary to our assumptions. Therefore $`det(d)=0`$ which means that the rows of $`d`$ are linearly dependent. Similarly, the columns of $`h`$ are linearly dependent as well. We denote them by $`p^1c`$ and $`p^2c`$ (where $`c`$ is a column vector). Thus $`h=cp^T`$, where $`p^T:=(p^1,p^2)`$.
$`dh=0`$ means that the columns of $`h`$ are orthogonal to the rows of $`d`$. Therefore these rows are of the form $`q^1c^{}`$, $`q^2c^{}`$, where $`c^{}`$ is a vector orthogonal to $`c`$, and, finally $`d=qc^{}`$. Thus we obtained (44).
Taking into account the condition $`dh,_\mu =0`$ we get
$$0=qc^{}(c,_\mu p^T+cp^T,_\mu )=qc^{}c,_\mu p^Tc^{}c,_\mu =0$$
This means that $`c^2c^1,_\mu =c^1c^2,_\mu `$, or $`c^2/c^1`$ is a constant. In other words, $`c^1=fc^{10}`$, $`c^2=fc^{20}`$ ($`f`$ is a function, $`c^{10}`$, $`c^{20}`$ are constants. To complete the proof we redefine $`pfp`$, $`qfq`$, and $`c^{k0}c^k`$. $`\mathrm{}`$
Therefore,
$$D(\lambda _k)=\mathrm{\Psi }(\lambda _k)q_kc_k^{}\mathrm{\Psi }^1(\lambda _k),$$
(45)
where $`c_k`$ are given constant column unit vectors, $`c_k^{}`$ is a row vector orthogonal to $`c_k`$ and $`q_k`$ are some vector-valued functions (column vectors). We keep the notation $`q_kc_k^{}d_k`$, but now in general $`d_k^20`$.
We notice that the freedom concerning the choice of $`q_k`$ corresponds to the arbitrariness of the normalization matrix. In particular, the condition (5) imposes strong constraints on $`𝒩`$. The condition (5) can be rewritten as $`q_k=\phi _kc_k`$
The constraint (39) implies $`detD(\lambda _k)=0`$. In the case of $`2\times 2`$ matrices the equation $`detD(\lambda )=0`$ (where $`D`$ is given by (38)) has $`2N`$ roots (at most): $`\lambda _1`$, …$`\lambda _{2N}`$.
Taking any $`N+1`$ pairwise different roots (say $`\lambda _1,\mathrm{},\lambda _{N+1}`$) and using Lagrange’s interpolation formula for polynomials, we get the generalization of the formula (9):
$$D(\lambda )=\underset{k=1}{\overset{N+1}{}}\left(\underset{\stackrel{j=1}{jk}}{\overset{N+1}{}}\frac{(\lambda \lambda _j)}{(\lambda _k\lambda _j)}\right)\mathrm{\Psi }(\lambda _k)q_kc_k^{}\mathrm{\Psi }^1(\lambda _k).$$
(46)
We have also $`N1`$ matrix constraints which result from evaluating the formula (46) at $`\lambda _{N+2},\mathrm{},\lambda _{2N}`$:
$$\underset{k=0}{\overset{N+1}{}}\frac{\mathrm{\Psi }(\lambda _k)q_kc_k^{}\mathrm{\Psi }^1(\lambda _k)}{(\lambda _k\lambda _0)\mathrm{}(\lambda _k\lambda _{k1})(\lambda _k\lambda _{k+1})\mathrm{}(\lambda _k\lambda _{N+1})}=0,$$
(47)
where $`\lambda _0=\lambda _{N+2},\mathrm{},\lambda _{2N}`$.
We denote
$$Q_k:=\mathrm{\Psi }(\lambda _k)q_k,C_k^{}:=c_k^{}\mathrm{\Psi }^1(\lambda _k)$$
(48)
The Darboux matrix is parameterized by $`2N`$ constants $`\lambda _k`$, $`2N`$ vector functions $`q_k`$ and $`2N`$ constant vectors $`c_k`$ subject to the constraints (47).
The crucial point consists in solving the system (47) in order to get parameterization of the Darboux matrix by a set of independent quantities. We plan to express $`2N2`$ functions from among $`Q_1,\mathrm{},Q_{2N}`$ by other data. For instance, we choose $`Q_1,Q_2`$ as independent functions (they correspond to the normalization matrix $`𝒩`$).
We rewrite the system (47) as
$$\sigma _{\nu 0}Q_\nu C_\nu ^{}+\underset{k=1}{\overset{N+1}{}}\sigma _{\nu k}Q_kC_k^{}=0,(\nu =N+2,\mathrm{},2N),$$
(49)
where
$$\begin{array}{c}\sigma _{\nu k}=\frac{1}{(\lambda _k\lambda _\nu )(\lambda _k\lambda _1)\mathrm{}(\lambda _k\lambda _{k1})(\lambda _k\lambda _{k+1})\mathrm{}(\lambda _k\lambda _{N+1})}\hfill \\ \sigma _{\nu 0}=\frac{1}{(\lambda _\nu \lambda _1)\mathrm{}(\lambda _\nu \lambda _N)(\lambda _\nu \lambda _{N+1})}.\hfill \end{array}$$
Thus we have a system (49) linear with respect to $`Q_k`$. We are going to express $`2N2`$ vector functions $`Q_3,\mathrm{},Q_{2N}`$ by $`Q_1,Q_2`$ and the other parameters: $`C_k,\lambda _k`$. Then, using (48), we could get $`q_3,\mathrm{},q_{2N}`$, etc. However, it is better to write (46) in terms of $`Q_k`$:
$$D(\lambda )=\underset{k=1}{\overset{N+1}{}}\left(\underset{\stackrel{j=1}{jk}}{\overset{N+1}{}}\frac{(\lambda \lambda _j)}{(\lambda _k\lambda _j)}\right)Q_kc_k^{}\mathrm{\Psi }^1(\lambda _k).$$
(50)
Taking the scalar product of (49) by $`C_1`$ we get
$$Q_\nu =\underset{k=2}{\overset{N+1}{}}\frac{\sigma _{\nu k}C_k^{}C_1}{\sigma _{\nu 0}C_\nu ^{}C_1}Q_k,(\nu =N+2,\mathrm{},2N),$$
(51)
and the scalar product of $`\nu `$th equation of (49) by $`C_\mu `$ yields
$$\underset{k=1}{\overset{N+1}{}}\sigma _{\nu k}C_k^{}C_\nu Q_k=0,(\nu =N+2,\mathrm{},2N).$$
(52)
This is a system of $`N1`$ linear equations with respect to $`Q_1,\mathrm{},Q_{N+1}`$. Therefore, we can (for instance) express $`Q_3,\mathrm{},Q_{N+1}`$ in terms of $`Q_1,Q_2`$. Then, using (51), we have $`Q_{N+2},\mathrm{},Q_{2N}`$ expressed in the similar way.
Our method is closely related to the Neugebauer-Meinel approach . Let $`D`$ is given by (38). We denote by $`F(D(\lambda ))`$ the adjugate (or adjoint) matrix of $`D`$ which is, obviously, a polynom in $`\lambda `$. Thus
$$D(\lambda )F(D(\lambda ))=w(\lambda )I$$
(53)
where $`w(\lambda )=det(D(\lambda ))`$ is a scalar polynom and $`I`$ is the identity matrix.
Therefore, we can put $`T(\lambda )=F(D(\lambda ))`$ in the formula (39) and identify $`\lambda _k`$ with zeros of $`detD(\lambda )`$.
In the Neugebauer approach the matrix coefficients $`A_k`$ of the Darboux matrix are obtained by solving the following system
$$D(\lambda _k)\mathrm{\Psi }(\lambda _k)c_k=0,(k=1,\mathrm{},nN)$$
(54)
where $`\lambda _k`$ and constant vectors $`c_k`$ are treated as given parameters. Thus one has $`n^2N`$ scalar equations for $`(N+1)n^2`$ scalar variables. One of the matrices $`A_k`$, say $`A_N`$, is considered as undetermined normalization matrix.
We point out that $`D(\lambda _k)`$ given by the formula (45) satisfy (54).
## 5 The discrete case
The discrete analogue of (1) is the following system of linear difference equation
$$T_\mu \mathrm{\Psi }=U_\mu \mathrm{\Psi },(\mu =1,\mathrm{},m),$$
(55)
where $`T_\nu `$ denotes the shift in $`\nu `$th variable, i.e., $`(T_\nu \mathrm{\Psi })(x^1,\mathrm{},x^\nu ,\mathrm{},x^m):=\mathrm{\Psi }(x^1,\mathrm{},x^\nu +1,\mathrm{},x^m)`$. The Darboux transformation is defined in the standard way:
$$\stackrel{~}{\mathrm{\Psi }}=D\mathrm{\Psi },T_\mu \stackrel{~}{\mathrm{\Psi }}=\stackrel{~}{U}_\mu \stackrel{~}{\mathrm{\Psi }}.$$
(56)
Therefore $`(T_\mu D)(T_\mu \mathrm{\Psi })=\stackrel{~}{U}_\mu D\mathrm{\Psi }`$, and, finally
$$(T_\mu D)U_\mu =\stackrel{~}{U}_\mu D$$
(57)
If $`D^2(\lambda _1)=0`$, then multiplying (57) by $`D(\lambda )`$ from the right, and evaluating the obtained equation at $`\lambda =\lambda _1`$ we see that the right hand side vanishes and we get:
$$(T_\mu D_1)U_\mu (\lambda _1)D_1=0$$
(58)
where we denote $`D_1:=D(\lambda _1)`$. In order to solve (58) we put
$$D_1=\phi _1\mathrm{\Psi }_1d_1\mathrm{\Psi }_1^1$$
where $`\mathrm{\Psi }_1:=\mathrm{\Psi }(\lambda _1)`$. Then (58) takes the form:
$$\phi _1T_\mu (\phi _1)(T_\mu \mathrm{\Psi }_1)(T_\mu d_1)d_1\mathrm{\Psi }_1^1=0.$$
Therefore, if
$$(T_\mu d_1)d_1=0$$
(59)
then the equation (58) is satisfied. The condition (59) can be rewritten (at least in the matrix case) as
$$\mathrm{Im}d_1\mathrm{ker}(T_\mu d_1)$$
In other words, the sequence of linear operators
$$\mathrm{}T_\mu ^1d_1d_1T_\mu d_1T_\mu ^2d_1\mathrm{}$$
is an exact sequence .
Similarly as in the smooth case we mostly confine ourselves to the simplest solution of (59), i.e., $`d_1=\mathrm{const}`$ which implies $`d_1^2=0`$. The Darboux matrix has the same form (9) as in the continuum case.
Summary. In this paper we developed the approach of considering explicitly the most important reductions, extending our results on the $`N`$-soliton case, and showing that the discrete case is, as usual, very similar to the continuous one.
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# Reciprocals of Binary Power Series
## 1 Introduction
There is a unique set $`B`$ of nonnegative integers with the property that each positive integer can be written in the form $`s^2+b`$ ($`s:=\{0,1,2,\mathrm{}\},bB`$) in an even number of ways. Specifically,
$$B=\{0,1,2,3,5,7,8,9,13,17,18,23,27,29,31,32,35,\mathrm{}\}.$$
Are the even numbers in $`B`$ exactly those of the form $`2k^2`$? Does $`B`$ have positive density?
Before addressing these two questions, we restate and motivate the problem in greater generality. Given any sets $`A,B`$, the asymmetric additive representation function is defined by
$$R(n):=\mathrm{\#}\{(a,b):n=a+b,aA,bB\};$$
equivalently, we could define $`R`$ by noting that
$$\left(\underset{aA}{}q^a\right)\left(\underset{bB}{}q^b\right)=\underset{n=0}{\overset{\mathrm{}}{}}R(n)q^n.$$
We are interested in the situation where $`R(0)=1`$ and $`R(n)0(mod2)`$ for $`n>0`$, i.e., the situation where $`_nR(n)q^n=1`$ in the ring of power series $`𝔽_2[[q]]`$. In this case, we say that $`A`$ and $`B`$ are *reciprocals*, and we write $`\overline{A}=B`$ and $`\overline{B}=A`$. The general problem of this paper is to find the reciprocals of several special sets $`A`$, and to draw some conclusions about “typical” properties of reciprocals. We are particularly concerned with the relative density,
$$\delta (\overline{A},n):=\frac{|\overline{A}[0,n]|}{n+1},$$
and the density $`\delta (\overline{A}):=lim_n\mathrm{}\delta (\overline{A},n)`$ (when the limit exists).
We began studying this problem after reading two articles by Berndt, Yee, and Zaharescu \[MR2039324, MR1984662\], where bounds on the density of the set $`P_{\text{odd}}:=\{n:p(n)1(mod2)\},`$ with $`p(n)`$ being the ordinary partition function<sup>a</sup><sup>a</sup>a$`p(n)`$ is the number of ways to write $`n`$ as a sum of nonincreasing positive integers. For example, $`4=3+1=2+2=2+1+1=1+1+1+1`$, so $`p(4)=5`$., are proved. The starting point for their work is Euler’s pentagonal number theorem \[MR1083765\]\*Theorem 10.9, and in particular that the reciprocal of $`P_{\text{odd}}`$ is the set $`\{n(3n+1)/2:n\}`$ of pentagonal numbers. Since the known bounds (see \[MR1816213\] and \[MR1657968\] for the currently-best results) on the thickness of $`P_{\text{odd}}`$ are so strikingly far from what is believed to be true, we felt that it would be beneficial to study the “reciprocal” notion in a more general setting.
$`P_{\text{odd}}`$ is pictured in Figure 1, where not only does it appear to have density $`1/2`$, but the walk defined by $`w(n):=2\left|P_{\text{odd}}[0,n]\right|n`$ visually appears to be a simple random walk. See \[MR0227126\] for a report of more elaborate statistical tests on the set $`P_{\text{odd}}`$. We note that while $`p(n)`$ appears to be uniformly distributed modulo 2 and 3, it has been known since the time of Ramanujan to not be uniformly distributed modulo 5, 7, or 11.
In contrast to that of the pentagonal numbers, the density of the reciprocal of the squares appears to drop off steadily to 0. Set $`S:=\{n^2:n\}`$, with reciprocal $`\overline{S}`$. The relative density of $`\overline{S}`$ is pictured in Figure 2. In Section 6.1, we prove that the even numbers in $`\overline{S}`$ are precisely $`\{2n^2:n\}`$, and we characterize the $`n\overline{S}`$ with $`n1(mod4)`$ as those $`n`$ whose prime factorization has a particular shape. Those $`n\overline{S}`$ with $`n3(mod4)`$ are characterized in terms of the number of representations of $`n`$ by certain quadratic forms.
Generalizing the squares and pentagonal numbers, we treat
$$\mathrm{\Theta }(c_1,c_2):=\{c_1n+c_2\frac{n(n1)}{2}:n\}$$
for general $`c_1`$ and $`c_2`$ in Section 6. A few interesting special cases are the binomial coefficients $`\mathrm{\Theta }(0,1)=\{\left(\genfrac{}{}{0pt}{}{n}{2}\right):n\}`$, the squares $`\mathrm{\Theta }(1,2)`$, and the pentagonal numbers $`\mathrm{\Theta }(1,3)`$.
###### Conjecture 1.1.
The reciprocal of the set $`\mathrm{\Theta }(c_1,c_2)`$, where $`02c_1c_2`$ and $`\mathrm{gcd}(c_1,c_2)=1`$, has density $`0`$ if $`c_22(mod4)`$, and otherwise has density $`1/2`$. More precisely, if $`c_22(mod4)`$, then
$$\underset{n\mathrm{}}{lim}\frac{\left|\overline{\mathrm{\Theta }(c_1,c_2)}[0,n]\right|}{n/\mathrm{log}n}=C,$$
for some positive constant $`C`$ depending only on $`c_2`$. If $`c_22(mod4)`$, then
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{\left|\overline{\mathrm{\Theta }(c_1,c_2)}[0,n]\right|n/2}{\sqrt{n\mathrm{log}\mathrm{log}(n)/2}}\right|=1.$$
Numerically, it seems that the constant $`C`$ is 2 if $`c_2=2`$ or 6, and $`C=4`$ if $`c_2=10`$. We lack sufficient data to guess the other values. The authors believe that the non-effective $`c_22(mod4)`$ case might be provable by showing that the generating function of $`\mathrm{\Theta }(c_1,c_2)`$ is congruent modulo 2 to an integer-weight modular form, which has almost all of its Fourier coefficients even. This is outside the scope of this paper, and we leave it as an area for further study.
The $`c_22(mod4)`$ case is motivated by the celebrated law of the iterated logarithm. Let $`X_1,X_2,\mathrm{}`$ be independent random variables taking the values 0 and 1 with probability $`1/2`$. The law of the iterated logarithm states that
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{_{i=1}^nX_in/2}{\sqrt{n\mathrm{log}\mathrm{log}(n)/2}}\right|=1$$
with probability 1. What we actually would like to conjecture is that reciprocal of $`\mathrm{\Theta }(c_1,c_2)`$, with $`c_22(mod4)`$, is statistically indistinguishable from a truly random set with density $`1/2`$. The phrase “statistically indistinguishable” is too vague, however, so in Conjecture 1.1 we have settled for this one specific statistic.
The natural expectation is that, barring some cosmic coincidence or obvious structure, the reciprocal of a set should have density $`1/2`$. This is affirmed by the case of a random set, which we handle in detail in Section 3: let $`X_1,X_2,\mathrm{}`$ be independent random variables taking the values 0 and 1, with probabilities bounded away from 0 and 1, and set $`F:=\{0\}\{n:X_n=1\}`$. Theorem 3.1 states that the reciprocal of $`F`$ has density $`1/2`$ with probability 1. This makes the sets whose reciprocals do not have density $`1/2`$ the interesting ones.
Our purpose is to identify relevant properties of those sets whose reciprocals have density different from $`1/2`$. Specifically, in addition to random sets and $`\mathrm{\Theta }(c_1,c_2)`$, we consider finite sets, the set of powers of two, and the set of Prouhet-Thue-Morse numbers.
* Finite sets: the reciprocal has a rational density, and appears to typically have density slightly below $`1/2`$. We identify through algebraic properties two infinite classes of polynomials, one whose reciprocals have density strictly larger than $`1/2`$, and one whose reciprocals have density at most $`1/2`$.
* Powers of 2: the reciprocal of the thin set $`\{0\}\{2^n:n\}`$ is the thin set $`\{2^n1:n\}`$. In particular, we describe the reciprocal of $`\{0\}\{2^{mn}:n\}`$ for every $`m`$.
* Prouhet-Thue-Morse numbers<sup>b</sup><sup>b</sup>b$`T=\{0,3,5,6,9,10,12,15,17,18,20,\mathrm{}\}`$: the reciprocal of
$$T:=\{n:\text{the binary expansion of }n\text{ contains an even number of “1”s}\}.$$
has density $`1/3`$. Specifically, we prove that $`k\overline{T}`$ if and only if $`k=0`$ or $`(k\pm 1)/4`$ is an integer whose binary expansion ends in an even number of zeros.
The strongest conjecture that is consistent with our theorems, our experiments, and Conjecture 1.1, is Conjecture 1.2.
###### Conjecture 1.2.
If a set contains 0, is not periodic, and is uniformly distributed modulo every power of 2, then its reciprocal has positive density.
We now include a section-by-section agenda for the remainder of the paper.
Motivate and contextualize reciprocals of sets.
Introduce notation and derive general expressions for reciprocals.
Consider reciprocals of random sets with positive density.
Consider reciprocals of finite sets.
Consider the reciprocal of the powers of 2, and similar sets.
Consider the reciprocal of $`\mathrm{\Theta }(c_1,c_2)`$, particularly the squares.
Consider the Prouhet-Thue-Morse sequence.
## 2 Notation and General Formulas
Throughout this paper, we let
$$(q)=f_0+f_1q+f_2q^2+\mathrm{}\text{and}\overline{}(q)=\overline{f}_0+\overline{f_1}q+\overline{f_2}q^2+\mathrm{}$$
(1)
be elements of $`𝔽_2[[q]]`$ that satisfy the equation
$$(q)\overline{}(q)=1.$$
(2)
In particular, $`f_0=\overline{f}_0=1`$. We define the integer sets $`F:=\{n0:f_n=1\}`$ and $`\overline{F}:=\{n0:\overline{f}_n=1\}`$.
Note that (2) implies (for all $`k1`$) that $`(q^k)\overline{}(q^k)=1`$ also. This corresponds to noting that multiplying everything in $`F`$ by $`k`$ has the effect of multiplying everything in $`\overline{F}`$ by $`k`$. With this in mind, we sometimes make the convenient assumption that $`\mathrm{gcd}F=1`$.
Our next lemma is a fundamental identity in $`𝔽_2[[q]]`$, and has a number of remarkable consequences. We use it frequently throughout this paper.
###### Lemma 2.1.
The reciprocal of $`(q)`$ is $`(q)(q^2)(q^4)(q^8)\mathrm{}`$. That is,
$$1=(q)\underset{k=0}{\overset{\mathrm{}}{}}(q^{2^k}).$$
(3)
###### Proof.
First, notice that both sides of this equation have constant term equal to $`1`$. Also notice that for any fixed $`n>0`$, only finitely many terms of the infinite product affect the coefficient of $`q^n`$. Thus, the coefficient of $`q^n`$ on the right hand side of (3) is also the coefficient of $`q^n`$ in
$`(q){\displaystyle \underset{k=0}{\overset{\mathrm{log}_2n}{}}}(q^{2^k}).`$
By the so-called children’s binomial theorem<sup>c</sup><sup>c</sup>c$`(a+b)^2=a^2+b^2(mod2)`$, $`(q)(q)=(q^2)`$. Multiplying by $`(q^2)`$, we see that $`(q)(q)(q^2)=(q^2)(q^2)=(q^4)`$, and continuing we get
$`(q)(q)(q^2)(q^4)\mathrm{}(q^{2^{\mathrm{log}_2n}})=\left(q^{2^{\mathrm{log}_2n+1}}\right).`$ (4)
Now notice that since $`0<n<2^{\mathrm{log}_2n+1}`$, the coefficient of $`q^n`$ on the right hand side of (4) is $`0`$, and our result follows. ∎
We now give a list of recurrences for $`\overline{f}_n`$, discuss the usefulness of each, and prove them.
###### Lemma 2.2.
If $`(q)\overline{}(q)=1`$, then $`\overline{f}_0=1`$ and for $`n>0`$,
1. $`\overline{f}_n={\displaystyle \underset{j=1}{\overset{n}{}}}f_j\overline{f}_{nj}`$;
2. $`\overline{f}_n=1`$ if and only if $`\mathrm{\#}\{(x_0,x_1,\mathrm{}):x_iF,n={\displaystyle \underset{i0}{}}x_i2^i\}`$ is odd;
3. $`\overline{f}_n={\displaystyle \underset{\stackrel{}{x}}{}}f_{x_1}f_{x_2}\mathrm{}f_x_{\mathrm{}}`$, where the summation extends over all tuples $`\stackrel{}{x}=(x_1,\mathrm{},x_{\mathrm{}})`$ with $`n=_{i=1}^{\mathrm{}}x_i`$ and each $`x_i>0`$ ($`\mathrm{}`$ is allowed to vary);
4. $`\overline{f}_n={\displaystyle \underset{0i<n/4}{}}f_{n2i}\overline{f}_i+G(f_1,f_2,\mathrm{},f_{n/2})`$, for some function $`G`$.
Lemma 2.2(i ) is valuable because of its simplicity. For instance, it is immediately apparent from this recurrence relation that $`\overline{}`$ is uniquely defined and always exists (provided $`f_0=1`$).
In several of the examples we consider, the set $`F`$ has some special properties modulo a power of 2. Lemma 2.2(ii ) facilitates our exploitation of these special properties.
Lemma 2.2(iii ) is useful because of its symmetry, and because its right-hand side does not expressly reference the $`\overline{f}`$ sequence. As a specific example, let $`r(n)`$ be the number of ways to write $`n`$ as a sum of positive pentagonal numbers (counting order). Then, by Lemma 2.2(iii ), $`p(n)r(n)(mod2)`$. We also use Lemma 2.2(iii ), for example, to prove Lemma 2.2(iv ).
If one lets the $`f_i`$ be independent random variables, then the expression in Lemma 2.2(iv ) contains a summation of weakly dependent random variables, and a deterministic function of $`f_1,\mathrm{},f_{n/2}`$. This allows us to say something explicit about the distribution of the resulting random variable $`\overline{f}_n`$ (see Theorem 3.1).
Another remarkable aspect of Lemma 2.2(iv ) is that $`\overline{f}_n`$ does not depend in any way on $`f_{n1},f_{n3},\mathrm{},f_{nc}`$, where $`c`$ is the largest odd number strictly less than $`n/2`$. For example,
$$\overline{f}_{11}=f_{11}+f_9f_1+f_7f_2+f_7f_1+f_5f_3+f_4f_3+f_3f_2f_1+f_3f_2+f_2f_1+f_1$$
does not depend on $`f_{10},f_8,`$ or $`f_6`$.
###### Proof.
Comparing the coefficients of $`q^n`$ on the left- and right-hand sides of equation (2) yields $`_{j=0}^nf_j\overline{f}_{nj}=0`$. Lemma 2.2(i ) is this expression rearranged, using the fact that $`f_0=1`$.
Similarly, Lemma 2.2(ii ) equates the coefficients of $`q^n`$ on the left- and right-hand sides of equation (3), with the right-hand side interpreted as a product in $``$.
One can prove Lemma 2.2(iii ) by induction, using Lemma 2.2(i ) to complete the induction step. Alternatively, one may simply compare the coefficients of $`q^n`$ on the left- and right-hand sides of
$$\overline{}=\frac{1}{}=\frac{1}{1(1)}=1+(1)+(1)^2+(1)^3+\mathrm{},$$
which is valid because we are working over $`𝔽_2`$.
Recall Kummer’s result that the multinomial coefficient $`\left({\displaystyle \genfrac{}{}{0pt}{}{m_1+\mathrm{}+m_k}{m_1,m_2,\mathrm{},m_k}}\right)={\displaystyle \frac{(m_1+\mathrm{}+m_k)!}{m_1!m_2!\mathrm{}m_k!}}`$ is relatively prime to a prime $`p`$ if and only if $`m_1,\mathrm{},m_k`$ can be added in base $`p`$ without carrying \[Kummer\]. We are working with $`p=2`$, so our condition is: $`\left(\genfrac{}{}{0pt}{}{m_1+\mathrm{}+m_k}{m_1,\mathrm{},m_k}\right)`$ is odd if and only if no two of the binary expansions of $`m_1,\mathrm{},m_k`$ have a “1” in the same position. We call such a list of positive integers $`m_1,\mathrm{},m_k`$ non-overlapping.
Let $`\pi (n)`$ be the set of partitions of $`n`$ whose distinct parts $`x_1,\mathrm{},x_k`$ have non-overlapping multiplicities $`m_1,\mathrm{},m_k`$. Continuing from 2.2(iii ), we have
$`\overline{f}_n`$ $`={\displaystyle \underset{\begin{array}{c}x_1+\mathrm{}+x_{\mathrm{}}=n\\ x_i>0\end{array}}{}}f_{x_1}f_{x_2}\mathrm{}f_x_{\mathrm{}}`$
$`={\displaystyle \underset{\begin{array}{c}m_1a_1+\mathrm{}+m_ka_k=n\\ a_1>\mathrm{}>a_k>0\\ m_i>0\end{array}}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{m_1+\mathrm{}+m_k}{m_1,\mathrm{},m_k}}\right)f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}`$
$`={\displaystyle \underset{\pi (n)}{}}f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}`$
$`={\displaystyle \underset{\begin{array}{c}\pi (n)\\ a_1>n/2\end{array}}{}}f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}+{\displaystyle \underset{\begin{array}{c}\pi (n)\\ a_1n/2\end{array}}{}}f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}`$
If $`a_1>n/2`$, then it must have multiplicity $`m_1=1`$, and if $`m_1,\mathrm{},m_k`$ are non-overlapping then the other $`m_i`$ are even:
$$f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}=f_{a_1}\left(f_{a_2}^{m_2/2}\mathrm{}f_{a_k}^{m_k/2}\right)^2.$$
This implies that $`na_1`$ is even, and $`a_2\frac{m_2}{2}+\mathrm{}+a_k\frac{m_k}{2}`$ is a partition of $`(na_1)/2`$. Setting $`2i=na_1`$, we get
$$\overline{f}_n=\underset{\pi (n)}{}f_{a_1}^{m_1}f_{a_2}^{m_2}\mathrm{}f_{a_k}^{m_k}=\underset{0i<n/4}{}f_{n2i}\underset{\pi (i)}{}f_{a_1}^{m_1}\mathrm{}f_{a_k}^{m_k}+\underset{\begin{array}{c}\pi (n)\\ a_in/2\end{array}}{}f_{a_1}^{m_1}\mathrm{}f_{a_k}^{m_k}.$$
Using Lemma 2.2(iii ) again, this becomes
$$\overline{f}_n=\underset{0i<n/4}{}f_{n2i}\overline{f}_i+G(f_1,\mathrm{},f_{n/2})$$
for a specific function $`G`$. ∎
## 3 Random power series
In this section, we consider the reciprocal of a random power series in $`𝔽_2[[q]]`$. The results of this section are strong evidence that the density of $``$ plays little to no role in determining the density of $`\overline{}`$, and that unless the coefficients of $``$ have some structure, the density of $`\overline{}`$ is $`1/2`$.
Recall that a Bernoulli variable is a random variable that takes values in $`\{0,1\}`$.
###### Theorem 3.1.
Suppose that $`f_1,f_2,`$ are independent Bernoulli variables, with
$$\underset{n}{inf}\mathrm{min}\{\left[f_n=0\right],\left[f_n=1\right]\}>0.$$
Then $`\delta (\overline{})=1/2`$ with probability 1.
We need the following two lemmas.
###### Lemma 3.2 (Lévy’s Borel-Cantelli lemma).
Let $`E_1,E_2,\mathrm{},`$ be events, and define $`Z_n:=_{k=1}^nI_{E_k}`$, the random variable that records the number of $`E_1,E_2,\mathrm{},E_n`$ that occur. Define
$$\xi _k:=\left[E_kE_1,E_2,\mathrm{},E_{k1}\right].$$
If $`_{k=1}^{\mathrm{}}\xi _k`$ diverges, then $`Z_n`$ is asymptotically equal to $`_{k=1}^n\xi _k`$ with probability 1.
For a proof of Lévy’s Borel-Cantelli lemma, we refer the reader to \[MR1155402\]\*Sec 12.15.
###### Lemma 3.3 (Binary Central Limit Theorem).
Let $`X_i`$ be 0 with probability $`\gamma _i`$ and 1 with probability $`1\gamma _i`$, and suppose that $`X_1,X_2,`$ are independent. Then, as $`n\mathrm{}`$,
$$\left[\underset{i=1}{\overset{n}{}}X_i0(mod2)\right]\frac{1}{2}$$
if and only if some $`\gamma _i=1/2`$ or $`_{i=1}^n\mathrm{min}\{\gamma _i,1\gamma _i\}`$ diverges.
###### Proof.
Let $`S_n:=_{i=1}^nX_i`$, and define $`p_i`$ by $`\left[S_n0(mod2)\right]=p_i`$. Clearly $`S_n`$ is even if and only if $`S_{n1}`$ and $`X_n`$ are both even or both odd:
$$p_n=p_{n1}\gamma _n+(1p_{n1})(1\gamma _n).$$
Clearly $`2p_11=2\gamma _11`$, and
$$2p_n1=2\left(p_{n1}\gamma _n+(1p_{n1})(1\gamma _n)\right)1=(2p_{n1}1)(2\gamma _n1),$$
which provides the base case and inductive step for the equality
$$2p_n1=\underset{i=1}{\overset{n}{}}(2\gamma _i1).$$
By the standard results for infinite products, we now see that $`2p_n10`$ if and only if $`2\gamma _i1=0`$ for some $`i`$ or $`_{i=1}^n\mathrm{min}\{\gamma _i,1\gamma _i\}`$ diverges. ∎
###### Proof of Theorem 3.1.
We begin with some notation:
$`\alpha _n`$ $`:=\left[f_n=0\right],`$
$`\beta _n`$ $`:=\mathrm{min}\{\alpha _n,1\alpha _n\},`$
$`\beta `$ $`:=\underset{n\mathrm{}}{inf}\beta _n,`$
and note that $`0<\beta 1/2`$. We will show first that $`\left[\overline{f}_n=0\right]1/2`$ as $`n\mathrm{}`$, and then will show that $`\delta (\overline{F})=1/2`$ with probability 1.
Lemma 2.2(i ) says that $`\overline{f}_n=f_n+_{j=1}^{n1}f_j\overline{f}_{nj}`$, whence
$$\left[\overline{f}_n=0\right]=\left[f_n=0\right]\left[\underset{j=1}{\overset{n1}{}}f_j\overline{f}_{nj}=0\right]+\left[f_n=1\right]\left[\underset{j=1}{\overset{n1}{}}f_j\overline{f}_{nj}=1\right]$$
is a weighted average of $`\left[f_n=0\right]=\alpha _n`$ and $`\left[f_n=1\right]=1\alpha _n`$. Consequently, $`\left[\overline{f}_n=0\right]\beta _n\beta `$ and $`\left[\overline{f}_n=1\right]\beta _n\beta `$.
Set $`B_n:=\{i:0i<n/4,\overline{f}_i=1\}`$, and set $`G_n:=G(f_1,\mathrm{},f_{n/2})`$, where $`G`$ is the function from Lemma 2.2(iv ). We have, from Lemma 2.2(iv ),
$$\overline{f}_n=\underset{iB_n}{}f_{n2i}+G_n.$$
From the previous paragraph, we know that $`𝔼\left[|B_n|\right]`$ is at least $`_{i=0}^{n/4}\beta _i\beta n/4`$. In particular, a routine calculation shows that $`|B_n|\mathrm{}`$ with probability 1. Thus, $`\left[|B_n|>K_n\right]1`$ if $`K_n`$ goes to infinity sufficiently slowly. We have
$$\begin{array}{c}\left[\overline{f}_n=0\right]=\hfill \\ \hfill [\overline{f}_n=0||B_n|>K_n][|B_n|>K_n]+[\overline{f}_n=0||B_n|K_n][|B_n|K_n],\end{array}$$
which for large $`n`$ becomes $`[\overline{f}_n=0]=[\overline{f}_n=0||B_n|>K_n]`$.
We now observe that $`\overline{f}_n=0`$ if and only if $`G_n=_{iB_n}f_{n2i}`$ (call this sum $`\sigma _n`$), so that
$$\begin{array}{c}[\overline{f}_n=0||B_n|>K_n]=\hfill \\ \hfill [G_n=0||B_n|>K_n,\sigma _n=0][\sigma _n=0||B_n|>K_n]+\\ \hfill [G_n=1||B_n|>K_n,\sigma _n=0][\sigma _n=1||B_n|>K_n].\end{array}$$
This is a weighted average of $`[\sigma _n=0||B_n|>K_n]`$ and $`[\sigma _n=1||B_n|>K_n]`$, both of which go to $`1/2`$ as $`n\mathrm{}`$ by the Binary Central Limit Theorem. Thus,
$$[\overline{f}_n=0][\overline{f}_n=0||B_n|>K_n]\frac{1}{2}$$
with each of the “$``$” becoming “$`=`$” as $`n\mathrm{}`$.
Now that we have shown that $`\left[\overline{f}_n=0\right]1/2`$, we know that $`𝔼\left[\delta (\overline{F},n)\right]1/2`$, but this does not imply that $`\delta (\overline{F},n)1/2`$ ever, much less with probability 1. This last step again requires the at-least-weak independence of $`\overline{f}_n`$ from $`\overline{f}_1,\mathrm{},\overline{f}_{n1}`$, and the technicalities are handled for us by Lévy’s Borel-Cantelli lemma.
Let $`E_k`$ be the event $`\{\overline{f}_k=0\}`$, and set $`\xi _k:=\left[\overline{f}_k=0\overline{f}_1,\overline{f}_2,\mathrm{},\overline{f}_n\right]`$. By the comment above, $`0<\beta \xi _k`$, so $`_{k=1}^{\mathrm{}}\xi _k=\mathrm{}`$. Thus, by Lemma 3.2,
$$\underset{n\mathrm{}}{lim}\frac{\delta (\overline{F},n)}{\frac{1}{n}_{k=1}^n\xi _k}=1$$
(5)
with probability 1. For every $`ϵ>0`$ there is an $`n_0`$ such that for all $`n>n_0`$
$$(1ϵ)\frac{1}{n}\underset{k=1}{\overset{n}{}}\xi _k\delta (\overline{F},n)(1+ϵ)\frac{1}{n}\underset{k=1}{\overset{n}{}}\xi _k.$$
These upper and lower bounds on $`\delta (\overline{F},n)`$ are non-random, so we may take expectations (for large $`n`$) to get
$$(1ϵ)\frac{1}{n}\underset{k=1}{\overset{n}{}}\xi _k\frac{1}{2}(1+ϵ)\frac{1}{n}\underset{k=1}{\overset{n}{}}\xi _k,$$
where we have used the linearity of expectation and the previously proved $`𝔼\left[\overline{f}_n=0\right]=\left[\overline{f}_n=0\right]1/2`$. This implies that $`\frac{1}{n}_{k=1}^n\xi _k1/2`$ also. Consequently, (5) now implies that
$$\delta (\overline{F}):=\underset{n\mathrm{}}{lim}\delta (\overline{F},n)=\frac{1}{2}$$
with probability 1. ∎
## 4 Polynomials
In this section, we study the reciprocals of polynomials in $`𝔽_2[q]`$. The coefficients of such a reciprocal are periodic (see Proposition 4.3 below), and so the reciprocal has rational density. We also give some indication of how the densities of reciprocals of polynomials are distributed, beginning in Subsection 4.1. In Subsection 4.2, we use the theory of de Bruijn cycles to exhibit an infinite family of polynomials whose reciprocals have densities strictly larger than $`1/2`$; in Subsection 4.3, we show that if two polynomials have product $`1+q^D`$ ($`D4`$), then at least one of them has a reciprocal with density at most $`1/2`$. In Subsection 4.4, we show that the reciprocal of an eventually periodic set<sup>d</sup><sup>d</sup>dMore precisely, a set whose indicator function is eventually periodic. containing 0 is an eventually periodic set containing 0.
Let $`_{i=0}^{\mathrm{}}b_i2^i`$ be the binary expansion of $`n`$; we define the polynomial $`𝒫_n(q):=_{i=0}^{\mathrm{}}b_iq^i𝔽_2[q]`$. Clearly this indexes all polynomials, and the invertible polynomials are precisely those with $`n`$ odd. For a polynomial $`𝒬𝔽_2[q]`$, we let $`\widehat{}𝒬`$ be the same polynomial with coefficients (all 0 or 1) in $``$. For instance, $`𝒬=𝒫_{\widehat{}𝒬(2)}`$ for every polynomial $`𝒬`$.
We denote by $`\mathrm{}(𝒫)`$ the length of the polynomial, i.e., $`\mathrm{}(𝒫)=\widehat{𝒫}(1)`$, and by $`\mathrm{deg}(𝒫)`$ the degree of the polynomial. Also, $`\mathrm{ord}(𝒫)`$ is the least positive $`D`$ such that $`𝒫`$ divides $`1+q^D`$. It is not immediately obvious that $`\mathrm{ord}(𝒫)`$ is well defined: it is for invertible $`𝒫`$, and this is the content of Proposition 4.1 below. For each polynomial $`𝒫`$, we define $`𝒫^{}`$ by $`𝒫𝒫^{}=1+q^{\mathrm{ord}𝒫}`$. We shall see that the properties of $`𝒫`$ and $`𝒫^{}`$ are intimately related (Propositions 4.2 and 4.5).
If $`\mathrm{ord}(𝒫)=2^{\mathrm{deg}(𝒫)}1`$, then $`𝒫`$ is called primitive, and $`𝔽_2[q]/()`$ is isomorphic to $`𝔽_{2^{\mathrm{deg}(𝒫)}}`$, with multiplicative generator $`q`$. All primitive polynomials are irreducible, but not vice versa; for example $`1+q^3+q^6`$ and $`1+q+q^2+q^3+q^4`$ are irreducible but not primitive.
Figure 3 tabulates properties of $`𝒫_n`$ for odd $`n<256`$, including factorizations, $`𝒫_n^{}`$, and densities of reciprocals.
In Figure 4, we plot the points $`(n,\delta (\overline{}𝒫_n))`$ for odd $`n`$ less than $`2^{12}`$. We note that $`\delta (\overline{}𝒫_n)`$ tends to be near $`1/2`$, but is biased toward being below $`1/2`$. This is also suggested, but not proven, by Proposition 4.5 below. In Proposition 4.4, we give an algebraically-described infinite set of $`n`$ such that $`\delta (\overline{}𝒫_n)>1/2`$. Note that $`\delta (\overline{}𝒫_n(q^k))=\frac{1}{k}\delta (\overline{}𝒫_n(q))`$, i.e., if $`𝒫_n`$ is a polynomial in $`q^2`$, $`q^3`$, etc, then its density is a priori less than $`1/2`$, $`1/3`$, etc. These points have been plotted with squares.
In Figure 5, we plot the empirical distribution function of $`\delta (\overline{}𝒫_n)`$. The large discontinuities near $`1/2`$ mean that these densities occur with large frequency (fully 421 of the 2048 polynomials $`𝒫_1,𝒫_3,\mathrm{},𝒫_{4095}`$ have reciprocals with density exactly $`1/2`$). Again visible in Figure 5 is the preference of $`𝒫`$ to have reciprocal with density less than $`1/2`$. The most interesting issue raised in this section, which remains unanswered, is to describe the set
$$\{\delta (\overline{}𝒫):𝒫\text{ is a polynomial}\}.$$
For example, is there an $`n`$ with $`\delta (\overline{}𝒫_n)=3/4`$?
### 4.1 Order and Density
Our first proposition demonstrates that $`\mathrm{ord}(𝒫)`$ is well-defined, and our next proposition shows the connection between $`\delta (\overline{}𝒫)`$, $`\mathrm{ord}(𝒫)`$, and $`𝒫^{}`$.
###### Proposition 4.1.
If $`𝒫`$ is a polynomial, then $`\mathrm{ord}(𝒫)`$ is finite.<sup>e</sup><sup>e</sup>eActually, the proof can be refined to show that $`\mathrm{ord}(𝒫)2^{\mathrm{deg}(𝒫)}1`$ if $`𝒫`$ is irreducible, and otherwise $`\mathrm{ord}(𝒫)=2^i\mathrm{lcm}\{\mathrm{ord}(𝒱_1),\mathrm{},\mathrm{ord}(𝒱_k)\}`$ for some $`12^ik`$, where $`𝒫=𝒱_1\mathrm{}𝒱_k`$.
###### Proof.
Let $`𝒱_1,\mathrm{},𝒱_k`$ be the irreducible factors of $`𝒫`$, and let $`d_i`$ be the multiplicative order of $`q`$ in the field $`𝔽_2[q]/(𝒱_i)`$. In particular, $`1+q^{xd_i}`$ is a multiple of $`𝒱_i`$ for each $`x`$. Set $`L:=\mathrm{lcm}\{d_1,\mathrm{},d_k\}`$ and define $`𝒱_i^{}`$ by $`𝒱_1𝒱_1^{}=1+q^L`$ and for $`1<ik`$ by $`𝒱_i𝒱_i^{}=1+q^{2^{i2}L}`$. Now
$$\begin{array}{c}𝒫\underset{i=1}{\overset{k}{}}𝒱_i^{}=(1+q^L)(1+q^L)(1+q^{2L})\mathrm{}(1+q^{2^{k2}L})\hfill \\ \hfill =(1+q^{2L})(1+q^{2L})(1+q^{4L})\mathrm{}(1+q^{2^{k2}L})=(1+q^{2^{k1}L}),\end{array}$$
by repeated use of the children’s binomial theorem. ∎
We emphasize that, given $`𝒫`$ and the equality $`𝒫=1+q^D`$ for some $``$, $`D`$ is not uniquely determined. For example, $`𝒫(q)𝒫(q)(q^2)=1+q^{2D}`$. Nor does the proof given above always provide the minimal $`D`$.
###### Proposition 4.2.
$`\delta (\overline{}𝒫)=\mathrm{}(𝒫^{})/\mathrm{ord}(𝒫)`$.
###### Proof.
Since $`𝒫\frac{𝒫^{}}{1+q^{\mathrm{ord}𝒫}}=1`$, we see that the reciprocal of $`𝒫`$ is periodic with period $`\mathrm{ord}𝒫`$ (although this may not be the minimal period), and in each period has density $`\mathrm{}(𝒫^{})/\mathrm{ord}𝒫`$. ∎
### 4.2 de Bruijn cycle algebra
Our next proposition shows that the reciprocal of a polynomial is a special case of a linear-shift register. Fortunately, there is an enormous literature on linear-shift registers (see \[ShiftRegisterBook\], for example).
###### Proposition 4.3.
If $``$ is a polynomial with degree $`d`$, then (letting $`\overline{f}_j=0`$ for negative $`j`$)
$$\overline{f}_n=\underset{j=1}{\overset{d}{}}f_j\overline{f}_{nj}.$$
(6)
Alternatively, $`\overline{f}_n`$ is the constant term of $`q^nmod`$.
###### Proof.
Since $`f_j=0`$ for all $`j>d`$, the recurrence (6) follows immediately from Lemma 2.2(i ).
Let $`M`$ be the matrix whose $`k^{\text{th}}`$ row is the elementary vector supported in coordinate $`k+1`$, for $`k=1,\mathrm{},d1`$, and whose last row is the vector $`(f_0,\mathrm{},f_{d1})`$, i.e., $`M`$ is the companion matrix of $``$. Write $`c_n`$ for the constant coefficient of $`q^nmod`$. We claim that
$$M\left(\begin{array}{c}c_s\\ \mathrm{}\\ c_{s+d1}\end{array}\right)=\left(\begin{array}{c}c_{s+1}\\ \mathrm{}\\ c_{s+d}\end{array}\right).$$
(7)
To see this, let $`Y_k`$ denote scalar projection of elements of $`𝔽_2[q]/()`$ onto $`q^k`$, and let $`X`$ denote multiplication by $`q`$ in $`𝔽_2[q]/()`$. Both of these maps are linear, and it is easy to see that $`Y_k=Y_{k1}X^1+f_kY_0,`$ for $`1kd1`$. Therefore,
$`Y_0`$ $`=Y_{d1}X^1`$
$`=Y_{d2}X^2+f_{d1}Y_0X^1`$
$`=Y_{d3}X^3+f_{d2}Y_0X^2+f_{d1}Y_0X^1`$
$`\mathrm{}`$
$`={\displaystyle \underset{j=0}{\overset{d1}{}}}f_jY_0X^{dj}.`$
Applying $`Y_0`$ to $`q^{s+d}`$ yields $`c_{s+d}=_{j=0}^{d1}f_jc_{s+j}`$, which implies (7). Set $`a_n:=c_n`$ (define $`c`$ on negative subscripts by using the recurrence). Thus, the sequences $`(a_n)`$ and $`(\overline{f}_n)`$ satisfy the same recurrence, with initial conditions $`a_0=\overline{f}_0=1,c_i=a_i=\overline{f}_i=0`$ (for $`d<i<0`$). ∎
Our next proposition computes the density of the reciprocal of every primitive polynomial, and thereby produces an infinite family of polynomials whose reciprocals have density greater than $`1/2`$.
###### Proposition 4.4.
If $`𝒫`$ is a primitive polynomial with degree $`d`$, then $`\delta (\overline{}𝒫)={\displaystyle \frac{2^{d1}}{2^d1}}`$.
###### Proof.
A de Bruijn cycle of order $`d`$ is a binary sequence $`\{S(n)\}_{n=1}^{q^d}`$ in which every binary $`d`$-word appears in a “window” $`(S(n+1),\mathrm{},S(n+d))`$ for some $`j`$ (indices taken modulo $`q^d`$). A reduced de Bruijn cycle is a string of length $`q^d1`$ which achieves every $`d`$-word in some window, except for the word $`0^d`$. Note that a reduced de Bruijn cycle may always be turned into an ordinary de Bruijn cycle by inserting an extra “0” into its longest run of 0’s.
If $`𝒫`$ is primitive, then $`q`$ is a generator of $`𝔽_{2^d}^\times `$, and it is a classical result that the sequence of constant coefficients of the powers of a multiplicative generator yield a reduced de Bruijn cycle. Thus, by Proposition 4.3 the first $`2^d1`$ coefficients of $`\overline{}𝒫`$ are a reduced binary de Bruijn cycle of order $`d`$. The reader wishing to explore de Bruijn cycles further can find the basics in \[MR652466, ShiftRegisterBook\].
Since every string except $`0^d`$ appears in $`\overline{}𝒫`$, there are exactly $`2^{d1}`$ ones in any period. ∎
### 4.3 Polynomials with non-high density reciprocals
We see in Figure 4 that polynomials typically have reciprocals with density near $`1/2`$. In Figure 6, it is apparent that there is a connection between the density of $`\overline{}𝒫`$ and $`\overline{}𝒫^{}`$. Our next theorem elucidates the connection.
###### Proposition 4.5.
If $`\mathrm{ord}(𝒫)4`$, then $`\mathrm{min}\{\delta (\overline{}𝒫),\delta (\overline{}𝒫^{})\}1/2.`$
This proposition is best possible in that $`𝒫_{51}𝒫_{15}=1+q^8`$, and $`\delta (\overline{}𝒫_{51})=\delta (\overline{}𝒫_{15})=1/2`$.
###### Proof.
Set $`D:=\mathrm{ord}(𝒫)`$. We assume without loss of generality that $`\mathrm{deg}(𝒫)D/2\mathrm{deg}(𝒫^{})`$. If $`\mathrm{deg}(𝒫)<3`$, then we appeal to the following table of calculations:
$$\begin{array}{cccc}& & & \\ 𝒫& 𝒫^{}& \overline{}𝒫^{}& \delta (\overline{}𝒫^{})\\ & & & \\ & & & \\ 1& 1+q^D& _{n=0}^{\mathrm{}}q^{nD}& 1/D\\ 1+q& _{n=0}^{D1}q^n& _{n=0}^{\mathrm{}}(q^{nD}+q^{nD+1})& 2/D\\ 1+q^2& _{n=0}^{D1}q^{2n}& _{n=0}^{\mathrm{}}(q^{2nD}+q^{2nD+2})& 1/D\\ 1+q+q^2& (1+q)_{n=0}^{D/31}q^{3n}& _{n=0}^{\mathrm{}}(q^{nD}+q^{nD+1}+q^{nD+2})& 3/D\end{array}$$
In the case $`𝒫=1+q+q^2=\frac{1+q^3}{1+q}`$, we see also that $`D0(mod3)`$, and by hypothesis $`D4`$, so that $`D/31/2`$.
We assume now that $`\mathrm{deg}(𝒫)3`$. Since
$$\overline{}𝒫^{}=\frac{𝒫}{1+q^D}=𝒫+q^D𝒫+q^{2D}𝒫+q^{3D}𝒫+\mathrm{}$$
and $`\mathrm{deg}(𝒫)<D`$, we have $`\delta (\overline{}𝒫^{})=\mathrm{}(𝒫)/D.`$ If $`𝒫`$ has any zero coefficients, then $`\mathrm{}(𝒫)\mathrm{deg}(𝒫)D/2`$ and so $`\delta (\overline{}𝒫^{})1/2`$. If $`𝒫`$ has no zero coefficients, then $`𝒫=(1+q^{\mathrm{deg}(𝒫)+1})/(1+q)`$, in which case $`\overline{}𝒫=(1+q)/(1+q^{\mathrm{deg}(𝒫)+1})`$, a series which has density $`2/(\mathrm{deg}(𝒫)+1)`$. Since $`\mathrm{deg}(𝒫)3`$, this quantity is $`1/2`$. ∎
###### Corollary 4.6.
If $`𝒫\{1,1+q,1+q+q^2\}`$ is a polynomial and $`𝒫^{}`$ is primitive, then $`\delta (\overline{}𝒫)1/2`$.
### 4.4 Eventually periodic sets
An eventually periodic set is one whose generating function has the form $`(q)+\frac{𝒫(q)}{1+q^D}`$, for some polynomials $`,𝒫`$ with $`\mathrm{deg}(𝒫)<D`$, and exactly one of $`,𝒫`$ has constant term 1. The finite sets containing 0 are examples. Another example is given by the set $`\{n:n2(mod4)\}`$ (which has density $`3/4`$), whose reciprocal is the set $`\{1\}\{n:n\text{ is congruent to 0, 2, 5, or 6 modulo 7}\}`$ (which has density $`4/7`$).
###### Proposition 4.7.
The reciprocal of an eventually periodic set is an eventually periodic set.
This proposition is essentially the same as that which asserts that rational numbers have eventually periodic decimal expansions.
###### Proof.
Obviously, the reciprocal of a ratio of polynomials (each with constant term 1) is a ratio of polynomials (each with constant term 1). All that we need to observe is that such a ratio $`/𝒮`$ can be written in the form
$$\frac{}{𝒮}=+\frac{𝒬}{1+q^D},$$
with $`\mathrm{deg}(𝒬)<D`$. By long division, we can write $`/𝒮`$ in the form $`+𝒫/𝒮`$ with $`\mathrm{deg}(𝒫)<\mathrm{deg}(𝒮)`$. But this is the same as $`+\frac{𝒫𝒮^{}}{1+q^D}`$, where $`𝒮𝒮^{}=1+q^D`$, and $`\mathrm{deg}(𝒫𝒮^{})<\mathrm{deg}(𝒮𝒮^{})=D`$. ∎
## 5 The powers of two
We saw in Section 4 that the reciprocal of a polynomial (other than $`𝒫_1`$) has positive density. One might wonder if the reciprocal of any set with zero density has positive density. Our next theorem shows that this is not the case.
We note the $`m=1`$ case of Theorem 5.1: the reciprocal of $`A_1=\{0\}\{2^n:n\}`$ is $`\overline{A}_1=\{2^n1:n\}`$. This is easily proved directly by considering the following sum-preserving involution on $`A_1\times \overline{A}_1`$. For $`s,t`$ and distinct, define $`\mu (0,0)=(0,0)`$, $`\mu (0,2^{t+1}1)=(2^t,2^t1)`$, $`\mu (2^s,2^t1)=(2^t,2^s1)`$, $`\mu (2^t,2^t1)=(0,2^{t+1}1)`$. The existence of this sum-preserving fixed-point-free involution proves that every positive integer $`n`$ can be written in the form $`a+\overline{a}`$, where $`aA_1`$ and $`\overline{a}\overline{A}_1`$, in an even number of ways. A similar proof can be given for $`m=2`$, and presumably for any $`m`$, but quickly grows tedious. We now give an algebraic proof that does not depend on $`m`$.
###### Theorem 5.1.
Let $`m1`$. The reciprocal of the set $`A_m:=\{0\}\{2^{mn}:n\}`$ is the set
$$\overline{A}_m:=\{1+\underset{i=0}{\overset{m1}{}}x_i2^{i+mn_i}:x_i\{0,1\},\stackrel{}{x}\stackrel{}{0},n_i\}.$$
In particular, both $`\delta (A_m,n)`$ and $`\delta (\overline{A}_m,n)`$ are $`O_m\left(\frac{\mathrm{log}n}{n}\right)`$.
###### Proof.
Set $`(q)=_{n0}q^{2^{mn}}`$. By the children’s binomial theorem $`(q^2)=(q)^2`$, and consequently by induction we see that $`(q^{2^m})=(q)^{2^m}`$.
Now, by the definition of $``$, $`(q^{2^m})=(q)+q`$ and so
$`q`$ $`=(q)^{2^m}+(q)`$
$`=\left(1+(q)\right)\left((q)+(q)^2+(q)^3+\mathrm{}+(q)^{2^m1}\right)`$
$`=\left(1+(q)\right)\left(1+{\displaystyle \underset{i=0}{\overset{m1}{}}}\left(1+(q)^{2^i}\right)\right)`$
$`=\left(1+(q)\right)\left(1+{\displaystyle \underset{i=0}{\overset{m1}{}}}\left(1+(q^{2^i})\right)\right)`$
The series $`1+(q)`$ is the generating function of $`\{0\}\{2^{mn}:n\}`$, and $`1+_{i=0}^{m1}(1+(q^{2^i}))`$ is the generating function of $`\{_{i=0}^{m1}x_i2^{i+mn_i}:x_i\{0,1\},\stackrel{}{x}\stackrel{}{0},n_i\}`$, so this identity is equivalent to the theorem. ∎
The reader may be interested to note that the reciprocal of the extremely thick set $`\{2^n:n\}`$ is the thin set $`\{0,3\}\{2^n1,2^n3:n3\}`$, whereas the reciprocal of $`\{4^n:n\}`$ appears to have density $`1/2`$.
Our next theorem shows that the examples given by Theorem 5.1 are extremal. It is impossible for a set and its reciprocal to both grow sub-logarithmically. This result was suggested to us by Ernest Croot \[personal communication\].
###### Theorem 5.2.
Let $`F,\overline{F}`$ be reciprocals (not both $`\{0\}`$), and suppose that $`r`$ is the least positive integer in $`F\overline{F}`$. Then
$$\left|F[0,n]\right|+\left|\overline{F}[0,n]\right|2+\mathrm{log}_2(n/r).$$
###### Proof.
First, note that $`rF\overline{F}`$. Let $`Nr`$, so that neither $`F[1,N)`$ nor $`\overline{F}[1,N)`$ is empty, and let $`m,\overline{m}`$ be the largest elements of those sets. Since $`q^{m+\overline{m}}`$ occurs in the product $`\overline{}`$ at least once, it must occur at least twice. Since $`Nm+\overline{m}<2N`$, we see that
$$\left|F[N,2N)\right|+\left|\overline{F}[N,2N)\right|1.$$
Straightforward counting concludes the proof, since $`F\overline{F}`$ contains 0 twice, and must intersect each of the intervals $`[r,2r)`$, $`[2r,2^2r)`$, $`[2^2r,2^3r)`$, $`\mathrm{}`$. ∎
## 6 Theta functions
Every quadratic that takes integers to integers can be written in the form $`c_0+c_1n+c_2\frac{n(n1)}{2}`$ with $`c_i`$. We wish to study the ranges of such quadratics, but we only wish to consider sets that contain 0; without loss of generality we may take $`c_0=0`$. Thus, we set
$$\mathrm{\Theta }(c_1,c_2):=\{c_1n+c_2\frac{n(n1)}{2}:n\}.$$
Moreover, we are only interested in those sets that consist of nonnegative integers, so we may assume that $`c_2c_10`$. And since $`\mathrm{\Theta }(c_1,c_2)=\mathrm{\Theta }(c_2c_1,c_2)`$ we may also assume that $`c_22c_1`$. Finally, we are only interested in those sets whose $`\mathrm{gcd}`$ is 1: we can assume that $`\mathrm{gcd}(c_1,c_2)=1`$. The only set with $`c_1=0`$ not excluded is $`\mathrm{\Theta }(0,1)=\{\left(\genfrac{}{}{0pt}{}{n}{2}\right):n1\}`$, and the only set with $`c_2=2c_1`$ that is not excluded is $`\mathrm{\Theta }(1,2)=\{n^2:n0\}`$. Otherwise, we have $`c_2>2c_1>0`$, and $`\mathrm{gcd}(c_1,c_2)=1`$.
In Figure 7, we give the number of elements in the reciprocal of $`\mathrm{\Theta }(c_1,c_2)`$ (with $`c_218`$) that are at most $`10^5`$. We note that none of the entries of this table are larger than 50450, and the entries that are less than 49750 are exactly those with $`c_22(mod4)`$. This computation partially justifies Conjecture 1.1.
There is another property of $`\mathrm{\Theta }(c_1,c_2)`$ that happens exactly when $`c_22(mod4)`$: the set $`\mathrm{\Theta }(c_1,c_2)`$ is not uniformly distributed modulo 4.
###### Proposition 6.1.
Let $`\mathrm{gcd}(c_1,c_2)=1`$. The set $`\mathrm{\Theta }(c_1,c_2)`$ is uniformly distributed modulo every power of 2 if and only if $`c_22(mod4)`$.
###### Proof.
First, suppose that $`c_1=2k+1`$ and $`c_2=4\mathrm{}+2`$. Set
$$f(n):=c_1n+c_2\frac{n(n1)}{2}=(2\mathrm{}+1)n^2+2(k\mathrm{})n.$$
If $`k`$ and $`\mathrm{}`$ have the same parity, then $`f(n)(2\mathrm{}+1)n^2(mod4)`$, and since $`n^2`$ takes only two values modulo 4, the set is not uniformly distributed modulo 4. If $`k`$ and $`\mathrm{}`$ have different parity, then
$$(2\mathrm{}+1)n^2+2(k\mathrm{})n(2\mathrm{}+1)n^2+2n(mod4)$$
only takes on the values $`0,3`$ modulo 4. Thus, if $`c_22(mod4)`$, then $`\mathrm{\Theta }(c_1,c_2)`$ is not uniformly distributed modulo 4.
Now suppose that $`c_2=4\mathrm{}`$, and since $`\mathrm{gcd}(c_1,c_2)=1`$, we know that $`c_1`$ is odd. We have
$$f(n):=c_1n+c_2\frac{n(n1)}{2}=2\mathrm{}n^2+(c_12\mathrm{})nn(mod2).$$
The formal derivative of $`f(n)`$ is $`4\mathrm{}n+c_12\mathrm{}0(mod2)`$. By Hensel’s Lemma<sup>f</sup><sup>f</sup>fHensel’s Lemma: If $`f(n)`$ is a polynomial with integer coefficients, and the two congruences $`f(n)a(modp),f^{}(n)0(modp)`$ have a simultaneous solution, then $`f(n)a`$ has a unique solution modulo every power of the prime $`p`$., the range of the polynomial $`f(n)`$ hits every congruence class modulo every power of $`2`$. Since for every $`j`$, $`f(n)`$ is periodic modulo $`2^j`$ with period $`2^j`$, we see that it is uniformly distributed modulo $`2^j`$.
Now suppose that $`c_2=2\mathrm{}+1`$ is odd. Set
$`G`$ $`:=\{(2\mathrm{}+1)m(2m1)+c_1(2m):m\}`$
$`H`$ $`:=\{(2\mathrm{}+1)(2m+1)m+c_1(2m+1):m\}`$
so that $`\mathrm{\Theta }(c_1,2\mathrm{}+1)=GH`$. The set $`G`$ is the range of $`g(m):=f(2m)=(2\mathrm{}+1)m(2m1)+2c_1mm(mod2)`$, which has derivative $`g^{}(m)1(mod2)`$, and the set $`H`$ is the range of $`h(m):=(2\mathrm{}+1)(2m+1)m+c_1(2m+1)m+c_1(mod2)`$, which has derivative $`h^{}(m)1(mod2)`$. Thus, by Hensel’s Lemma, both $`G`$ and $`H`$ exhaust every congruence class modulo $`2^j`$, and by periodicity of $`g(m)`$ and $`h(m)`$ are therefore uniformly distributed modulo $`2^j`$. ∎
### 6.1 The squares
Let $`𝒮(q)=_{n=0}^{\mathrm{}}q^{n^2}`$, and $`S=\{0,1,4,9,16,25,\mathrm{}\}`$. Figure 2 shows $`\delta (\overline{S},x)`$ for two ranges of $`x`$. On the small scale, we see that the relative density behaves irregularly, with many small increases and decreases. On the larger scale, we see that the relative density seems to decrease inexorably.
We characterize completely the values of $`\overline{S}`$ in the residue classes $`0,1,2(mod4)`$.
Let $`\nu _p(n)`$ be the integer such that $`p^{\nu _p(n)}n`$ and $`p^{\nu _p(n)+1}n`$, so that
$$n=\underset{p\text{ prime}}{}p^{\nu _p(n)}$$
for every $`n`$. Let $`r_2(n)`$ be the number of representations of $`n`$ in the form $`y^2+z^2`$, where $`y`$ and $`z`$ are integers.
###### Theorem 6.2.
Let $`n`$. If $`n`$ is even, then $`n\overline{S}`$ if and only if $`n`$ is twice a square. If $`n1(mod4)`$ is not a square, then $`n\overline{S}`$ if and only if $`\nu _p(n)`$ is even for every prime $`p`$ except one, and that prime $`p`$ and $`\nu _p(n)`$ are both congruent to 1 modulo 4. If $`n1(mod4)`$ is a square, then $`n\overline{S}`$ if and only if $`\nu _p(n)2(mod4)`$ for an even number of primes $`p1(mod4)`$.
We will need the following lemmas. The first expresses $`\overline{s}_n`$ in terms of the number of representations of $`n`$ by a particular (depending on $`n`$) quadratic form. The second is quoted without proof from \[MR1083765\], and gives a formula for $`r_2(n)`$.
###### Lemma 6.3.
Let $`n`$, and let $`j`$ satisfy $`n2^j1(mod2^{j+1})`$. Then $`\overline{s}_n=1`$ if and only if
$$\mathrm{\#}\{(k_0,\mathrm{},k_{j1},k_{j+1}):k_i,n=2^{j+1}k_{j+1}^2+\underset{i=0}{\overset{j1}{}}2^ik_i^2\}$$
is odd.
###### Proof.
By Lemma 2.2(ii ), $`\overline{s}_n=1`$ exactly if there are an odd number of tuples $`(k_0,k_1,\mathrm{})`$ with weight
$$n=k_0^2+2k_1^2+4k_2^2+8k_3^2+\mathrm{}.$$
(8)
Let $`w(n)`$ be the number of such tuples. We give a weight-preserving involution $`\mu `$ of such tuples, and deduce the lemma from
$$w(n)\mathrm{\#}(\text{fixed points of }\mu \text{ with weight }n)(mod2).$$
Since $`n2^i1(mod2^{i+1})`$ for $`0i<j`$, reducing (8) modulo $`2,4,\mathrm{},2^j`$ successively tells us that $`k_0,k_1,\mathrm{},k_{j1}`$ are odd, while $`n2^j1(mod2^{j+1})`$ tells us that $`k_j`$ is even. Now define $`J`$ to be the least integer with the two properties: $`Jj+2`$; and $`2k_Jk_j`$.
We define
$$\mu (k_0,k_1,k_2,\mathrm{})=(k_0,k_1,\mathrm{},k_{j1},2k_J,k_{j+1},k_J,k_J,\mathrm{},k_J,k_j/2,k_{J+1},k_{J+2},\mathrm{}),$$
where $`k_J`$ is repeated $`Jj2`$ times. That this is a weight-preserving involution is a routine calculation.
The fixed points of $`\mu `$ are those tuples with $`0=k_j=k_{j+2}=k_{j+3}=\mathrm{}`$. In other words, there is a fixed point for each solution to
$$n=k_0^2+2k_1^2+\mathrm{}+2^{j1}k_{j1}^2+2^{j+1}k_{j+1}^2.\mathit{}$$
###### Lemma 6.4 (\[MR1083765\]\*Theorem 3.22).
If $`\nu _p(n)`$ is odd for any prime $`p`$ congruent to 3 (modulo 4), then $`r_2(n)=0`$. Otherwise, $`r_2(n)=4_p(\nu _p(n)+1)`$, where the product extends over all primes congruent to 1 (modulo 4).
###### Proof of Theorem 6.2.
If $`n`$ is even, then $`n2^01(mod2^{0+1})`$, so we can apply Lemma 6.3 with $`j=0`$ to arrive at $`\overline{s}_n=1`$ if and only if $`n`$ has an odd number of representations of the form $`2k_1^2`$ (with $`k_10`$). Clearly there cannot be more than one such representation, and there is one exactly if $`n`$ is twice a perfect square.
If $`n1(mod4)`$, then we may apply Lemma 6.3 with $`j=1`$ to arrive at $`\overline{s}_n=1`$ if and only if $`n`$ has an odd number of representations of the form $`k_0^2+4k_2^2`$ (with $`k_0`$ and $`k_2`$ nonnegative).
We assume for now that $`n`$ is not a square. Since $`n`$ is odd, there are no such representations with $`k_0=0`$, and since $`n`$ is not a square, there are no such representations with $`k_2=0`$. Thus, every such representation $`k_0^2+4k_2^2`$ gives rise to 8 representations $`\{(\pm k_0)^2+(\pm 2k_2)^2,(\pm 2k_2)^2+(\pm k_0)^2\}`$ of $`n`$ in the form $`y^2+z^2`$. Moreover, any solution to $`n=y^2+z^2`$ must have one of $`y`$ or $`z`$ even and the other odd since $`n`$ is odd, and $`yz`$ since $`n`$ is odd. Since $`n`$ is not a square, neither $`y`$ nor $`z`$ is zero. Every representation $`(y,z)`$ occurs as one of a family of 8 such representations, and one of these has $`n=y^2+z^2=y^2+4(z/2)^2`$ with $`y>0`$ and $`z>0`$. Thus, $`\overline{s}_n=1`$ if and only if $`r_2(n)/8`$ is odd.
By Lemma 6.4, $`r_2(n)/8=0`$ if $`\nu _p(n)`$ is odd for any prime $`p`$ congruent to 3 modulo 4. Otherwise, $`r_2(n)/8=\frac{1}{2}_p(\nu _p(n)+1)`$, where the product extends over those primes that are congruent to 1 modulo 4 (in the remainder of this paragraph, $`p`$ is always 1 modulo 4). First, note that $`\nu _p(n)`$ is odd for some prime $`p`$ since $`n`$ is not a square. If some $`\nu _p(n)`$ is 3 modulo 4 for some $`p`$, then $`r_2(n)/8`$ is even, and similarly if $`\nu _p(n)`$ is 1 modulo 4 for two primes $`p`$. Thus, $`r_2(n)/8`$ is odd precisely if $`\nu _p(n)`$ is odd for exactly one prime, and both that prime and $`\nu _p(n)`$ are 1 modulo 4.
Now we assume that $`n1(mod4)`$ is a square, say $`n=x^2`$. Then, as above, most representations of $`n`$ in the form $`k_0^2+4k_2^2`$ correspond to 8 representations of $`n`$ in the form $`y^2+z^2`$, but the representation $`n=x^2+40^2`$ only corresponds to 4 representations in the form $`y^2+z^2`$. Since $`n`$ is a square, we know that $`\nu _p(n)`$ is even for every prime $`p`$. Thus, $`\overline{s}_n=1`$ if and only if
$$\frac{r_2(n)4}{8}+1$$
is odd. Using the formula from Lemma 6.4, this happens exactly if $`1_p(\nu _p(n)+1)(mod4)`$, where the product extends over primes that are 1 modulo 4. This, in turn, happens exactly when $`\nu _p(n)2(mod4)`$ for an even number of primes $`p1(mod4)`$. ∎
We suspect that $`\delta (\overline{S})=0`$ and that this may follow from the theory of modular forms, but again, this is outside the scope of this paper. We emphasize in Corollary 6.5 that our characterization of $`\overline{S}`$ is consistent with Conjecture 1.1.
###### Corollary 6.5.
The set $`\{n:n\overline{S},n3(mod4)\}`$ has zero density.
###### Proof.
By Theorem 6.2, the set $`\overline{S}`$ clearly has no density in $`0mod2`$. We will use the description given in Theorem 6.2 to show that $`\overline{S}`$ also has zero density in $`1mod4`$.
By the Wiener-Ikehara Theorem (see \[MR2111739\]\*Section 7.2), we have for any set $`A`$ of positive integers
$$\underset{n\mathrm{}}{lim}\delta (A,n)=\underset{s1+}{lim}(s1)\underset{aA}{}a^s.$$
Set $`A=\{n^2p:1n,p\text{ prime}\}`$, and observe that $`\delta (A)=0`$ since
$`\delta (A)`$ $`\underset{s1+}{lim}(s1){\displaystyle \underset{aA}{}}a^s`$
$`=\underset{s1+}{lim}(s1)\left({\displaystyle \underset{p\text{ prime}}{}}(1p^{2s})^1\right)\left({\displaystyle \underset{p\text{ prime}}{}}p^s\right)`$
$`=\underset{s1+}{lim}(s1)\zeta (2s)\left({\displaystyle \underset{p\text{ prime}}{}}p^s\right)`$
$`=\zeta (2)\delta (\text{primes})=0`$
Note that the subset of $`\overline{S}`$ whose elements are even has density 0, and the subset whose elements are congruent to 1 modulo 4 is (except for some squares) contained in $`A`$. Thus
$$\delta (\{n\overline{S}:n3(mod4)\})\delta (\text{squares})+\delta (A)=0.\mathit{}$$
## 7 Prouhet-Thue-Morse numbers
Set $`t_n=1`$ if the binary expansion of $`n`$ contains an even number of “1”s, and set $`t_n=0`$ otherwise. The set $`T:=\{n:t_n=1\}=\{0,3,5,6,9,\mathrm{}\}`$ is called the Prouhet-Thue-Morse sequence. This sequence frequently arises because it simultaneously has enough structure to analyze, and enough “random-like” behavior to be interesting. The survey \[MR1843077\] details four of the occasions that the sequence has been independently rediscovered: first in number theory (Prouhet), then combinatorics (Thue), then in differential geometry (Morse), and finally chess grandmaster Max Euwe rediscovered it to demonstrate that the rules then in use did not imply that chess is a finite game.
For every $`n`$, $`2nT`$ if and only if $`2n+1T`$; thus $`𝒯(q):=_{n=0}^{\mathrm{}}t_nq^n`$ has $`\delta (𝒯)=1/2`$. The sequence $`t_0,t_1,\mathrm{}`$ is not eventually periodic (in fact, the real number with binary expansion $`0.t_0t_1t_2\mathrm{}`$ is transcendental \[MR0457363, MR1869317\]), so $`\overline{}𝒯`$ is not a polynomial. A counting argument \[MR1843077\] reveals the interesting identity:
$$(1+q)^3𝒯(q)^2+(1+q)^2𝒯(q)=q.$$
(9)
Multiplying by $`\overline{}𝒯(q)`$ yields $`q\overline{}𝒯(q)=(1+q+q^2+q^3)𝒯(q)+1+q^2`$, whence for $`n2`$
$$\overline{t}_n=t_{n+1}+t_n+t_{n1}+t_{n2}.$$
This leads reasonably directly (albeit with the modest labor involved in deriving (9)) to a proof of Theorem 7.1. Instead, we give a proof which does not rely on the special form of the functional equation (9), and so is more representative of the process of finding reciprocals.
###### Theorem 7.1.
The reciprocal of the set $`T`$ of Prouhet-Thue-Morse numbers is
$$\overline{T}=\{0\}\{4k\pm 1:\text{the binary expansion of }k1\text{ ends in an even number of “}\text{1}\text{”s}\}.$$
Consequently, $`\delta (\overline{}𝒯)=1/3`$.
If (the binary expansion of) $`k`$ ends in an even number of “1”s, then $`4k+1`$ ends with a string 10<sup>2k+1</sup>1 (a “1” followed by an odd number of “0”s followed by a single “1”), while $`4k1`$ ends with a string 01<sup>2k</sup> (a “0” followed by a positive even number of ”1”s).
###### Proof.
By Lemma 2.1, $`\overline{}𝒯(q)=_{n=0}^{\mathrm{}}r(n)q^n`$, where $`r(n)`$ is the number of ways to write $`n`$ as
$$n=s_0+2s_1+4s_2+8s_3+\mathrm{}+2^ks_k+\mathrm{}$$
where the $`s_k`$ are Prouhet-Thue-Morse numbers. We will build an involution $`\tau `$ on the set of such representations, and $`r(n)`$ will have the same parity as the number of fixed points of $`\tau `$.
By a tuple, we mean an infinite list of Prouhet-Thue-Morse numbers which is 0 from some point on. The weight of a tuple $`(s_0,s_1,\mathrm{})`$ is $`_{n=0}^{\mathrm{}}s_n2^n`$.
We now give the weight-preserving permutation $`\tau `$ of the set of tuples which is actually an involution. The permutation $`\tau `$ has an odd number of fixed points with weight $`n>0`$ if and only if the binary expansion of $`n`$ ends with a string 10<sup>2k+1</sup>1 (a “1” followed by an odd number of “0”s followed by a single “1”) or ends with a string 01<sup>2k</sup> (a “0” followed by a positive even number of ”1”s). These are exactly the numbers of the form $`4k\pm 1`$, where the binary expansion of $`k`$ ends in an even number of “0”s, and this will conclude the proof.
#### Defining the permutation $`\tau `$:
Suppose that $`s_0`$ is even. If $`s_02s_1`$, then set
$$\tau (s_0,s_1,s_2,\mathrm{}):=(2s_1,s_0/2,s_2,s_3,\mathrm{}).$$
If $`s_0=2s_1`$, then let $`i`$ be minimal with $`s_1s_i`$, and set
$$\tau (s_0,s_1,s_2,\mathrm{}):=(2s_i,s_i,s_i,\mathrm{},s_i,s_1,s_{i+1},s_{i+2},\mathrm{}),$$
where $`s_i`$ is repeated $`i1`$ times. The only fixed point with $`s_0`$ even is $`(0,0,\mathrm{})`$ with weight 0.
Now suppose that $`s_03(mod4)`$. Since $`(s_03)/2`$ is even, we can define $`v_0,v_2,v_3,\mathrm{}`$ by
$$(v_0,v_2,v_3,\mathrm{}):=\tau ((s_03)/2,s_2,s_3,\mathrm{}),$$
where the action of $`\tau `$ has already been defined above. Note that $`v_1`$ is not defined, and that $`s_1`$ has not been used. We now set
$$\tau (s_0,s_1,s_2,\mathrm{}):=(2v_0+3,s_1,v_2,v_3,\mathrm{}).$$
The only fixed points with $`s_03(mod4)`$ are the tuples of the form $`(3,s_1,0,0,\mathrm{})`$, where $`s_1`$ is a Prouhet-Thue-Morse number. These fixed points have weight $`3+2s_1`$.
Now suppose that $`s_01(mod4)`$. If there exists an $`L`$ such that $`s_i=(s_0+1)/2`$ for $`0<iL`$ and $`s_i=0`$ for $`i>L`$, then we let $`\tau `$ fix the tuple. These will be the only fixed points of $`\tau `$ with $`s_01(mod4)`$, and will have weight $`2^Ls_0+2^L1`$. Otherwise, if any $`s_i`$ is even (except for the tail of zeros in the tuple $`(s_0,s_1,s_2,\mathrm{})`$), then let $`K:=\mathrm{min}\{i:s_i\text{ even}\}`$, and set
$$\tau (s_0,s_1,s_2,\mathrm{}):=(s_0,s_1,s_2,\mathrm{},s_{K1},\tau (s_K,s_{K+1},s_{K+2},\mathrm{})).$$
If on the other hand all $`s_i`$ are odd (except for the ending string of zeros), then define $`v_0,v_1,v_2,\mathrm{}`$ by
$$(v_0,v_1,v_2,\mathrm{}):=\tau (s_0+1,s_1,s_2,s_3,\mathrm{})$$
and set
$$\tau (s_0,s_1,s_2,\mathrm{}):=(v_01,v_1,v_2,\mathrm{}).$$
That $`\tau `$ is an involution with precisely the claimed fixed points is simply a matter of checking the various cases; we cheerfully leave this important tedium to the reader.
#### Analysis of $`\tau `$’s fixed points with weight $`n`$:
Suppose that $`n`$ is even. By parity considerations, we see that all tuples $`(s_0,s_1,\mathrm{})`$ with weight $`_{i=0}^{\mathrm{}}s_i2^i=n`$ have $`s_0`$ even. Since the only fixed point with $`s_0`$ even is $`(0,0,\mathrm{})`$, we see that $`r(0)=1`$ and $`r(n)`$ is even for all even $`n>0`$. From this point on we assume that $`n`$ is odd.
Suppose that $`n1(mod4)`$, and $`(s_0,s_1,\mathrm{})`$ is a fixed point of $`\tau `$ with weight $`n`$. Since $`n`$ is odd, $`s_0`$ is either 1 or 3 modulo 4. If $`s_01(mod4)`$, then $`s_11(mod2)`$, and such a tuple can be fixed by $`\tau `$ only if $`n=s_0`$, and $`n`$ is a Prouhet-Thue-Morse number. If $`s_03(mod4)`$, then $`s_11(mod2)`$, and such a tuple can be fixed by $`\tau `$ only if $`n=3+2s_1`$, i.e., if $`(n3)/2`$ is a Prouhet-Thue-Morse number (and in this case there is exactly one such tuple). Thus $`\tau `$ has either 0, 1, or 2 fixed points, and we care about when it has an odd number of fixed points. Since $`n1(mod4)`$, the binary expansion of $`n`$ can be written as $`(\text{x10}\text{k}\text{1})_2`$ for some binary string x and positive integer $`k`$. We see that the binary expansion of $`(n3)/2`$ is $`(\text{x01}\text{k})_2`$. Thus, if $`k`$ is even, then either both $`n`$ and $`(n3)/2`$ are Prouhet-Thue-Morse numbers or neither is. If $`k`$ is odd, then exactly one of $`n`$ and $`(n3)/2`$ are Prouhet-Thue-Morse numbers. Hence, $`\tau `$ has an odd number of fixed points exactly if the binary expansion of $`n`$ ends in 10<sup>k</sup>1, with $`k`$ an odd number.
Now suppose that $`n3(mod4)`$, and $`(s_0,s_1,\mathrm{})`$ is a fixed point of $`\tau `$ with weight $`n`$. Since $`n`$ is odd, $`s_0`$ is either 1 or 3 modulo 4. If $`s_01(mod4)`$, then $`s_11(mod2)`$, and such a tuple can be fixed by $`\tau `$ only if $`s_i=(s_0+1)/2`$ for all $`0<iL`$ and $`s_i=0`$ for $`i>L`$. In this case, $`n=2^Ls_0+(2^L1)`$. Since $`s_01(mod4)`$, this implies that the binary expansion of $`n`$ ends with $`L+1`$1”s (in particular, at most one value of $`L`$ can lead to such a fixed point). Moreover, $`2^Ls_0+(2^L1)`$ is a Prouhet-Thue-Morse number if and only if $`L`$ is even. If $`s_03(mod4)`$, then $`s_10(mod2)`$, and such a tuple is fixed if and only if it is of the form $`(3,s_1,0,0,\mathrm{})`$. This can happen exactly if $`(n3)/2`$ is a Prouhet-Thue-Morse number.
Suppose that $`n`$ is a Prouhet-Thue-Morse number. If the binary expansion of $`n`$ ends in exactly $`2k>0`$ “1”s, then $`(3,(n3)/2,0,0,\mathrm{})`$ is the only fixed point of $`\tau `$. If the binary expansion of $`n`$ ends in $`2k+1>0`$1”s, then both $`(3,(n3)/2,0,0,\mathrm{})`$ and
$$(\frac{n2^{2k}+1}{2^{2k}},\frac{n+1}{2^{2k+1}},\frac{n+1}{2^{2k+1}},\mathrm{},0,0,\mathrm{})$$
(the term $`(n+1)/2^{2k+1}`$ is repeated $`2k`$ times) are fixed points.
Now suppose that $`n`$ is not a Prouhet-Thue-Morse number. If the binary expansion of $`n`$ ends in exactly $`2k>0`$1”s, then
$$(\frac{n2^{2k1}+1}{2^{2k1}},\frac{n+1}{2^{2k}},\frac{n+1}{2^{2k}},\mathrm{},0,0,\mathrm{})$$
(the term $`(n+1)/2^{2k}`$ is repeated $`2k1`$ times) is the only fixed point. If the binary expansion of $`n`$ ends in $`2k+1>0`$1”s, then there are no fixed points. ∎
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# Electromagnetic field at Finite Temperature: A first order approach
## 1 Introduction
Quantum Field Theory at Finite Temperature was motivated by the increasing interest in studying the properties of matter under extreme conditions as, for example, at very high temperature or density. The pioneering works joining together the Statistical and Quantum Field Theory were developed mainly by Matsubara in a non relativistic context and, the relativistic case by Fradkin who, via the functional approach, studied the different methods for calculating the Thermal Green’s functions as well as the structure of the Ward identities in QED<sub>4</sub>. Later other works within the Thermal Field Theory appeared whose principal interest was to explore the possibility of restoring some broken symmetries that occur at zero temperature, as for example the $`SU(2)\times U(1)`$ symmetry of weak interaction. The Finite Temperature gauge theories and the problems concerning to the choice of a physical gauge and its dependence was analyzed by Bernard , in particular, the free electromagnetic field.
On the other hand, at zero temperature, there is an alternative way to study the properties for the electromagnetic field which is known as the massless Duffin-Kemmer-Petiau (DKP) theory that is not a trivial limit of the massive Duffin-Kemmer-Petiau theory (DKP) that appears as an alternative formalism for the description of the spin 0 and spin 1 particles in a unified formulation. The DKP theory gives a first order linear equation and it is very similar to the Dirac one but the $`\beta ^\mu `$ matrices satisfy a different algebraic relation. The massive and massless case of the theory were considered in the works where the equivalence of the DKP theory with the theories like Klein-Gordon-Fock (KGF) and Maxwell was proved in Minkokwski space-time and studied in curved space-time such as Riemann and Riemann-Cartan .
At Finite Temperature the massive case is treated in the work where the Bose-Einstein condensation is investigated in the spin 0 sector. The equivalence of many-photons Thermal Green’s functions of the DKP and KGF theories was also proved for the scalar sector calculating the polarization operator at 1-loop order and, in is shown the equivalence of many-gluons Green’s functions in the DKP and KGF Statistical Quantum Field Theories
As above mentioned, all accomplished studies on the massless DKP theory were made at zero temperature. The aim of this work is to study the thermodynamics of the electromagnetic field by using the massless DKP theory. The paper is organized as following: In section 2, we give a brief review of the massless DKP formalism considering a theory with only one real DKP field. In section 3, the constraint analysis is performed for the DKP theory and it has been shown that the model has two first class constraints. In section 4, we calculate the partition function and, finally, we give our conclusions and commentaries.
## 2 The Massless DKP Theory
The massless DKP theory is described in Minkowski space-time by the following Lagrangian density
$$_M=i\overline{\psi }\gamma \beta ^\mu _\mu \psi i_\mu \overline{\psi }\beta ^\mu \gamma \psi \overline{\psi }\gamma \psi ,$$
(1)
where $`\overline{\psi }=\psi ^{}\eta _0`$ with $`\eta _0=\left(2\beta _0^21\right)`$. The $`\beta ^\mu `$ and $`\gamma `$ are singular square matrices satisfying the following algebra
$$\beta ^\mu \beta ^\lambda \beta ^\nu +\beta ^\nu \beta ^\lambda \beta ^\mu =\eta ^{\mu \lambda }\beta ^\nu +\eta ^{\nu \lambda }\beta ^\mu ,$$
(2)
$$\beta ^\mu \gamma +\gamma \beta ^\mu =\beta ^\mu ,\gamma ^2=\gamma .$$
(3)
Due to the singular character of the $`\beta ^\mu `$ matrices the transition of the massive case to the massless theory is non trivial and demands a different treatment. On the other hand the representations of the DKP algebra are reducible and contain the sectors of spin 0 and spin 1 in their structure.
From the Lagrangian (1) we obtain the equation of motion for the massless DKP field
$$i\beta ^\mu _\mu \psi \gamma \psi =0.$$
(4)
It can be shown that the Lagrangian (1) and the massless DKP equation remain invariants under the following gauge transformation
$$\psi \psi ^{}=\psi +\left(1\gamma \right)\mathrm{\Phi },$$
(5)
iff the field $`\mathrm{\Phi }`$ satisfies the condition
$$i\beta ^\mu _\mu \left(1\gamma \right)\mathrm{\Phi }=0.$$
(6)
When the fields under consideration are no charged we have a real DKP field<sup>1</sup><sup>1</sup>1When the DKP field is real, the $`\beta ^\mu `$ matrices must been satisfied $`\beta _0^T=\beta _0,\beta _k^T=\beta _k`$. The representation of the $`\beta ^\mu `$ matrices for the spin 1 sector is given in the appendix A, where we also included the $`\gamma `$ matrix. $`\psi `$, in such situation the Lagrangian (1) takes the following form
$$=\frac{i}{2}\psi ^T\left(\eta ^0\gamma \beta ^\mu \right)_\mu \psi \frac{i}{2}_\mu \psi ^T\left(\eta ^0\beta ^\mu \gamma \right)\psi \frac{1}{2}\psi ^T\left(\eta ^0\gamma \right)\psi $$
(7)
## 3 Constraint Analysis
We proceed the study of the constraint analysis to the real massless DKP theory from the Lagrangian (7) which is written as
$$=\frac{i}{2}\psi ^a\left(\eta ^0\gamma \beta ^\mu \right)_{ab}_\mu \psi ^b\frac{i}{2}_\mu \psi ^a\left(\eta ^0\beta ^\mu \gamma \right)_{ab}\psi ^b\frac{1}{2}\psi ^a\left(\eta ^0\gamma \right)_{ab}\psi ^b.$$
(8)
as usual we define the canonical momentum $`\pi _a`$ as
$$\pi _a=\frac{\delta L}{\delta \dot{\psi }^a}=i\left(\beta ^0\gamma \right)_{ab}\psi ^b,$$
(9)
from which a set of primary constraints appear $`\theta `$
$$\theta _a=\pi _a+i\left(\beta ^0\gamma \right)_{ab}\psi ^b,$$
(10)
because the two different representations for the $`\beta ^\mu `$ matrices we have for the spin $`0`$ sector that $`a=\phi ,0,1,2,3`$ and for spin $`1`$ sector that $`a=`$ $`0,1,2,\mathrm{},9`$.
The canonical Hamiltonian density $`_C`$ that follows from the Lagrangian (8) is given by
$$_C=\frac{i}{2}_k\psi ^a\left(\eta ^0\beta ^k\gamma \right)_{ab}\psi ^b\frac{i}{2}\psi ^a\left(\eta ^0\gamma \beta ^k\right)_{ab}_k\psi ^b+\frac{1}{2}\psi ^a\left(\eta ^0\gamma \right)_{ab}\psi ^b$$
(11)
and considering the set of constraints (10) we have the primary Hamiltonian density $`_P`$as
$$_P=_C+\lambda ^a\theta _a,$$
(12)
where $`\lambda ^a`$ are the Lagrange multiplier.
The Poisson bracket (PB) for the primary constraints results in
$$\{\theta _a\left(𝐱\right),\theta _b\left(𝐲\right)\}=i(\beta ^0)_{ab}\delta \left(𝐱𝐲\right).$$
(13)
To investigate the possibility of obtaining more constraints in the theory we apply the preservation in time of $`\theta _a,`$ i.e.
$$\dot{\theta }_a\left(𝐱\right)=\{\theta _a\left(𝐱\right),H_P\},$$
(14)
where $`H_P={\displaystyle d^3𝐳_P}`$ is the primary Hamiltonian. Thus, the stability condition provide
$$\dot{\theta }_a=i\left(\eta ^0\beta ^k\right)_{ab}_k\psi ^b\left(\eta ^0\gamma \right)_{ab}\psi ^b+i\beta _{ab}^0\lambda ^b0,$$
(15)
being $`\beta ^\mu `$ singular matrices we conclude that not all $`\lambda `$ coefficients can be obtained from the relation (15) and more constraints appear. These new constraints are selected by means of the projector $`𝐌=\left(1\beta _0^2\right)`$, thus, we obtain
$$\theta _a^{(2)}=\left(𝐌\beta ^k\right)_{ac}i_k\psi ^c\left(𝐌\gamma \right)_{ac}\psi ^c$$
(16)
that are a set of secondary constrains. When the preservation in time of these secondary constraints is imposed no more constraints appear in the theory.
Now we calculate the PB for all primary and secondary constraints
$`\{\theta _a\left(𝐱\right),\theta _b^{(2)}\left(𝐲\right)\}`$ $`=`$ $`\left[\left(i\beta ^k_k^x+\gamma \right)𝐌\right]_{ab}\delta \left(𝐱𝐲\right),`$
$`\{\theta _b^{(2)}\left(𝐱\right),\theta _a\left(𝐲\right)\}`$ $`=`$ $`\left[𝐌\left(i\beta ^k_k^x\gamma \right)\right]_{ab}\delta \left(𝐱𝐲\right),`$ (17)
$`\{\theta _a^{(2)}\left(𝐱\right),\theta _b^{(2)}\left(𝐲\right)\}`$ $`=`$ $`0.`$
We can write all set of constraints as $`\zeta _a=\{\theta _a,\theta _a^{(2)}\}`$ such that its matrix is
$$\{\zeta _a\left(𝐱\right),\zeta _b\left(𝐲\right)\}=\left[\begin{array}{ccc}i\beta ^0& & (i\beta ^k_k^x+\gamma )𝐌\\ & & \\ 𝐌(i\beta ^k_k^x\gamma )& & 0\end{array}\right]\delta (𝐱𝐲)$$
(18)
and the determinant of $`\{\zeta _a,\zeta _b\}`$ is zero. But the constraint structure depends of the chosen representation for the DKP algebra as we will see go on.
### 3.1 Spin 1 sector<sup>2</sup><sup>2</sup>2For the spin 0, the representation for the $`\beta ^\mu `$ matrices is $`5\times 5`$. The rank of the constraint matrix $`\{\zeta _a,\zeta _b\}`$ is 8, the null space has two trivial constraints which are irrelevant, and the set of constraints is second class.
For this sector the representation of the $`\beta ^\mu `$ matrices is $`10\times 10`$ and the spin 1 DKP field is a column matrix with ten real components
$$\psi =(\psi ^0,\psi ^1,\psi ^2,\psi ^3,\psi ^4,\psi ^5,\psi ^6,\psi ^7,\psi ^8,\psi ^9)^T.$$
(19)
Consequently, in this sector, from (10) we observe that there are ten primary constraints and write explicitly
$`\theta _0=\pi _0,`$ (20)
$`\theta _1=\pi _1+\psi ^7,\theta _2=\pi _2+\psi ^8,\theta _3=\pi _3+\psi ^9,`$ (21)
$`\theta _n=\pi _n,n=4,5,6,7,8,9.`$ (22)
And from (16) we obtain four secondary constraints given by
$`\theta _0^{(2)}=_1\psi ^7_2\psi ^8_3\psi ^9,`$ (23)
$`\theta _4^{(2)}=_3\psi ^2_2\psi ^3\psi ^4,\theta _5^{(2)}=_1\psi ^3_3\psi ^1\psi ^5,\theta _6^{(2)}=_2\psi ^1_1\psi ^2\psi ^6,`$ (24)
To classify these constraints as first and second class we perform the calculation of the PB between all these primary and secondary constraints such as it is shown by the matrix $`\{\zeta _a\left(𝐱\right),\zeta _b\left(𝐲\right)\}`$ in (18). In our case, the rank of the matrix $`\{\zeta _a,\zeta _b\}`$ is 12 which is the number of second class constraints and, the dimension of the null space is 2 that gives the number of the first class constraints. The null space is formed by the constraint $`\theta _0=\pi _0`$ and by the linear combination of second class constraints $`_1\theta _1+_2\theta _2+_3\theta _3+\theta _0^{\left(2\right)}`$ that defines another first class constraint $`G=_k\pi _k`$. Thus, we obtain two first class constraints
$$\theta _0=\pi _0,G=_k\pi _k,k=1,2,3$$
(25)
and twelve second class constraints given by the equations (21), (22) and (24).
The projectors of the spin 1 sector are defined as
$`R^\mu `$ $`=`$ $`\left(\beta ^1\right)^2\left(\beta ^2\right)^2\left(\beta ^3\right)^2\left[\beta ^\mu \beta ^0\eta ^{\mu 0}\right]`$
$`R^{\mu \nu }`$ $`=`$ $`R^\mu \beta ^\nu ,\mu ,\nu =0,1,2,3.`$
such that the field $`\psi ^\mu =R^\mu \psi `$ is a Lorentz vector and $`\psi ^{\mu \nu }=R^{\mu \nu }\psi `$ is an antisymmetric second-rank Lorentz tensor; when multiplied with the $`\gamma `$ matrix we also get $`R^\mu \gamma \psi =0`$ and $`R^{\mu \nu }\gamma \psi =R^{\mu \nu }\psi `$.
Using the projectors $`R_\mu `$ and $`R_{\mu \nu }`$ and, from the relation (5) we conclude that only the vector components of the DKP field are transformed as it is shown to follow
$`\psi _\mu ^{}`$ $`=`$ $`\psi _\mu +\mathrm{\Phi }_\mu ,`$
$`\psi _n^{}`$ $`=`$ $`\psi _n,n=4,5,6,7,8,9,`$
and from (6) we get $`R_\mu \mathrm{\Phi }=\mathrm{\Phi }_\mu =\pm _\mu \mathrm{\Lambda },`$ being $`\mathrm{\Lambda }`$ an arbitrary scalar function, thus, we can conclude that the theory under consideration is a local $`U\left(1\right)`$ gauge field theory. Then, we impose the following gauge fixing conditions
$$\mathrm{\Omega }_1=_k\psi ^k,\mathrm{\Omega }_2=\psi ^0,$$
(28)
such that the set $`\chi _A^{}=\{\theta _0,G,\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ is second class, thus, the PB matrix of the set is given by
$$D_{A^{}B^{}}=\{\chi _A^{}\left(𝐱\right),\chi _A^{}\left(𝐲\right)\}=\left[\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& _x& 0\\ 0& _x& 0& 0\\ 1& 0& 0& 0\end{array}\right]\delta \left(𝐱𝐲\right)$$
(29)
where $`=\left(_k\right)^2=\left(_1\right)^2+\left(_2\right)^2+\left(_3\right)^2`$, computing the functional determinant we get
$$detD_{A^{}B^{}}=\left[det\right]^2.$$
(30)
## 4 The Partition Function
Now we study the thermodynamic equilibrium of the electromagnetic field using the DKP formalism. Such as we see the constraint analysis of the theory gives for the spin 1 sector two first class constraints characterizing a local $`U(1)`$ gauge theory. Then, we write the partition function for the massless real DKP field using the Hamiltonian formalism
$$Z=\underset{periodic}{}𝒟\psi 𝒟\pi \delta \left(\mathrm{\Theta }_A\right)\delta \left(\chi _A^{}\right)\left(detC_{AB}\right)^{1/2}\left(detD_{A^{}B^{}}\right)^{1/2}\mathrm{exp}\left\{_\beta d^4x\left(i\pi _a_\tau \psi ^a_C\right)\right\},$$
(31)
where $`_C`$ is given by (11) and
$$𝒟\psi =𝒟\psi ^a,𝒟\pi =𝒟\pi _a,a=0,1,2,\mathrm{},9.$$
(32)
The fields $`\psi `$ are restricted by the periodicity condition
$$\psi (0,𝐱)=\psi (\beta ,𝐱),$$
(33)
where $`\mathrm{\Theta }_A=\{\theta _a,\theta _b^{(2)}\}`$ is the set of second class constraints given by the equations (21), (22) and (24); the set $`\chi _A^{}=\{\theta _0,G,\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ is given by the set of first class constraints (25) and its respective gauge fixing conditions (28). The matrix $`C_{AB}=\{\mathrm{\Theta }_A,\mathrm{\Theta }_B\}`$ can be obtained from (18) and its determinant is $`detC_{AB}=1`$ . The matrix $`D_{A^{}B^{}}`$ $`=\{\chi _A^{},\chi _B^{}\}`$ and its determinant $`detD_{A^{}B^{}}`$ are given by the equations (29) and (30), respectively.
But, it is interesting to perform the calculation of the partition function in a manifest covariant way. Thus, it is possible to show that the equation (31) becomes
$$Z=N(\beta )\underset{periodic}{}𝒟\psi \delta \left(F\left[\psi ^A\right]\right)det\left|\frac{F^g}{\mathrm{\Lambda }}\right|\mathrm{exp}\left\{_\beta d^4x\left[\frac{1}{2}\psi ^T\eta ^0\left(i\beta ^A_A\gamma \right)\psi \right]\right\},$$
(34)
where $`F\left[\psi ^A\right]`$ is an arbitrary gauge fixing condition. Here we consider
$$F\left[\psi ^A\right]=\frac{1}{\sqrt{\zeta }}_A\psi ^Af,F^g\left[\psi ^A\right]=F\left[\psi ^A\right]\frac{1}{\sqrt{\zeta }}\mathrm{}\mathrm{\Lambda }.$$
(35)
with the gauge transformation $`\psi ^A\psi ^A^A\mathrm{\Lambda }`$ , and $`f`$ is an arbitrary scalar function.
It is worthwhile to note that (34) is exactly the Faddeev-Popov technique used to quantize a local gauge theory. Consequently, the equation (34) can be expressed as being
$$Z=N(\beta )\underset{periodic}{}𝒟\psi det\left|\frac{1}{\sqrt{\zeta }}\mathrm{}\right|\mathrm{exp}\left\{_\beta d^4x\left[\frac{1}{2}\psi ^T\eta ^0\left(i\beta ^A_A\gamma \right)\psi \frac{1}{2\zeta }\left(_A\psi ^A\right)^2\right]\right\}.$$
(36)
where the index $`A=\tau ,1,2,3`$ and $`\beta ^A_A=i\beta ^0_\tau +\beta ^k_k`$. We remark that at zero temperature the $`det\left(\mathrm{}\right)`$ is a constant that can be ignored, however, at Finite Temperature it turns out a very important temperature dependent term. Using the projectors $`R^A`$ we rewrite the gauge fixing term in matrix form such that the partition function reads as
$$Z=N\left(\beta \right)det\left|\frac{1}{\sqrt{\zeta }}\mathrm{}\right|Z^{},$$
(37)
where
$$Z^{}=\underset{periodic}{}𝒟\psi \mathrm{exp}\left\{_\beta d^4x\left[\frac{1}{2}\psi ^T\eta ^0\left(i\beta ^A_A\gamma +\frac{1}{\zeta }\eta ^0\left(R^A\right)^TR^B_A_B\right)\psi \right]\right\},$$
(38)
and $`R^A_A=iR^0_\tau +R^k_k`$.
To perform the calculation of the functional integral we use the Fourier series of the DKP field
$$\psi (\tau ,𝐱)=\frac{1}{\beta }\underset{n}{}\frac{d^3𝐩}{\left(2\pi \right)^3}\stackrel{~}{\psi }(\omega _n,𝐩)e^{i\left(\mathrm{𝐩𝐱}+\omega _n\tau \right)},$$
(39)
with $`\omega _n=2\pi n/\beta `$ and the periodicity conditions (33) for the DKP field are imposed.
Then, substituting (39) in (38) we get
$`Z^{}`$ $`=`$ $`{\displaystyle \underset{n,𝐩}{}}det\left[i\omega _n\beta ^0+p_k\beta ^k+\gamma +{\displaystyle \frac{1}{\zeta }}\eta ^0\left(i\omega _nR^0+p_kR^k\right)^T\left(i\omega _nR^0+p_kR^k\right)\right]^{1/2}`$ (40)
$`=`$ $`{\displaystyle \underset{n,𝐩}{}}\left[{\displaystyle \frac{\left(\omega _n^2+\omega _𝐩^2\right)^4}{\zeta }}\right]^{1/2}={\displaystyle \underset{n,𝐩}{}}\left(\omega _n^2+\omega _𝐩^2\right)^2\sqrt{\zeta },`$
where $`\omega _𝐩=\left|𝐩\right|`$. The determinant in (37) is
$$det\left|\frac{1}{\sqrt{\zeta }}\mathrm{}\right|=\underset{n,𝐩}{}\frac{\left(\omega _n^2+\omega _𝐩^2\right)^1}{\sqrt{\zeta }}.$$
(41)
Finally, the partition function reads
$$Z=N\left(\beta \right)\underset{n,𝐩}{}\left(\omega _n^2+\omega _𝐩^2\right)^1,$$
(42)
as we are interesting in $`\mathrm{ln}Z`$, then, from the last expression we obtain
$$\mathrm{ln}Z=\mathrm{ln}N\left(\beta \right)+2V\underset{n}{}\frac{d^3𝐩}{\left(2\pi \right)^3}\mathrm{ln}\beta V\underset{n}{}\frac{d^3𝐩}{\left(2\pi \right)^3}\mathrm{ln}\beta ^2\left(\omega _n^2+\omega _𝐩^2\right).$$
(43)
The value for the $`N\left(\beta \right)`$ is selected in a manner that it cancels the divergent term, i.e we choose
$$\mathrm{ln}N\left(\beta \right)=2V\underset{n}{}\frac{d^3𝐩}{\left(2\pi \right)^3}\mathrm{ln}\beta ,$$
(44)
next, we perform the sum (see, for example, ) and get the following expression for the partition function
$$\mathrm{ln}Z=2V\frac{d^3𝐩}{\left(2\pi \right)^3}\left[\frac{\beta \omega _𝐩}{2}+\mathrm{ln}\left(1e^{\beta \omega _𝐩}\right)\right],$$
(45)
which describe a massless bosonic field with two polarization states that is the characteristic behavior of the electromagnetic field in thermodynamical equilibrium.
## 5 Conclusions
In this work we study the massless DKP theory at Finite Temperature, and it is shown the constraint structure of the model leads to conclude that it is a local $`U(1)`$ gauge theory in its spin 1 sector. Such analysis allow to construct the correct partition function using the Hamiltonian procedure. Also, we show that it is possible to arrive to the covariant expression which is exactly the covariant quantization of a gauge theory using the Faddeev-Popov approach. Consequently, the partition function of the spin 1 sector gives the partition function of a zero-mass Bose gas with two polarization states, i.e. the electromagnetic field modes in thermodynamical equilibrium.
The perspectives to be followed are to study the Finite Temperature properties of the massless DKP field in curved space-time and, consequently, analyze the curvature effects in the thermodynamics of the electromagnetic field via the DKP formalism. Advances in this directions will be reported elsewhere.
### Acknowledgements
RC (grant 01/12611-7) thanks FAPESP for full support. BMP thanks CNPq and FAPESP (grant 02/00222-9) for partial support.
## Appendix A Spin 1 representation for the massless DKP algebra
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# Affine Toda field theories related to Coxeter groups of non-crystallographic type
## 1 Introduction
The Ising model is generally considered as the prime example of integrable models. When viewed in the continuous limit as a $`c=1/2`$ conformal field theory , it is a well known fact that it can be realized as an $`E_8^{(1)}E_8^{(1)}/E_8^{(2)}`$-coset model . Even when the conformal symmetry is broken, by perturbing the theory with a primary field of scaling dimension ($`1/16,1/16`$) , the $`E_8`$ structure survives in form of a (minimal) $`E_8`$-affine Toda field theory (ATFT) . It will be one of the results in this paper to show that there is an even more fundamental structure than $`E_8`$ underlying this particular model, the non-crystallographic Coxeter group $`H_4`$. We draw here on the observation made first by Sherbak in 1988 , namely that $`H_4`$ can be embedded into $`E_8`$, see also for further developments of these mathematical structures. Loosely speaking, one may regard the $`E_8`$-theory as two copies of $`H_4`$-theories. We get a first glimpse of this structure from a more physical point of view when we bring the mass spectrum of minimal $`E_8`$-affine Toda field theory found originally in into the form
$$\begin{array}{cccc}m_1=1,\hfill & m_2=2\mathrm{cos}\frac{\pi }{30},\hfill & m_3=\sqrt{\mathrm{sin}\frac{11\pi }{30}/\mathrm{sin}\frac{\pi }{30}},\hfill & m_4=2\varphi \mathrm{cos}\frac{7\pi }{30},\hfill \\ m_5=\varphi m_1,\hfill & m_6=\varphi m_2,\hfill & m_7=\varphi m_3,\hfill & m_8=\varphi m_4.\hfill \end{array}$$
(1)
We have set here the overall mass scale to one. Remarkably, these mass ratios are the same in the classical as well as in the quantum theory, as all masses renormalize with an overall factor, see and references therein. We observe here that there are four “fundamental” masses present in the theory, whereas the other ones can be obtained simply by a multiplication with the golden ratio
$$\varphi =\frac{1}{2}(1+\sqrt{5})=\varphi ^21.$$
(2)
It will turn out that each of the sets ($`m_1,m_2,m_3,m_4`$) and ($`m_5,m_6,m_7,m_8`$) can be associated with an $`H_4`$-ATFT.
The other popular integrable quantum field theory is the sine-Gordon model, see e.g. and references therein. It is a well known fact that once the coupling constant $`\nu `$ is taken to be $`1/\nu =n`$, with $`n`$ being an integer, the backscattering amplitudes vanish and the theory reduces to a minimal $`D_{n+1}`$-ATFT. In particular for $`n=5`$ we find a similar pattern for the mass ratios as discussed above. The $`D_6`$-ATFT mass spectrum reads up to an overall mass scale
$$\begin{array}{ccc}m_1=\varphi ^1,\hfill & m_2=\sqrt{1+\varphi ^2},\hfill & m_3=1,\hfill \\ m_4=\varphi m_1,\hfill & m_5=\varphi m_2,\hfill & m_6=\varphi m_3.\hfill \end{array}$$
(3)
It this case the sets ($`m_1,m_2,m_3`$) and ($`m_4,m_5,m_6`$) can be associated with an $`H_3`$-ATFT.
The above mentioned structure can be explained simply by the fact that $`H_4`$ can be embedded into $`E_8`$ and $`H_3`$ into $`D_6`$, such that the non-crystallographic structure is “visible” inside the theories related to crystallographic Coxeter groups.
Besides having a “non-crystallographic pattern” inside theories related to crystallographic Coxeter groups, it is interesting to ask the question whether it is possible to construct theories purely based on these latter groups. In particular, $`H_3`$ being a three-dimensional symmetry group of the icosahedron, a regular solid with 20 triangle faces, finds a natural application in physical , chemical an even biological systems . In the context of integrable (solvable) models, Calogero-Sutherland models have been formulated based also on $`H_3`$ and it should be possible to extend these investigations to other non-crystallographic Coxeter groups. However, so far no ATFT for such type of group has been considered, the main reason being that unlike for crystallographic ones, in this case there is no Lie algebra at disposal, which is vital in that context. This deficiency can be overcome by exploiting the embedding structure and reduce the theories associated to crystallographic Coxeter groups to new types of theories related entirely to non-crystallographic ones. This is somewhat similar in spirit to the folding procedure carried out by Olive and Turok , who constructed ATFT for non-simply laced Lie algebras from those related to simply laced algebras, exploiting the embedding of the former into the latter. However, in comparison, there is one crucial difference. Whereas in the folding scenario the reduced models are identical to a formulation purely in terms of non-simply laced algebras, the models we obtain here vitally rely in their construction on the embedding and can not be formulated directly in terms of non-crystallographic Coxeter groups on the level of the Lagrangian.
Our manuscript is organized as follows: In section 2 we review and develop the mathematics associated to the embedding of non-crystallographic into crystallographic Coxeter groups. In section 3 we apply these notions to affine Toda field theory and construct in particular their classical mass spectra and fusing structures. In addition we start the development of a quantum field theory by computing the mass renormalisations. We state our conclusions in section 4. In the appendix we present explicit computations of various orbits of coloured simple roots related to non-crystallographic and crystallographic Coxeter groups and exhibit how they can be embedded into one another.
## 2 Embedding of non-crystallographic into crystallographic Coxeter groups
Coxeter graphs are finite graphs, whose edges are labelled by some integers $`m_{ij}`$ joining the vertices $`i`$ and $`j`$ . To each of these graphs one can associate a finite reflection group. When the crystallographic condition is satisfied, that is $`m_{ij}`$ for $`i`$ $`j`$ takes only the values $`2,3,4`$ or $`6`$, these groups are Weyl groups. In contrast to the non-crystallographic groups, the crystallographic ones can be related to Lie algebras and Lie groups. Lie theory is exploited largely in the context of integrable models, which is one of the main reasons why crystallographic groups enjoy wider applications. Here we use the embedding of non-crystallographic groups into crystallographic ones, such that we can still exploit the Lie structure related to the larger groups.
### 2.1 Root systems
In order to assemble the necessary mathematical notions, we start by introducing a map $`\omega `$ from a root system$`\mathrm{\Delta }^c`$ which is invariant under the action of a crystallographic Coxeter group $`𝒲`$ of rank $`\mathrm{}`$ into the union of two sets $`\stackrel{~}{\mathrm{\Delta }}^{nc}`$ related to a non-crystallographic group $`\stackrel{~}{𝒲}`$ of rank $`\stackrel{~}{\mathrm{}}=\mathrm{}/2`$
$$\omega :\mathrm{\Delta }^c\stackrel{~}{\mathrm{\Delta }}^{nc}\varphi \stackrel{~}{\mathrm{\Delta }}^{nc}.$$
(4)
Throughout this manuscript we adopt the notation that quantities related to crystallographic and non-crystallographic groups are specified with the same symbol, distinguished by an additional tilde, e.g. $`\alpha \mathrm{\Delta }^c=\mathrm{\Delta }`$ and $`\stackrel{~}{\alpha }\stackrel{~}{\mathrm{\Delta }}^{nc}=\stackrel{~}{\mathrm{\Delta }}`$. Introducing a special labelling for the vertices on the Coxeter graphs, or equivalently the simple roots, we can always realize this map as
$$\alpha _i\omega (\alpha _i)=\{\begin{array}{cc}\stackrel{~}{\alpha }_i\hfill & \text{for }1i\stackrel{~}{\mathrm{}}=\mathrm{}/2\hfill \\ \varphi \stackrel{~}{\alpha }_{i\stackrel{~}{\mathrm{}}}\hfill & \text{for }\stackrel{~}{\mathrm{}}<i\mathrm{}.\hfill \end{array}$$
(5)
Our labelling allows for a generic treatment of the embedding and differs for instance from the one used in . A further important property guaranteed by our conventions is $`\alpha _i\alpha _{i+\stackrel{~}{\mathrm{}}}=0`$. Both types of root systems $`\mathrm{\Delta }`$ and $`\stackrel{~}{\mathrm{\Delta }}`$ are equipped with a symmetric bilinear form or inner product. In Moody and Patera noticed the remarkable fact that the map $`\omega `$ is an isometric isomorphism, such that we may compute inner products in the root system $`\mathrm{\Delta }`$ from inner products in $`\stackrel{~}{\mathrm{\Delta }}`$
$$\alpha \beta =R(\omega (\alpha )\omega (\beta )).$$
(6)
Here the map $`R`$, called a rational form relative to $`\varphi `$, extracts from a number of the form<sup>1</sup><sup>1</sup>1Note that higher powers of $`\varphi `$ can be reduced to that form by a repeated use of (2), e.g. $`\varphi ^2=1+\varphi `$, $`\varphi ^3=1+2\varphi `$, $`\varphi ^4=2+3\varphi `$, $`\varphi ^5=3+5\varphi `$,$`\mathrm{}`$, $`\varphi ^n=f_{n1}+\varphi f_n`$ where $`f_n`$ is the $`n`$-th Fibonacci number obeying the recursive relation $`f_{n+1}=`$ $`f_n+`$ $`f_{n1}`$. $`a+\varphi b`$ with $`a,b`$ the rational part $`a`$
$$R(a+\varphi b)=a.$$
(7)
We normalize all our roots to have length 2, such that $`\alpha ^2=`$ $`\stackrel{~}{\alpha }^2=2`$ for $`\alpha \mathrm{\Delta }`$ and $`\stackrel{~}{\alpha }\stackrel{~}{\mathrm{\Delta }}`$. According to (6) we may therefore compute the Cartan matrix related to $`\mathrm{\Delta }`$ entirely from inner products in $`\stackrel{~}{\mathrm{\Delta }}`$
$$K_{ij}=\frac{2\alpha _i\alpha _j}{\alpha _j\alpha _j}=R(\omega (\alpha _i)\omega (\alpha _j))=R(\varphi ^{t_i+t_j}\stackrel{~}{\alpha }_{it_i\stackrel{~}{\mathrm{}}}\stackrel{~}{\alpha }_{jt_j\stackrel{~}{\mathrm{}}}).$$
(8)
where $`t_i=0`$ for $`1i\stackrel{~}{\mathrm{}}`$ and $`t_i=1`$ for $`\stackrel{~}{\mathrm{}}<i\mathrm{}`$. For our purposes it will be most important to achieve also the opposite, which can not be found in , namely to compute inner products in $`\stackrel{~}{\mathrm{\Delta }}`$ from those in $`\mathrm{\Delta }`$. For this aim we introduce here the map
$$\stackrel{~}{\omega }:\stackrel{~}{\mathrm{\Delta }}\mathrm{\Delta }\varphi \mathrm{\Delta },$$
(9)
which acts on the simple roots in $`\stackrel{~}{\mathrm{\Delta }}`$ as
$$\stackrel{~}{\alpha }_i\stackrel{~}{\omega }(\stackrel{~}{\alpha }_i)=(\alpha _i+\varphi \alpha _{i+\stackrel{~}{\mathrm{}}})\text{for }1i\stackrel{~}{\mathrm{}}.$$
(10)
Note that $`\stackrel{~}{\omega }(\stackrel{~}{\alpha }_i)^2=2+2\varphi ^2`$, such that $`\stackrel{~}{\omega }`$ is not an isometry. Instead, we find that inner products in $`\stackrel{~}{\mathrm{\Delta }}`$ are related to inner products in $`\mathrm{\Delta }`$ by means of
$$\omega (\alpha _i)\stackrel{~}{\alpha }_j=\alpha _i\stackrel{~}{\omega }(\stackrel{~}{\alpha }_j)\text{for }1j\stackrel{~}{\mathrm{}},1i\mathrm{}.$$
(11)
Expanding (11) yields immediately a relation between the Cartan matrices of $`\mathrm{\Delta }`$ and $`\stackrel{~}{\mathrm{\Delta }}`$
$$K_{ij}+\varphi K_{i(j+\stackrel{~}{\mathrm{}})}=\{\begin{array}{cc}\stackrel{~}{K}_{ij}\hfill & \text{for }1i\stackrel{~}{\mathrm{}}\hfill \\ \varphi \stackrel{~}{K}_{(i\stackrel{~}{\mathrm{}})j}\hfill & \text{for }\stackrel{~}{\mathrm{}}<i\mathrm{}.\hfill \end{array}$$
(12)
Noting that for the Coxeter groups we consider the Cartan matrices are symmetric, such that the relations (12) also hold with $`ij`$. As the right hand side involves only one inner product in $`\stackrel{~}{\mathrm{\Delta }}`$, this formula achieves our objective and we may now express inner products in $`\stackrel{~}{\mathrm{\Delta }}`$ in terms of those in $`\mathrm{\Delta }`$. Recalling that simple roots and fundamental weights $`\lambda _i`$ are related as $`\alpha _i=_jK_{ij}\lambda _j`$ and $`\stackrel{~}{\alpha }_i=_j\stackrel{~}{K}_{ij}\stackrel{~}{\lambda }_j`$, it follows directly that (11) also holds when we replace simple roots by fundamental weights. As $`K_{ij}^1=\lambda _i\lambda _j`$, this means that (12) also holds for the inverse matrices
$$K_{ij}^1+\varphi K_{i(j+\stackrel{~}{\mathrm{}})}^1=\{\begin{array}{cc}\stackrel{~}{K}_{ij}^1\hfill & \text{for }1i\stackrel{~}{\mathrm{}}\hfill \\ \varphi \stackrel{~}{K}_{(i\stackrel{~}{\mathrm{}})j}^1\hfill & \text{for }\stackrel{~}{\mathrm{}}<i\mathrm{}\hfill \end{array}$$
(13)
and (5), (10) for $`\alpha _i\lambda _i`$, $`\stackrel{~}{\alpha }_i`$ $`\stackrel{~}{\lambda }_i`$. With regard to our application it will be particularly important to relate the eigensystems $`K`$ and $`\stackrel{~}{K}`$. Abbreviating $`\kappa _{ij}=K_{ij}`$ and $`\widehat{\kappa }_{ij}=K_{i(j+\stackrel{~}{\mathrm{}})}`$ for $`1i,j\stackrel{~}{\mathrm{}}`$ we can block-decompose the Cartan matrix $`K`$ further and verify that
$$U^1KU=U^1\left(\begin{array}{cc}\kappa & \widehat{\kappa }\\ \widehat{\kappa }& \kappa +\widehat{\kappa }\end{array}\right)U=\left(\begin{array}{cc}\stackrel{~}{K}(\varphi )& 0\\ 0& \stackrel{~}{K}(\varphi ^1)\end{array}\right)\text{with }U=\left(\begin{array}{cc}\hfill 𝕀& \hfill 𝕀\\ \hfill \varphi 𝕀& \hfill 𝕀/\varphi \end{array}\right),$$
(14)
where $`𝕀`$ is the $`\stackrel{~}{\mathrm{}}\times \stackrel{~}{\mathrm{}}`$ unit matrix. This means of course that the $`\mathrm{}`$ eigenvalues of $`K`$, which are known to be of the form $`e_n=4\mathrm{sin}^2(\pi s_n/2h)`$ with $`h`$ being the Coxeter number (see below for more details). The $`s_n`$ label the $`\mathrm{}`$ exponents of $`𝒲`$ and are identical to the union of the $`\stackrel{~}{\mathrm{}}`$eigenvalues of $`\stackrel{~}{K}(\varphi )`$ and $`\stackrel{~}{K}(\varphi ^1)`$
$$𝒮=\{s_1,s_2,\mathrm{},s_{\mathrm{}}\}=\stackrel{~}{𝒮}(\varphi )\stackrel{~}{𝒮}^{}(\varphi ^1).$$
(15)
The eigenvalues are invariant under the change $`\varphi \varphi `$. Labelling the $`\mathrm{}`$ eigenvectors of $`K`$ by $`y_n`$, such that $`Ky_n=e_ny_n`$, we can construct the eigenvectors $`\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )`$ and $`\stackrel{~}{y}_{\stackrel{~}{s}^{}}(\varphi ^1)`$ of $`\stackrel{~}{K}(\varphi )`$ and $`\stackrel{~}{K}(\varphi ^1)`$, respectively, from (14) as
$$U^1y_s=\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )\stackrel{~}{y}_{\stackrel{~}{s}^{}}(\varphi ^1).$$
(16)
Conversely, we can construct the eigenvectors of $`K`$ from the knowledge of the eigenvectors of $`\stackrel{~}{K}(\varphi )`$ and $`\stackrel{~}{K}(\varphi ^1)`$
$$y_{\stackrel{~}{s}}=\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )\varphi \stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )\text{and }y_{\stackrel{~}{s}^{}}=\stackrel{~}{y}_{\stackrel{~}{s}^{}}(\varphi ^1)(\varphi ^1)\stackrel{~}{y}_{\stackrel{~}{s}^{}}(\varphi ^1).$$
(17)
The first identity follows by exploiting (12)
$$Ky_{\stackrel{~}{s}}=K\left(\genfrac{}{}{0pt}{}{\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}{\varphi \stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}\right)=\left(\begin{array}{cc}\kappa & \widehat{\kappa }\\ \widehat{\kappa }& \kappa +\widehat{\kappa }\end{array}\right)\left(\genfrac{}{}{0pt}{}{\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}{\varphi \stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}\right)=\left(\begin{array}{cc}\stackrel{~}{K}(\varphi )& 0\\ 0& \stackrel{~}{K}(\varphi )\end{array}\right)\left(\genfrac{}{}{0pt}{}{\stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}{\varphi \stackrel{~}{y}_{\stackrel{~}{s}}(\varphi )}\right).$$
(18)
The second relation in (17) is obtained by the same argumentation with $`\varphi \varphi ^1`$. These facts will not only be crucial to formulate new types of ATFT, but also to explain patterns in well studied models. The knowledge of the distribution of the exponents with respect to the embedding is important as they grade the conserved charges. The classical masses are known to organise as components of the Perron-Frobenius vector $`y_1`$, such that (17) explains the aforementioned mass patterns (1) and (3).
### 2.2 Coxeter groups
Having related the root systems $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }`$, let us see next how to relate the corresponding Coxeter groups $`𝒲`$ and $`\stackrel{~}{𝒲}`$, which leave them invariant. We recall that Coxeter groups are generated by reflections on the hyperplane through the origin orthogonal to simple roots $`\alpha _i`$
$$\sigma _i(x)=x2\frac{x\alpha _i}{\alpha _i\alpha _i}\alpha _i\text{for }1i\mathrm{},x^{\mathrm{}}.$$
(19)
We can then think of the Coxeter group $`𝒲`$ as the set of all words in the generators $`\{\sigma _i\}`$ subject to the relations
$$(\sigma _i\sigma _j)^{m_{ij}}=1,1i,j\mathrm{},$$
(20)
where the $`m_{ij}`$ can take the values $`1,2,3,4,5`$ or $`6`$. Here we focus on the Coxeter groups for which the Cartan matrix is symmetric, in which case $`m_{ij}=\pi \mathrm{arccos}^1(K_{ij}/2)`$ are the integers labelling the edges of the Coxeter graph mentioned at the beginning of this section. $`\stackrel{~}{𝒲}`$ is constructed analogously when replacing $`K`$ by $`\stackrel{~}{K}(\varphi )`$. Note that when we use instead of $`\stackrel{~}{K}(\varphi )`$ the matrix $`\stackrel{~}{K}(\varphi ^1)`$, which naturally emerges through (14), this will lead to the same Coxeter relations.
From a group theoretical point of view we can identify $`\stackrel{~}{\sigma }_i\sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}}`$. With the help of the map $`\omega `$ we can relate the reflections $`\stackrel{~}{\sigma }_i`$ building up $`\stackrel{~}{𝒲}`$ to those constituting $`𝒲`$ as
$$\stackrel{~}{\sigma }_i\omega =\omega \sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}}\text{for }1i\stackrel{~}{\mathrm{}}.$$
(21)
This is seen easily by acting on a simple root in $`\mathrm{\Delta }`$
$`\omega \sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}}(\alpha _j)`$ $`=`$ $`\omega \left[\alpha _j(\alpha _j\alpha _{i+\stackrel{~}{\mathrm{}}})\alpha _{i+\stackrel{~}{\mathrm{}}}(\alpha _j\alpha _i)\alpha _i\right]`$ (22)
$`=`$ $`\omega (\alpha _j)\varphi (\alpha _j\alpha _{i+\stackrel{~}{\mathrm{}}})\stackrel{~}{\alpha }_i(\alpha _j\alpha _i)\stackrel{~}{\alpha }_i,`$ (23)
where in (22) we used twice (19) and the fact that in our labelling we always have $`\alpha _i\alpha _{i+\stackrel{~}{\mathrm{}}}=0`$. Then (23) simply follows upon using (5). On the other hand using (19) for $`\stackrel{~}{𝒲}`$ and (11) thereafter, we obtain
$`\stackrel{~}{\sigma }_i\omega (\alpha _j)`$ $`=`$ $`\omega (\alpha _j)(\omega (\alpha _j)\stackrel{~}{\alpha }_i)\stackrel{~}{\alpha }_i`$ (24)
$`=`$ $`\omega (\alpha _j)\varphi (\alpha _j\alpha _{i+\stackrel{~}{\mathrm{}}})\stackrel{~}{\alpha }_i(\alpha _j\alpha _i)\stackrel{~}{\alpha }_i.`$ (25)
In a similar identification has been made for the specific case of the embedding $`H_4E_8`$, which relies on the property of an inflation map which mimics the action of $`\omega `$ entirely inside $`\mathrm{\Delta }`$. Here we avoid the introduction of such a quantity.
Furthermore, for the second map $`\stackrel{~}{\omega }`$ we have the supplementary identity
$$\stackrel{~}{\omega }\stackrel{~}{\sigma }_i=\sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}}\stackrel{~}{\omega }\text{for }1i\stackrel{~}{\mathrm{}},$$
(26)
which follows from a similar argument as (21) upon using (19), (10) and (12). As we saw already in (11), we note here that $`\stackrel{~}{\omega }`$ plays the role of the “inverse” of $`\omega `$. There is no analogue to the inflation map in this case. More precisely we find
$$\omega \stackrel{~}{\omega }=(1+\varphi ^2)𝕀\text{and }\stackrel{~}{\omega }\omega =\left(\begin{array}{cc}𝕀\hfill & \varphi 𝕀\hfill \\ \varphi 𝕀\hfill & \varphi ^2𝕀\hfill \end{array}\right),$$
(27)
where $`𝕀`$ is the $`\stackrel{~}{\mathrm{}}\times \stackrel{~}{\mathrm{}}`$ unit matrix already encountered in (14).
For the application we have in mind, it is important to note that the entire root system $`\mathrm{\Delta }`$ can be separated into orbits $`\mathrm{\Omega }_i`$, such that $`\mathrm{\Delta }=_{i=1}^{\mathrm{}}\mathrm{\Omega }_i`$, each containing $`h`$ roots. Here $`h`$ is the order of the Coxeter element $`\sigma `$, i.e. the Coxeter number already introduced after (14), which is a product of $`\mathrm{}`$ simple reflections with the property $`\sigma ^h=1`$. As the reflections do not commute in general, such that Coxeter elements only form conjugacy classes, we have to specify our conventions. For this purpose we attach values $`c_i=\pm 1`$ to the vertices $`i`$ of the Coxeter graph, in such a way that no two vertices related to the same value are linked together. The vertices then separate into two disjoint sets $`V_\pm `$ and the Coxeter element $`\sigma =_{iV_{}}\sigma _i`$ $`_{iV_+}\sigma _i`$ is uniquely defined. Introducing “coloured” simple roots as $`\gamma _i=c_i\alpha _i`$, each orbit $`\mathrm{\Omega }_i(\stackrel{~}{\mathrm{\Omega }}_i)`$ is then generated by $`h(\stackrel{~}{h})`$ successive actions of $`\sigma (\stackrel{~}{\sigma })`$ on $`\gamma _i(\stackrel{~}{\gamma }_i)`$ . Since we know how to relate the simple reflections of $`𝒲`$ and $`\stackrel{~}{𝒲}`$ by means of (21) and (26), it is obvious that the Coxeter elements are intertwined as
$$\stackrel{~}{\sigma }\omega =\omega \sigma \text{and }\stackrel{~}{\omega }\stackrel{~}{\sigma }=\sigma \stackrel{~}{\omega }.$$
(28)
Note further that
$$\stackrel{~}{c}_i=c_i=c_{i+\stackrel{~}{\mathrm{}}}.$$
(29)
which is important for (28) to work, as it guarantees that (21) is not an obstacle for the above mentioned separation of the product into different colours. We then find that
$$\omega (\mathrm{\Omega }_i)=\{\begin{array}{cc}\stackrel{~}{\mathrm{\Omega }}_i\hfill & \text{for }1i\stackrel{~}{\mathrm{}}=\mathrm{}/2\hfill \\ \varphi \stackrel{~}{\mathrm{\Omega }}_i\hfill & \text{for }\stackrel{~}{\mathrm{}}<i\mathrm{}.\hfill \end{array}$$
(30)
Thus we can realize the map (4) orbit by orbit.
After this generic preliminaries let us discuss in detail the concrete examples of the embeddings $`H_2A_4`$, $`H_3D_6`$ and $`H_4E_8`$.
### 2.3 The embedding $`H_2A_4`$
We start with the most simple example, that is the embedding of $`H_2`$ (also referred to as $`I(5)`$, see e.g. ) into $`A_4`$. First of all we have to fix our conventions for naming the simple roots, which we do by means of the following Coxeter graph (where we adopt the common rule that the label $`m_{ij}=3`$ corresponds to one lace)
These conventions guarantee that we can realize the map $`\omega `$ as defined in (5), which is also indicated in the above diagrams. Accordingly, the Cartan matrix of $`A_4`$, as defined in general in (8), reads
$$K=\left(\begin{array}{cccc}\hfill 2& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 2& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1& \hfill 2\end{array}\right)=R\left(\begin{array}{cc}\hfill \stackrel{~}{K}& \hfill \varphi \stackrel{~}{K}\\ \hfill \varphi \stackrel{~}{K}& \hfill \varphi ^2\stackrel{~}{K}\end{array}\right),$$
(31)
where the Cartan matrix of $`H_2`$ is
$$\stackrel{~}{K}=\left(\begin{array}{cc}\hfill 2& \hfill \varphi \\ \hfill \varphi & \hfill 2\end{array}\right).$$
(32)
We may read off relation (12) directly by comparing (31) with (32). Furthermore, we note from the second relation in (31) that the identity (8) holds by reducing the higher powers of $`\varphi `$ as indicated above. The Coxeter numbers are $`h=\stackrel{~}{h}=5`$ and the set of exponents separate according to (15) into
$$\{1,2,3,4\}=\{1,4\}\{2,3\}.$$
(33)
Let us now see in detail how $`\omega `$ acts on the orbits $`\mathrm{\Omega }_i`$ and how the map (30) is realized. We choose $`\sigma =`$ $`\sigma _1\sigma _3\sigma _2\sigma _4`$ and $`\stackrel{~}{\sigma }=`$ $`\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2`$ for the Coxeter element of $`A_4`$ and $`H_2`$, respectively. The corresponding orbits $`\mathrm{\Omega }_i`$ and $`\stackrel{~}{\mathrm{\Omega }}_i`$ are computed by successive actions of $`\sigma `$ and $`\stackrel{~}{\sigma }`$, respectively, on the simple roots. One realizes that the map $`\omega `$ relates them indeed as specified in (30). Indicating in the first column the elements $`\sigma ^p`$ $`(\stackrel{~}{\sigma }^p)`$ for $`1ph=\stackrel{~}{h}`$ with which we act on the roots reported in the second row, we find for instance
| | $`\mathrm{\Omega }_1`$ | $`\omega (\mathrm{\Omega }_1)=\stackrel{~}{\mathrm{\Omega }}_1`$ | $`\mathrm{\Omega }_3`$ | $`\omega (\mathrm{\Omega }_3)=\varphi \stackrel{~}{\mathrm{\Omega }}_1`$ |
| --- | --- | --- | --- | --- |
| $`\sigma ^0,\stackrel{~}{\sigma }^0`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _3`$ | $`\varphi \stackrel{~}{\alpha }_1`$ |
| $`\sigma ^1,\stackrel{~}{\sigma }^1`$ | $`\alpha _3+\alpha _4`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _4`$ | $`\varphi ^2(\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)`$ |
| $`\sigma ^2,\stackrel{~}{\sigma }^2`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_2`$ |
| $`\sigma ^3,\stackrel{~}{\sigma }^3`$ | $`\alpha _2\alpha _3`$ | $`\varphi \stackrel{~}{\alpha }_1\stackrel{~}{\alpha }_2`$ | $`\alpha _1\alpha _3\alpha _4`$ | $`\varphi ^2\stackrel{~}{\alpha }_1\varphi \stackrel{~}{\alpha }_2`$ |
| $`\sigma ^4,\stackrel{~}{\sigma }^4`$ | $`\alpha _1\alpha _4`$ | $`\stackrel{~}{\alpha }_1\varphi \stackrel{~}{\alpha }_2`$ | $`\alpha _2\alpha _3\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_1\varphi ^2\stackrel{~}{\alpha }_2`$ |
| $`\sigma ^5,\stackrel{~}{\sigma }^5`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _3`$ | $`\varphi \stackrel{~}{\alpha }_1`$ |
In order to establish the last identification $`\omega (\mathrm{\Omega }_3)=\varphi \stackrel{~}{\mathrm{\Omega }}_1`$ we simply need to make use of relation (2). Furthermore, we obtain for the remaining orbits
| | $`\mathrm{\Omega }_2`$ | $`\omega (\mathrm{\Omega }_2)=\stackrel{~}{\mathrm{\Omega }}_2`$ | $`\mathrm{\Omega }_4`$ | $`\omega (\mathrm{\Omega }_4)=\varphi \stackrel{~}{\mathrm{\Omega }}_2`$ |
| --- | --- | --- | --- | --- |
| $`\sigma ^0,\stackrel{~}{\sigma }^0`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_2`$ |
| $`\sigma ^1,\stackrel{~}{\sigma }^1`$ | $`\alpha _2+\alpha _3`$ | $`\varphi \stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2`$ | $`\alpha _1+\alpha _3+\alpha _4`$ | $`\varphi ^2\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2`$ |
| $`\sigma ^2,\stackrel{~}{\sigma }^2`$ | $`\alpha _1+\alpha _4`$ | $`\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2`$ | $`\alpha _2+\alpha _3+\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_1+\varphi ^2\stackrel{~}{\alpha }_2`$ |
| $`\sigma ^3,\stackrel{~}{\sigma }^3`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _3`$ | $`\varphi \stackrel{~}{\alpha }_1`$ |
| $`\sigma ^4,\stackrel{~}{\sigma }^4`$ | $`\alpha _3\alpha _4`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)`$ | $`\alpha _1\alpha _2\alpha _3\alpha _4`$ | $`\varphi ^2(\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)`$ |
| $`\sigma ^5,\stackrel{~}{\sigma }^5`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_2`$ |
The identification $`\omega (\mathrm{\Omega }_4)=\varphi \stackrel{~}{\mathrm{\Omega }}_2`$ follows upon using (2). Having established how the root system $`\mathrm{\Delta }`$ can be mapped into $`\stackrel{~}{\mathrm{\Delta }}\varphi \stackrel{~}{\mathrm{\Delta }}`$ orbit by orbit, we want to see next how $`𝒲`$ relates to $`\stackrel{~}{𝒲}`$. In principle we have to check all $`|\stackrel{~}{𝒲}|`$ relations in (20). However, it suffices to establish the identification (21) for the generating relations of $`\stackrel{~}{𝒲}`$. It is known that $`H_2`$ can be generated entirely from $`\stackrel{~}{\sigma }_i^2=1`$ for $`i=1,2,3`$ together with
$$(\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2)^5=1,$$
(34)
which follows directly from the previous tables, since $`\sigma =\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2`$. In $`A_4`$ this corresponds to
$$(\sigma _1\sigma _3\sigma _2\sigma _4)^5=1.$$
(35)
The remaining relations are trivially satisfied. We know that by definition $`\sigma _i^2=1`$ and therefore this also holds when squaring the product of the embedding on the right hand side of (21) $`(\sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}})^2=\sigma _i^2\sigma _{i+\stackrel{~}{\mathrm{}}}^2=1`$. We used here the last equality in (29), such that $`\sigma _i`$ and $`\sigma _{i+\stackrel{~}{\mathrm{}}}`$ commute.
### 2.4 The embedding $`H_3D_6`$
In this case we fix our conventions as
The Cartan matrix (8) of $`D_6`$ then reads
$$K=\left(\begin{array}{cccccc}\hfill 2& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 2& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 2& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1& \hfill 2& \hfill 1\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 2\end{array}\right)=R\left(\begin{array}{cc}\hfill \stackrel{~}{K}& \hfill \varphi \stackrel{~}{K}\\ \hfill \varphi \stackrel{~}{K}& \hfill \varphi ^2\stackrel{~}{K}\end{array}\right),$$
(36)
where we also exhibit relation (8). Noting further that the Cartan matrix of $`H_3`$ is
$$\stackrel{~}{K}=\left(\begin{array}{ccc}\hfill 2& \hfill 1& \hfill 0\\ \hfill 1& \hfill 2& \hfill \varphi \\ \hfill 0& \hfill \varphi & \hfill 2\end{array}\right),$$
(37)
the relation (12) is read off directly by comparing (36) with (37). The Coxeter numbers are $`h=\stackrel{~}{h}=10`$ and the set of exponents separate according to (15) into
$$\{1,3,5,5,7,9\}=\{1,5,9\}\{3,5,7\}.$$
(38)
Let us now see how $`\omega `$ acts on the orbits $`\mathrm{\Omega }_i`$ and how the map (30) is realized. We choose $`\sigma =`$ $`\sigma _1\sigma _4\sigma _3\sigma _6\sigma _2\sigma _5`$ and $`\stackrel{~}{\sigma }=`$ $`\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_2`$ for the Coxeter element of $`D_6`$ and $`H_3`$, respectively. The corresponding orbits $`\mathrm{\Omega }_i`$ and $`\stackrel{~}{\mathrm{\Omega }}_i`$ are computed by successive actions of $`\sigma `$ and $`\stackrel{~}{\sigma }`$, respectively, on the simple roots. One realizes that the map $`\omega `$ relates them indeed as specified in (30). See appendix A for the explicit computation of the orbits.
Once more we may check how the root system $`\mathrm{\Delta }`$ can be mapped into $`\stackrel{~}{\mathrm{\Delta }}\varphi \stackrel{~}{\mathrm{\Delta }}`$ orbit by orbit. First we relate $`𝒲`$ to $`\stackrel{~}{𝒲}`$ and verify (21) for the generating relations of $`\stackrel{~}{𝒲}`$, which for $`H_3`$ are known to be $`\stackrel{~}{\sigma }_i^2=\stackrel{~}{\sigma }_{i+\stackrel{~}{\mathrm{}}}^2=1`$ for $`i=1,2,3`$ together with
$$(\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_3)^5=(\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2)^3=(\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_3)^2=1.$$
(39)
We verify that (39) corresponds to
$$(\sigma _2\sigma _5\sigma _3\sigma _6)^5=(\sigma _1\sigma _4\sigma _2\sigma _5)^3=(\sigma _1\sigma _4\sigma _3\sigma _6)^2=1.$$
(40)
By the same reasoning as in the $`H_2`$-case it also follows that $`(\sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}})^2=1`$.
### 2.5 The embedding $`H_4E_8`$
Also in this case we first fix our conventions for naming the simple roots in order to guarantee that we can realize the map $`\omega `$ as defined in (5). Here we label according to the Coxeter graph
The corresponding Cartan matrix of $`E_8`$ together with its construction from inner products in $`\stackrel{~}{\mathrm{\Delta }}`$ in agreement with (8) is
$$K=\left(\begin{array}{cccccccc}\hfill 2& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 2& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 2& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 2& \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 2& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 2& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 2\end{array}\right)=R\left(\begin{array}{cc}\hfill \stackrel{~}{K}& \hfill \varphi \stackrel{~}{K}\\ \hfill \varphi \stackrel{~}{K}& \hfill \varphi ^2\stackrel{~}{K}\end{array}\right).$$
(41)
The Cartan matrix of $`H_4`$ reads
$$\stackrel{~}{K}=\left(\begin{array}{cccc}\hfill 2& \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 2& \hfill \varphi \\ \hfill 0& \hfill 0& \hfill \varphi & \hfill 2\end{array}\right).$$
(42)
Now the Coxeter numbers are $`h=\stackrel{~}{h}=30`$ and the set of exponents separate according to (15) into
$$\{1,7,11,13,17,19,23,29\}=\{1,11,19,29\}\{7,13,17,23\}.$$
(43)
In order to see how $`\omega `$ acts on the orbits $`\mathrm{\Omega }_i`$ and how the map (30) is realized, we choose $`\sigma =`$ $`\sigma _1\sigma _5\sigma _3\sigma _7\sigma _2\sigma _6\sigma _4\sigma _8`$ and $`\stackrel{~}{\sigma }=`$ $`\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_4`$ for the Coxeter element of $`E_8`$ and $`H_4`$, respectively. We may then compute the corresponding orbits $`\mathrm{\Omega }_i`$, $`\stackrel{~}{\mathrm{\Omega }}_i`$ by successive actions of $`\sigma `$, $`\stackrel{~}{\sigma }`$ and the simple roots and realize that the map $`\omega `$ relates them indeed as specified in (30). See appendix for the orbits. The individual reflections are related as (21) and the generating relations for the Coxeter group are
$$\left(\stackrel{~}{\sigma }_6\stackrel{~}{\sigma }_7\right)^5=(\stackrel{~}{\sigma }_4\stackrel{~}{\sigma }_5)^3=(\stackrel{~}{\sigma }_4\stackrel{~}{\sigma }_6)^3=1.$$
(44)
which correspond to
$$(\sigma _2\sigma _6\sigma _3\sigma _7)^5=(\sigma _4\sigma _8\sigma _1\sigma _5)^3=(\sigma _4\sigma _8\sigma _2\sigma _6)^3=1.$$
(45)
By the same reasoning as in the $`H_2`$-case it also follows that $`(\sigma _i\sigma _{i+\stackrel{~}{\mathrm{}}})^2=1`$ .
## 3 Affine Toda field theories
### 3.1 Generalities
In this section we will demonstrate how one may construct an affine Toda field theory related to a non-crystallographic Coxeter group $`\stackrel{~}{𝒲}`$ from a theory related to a crystallographic Coxeter group $`𝒲`$ by means of the discussed embedding. We start by taking $`G`$ to be a Lie group with $`HG`$ a maximal Torus and $`𝐡`$ its Cartan subalgebra. Then the affine Toda field theory Lagrangian can be expressed as
$$=\frac{1}{\beta ^2}\mathrm{T}r\left(\frac{1}{2}_\mu g^1^\mu gm^2gEg^1E^{}\right),$$
(46)
where $`g=\mathrm{exp}(\beta \mathrm{\Phi })H`$, $`\beta `$ is a coupling constant and $`m`$a mass scale. The regular element $`E=\eta h^{}`$ with conjugate $`E^{}`$ can be expanded in the Cartan subalgebra in apposition $`h_i^{}H^{}`$ . We normalize the trace according to the Cartan-Weyl basis, that is $`\mathrm{T}r\left(E_\alpha E_\alpha \right)=1`$. Notice that one can not write down a Lagrangian of the type (46) and relate it directly to a non-crystallographic Coxeter group, since there is no proper Lie group and Lie algebra associated to them.
Conventionally one introduces $`\mathrm{}`$ scalar fields $`\varphi _i`$ by expanding the field $`\mathrm{\Phi }`$ in the Cartan subalgebra $`\mathrm{\Phi }=_{i=1}^{\mathrm{}}\varphi _iH_i`$ with Cartan Weyl generators $`H_i𝐡`$ . Developing now the Lagrangian (46) in powers of the coupling constant, it follows that the term of zeroth order in $`\beta `$, i.e. the quadratic term in $`\varphi _i`$ becomes
$$\frac{1}{2}\underset{i=1}{\overset{\mathrm{}}{}}m_i^2|\varphi _i|^2=\frac{1}{2}m^2\underset{i=1}{\overset{\mathrm{}}{}}|\eta \alpha _i|^2|\varphi _i|^2,$$
(47)
such that the mass of particle $`i`$ can be identified as $`m_i=m|\eta \alpha _i|`$. Here $`\eta `$ is the eigenvector of the Coxeter element with eigenvalue $`\mathrm{exp}(2\pi i/h)`$. Proceeding to the first order term in $`\beta `$, the constant in front of the cubic terms in the fields divided by the symmetry factor $`3!`$ is taken to define the three-point coupling. It is computed to
$$C_{ijk}=\frac{4\beta m^2}{\sqrt{h}}\epsilon _{ijq}\mathrm{\Delta }_{ijk}$$
(48)
where $`\mathrm{\Delta }_{ijk}=\sqrt{s(sm_i)(sm_j)(sm_k)}`$ with $`s=(m_i+m_j+m_k)/2`$ is Heron’s formula for the area of the triangle formed by the masses of the three fusing particles $`i,j,k`$. The structure constants $`\epsilon _{ijq}`$ result from the Lie algebraic commutator $`[E_{\gamma _i},E_{\sigma ^q\gamma _j}]=\epsilon _{ijq}E_{\gamma _i+\sigma ^q\gamma _j}=\epsilon _{ijq}E_{\sigma ^{\stackrel{~}{q}}\gamma _{\overline{k}}}`$ and for simply laced Lie algebras are normalized to $`\epsilon _{ijq}=\pm 1,0`$. The $`\gamma _i=c_i\alpha _i`$ are the coloured simple roots introduced in section 2.2. In other words the three-point coupling $`C_{ijk}`$ is non-vanishing if and only if the commutator of step operators related to roots in the orbits $`\mathrm{\Omega }_i`$ and $`\mathrm{\Omega }_j`$ is non-zero. Through this reasoning the fusing rule
$$C_{ijk}0\gamma _i+\sigma ^q\gamma _j+\sigma ^p\gamma _k=0.$$
(49)
can be related to the ATFT-Lagrangian.
Having computed these classical quantities one may ask the question towards a formulation of the corresponding quantum field theory. A first glimpse of its nature is obtained from a perturbative computation of the mass corrections. For ATFT the logarithmic divergencies of the self-energy corrections can be removed simply by normal ordering . Then the one loop corrections to the mass of the particle $`k`$ are just the finite contributions resulting from a bubble graph, which were found to be
$$\delta m_k^2=i\frac{_{ij}(m_k^2)}{(2\pi )^2},$$
(50)
where the sum extends over all intermediate contributions
$$_{ij}(m_k^2)=i\pi \frac{C_{ijk}^2(\pi \theta _{ij}^k)}{m_im_j\mathrm{sin}\theta _{ij}^k},$$
(51)
with $`\theta _{ij}^k`$ being the fusing angle for the process $`i+j\overline{k}`$.
One should note that instead of the Lie algebraic formulation for the affine Toda field theory Lagrangian (46), one can in principle also start with a Lagrangian for which the trace in (46) is already computed in the adjoint representation, such that it involves only roots rather than Lie algebraic quantities
$$=\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }\frac{m^2}{\beta ^2}\underset{i=0}{\overset{\mathrm{}}{}}n_i\mathrm{exp}(\beta \alpha _i\mathrm{\Phi }).$$
(52)
Here $`\alpha _0=`$ $`_{i=0}^{\mathrm{}}n_i\alpha _i`$ is the negative of the highest root and $`n_i`$ the Kac labels. In this case the expansion of the Lagrangian in $`\beta `$ yields
$$\left(M^2\right)_{ij}=m^2\underset{p=0}{\overset{\mathrm{}}{}}n_p\alpha _p^i\alpha _p^j\text{and }C_{ijk}=\beta m^2\underset{p=0}{\overset{\mathrm{}}{}}n_p\alpha _p^i\alpha _p^j\alpha _p^k.$$
(53)
In most cases the two formulations are equivalent, however one is often more advantageous than the other. For instance, (52) does not allow for a generic case independent treatment as it relies on a special choice of the basis for the roots $`\alpha _i`$ needed to ensure that the mass matrix becomes diagonal. In addition (46) yields an explanation for the fact that the masses of ATFT organize into the Perron-Frobenius vector of the Cartan matrix and a generic derivation of the fusing rule (49). For our present purposes it is important to note that (46) relies on the existence of a Lie algebra, whereas (52) only requires a root system. Thus in principle we could write down a Lagrangian of the type (52) for non-crystallographic Coxeter groups, but the formulation of an equivalent Lie algebraic version would be impossible. In addition, classical integrability for ATFT is established by means of a Lax pair formulation in terms of Lie algebraic quantities . Therefore it would remain an open issue whether for non-crystallographic Coxeter groups (52) corresponds to integrable models or not. Here our formulation will not be a Lagrangian of the form (52) involving a non-crystallographic root system, but rather a reduced system which exploits the previously discussed embedding structure with regard to the Cartan subalgebra in apposition. The crystallographic Lie algebra structure is preserved in the reduction procedure, from which we can see immediately that the new theories also possess Lax pairs ensuring the classical integrability.
### 3.2 Reduction from crystallographic to non-crystallographic ATFT
Somewhat analogue to the folding procedure , we may now alter the above expansion of $`\mathrm{\Phi }`$ in order to define a new theory, which in this case only contains $`\stackrel{~}{\mathrm{}}=\mathrm{}/2`$ scalar fields $`\stackrel{~}{\phi }_i`$. This is achieved by expanding
$$\mathrm{\Phi }=\underset{i=1}{\overset{\mathrm{}}{}}\mu (\phi _i)H_i,\text{with }\mu (\phi _i)=\frac{1}{\varphi \sqrt{3}}\{\begin{array}{cc}\stackrel{~}{\phi }_i\hfill & \text{for }1i\stackrel{~}{\mathrm{}}=\mathrm{}/2\hfill \\ \varphi \stackrel{~}{\phi }_{i\stackrel{~}{\mathrm{}}}\hfill & \text{for }\stackrel{~}{\mathrm{}}<i\mathrm{}.\hfill \end{array}$$
(54)
Here the map $`\mu `$ is inspired by the previously introduced map $`\stackrel{~}{\omega }`$ (10). At this point one could have also defined $`\mu `$ in such a way that it multiplies the fields $`\phi _i`$ for $`\stackrel{~}{\mathrm{}}<i\mathrm{}`$ by $`\varphi ^1`$. However, the choice (54) is distinct by the structure emerging below for the three-point coupling. Other possibilities are also conceivable. Having an alternative expansion for the fields $`\mathrm{\Phi }`$ as in (54), we re-consider first the quadratic term (47) in the expansion of the Lagrangian, which now becomes
$$\frac{1}{2}\underset{i=1}{\overset{\mathrm{}}{}}m_i^2|\mu (\stackrel{~}{\phi }_i)|^2=\frac{1}{2}\underset{i=1}{\overset{\stackrel{~}{\mathrm{}}}{}}\frac{1}{3}(\varphi ^2m_i^2+m_{i+\stackrel{~}{\mathrm{}}}^2)|(\stackrel{~}{\phi }_i)|^2.$$
(55)
We take this identity as the defining relation for the masses of the $`\stackrel{~}{\mathrm{}}`$ new scalar fields $`\stackrel{~}{\varphi }_i`$
$$\stackrel{~}{m}_i^2=\frac{1}{3}(\varphi ^2m_i^2+m_{i+\stackrel{~}{\mathrm{}}}^2)=m_i^2.$$
(56)
In (56) we made use of the fact that the masses $`m_i`$ can be identified as the components of the Perron-Frobenius eigenvector of $`K`$, that is $`m_i=(y_1)_i`$. Subsequently we employ (17) which implies that $`m_i=\varphi m_{i+\stackrel{~}{\mathrm{}}}`$ and then $`\varphi ^2+\varphi ^2=3`$. We also notice the fact that the masses $`\stackrel{~}{m}_i`$ admit a “genuine” interpretation as belonging to an affine Toda field theory related to a non-crystallographic Coxeter group, since (17) also implies that the masses $`\stackrel{~}{m}_i`$ are the components of the Perron-Frobenius eigenvector of $`\stackrel{~}{K}(\varphi )`$.
Similarly we proceed further and read off the next order term in $`\beta `$ to define a new three-point coupling $`\stackrel{~}{C}_{ijk}`$. We compute
$`\stackrel{~}{C}_{ijk}`$ $`=`$ $`\varphi ^3C_{(i+\mathrm{})(j+\stackrel{~}{\mathrm{}})(k+\stackrel{~}{\mathrm{}})}+\varphi ^2\left(C_{i(j+\stackrel{~}{\mathrm{}})(k+\stackrel{~}{\mathrm{}})}+C_{(i+\stackrel{~}{\mathrm{}})j(k+\stackrel{~}{\mathrm{}})}+C_{(i+\stackrel{~}{\mathrm{}})(j+\stackrel{~}{\mathrm{}})k}\right)`$ (57)
$`+\varphi \left(C_{(i+\stackrel{~}{\mathrm{}})jk}+C_{i(j+\stackrel{~}{\mathrm{}})k}+C_{ij(k+\stackrel{~}{\mathrm{}})}\right)+C_{ijk},`$
for $`1i,j,k,\stackrel{~}{\mathrm{}}`$. The identification of the Coxeter element (21) translates the fusing rule (49) into
$`\stackrel{~}{C}_{ijk}0`$ $``$ $`\omega \gamma _i+\stackrel{~}{\sigma }^q\omega \gamma _j+\stackrel{~}{\sigma }^p\omega \gamma _k=0,`$ (58)
$``$ $`\varphi ^{t_i}\stackrel{~}{\gamma }_i+\varphi ^{t_j}\stackrel{~}{\sigma }^q\stackrel{~}{\gamma }_j+\varphi ^{t_k}\stackrel{~}{\sigma }^p\stackrel{~}{\gamma }_k=0,`$ (59)
where the $`t_i`$ are the integers introduced in equation (8). The relation (59) is entirely expressed in quantities belonging to the non-crystallographic Coxeter group. In comparison with the fusing rule related to ATFT for simply laced Lie algebras (49), it had to be “$`\varphi `$-deformed”, somewhat similar to the q-deformed versions of the fusing rule needed for ATFT associated with non-simply laced Lie algebras . Note that besides the two non-equivalent solutions related as $`qh+1/2(c_ic_j)q`$ and $`ph+1/2(c_ic_j)p`$ , in (59) there could be more solutions corresponding to different values of the integers $`t_i,t_j,t_k`$.
Alternatively, we may compute the masses and three-point couplings from the Lagrangian (52). Using the expansion (54) in (52) yields the same quantities. We obtain now
$$\left(\stackrel{~}{M}^2\right)_{ij}=m^2\underset{p=0}{\overset{\mathrm{}}{}}n_p\widehat{\alpha }_p^i\widehat{\alpha }_p^j\text{and }\stackrel{~}{C}_{ijk}=\beta m^2\underset{p=0}{\overset{\mathrm{}}{}}n_p\widehat{\alpha }_p^i\widehat{\alpha }_p^j\widehat{\alpha }_p^k.$$
(60)
where we have folded the $`\mathrm{}`$-components of the root $`\alpha _p^i`$ with $`1i\mathrm{}`$ into an $`\stackrel{~}{\mathrm{}}`$-component vector
$$\widehat{\alpha }_p^i=\frac{1}{3}(\varphi ^1\alpha _p^i+\alpha _p^{i+\stackrel{~}{\mathrm{}}}),$$
(61)
for $`1i\stackrel{~}{\mathrm{}}`$. We stress that this is not the same as writing down (52) for a non-crystallographic root system, since it still involves $`\mathrm{}`$ roots, albeit now represented in $`^\stackrel{~}{\mathrm{}}`$. Having computed the three-point couplings for the reduced theory, we may compute the mass renormalisation from (51) as
$$_{ij}(\stackrel{~}{m}_k^2)=i\pi \frac{\stackrel{~}{C}_{ijk}^2(\pi \theta _{ij}^k)}{\stackrel{~}{m}_i\stackrel{~}{m}_j\mathrm{sin}\theta _{ij}^k}.$$
(62)
This will shed light on the possible form of the scattering matrix for ATFT related to non-crystallographic Coxeter groups.
### 3.3 From $`A_4`$ to $`H_2`$-affine Toda field theory
Let us now make the above general formulae more explicit. When ignoring the overall mass scale, the masses of $`A_4`$-ATFT can be brought into the simple form
$$m_1=m_2=1\text{and }m_3=m_4=\varphi .$$
(63)
Keeping the same normalization for the overall mass scale, the identity (56) yields for the classical masses of the $`H_2`$-ATFT
$$\stackrel{~}{m}_1=m_1\text{and }\stackrel{~}{m}_2=m_2.$$
(64)
According to (48) and (57) we then compute the three point couplings $`C_{ijk}`$ and $`\stackrel{~}{C}_{ijk}`$ together with their corresponding fusing rules, which result from the tables provided in section 2.3
$$\begin{array}{ccc}& & \\ C_{113}=\frac{4i}{\sqrt{5}}\mathrm{\Delta }_{113}\hfill & \gamma _1+\sigma \gamma _1+\sigma ^3\gamma _3=0\hfill & \stackrel{~}{C}_{111}=3\varphi C_{113}=\frac{4i}{\sqrt{5}}\sqrt{9+12\varphi }\stackrel{~}{\mathrm{\Delta }}_{111}\hfill \\ & & \\ C_{144}=\varphi C_{113}\hfill & \gamma _1+\sigma ^2\gamma _4+\sigma ^4\gamma _4=0\hfill & \stackrel{~}{C}_{122}=\varphi ^3C_{113}\hfill \\ & & \\ C_{224}=C_{113}\hfill & \gamma _2+\sigma \gamma _2+\sigma ^3\gamma _4=0\hfill & \stackrel{~}{C}_{222}=\stackrel{~}{C}_{111}\hfill \\ & & \\ C_{233}=\varphi C_{113}\hfill & \gamma _2+\sigma \gamma _3+\sigma ^3\gamma _3=0\hfill & \stackrel{~}{C}_{211}=\stackrel{~}{C}_{122}.\hfill \end{array}$$
(65)
We did not report the factor $`\beta m^2`$ in each of the couplings. Here $`\stackrel{~}{\mathrm{\Delta }}_{ijk}`$ is the area of the triangle bounded by the masses $`\stackrel{~}{m}_i,\stackrel{~}{m}_j,\stackrel{~}{m}_k`$. Note that now the factor of proportionality between $`|\stackrel{~}{C}_{ijk}|`$ and $`\stackrel{~}{\mathrm{\Delta }}_{ijk}`$ is no longer universal as in (48). This is sufficient information to carry out the perturbation theory up to order $`\beta ^2`$. The mass corrections to $`\stackrel{~}{m}_1`$ according to (62) are computed by summing the convergent contributions to the one-loop corrections, which up to the symmetry factors are
We omitted the usual arrows indicating the time direction and assume throughout this paper that they all run to the right hand side. We computed the symmetry factors by applying Wick’s theorem and include them into formula (62), such that
$$_{ij}(\stackrel{~}{m}_1^2)=i\pi \frac{\frac{\pi }{3}}{\mathrm{sin}\frac{2\pi }{3}}\left(18\stackrel{~}{C}_{111}^2+36\stackrel{~}{C}_{112}^2+18\stackrel{~}{C}_{221}^2\right).$$
(66)
where we assume that the particles are conjugate to each other, that is $`\overline{1}=2`$, $`\overline{2}=1`$ where the bar indicates the anti-particle. Similarly, the mass corrections to $`\stackrel{~}{m}_2`$ result from summing
which gives
$$_{ij}(\stackrel{~}{m}_2^2)=i\pi \frac{\frac{\pi }{3}}{\mathrm{sin}\frac{2\pi }{3}}\left(18\stackrel{~}{C}_{222}^2+36\stackrel{~}{C}_{221}^2+18\stackrel{~}{C}_{112}^2\right).$$
(67)
Assembling this according to (50) yields the important fact that the classical mass ratios are conserved in the quantum field theory
$$\frac{\delta m_1^2}{m_1^2}=\frac{\delta m_2^2}{m_2^2}=\frac{\sqrt{3}}{4}(5+3\varphi )\stackrel{~}{C}_{111}^2.$$
(68)
This means at first order perturbation theory the masses renormalise equally. As the masses of both particles coincide and they undergo the same fusing processes, they appear to be indistinguishable at this stage. This possibly hints towards a non-diagonal scattering theory, that means a theory in which backscattering is possible. However, it remains to be seen whether there exist higher charges in this theory which make the particles distinct.
### 3.4 From $`D_6`$ to $`H_3`$-affine Toda field theory
The masses of $`D_6`$-ATFT are known for a long time and can be brought into the form (3). Keeping the same normalization for the overall mass scale, the identity (56) yields for the classical masses of the $`H_3`$-ATFT
$$\stackrel{~}{m}_1=\varphi ^1,\stackrel{~}{m}_2=\sqrt{1+\varphi ^2},\stackrel{~}{m}_3=1.$$
(69)
According to (48) and (57) we then compute the three point couplings $`C_{ijk}`$ and $`\stackrel{~}{C}_{ijk}`$ together with their corresponding fusing rules, which result from the tables provided in the appendix
$$\begin{array}{ccc}& & \\ \begin{array}{c}C_{112}=4i/\sqrt{10}\mathrm{\Delta }_{112}\hfill \\ C_{442}=\varphi ^3C_{112}\hfill \\ C_{445}=\varphi ^2C_{112}\hfill \end{array}\hfill & \begin{array}{c}\gamma _1+\sigma \gamma _1+\sigma ^6\gamma _2=0\hfill \\ \gamma _4+\sigma ^3\gamma _4+\sigma ^7\gamma _2=0\hfill \\ \gamma _4+\sigma \gamma _4+\sigma ^6\gamma _5=0\hfill \end{array}\hfill & \stackrel{~}{C}_{112}=(10\varphi +7)\stackrel{~}{\mathrm{\Delta }}_{112}\hfill \\ & & \\ \begin{array}{c}C_{665}=\varphi ^5C_{112}\hfill \\ C_{332}=\varphi ^3C_{112}\hfill \\ C_{335}=\varphi ^2C_{112}\hfill \end{array}\hfill & \begin{array}{c}\gamma _6+\sigma ^3\gamma _6+\sigma ^7\gamma _5=0\hfill \\ \gamma _3+\sigma ^3\gamma _3+\sigma ^7\gamma _2=0\hfill \\ \gamma _3+\sigma \gamma _3+\sigma ^6\gamma _5=0\hfill \end{array}\hfill & \stackrel{~}{C}_{332}=(1+5\varphi )\stackrel{~}{\mathrm{\Delta }}_{332}\hfill \\ & & \\ \begin{array}{c}C_{126}=\varphi ^2C_{112}\hfill \\ C_{156}=\varphi ^3C_{112}\hfill \end{array}\hfill & \begin{array}{c}\gamma _1+\sigma ^2\gamma _2+\sigma ^6\gamma _6=0\hfill \\ \gamma _1+\sigma ^4\gamma _5+\sigma ^8\gamma _6=0\hfill \end{array}\hfill & \stackrel{~}{C}_{123}=5.56758\stackrel{~}{\mathrm{\Delta }}_{123}\hfill \\ & & \\ \begin{array}{c}C_{225}=(\varphi ^41)C_{112}\hfill \\ C_{255}=(\varphi \varphi ^5)C_{112}\hfill \end{array}\hfill & \begin{array}{c}\gamma _2+\sigma ^2\gamma _2+\sigma ^6\gamma _5=0\hfill \\ \gamma _2+\sigma ^3\gamma _5+\sigma ^7\gamma _5=0\hfill \end{array}\hfill & \stackrel{~}{C}_{222}=8.6253\stackrel{~}{\mathrm{\Delta }}_{222}\hfill \\ & & \\ C_{134}=\varphi ^2C_{112}\hfill & \gamma _1+\sigma ^3\gamma _3+\sigma ^7\gamma _4=0\hfill & \stackrel{~}{C}_{113}=2\varphi ^2\stackrel{~}{\mathrm{\Delta }}_{113}\hfill \\ & & \\ C_{346}=\varphi ^3C_{112}\hfill & \gamma _3+\sigma ^8\gamma _4+\sigma ^4\gamma _6=0\hfill & \stackrel{~}{C}_{331}=2\varphi ^3\stackrel{~}{\mathrm{\Delta }}_{331}.\hfill \end{array}$$
(70)
Here we did not report the factor of $`i\beta m^24/\sqrt{10}`$ in $`\stackrel{~}{C}`$. The fusing rules reduce according to (58), for instance
$`\gamma _1+\sigma \gamma _1+\sigma ^6\gamma _2`$ $`=`$ $`0\stackrel{~}{\gamma }_1+\stackrel{~}{\sigma }\stackrel{~}{\gamma }_1+\stackrel{~}{\sigma }^6\stackrel{~}{\gamma }_2=0,`$ (71)
$`\gamma _4+\sigma ^3\gamma _4+\sigma ^7\gamma _2`$ $`=`$ $`0\varphi \stackrel{~}{\gamma }_1+\varphi \stackrel{~}{\sigma }^3\stackrel{~}{\gamma }_1+\stackrel{~}{\sigma }^7\stackrel{~}{\gamma }_2=0,`$ (72)
$`\gamma _4+\sigma \gamma _4+\sigma ^6\gamma _5`$ $`=`$ $`0\varphi \stackrel{~}{\gamma }_1+\varphi \stackrel{~}{\sigma }\stackrel{~}{\gamma }_1+\varphi \sigma ^6\gamma _2=0.`$ (73)
Note that we can construct the solution (73) trivially from (71) simply by multiplying it with $`\varphi `$. However, (72) can not be obtained from (71) or (73) in such a manner and has to be regarded as independent.
As described in the previous section, we compute the mass renormalisation to
$`{\displaystyle _{ij}}(\stackrel{~}{m}_1^2)`$ $`=`$ $`i\pi \left({\displaystyle \frac{36\stackrel{~}{C}_{112}^2\frac{\pi }{10}}{\stackrel{~}{m}_1\stackrel{~}{m}_2\mathrm{sin}\frac{9\pi }{10}}}+{\displaystyle \frac{36\stackrel{~}{C}_{123}^2\widehat{\varphi }}{\stackrel{~}{m}_2\stackrel{~}{m}_3\mathrm{sin}(\pi \widehat{\varphi })}}+{\displaystyle \frac{18\stackrel{~}{C}_{133}^2\frac{2\pi }{10}}{\stackrel{~}{m}_3\stackrel{~}{m}_3\mathrm{sin}\frac{8\pi }{10}}}+{\displaystyle \frac{36\stackrel{~}{C}_{113}^2\frac{2\pi }{10}}{\stackrel{~}{m}_1\stackrel{~}{m}_3\mathrm{sin}\frac{8\pi }{10}}}\right)`$ (74)
$`{\displaystyle _{ij}}(\stackrel{~}{m}_2^2)`$ $`=`$ $`i\pi \left({\displaystyle \frac{18\stackrel{~}{C}_{222}^2\frac{\pi }{3}}{\stackrel{~}{m}_2\stackrel{~}{m}_2\mathrm{sin}\frac{2\pi }{3}}}+{\displaystyle \frac{36\stackrel{~}{C}_{123}^2\frac{\pi }{2}}{\stackrel{~}{m}_1\stackrel{~}{m}_3\mathrm{sin}\frac{\pi }{2}}}+{\displaystyle \frac{18\stackrel{~}{C}_{112}^2\frac{8\pi }{10}}{\stackrel{~}{m}_1\stackrel{~}{m}_1\mathrm{sin}\frac{2\pi }{10}}}+{\displaystyle \frac{18\stackrel{~}{C}_{233}^2\frac{4\pi }{10}}{\stackrel{~}{m}_3\stackrel{~}{m}_3\mathrm{sin}\frac{6\pi }{10}}}\right)`$ (75)
$`{\displaystyle _{ij}}(\stackrel{~}{m}_3^2)`$ $`=`$ $`i\pi \left({\displaystyle \frac{36\stackrel{~}{C}_{233}^2\frac{3\pi }{10}}{\stackrel{~}{m}_2\stackrel{~}{m}_3\mathrm{sin}\frac{7\pi }{10}}}+{\displaystyle \frac{36\stackrel{~}{C}_{123}^2\stackrel{ˇ}{\varphi }}{\stackrel{~}{m}_2\stackrel{~}{m}_3\mathrm{sin}(\pi \stackrel{ˇ}{\varphi })}}+{\displaystyle \frac{18\stackrel{~}{C}_{113}^2\frac{4\pi }{10}}{\stackrel{~}{m}_1\stackrel{~}{m}_1\mathrm{sin}\frac{6\pi }{10}}}+{\displaystyle \frac{36\stackrel{~}{C}_{133}^2\frac{4\pi }{10}}{\stackrel{~}{m}_1\stackrel{~}{m}_3\mathrm{sin}\frac{6\pi }{10}}}\right)`$ (76)
where we abbreviate $`\widehat{\varphi }=\mathrm{arctan}\varphi ^1`$, $`\stackrel{ˇ}{\varphi }=\mathrm{arctan}\varphi `$. From this it follows then that classical mass ratios are not conserved in the quantum field theory
$$\frac{\delta m_1^2}{m_1^2}=196.996\mathrm{}\frac{\delta m_2^2}{m_2^2}=647.392\mathrm{}\text{and }\frac{\delta m_3^2}{m_3^2}=924.343\mathrm{}$$
(77)
This means that the scattering matrix for $`H_3`$-ATFT can not be of the simple form as for ATFT related to simply laced Lie algebras.
### 3.5 From $`E_8`$ to $`H_4`$ affine Toda field theory
The masses of $`E_8`$-ATFT in the form (1) indicate the underlying $`H_4`$ structure and (56) yields for the classical masses for the $`H_4`$-ATFT
$$\stackrel{~}{m}_1=m_1,\stackrel{~}{m}_2=m_2,\stackrel{~}{m}_3=m_3,\stackrel{~}{m}_4=m_4.$$
(78)
Similarly as in the previous section we compute the three point couplings $`C_{ijk}`$ and $`\stackrel{~}{C}_{ijk}`$ from (48) and (57) together with their corresponding fusing rules
| $`\begin{array}{c}C_{111}=\frac{4i}{\sqrt{30}}0.433013\\ C_{511}=\frac{4i}{\sqrt{30}}0.475528\\ C_{551}=\varphi C_{511}\\ C_{555}=\varphi ^2C_{111}\end{array}`$ | $`\begin{array}{c}\gamma _1+\sigma ^{10}\gamma _1+\sigma ^{20}\gamma _1=0\\ \gamma _5+\sigma ^{12}\gamma _1+\sigma ^{18}\gamma _1=0\\ \gamma _5+\sigma ^{12}\gamma _5+\sigma ^{21}\gamma _1=0\\ \gamma _5+\sigma ^{10}\gamma _5+\sigma ^{20}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{111}=31.3768\stackrel{~}{\mathrm{\Delta }}_{111}`$ |
| --- | --- | --- |
| $`\begin{array}{c}C_{211}=\frac{4i}{\sqrt{30}}0.103956\\ C_{521}=\varphi ^2C_{321}\\ C_{655}=\varphi ^2C_{211}\end{array}`$ | $`\begin{array}{c}\gamma _2+\sigma ^{14}\gamma _1+\sigma ^{15}\gamma _1=0\\ \gamma _5+\sigma ^{13}\gamma _2+\sigma ^{23}\gamma _1=0\\ \gamma _6+\sigma ^{14}\gamma _5+\sigma ^{15}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{211}=37.1363\stackrel{~}{\mathrm{\Delta }}_{211}`$ |
| $`\begin{array}{c}C_{321}=\frac{4i}{\sqrt{30}}0.307324\\ C_{631}=C_{432}/\varphi \\ C_{653}=C_{432}\\ C_{765}=\varphi ^2C_{321}\end{array}`$ | $`\begin{array}{c}\gamma _3+\sigma ^{15}\gamma _2+\sigma ^{16}\gamma _1=0\\ \gamma _6+\sigma ^{13}\gamma _3+\sigma ^{20}\gamma _1=0\\ \gamma _6+\sigma ^9\gamma _5+\sigma ^{17}\gamma _3=0\\ \gamma _7+\sigma ^{15}\gamma _6+\sigma ^{16}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{321}=22.6132\stackrel{~}{\mathrm{\Delta }}_{321}`$ |
| $`\begin{array}{c}C_{421}=\frac{4i}{\sqrt{30}}0.972789\\ C_{641}=\varphi ^2C_{321}\\ C_{861}=\varphi ^3C_{321}\\ C_{865}=\varphi ^2C_{421}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^{13}\gamma _2+\sigma ^{19}\gamma _1=0\\ \gamma _6+\sigma ^{14}\gamma _4+\sigma ^{17}\gamma _1=0\\ \gamma _8+\sigma ^{14}\gamma _6+\sigma ^{18}\gamma _1=0\\ \gamma _8+\sigma ^{13}\gamma _6+\sigma ^{19}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{421}=9.92482\stackrel{~}{\mathrm{\Delta }}_{421}`$ |
| $`\begin{array}{c}C_{431}=\frac{4i}{\sqrt{30}}1.09848\\ C_{831}=C_{421}/\varphi \\ C_{871}=C_{421}\\ C_{875}=\varphi ^2C_{431}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^{13}\gamma _3+\sigma ^{24}\gamma _1=0\\ \gamma _8+\sigma ^{14}\gamma _3+\sigma ^{16}\gamma _1=0\\ \gamma _8+\sigma ^{14}\gamma _7+\sigma ^{27}\gamma _1=0\\ \gamma _8+\sigma ^{13}\gamma _7+\sigma ^{24}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{431}=8.65727\stackrel{~}{\mathrm{\Delta }}_{431}`$ |
| $`\begin{array}{c}C_{441}=\frac{4i}{\sqrt{30}}1.17616\\ C_{854}=C_{421}\\ C_{885}=\varphi ^2C_{441}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^{13}\gamma _4+\sigma ^{21}\gamma _1=0\\ \gamma _8+\sigma ^{13}\gamma _5+\sigma ^{16}\gamma _4=0\\ \gamma _8+\sigma ^{13}\gamma _8+\sigma ^{21}\gamma _5=0\end{array}`$ | $`\stackrel{~}{C}_{441}=14.4209\stackrel{~}{\mathrm{\Delta }}_{441}`$ |
| $`\begin{array}{c}C_{541}=C_{421}/\varphi \\ C_{554}=\varphi ^3C_{321}\end{array}`$ | $`\begin{array}{c}\gamma _5+\sigma ^{14}\gamma _4+\sigma ^{26}\gamma _1=0\\ \gamma _5+\sigma ^7\gamma _5+\sigma ^{19}\gamma _4=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{411}=5.35386\end{array}`$ |
| $`\begin{array}{c}C_{771}=C_{432}\end{array}`$ | $`\begin{array}{c}\gamma _7+\sigma ^{14}\gamma _7+\sigma ^{22}\gamma _1=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{331}=4.27444\stackrel{~}{\mathrm{\Delta }}_{331}\end{array}`$ |
| $`\begin{array}{c}C_{222}=\frac{4i}{\sqrt{30}}1.71313\\ C_{622}=\frac{4i}{\sqrt{30}}1.88133\\ C_{662}=\varphi C_{622}\\ C_{666}=\varphi ^2C_{222}\end{array}`$ | $`\begin{array}{c}\gamma _2+\sigma ^{10}\gamma _2+\sigma ^{20}\gamma _2=0\\ \gamma _6+\sigma ^{18}\gamma _2+\sigma ^{12}\gamma _2=0\\ \gamma _6+\sigma ^{12}\gamma _6+\sigma ^{21}\gamma _2=0\\ \gamma _6+\sigma ^{10}\gamma _6+\sigma ^{20}\gamma _6=0\end{array}`$ | $`\stackrel{~}{C}_{222}=9.1965\stackrel{~}{\mathrm{\Delta }}_{222}`$ |
| $`\begin{array}{c}C_{322}=\frac{4i}{\sqrt{30}}1.96731\\ C_{762}=C_{432}\\ C_{766}=\varphi ^2C_{322}\end{array}`$ | $`\begin{array}{c}\gamma _3+\sigma ^{12}\gamma _2+\sigma ^{19}\gamma _2=0\\ \gamma _7+\sigma ^{14}\gamma _6+\sigma ^{18}\gamma _2=0\\ \gamma _7+\sigma ^{12}\gamma _6+\sigma ^{19}\gamma _6=0\end{array}`$ | $`\stackrel{~}{C}_{322}=3.75947\stackrel{~}{\mathrm{\Delta }}_{322}`$ |
| $`\begin{array}{c}C_{522}=C_{432}/\varphi \\ C_{652}=\varphi ^3C_{321}\end{array}`$ | $`\begin{array}{c}\gamma _5+\sigma ^{10}\gamma _2+\sigma ^{21}\gamma _2=0\\ \gamma _6+\sigma ^{12}\gamma _5+\sigma ^{17}\gamma _2=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{122}=9.55253\stackrel{~}{\mathrm{\Delta }}_{122}\end{array}`$ |
| --- | --- | --- |
| $`\begin{array}{c}C_{822}=\varphi ^2C_{321}\end{array}`$ | $`\begin{array}{c}\gamma _8+\sigma ^{14}\gamma _2+\sigma ^{16}\gamma _2=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{422}=0.683318\stackrel{~}{\mathrm{\Delta }}_{422}\end{array}`$ |
| $`\begin{array}{c}C_{432}=\frac{4i}{\sqrt{30}}2.37859\\ C_{832}=\varphi ^2C_{431}\\ C_{863}=\varphi ^3C_{431}\\ C_{876}=\varphi ^2C_{432}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^{11}\gamma _3+\sigma ^{22}\gamma _2=0\\ \gamma _8+\sigma ^{12}\gamma _3+\sigma ^{19}\gamma _2=0\\ \gamma _8+\sigma ^{11}\gamma _6+\sigma ^{19}\gamma _3=0\\ \gamma _8+\sigma ^{11}\gamma _7+\sigma ^{22}\gamma _6=0\end{array}`$ | $`\stackrel{~}{C}_{432}=15.2555\stackrel{~}{\mathrm{\Delta }}_{432}`$ |
| $`\begin{array}{c}C_{732}=C_{432}/\varphi \\ C_{772}=\varphi ^3C_{431}\end{array}`$ | $`\begin{array}{c}\gamma _7+\sigma ^{14}\gamma _3+\sigma ^{17}\gamma _2=0\\ \gamma _7+\sigma ^{13}\gamma _7+\sigma ^{22}\gamma _2=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{332}=2.68177\stackrel{~}{\mathrm{\Delta }}_{332}\end{array}`$ |
| $`\begin{array}{c}C_{333}=\frac{4i}{\sqrt{30}}3.78439\\ C_{733}=\frac{4i}{\sqrt{30}}4.15597\\ C_{773}=\varphi C_{733}\\ C_{777}=\varphi ^2C_{333}\end{array}`$ | $`\begin{array}{c}\gamma _3+\sigma ^{10}\gamma _3+\sigma ^{20}\gamma _3=0\\ \gamma _7+\sigma ^{12}\gamma _3+\sigma ^{18}\gamma _3=0\\ \gamma _7+\sigma ^{12}\gamma _7+\sigma ^{21}\gamma _3=0\\ \gamma _7+\sigma ^{10}\gamma _7+\sigma ^{20}\gamma _7=0\end{array}`$ | $`\stackrel{~}{C}_{333}=7.1965\stackrel{~}{\mathrm{\Delta }}_{333}`$ |
| $`\begin{array}{c}C_{433}=\frac{4i}{\sqrt{30}}3.24742\\ C_{743}=\varphi ^2C_{431}\\ C_{877}=\varphi ^2C_{433}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^9\gamma _3+\sigma ^{20}\gamma _3=0\\ \gamma _7+\sigma ^{13}\gamma _4+\sigma ^{17}\gamma _3=0\\ \gamma _8+\sigma ^9\gamma _7+\sigma ^{20}\gamma _7=0\end{array}`$ | $`\stackrel{~}{C}_{433}=9.22437\stackrel{~}{\mathrm{\Delta }}_{433}`$ |
| $`\begin{array}{c}C_{553}=C_{421}\end{array}`$ | $`\begin{array}{c}\gamma _5+\sigma ^4\gamma _5+\sigma ^{17}\gamma _3=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{113}=2.5468\end{array}`$ |
| $`\begin{array}{c}C_{444}=\frac{4i}{\sqrt{30}}2.50428\\ C_{844}=\frac{4i}{\sqrt{30}}2.75016\\ C_{884}=\varphi C_{844}\\ C_{888}=\varphi ^2C_{444}\end{array}`$ | $`\begin{array}{c}\gamma _4+\sigma ^{10}\gamma _4+\sigma ^{20}\gamma _4=0\\ \gamma _8+\sigma ^{12}\gamma _4+\sigma ^{18}\gamma _4=0\\ \gamma _8+\sigma ^{12}\gamma _8+\sigma ^{21}\gamma _4=0\\ \gamma _8+\sigma ^{10}\gamma _8+\sigma ^{20}\gamma _8=0\end{array}`$ | $`\stackrel{~}{C}_{444}=29.3768\stackrel{~}{\mathrm{\Delta }}_{444}`$ |
| $`\begin{array}{c}C_{644}=\varphi ^2C_{431}\end{array}`$ | $`\begin{array}{c}\gamma _6+\sigma ^{11}\gamma _4+\sigma ^{19}\gamma _4=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{244}=2.13686\stackrel{~}{\mathrm{\Delta }}_{244}\end{array}`$ |
| $`\begin{array}{c}C_{744}=C_{421}/\varphi \\ C_{874}=\varphi ^3C_{431}\end{array}`$ | $`\begin{array}{c}\gamma _7+\sigma ^{15}\gamma _4+\sigma ^{16}\gamma _4=0\\ \gamma _8+\sigma ^{12}\gamma _7+\sigma ^{23}\gamma _4=0\end{array}`$ | $`\begin{array}{c}\stackrel{~}{C}_{344}=8.34233\stackrel{~}{\mathrm{\Delta }}_{344}\end{array}`$ |
Also here we did not report the overall factor of $`i\beta m^24/\sqrt{30}`$ in $`\stackrel{~}{C}`$. Note that $`\stackrel{~}{C}_{411}`$ and $`\stackrel{~}{C}_{311}`$ have no classical mass triangle associated to them and therefore yield no poles in the propagators. The nonvanishing one-loop contributions are therefore
From this it follows that the classical mass ratios are not conserved in the quantum field theory
$$\frac{\delta m_1^2}{m_1^2}=54045.1\mathrm{}\frac{\delta m_2^2}{m_2^2}=68239.3\mathrm{}\frac{\delta m_3^2}{m_3^2}=11488.2\mathrm{}\frac{\delta m_4^2}{m_4^2}=2914.28\mathrm{}$$
(79)
Hence we have the same type of behaviour under renormalisation as in the $`H_3`$-ATFT obtained from the reduction of the $`D_6`$-ATFT.
## 4 Conclusions
With regard to previously studied ATFT, the embedding of non-crystallographic into crystallographic Coxeter groups leads to an explanation for the fact that in some theories the masses can be organised into pairs such that one mass differs from the other only by a factor of $`\varphi `$. This also holds for the higher charges.
We showed that it is possible to construct ATFT related to non-crystallographic Coxeter groups despite the fact that there is no Lie algebra associated to them. The construction is possible since one may exploit the embedding of non-crystallographic into crystallographic Coxeter groups, making use of the fact that the latter do possess a Lie algebraic structure and thus preserving integrability. Unlike the folding from simply laced Lie algebras to non-simply laced Lie algebras , the resulting theories we obtain here are not equivalent to a direct formulation of the theory in terms of non-crystallographic Coxeter groups. In this context the reduction procedure is vital for consistency and not merely an additional structure. It is of course possible to write down a Lagrangian of the form (52) involving directly roots of $`\stackrel{~}{\mathrm{\Delta }}`$, but it remains to be seen whether such theories are consistent and especially if they are classically integrable.
With regard to the quantum field theory, our computations showed for the $`H_3,H_4`$-ATFT that the masses of these new theories do not renormalise with an overall factor, i.e. $`\delta m_k^2/m_k^2`$ is not a universal constant for all particle types $`k`$, preventing that the classical mass ratios can be maintained in the quantum field theory. Remarkably this was true for ATFT related to simply laced Lie algebras, which allowed for a relatively straightforward construction of the scattering matrices . For ATFT related to non-simply laced Lie algebras this was found no longer to be true, such that different types of scenarios had to be devised. One is to have floating masses such that depending on the coupling constant the masses flow from one Lie algebra in the weak limit to its Langlands dual in the strong coupling limit . The other alternative proposal was to introduce additional Fermions into the model , which compensate for the unequal mass shifts. From our analysis it is clear that the construction of a consistent quantum field theory for the proposed $`H_3,H_4`$-theories has to be modelled along the line of the construction of theories related to non-simply laced algebras due to unequal mass renormalisations for each individual particle. It remains to be seen in future work, which of the prescriptions will be successful in this context.
As we showed, the behaviour under renormalisation is different for the $`H_2`$-ATFT, where the classical mass ratios remain preserved up to first order perturbation theory. Despite this, it is not immediately obvious how to write down a scattering matrix to all orders in perturbation theory.
Our detailed analysis of the embedding of non-crystallographic into crystallographic Coxeter groups allows one to apply the aforementioned reduction method to a wide range of application in physics, chemistry and biology, where Coxeter groups play a role. For instance in our forthcoming publication we will apply this method to the generalized Calogero-Moser models.
Acknowledgments. C.K. is financially supported by a University Research Fellowship of the Royal Society.
Appendix
## Appendix A The orbits of $`H_3`$ and $`D_6`$
Successive action of $`\sigma =`$ $`\sigma _1\sigma _4\sigma _3\sigma _6\sigma _2\sigma _5`$ and $`\stackrel{~}{\sigma }=`$ $`\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_2`$ yields
| | $`\mathrm{\Omega }_1`$ | $`\omega (\mathrm{\Omega }_1)=\stackrel{~}{\mathrm{\Omega }}_1`$ | $`\mathrm{\Omega }_4`$ | $`\omega (\mathrm{\Omega }_4)=\varphi \stackrel{~}{\mathrm{\Omega }}_1`$ |
| --- | --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_1`$ |
| $`\sigma ^1`$ | $`\alpha _2+\alpha _6`$ | $`\stackrel{~}{\alpha }_2+\varphi \stackrel{~}{\alpha }_3`$ | $`\alpha _3+\alpha _5+\alpha _6`$ | $`\varphi \stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3`$ |
| $`\sigma ^2`$ | $`\alpha _3+\alpha _4+\alpha _5`$ | $`\varphi \stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3`$ | $`\alpha _1+\alpha _2+\alpha _4+\alpha _5+\alpha _6`$ | $`\varphi ^2\stackrel{~}{\alpha }_1+\varphi ^2\stackrel{~}{\alpha }_2+\varphi \stackrel{~}{\alpha }_3`$ |
| $`\sigma ^3`$ | $`\alpha _5+\alpha _6`$ | $`\varphi \left(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3\right)`$ | $`\alpha _2+\alpha _3+\alpha _5+\alpha _6`$ | $`\varphi ^2\left(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3\right)`$ |
| $`\sigma ^4`$ | $`\alpha _1+\alpha _2`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2`$ | $`\alpha _4+\alpha _5`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)`$ |
| $`\sigma ^5`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _4`$ | $`\varphi \stackrel{~}{\alpha }_1`$ |
| | $`\mathrm{\Omega }_2`$ | $`\omega (\mathrm{\Omega }_2)=\stackrel{~}{\mathrm{\Omega }}_2`$ | $`\mathrm{\Omega }_5`$ |
| --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _5`$ |
| $`\sigma ^1`$ | $`\alpha _1+\alpha _2+\alpha _6`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\varphi \stackrel{~}{\alpha }_3`$ | $`\alpha _3+\alpha _4+\alpha _5+\alpha _6`$ |
| $`\sigma ^2`$ | $`\alpha _2+\alpha _3+\alpha _4+\alpha _5+\alpha _6`$ | $`\varphi \stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _4+2\alpha _5+2\alpha _6`$ |
| $`\sigma ^3`$ | $`\alpha _3+\alpha _4+2\alpha _5+\alpha _6`$ | $`\varphi \stackrel{~}{\alpha }_1+2\varphi \stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3`$ | $`\alpha _1+2\alpha _2+\alpha _3+\alpha _4+2\alpha _5+2\alpha _6`$ |
| $`\sigma ^4`$ | $`\alpha _1+\alpha _2+\alpha _5+\alpha _6`$ | $`\stackrel{~}{\alpha }_1+\varphi ^2\stackrel{~}{\alpha }_2+\varphi \stackrel{~}{\alpha }_3`$ | $`\alpha _2+\alpha _3+\alpha _4+2\alpha _5+\alpha _6`$ |
| $`\sigma ^5`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _5`$ |
| | $`\mathrm{\Omega }_3`$ | $`\omega (\mathrm{\Omega }_3)=\stackrel{~}{\mathrm{\Omega }}_3`$ | $`\mathrm{\Omega }_6`$ |
| $`\sigma ^0`$ | $`\alpha _3`$ | $`\stackrel{~}{\alpha }_3`$ | $`\alpha _6`$ |
| $`\sigma ^1`$ | $`\alpha _4+\alpha _5+\alpha _6`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _4+\alpha _5+\alpha _6`$ |
| $`\sigma ^2`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _5+\alpha _6`$ | $`\stackrel{~}{\alpha }_1+\varphi ^2\stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3`$ | $`\alpha _2+\alpha _3+\alpha _4+2\alpha _5+2\alpha _6`$ |
| $`\sigma ^3`$ | $`\alpha _2+\alpha _4+\alpha _5+\alpha _6`$ | $`\varphi (\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _4+2\alpha _5+\alpha _6`$ |
| $`\sigma ^4`$ | $`\alpha _3+\alpha _5`$ | $`\varphi \stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3`$ | $`\alpha _2+\alpha _5+\alpha _6`$ |
| $`\sigma ^5`$ | $`\alpha _3`$ | $`\stackrel{~}{\alpha }_3`$ | $`\alpha _6`$ |
We did not write here the additional $`\stackrel{~}{\sigma }^i`$ in the first column. The identities $`\omega (\mathrm{\Omega }_4)=\varphi \stackrel{~}{\mathrm{\Omega }}_1`$, $`\omega (\mathrm{\Omega }_4)=\varphi \stackrel{~}{\mathrm{\Omega }}_1`$and $`\omega (\mathrm{\Omega }_6)=\varphi \stackrel{~}{\mathrm{\Omega }}_3`$ follow upon using (2).
## Appendix B The orbits of $`H_4`$ and $`E_8`$
Successive action of $`\sigma =`$ $`\sigma _1\sigma _5\sigma _3\sigma _7\sigma _2\sigma _6\sigma _4\sigma _8`$ and $`\stackrel{~}{\sigma }=`$ $`\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_4`$ yields
| | $`\mathrm{\Omega }_1`$ | $`\omega (\mathrm{\Omega }_1)=\stackrel{~}{\mathrm{\Omega }}_1`$ | $`\mathrm{\Omega }_5`$ |
| --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _5`$ |
| $`\sigma ^1`$ | $`\alpha _2+\alpha _3`$ | $`\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3`$ | $`\alpha _6+\alpha _7`$ |
| $`\sigma ^2`$ | $`\alpha _7+\alpha _8`$ | $`\varphi (\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4)`$ | $`\alpha _3+\alpha _4+\alpha _7+\alpha _8`$ |
| $`\sigma ^3`$ | $`\alpha _4+\alpha _5+\alpha _7+\alpha _8`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)+\stackrel{~}{\alpha }_4`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _5\hfill \\ +\alpha _6+\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^4`$ | $`\alpha _3+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi \stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3+\varphi \stackrel{~}{\alpha }_4`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _6\hfill \\ +2\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^5`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _4+\alpha _7+\alpha _8`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\varphi ^2(\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _5+\alpha _6\hfill \\ +2\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^6`$ | $`\alpha _2+\alpha _3+\alpha _5+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_4)+\varphi ^2(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\hfill \\ +\alpha _5+2\alpha _6+2\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^7`$ | $`\alpha _4+\alpha _6+2\alpha _7+\alpha _8`$ | $`\varphi \stackrel{~}{\alpha }_2+2\varphi \stackrel{~}{\alpha }_3+\varphi ^2\stackrel{~}{\alpha }_4`$ | $`\begin{array}{c}\alpha _2+2\alpha _3+\alpha _4+\alpha _5\hfill \\ +\alpha _6+2\alpha _7+2\alpha _8\hfill \end{array}`$ |
| $`\sigma ^8`$ | $`\alpha _3+\alpha _4+\alpha _5+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2)+\varphi ^2(\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\hfill \\ +\alpha _5+\alpha _6+2\alpha _7+2\alpha _8\hfill \end{array}`$ |
| $`\sigma ^9`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _6+\alpha _7+\alpha _8`$ | $`\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_4+\varphi ^2(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\hfill \\ +2\alpha _6+2\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^{10}`$ | $`\alpha _2+\alpha _3+\alpha _4+\alpha _6+\alpha _7+\alpha _8`$ | $`\stackrel{~}{\alpha }_2+\varphi ^2(\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _6+2\alpha _7\hfill \\ +2\alpha _8\hfill \end{array}`$ |
| $`\sigma ^{11}`$ | $`\alpha _5+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi (\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\hfill \\ +\alpha _5+\alpha _6+\alpha _7+\alpha _8\hfill \end{array}`$ |
| $`\sigma ^{12}`$ | $`\alpha _4+\alpha _6+\alpha _7`$ | $`\varphi (\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)+\stackrel{~}{\alpha }_4`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _6+\alpha _7\hfill \\ +\alpha _8\hfill \end{array}`$ |
| $`\sigma ^{13}`$ | $`\alpha _3+\alpha _8`$ | $`\stackrel{~}{\alpha }_3+\varphi \stackrel{~}{\alpha }_4`$ | $`\alpha _4+\alpha _7+\alpha _8`$ |
| $`\sigma ^{14}`$ | $`\alpha _1+\alpha _2`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2`$ | $`\alpha _5+\alpha _6`$ |
| $`\sigma ^{15}`$ | $`\alpha _1`$ | $`\stackrel{~}{\alpha }_1`$ | $`\alpha _5`$ |
| | $`\mathrm{\Omega }_2`$ | $`\omega (\mathrm{\Omega }_2)=\stackrel{~}{\mathrm{\Omega }}_2`$ | $`\mathrm{\Omega }_6`$ |
| --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _6`$ |
| $`\sigma ^1`$ | $`\alpha _1+\alpha _2+\alpha _3`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3`$ | $`\alpha _5+\alpha _6+\alpha _7`$ |
| $`\sigma ^2`$ | $`\alpha _2+\alpha _3+\alpha _7+\alpha _8`$ | $`\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{3}^{}+`$$`\varphi `$$`\stackrel{~}{\alpha }_4`$ | $`\alpha _3+\alpha _4+\alpha _6+2\alpha _7+\alpha _8`$ |
| $`\sigma ^3`$ | $`\alpha _4+\alpha _5+\alpha _6+2\alpha _7+\alpha _8`$ | $`\varphi `$$`(\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+2\stackrel{~}{\alpha }_3+`$$`\varphi `$$`\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2\alpha _3+\alpha _4\\ +\alpha _5+\alpha _6+2\alpha _7+2\alpha _8\end{array}`$ |
| $`\sigma ^4`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _5\\ \hfill +2(\alpha _6+\alpha _7)+\alpha _8\end{array}`$ | $`\varphi `$$`(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_2)+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{3}^{}+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+2\alpha _3+\alpha _4\\ +\alpha _5+2\alpha _6+3\alpha _7+2\alpha _8\end{array}`$ |
| $`\sigma ^5`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _4+\alpha _6\\ \hfill +2(\alpha _3+\alpha _7+\alpha _8)\end{array}`$ | $`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_2+2\stackrel{~}{\alpha }_3)+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _2+2\alpha _3+2\alpha _4+\alpha _5\\ +2\alpha _6+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^6`$ | $`\begin{array}{c}\alpha _1+\alpha _4+\alpha _5+\alpha _6\\ +2(\alpha _2+\alpha _3+\alpha _7+\alpha _8)\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_3+`$$`\varphi `$$`\stackrel{~}{\alpha }_4)+(2+`$$`\varphi `$$`)\stackrel{~}{\alpha }_2`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2\alpha _3+2\alpha _4\\ +2\alpha _5+3\alpha _6+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^7`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\\ +3\alpha _7+2(\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+(`$$`\varphi `$$`{}_{}{}^{4}1)\stackrel{~}{\alpha }_3+`$$`\varphi `$$`{}_{}{}^{3}(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+3\alpha _3+2\alpha _4\\ +\alpha _5+3\alpha _6+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^8`$ | $`\begin{array}{c}\alpha _3+\alpha _5+3\alpha _7\\ +2(\alpha _4+\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`(\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+2`$$`\varphi `$$`\stackrel{~}{\alpha }_4)+(`$$`\varphi `$$`{}_{}{}^{4}1)\stackrel{~}{\alpha }_3`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+3\alpha _3+2\alpha _4\\ +\alpha _5+2\alpha _6+4\alpha _7+4\alpha _8\end{array}`$ |
| $`\sigma ^9`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _4+\alpha _5\\ +2(\alpha _3+\alpha _6+\alpha _7+\alpha _8)\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{1}{}^{}+2\stackrel{~}{\alpha }{}_{3}{}^{})+`$$`\varphi `$$`{}_{}{}^{3}(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+2\alpha _3+2\alpha _4\\ +2\alpha _5+3\alpha _6+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^{10}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _3)+\alpha _4\\ +\alpha _6+2(\alpha _7+\alpha _8)\end{array}`$ | $`\stackrel{~}{\alpha }_1+\left(2+\varphi \right)\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}(2\stackrel{~}{\alpha }_3+`$$`\varphi `$$`\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _2+2\alpha _3+2\alpha _4+\alpha _5\\ +3\alpha _6+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^{11}`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\\ +\alpha _6+2\alpha _7+2\alpha _8\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{2}^{}+`$$`\varphi `$$`{}_{}{}^{3}(\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2\alpha _3+2\alpha _4\\ +\alpha _5+2\alpha _6+3\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^{12}`$ | $`\alpha _4+\alpha _5+2(\alpha _6+\alpha _7)+\alpha _8`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+2`$$`\varphi `$$`(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3)+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+2\alpha _3+\alpha _4\\ +\alpha _5+2\alpha _6+2\alpha _7+2\alpha _8\end{array}`$ |
| $`\sigma ^{13}`$ | $`\alpha _3+\alpha _4+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi `$$`\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\alpha _2+\alpha _3+\alpha _4+\alpha _6+2(\alpha _7+\alpha _8)`$ |
| $`\sigma ^{14}`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _8`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3+`$$`\varphi `$$`\stackrel{~}{\alpha }_4`$ | $`\alpha _4+\alpha _5+\alpha _6+\alpha _7+\alpha _8`$ |
| $`\sigma ^{15}`$ | $`\alpha _2`$ | $`\stackrel{~}{\alpha }_2`$ | $`\alpha _6`$ |
| | $`\mathrm{\Omega }_3`$ | $`\omega (\mathrm{\Omega }_3)=\stackrel{~}{\mathrm{\Omega }}_3`$ | $`\mathrm{\Omega }_7`$ |
| --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _3`$ | $`\stackrel{~}{\alpha }_3`$ | $`\alpha _7`$ |
| $`\sigma ^1`$ | $`\alpha _1+\alpha _2+\alpha _3+\alpha _7+\alpha _8`$ | $`\stackrel{~}{\alpha }_1+\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{3}^{}+`$$`\varphi `$$`\stackrel{~}{\alpha }_4`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _5\\ +\alpha _6+2\alpha _7+\alpha _8\end{array}`$ |
| $`\sigma ^2`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\\ +\alpha _6+2\alpha _7+\alpha _8\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{2}^{}+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{3}^{}+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _4+\alpha _5\\ +3\alpha _7+2(\alpha _3+\alpha _6+\alpha _8)\end{array}`$ |
| $`\sigma ^3`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _5\\ +3\alpha _7+2(\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_2)+(`$$`\varphi `$$`{}_{}{}^{4}1)\stackrel{~}{\alpha }_3+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _6)\\ +\alpha _5+4\alpha _7+3(\alpha _3+\alpha _8)\end{array}`$ |
| $`\sigma ^4`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2(\alpha _3+\alpha _4)\\ +\alpha _5+3\alpha _7+2(\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_1+`$$`\varphi `$$`\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{3}^{}+2\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _5)\\ +3(\alpha _3+\alpha _6)+5\alpha _7+4\alpha _8\end{array}`$ |
| $`\sigma ^5`$ | $`\begin{array}{c}\alpha _1+\alpha _4+2(\alpha _2+\alpha _6)\\ +\alpha _5+3(\alpha _3+\alpha _7+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }{}_{1}{}^{}+2\stackrel{~}{\alpha }{}_{2}{}^{}+3\stackrel{~}{\alpha }{}_{3}{}^{})\\ +\left(\varphi ^41\right)\stackrel{~}{\alpha }_4\end{array}`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+3(\alpha _3+\alpha _4)\\ +2\alpha _5+6\alpha _7+4(\alpha _6+\alpha _8)\end{array}`$ |
| $`\sigma ^6`$ | $`\begin{array}{c}2(\alpha _2+\alpha _3+\alpha _4+\alpha _6)\\ +\alpha _1+\alpha _5+4\alpha _7+3\alpha _8\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_2+2`$$`\varphi `$$`\stackrel{~}{\alpha }_3+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{4}^{})`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _5)+3\alpha _4\\ +4(\alpha _3+\alpha _6)+6\alpha _7+5\alpha _8\end{array}`$ |
| $`\sigma ^7`$ | $`\begin{array}{c}\alpha _2+2(\alpha _3+\alpha _4+\alpha _5)\\ +4\alpha _7+3(\alpha _6+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi (\stackrel{~}{\alpha }_1+2\varphi ^2\stackrel{~}{\alpha }_3+\varphi ^3\stackrel{~}{\alpha }_4)\\ +\left(\varphi ^41\right)\stackrel{~}{\alpha }_2\end{array}`$ | $`\begin{array}{c}2(\alpha _1+\alpha _5)+3(\alpha _2+\alpha _4)\\ +4(\alpha _3+\alpha _6)+6\alpha _7+5\alpha _8\end{array}`$ |
| $`\sigma ^8`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2(\alpha _3+\alpha _4)\\ +\alpha _5+4\alpha _7+3(\alpha _6+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }_1+2\varphi \stackrel{~}{\alpha }_3+\varphi ^2\stackrel{~}{\alpha }_4)\\ +\left(\varphi ^41\right)\stackrel{~}{\alpha }_2\end{array}`$ | $`\begin{array}{c}\alpha _1+3(\alpha _2+\alpha _4)+2\alpha _5\\ +4(\alpha _3+\alpha _6)+6\alpha _7+5\alpha _8\end{array}`$ |
| $`\sigma ^9`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _6)\\ +\alpha _5+3(\alpha _3+\alpha _7+\alpha _8)\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{1}{}^{}+2\stackrel{~}{\alpha }{}_{2}{}^{}+3\stackrel{~}{\alpha }{}_{3}{}^{})+`$$`\varphi `$$`{}_{}{}^{4}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _5)+4\alpha _6\\ +3(\alpha _3+\alpha _4)+6\alpha _7+5\alpha _8\end{array}`$ |
| $`\sigma ^{10}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _3+\alpha _6)\\ +\alpha _4+\alpha _5+3(\alpha _7+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3)\\ +\left(\varphi ^41\right)\stackrel{~}{\alpha }_4\end{array}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _5)+4\alpha _6\\ +3(\alpha _3+\alpha _4)+5\alpha _7+4\alpha _8\end{array}`$ |
| $`\sigma ^{11}`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _5+3\alpha _7\\ +2(\alpha _4+\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{2}^{}+(`$$`\varphi `$$`{}_{}{}^{4}1)\stackrel{~}{\alpha }_3+2`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{4}^{}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4)+\alpha _5\\ +3(\alpha _3+\alpha _6)+4(\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{12}`$ | $`\begin{array}{c}\alpha _3+\alpha _4+\alpha _5\\ +2(\alpha _6+\alpha _7+\alpha _8)\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+2`$$`\varphi `$$`\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{3}(\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _3+\alpha _4)\\ +\alpha _5+2\alpha _6+3(\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{13}`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _6+\alpha _7+\alpha _8\end{array}`$ | $`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{2}{}^{}+\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\\ +2(\alpha _6+\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{14}`$ | $`\alpha _2+\alpha _3+\alpha _8`$ | $`\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_3+`$$`\varphi `$$`\stackrel{~}{\alpha }_4`$ | $`\alpha _4+\alpha _6+\alpha _7+\alpha _8`$ |
| $`\sigma ^{15}`$ | $`\alpha _3`$ | $`\stackrel{~}{\alpha }_3`$ | $`\alpha _7`$ |
| | $`\mathrm{\Omega }_4`$ | $`\omega (\mathrm{\Omega }_4)=\stackrel{~}{\mathrm{\Omega }}_4`$ | $`\mathrm{\Omega }_8`$ |
| --- | --- | --- | --- |
| $`\sigma ^0`$ | $`\alpha _4`$ | $`\stackrel{~}{\alpha }_4`$ | $`\alpha _8`$ |
| $`\sigma ^1`$ | $`\alpha _4+\alpha _7`$ | $`\varphi `$$`\stackrel{~}{\alpha }_3+\stackrel{~}{\alpha }_4`$ | $`\alpha _3+\alpha _7+\alpha _8`$ |
| $`\sigma ^2`$ | $`\alpha _3+\alpha _5+\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi `$$`(\stackrel{~}{\alpha }{}_{1}{}^{}+\stackrel{~}{\alpha }{}_{2}{}^{}+\varphi \stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _5+\alpha _6+2\alpha _7+\alpha _8\end{array}`$ |
| $`\sigma ^3`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _6+2\alpha _7+\alpha _8\end{array}`$ | $`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_4)+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{3}^{}`$ | $`\begin{array}{c}\alpha _2+\alpha _4+\alpha _5+3\alpha _7\\ +2(\alpha _3+\alpha _6+\alpha _8)\end{array}`$ |
| $`\sigma ^4`$ | $`\begin{array}{c}\alpha _2+\alpha _4+\alpha _5+\alpha _6\\ +2(\alpha _3+\alpha _7+\alpha _8)\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{2}{}^{}+2\stackrel{~}{\alpha }{}_{3}{}^{}+\varphi \stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2(\alpha _3+\alpha _4+\alpha _6)\\ +\alpha _5+4\alpha _7+3\alpha _8\end{array}`$ |
| $`\sigma ^5`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _5+3\alpha _7+2(\alpha _6+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2+\varphi \stackrel{~}{\alpha }_4)\\ +(\varphi ^41)\stackrel{~}{\alpha }_3\end{array}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _5)\\ +3(\alpha _3+\alpha _6+\alpha _8)+4\alpha _7\end{array}`$ |
| $`\sigma ^6`$ | $`\begin{array}{c}\alpha _2+2(\alpha _3+\alpha _4)+\alpha _5\\ +3\alpha _7+2(\alpha _6+\alpha _8)\end{array}`$ | $`\varphi `$$`(\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{2}^{}+`$$`\varphi `$$`{}_{}{}^{3}\stackrel{~}{\alpha }_{3}^{}+`$$`\varphi `$$`\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4)+\alpha _5\\ +3(\alpha _3+\alpha _6)+5\alpha _7+4\alpha _8\end{array}`$ |
| $`\sigma ^7`$ | $`\begin{array}{c}\alpha _1+\alpha _2+2(\alpha _3+\alpha _6)\\ +\alpha _4+\alpha _5+3(\alpha _7+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }_1+\varphi \stackrel{~}{\alpha }_2+\varphi ^2\stackrel{~}{\alpha }_3)\\ +(\varphi ^41)\stackrel{~}{\alpha }_4\end{array}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _5)+5\alpha _7\\ +3(\alpha _3+\alpha _4+\alpha _6)+4\alpha _8\end{array}`$ |
| $`\sigma ^8`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _3+\alpha _4)\\ +\alpha _5+2(\alpha _6+\alpha _8)+3\alpha _7\end{array}`$ | $`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_2+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{3}^{}+2\stackrel{~}{\alpha }_4)`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _5)\\ +3\alpha _3+5\alpha _7+4(\alpha _6+\alpha _8)\end{array}`$ |
| $`\sigma ^9`$ | $`\begin{array}{c}\alpha _2+\alpha _4+\alpha _5\\ +2(\alpha _6+\alpha _3)+3(\alpha _7+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi \stackrel{~}{\alpha }_1+\varphi ^3\stackrel{~}{\alpha }_2+\varphi ^4\stackrel{~}{\alpha }_3\\ +(\varphi ^41)\stackrel{~}{\alpha }_4\end{array}`$ | $`\begin{array}{c}\alpha _1+2\alpha _2+\alpha _5+5\alpha _7\\ +3(\alpha _3+\alpha _4+\alpha _6)+4\alpha _8\end{array}`$ |
| $`\sigma ^{10}`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _5\\ +2(\alpha _4+\alpha _6+\alpha _8)+3\alpha _7\end{array}`$ | $`\begin{array}{c}\varphi ^2(\stackrel{~}{\alpha }_1+2\stackrel{~}{\alpha }_4)+\varphi ^3\stackrel{~}{\alpha }_2\\ +\left(\varphi ^41\right)\stackrel{~}{\alpha }_3\end{array}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _4+\alpha _5)\\ +3(\alpha _3+\alpha _6)+4(\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{11}`$ | $`\begin{array}{c}\alpha _2+\alpha _4+\alpha _5\\ +2(\alpha _3+\alpha _6+\alpha _7+\alpha _8)\end{array}`$ | $`\begin{array}{c}\varphi \stackrel{~}{\alpha }_1+\varphi ^3(\stackrel{~}{\alpha }_2+\stackrel{~}{\alpha }_4)\\ +2\varphi ^2\stackrel{~}{\alpha }_3\end{array}`$ | $`\begin{array}{c}\alpha _1+2(\alpha _2+\alpha _3+\alpha _4)\\ +\alpha _5+4\alpha _7+3(\alpha _6+\alpha _8)\end{array}`$ |
| $`\sigma ^{12}`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _6+2(\alpha _7+\alpha _8)\end{array}`$ | $`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}\stackrel{~}{\alpha }_{2}^{}+`$$`\varphi `$$`{}_{}{}^{3}(\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _2+2(\alpha _3+\alpha _4+\alpha _6)\\ +\alpha _5+3(\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{13}`$ | $`\begin{array}{c}\alpha _2+\alpha _3+\alpha _4+\alpha _5\\ +\alpha _6+\alpha _7+\alpha _8\end{array}`$ | $`\varphi `$$`\stackrel{~}{\alpha }_1+`$$`\varphi `$$`{}_{}{}^{2}(\stackrel{~}{\alpha }{}_{2}{}^{}+\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\begin{array}{c}\alpha _1+\alpha _2+\alpha _3+\alpha _4\\ +\alpha _5+2(\alpha _6+\alpha _7+\alpha _8)\end{array}`$ |
| $`\sigma ^{14}`$ | $`\alpha _6+\alpha _7+\alpha _8`$ | $`\varphi `$$`(\stackrel{~}{\alpha }{}_{2}{}^{}+\stackrel{~}{\alpha }{}_{3}{}^{}+\stackrel{~}{\alpha }{}_{4}{}^{})`$ | $`\alpha _2+\alpha _3+\alpha _4+\alpha _6+\alpha _7+\alpha _8`$ |
| $`\sigma ^{15}`$ | $`\alpha _4`$ | $`\stackrel{~}{\alpha }_4`$ | $`\alpha _8`$ |
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# Cohomology operations and the Deligne conjecture
## Introduction
In this note, all algebraic objects will be defined over a fixed field $`𝐤`$ of characteristic zero. An algebra means an algebra over a quadratic Koszul operad $`𝒫`$ \[26, II.3.3\]. This generality covers all “reasonable” algebras – associative, Lie, commutative associative, Poisson, Gerstenhaber, Leibniz, &c.
By the cohomology of a $`𝒫`$-algebra $`A`$ we mean the operadic cohomology $`H_𝒫^{}(A;A)`$ of $`A`$ with coefficients in itself \[26, II.3.100\], defined as the cohomology of the cochain complex $`C_𝒫^{}(A;A)=(C_𝒫^{}(A;A),d_𝒫)`$ recalled in A.6 of the appendix to this note. The complex $`C_𝒫^{}(A;A)`$ generalizes the “standard constructions” and $`H_𝒫^{}(A;A)`$ the “classical” cohomology (Hochschild for associative algebras, Chevalley-Eilenberg for Lie algebras, Harrison for associative commutative algebras, &c.) In general, $`H_𝒫^{}(A;A)`$ agrees with the triple cohomology \[8, Proposition 8.6\] and governs deformations of $`A`$ in the category of $`𝒫`$-algebras.
By a natural operation we mean a multilinear operation on $`H_𝒫^{}(A;A)`$ induced by a natural multilinear cochain operation on $`C_𝒫^{}(A;A)`$. Naturality means being defined using data that do not depend on a concrete algebra $`A`$ only. An example is the classical cup product $`f,gfg`$ of Hochschild cochains, resp. the induced graded commutative associative multiplication on the Hochschild cohomology of associative algebras . Our definition excludes some operations that are also “natural” in some sense, such as the degree zero unary operation defined as the projection $`\pi _n:H_𝒫^{}(A;A)H_𝒫^n(A;A)`$, $`n0`$, because this operation is not induced by any natural cochain map. A precise definition of natural operations is given in Section 7. Our aim is to describe the homotopy type of the dg operad $`_𝒫=\{_𝒫(n)\}_{n0}`$ of all these natural operations, see Problem 1 and its baby version Problem 20. The reward would be an ultimate understanding of the structure of the cohomology of a given type of algebras.
Our original hope was that the homotopy type of $`_𝒫`$ would be that of another Koszul quadratic operad $`𝒬_𝒫`$ determined by $`𝒫`$ in an explicit and simple manner. Examples we had the in mind were $`𝒫=𝒜\mathrm{𝑠𝑠}`$ for which probably $`𝒬_𝒫=𝒢\mathrm{𝑒𝑟}`$, the operad for Gerstenhaber algebras, and $`𝒫=\mathrm{𝑖𝑒}`$ for which probably $`𝒬_𝒫=\mathrm{𝑖𝑒}`$, the operad for Lie algebras. Calculations presented in this note however show that the homotopy type of $`_𝒫`$ is in general more complicated, therefore the property that makes the homotopy type of $`_𝒫`$ for $`𝒫=𝒜\mathrm{𝑠𝑠}`$ or $`𝒫=\mathrm{𝑖𝑒}`$ so nice must be finer than just the Koszulity of $`𝒫`$. We have no idea what this property is.
We feel that our formulations are somehow unsatisfactory – we would certainly prefer a concept that would not depend on a “representation” of the cohomology by a concrete cochain complex. In an ideal world, we should be working with natural operations in an appropriate “derived” category in which the cohomology is the hom-functor. The possibility of such a more conceptual approach for associative algebras and their Hochschild cohomology was indicated by , see also .
Another possibility could be to consider $`H_𝒫^{}(A;A)`$ as the cohomology of the cotangent complex of a suitable suitably derived stack of the variety of structure constants of $`𝒫`$-algebras, see , and study automorphisms of the point of this stack representing the algebra $`A`$. Our feeling is, however, that these fancier approaches, despite their beauty and generality, are still not developed enough to give concrete answers to concrete questions.
Let us explain the title of this note. In his famous letter , P. Deligne asked whether the Gerstenhaber algebra structure on the Hochschild cohomology of an associative algebra given by the cup product and the intrinsic bracket is induced by a natural action of singular chains on the little discs operad. There are several proofs of this so-called Deligne conjecture today . Assume one can prove that the operad of all natural operations on the Hochschild complex (that is, $`_{𝒜\mathrm{𝑠𝑠}}`$ in our notation) has the homotopy type of the operad for Gerstenhaber algebras. The formality of the little discs operad would then immediately imply the Deligne conjecture by simple homological considerations.
In fact, most of the proofs of the Deligne conjecture we are aware of , involve a conveniently chosen suboperad of $`_{𝒜\mathrm{𝑠𝑠}}`$ whose homotopy type is detected by Fiedorowicz’ recognition principle for $`E_2`$-operads . We will discuss these proofs in Section 6. Other proofs based on the Etingof-Kazhdan (de)quantization were given in . Several attempts have also been made to find a suitable filtration of the Fulton-MacPherson compactification of the configuration space of points in the plane to prove the conjecture . The Deligne conjecture has surprising implications for the existence of the deformation quantization of Poisson manifolds .
Acknowledgments. I would like to express my thanks to F. Chapoton, E. Getzler, V. Hinich M. Livernet, P. Somberg and A.A. Voronov for many useful comments and remarks. My special thanks are due to D. Tamarkin for inspiring discussions during my stay at the Northwestern University in April 2004.
## 1. Formulation of the problem
In this section we state the problems sketched out in the introduction more concretely and formulate also some conjectures. Let $`_𝒫=(_𝒫,\delta _𝒫)`$ be the dg-operad of all natural multilinear operations on the cochain complex $`C_𝒫^{}(A;A)=(C_𝒫^{}(A;A),d_𝒫)`$. The $`n`$-th component $`_𝒫(n)`$ of $`_𝒫`$ is the space of all $`n`$-linear natural operations $`C_𝒫^{}(A;A)^nC_𝒫^{}(A;A)`$ with the grading induced by the grading of $`C_𝒫^{}(A;A)`$: $`U_𝒫(n)`$ has degree $`d`$ if
$$U(f_1,\mathrm{},f_n)𝒞_𝒫^{m_1+\mathrm{}+m_n+d}(A;A),$$
whenever $`f_i𝒞_𝒫^{m_i}(A;A)`$ for $`1in`$. In this case we write $`U_𝒫^d(n)`$. Each $`_𝒫(n)`$ is equipped with the degree $`+1`$ differential $`\delta _𝒫`$ induced by the differential $`d_𝒫`$ of $`C_𝒫^{}(A;A)`$ in the usual way.
A precise definition of the operad $`_𝒫`$ is given in Section 7. Here we emphasize only that $`_𝒫^d(n)=0`$ for $`d<0`$ and that $`_𝒫(0)0`$ for any nontrivial $`𝒫`$. The central problem of the paper reads:
###### Problem 1.
Describe the homotopy type (in the non-abelian derived category) of the dg operad $`_𝒫`$. In particular, calculate the cohomology of $`_𝒫`$.
A baby-version of this problem is Problem 20 of Section 3. Closely related is:
###### Problem 2.
Find a property characterizing operads $`𝒫`$ for which $`_𝒫`$ is formal and has the homotopy type of some Koszul quadratic operad.
We will see, in Example 15, a simple quadratic Koszul operad $`𝒟`$ such that $`H^{}(_𝒟(0),\delta _𝒟)0`$. This clearly means that $`_𝒟`$ does not have the homotopy type of a quadratic Koszul operad, therefore the property answering Problem 2 must be stronger than Koszulness of $`𝒫`$.
Suppose that $`𝒫`$ is the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$ \[26, Remark II.1.15\]. In this case there exists a dg-suboperad $`_{\underset{¯}{𝒫}}`$ of $`_𝒫`$ consisting of natural operations that preserve the order of inputs of $`𝒫`$-cochains. For example, the classical cup product $`fgC_{𝒜\mathrm{𝑠𝑠}}^1(A;A)\mathrm{𝐿𝑖𝑛}(A^2,A)`$ of Hochschild cochains $`f,gC_{𝒜\mathrm{𝑠𝑠}}^0(A;A)\mathrm{𝐿𝑖𝑛}(A,A)`$ defined as
$$(fg)(ab):=f(a)g(b)\text{ for }abAA,$$
with $``$ denoting the associative multiplication of $`A`$, belongs to $`_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}`$, while the operation
$$U(f,g):=f(b)g(a)\text{ for }abAA,$$
does not, see Definition 34 of Section 7 for details.
Since $`_{\underset{¯}{𝒫}}(n)`$ is a $`\mathrm{\Sigma }_n`$-closed subspace of $`_𝒫(n)`$, $`n0`$, $`_{\underset{¯}{𝒫}}`$ is a usual, not only a non-$`\mathrm{\Sigma }`$, operad. We will see in Example 17 that, surprisingly, the homotopy type of $`_{\underset{¯}{𝒫}}`$ in general differs from the homotopy type of $`_𝒫`$. We therefore formulate:
###### Problem 3.
Let $`\underset{¯}{𝒫}`$ a non-$`\mathrm{\Sigma }`$ quadratic Koszul operad. Describe the homotopy type of the dg operad $`_{\underset{¯}{𝒫}}`$. In particular, calculate the cohomology of $`_{\underset{¯}{𝒫}}`$.
In Section 6(i) we give some indications that the operad $`_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}`$ has the homotopy type of the operad $`𝒢\mathrm{𝑒𝑟}`$ for Gerstenhaber algebras, see A.4 for a definition of $`𝒢\mathrm{𝑒𝑟}`$.
One may consider also strongly homotopy versions of the above problems. Recall that a strongly homotopy $`𝒫`$-algebra is, by , an algebra over the minimal model $`sh𝒫`$ of the operad $`𝒫`$. Let us denote by $`sh_𝒫=_{sh𝒫}`$ the dg-operad of natural operations on the cochain complex $`C_{sh𝒫}^{}(A,A)`$ for the cohomology of a strongly homotopy algebra $`A`$ with coefficients in itself. An example of this type of operad is the “minimal operad” $`M`$ considered in , which is a certain suboperad of $`_{sh𝒜\mathrm{𝑠𝑠}}`$, see Section 6(iii).
It is clear that there exists a canonical map $`_𝒫_{sh𝒫}`$, but simple examples show that, again rather surprisingly, this map is in general not a homotopy equivalence. Let us formulate:
###### Problem 4.
Describe the homotopy type of the dg-operad $`_{sh𝒫}`$ of natural operations on the cohomology of strongly homotopy $`𝒫`$-algebras.
Other problems formulated in this paper are Problem 16 of Section 2 and Problems 20,21 of Section 3.
Let us finally formulate also some conjectures. Although the operads $`_𝒫`$ and $`_{𝒫^!}`$ are not isomorphic (see Section 7), computational evidences together with an equivalence between the derived category of $`𝒫`$ algebras and the derived category of $`𝒫^!`$-algebras lead us to believe in:
###### Conjecture 5.
The homotopy type of the operad $`_𝒫`$ is the same as the homotopy type of $`_{𝒫^!}`$.
The following two conjectures concern the homotopy type of $`_𝒫`$ for $`𝒫=𝒜\mathrm{𝑠𝑠}`$ and $`𝒫=\mathrm{𝑖𝑒}`$.
###### Conjecture 6.
The operad $`_{𝒜\mathrm{𝑠𝑠}}`$ has the homotopy type of the operad $`𝒢\mathrm{𝑒𝑟}`$ for Gerstenhaber algebras.
Some results which may be helpful in the proof of the above conjecture are recalled in Section 6.
###### Conjecture 7.
The operad $`_{\mathrm{𝑖𝑒}}`$ has the homotopy type of the operad $`\mathrm{𝑖𝑒}`$.
According to a formality theorem \[24, Proposition 3.4\], it is enough to prove that
$$H^{}(_{\mathrm{𝑖𝑒}},\delta _{\mathrm{𝑖𝑒}})\mathrm{𝑖𝑒}.$$
Since $`H^0(_{\mathrm{𝑖𝑒}},\delta _{\mathrm{𝑖𝑒}})\mathrm{𝑖𝑒}`$ (see Section 4), Conjecture 7 is equivalent to the acyclicity of $`_{\mathrm{𝑖𝑒}}`$ in positive degrees. Another conjecture, Conjecture 22, is given in Section 4.
Let us finish this section with one exceptional example. The trivial operad $`\mathrm{𝟏}`$ is a Koszul quadratic self-dual operad. A $`\mathrm{𝟏}`$-algebra is a vector space $`A`$ with no operations. Clearly $`C_\mathrm{𝟏}^{}(A;A)`$ is just the space $`\mathrm{𝐿𝑖𝑛}(A,A)`$ of linear maps $`f:AA`$ concentrated in degree zero with trivial differential, thus $`H_\mathrm{𝟏}^{}(A;A)=\mathrm{𝐿𝑖𝑛}(A,A)`$. It is also clear that all natural operations on $`\mathrm{𝐿𝑖𝑛}(A,A)`$ are the identity $`11_A\mathrm{𝐿𝑖𝑛}(A,A)`$ considered as a degree zero constant, and iterated compositions
$$\mathrm{𝐿𝑖𝑛}(A,A)f_1,f_2,\mathrm{},f_nf_1f_2\mathrm{}f_n\mathrm{𝐿𝑖𝑛}(A,A),n1.$$
Therefore
$$_\mathrm{𝟏}U𝒜\mathrm{𝑠𝑠},$$
the operad for unital associative algebras. This example is pathological in that the canonical element introduced in Definition 8 equals zero. Therefore, from now on all quadratic Koszul operads in this note will be nontrivial in the sense that $`𝒫\mathrm{𝟏}`$.
## 2. The constants $`_𝒫(0)`$ – soul without body
This section, as well as the rest of the paper, relies on terminology and notation recalled in the Appendix. The main result of this part is Proposition 9 which describes the dg-vector space $`_𝒫(0)=(_𝒫(0),\delta _𝒫)`$ of “constants.” It is not difficult to see (compare also Example 35 of Section 7) that
$$_𝒫^{m1}(0)𝐬(𝒫(m)𝒫^!(m))^{\mathrm{\Sigma }_m},m1,$$
with the action $`_𝒫^{m1}(0)C_𝒫^{m1}(A;A)`$ given as the composition
$`𝐬(𝒫(m)𝒫^!(m))^{\mathrm{\Sigma }_m}\stackrel{}{}(𝐬𝒫(m)𝒫^!(m))^{\mathrm{\Sigma }_m}\stackrel{𝐬\alpha 11}{}(𝐬\mathrm{𝑛𝑑}_A(m)𝒫^!(m))^{\mathrm{\Sigma }_m}`$
$`(\mathrm{𝑛𝑑}_A(m)𝒫^!(m))^{\mathrm{\Sigma }_m}=[\mathrm{𝐿𝑖𝑛}((A)^m,A)𝒫^!(m)]^{\mathrm{\Sigma }_m}=C_𝒫^{m1}(A;A).`$
Since composition (2) is monic for all “generic” $`𝒫`$-algebras $`A`$, $`(_𝒫^{}(0),\delta _𝒫)`$ is “morally” the subcomplex of natural elements in $`(C_𝒫^{}(A;A),d_𝒫)`$. Before going further, we must recall the following general construction. Let $`𝒯`$ be an operad. It is well-known that the formula
$$[f,g]:=fg(1)^{(m1)(n1)}gf,$$
where $`fg`$ is, for $`f𝒯(m)`$ and $`g𝒯(n)`$, defined by
$$fg:=\underset{1im}{}(1)^{(n1)(i1)}f_ig,$$
makes the direct sum $`𝒯_{}=_{m0}𝒯_{}`$, with
$$𝒯_{m1}:=^{m1}𝒯(m)=𝐬𝒯(m),$$
a graded Lie algebra. Another standard fact is that each element $`\omega 𝒯_1=𝐬𝒯(2)`$ satisfying $`[\omega ,\omega ]=0`$ defines a degree $`+1`$ differential $`\delta _\omega :𝒯_{}𝒯_{+1}`$ by
$$\delta _\omega (t):=[t,\omega ],\text{ for }t𝒯_{}.$$
It is helpful to observe that the condition $`[\omega ,\omega ]=0`$ means the associativity:
(2)
$$\omega _1\omega =\omega _2\omega $$
and that the differential $`\delta _\omega `$ in terms of $`_i`$-operations equals
$$\delta _\omega (t)=t_1\omega t_2\omega +\mathrm{}(1)^mt_m\omega +(1)^m\omega _1t\omega _2t,\text{ for }t𝒯(m).$$
As proved in , the graded Lie algebra structure $`(𝒯_{},[,])`$ descents to the space of coinvariants therefore it induces, via the canonical isomorphism between invariants and coinvariants, a Lie bracket, denoted again $`[,]`$, on the graded vector space $`𝒯_{}^\mathrm{\Sigma }=_{m0}𝒯_m^\mathrm{\Sigma }`$ with pieces
$$𝒯_{m1}^\mathrm{\Sigma }:=^{m1}(𝒯(m)\mathrm{𝑠𝑔𝑛}_m)^{\mathrm{\Sigma }_m}=𝐬𝒯(m)^{\mathrm{\Sigma }_m}.$$
As usual, an element $`\varphi 𝒯_1^\mathrm{\Sigma }=𝐬𝒯(2)^{\mathrm{\Sigma }_2}`$ satisfying $`[\varphi ,\varphi ]=0`$ induces a degree $`+1`$ differential $`\delta _\varphi ^\mathrm{\Sigma }:𝒯_{}^\mathrm{\Sigma }𝒯_{+1}^\mathrm{\Sigma }`$ by
(3)
$$\delta _\varphi ^\mathrm{\Sigma }t:=[\varphi ,t],\text{ for }t𝒯_{}^\mathrm{\Sigma }.$$
In Proposition 9 below we put $`𝒯:=(𝒫𝒫^!)`$ and define the differential (3) by taking as $`\varphi `$ the canonical element $`\chi `$ introduced in the following definition in which $`\mathrm{\#}`$ denotes the linear dual.
###### Definition 8.
Let $`𝒫`$ be a quadratic Koszul operad. The canonical element $`\chi `$ is the element of $`𝐬(𝒫𝒫^!)(2)^{\mathrm{\Sigma }_2}`$ corresponding, under the standard identification
$$𝐬(𝒫𝒫^!)(2)(𝒫𝒫^\mathrm{\#})(2)(𝒫(2)𝒫(2)^\mathrm{\#})\mathrm{𝐿𝑖𝑛}(𝒫(2),𝒫(2)),$$
to the suspension of the identity map $`11_{𝒫(2)}\mathrm{𝐿𝑖𝑛}(𝒫(2),𝒫(2))`$.
Observe that $`\chi `$ is symmetric,
(4)
$$\chi \tau =\chi \text{ for }\tau \mathrm{\Sigma }_2\text{,}$$
therefore indeed $`\chi 𝐬(𝒫𝒫^!)(2)^{\mathrm{\Sigma }_2}`$. The condition $`[\chi ,\chi ]=0`$ is equivalent to the Jacobi identity (17) for $`\chi `$ which follows from \[13, Corollary 2.2.9(b)\], see also the proof of Proposition 26.
###### Proposition 9.
There is a natural isomorphism of cochain complexes
$$(_𝒫^{}(0),\delta _𝒫)((𝒫𝒫^!)_{}^\mathrm{\Sigma },\delta _\chi ^\mathrm{\Sigma }).$$
If $`𝒫`$ is the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$ \[26, Remark II.1.15\], then there is a similar description of the chain complex $`(_{\underset{¯}{𝒫}}^{}(0),\delta _{\underset{¯}{𝒫}})`$ obtained as follows. The definition of the graded Lie algebra $`(𝒯_{},[,])`$ given above clearly makes sense also when $`𝒯`$ is a non-$`\mathrm{\Sigma }`$ operad. Observe also that there exists the non-$`\mathrm{\Sigma }`$ quadratic dual $`\underset{¯}{𝒫}^!`$ of $`\underset{¯}{𝒫}`$ and that one may introduce the non-$`\mathrm{\Sigma }`$ canonical element $`\underset{¯}{\chi }𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(2)`$ in exactly the same manner as its symmetric version. The element $`\underset{¯}{\chi }`$ obviously satisfies the associativity condition (2):
$$\underset{¯}{\chi }_1\underset{¯}{\chi }=\underset{¯}{\chi }_2\underset{¯}{\chi }.$$
Our non-$`\mathrm{\Sigma }`$ version of Proposition 9 reads:
###### Proposition 10.
Let $`𝒫`$ be the symmetrization of a quadratic Koszul non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$. Then
$$(_{\underset{¯}{𝒫}}^{}(0),\delta _{\underset{¯}{𝒫}})((\underset{¯}{𝒫}\underset{¯}{𝒫}^!)_{},\delta _{\underset{¯}{\chi }}).$$
Let us make a comment on the meaning of the cohomology $`H^{}(_𝒫(0),\delta _𝒫)`$. The natural morphism
$$M:H^{}(_𝒫(0),\delta _𝒫)H_𝒫^{}(A;A)$$
induced by action (2) is monic for any “generic” $`𝒫`$-algebra $`A`$, therefore elements $`H^{}(_𝒫(0),\delta _𝒫)`$ represent natural generically nontrivial homology classes in the cohomology of $`𝒫`$-algebras. This leads one to believe that $`H^{}(_𝒫(0),\delta _𝒫)=0`$ for all well-behaved operads, since otherwise people would stumble upon nontrivial natural classes – compare the Casimir element in the cohomology of simple Lie algebras. Example 15 however contradict this reasonable assumption. We believe that $`H^{}(_𝒫(0),\delta _𝒫)`$ is an important invariant of the operad $`𝒫`$ that deserves its own name:
###### Definition 11.
We call the graded vector space $`H^{}(_𝒫(0),\delta _𝒫)`$ described in Proposition 9 the soul of the cohomology of $`𝒫`$-algebras.
It is easy to prove that $`H^0(_𝒫(0),\delta _𝒫)`$ is always trivial.
###### Example 12.
Let us describe the complex calculating the soul $`H^{}(_{𝒜\mathrm{𝑠𝑠}}(0),\delta _{𝒜\mathrm{𝑠𝑠}})`$ of the Hochschild cohomology. Clearly
$$(𝒫𝒫^!)_{m1}^\mathrm{\Sigma }=(𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})_{m1}^\mathrm{\Sigma }𝐬(𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})(m)^{\mathrm{\Sigma }_m}𝐬𝒜\mathrm{𝑠𝑠}(m),$$
therefore the complex $`((𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})^\mathrm{\Sigma },\delta _\chi ^\mathrm{\Sigma })`$ has the form
(5)
$$𝐤\stackrel{\delta _\chi ^\mathrm{\Sigma }}{}𝐤[\mathrm{\Sigma }_2]\stackrel{\delta _\chi ^\mathrm{\Sigma }}{}𝐤[\mathrm{\Sigma }_3]\stackrel{\delta _\chi ^\mathrm{\Sigma }}{}𝐤[\mathrm{\Sigma }_4]\stackrel{\delta _\chi ^\mathrm{\Sigma }}{}\mathrm{}$$
Is also easy to describe the differential $`\delta _\chi ^\mathrm{\Sigma }`$; on a permutation $`\sigma \mathrm{\Sigma }_m`$ it acts as
$$\delta _\chi ^\mathrm{\Sigma }(\sigma ):=d_0(\sigma )d_1(\sigma )+d_2(\sigma )\mathrm{}+(1)^{m+1}d_{m+1}(\sigma )𝐤[\mathrm{\Sigma }_{m+1}],$$
where $`d_0(\sigma ):=11\times \sigma `$, $`d_{m+1}(\sigma ):=\sigma \times 11`$ and $`d_i(\sigma )\mathrm{\Sigma }_{m+1}`$ is the permutation obtained by doubling the $`i`$th input of $`\sigma `$. In Theorem 13 below we prove that (5) is acyclic.
Since $`𝒜\mathrm{𝑠𝑠}`$ is the symmetrization of the non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒜\mathrm{𝑠𝑠}}`$, it makes sense to consider also the subcomplex $`(_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}^{}(0),\delta _{\underset{¯}{𝒜\mathrm{𝑠𝑠}}})`$ of $`(_{𝒜\mathrm{𝑠𝑠}}^{}(0),\delta _{𝒜\mathrm{𝑠𝑠}})`$ described in Proposition 10. This subcomplex is obviously isomorphic to the acyclic complex
(6)
$$𝐤\stackrel{d_0}{}𝐤\stackrel{d_1}{}𝐤\stackrel{d_2}{}𝐤\stackrel{d_3}{}\mathrm{}$$
in which $`d_{2i}=11_𝐤`$ and $`d_{2i+1}=0`$, $`i0`$. The inclusion $`(_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}^{}(0),\delta _{\underset{¯}{𝒜\mathrm{𝑠𝑠}}})(_{𝒜\mathrm{𝑠𝑠}}^{}(0),\delta _{𝒜\mathrm{𝑠𝑠}})`$ sends the generator $`1𝐤`$ of the $`n`$th piece of (6) into the identity permutation $`11_{n1}𝐤[\mathrm{\Sigma }_{n1}]`$ in complex (5).
###### Theorem 13.
The soul $`H^{}(_{𝒜\mathrm{𝑠𝑠}}(0),\delta _{𝒜\mathrm{𝑠𝑠}})`$ of the Hochschild cohomology is trivial.
Proof. We must prove that (5) is an acyclic complex. The idea will be to show that it decomposes into a direct sum of acyclic subcomplexes indexed by primitive, in the sense introduced below, permutations.
We define first, for each $`\sigma \mathrm{\Sigma }_n`$, a natural number $`g(\sigma )`$, $`0g(\sigma )n`$, which we call the grade of $`\sigma `$, as follows. The grade of the unit $`11_n\mathrm{\Sigma }_n`$ is $`n1`$, $`g(11_n):=n1`$. For $`\sigma 11_n`$, let
$`a(\sigma )`$ $`:=`$ $`\mathrm{max}\{i;\sigma =11_i\times \tau \text{ for some }\tau \mathrm{\Sigma }_{ni}\},\text{ and}`$
$`c(\sigma )`$ $`:=`$ $`\mathrm{max}\{j;\sigma =\lambda \times 11_j\text{ for some }\lambda \mathrm{\Sigma }_{ni}\}.`$
There clearly exists a unique $`\omega =\omega (\sigma )\mathrm{\Sigma }_{na(\sigma )c(\sigma )}`$ such that $`\sigma =11_{a(\sigma )}\times \omega (\sigma )\times 11_{c(\sigma )}`$. Let, finally, $`b(\sigma )`$ be the number of “doubled strings” in $`\omega (\sigma )`$,
$$b(\sigma ):=\{1s<k;\omega (s+1)=\omega (s)+1\}.$$
The grade of $`\sigma `$ is then defined by
$$g(\sigma ):=a(\sigma )+b(\sigma )+c(\sigma ),$$
see Figure 1 for examples. Observe that the differential $`\delta _\chi ^\mathrm{\Sigma }`$ of (5) raises the grade by $`+1`$.
Let us call $`\chi \mathrm{\Sigma }_k`$, $`k1`$, primitive if $`g(\sigma )=0`$. Observe that, according to our definitions, $`11_n\mathrm{\Sigma }_n`$ is primitive if and only if $`n=1`$. For each $`\sigma \mathrm{\Sigma }_n`$, $`\sigma 11_n`$, we define a unique primitive $`\kappa =\kappa (\sigma )\mathrm{\Sigma }_k`$, $`k=ng(\sigma )`$, by contracting all “multiple strings” of $`\omega (\sigma )`$ into “simple” ones, see Figure 1. We put $`\chi (11_n):=11_1`$.
For a primitive $`\chi `$, let $`P^{}(\chi )`$ be the graded subspace of (5) spanned by all permutations $`\sigma `$ with $`\chi =\kappa (\sigma )`$. The following statements can be easily verified:
(i) Each $`P^{}(\chi )`$ is a subcomplex of (5).
(ii) Complex (5) decomposes as the summation $`_\chi P^{}(\chi )`$ over all primitive permutations $`\chi `$.
(iii) For each primitive $`\chi \mathrm{\Sigma }_n`$,
$$P^{}(\chi )P^{}(11_1)\mathrm{}P^{}(11_1)\text{ (}n+2\text{ times).}$$
The proof is finished by observing that $`P^{}(11_1)`$ is isomorphic to the acyclic complex (6) and applying the Künneth formula. mm
###### Example 14.
In this example we describe the soul of the Chevalley-Eilenberg cohomology of Lie algebras which is, due to the obvious self-duality of Proposition 9, the same as the soul of the Harrison cohomology of commutative associative algebras. In both cases
$$(𝒫𝒫^!)_{m1}^\mathrm{\Sigma }=(𝒞\mathrm{𝑜𝑚}\mathrm{𝑖𝑒})_{m1}^\mathrm{\Sigma }𝐬\mathrm{𝑖𝑒}(m)^{\mathrm{\Sigma }_m}=^{m1}(\mathrm{𝑖𝑒}(m)\mathrm{𝑠𝑔𝑛}_m)^{\mathrm{\Sigma }_m}𝐤$$
(see ) and one may identify $`((𝒞\mathrm{𝑜𝑚}\mathrm{𝑖𝑒})^\mathrm{\Sigma },\delta _\chi ^\mathrm{\Sigma })`$ with the acyclic complex (6). Therefore the souls of both the Chevalley-Eilenberg cohomology and the Harrison cohomology are trivial.
###### Example 15.
This example presents a Koszul quadratic operad with a nontrivial soul. Let $`𝒟:=𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠}`$ be the free product of two copies of the associative operad. Operad $`𝒟`$ is a Koszul quadratic operad, whose quadratic dual $`𝒟^!`$ equals the coproduct $`𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠}`$ defined by
$$(𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})(m):=\{\begin{array}{cc}𝐤,\hfill & \text{if }m=1\text{ and}\hfill \\ 𝒜\mathrm{𝑠𝑠}(m)𝒜\mathrm{𝑠𝑠}(m),\hfill & \text{if }m2\text{.}\hfill \end{array}$$
Obviously, $`𝒟`$-algebras are triples $`A=(A,\mu _1,\mu _2)`$ consisting of two independent associative multiplications $`\mu _1,\mu _2:AAA`$. The cohomology of these algebras is the cohomology of the total complex $`(C_𝒟^{}(A;A),d_𝒟)`$ of the “meager” bicomplex in Figure 2. The horizontal line is the Hochschild cochain complex of the associative algebra $`A_1:=(A,\mu _1)`$, the vertical line the Hochschild complex of $`A_2:=(A,\mu _2)`$.
Let $`e`$ denote the identity $`11\mathrm{𝐿𝑖𝑛}(A,A)`$ considered as a natural element of $`C_𝒟^0(A;A)`$. Clearly
$$d_𝒟(d_1e)=d_𝒟(d_2e)=0$$
therefore both $`d_1e`$ and $`d_2e`$ are natural cochains in $`C_𝒟^1(A;A)`$ thus representing $`\delta _𝒟`$-cochains in $`_𝒟^1(0)`$. The equality
$$d_𝒟e=d_1e+d_2e$$
implies that $`d_1e+d_2e`$ is $`\delta _𝒟`$-homologous to zero in $`_𝒟^1(0)`$. We conclude that
$$H^1(_𝒟(0),\delta _𝒟)\mathrm{𝑆𝑝𝑎𝑛}([d_1e])\mathrm{𝑆𝑝𝑎𝑛}([d_2e])𝐤.$$
We saw that the souls of the Hochschild ($`𝒫=𝒜\mathrm{𝑠𝑠})`$, Chevalley-Eilenberg ($`𝒫=\mathrm{𝑖𝑒}`$) and Harrison ($`𝒫=𝒜\mathrm{𝑠𝑠}`$) cohomologies were trivial, while the soul of the cohomology for $`𝒟`$-algebras analyzed in Example 15 was not. This leads us to formulate:
###### Problem 16.
Which property of a quadratic Koszul operad $`𝒫`$ implies the triviality of the soul $`H^{}((𝒫𝒫^!)^\mathrm{\Sigma },\delta _\chi ^\mathrm{\Sigma })`$ of the $`𝒫`$-cohomology?
###### Example 17.
In this example we describe a non-$`\mathrm{\Sigma }`$ quadratic Koszul operad with the property that $`(_{\underset{¯}{𝒫}}^{}(0),\delta _{\underset{¯}{𝒫}})`$ is acyclic but the soul $`(_𝒫^{}(0),\delta _𝒫)`$ is not. This shows that the homotopy type of $`_𝒫`$ is in general different from the homotopy type of $`_{\underset{¯}{𝒫}}`$.
Let $`\underset{¯}{a}g`$ be the free non-$`\mathrm{\Sigma }`$ operad generated by one bilinear operation, $`\underset{¯}{a}g:=\underset{¯}{\mathrm{\Gamma }}(\underset{¯}{\mu })`$, and $`\mathrm{𝑎𝑔}`$ its symmetrization. The corresponding cochain complex $`C_{\mathrm{𝑎𝑔}}^{}(A;A)`$ is the truncation
$$\mathrm{𝐿𝑖𝑛}(A,A)\stackrel{d}{}\mathrm{𝐿𝑖𝑛}(A^2,A)$$
of the Hochschild complex. The complex $`(_{\mathrm{𝑎𝑔}}^{}(0),\delta _{\mathrm{𝑎𝑔}})`$ defining the soul of $`\mathrm{𝑎𝑔}`$ is the truncation
$$𝐤\stackrel{\delta _\chi ^\mathrm{\Sigma }}{}𝐤[\mathrm{\Sigma }_2]$$
of (5), and is manifestly not acyclic. On the other hand, $`(_{\underset{¯}{a}g}^{}(0),\delta _{\underset{¯}{a}g})`$ is acyclic, isomorphic to the truncation $`𝐤\stackrel{d_0}{}𝐤`$ of (6). We conclude that $`H^{}(_{\underset{¯}{a}g}^{}(0),\delta _{\underset{¯}{a}g})=0`$ while
$$H^{}(_{\mathrm{𝑎𝑔}}^{}(0),\delta _{\mathrm{𝑎𝑔}})=H^1(_{\mathrm{𝑎𝑔}}^{}(0),\delta _{\mathrm{𝑎𝑔}})𝐤.$$
## 3. Homotopy type of $`(1)`$ – surprises continue
In this section we study, as a next step toward the understanding of $`_𝒫`$, the homotopy type of the associative dg-algebra $`_𝒫(1)=(_𝒫^{}(1),\delta _𝒫)`$. Since the operad $`𝒫^!`$ is a module, in the sense of , over itself, it makes sense to consider the space $`\mathrm{𝐸𝑛𝑑}_{𝒫^!}(𝒫^!)`$ of all $`𝒫^!`$-module endomorphisms $`\alpha :𝒫^!𝒫^!`$. Very crucially,
(7)
$$\mathrm{𝐸𝑛𝑑}_{𝒫^!}(𝒫^!)𝐤,$$
because each $`\alpha \mathrm{𝐸𝑛𝑑}_{𝒫^!}(𝒫^!)`$ is uniquely determined by the value $`\alpha _1(1)𝒫^!(1)𝐤`$ and, conversely, for each $`\phi 𝐤`$ the homomorphism $`\alpha :=\phi 11_{𝒫^!}`$ is such that $`\alpha _1(1)=\phi `$.
###### Proposition 18.
There is a canonical identification of associative unital algebras
(8)
$$H^0(_𝒫^{}(1),\delta _𝒫)\mathrm{𝐸𝑛𝑑}_{𝒫^!}(𝒫^!)𝐤.$$
Proof. Since there are no elements in negative degrees,
$$H^0(_𝒫(1))=\mathrm{𝐾𝑒𝑟}(\delta :_𝒫^0(1)_𝒫^1(1)).$$
By definition, elements of the kernel $`\mathrm{𝐾𝑒𝑟}(\delta )`$ are natural chain maps
$$\{\phi _m:C_𝒫^m(A;A)C_𝒫^m(A;A)\}_{m0}.$$
As explained in Example 36, the naturality of $`\phi _m`$ means that it is induced by a $`\mathrm{\Sigma }_{m+1}`$-equivariant map $`\alpha _{m+1}:𝒫^!(m+1)𝒫^!(m+1)`$ . It is easy to verify that the collection $`\{\alpha _m\}`$ determines a chain map if and only if it assembles into a $`𝒫^!`$-module endomorphism $`\alpha :𝒫^!𝒫^!`$. mm
The following example shows that the dg-associative algebra $`_𝒫(1)`$ might in general have nontrivial cohomology in positive degrees.
###### Example 19.
Let $`𝒮\mathrm{𝑦𝑚}`$ be the operad describing algebras with one commutative bilinear multiplication and no axioms. Explicitly, $`𝒮\mathrm{𝑦𝑚}`$ is the free operad generated by the trivial representation of $`\mathrm{\Sigma }_2`$ placed in arity two. It is a Koszul quadratic operad whose quadratic dual $`𝒮\mathrm{𝑦𝑚}^!`$ is given by $`𝒮\mathrm{𝑦𝑚}^!(1)=𝐤`$, $`𝒮\mathrm{𝑦𝑚}^!(2)=\mathrm{𝑠𝑔𝑛}_2`$ (the signum representation of $`\mathrm{\Sigma }_2`$), and $`𝒮\mathrm{𝑦𝑚}^!(m)=0`$ for $`m3`$.
The cohomology of a $`𝒮\mathrm{𝑦𝑚}`$-algebra $`A=(A,)`$ is the cohomology of the two-term complex (which should be interpreted as a truncation of the Harrison complex)
$$\mathrm{𝐿𝑖𝑛}(A,A)\stackrel{d}{}\mathrm{𝐿𝑖𝑛}(S^2A,A),$$
where $`S^2A`$ is the second symmetric power of $`A`$. The differential $`d`$ is given by the usual formula
$$(d\varphi )(a,b):=a\varphi (b)\varphi (ab)+\varphi (b)a,$$
for $`\varphi \mathrm{𝐿𝑖𝑛}(A,A)`$ and $`a,bA`$.
We are going to describe the dg-algebra $`_{𝒮\mathrm{𝑦𝑚}}^{}(1)`$. Let $`\alpha `$ be the projection of $`\mathrm{𝐿𝑖𝑛}(A,A)\mathrm{𝐿𝑖𝑛}(S^2A,A)`$ onto $`\mathrm{𝐿𝑖𝑛}(A,A)`$ and $`\beta `$ the projection onto $`\mathrm{𝐿𝑖𝑛}(S^2A,A)`$. Let $`u`$ and $`v`$ be degree $`+1`$ operations given by
$$u(\varphi )(a,b):=a\varphi (b)+\varphi (a)b\text{ and }v(\varphi )(a,b):=\varphi (ab),$$
for $`\varphi \mathrm{𝐿𝑖𝑛}(A,A)`$ and $`a,bA`$. Then clearly $`_{𝒮\mathrm{𝑦𝑚}}^0(1)`$ is the semisimple algebra $`𝐤𝐤`$ spanned by $`\alpha `$ and $`\beta `$, and the space $`_{𝒮\mathrm{𝑦𝑚}}^1(1)`$ is two-dimensional, spanned by $`u`$ and $`v`$. The higher $`_{𝒮\mathrm{𝑦𝑚}}^i(1)`$ are, for $`i2`$, trivial. To describe the multiplication in $`_{𝒮\mathrm{𝑦𝑚}}^{}(1)`$, it is enough to specify how $`_{𝒮\mathrm{𝑦𝑚}}^0(1)`$ acts on $`_{𝒮\mathrm{𝑦𝑚}}^1(1)`$. This action is given by
$$\alpha b=0=b\beta \text{ and }b\alpha =b=\beta b,\text{ for }b_{𝒮\mathrm{𝑦𝑚}}^1(1)\text{.}$$
The differential $`\delta _{𝒮\mathrm{𝑦𝑚}}`$ of $`_{𝒮\mathrm{𝑦𝑚}}^{}(1)`$ acts by
$$\delta \alpha =\delta \beta =uv,\delta u=\delta v=0.$$
The cohomology of $`(_{𝒮\mathrm{𝑦𝑚}}^{}(1),\delta _{𝒮\mathrm{𝑦𝑚}})`$ can be easily calculated,
$$H^{}(_{𝒮\mathrm{𝑦𝑚}}^{}(1))𝐤W,$$
where $`W`$ is the vector space spanned by the class $`[u]`$. We leave as a simple exercise to prove that there exist a quasi-isomorphism $`H^{}(_{𝒮\mathrm{𝑦𝑚}}^{}(1))_{𝒮\mathrm{𝑦𝑚}}^{}(1)`$. The dg-associative algebra $`_{𝒮\mathrm{𝑦𝑚}}(1)`$ is therefore formal.
Here is a baby version of Problem 1:
###### Problem 20.
Describe the homotopy type of the unital differential graded associative algebra $`_𝒫(1)=(_𝒫^{}(1),\delta _𝒫)`$. In particular, calculate the cohomology of $`_𝒫(1)`$.
We expect that the homotopy type of $`_𝒫(1)`$ is that of $`𝐤`$ for all “reasonable” operads, though we do not know what “reasonable” means – the operad $`𝒮\mathrm{𝑦𝑚}`$ of Example 19 seems reasonable enough, yet the homotopy type of $`_{𝒮\mathrm{𝑦𝑚}}(1)`$ is nontrivial. Let us close this section by formulating:
###### Problem 21.
Which property of the operad $`𝒫`$ implies that the dg associative algebra $`(_𝒫^{}(1),\delta _𝒫)`$ has the homotopy type of $`𝐤`$?
## 4. The operad $`H^0(_𝒫)`$ and the intrinsic bracket
It is well-known that the chain complex $`C_𝒫^{}(A;A)`$ always carries a natural dg Lie algebra structure given by the intrinsic bracket. The easiest way to construct such a bracket is to identify $`C_𝒫^{}(A;A)`$ with the dg Lie algebra $`\mathrm{𝐶𝑜𝑑𝑒𝑟}^{}(_{𝒫^!}^c(A))`$ of coderivations of the cofree nilpotent $`𝒫^!`$-coalgebra cogenerated by the desuspension $`A`$ as it was done in \[26, Definition II.3.99\]. In this way we obtain a natural homomorphism
(9)
$$I:(\mathrm{𝑖𝑒},0)(_𝒫,\delta _𝒫)$$
of dg operads. If $`𝒫`$ is the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$, then $`\mathrm{𝐼𝑚}(I)_{\underset{¯}{𝒫}}`$, therefore the map $`I`$ of (9) factorizes through the natural inclusion $`(_{\underset{¯}{𝒫}},\delta _{\underset{¯}{𝒫}})(_𝒫,\delta _𝒫)`$. Computational evidences lead us to:
###### Conjecture 22.
The natural homomorphism $`I:(\mathrm{𝑖𝑒},0)(_𝒫,\delta _𝒫)`$ induces an isomorphism of operads
$$H^0(_𝒫)\mathrm{𝑖𝑒},$$
for an arbitrary nontrivial quadratic Koszul $`𝒫`$.
We were able to verify Conjecture 22 for $`𝒫=\mathrm{𝑖𝑒}`$, that is, to prove
(10)
$$H^0(_{\mathrm{𝑖𝑒}})\mathrm{𝑖𝑒}.$$
This isomorphism turned out to be related to a certain characterization of free Lie algebras inside free pre-Lie algebras. More precisely, let $`\mathrm{𝑝𝑟𝑒}𝕃(X)`$ denote the free pre-Lie algebra generated by a set $`X`$. The commutator of the pre-Lie product makes $`\mathrm{𝑝𝑟𝑒}𝕃(X)`$ a Lie algebra. Let $`𝕃(X)\mathrm{𝑝𝑟𝑒}𝕃(X)`$ be the Lie algebra generated by $`X`$ in $`\mathrm{𝑝𝑟𝑒}𝕃(X)`$. It is not hard to see that $`𝕃(X)`$ is in fact isomorphic to the free Lie algebra on $`X`$, see also . Then (10) is implied by a very explicit characterization of the subspace $`𝕃(X)`$ of $`\mathrm{𝑝𝑟𝑒}𝕃(X)`$.
Similarly, the conjectured isomorphism $`H^0(_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}})\mathrm{𝑖𝑒}`$ can be translated into a certain characterization of free Lie algebras inside free brace algebras. We were also able to prove that, for an arbitrary quadratic Koszul operad
(11)
$$H^0(_𝒫(2))\mathrm{𝑠𝑔𝑛}_2,$$
the signum representation of $`\mathrm{\Sigma }_2`$, by describing $`H^0(_𝒫(2))`$ in terms of suitably defined pairings $`𝒫^!𝒫^!𝒫^!`$.
Let us close this section by a couple of remarks which will be useful in the proof of Proposition 26. As we recalled at the beginning of this section, there is a canonical isomorphism $`C_𝒫^{}(A;A)\mathrm{𝐶𝑜𝑑𝑒𝑟}^{}(_{𝒫^!}^c(A))`$. It is well-known that coderivations of a cofree nilpotent algebra form a natural pre-Lie algebra \[26, II.3.9\], therefore one has a natural homomorphism of non-dg operads
(12)
$$\mathrm{𝑝𝑟𝑒}I:\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}_𝒫.$$
The map (9) is then the composition
$$\mathrm{𝑖𝑒}\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}\stackrel{\mathrm{𝑝𝑟𝑒}I}{}_𝒫$$
of $`\mathrm{𝑝𝑟𝑒}I`$ with the anti-symmetrization map $`\mathrm{𝑖𝑒}\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}`$.
## 5. The cup products
The central statement of this section is Theorem 23 that claims that the suspension $`𝐬(𝒫𝒫^!)`$ (see A.3) of the operad $`(𝒫𝒫^!)`$ acts on $`C_𝒫^{}(A;A)`$, and Theorem 24 that characterizes which elements of $`𝐬(𝒫𝒫^!)`$ act via chain maps. Observe that the operad $`𝐬(𝒫𝒫^!)`$ need not be quadratic even when $`𝒫`$ is.
###### Theorem 23.
There is a canonical action of the operad $`𝐬(𝒫𝒫^!)`$ on the graded vector space $`C_𝒫^{}(A;A)`$, via natural operations. This action can be interpreted as an inclusion of non-differential graded operads
(13)
$$\mathrm{𝑐𝑢𝑝}:𝐬(𝒫𝒫^!)_𝒫.$$
Proof. The proof relies on the notation introduced/recalled in A.6 and A.5. The “tautological” action of the endomorphism operad $`\mathrm{𝑛𝑑}_A`$ on $`A`$ tensored with the action of $`𝒫^!`$ on itself makes the graded vector space $`\stackrel{~}{C}_𝒫^{}(A;A)=_{m0}\stackrel{~}{C}_𝒫^m(A;A)`$ an $`𝐬(\mathrm{𝑛𝑑}_A𝒫^!)`$-algebra. It is straightforward to prove that this action induces, via
(14)
$$t(f_1,\mathrm{},f_n):=\mathrm{𝐴𝑣𝑒𝑟}\left(\text{}t(\iota (f_1),\mathrm{},\iota (f_n))\right),$$
for $`t𝐬(\mathrm{𝑛𝑑}_A𝒫^!)(n)`$ and $`f_1,\mathrm{},f_nC_𝒫^{}(A;A)`$, an action of $`𝐬(\mathrm{𝑛𝑑}_A𝒫^!)`$ on the graded vector space $`C_𝒫^{}(A;A)`$. Suppose that $`A`$ is a $`𝒫`$-algebra, with the structure given by $`\alpha :𝒫\mathrm{𝑛𝑑}_A`$. Action (13) is obtained by composing action (14) with the homomorphism $`𝐬(\alpha 11):𝐬(𝒫𝒫^!)𝐬(\mathrm{𝑛𝑑}_A𝒫^!)`$. An alternative description of (13) is given in Example 38 of Section 7.mm
We use inclusion (13) to view $`𝐬(𝒫𝒫^!)`$ as a suboperad of $`_𝒫`$. Elements of $`𝐬(𝒫𝒫^!)`$ need not be $`\delta _𝒫`$-closed in $`_𝒫`$; let $`𝒵_𝒫𝐬(𝒫𝒫^!)`$ denote the suboperad of $`\delta _𝒫`$-cocycles. In Proposition 24, which describes $`𝒵_𝒫`$ explicitly, we use the canonical element $`\chi `$ introduced in Definition 8.
###### Theorem 24.
The suboperad $`𝒵_𝒫`$ of $`\delta _𝒫`$-closed elements in $`𝐬(𝒫𝒫^!)`$ is characterized as follows: $`t𝐬(𝒫𝒫^!)(n)`$ belongs to $`𝒵_𝒫(n)`$ if and only if
(15)
$$\chi _2t+t_1\chi +(t_2\chi )(12)+(t_3\chi )(123)+\mathrm{}+(t_n\chi )(123\mathrm{}n)=0,$$
where $`(123\mathrm{}k)\mathrm{\Sigma }_{n+1}`$ is the cycle
$$\left(\begin{array}{cccccccc}1& 2& 3& \mathrm{}& k& k+1& \mathrm{}& n+1\\ 2& 3& 4& \mathrm{}& 1& k+1& \mathrm{}& n+1\end{array}\right).$$
The proof is a completely straightforward calculation. We recommend as an exercise to verify that solutions of (15) are indeed closed under operadic composition. The meaning of equation (15) should be clear from Figure 3. The importance of the operad $`𝒵_𝒫`$ is explained by
###### Corollary 25.
The map $`\mathrm{𝑐𝑢𝑝}`$ of (13) induces a canonical map (denoted again $`\mathrm{𝑐𝑢𝑝}`$)
(16)
$$\mathrm{𝑐𝑢𝑝}:𝒵_𝒫H^{}(_𝒫,\delta _𝒫),$$
therefore $`H_𝒫^{}(A;A)`$ is a natural $`𝒵_𝒫`$-algebra.
From reasons which become clear later we call operations induced by elements of $`𝒵_𝒫`$ the cup products. The following proposition in which $`\mathrm{𝑖𝑒}`$ is the operad for Lie algebras (see A.2) shows that the operad $`𝒵_𝒫`$ is always nontrivial (provided $`𝒫\mathrm{𝟏}`$) while the map (16) is never monic.
###### Proposition 26.
The operad $`𝒵_𝒫`$ contains the canonical element $`\chi `$. There exists a unique map $`L:𝐬\mathrm{𝑖𝑒}𝒵_𝒫`$ that sends the generator $`𝐬\lambda 𝐬\mathrm{𝑖𝑒}(2)`$ into $`\chi 𝒵_𝒫(2)`$. All elements in the image of this map are $`\delta _𝒫`$-cohomologous to zero in $`_𝒫`$.
Proof. Recall \[13, Corollary 2.2.9(b)\] that, for each quadratic operad $`𝒫`$, there exits a morphism of operads $`\mathrm{𝑖𝑒}𝒫𝒫^!`$ that takes the generator $`\lambda \mathrm{𝑖𝑒}(2)`$ into the identity operator in $`𝒫(2)𝒫(2)^\mathrm{\#}𝒫(2)𝒫^!(2)`$. Let $`L:𝐬\mathrm{𝑖𝑒}𝐬(𝒫𝒫^!)`$ be the suspension of this morphism. Let us prove, using Theorem 24, that $`\chi 𝒵_𝒫(2)`$. Equation (15) for $`t=\chi `$ reads
$$\chi _2\chi +\chi _1\chi +(\chi _2\chi )(12)=0,$$
which can be written, due to the symmetry (4) of $`\chi `$, as the Jacobi identity for a degree $`1`$ “multiplication” $`\chi `$:
(17)
$$\chi _1\chi +(\chi _1\chi )(123)+(\chi _1\chi )(132)=0,$$
or, pictorially,
But (17) is satisfied, because $`\chi =L(𝐬\lambda )`$ by definition, and $`𝐬\lambda 𝐬\mathrm{𝑖𝑒}(2)`$ satisfies the same condition in $`𝐬\mathrm{𝑖𝑒}`$. The inclusion $`\mathrm{𝐼𝑚}(L)𝒵_𝒫`$ follows from the fact that $`\mathrm{𝐼𝑚}(L)`$ is generated by $`\chi `$ and that $`𝒵_𝒫`$ is a suboperad of $`𝐬(𝒫𝒫^!)`$.
Let us prove that all elements in the image of $`L`$ are $`\delta _𝒫`$-cohomologous to zero. Let $`\mathrm{}\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}(2)`$ be the generator of the quadratic operad $`\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}`$ for pre-Lie algebras and let $`:=\mathrm{𝑝𝑟𝑒}I(\mathrm{})_𝒫^0(2)`$, where $`\mathrm{𝑝𝑟𝑒𝐼}:\mathrm{𝑝𝑟𝑒}\mathrm{𝑖𝑒}_𝒫`$ is the map considered in (12) at the end of Section 4. It is easy to verify that then $`\chi =\delta _𝒫()`$. This finishes the proof of Proposition 26, because $`\mathrm{𝐼𝑚}(L)`$ is generated by $`\chi `$.mm
Suppose that $`𝒫`$ is the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$. Given $`t𝐬(𝒫𝒫^!)(n)`$ as in Theorem 24, $`\mathrm{𝑐𝑢𝑝}(t)_{\underset{¯}{𝒫}}(n)`$ if and only if $`t`$ belongs to the $`\mathrm{\Sigma }_n`$-closure of $`𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(n)`$ in $`𝐬(𝒫𝒫^!)(n)`$, that is, if $`t=\underset{¯}{t}\sigma `$ for some $`\underset{¯}{t}𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(n)`$ and $`\sigma \mathrm{\Sigma }_n`$. In the following non-$`\mathrm{\Sigma }`$ version of Theorem 24, $`\underset{¯}{\chi }𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(2)`$ is the non-$`\mathrm{\Sigma }`$ canonical element introduced in Section 2.
###### Theorem 27.
An element $`\underset{¯}{t}𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(n)𝐬(𝒫𝒫^!)(n)`$ belongs to $`𝒵_𝒫(n)`$ if and only if
(18)
$$\underset{¯}{\chi }_2\underset{¯}{t}=\underset{¯}{t}_1\underset{¯}{\chi }=\underset{¯}{t}_2\underset{¯}{\chi }=\mathrm{}=\underset{¯}{t}_n\underset{¯}{\chi }=\underset{¯}{\chi }_1\underset{¯}{t},$$
see Figure 4.
Proof of Theorem 27 is a straightforward verification. The proof of the following proposition is similar to that of Proposition 26.
###### Proposition 28.
Let $`𝒫`$ be the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$. Then $`\underset{¯}{\chi }𝒵_𝒫`$ and there exists a unique map $`A:𝐬𝒜\mathrm{𝑠𝑠}𝒵_𝒫`$ defined by $`A(𝐬\mu ):=\underset{¯}{\chi }`$, where $`𝐬\mu 𝐬𝒜\mathrm{𝑠𝑠}(2)`$ is the suspension of the generator $`\mu `$ (see A.2). Moreover, the diagram
where $`L`$ is as in Proposition 26, with the vertical map given by the anti-commutator of the associative product, commutes.
###### Example 29.
– Hochschild cohomology. Let $`𝒫=𝒜\mathrm{𝑠𝑠}`$ be the operad for associative algebras. Then $`𝐬(𝒫𝒫^!)=𝐬(𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})`$ and a simple calculation reveals that the map $`A`$ of Proposition 28 is the suspended diagonal $`𝐬\mathrm{\Delta }:𝐬𝒜\mathrm{𝑠𝑠}𝐬(𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠})`$ and that $`𝒵_{𝒜\mathrm{𝑠𝑠}}=\mathrm{𝐼𝑚}(A)`$. Therefore
$$𝒵_{𝒜\mathrm{𝑠𝑠}}𝐬𝒜\mathrm{𝑠𝑠}.$$
The generator $`𝐬\mu 𝐬𝒜\mathrm{𝑠𝑠}(2)`$ is mapped to the “classical” cup product $`f,gfg`$ of Hochschild cochains , and the generator $`𝐬\lambda 𝐬\mathrm{𝑖𝑒}(2)`$ to the anti-commutator of this cup product:
$$f,gfg+(1)^{|f||g|}gf$$
which cohomologous to zero, because the cup product of Hochschild cochains is homotopy commutative \[9, Theorem 3\].
###### Example 30.
– Chevalley-Eilenberg cohomology. If $`𝒫=\mathrm{𝑖𝑒}`$ is the operad for Lie algebras, then $`𝐬(𝒫𝒫^!)=𝐬(\mathrm{𝑖𝑒}𝒞\mathrm{𝑜𝑚})𝐬\mathrm{𝑖𝑒}`$ and we see immediately that
(19)
$$𝒵_{\mathrm{𝑖𝑒}}𝐬\mathrm{𝑖𝑒}=\mathrm{𝐼𝑚}(L).$$
The generator $`𝐬\lambda \mathrm{𝑖𝑒}(2)`$ is mapped to the product $`f,g\{f,g\}`$, which is cohomologous to zero, see \[22, Exercise 7\].
###### Example 31.
– Harrison cohomology. Here $`𝒫=𝒞\mathrm{𝑜𝑚}`$ is the operad for commutative associative algebras and $`𝐬(𝒫𝒫^!)=𝐬(𝒞\mathrm{𝑜𝑚}\mathrm{𝑖𝑒})𝐬\mathrm{𝑖𝑒}`$, therefore, as in Example 30,
(20)
$$𝒵_{𝒞\mathrm{𝑜𝑚}}𝐬\mathrm{𝑖𝑒}=\mathrm{𝐼𝑚}(L).$$
Equations (19) and (20) illustrate the obvious self-duality of the space of cup products:
$$𝒵_{𝒫^!}𝒵_𝒫,$$
compare Conjecture 5.
###### Example 32.
If $`𝒟=𝒜\mathrm{𝑠𝑠}𝒜\mathrm{𝑠𝑠}`$ is as in Example 15, then
$$𝒵_𝒟=𝐬𝒜\mathrm{𝑠𝑠}𝐬𝒜\mathrm{𝑠𝑠}.$$
Let us describe products corresponding to the generators of the $`4`$-dimensional vector space
$$(𝐬𝒜\mathrm{𝑠𝑠}𝐬𝒜\mathrm{𝑠𝑠})(2)=𝐬𝒜\mathrm{𝑠𝑠}(2)𝐬𝒜\mathrm{𝑠𝑠}(2)𝐤[\mathrm{\Sigma }_2]𝐤[\mathrm{\Sigma }_2].$$
Recall that $`C_𝒟^{}(A;A)`$ is the total complex of the meager bicomplex in Figure 2. Let $`^1`$ (resp. $`^2`$) be the cup product in the horizontal (resp. vertical) subcomplex in Figure 2. Let $`\pi _1`$ (resp. $`\pi _2`$) be the projection of $`C_𝒟^{}(A;A)`$ onto the horizontal (vertical) subcomplex. Likewise, let $`\iota _1`$ (resp. $`\iota _2`$) be the inclusion. Although neither $`\pi _i`$, $`\iota _i`$ nor $`^i`$ are chain maps ($`i=1,2`$), the compositions
$$f_1g:=\iota _1(\pi _1f^1\pi _1g)\text{ and }f_2g:=\iota _2(\pi _2f^2\pi _2g)$$
are chain operations. The generators of $`𝐬𝒜\mathrm{𝑠𝑠}(2)𝐬𝒜\mathrm{𝑠𝑠}(2)`$ then correspond to the four operations
$$f,gf_1g,f,gf_2g,f,gg_1f\text{ and }f,gg_2f.$$
The combination
$$(f_1g+f_2g)+(1)^{|f||g|}(g_1f+g_2f)$$
is cohomologous to zero and the image $`T(𝒵_𝒟(2))`$ of $`𝒵_𝒟`$ in $`H^1(_𝒟(2))`$ is easily seen to be $`3`$-dimensional.
## 6. Operad $`_{𝒜\mathrm{𝑠𝑠}}`$ and the Deligne conjecture
In this section we recall some results related to $`_{𝒜\mathrm{𝑠𝑠}}`$ and the Deligne conjecture. Let us make a necessary comment about our degree convention. We use the grading such that the intrinsic bracket of Section 4 has degree $`0`$ in $`_𝒫^{}(2)`$, while the $`n`$-fold cup products of Section 5 are of degree $`n1`$ in $`_𝒫^{}(n)`$. In the literature related to the Deligne conjecture, the convention under which the intrinsic bracket has degree $`1`$ and the $`n`$-fold cup products are of degree $`0`$ is often used. These two conventions are tied by the following regrading operator:
$$\mathrm{𝑅𝑒𝑔}(_𝒫^{}(n)):=_𝒫^{n1}(n).$$
In what follows we identify operads that differ only by the above regrading. In particular, the operad $`𝒢\mathrm{𝑒𝑟}`$ for Gerstenhaber algebras becomes identified with the operad $`\mathrm{𝑟𝑎𝑖𝑑}`$ for braid algebras (also called 1-algebras), see A.4.
Let us recall that a topological operad $`𝒜`$ is an $`E_2`$-operad if it has the homotopy type of the little discs operad $`𝒟_2`$ . According to the Formality Theorem , the operad $`S_{}(𝒜)`$ of singular chains on such an operad has the homotopy type of the operad $`\mathrm{𝑟𝑎𝑖𝑑}`$ for braid algebras.
(i) D. Tamarkin and B. Tsygan studied in \[32, Section 3\] a certain operad $`F=\{F(n)\}_{n1}`$ of natural operations on the cosimplicial Hochschild complex $`C^{}(X,X)`$ of a topological unital monoid $`X`$. The $`n`$-th space of this cosimplicial set is the space $`\mathrm{𝐶𝑜𝑛𝑡}(X^{\times n},X)`$ of continuous maps from the $`n`$-th cartesian power of $`X`$ to $`X`$. For each $`n1`$, $`F(n)`$ is a functor $`(\mathrm{\Delta }^{\mathrm{𝑜𝑝}})^n\times \mathrm{\Delta }\mathrm{𝖲𝖾𝗍𝗌}`$. They then considered a topological operad $`E=\{E(n)\}_{n1}`$ whose pieces are the topological realizations of these functors and claimed that $`E`$ is an $`E_2`$-operad.
It is not difficult to see that the operad $`\mathrm{𝐶𝑁}_{}(F)`$ of normalized chains of $`F`$ coincides with the operad $`\{_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}(n)\}_{n1}`$ (our $`_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}`$ without constants). Since $`(_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}^{}(0),\delta _{\underset{¯}{𝒜\mathrm{𝑠𝑠}}})`$ is acyclic (see Example 12), we could conclude that $`_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}`$ has the homotopy type of $`𝒢\mathrm{𝑒𝑟}`$, but we must bear in mind that the arguments in were merely sketched.
(ii) J.E. McClure and J.H. Smith considered in a dg-suboperad $`𝒮_2`$ of their “sequence” operad $`𝒮`$ and proved that $`𝒮_2`$ naturally acts on the Hochschild cochain complex of an associative algebra. In our terminology this means that they constructed a canonical map $`𝒮_2_{𝒜\mathrm{𝑠𝑠}}`$. They then verified the Deligne conjecture by showing, using a result of , that $`𝒮_2`$ has the homotopy type of the singular chain complex $`S_{}(𝒟_2)`$ of the little discs operad. Their proof is a very reliable one.
(iii) M. Kontsevich and Y. Soibelman introduced a “minimal operad” $`M`$ naturally acting on the Hochschild cochain complex of an $`A_{\mathrm{}}`$-algebra. In our terminology, $`M`$ was a suboperad of $`_{sh𝒜\mathrm{𝑠𝑠}}`$ generated by braces and cup products. They then argued that $`M`$ has the homotopy type of the operad of suitably defined piecewise algebraic chains on the operad $`FM_2`$ of the Fulton-MacPherson compactification of the configuration space of points in $`^2`$. Since $`FM_2`$ is, by \[30, Proposition 3.9\], an $`E_2`$-operad, they concluded that $`M`$ has the homotopy type of $`𝒢\mathrm{𝑒𝑟}`$.
(iv) R.M. Kaufmann realized in that the cellular chains $`\mathrm{𝐶𝐶}_{}(𝒞\mathrm{𝑎𝑐𝑡}^1)`$ on his operad-up-to-homotopy $`𝒞\mathrm{𝑎𝑐𝑡}^1`$ of spineless normalized cacti is a honest operad which naturally acts on the Hochschild cochain complex, via braces and cup products. By comparing $`𝒞\mathrm{𝑎𝑐𝑡}^1`$ to the operad $`𝒞\mathrm{𝑎𝑐𝑡}`$ of spineless (non-normalized) cacti, he concluded that $`CC_{}(𝒞\mathrm{𝑎𝑐𝑡}^1)`$ is a model for chains on the little discs operad $`𝒟_2`$.
All the proofs of the Deligne conjecture mentioned above use some special features of associative algebras and $`E_2`$-operads, such as the cosimplicial structure of the Hochschild cochain complex, Fiedorowicz’ detection principle, or a relation to the Fulton-MacPherson and cacti operads. None of these features are available for a general operad $`𝒫`$, we therefore think that the analysis of the homotopy type of $`_𝒫`$ for a general Koszul quadratic $`𝒫`$ is substantially more difficult than the analysis of $`_{𝒜\mathrm{𝑠𝑠}}`$.
Let us mention that there are other approaches to the Deligne conjecture, as D. Tamarkin’s proofs that use the Etingof-Kazhdan quantization , or those based on a suitable filtration of the Fulton-MacPherson compactification $`FM_2`$, see E. Getzler and J.D.S. Jones or A.A. Voronov .
## 7. Natural operations
Let use recall the following definitions which can be found for example in \[26, Section II.1.5\]. By a tree we mean a connected graph $`T`$ without loops. A valence of a vertex $`v`$ of $`T`$ is the number of edges adjacent to $`v`$. A leg or leaf of $`T`$ is an edge adjacent to a vertex of valence one, other edges of $`T`$ are interior. We in fact discard vertices of valence one at the endpoints of the legs, therefore the legs become “half-edges” having only one vertex. By a rooted or directed tree we mean a tree with a distinguished output leg called the root. The remaining legs are called the input legs of the tree. A tree with $`a`$ input legs labelled by elements of the set $`\{1,2,\mathrm{},a\}`$ is called an $`a`$-tree. A rooted tree is automatically oriented, each edge pointing towards the root. The edges pointing towards a given vertex $`v`$ are called the input edges of $`v`$, the number of these input edges is then the arity of $`v`$ denoted $`\mathrm{𝑎𝑟}(v)`$. Vertices of arity one are called unary, vertices of arity two binary, vertices of arity three ternary, etc.
Notation. Let $`n,m`$ and $`m_1,\mathrm{},m_n`$ be non-negative integers. In the rest of this section, $`i`$ will always denote an integer between $`0`$ and $`n`$, $`a:=m+1`$ and $`a_i:=m_i+1`$. We will also assume the notation introduced in A.6.
An $`n`$-linear natural operation
$$U:C_𝒫^{m_1}(A;A)\mathrm{}C_𝒫^{m_n}(A;A)C_𝒫^m(A;A)$$
is given by the following data.
(i) A rooted $`a`$-tree $`T`$ with $`n`$ white vertices $`w_1,\mathrm{},w_n`$ of arities $`a_1,\mathrm{},a_n`$, and $`k`$ at least binary black vertices, $`k0`$.
(ii) A linear order on the set of input edges of each white vertex of $`T`$.
(iii) A decoration of black vertices of $`T`$ by elements of $`𝒫`$.
(iv) A linear map $`\mathrm{\Phi }:𝐬𝒫^!(a_1)\mathrm{}𝐬𝒫^!(a_n)𝐬𝒫^!(a)`$.
Given the above data and $`f_iC_𝒫^{m_i}(A;A)`$, the value $`U(f_1,\mathrm{},f_n)C_𝒫^m(A;A)`$ is defined as follows. Let us decompose
$$f_i=\underset{\kappa _i}{}\varphi _i^{\kappa _i}q_{\kappa _i}^i[\mathrm{𝐿𝑖𝑛}(A^{a_i},A)𝐬𝒫^!(a_i)]^{\mathrm{\Sigma }_{a_i}}C_𝒫^{m_i}(A;A),$$
where $`\varphi _i^{\kappa _i}\mathrm{𝐿𝑖𝑛}(A^{a_i},A)`$, $`q_{\kappa _i}^i𝐬𝒫^!(a_i)`$ and $`\kappa _i`$ is a summation index. Since the inputs of white vertices are linearly ordered, each $`\varphi _i^{\kappa _i}`$ determines a decoration of the white vertex $`w_i`$ by an element of $`\mathrm{𝐿𝑖𝑛}(A^{a_i},A)=\mathrm{𝑛𝑑}_A(a_i)`$. Recall that $`A`$ is a $`𝒫`$-algebra with the structure homomorphism $`\alpha :𝒫\mathrm{𝑛𝑑}_A`$. Applying $`\alpha `$ to the decorations of the black vertices we decorate also black vertices with elements of $`\mathrm{𝑛𝑑}_A`$. So $`T`$ is now a tree with all vertices decorated by $`\mathrm{𝑛𝑑}_A`$. The composition in the operad $`\mathrm{𝑛𝑑}_V`$ along $`T`$ determines, for each $`k_1,\mathrm{},k_n`$, the element
$$T(\varphi _1^{\kappa _1},\mathrm{},\varphi _n^{\kappa _n})\mathrm{𝐿𝑖𝑛}(A^a,A).$$
Let
$$\stackrel{~}{U}(f_1,\mathrm{},f_n):=\underset{\kappa _1,\mathrm{},\kappa _n}{}T(\varphi _1^{\kappa _1},\mathrm{},\varphi _n^{\kappa _n})\mathrm{\Phi }(q_{\kappa _1}^1,\mathrm{},q_{\kappa _n}^n)\mathrm{𝐿𝑖𝑛}(A^a,A)𝐬𝒫^!(a)\stackrel{~}{C}_𝒫^m(A;A).$$
Finally, let $`U(f_1,\mathrm{},f_n):=\mathrm{𝐴𝑣𝑒𝑟}(\stackrel{~}{U}(f_1,\mathrm{},f_n))C_𝒫^m(A;A)`$. It follows from an elementary combinatorics of trees that
$$\mathrm{deg}(U)=\mathrm{𝑎𝑟}(b_1)+\mathrm{}+\mathrm{𝑎𝑟}(b_k)k,$$
therefore $`\mathrm{deg}(U)`$ is always non-negative and $`\mathrm{deg}(U)=0`$ if and only if $`T`$ has no black vertex.
###### Definition 33.
Let $`_𝒫:=\{_𝒫(n)\}_{n0}`$ be the operad spanned by all natural operations $`U=U_{(T,\mathrm{\Phi })}`$ in the above sense. Since the differential $`d_𝒫`$ of $`C_𝒫^{}(A;A)`$ is itself a natural operation living in $`_𝒫^1(1)`$, it induces a differential $`\delta _𝒫`$ on $`_𝒫`$ by the standard formula
$$\delta _𝒫(U)(f_1,\mathrm{},f_n):=d_𝒫U(f_1,\mathrm{},f_n)(1)^{|U|}\underset{1in}{}(1)^{|f_1|+\mathrm{}+|f_{i1}|}U(f_1,\mathrm{},d_𝒫f_i,\mathrm{},f_n),$$
making $`_𝒫=(_𝒫^{},\delta _𝒫)`$ a dg-operad.
Heuristically, the value $`U_{(T,\mathrm{\Phi })}(f_1,\mathrm{},f_n)`$ is given by inserting $`f_i`$ at the vertex $`w_i`$ of $`T`$, $`1in`$, and then performing the composition along $`\mathrm{\Phi }`$. The operadic composition of $`_𝒫`$ is the vertex insertion similar to that of and the symmetric group permutes the labels of white vertices. In the following definition we introduce a non-$`\mathrm{\Sigma }`$ version of $`_𝒫`$.
###### Definition 34.
Suppose $`𝒫`$ is the symmetrization of a non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$. Let $`_{\underset{¯}{𝒫}}`$ be the dg-suboperad of $`_𝒫`$ spanned by natural operations $`U_{(T,\mathrm{\Phi })}`$ as in Definition 33 such that the tree $`T`$ is planar, with black vertices decorated by elements of $`\underset{¯}{𝒫}`$, and the map $`\mathrm{\Phi }`$ such that
$$\mathrm{\Phi }(𝐬\underset{¯}{𝒫}^!(a_1)\mathrm{}𝐬\underset{¯}{𝒫}^!(a_n))𝐬\underset{¯}{𝒫}^!(a).$$
###### Example 35.
– Constants. Let us see what happens if $`T`$ is the $`a`$-corolla with one black vertex decorated by $`p𝒫(a)`$ and no white vertices as in Figure 5.
The map $`\mathrm{\Phi }:𝐤𝐬𝒫^!(a)`$ is given by specifying an element $`\phi :=\mathrm{\Phi }(1)𝐬𝒫^!(a)`$ and $`\stackrel{~}{U}`$ determined by this $`\mathrm{\Phi }`$ equals $`\alpha (p)\phi \stackrel{~}{C}_𝒫^m(A;A)`$. Since $`\alpha `$ is equivariant,
$$\mathrm{𝐴𝑣𝑒𝑟}(\alpha (p)\phi )=(\alpha 11)(\mathrm{𝐴𝑣𝑒𝑟}(p\phi ))$$
therefore $`U:=\mathrm{𝐴𝑣𝑒𝑟}(\alpha (p)\phi )C_𝒫^m(A;A)`$ is parametrized by an element in the image of the averaging map
$$\mathrm{𝐴𝑣𝑒𝑟}:𝒫(a)𝐬𝒫^!(a)(𝒫(a)𝐬𝒫^!(a))^{\mathrm{\Sigma }_a},$$
in other words,
$$_𝒫^m(0)𝐬(𝒫𝒫^!)(a)^{\mathrm{\Sigma }_a},m0.$$
It is equally easy to see that, for a quadratic Koszul non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$,
$$_{\underset{¯}{𝒫}}^m(0)𝐬(\underset{¯}{𝒫}\underset{¯}{𝒫}^!)(a),m0.$$
###### Example 36.
Unary operations of degree $`0`$. Now $`T`$ is an $`a`$-corolla with one white planar vertex and no black vertices, with input legs labelled $`\sigma (1),\mathrm{},\sigma (a)`$, $`\sigma \mathrm{\Sigma }_a`$, as shown in Figure 5, and $`\mathrm{\Phi }:𝐬𝒫^!(a)𝐬𝒫^!(a)`$ is a linear map. If
$$f=\underset{\kappa }{}\varphi ^\kappa q_\kappa [\mathrm{𝐿𝑖𝑛}(A^a,A)𝐬𝒫^!(a)]^{\mathrm{\Sigma }_a}C_𝒫^m(A;A),$$
then $`U(f)=\mathrm{𝐴𝑣𝑒𝑟}(_\kappa \varphi ^\kappa \sigma ^1\mathrm{\Phi }(q_\kappa ))`$. Since $`f=\varphi ^\kappa q_\kappa `$ is $`\mathrm{\Sigma }_a`$-stable,
$$U(f)=\varphi ^\kappa \mathrm{𝐴𝑣𝑒𝑟}(\mathrm{\Phi }_\kappa )(q_\kappa ).$$
Therefore $`U(f)`$ is given by a $`\mathrm{\Sigma }_a`$-equivariant map $`\mathrm{\Psi }:=^{a1}\mathrm{𝐴𝑣𝑒𝑟}(\mathrm{\Phi }):𝒫^!(a)𝒫^!(a)`$, thus
(21)
$$_𝒫^0(1)\mathrm{𝐿𝑖𝑛}_\mathrm{\Sigma }(𝒫^!,𝒫^!),$$
the space of all collections $`\{\psi _n:𝒫^!(n)𝒫^!(n)\}_{n0}`$ of equivariant maps. We leave as an exercise to verify that, for a non-$`\mathrm{\Sigma }`$ quadratic Koszul operad $`\underset{¯}{𝒫}`$,
$$_{\underset{¯}{𝒫}}^0(1)\mathrm{𝐿𝑖𝑛}(\underset{¯}{𝒫}^!,\underset{¯}{𝒫}^!).$$
###### Example 37.
Projections. Let $`p_m_𝒫^0(1)`$ be given, in identification (21), by $`\mathrm{\Psi }\mathrm{𝐿𝑖𝑛}_\mathrm{\Sigma }(𝒫^!,𝒫^!)`$ defined as
$$\mathrm{\Psi }|_{𝒫^!(a)}=\{\begin{array}{cc}11_{𝒫^!(a)},\hfill & \text{for }a=m+1\text{ and}\hfill \\ 0,\hfill & \text{othewise.}\hfill \end{array}$$
Clearly, $`p_m`$ is the projection $`C_𝒫^{}(A;A)C_𝒫^m(A;A)`$. The system of all these projections makes $`_𝒫`$ an $`^0`$-colored operad, where $`^0`$ is the set of non-negative integers. Since these projections do not commute with $`d_𝒫`$ (that is $`\delta _𝒫(p_m)0`$ for a generic $`𝒫`$), $`(_𝒫,\delta _𝒫)`$ is not a dg $`^0`$-colored operad.
###### Example 38.
Cup products. In this example we explain how an element
$$t=p𝐬q𝒫(n)𝐬𝒫^!(n)𝐬(𝒫𝒫^!)(n)$$
determines a natural operation in $`_𝒫^{n1}(n)`$. Let $`T`$ be as in Figure 6, with the black vertex decorated by $`p𝒫(n)`$,
and let the linear map $`\mathrm{\Phi }:𝐬𝒫^!(a_1)\mathrm{}𝐬𝒫^!(a_n)𝐬𝒫^!(a)`$ be given by the operadic composition in $`𝐬𝒫^!`$:
$$\mathrm{\Phi }(𝐬q_1,\mathrm{},𝐬q_n):=𝐬q(𝐬q_1,\mathrm{},𝐬q_n),q_i𝒫^!(a_i),1in.$$
It is more or less clear that the natural operation $`U_{(T,\mathrm{\Phi })}`$ determined by the above data agrees with the cup product $`\mathrm{𝑐𝑢𝑝}(t)`$ of Theorem 23. We recommend as another exercise to verify that also the intrinsic bracket described in (9) is given by natural operations in the sense of this section.
###### Example 39.
Let us describe all natural operations $`C_𝒫^1(A;A)C_𝒫^1(A;A)C_𝒫^2(A;A)`$ for some particular choices of $`𝒫`$.
(i) Hochschild cohomology. For $`𝒫=𝒜\mathrm{𝑠𝑠}`$, $`C_𝒫^1(A;A)=\mathrm{𝐿𝑖𝑛}(A^2,A)`$, $`C_𝒫^2(A;A)=\mathrm{𝐿𝑖𝑛}(A^3,A)`$, and the only natural operations $`C_𝒫^1(A;A)C_𝒫^1(A;A)C_𝒫^2(A;A)`$ are linear combinations of
$$f,g(f_1g)\sigma ,f,g(f_2g)\sigma ,f,g(g_1f)\sigma ,f,g(g_2f)\sigma ,\sigma \mathrm{\Sigma }_3,$$
where $`_1`$, $`_2`$ are Gerstenhaber-type products given by
(22)
$$(u_1v)(a,b,c):=u(v(a,b),c)),(u_2v)(a,b,c):=u(a,v(b,c)),$$
for $`u,vC_𝒫^2(A;A)`$, $`a,b,cA`$, and $`\sigma \mathrm{\Sigma }_3`$ permutes the factors of $`A^3`$ in the usual way. Operations belonging to $`_{\underset{¯}{𝒜\mathrm{𝑠𝑠}}}^0(2)`$ are linear combinations of the operations (22) with $`\sigma =11_3`$, the unit of $`\mathrm{\Sigma }_3`$.
(ii) Chevalley-Eilenberg cohomology. If $`𝒫=\mathrm{𝑖𝑒}`$, then $`C_𝒫^1(A;A)=\mathrm{𝐿𝑖𝑛}(^2,A)`$, $`C_𝒫^2(A;A)=\mathrm{𝐿𝑖𝑛}(^3A,A)`$, where $`^nA`$ denotes the $`n`$-th exterior power. The only natural operations $`C_𝒫^1(A;A)C_𝒫^1(A;A)C_𝒫^2(A;A)`$ are linear combinations of
$$f,gfg\text{ and }f,ggf,$$
where
$$(uv)(a,b,c):=u(v(a,b),c)+u(v(b,c),a)+u(v(c,a),b)$$
for $`u,vC_𝒫^2(A;A)`$ and $`a,b,cA`$.
(iii) Harrison cohomology. If $`𝒫=𝒞\mathrm{𝑜𝑚}`$, then
$$C_𝒫^1(A;A)=\{u\mathrm{𝐿𝑖𝑛}(A^2,A);u(a,b)u(b,a)=0\}$$
and $`C_𝒫^2(A;A)`$ consists of all $`w\mathrm{𝐿𝑖𝑛}(A^3,A)`$ such that
$$w(a,b,c)w(b,a,c)+w(b,c,a)=w(a,b,c)w(a,c,b)+w(c,a,b)=0,$$
for $`a,b,cA`$. Natural operations $`C_𝒫^1(A;A)C_𝒫^1(A;A)C_𝒫^2(A;A)`$ are linear combinations of
$$f,gfg\text{ and }f,ggf,$$
where
$$uv:=u(v(a,b)c)u(v(b,c),a),$$
for $`u,vC_𝒫^2(A;A)`$ and $`a,b,cA`$.
## Appendix: Notations, conventions and background material
###### A.1.
In this note, an operad means an operad in the category of differential graded (dg) vector spaces, that is, a sequence $`𝒫=\{𝒫(n)\}_{n0}`$ of right $`\mathrm{\Sigma }_n`$-modules with structure operations
$$\gamma :𝒫(n)𝒫(k_1)\mathrm{}𝒫(k_n)𝒫(k_1+\mathrm{}+k_n),$$
for $`n1`$ and $`k_1,\mathrm{},k_n0`$, and a unit map $`\eta :𝐤𝒫(1)`$ that satisfy the usual axioms . Instead of $`\gamma (pp_1\mathrm{}p_n)`$ we will often write $`\gamma (p,p_1,\mathrm{},p_n)`$ or $`p(p_1,\mathrm{},p_n)`$. Recall that operads can be equivalently defined using the $`_i`$-operations
$$_i:𝒫(m)𝒫(n)𝒫(m+n1)$$
defined, for $`m,n0`$, $`1im`$, by
$$p_iq:=\gamma (pe^{(i1)}qe^{mi}),$$
where $`e:=\eta (1)`$.
If we remove from the above definition all references to the symmetric group actions, we get a definition of a non-$`\mathrm{\Sigma }`$ operad. Each non-$`\mathrm{\Sigma }`$ operad $`\underset{¯}{𝒫}`$ generates a unique (usual) operad $`𝒫`$ such that $`𝒫(n)\underset{¯}{𝒫}(n)𝐤[\mathrm{\Sigma }_n]`$, $`n0`$.
###### A.2.
For each set of operations $`E`$, there exists the free operad $`\mathrm{\Gamma }(E)`$ generated by $`E`$ \[26, Proposition II.1.92\]. Let $`\mu `$ denote a bilinear operation placed in degree $`0`$. The operad $`𝒜\mathrm{𝑠𝑠}`$ for associative algebras is the quotient
$$𝒜\mathrm{𝑠𝑠}:=\mathrm{\Gamma }(\mu )/(\mu _1\mu \mu _2\mu ),$$
where $`(\mu _1\mu \mu _2\mu )`$ denotes the operadic ideal generated by the associativity axiom for $`\mu `$.
If $`\lambda `$ is a skew-symmetric bilinear operation, then the operad for Lie algebras is the quotient
$$\mathrm{𝑖𝑒}:=\mathrm{\Gamma }(\lambda )/(\mathrm{𝐽𝑎𝑐𝑜𝑏𝑖}(\lambda )),$$
where
$$\mathrm{𝐽𝑎𝑐𝑜𝑏𝑖}(\lambda ):=\underset{\sigma C_3}{}(\lambda _1\lambda )\sigma $$
with the summation taken over the order $`3`$ cyclic subgroup $`C_3`$ of $`\mathrm{\Sigma }_3`$, denotes the Jacobi identity for $`\lambda `$.
Finally, for an arbitrary differential graded vector space $`V`$, there is the endomorphism operad $`\mathrm{𝑛𝑑}_V=\{Lin(V^n,V)\}_{n0}`$, with structure operations given as the usual composition of multilinear maps. A $`𝒫`$-algebra is then a homomorphism $`\alpha :𝒫\mathrm{𝑛𝑑}_V`$. We sometimes call $`\alpha `$ also an action of $`𝒫`$ on $`V`$.
###### A.3.
The suspension $`𝐬A=\{𝐬A(n)\}_{n0}`$ of a $`\mathrm{\Sigma }`$-module $`A=\{A(n)\}_{n0}`$ is defined by
$$𝐬A(n):=^{n1}A(n)\mathrm{𝑠𝑔𝑛}_n,$$
where $`\mathrm{𝑠𝑔𝑛}_n`$ denotes the signum representation of $`\mathrm{\Sigma }_n`$, see \[26, Definition II.3.15\]. If $`𝒫`$ is an operad, then the collection $`𝐬𝒫`$ carries a canonical induced operad structure and the operad $`𝐬𝒫`$ is called the operadic suspension of $`𝒫`$. For any two operads $`𝒫`$ and $`𝒬`$,
$$𝐬(𝒫𝒬)𝐬𝒫𝒬𝒫𝐬𝒬.$$
###### A.4.
An $`(m,n)`$-algebra is \[8, Example 9.4\] a graded vector space $`A`$ together with two bilinear maps, $`:AAA`$ of degree $`m`$, and $`[,]:AAA`$ of degree $`n`$ ($`m`$ and $`n`$ are natural numbers), such that, for any homogeneous $`a,b,cA`$,
* $`ab=(1)^{|a||b|+m}ba`$,
* $`[a,b]=(1)^{|a||b|+n}[b,a]`$,
* $``$ is associative in the sense that
$$a(bc)=(1)^{m(|a|+1)}(ab)c,$$
* $`[,]`$ satisfies the following form of the Jacobi identity:
$$(1)^{|a|(|c|+n)}[a,[b,c]]+(1)^{|b|(|a|+n)}[b,[c,a]]+(1)^{|c|(|b|+n)}[c,[a,b]]=0,$$
* the operations $``$ and $`[,]`$ are compatible in the sense that
$$(1)^{m|a|}[a,bc]=[a,b]c+(1)^{(|b||c|+m)}[a,c]b.$$
$`(0,1)`$-algebras were considered in under the name $`2`$-algebras or braid algebras. The corresponding operad $`\mathrm{𝑟𝑎𝑖𝑑}`$ is isomorphic to the homology of the little discs operad $`𝒟_2`$, $`\mathrm{𝑟𝑎𝑖𝑑}H_{}(𝒟_2)`$. Following \[11, Section 10\], we call $`(1,0)`$-algebras Gerstenhaber algebras, though the terminology is not unique, compare for instance \[10, Subsection 10.2\] where a Gerstenhaber algebra means a $`(0,1)`$-algebra.
###### A.5.
Let $`M`$ be a right module over a finite group $`G`$. We denote, as usual
$$M^G:=\{mM;mg=g\text{ for all }gG\}\text{ and }M_G:=\frac{M}{(mmg;mM,gG)}.$$
Let $`\mathrm{𝐴𝑣𝑒𝑟}:MM^G`$ be the “averaging” map given by
$$\mathrm{𝐴𝑣𝑒𝑟}(m):=\frac{1}{|G|}\underset{gG}{}mg.$$
It is a standard fact that the composition $`\pi \iota `$ of the projection $`\pi :MM_G`$ with the inclusion $`\iota :M^GM`$ is the identity and that $`\mathrm{𝐴𝑣𝑒𝑟}`$ is a right inverse to $`\iota `$.
###### A.6.
Let us recall the operadic cochain complex and introduce some useful notations. As a graded vector space, the operadic cochain complex is defined by \[26, Definition II.3.99\]:
(23)
$$C_𝒫^{n1}(A;A)=[\mathrm{𝐿𝑖𝑛}((A)^n,A)𝒫^!(n)]^{\mathrm{\Sigma }_n},n1,$$
where $`A`$ denotes the desuspension of the graded vector space $`A`$. It will be convenient to denote
$$\stackrel{~}{C}_𝒫^{n1}(A;A):=\mathrm{𝐿𝑖𝑛}((A)^n,A)𝒫^!(n),$$
so that $`C_𝒫^{n1}(A;A)\stackrel{~}{C}_𝒫^{n1}(A;A)^{\mathrm{\Sigma }_n}`$. The averaging over the $`\mathrm{\Sigma }_n`$-action defines an epimorphism
$$\mathrm{𝐴𝑣𝑒𝑟}:\stackrel{~}{C}_𝒫^{n1}(A;A)C_𝒫^{n1}(A;A)$$
of graded modules which is a left inverse to the inclusion
$$\iota :C_𝒫^{n1}(A;A)\stackrel{~}{C}_𝒫^{n1}(A;A).$$
We will often use the following canonical isomorphisms of graded $`\mathrm{\Sigma }_n`$-modules:
$`\stackrel{~}{C}_𝒫^{n1}(A;A)`$ $`=`$ $`\mathrm{𝐿𝑖𝑛}((A)^n,A)𝒫^!(n)^{n1}(\mathrm{𝐿𝑖𝑛}(A^n,A)𝒫^!(n)\mathrm{𝑠𝑔𝑛}_n)`$
$``$ $`𝐬(\mathrm{𝐿𝑖𝑛}(A^n,A)𝒫^!(n))𝐬\mathrm{𝑛𝑑}_A(n)𝒫^!(n)`$
$``$ $`\mathrm{𝑛𝑑}_A(n)𝐬𝒫^!(n)\mathrm{𝑛𝑑}_A(n)𝒫^!(n).`$
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# Spin 3/2 Penta-quarks in anisotropic lattice QCD
## I Introduction
The recent discovery of the manifestly exotic baryon $`\mathrm{\Theta }^+(1540)`$ by the LEPS group at SPring-8 has made a great impact on the exotic hadron physics nakano . Apart from other pentaquark baryon candidates, $`\mathrm{\Xi }^{}(1862)`$ NA49 and $`\mathrm{\Theta }_c(3099)`$ H1 , several other candidates of exotic hadrons have also been discovered, such as $`X(3872)`$, $`D_s(2317)`$, $`S_0(3115)`$, $`X(3940)`$ and $`Y(3840)`$ exotics . They also receive an increasing interest from theoretical side as well. $`\mathrm{\Theta }^+(1540)`$ is supposed to have baryon number $`B=1`$, charge $`Q=+1`$ and strangeness $`S=+1`$. Since $`uudd\overline{s}`$ is the simplest quark content to implement this quantum number, $`\mathrm{\Theta }^+(1540)`$ is a manifestly exotic penta-quark(5Q) state. The penta-quark $`\mathrm{\Theta }^+`$ had been considered several times even before the experimental discovery diakonov ; jaffe-76 ; strottman ; weigel ; praszalowicz . In particular, LABEL:diakonov provided the direct motivation of the experimental search nakano . The discovered peak in the $`nK^+`$ invariant mass is centered at $`1.54\pm 0.01`$ GeV with a width smaller than 25 MeV. At the present stage, some groups confirmed the LEPS discoverydiana ; clas ; saphir ; experiments , while the others reported null resultsnull . It will still take a while to establish the existence or non-existence of $`\mathrm{\Theta }^+(1540)`$ experimentally hicks . LABEL:saphir claims that $`\mathrm{\Theta }^+`$ must be isoscalar, since no $`\mathrm{\Theta }^{++}`$ is observed in the $`pK^+`$ invariant mass spectrum.
Enormous theoretical efforts have been devoted to 5Q baryons diakonov ; jaffe-76 ; strottman ; weigel ; praszalowicz ; oka ; zhu.review ; cohen ; itzhaki ; kim ; hosaka ; jaffe ; lipkin ; carlson-positive ; stancu ; jennings ; glozman ; enyo ; zhu ; matheus ; sugiyama ; carlson ; shinozaki ; huang ; maezawa ; shlee ; narodetskii ; sugamoto ; suganuma ; okiharu ; bicudo ; oset ; hosaka32 ; nishikawa32 ; takeuchi32 ; sugiyama32 ; zhu32 ; inoue32 ; jaffe32 ; huang32 ; capstic32 ; dudek32 ; hyodo32 ; nam32 . One of the most challenging problems in understanding its structure is its extremely narrow decay width as $`\mathrm{\Gamma }1`$ MeV pdg . Several ideas have been proposed: (1) $`I=2`$ possibility capstic32 , (2) Jaffe-Wilczek’s diquark picture jaffe , (3) $`\pi KN`$ hepta-quark picture bicudo ; kishimoto , (4) string picture suganuma ; sugamoto , (5) $`J^P=3/2^{}`$ possibility hosaka ; enyo . Although each gives a mechanism to explain the narrow decay width, none of them can satisfy all the known properties of $`\mathrm{\Theta }^+(1540)`$ simultaneously.
In this paper, we are interested in the $`J=3/2`$ possibility, $`J^P=3/2^{}`$ in particular. Note that the spin of $`\mathrm{\Theta }^+(1540)`$ has not yet been determined experimentally. In the constituent quark picture, the narrow decay width of $`J^P=3/2^{}`$ penta-quarks can be understood in the following way hosaka ; enyo . We expect that the special configuration $`(0s)^5`$ is dominant in the 5Q ground-state in $`J^P=3/2^{}`$ channel. Although $`J^P=3/2^{}`$ penta-quarks can decay to KN in the d-wave, the spectroscopic factor to find d-wave KN states in the dominant $`(0s)^5`$ configuration vanishes. Since the decay is thus allowed only through its sub-dominant d-wave configuration, the decay width is suppressed. Note that it is further suppressed by the d-wave centrifugal barrier, leading to the significantly narrow decay width of $`J^P=3/2^{}`$ penta-quarks. A possible disadvantage of $`J^P=3/2^{}`$ assignment is that such a state tends to be massive due to the color-magnetic interaction in the constituent quark models, which seems to be one of the main reasons why there are only a limited number of effective model studies for spin 3/2 penta-quarks hosaka32 ; nishikawa32 ; takeuchi32 ; sugiyama32 ; zhu32 ; inoue32 ; jaffe32 ; huang32 ; capstic32 ; hyodo32 ; nam32 ; enyo . However, it is not clear whether these conventional framework is applicable to a new exotic 5Q system as $`\mathrm{\Theta }^+(1540)`$ without involving any modifications. Indeed, a model was proposed where a part of the role of the color-magnetic interaction can be played by the flavor-spin interaction, which makes the mass-splitting between the $`1/2^{}`$ and the $`3/2^{}`$ states smaller takeuchi32 .
There have been several lattice QCD calculations of 5Q states by todayscikor12 ; sasaki ; chiu ; kentacky ; ishii12 ; rabbit ; lasscock12 ; alexandrou12 ; csikor122 ; holland ; lasscock32 . However, these studies are restricted to $`J^P=1/2^\pm `$ channels except for a very recent one lasscock32 . Enormous efforts are being devoted to more accurate studies of $`J^P=1/2^\pm `$ states, using the variational technique to extract multiple excited states, among which a compact resonance state is sought for. Indeed, quite large scale calculations are planned and being performedfleming attempting to elucidate some of the mysterious natures of $`\mathrm{\Theta }^+(1540)`$ such as its diquark structure and/or non-localities desired in interpolating fields. Here, we emphasize again that these studies are aiming at $`J^P=1/2^\pm `$ states, not at $`3/2^\pm `$ states.
In this paper, we present anisotropic lattice QCD results on 5Q states in $`J^P=3/2^\pm `$ channels using a large number of gauge configurations as $`N_{\mathrm{conf}}=1000`$ as an attempt to search for a low-lying 5Q state in $`J^P=3/2^\pm `$ channel. We adopt the standard Wilson gauge action at $`\beta =5.75`$ on the $`12^3\times 96`$ lattice with the renormalized anisotropy $`a_s/a_t=4`$. The anisotropic lattice is known to serve as a powerful tool for high-precision measurements of temporal correlators klassen ; matsufuru ; nemoto ; ishii-gb . The large number of gauge configurations $`N_{\mathrm{conf}}=1000`$ plays a key role in our calculation, because 5Q correlators in $`J^P=3/2^\pm `$ channels are found to be quite noisy. For quark action, we adopt $`O(a)`$-improved Wilson (clover) action with four values of the hopping parameter as $`\kappa =0.1210(0.0010)0.1240`$. One of the purpose of our calculation is to examine how the results depend on the choice of interpolating field operators. We employ several types of interpolating fields as (a) the NK-type, (b) the (color-)twisted NK-type, (c) a diquark-type, and adopt a smeared source to enhance the low-lying spectra.
In $`J^P=3/2^{}`$ channel, we obtain massive states $`m_{5\mathrm{Q}}2.12.2`$ GeV except for the diquark-type interpolating field, which involves a considerable size of the statistical error. In $`J^P=3/2^+`$ channel, we obtain more massive states $`m_{5\mathrm{Q}}2.42.6`$ GeV. None of these 5Q states appear below the NK threshold. Note that the NK threshold is raised up by about $`200250`$ MeV due to the finite extent of the spatial lattice as $`L2.15`$ fm, from which we expect the penta-quark signal to appear below the (raised) NK threshold considering the empirical mass difference between N+K(1440) and $`\mathrm{\Theta }^+(1540)`$. To clarify whether our 5Q states are compact resonance or not, we perform an analysis with the hybrid boundary condition(HBC), which was recently proposed in LABEL:ishii12. HBC analysis indicates that no compact 5Q resonance is contained in our 5Q states both in $`J^P=3/2^\pm `$ channels.
The paper is organized as follows. In Sect. II, we discuss the general formalisms. We begin by introducing several types of interpolating fields, determining their parity transformation properties. We next consider the temporal correlator and its spectral decomposition. We finally discuss the two-particle scattering states involved in 5Q spectra, and introduce the hybrid boundary condition (HBC) to examine whether a state of our concern is a compact resonance state or not. Sect. III is devoted to the brief descriptions of our lattice action and parameters. In Sect. IV, we present our numerical results for $`J^P=3/2^\pm `$ channels in the standard periodic boundary condition(PBC). We show 5Q correlators of various interpolating fields, i.e., the NK-type, the (color-)twisted NK-type, the diquark-type. In Sect. V, we attempt to determine whether these 5Q states are compact 5Q resonance states or two-particle scattering states by using the HBC. In Sect. VI, we summarize our results.
## II General formalisms
### II.1 Interpolating fields
We consider an iso-scalar interpolating field of NK-type in Rarita-Schwinger form ioffe ; benmerrouche ; hemmert as
$`\psi _\mu `$ $``$ $`ϵ_{abc}\left(u_a^TC\gamma _5d_b\right)u_c\left(\overline{s}_d\gamma _\mu d_d\right)`$
$`ϵ_{abc}\left(u_a^TC\gamma _5d_b\right)d_c\left(\overline{s}_d\gamma _\mu u_d\right),`$
where $`\mu `$ denotes the Lorentz index, $`ad`$ refer to the color indices, and $`C=\gamma _4\gamma _2`$ denotes the charge conjugation matrix. Unless otherwise indicated, the gamma matrices are represented in the Euclidean form given in LABEL:montvay.
We are also interested in the (color-)twisted NK-type interpolating field as
$`\psi _\mu `$ $``$ $`ϵ_{abc}\left(u_a^TC\gamma _5d_b\right)u_d\left(\overline{s}_d\gamma _\mu d_c\right)`$
$`ϵ_{abc}\left(u_a^TC\gamma _5d_b\right)d_d\left(\overline{s}_d\gamma _\mu u_c\right),`$
which is an extension to the one originally proposed in LABEL:scikor12 to study $`J^P=1/2^P`$ 5Q states. It has a slightly more elaborate color-structure than Eq. (II.1), suggesting somewhat stronger coupling to a genuine 5Q state, if it exists, than simple NK states.
Another interpolating fields of our possible interests are diquark-type interpolating fields such as
$$\psi _\mu ϵ_{abc}ϵ_{def}ϵ_{cfg}\left(u_a^TC\gamma _5d_b\right)\left(u_d^TC\gamma _5\gamma _\mu d_e\right)C\gamma _5\overline{s}_g$$
(3)
which is an extension to the one proposed in Refs. sasaki ; sugiyama . The first factor corresponds to the scalar diquark (color $`\overline{\mathrm{𝟑}}`$, $`I=0`$, $`J^P=0^+`$), which is expected to play important roles in hadron physics jaffe-exotica . The second factor corresponds to the vector diquark (color $`\overline{\mathrm{𝟑}}`$, $`I=0`$, $`J^P=1^{}`$). Note that, although the axial-vector diquark (color $`\overline{\mathrm{𝟑}}`$, $`I=1`$, $`J^P=1^+`$) is considered to play a more important role than the vector diquark, it cannot replace the vector diquark due to its iso-vector nature. Unless otherwise indicated, we refer to Eq. (3) as the “diquark-type” interpolating field. We can also consider another interpolating field of diquark-type as
$$\psi _\mu ϵ_{abc}ϵ_{def}ϵ_{cfg}\left(u_a^TCd_b\right)\left(u_d^TC\gamma _5\gamma _\mu d_e\right)C\overline{s}_g,$$
(4)
which consists of the pseudo-scalar diquark (color $`\overline{3}`$, $`I=0`$, $`J^P=0^{}`$) and the vector diquark. However, actual lattice QCD calculation shows that its correlator is afflicted with quite a huge statistical error. A possible reason could be attributed to the fact that both of these diquark fields do not survive the non-relativistic limit. Hence, we do not consider this interpolating field in this paper.
Under the spatial reflection of the quark fields as
$$q(\tau ,\stackrel{}{x})\gamma _4q(\tau ,\stackrel{}{x}),$$
(5)
all of these interpolating fields transform as
$$\psi _i(\tau ,\stackrel{}{x})\gamma _4\psi _i(\tau ,\stackrel{}{x}),$$
(6)
for $`i=1,2,3`$.
### II.2 5Q correlators and parity projection
We consider the Euclidean temporal correlator as
$$G_{\mu \nu }(\tau )\underset{\stackrel{}{x}}{}\psi _\mu (\tau ,\stackrel{}{x})\overline{\psi }_\nu (0,\stackrel{}{0}),$$
(7)
where $`_\stackrel{}{x}`$ projects the total 5Q momentum to zero. Since the spin 3/2 contribution from the temporal component of Rarita-Schwinger spinor vanishes in the rest frame, we can restrict ourselves to the spatial parts, i.e., $`\mu \nu =1,2,3`$. Now, Eq. (7) is decomposed in the following way:
$`G_{ij}(\tau )`$ $`=`$ $`𝐏_{ij}^{(3/2)}G^{(3/2)}(\tau )+𝐏_{ij}^{(1/2)}G^{(1/2)}(\tau ),`$ (8)
where $`i,j=1,2,3`$ denote the spatial part of the Lorentz indices, $`𝐏^{(3/2)}`$ and $`𝐏^{(1/2)}`$ denote the projection matrices onto the spin 3/2 and 1/2 subspaces defined as
$`𝐏_{ij}^{(3/2)}`$ $``$ $`\delta _{ij}(1/3)\gamma _i\gamma _j,`$ (9)
$`𝐏_{ij}^{(1/2)}`$ $``$ $`(1/3)\gamma _i\gamma _j.`$
They satisfy the following relations as
$`𝐏_{ij}^{(3/2)}𝐏_{jk}^{(3/2)}`$ $`=`$ $`𝐏_{ik}^{(3/2)}`$ (10)
$`𝐏_{ij}^{(1/2)}𝐏_{jk}^{(1/2)}`$ $`=`$ $`𝐏_{ik}^{(1/2)}`$
$`𝐏_{ij}^{(1/2)}+𝐏_{ij}^{(3/2)}`$ $`=`$ $`\delta _{ij}`$
$`𝐏_{ij}^{(1/2)}𝐏_{jk}^{(3/2)}`$ $`=`$ $`𝐏_{ij}^{(3/2)}𝐏_{jk}^{(1/2)}=0.`$
Here, summations over repeated indices are understood. $`G^{(3/2)}(\tau )`$ and $`G^{(1/2)}(\tau )`$ in Eq. (7) denote the spin 3/2 and 1/2 contributions to $`G(\tau )`$, respectively, which can be derived by operating $`𝐏^{(3/2)}`$ and $`𝐏^{(1/2)}`$ on $`G(\tau )`$, respectively. (In our practical lattice QCD calculation, we construct the Rarita-Schwinger correlator $`G_{ij}(\tau )`$ for $`i=1,2,3`$ and $`j=3`$(fixed), and multiply $`𝐏^{(3/2)}`$ from the left to obtain $`G^{(3/2)}(\tau )`$.)
In the asymptotic region ($`0\tau N_t`$), contaminations of excited states are suppressed. Considering the parity transformation property Eq. (6), $`G^{(3/2)}(\tau )`$ and $`G^{(1/2)}(\tau )`$ are expressed in this region as
$`G^{(3/2)}(\tau )`$
$`=`$ $`P_+\left\{|\lambda _{3/2^{}}|^2e^{\tau m_{3/2^{}}}+|\lambda _{3/2^+}|^2e^{(N_t\tau )m_{3/2^+}}\right\}`$
$``$ $`P_{}\left\{|\lambda _{3/2^+}|^2e^{\tau m_{3/2^+}}+|\lambda _{3/2^{}}|^2e^{(N_t\tau )m_{3/2^{}}}\right\}`$
$`G^{(1/2)}(\tau )`$
$`=`$ $`P_+\left\{|\lambda _{1/2^{}}|^2e^{\tau m_{1/2^{}}}+|\lambda _{1/2^+}|^2e^{(N_t\tau )m_{1/2^+}}\right\}`$
$``$ $`P_{}\left\{|\lambda _{1/2^+}|^2e^{\tau m_{1/2^+}}+|\lambda _{1/2^{}}|^2e^{(N_t\tau )m_{1/2^{}}}\right\},`$
where $`P_\pm (1\pm \gamma _4)/2`$ denote the projection matrices onto the “upper” and “lower” Dirac subspaces, respectively. $`m_{3/2^\pm }`$ and $`m_{1/2^\pm }`$ denote the lowest-lying masses in $`J^P=3/2^\pm `$ and $`1/2^\pm `$ channels, respectively. $`\lambda _{3/2^\pm }`$ and $`\lambda _{1/2^\pm }`$ represent the couplings to the interpolating field Eq. (II.1) with $`J^P=3/2^\pm `$ and $`1/2^\pm `$ states, respectively. In Eq. (II.2), we adopt the anti-periodic boundary condition along the temporal direction. A brief derivation of Eq. (8) and Eq. (II.2) is presented in Appendix. A. The forward propagation is dominant in the region $`0<\tau N_t/2`$, while the backward propagation is dominant in the region $`N_t/2\tau <N_t`$. To separate the negative (positive) parity contribution, we restrict ourselves to the region $`0<\tau N_t/2`$, and examine the “upper” (“lower”) Dirac component.
### II.3 Scattering states involved in 5Q spectrum
We consider the (two-particle) scattering states involved in 5Q spectrum. For $`J^P=3/2^\pm `$ iso-scalar penta-quarks, NK and NK scattering states play an important role. ($`\mathrm{\Delta }`$K does not couple to the iso-scalar channel.) These states are expressed as
$$|N(\stackrel{}{p},s)K(\stackrel{}{p}),|N(\stackrel{}{p},s)K^{}(\stackrel{}{p},i),$$
(12)
where $`s`$ and $`i`$ denote the spin of the nucleon and K, and $`\stackrel{}{p}`$ denotes the spatial momentum allowed for a particular choice of the spatial boundary condition adopted. For instance, if these hadrons are subject to the spatially periodic boundary condition, their momenta are quantized as
$$p_i=2n_i\pi /L,n_i\text{Z}\text{Z},$$
(13)
where $`L`$ denotes the spatial extent of the lattice. In contrast, if they are subject to the spatially anti-periodic boundary condition, their momenta are quantized as
$$p_i=(2n_i+1)\pi /L,n_i\text{Z}\text{Z}.$$
(14)
We first perform the parity projections. The positive and the negative parity states are obtained in the following way:
$`|NK(\pm )`$
$`=`$ $`|N(\stackrel{}{p},s)K(\stackrel{}{p})|N(\stackrel{}{p},s)K(\stackrel{}{p})`$
$`|NK^{}(\pm )`$
$`=`$ $`|N(\stackrel{}{p},s)K^{}(\stackrel{}{p},i)|N(\stackrel{}{p},s)K^{}(\stackrel{}{p},i).`$
Assuming that the interactions between N and K and between N and K are weak, their energies are approximated as
$`E_{NK}`$ $``$ $`\sqrt{m_N^2+\stackrel{}{p}^2}+\sqrt{m_K^2+\stackrel{}{p}^2}`$ (17)
$`E_{N^{}K}`$ $``$ $`\sqrt{m_N^2+\stackrel{}{p}^2}+\sqrt{m_K^{}^2+\stackrel{}{p}^2},`$ (18)
respectively. The scattering states which couple to $`J^P=3/2^\pm `$ penta-quarks are obtained as spin-3/2 projections of Eq. (II.3) and Eq. (II.3). The d-wave NK states and the s-wave NK states can couple to the $`J^P=3/2^{}`$ channel, while the p-wave NK states and NK states can couple to the $`J^P=3/2^+`$ channel.
The scattering states with vanishing spatial momentum $`\stackrel{}{p}=\stackrel{}{0}`$ are exceptional in the following sense. On the one hand, the positive parity states vanish, because the first terms coincides with the second terms in Eq. (II.3) and Eq. (II.3) in the right hand side. On the other hand, the negative parity states are constructed only from the spin degrees of freedom, i.e., the spin degrees of freedom of the nucleon in Eq. (II.3), and the spin degrees of freedoms of the nucleon and K in Eq. (II.3). By counting the degeneracy of the resulting states, it is straightforward to see that no d-wave states are contained, i.e., Eq. (II.3) gives only s-wave NK states in $`J^P=1/2^{}`$ channel, and that Eq. (II.3) gives only s-wave NK states in $`J^P=1/2^{}`$ and $`3/2^{}`$ channels.
### II.4 Hybrid boundary condition(HBC)
In order to determine whether a state of our concern is a compact 5Q resonance state or a scattering state of two particles, we use two distinct spatial boundary conditions(BC), i.e., the (standard) periodic BC(PBC) and the hybrid BC(HBC), which is recently proposed in LABEL:ishii12. In PBC, one imposes the spatially periodic BC on u,d and s-quarks. As a result, all the hadrons are subject to the periodic BC. In this case, due to Eq. (13), all hadrons can take zero-momentum, and the smallest non-vanishing momentum $`\stackrel{}{p}_{\mathrm{min}}`$ is of the form as
$$(\pm 2\pi /L,0,0),(0,\pm 2\pi /L,0),(0,0,\pm 2\pi /L),$$
(19)
which gives
$$|\stackrel{}{p}_{\mathrm{min}}^{\mathrm{PBC}}|=2\pi /L.$$
(20)
On the other hand, in HBC, we impose the spatially anti-periodic BC on u and d-quarks, whereas the spatially periodic BC is imposed on s-quark. Since N($`uud,udd`$), K($`u\overline{s},d\overline{s}`$) and K($`u\overline{s},d\overline{s}`$) contain odd numbers of u and d quarks, they are subject to the anti-periodic BC. Therefore, due to Eq. (14), N, K and K cannot have a vanishing momentum in HBC. The smallest possible momentum $`\stackrel{}{p}_{\mathrm{min}}`$ is of the form as
$$(\pm \pi /L,\pm \pi /L,\pm \pi /L).$$
(21)
Hence, its norm $`|\stackrel{}{p}_{\mathrm{min}}|`$ is expressed as
$$|\stackrel{}{p}_{\mathrm{min}}^{\mathrm{HBC}}|=\sqrt{3}\pi /L.$$
(22)
In contrast, $`\mathrm{\Theta }^+`$($`uudd\overline{s}`$) is subject to the spatially periodic BC, since it contains even number of u and d quarks. Therefore, $`\mathrm{\Theta }^+`$ can have the vanishing momentum.
Switching from PBC, HBC affects the low-lying two-particle scattering spectrum. A drastic change is expected in the s-wave NK channel. In PBC, the energy of the lowest NK state is given as
$$E_{\mathrm{min}}^{\mathrm{PBC}}(NK^{}(\text{s-wave}))m_N+m_K^{}.$$
(23)
In contrast, in HBC, since both N and K are required to have non-vanishing momenta $`|\stackrel{}{p}_{\mathrm{min}}|=\sqrt{3}\pi /L`$, the energy of the lowest NK state is raised up as
$`E_{\mathrm{min}}^{\mathrm{HBC}}(NK^{}(\text{s-wave}))`$
$``$ $`\sqrt{m_N^2+3\pi ^2/L^2}+\sqrt{m_K^{}^2+3\pi ^2/L^2}.`$
Note that the shift amounts typically to a few hundred MeV for $`L2`$ fm.
HBC affects NK(d-wave), NK(p-wave), NK(p-wave) as well. However, these changes are not as drastic as that in NK(s-wave), because they are induced by the minor change in the minimum momentum from $`|\stackrel{}{p}_{\mathrm{min}}|=2\pi /L`$ to $`\sqrt{3}\pi /L`$. In PBC, the energies of the lowest two-particles states are expressed as
$`E_{\mathrm{min}}^{\mathrm{PBC}}(NK(\text{p/d-wave}))`$
$``$ $`\sqrt{m_N^2+4\pi ^2/L^2}+\sqrt{m_K^2+4\pi ^2/L^2}`$
$`E_{\mathrm{min}}^{\mathrm{PBC}}(NK^{}(\text{p-wave}))`$
$``$ $`\sqrt{m_N^2+4\pi ^2/L^2}+\sqrt{m_K^{}^2+4\pi ^2/L^2}.`$
In HBC, they are shifted down as
$`E_{\mathrm{min}}^{\mathrm{HBC}}(NK(\text{p/d-wave}))`$
$``$ $`\sqrt{m_N^2+3\pi ^2/L^2}+\sqrt{m_K^2+3\pi ^2/L^2}`$
$`E_{\mathrm{min}}^{\mathrm{HBC}}(NK^{}(\text{p-wave}))`$
$``$ $`\sqrt{m_N^2+3\pi ^2/L^2}+\sqrt{m_K^{}^2+3\pi ^2/L^2}.`$
Numerical values of NK and NK thresholds for each hopping parameter in spatial lattice of the size $`L2.15`$ fm for both PBC and HBC are summarized in Table 1.
In contrast to the scattering states, HBC is not expected to affect a compact 5Q resonance $`\mathrm{\Theta }^+`$ so much. Since $`\mathrm{\Theta }^+(uudd\overline{s})`$ can have vanishing momentum also in HBC, the shift of the penta-quark mass $`m_{5\mathrm{Q}}`$ originates only from the change in its intrinsic structure. In this case, the shift is expected to be less significant than the shift induced by the kinematic reason as is the case in N, K, and K. Now our way to find a compact 5Q resonance state is to seek for such a state which is not affected by HBC.
## III Lattice actions and parameters
To generate gauge field configurations, we use the standard plaquette action on the anisotropic lattice of the size $`12^3\times 96`$ as
$`S_\mathrm{G}`$ $`=`$ $`{\displaystyle \frac{\beta }{N_c}}{\displaystyle \frac{1}{\gamma _\mathrm{G}}}{\displaystyle \underset{x,i<j3}{}}\text{Re}\text{Tr}\left\{1P_{ij}(x)\right\}`$
$`+`$ $`{\displaystyle \frac{\beta }{N_c}}\gamma _\mathrm{G}{\displaystyle \underset{x,i3}{}}\text{Re}\text{Tr}\left\{1P_{i4}(x)\right\},`$
where $`P_{\mu \nu }(x)\text{SU(3)}`$ denotes the plaquette operator in the $`\mu `$-$`\nu `$-plane. The lattice parameter and the bare anisotropy parameter are fixed as $`\beta 2N_c/g^2=5.75`$ and $`\gamma _\mathrm{G}=3.2552`$, respectively. These values are determined to reproduce the renormalized anisotropy as $`\xi a_s/a_t=4`$ klassen . Adopting the pseudo-heat-bath algorithm, we pick up gauge field configurations every 500 sweeps after skipping 10,000 sweeps for the thermalization. We use totally 1000 gauge field configurations to construct the temporal correlators. Note that the high statistics of $`N_{\mathrm{conf}}=1000`$ is quite essential for our study, because the 5Q correlators for spin 3/2 states are found to be rather noisy. In fact, a preliminary analysis with less statistics $`N_{\mathrm{conf}}500`$ leads to a spurious resonance-like state ishii320 . The lattice spacing is determined from the static quark potential adopting the Sommer parameter $`r_0^1=395`$ MeV ($`r_00.5`$ fm) as $`a_s^1=1.100(6)`$ GeV ($`a_\mathrm{s}0.18`$ fm). Note that the lattice size $`12^3\times 96`$ amounts to $`(2.15\text{fm})^3\times (4.30\text{fm})`$ in the physical unit.
We adopt the $`O(a)`$-improved Wilson (clover) action on the anisotropic lattice for quark fields $`\psi `$ and $`\overline{\psi }`$ as matsufuru
$`S_\mathrm{F}`$ $``$ $`{\displaystyle \underset{x,y}{}}\overline{\psi }(x)K(x,y)\psi (y),`$ (28)
$`K(x,y)`$ $``$ $`\delta _{x,y}\kappa _\mathrm{t}\{\text{}\begin{array}{c}(1\gamma _4)U_4(x)\delta _{x+\widehat{4},y}\hfill \\ +(1+\gamma _4)U_4^{}(x\widehat{4})\delta _{x\widehat{4},y}\text{}\}\hfill \end{array}`$ (31)
$``$ $`\kappa _\mathrm{s}{\displaystyle \underset{i}{}}\{\text{}\begin{array}{c}(r\gamma _i)U_i(x)\delta _{x+\widehat{i},y}\hfill \\ +(r+\gamma _i)U_i^{}(x\widehat{i})\delta _{x\widehat{i},y}\text{}\}\hfill \end{array}`$
$``$ $`\kappa _\mathrm{s}c_E{\displaystyle \underset{i}{}}\sigma _{i4}F_{i4}\delta _{x,y}r\kappa _sc_B{\displaystyle \underset{i<j}{}}\sigma _{ij}F_{ij}\delta _{x,y},`$
where $`\kappa _\mathrm{s}`$ and $`\kappa _\mathrm{t}`$ denote the hopping parameters for the spatial and the temporal directions, respectively. The field strength $`F_{\mu \nu }`$ is defined through the standard clover-leaf-type construction. $`r`$ denotes the Wilson parameter. $`c_E`$ and $`c_B`$ denote the clover coefficients. To achieve the tadpole improvement, the link variables are rescaled as $`U_i(x)U_i(x)/u_\mathrm{s}`$ and $`U_4(x)U_4(x)/u_\mathrm{t}`$, where $`u_\mathrm{s}`$ and $`u_\mathrm{t}`$ denote the mean-field values of the spatial and temporal link variables, respectively matsufuru ; nemoto . This is equivalent to the redefinition of the hopping parameters as the tadpole-improved ones (with tilde), i.e., $`\kappa _\mathrm{s}=\stackrel{~}{\kappa }_\mathrm{s}/u_\mathrm{s}`$ and $`\kappa _\mathrm{t}=\stackrel{~}{\kappa }_\mathrm{t}/u_\mathrm{t}`$. The anisotropy parameter is defined as $`\gamma _F\stackrel{~}{\kappa }_\mathrm{t}/\stackrel{~}{\kappa }_\mathrm{s}`$, which coincides with the renormalized anisotropy $`\xi =a_\mathrm{s}/a_\mathrm{t}`$ for sufficiently small quark mass at the tadpole-improved level matsufuru . For given $`\kappa _\mathrm{s}`$, the four parameters $`r`$, $`c_E`$, $`c_B`$ and $`\kappa _\mathrm{s}/\kappa _\mathrm{t}`$ should be, in principle, tuned so that “Lorentz symmetry” holds up to discretization errors of $`O(a^2)`$. Here, $`r`$, $`c_E`$ and $`c_B`$ are fixed by adopting the tadpole improved tree-level values as
$$r=\frac{1}{\xi },c_E=\frac{1}{u_\mathrm{s}u_\mathrm{t}^2},c_B=\frac{1}{u_\mathrm{s}^3}.$$
(35)
Only the value of $`\kappa _\mathrm{s}/\kappa _\mathrm{t}\left(=\gamma _F(u_\mathrm{s}/u_\mathrm{t})\right)`$ is tuned nonperturbatively by using the meson dispersion relation matsufuru . It is convenient to define $`\kappa `$ as
$$\frac{1}{\kappa }\frac{1}{\stackrel{~}{\kappa }_\mathrm{s}}2\left(\gamma _F3r4\right).$$
(36)
Then the bare quark mass is expressed as $`m_0=\frac{1}{2}(1/\kappa 8)`$ in the spatial lattice unit in the continuum limit. This $`\kappa `$ plays the role of the hopping parameter “$`\kappa `$” in the isotropic formulation. For detail, see Refs. nemoto ; matsufuru , where we take the lattice parameters. The values of the lattice parameters are summarized in Table 2.
We adopt four values of the hopping parameter as $`\kappa =0.1210,0.1220,0.1230`$ and $`0.1240`$, which correspond to $`m_\pi /m_\rho =0.81,0.78,0.73`$ and $`0.66`$, respectively. These values roughly cover the region $`m_sm2m_s`$. For temporal direction, we impose anti-periodic boundary condition on all the quark fields. For spatial directions, we impose periodic boundary condition on all the quarks, unless otherwise indicated. We refer to this boundary condition as “periodic BC (PBC)”.
By keeping $`\kappa _s=0.1240`$ fixed for $`s`$ quark, and by changing $`\kappa =0.12100.1240`$ for $`u`$ and $`d`$ quarks, we perform the chiral extrapolation to the physical quark mass region. In the following part of the paper, we will use
$$(\kappa _s,\kappa )=(0.1240,0.1220),$$
(37)
as a typical set of hopping parameters in presenting correlators and effective mass plots. For convenience, we summarize masses of $`\pi `$, $`\rho `$, K, K, N and N($`J^P=1/2^{}`$ baryon) for each hopping parameter $`\kappa `$ together with their values at the physical quark mass in Table 3. Here, the chiral extrapolations of these particles are performed by a linear function in $`m_\pi ^2`$. Unless otherwise indicated, we adopt the jackknife prescription to estimate the statistical errors.
We use a smeared source to enhance the low-lying spectra. More precisely, we employ spatially extended interpolating fields of the gaussian size $`\rho 0.4`$ fm, which is obtained by replacing the quark fields $`q(x)`$ in 5Q interpolating fields by the smeared quark fields $`q_{\mathrm{smear}}(x)`$ in the Coulomb gauge as
$$q_{\mathrm{smear}}(\tau ,\stackrel{}{x})𝒩\underset{\stackrel{}{y}}{}\mathrm{exp}\left\{\frac{|\stackrel{}{x}\stackrel{}{y}|^2}{2\rho ^2}\right\}q(\tau ,\stackrel{}{y}),$$
(38)
where $`𝒩`$ is an appropriate normalization factor. For a practical use, we extend Eq. (38) appropriately so as to fit a particular choice of the spatial boundary condition. In this paper, we present correlators with a smeared source and a point sink.
## IV Numerical results on 5Q spectrum
We present our lattice QCD results on 5Q spectrum in the standard periodic boundary condition(PBC) in this section.
### IV.1 $`J^P=3/2^{}`$ 5Q spectrum in PBC
We consider 5Q spectrum in $`J^P=3/2^{}`$ channel. In Fig. 1, we show the effective mass plots in $`J^P=3/2^{}`$ channel for three interpolating fields, i.e., (a) the NK-type, (b) the twisted NK-type, (c) a diquark-type. The dotted lines indicate the s-wave NK and the d-wave NK thresholds, which happen to coincide accidentally in Fig. 1 for the spatial lattice size $`L2.15`$ fm.
We define the effective mass as a function of $`\tau `$ by
$$m_{\mathrm{eff}}(\tau )\mathrm{log}\left(\frac{G^{(3/2)}(\tau )}{G^{(3/2)}(\tau +1)}\right),$$
(39)
where $`G^{(3/2)}(\tau )`$ denotes the temporal correlator. At sufficiently large $`\tau `$, the correlator is dominated by the lowest-lying state with energy $`m`$ as $`G^{(3/2)}(\tau )Ae^{m\tau }`$. Then Eq. (39) gives a constant as $`m_{\mathrm{eff}}(\tau )m`$. Thus a plateau in $`m_{\mathrm{eff}}(\tau )`$ indicates that the correlator is saturated by a single-state. In such cases, we can perform a single-exponential fit in the plateau region.
Fig. 1 (a) shows the effective mass plot for the NK-type interpolating field. In the region $`0\tau 24`$, the contamination of the higher spectral contributions are gradually reduced, which is indicated by the decreases in $`m_{\mathrm{eff}}(\tau )`$. There is a plateau in the interval $`25\tau 35`$, where a single-state is expected to dominate the 5Q correlator. Beyond $`\tau 36`$, the statistical error becomes large. In addition, the effect of the backward propagation becomes gradually more significant as $`\tau 48`$ is approached. Hence, we simply neglect the data for $`\tau 36`$, and perform the single-exponential fit in the region $`25\tau 35`$. We obtain $`m_{5\mathrm{Q}}=2.90(2)`$ GeV, which is denoted by the solid line. One sees that the 5Q states appears above the s-wave NK and the d-wave NK thresholds.
Fig. 1 (b) shows the effective mass plot for the twisted NK-type interpolating field. There is a plateau in the interval $`24\tau 35`$, where the single-exponential fit is performed leading to $`m_{5\mathrm{Q}}=2.89(1)`$ GeV. The 5Q state is again above the s-wave NK and the d-wave NK thresholds.
Fig. 1 (c) shows the effective mass plot for the diquark-type interpolating field. We see that the statistical error is too large to identify the plateau unambiguously. Hence, we do not perform the fit. Note that this plot is obtained by using $`N_{\mathrm{conf}}=1000`$ gauge configurations. A possible reason for such a large noise is that the interpolating field Eq. (3) does not survive the non-relativistic limit due to the vector diquark.
Now, we perform the chiral extrapolation. As mentioned before, we keep $`\kappa =0.1240`$ fixed for $`s`$-quark, and vary $`\kappa =0.12100.1240`$ for $`u`$ and $`d`$ quarks. Fig. 2 shows the 5Q masses in $`J^P=3/2^{}`$ channel against $`m_\pi ^2`$. Circles and boxes denote the data obtained from the NK-type and the twisted NK-type 5Q correlators, respectively. Note that they agree with each other within the statistical error. The open symbols refer to the direct lattice QCD data. Since these data behave almost linearly in $`m_\pi ^2`$, we adopt the linear chiral extrapolation in $`m_\pi ^2`$ to obtain $`m_{5\mathrm{Q}}`$ in the physical quark mass region. Note that the ordinary non-PS mesons and baryons show similar linearity in $`m_\pi ^2`$ nemoto . The closed symbols denote the results of the chiral extrapolation. We see that all the 5Q states appear above the s-wave NK and the d-wave NK thresholds. As a result of the chiral extrapolation, we obtain only massive 5Q states as $`m_{5\mathrm{Q}}=2.17(4),2.11(4)`$ GeV from the NK-type and the twisted NK-type correlators, respectively, which is too heavy to be identified with the experimentally observed $`\mathrm{\Theta }^+(1540)`$. Numerical values of $`m_{5\mathrm{Q}}`$ at each hopping parameter together with their chirally extrapolated values are summarized in Table 4. To obtain a low-lying state at $`m_{5\mathrm{Q}}1540`$ MeV, a 5Q state should appear below these thresholds at least in the light quark mass region. In this case, a significantly large chiral effect is required. Of course, this point can be in principle clarified by an explicit lattice QCD calculation with chiral fermions.
### IV.2 $`J^P=3/2^+`$ 5Q spectrum in PBC
We consider 5Q spectrum in $`J^P=3/2^+`$ channel. $`J^P=3/2^+`$ is an interesting quantum number from the view point of the diquark picture of Jaffe and Wilczekjaffe . In this picture, the pair of diquarks has angular momentum one, which is combined with the spin 1/2 of $`\overline{s}`$ quark. Hence, there are two possibilities as $`J^P=1/2^+`$ and $`3/2^+`$, i.e., the diquark picture can support $`J^P=3/2^+`$ possibility as well. Its mass splits from the $`J^P=1/2^+`$ state depending on a particular form of the LS-interaction. If it is massive, it is expected to have a large decay width. If it is light enough, its exotic structure may work to implement the narrow decay width as in $`J^P=1/2^+`$ case.
In Fig. 3, we show the 5Q effective mass plots in PBC employing three types of interpolating fields, i.e., (a) the NK-type, (b) the twisted NK-type, (c) the diquark-type. The dotted lines indicate the s-wave NK, the p-wave NK and the p-wave NK thresholds in the spatial lattice of the size $`L2.15`$ fm, respectively.
Fig. 3 (a) shows the 5Q effective mass plot employing the NK-type interpolating field. In the region, $`0\tau 17`$, the contaminations of higher spectral contributions become gradually reduced. There is a flat region $`18\tau 30`$, which is still afflicted with slightly large statistical errors. The single-exponential fit in this region gives $`m_{5\mathrm{Q}}=3.34(3)`$ GeV. Note that this value agrees with the s-wave NK threshold $`E_{\mathrm{th}}3.27`$ GeV. (See Table 3 for $`m_N^{}`$.)
Fig. 3 (b) shows the 5Q effective mass plot corresponding to the twisted NK-type interpolating field. We have a rather stable plateau in the interval $`21\tau 27`$, where the single-exponential fit is performed. We obtain $`m_{5\mathrm{Q}}=3.11(4)`$ GeV. The result is denoted by the solid line.
Fig. 3 (c) shows the 5Q effective mass plot for the diquark-type interpolating field. We find a plateau in the interval $`19\tau 29`$, where the single-exponential fit is performed. We obtain $`m_{5\mathrm{Q}}=3.16(2)`$ GeV, which is denoted by the solid line.
Now, we perform the chiral extrapolation. In Fig. 4, $`m_{5\mathrm{Q}}`$ is plotted against $`m_\pi ^2`$. Circles, boxes and triangles denote the data obtained from the NK-type, the twisted NK-type and the diquark-type 5Q correlators, respectively. Note that the latter two agree with each other within the statistical errors.
As a result of the chiral extrapolation, we obtain $`m_{5\mathrm{Q}}=2.64(7)`$ GeV from the NK-type correlator, $`m_{5\mathrm{Q}}=2.48(10)`$ GeV from the twisted NK-type correlator, and $`m_{5\mathrm{Q}}=2.42(6)`$ GeV from the diquark-type correlator. Numerical values of $`m_{5\mathrm{Q}}`$ in $`J^P=3/2^+`$ channel at each hopping parameter together with their chirally extrapolated values are also summarized in Table 4. The two data from the twisted NK-type and the diquark-type correlators are considered to be almost consistent with the p-wave NK threshold, while the data from the NK-type correlator seems to correspond to a more massive state, which is likely to be consistent with the NK(s-wave) threshold. We see again that all of our data of $`m_{5\mathrm{Q}}`$ appear above the NK threshold(p-wave), which is located above the artificially raised NK threshold(p-wave) due to the finiteness of the spatial lattice as $`L2.15`$ fm. As a result, we are left only with such massive 5Q states.
Now, several comments are in order. (1) LABEL:lasscock32 reported the existence of a low-lying 5Q state in $`J^P=3/2^+`$ channel using NK-type interpolating field. However, we have not observed such a low-lying 5Q state in our calculation. (2) Recall that, except for a single calculationchiu , lattice QCD calculations indicate that $`J^P=1/2^+`$ state is heavy scikor12 ; sasaki ; kentacky ; ishii12 ; rabbit ; lasscock12 ; alexandrou12 ; csikor122 ; holland , for instance $`m_{5\mathrm{Q}}2.25`$ GeV in LABEL:ishii12. From the viewpoint of the diquark picture, it could be natural to obtain such massive 5Q states in $`J^P=3/2^+`$ channel. If there were a low-lying 5Q state in $`J^P=3/2^+`$ channel, then the diquark picture could suggest also a low-lying 5Q state in $`J^P=1/2^+`$ channel nearby.
## V Analysis with HBC
In the previous section, we have only massive 5Q states, which are obtained by using the linear chiral extrapolation in $`m_\pi ^2`$. However, the chiral behavior may deviate from a simple linear one in the light quark mass region, which could lead to somewhat less massive states. Considering this, we think it of worth at this stage to analyze whether our 5Q states are compact 5Q resonances or not. This is done by switching the spatial periodic BC to the hybrid BC(HBC) introduced in Sect. II.
### V.1 $`J^P=3/2^{}`$ 5Q spectrum in HBC
Fig. 5 shows the 5Q effective mass plots in HBC employing the three types of interpolating fields, i.e., (a) the NK-type, (b) the twisted NK-type, and (c) the diquark type. These figures should be compared with their PBC counterparts in Fig. 1. The dotted lines denote the s-wave NK and the d-wave NK thresholds. For the typical set of hopping parameters, i.e., Eq. (37), the s-wave NK threshold(the thick dotted line) is raised up by $`180`$ MeV, and the d-wave NK threshold(the thin dotted line) is lowered down by $`70`$ MeV due to HBC in the finite spatial extent as $`L2.15`$ fm. (See Table 1.)
Fig. 5 (a) shows the 5Q effective mass plot for the NK-type interpolating field in HBC. We find a plateau in the interval $`23\tau 35`$, where the single-exponential fit is performed leading to $`m_{5\mathrm{Q}}=2.98(1)`$ GeV, which is denoted by the solid line. We see that $`m_{5\mathrm{Q}}`$ is raised up by 80 MeV due to HBC. The value of $`m_{5\mathrm{Q}}`$ is consistent with the s-wave NK threshold within the statistical error. Therefore, we regard this state as an NK scattering state.
Fig. 5 (b) shows the 5Q effective mass plot for the twisted NK-type interpolating field. We find a plateau in the interval $`24\tau 35`$, where the single exponential fit is performed leading to $`m_{5\mathrm{Q}}=2.98(1)`$ GeV, which is denoted by the solid line. The situation is similar to the NK-interpolating field case. We see that $`m_{5\mathrm{Q}}`$ is raised up by 90 MeV due to HBC. Since the value is consistent with the s-wave NK threshold within the statistical error, we regard it as an NK scattering state.
Fig. 5 (c) shows the 5Q effective mass plot for the diquark-type interpolating field. We see that it is afflicted with considerable size of statistical errors as before, due to which the best-fit is not performed.
In this way, all of our 5Q states in $`J^P=3/2^{}`$ channel turn out to be NK scattering states. More precisely, we do not observe any compact 5Q resonance states in $`J^P=3/2^{}`$ channel below the raised s-wave NK threshold, i.e., in the following region:
$$E\sqrt{m_N^2+\stackrel{}{p}_{\mathrm{min}}^2}+\sqrt{m_K^{}^2+\stackrel{}{p}_{\mathrm{min}}^2},$$
(40)
with $`|\stackrel{}{p}_{\mathrm{min}}|499`$ MeV.
### V.2 $`J^P=3/2^+`$ 5Q spectrum in HBC
Fig. 6 shows the 5Q effective mass plots in HBC employing the three types of interpolating fields, i.e., (a) the NK-type, (b) the twisted NK-type, and (c) the diquark-type. These figures should be compared with their PBC counterparts in Fig. 3. The meanings of the dotted and the solid lines are the same as in Fig. 3.
HBC may not be useful in $`J^P=3/2^+`$ channel, since it induces only minor changes in the two-particle spectra. For the typical set of hopping parameters, i.e., Eq. (37), the p-wave NK threshold is lowered down only by $`60`$ MeV, and the p-wave NK threshold is lowered down only by $`70`$ MeV. We see that these shifts are rather small. This is because they are induced by the changes in the minimum non-vanishing momentum, i.e., $`|\stackrel{}{p}_{\mathrm{min}}|=2\pi /L576`$ MeV to $`\sqrt{3}\pi /L499`$ MeV as mentioned before. In $`J^P=3/2^+`$ channel, NK(s-wave) threshold shows the most drastic change, i.e., the upper shift by 170 MeV, which however plays a less significant role, since its location is at rather high energy.
Fig. 6 (a) shows the 5Q effective mass plot employing the NK-type interpolating field. There is a flat region $`16\tau 25`$, which is still afflicted with slightly large statistical errors. The single-exponential fit in this region leads to $`m_{5\mathrm{Q}}=3.38(2)`$ GeV, which is denoted by the solid line. We see that $`m_{5\mathrm{Q}}`$ is raised up by 40 MeV. Although the shift of 40 MeV is rather small, $`m_{5\mathrm{Q}}`$ is again almost consistent with the s-wave NK threshold. Considering its rather large statistical error, this 5Q state is likely to be an s-wave NK scattering state. To draw a more solid conclusion on this state, it is necessary to improve the statistics further more.
Fig. 6 (b) shows the 5Q effective mass plot employing the twisted NK-type interpolating field. There is a plateau in the interval $`23\tau 31`$, where we perform the single-exponential fit. The result $`m_{5\mathrm{Q}}=3.02(3)`$ GeV is denoted by the solid line. We see that $`m_{5\mathrm{Q}}`$ is lowered down by 90 MeV, which is considered to be consistent with the shift of the NK(p-wave) threshold. Therefore, this state is likely to be an NK(p-wave) scattering state.
Fig. 6 (c) shows the 5Q effective mass plot employing the diquark-type interpolating field. Although the data is slightly noisy, there is a plateau in the interval $`23\tau 28`$. A single-exponential fit in this plateau region leads to $`m_{5\mathrm{Q}}=3.08(4)`$ GeV, which is denoted by the solid line. $`m_{5\mathrm{Q}}`$ is lowered down by 80 MeV due to HBC. The situation is similar to Fig. 6 (b). This state is likely to be an NK(p-wave) scattering state.
In this way, all of our 5Q states are likely to be either NK(s-wave) or NK(p-wave) states rather than compact 5Q resonance states. Of course, because HBC induces only minor changes in the 5Q spectrum in $`J^P=3/2^+`$ channel, and also because 5Q correlators still involve considerable size of statistical error, more statistics is desirable to draw a more solid conclusion on the real nature of these 5Q states. Here, we can at least state that these 5Q states are all massive, which locate above NK(p-wave) threshold.
## VI Summary and conclusion
We have studied $`J^P=3/2^\pm `$ penta-quark(5Q) baryons in anisotropic lattice QCD at the quenched level with a large number of gauge field configurations as $`N_{\mathrm{conf}}=1000`$ for high precision measurements. We emphasize that the spin of $`\mathrm{\Theta }^+(1540)`$ has not yet determined experimentally, and that the $`J^P=3/2^{}`$ assignment provides us with one of the possible solutions to the puzzle of the narrow decay width of $`\mathrm{\Theta }^+(1540)`$ hosaka32 . We have employed the standard Wilson gauge action on the anisotropic lattice of the size $`12^3\times 96`$ with the renormalized anisotropy $`a_\mathrm{s}/a_\mathrm{t}=4`$ at $`\beta =5.75`$, which leads to $`a_\mathrm{s}0.18`$ fm and $`a_\mathrm{t}0.045`$ fm. We have found that correlators of 5Q baryons in $`J^P=3/2^\pm `$ channels are rather noisy. Hence, the large statistics as $`N_{\mathrm{conf}}=1000`$ has played a key role to get a solid result in our calculation. For the quark part, we have employed $`O(a)`$-improved Wilson (clover) action with four values of the hopping parameters as $`\kappa =0.1210(0.0010)0.1240`$, which roughly cover the quark mass region as $`m_sm2m_s`$. To avoid the contaminations of higher spectral contributions, we have employed the spatially extended source in the 5Q correlators.
We have examined several types of the 5Q interpolating fields as (a) the NK-type, (b) the (color-)twisted NK-type, (c) the diquark-type. In $`J^P=3/2^{}`$ channel, there are plateaus in the effective mass plots for the NK-type and the twisted NK-type interpolating field, whereas no plateau has been identified in that for the diquark-type interpolating field due to the significantly large statistical error. The former two give almost identical results. We have employed the linear chiral extrapolations in $`m_\pi ^2`$, which have lead to $`m_{5\mathrm{Q}}2.17`$ and $`2.11`$ GeV for the NK-type and the twisted NK-type 5Q correlators, respectively. In $`J^P=3/2^+`$ channel, we have recognized plateaus in all the three effective mass plots. However, the plateau for the NK-type interpolating field is located at a somewhat higher energy than the other two. The chiral extrapolations have lead to $`m_{5\mathrm{Q}}2.64`$ GeV for the NK-type correlator, $`m_{5\mathrm{Q}}2.48`$ GeV for the twisted NK-type correlator, and $`m_{5\mathrm{Q}}2.42`$ GeV for the diquark-type correlator. In this way, our data have not supported low-lying 5Q states in both $`J^P=3/2^\pm `$ channels. All the 5Q states have been observed to appear above the d/p-wave NK threshold, which is artificially raised up by a few hundred MeV due to the finiteness of the spatial lattice as $`L2.15`$ fm. Note that, to obtain low-lying 5Q states in $`J^P=3/2^\pm `$ channel, a 5Q state should appear below the raised NK threshold(p/d-wave) at least in the light quark mass region.
In order to clarify whether the observed states are compact 5Q resonances or not, we have performed an analysis with hybrid boundary condition(HBC), which was recently proposed by LABEL:ishii12. In $`J^P=3/2^{}`$ channel, our 5Q states observed in the NK-type and the twisted NK-type correlators have turned out to be s-wave NK scattering states. In $`J^P=3/2^+`$ channel, for the twisted NK-type and the diquark-type correlators, the observed 5Q states are most likely to be NK(p-wave) scattering states. For the other one, i.e., the NK-type interpolating field, although more statistics is needed to draw a definite conclusion, it is most likely to be an s-wave NK scattering state. Note that, since HBC does not affected the two-particle spectra so much in $`J^P=3/2^+`$ channel, it is not easy to elucidate the natures of the 5Q states only with HBC. At any rate, whatever the real nature of these 5Q states may be, they result in a considerably massive states in the physical quark mass region, which cannot be identified as $`\mathrm{\Theta }^+(1540)`$ without involving a significantly large chiral contribution.
In this way, we have not obtained any relevant signals for low-lying compact 5Q resonance states in $`J^P=3/2^\pm `$ channel below 2.1 GeV in this paper, although the $`J^P=3/2^{}`$ possibility provides us with one of the possible solutions to the puzzle of the narrow decay width of $`\mathrm{\Theta }^+(1540)`$. To get more solid conclusion on the pentaquark, it is important to perform the systematic studies of the 5Q states with the various quantum numbers in lattice QCD with more sophisticated conditions. For instance, it is desired to use (1) unquenched full lattice QCD, (2) finer and larger volume lattice, (3) chiral fermion with small mass, (4) more sophisticated interpolating field corresponding to the diquark picture and so on. In any case, the mysterious exotic hadron of the pentaquark would be much clarified in future studies of lattice QCD as well as in the future experiments.
###### Acknowledgements.
We thank A. Hosaka, J. Sugiyama, T. Shinozaki for useful information and discussions. M. O and H. S are supported in part by Grant for Scientific Research ((B) No. 15340072 and (C) No. 16540236) from the Ministry of Education, Culture, Sports, Science and Technology, Japan. T. D. is supported by Special Postdoctoral Research Program of RIKEN. Y. N. is supported by 21st Century COE Program of Nagoya University. The lattice QCD Monte Carlo calculations have been performed on NEC-SX5 at Osaka University.
## Appendix A Spectral representation
Considering the importance of the parity determination of $`\mathrm{\Theta }^+`$, we present a brief derivation of the spectral representation of Rarita-Schwinger correlators, i.e., Eq. (8), with Eq. (II.2). In this section, gamma matrices are represented in Minkowskian form (See LABEL:itzykson). To avoid unnecessary complexities, we derive only $`J^P=3/2^\pm `$ parts. The $`J^P=1/2^\pm `$ parts can be obtained as a slight modification, on which we will make a comment at the end of the section.
We first consider the coupling of our interpolating fields $`\psi _\mu `$ to $`J^P=3/2^\pm `$ (anti-)baryon states. Due to Eq. (6), our interpolating fields, i.e., Eqs. (II.1),(II.1) and (3) have the negative intrinsic parity. Hence, their couplings to $`J^P=3/2^{}`$ (anti-)baryons are parameterized in the following way:
$`0|\psi _\mu (0)|B_{3/2^{}}(k,\alpha )`$ $`=`$ $`\lambda _{3/2^{}}u_\mu (m_{3/2^{}};k,\alpha )`$ (41)
$`0|\overline{\psi }_\mu (0)|\overline{B}_{3/2^{}}(k,\alpha )`$ $`=`$ $`\lambda _{3/2^{}}^{}\overline{v}_\mu (m_{3/2^{}};k,\alpha ),`$
where $`|B_{3/2^{}}(k,\alpha )`$ and $`|\overline{B}_{3/2^{}}(k,\alpha )`$ denote $`J^P=3/2^{}`$ (anti-)baryon states with momentum $`k`$, helicity $`\alpha `$, and mass $`m_{3/2^+}`$. $`u_\mu (m;k,\alpha )`$ and $`v_\mu (m;k,\alpha )`$ denote the Rarita-Schwinger spinors for $`J=3/2`$ particles with momentum $`k`$, helicity $`\alpha `$ and mass $`m`$ ioffe ; benmerrouche ; hemmert . Eqs. (II.1),(II.1) and (3) couples to $`J^P=3/2^+`$ (anti-)baryons as well. In this case, their couplings involve $`\gamma _5`$ in the following way:
$`0|\psi _\mu (0)|B_{3/2^+}(k,\alpha )`$ $`=`$ $`\lambda _{3/2^+}\gamma _5u_\mu (m_{3/2^+};k,\alpha )`$ (42)
$`0|\overline{\psi }_\mu (0)|\overline{B}_{3/2^+}(k,\alpha )`$ $`=`$ $`\lambda _{3/2^+}^{}\overline{v}_\mu (m_{3/2^+};k,\alpha )\gamma _5,`$
where $`|B_{3/2^+}(k,\alpha )`$ and $`|\overline{B}_{3/2^+}(k,\alpha )`$ denote $`J^P=3/2^+`$ (anti-)baryon states with momentum $`k`$, helicity $`\alpha `$ and mass $`m_{3/2^+}`$. “$``$” originates from the anti-commutativity of $`\gamma _0`$ and $`\gamma _5`$.
To derive the spectral representation, the best way would be to express it in the operator representation in the following way:
$$G_{\mu \nu }(\tau ,\stackrel{}{x})=Z^1\text{Tr}\left(e^{\beta H}T_\tau [\psi _\mu ^{}(\tau ,\stackrel{}{x})\overline{\psi }_\mu (0)]\right),$$
(43)
where $`\beta `$ denotes the temporal extent of the lattice, $`HH_{\mathrm{QCD}}`$ denotes the QCD Hamiltonian, $`Z\text{Tr}(e^{\beta H})`$ denotes the partition function, and $`T_\tau []`$ represents the time-ordered product along the imaginary time direction. The interpolating fields are represented in the Heisenberg picture in imaginary-time, i.e., $`\psi _\mu (\tau ,\stackrel{}{x})=e^{\tau H}\psi _\mu (0,\stackrel{}{x})e^{\tau H}`$ and $`\overline{\psi }_\mu (\tau ,\stackrel{}{x})=e^{\tau H}\overline{\psi }_\mu (0,\stackrel{}{x})e^{\tau H}`$. By restricting ourselves to the interval $`0\tau <\beta `$, Eq. (43) reduces to
$$G_{\mu \nu }(\tau ,\stackrel{}{x})=\text{Tr}\left(\frac{e^{\beta H}}{Z}\psi _\mu ^{}(\tau ,\stackrel{}{x})\overline{\psi }_\mu (0)\right).$$
(44)
Note that it can be equivalently expressed as
$$G_{\mu \nu }(\tau ,\stackrel{}{x})=\text{Tr}\left(\psi _\mu ^{}(\tau \beta ,\stackrel{}{x})\frac{e^{\beta H}}{Z}\overline{\psi }_\mu (0)\right).$$
(45)
In the large $`\beta `$ limit, we have $`e^{\beta H}/Z|00|`$, which is inserted into Eq. (44) and Eq. (45). Note that the resulting two expressions serve as independent contributions to the original “Tr”, i.e., Eq. (44) (or Eq. (45)). Hence, we keep these two contributions to obtain
$`G_{\mu \nu }(\tau ,\stackrel{}{x})`$
$``$ $`0|\psi _\mu ^{}(\tau ,\stackrel{}{x})\overline{\psi }_\mu (0)|0+0|\overline{\psi }_\mu (0)\psi _\mu ^{}(\tau \beta ,\stackrel{}{x})|0.`$
Note that the 1st term corresponds to the forward propagation, whereas the 2nd term to the backward propagation. By inserting single-(anti-)baryon intermediate states, and by using Eq. (41) and Eq. (42), we are left with
$`G_{\mu \nu }(\tau ,\stackrel{}{x})`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{m_{3/2^+}}{k_0}e^{\tau k_0}0|\psi _\mu ^{}(\stackrel{}{x})|B_{3/2^+}(k,\alpha )B_{3/2^+}(k,\alpha )|\overline{\psi }_\mu (0)|0}+\mathrm{}.`$
$`=`$ $`|\lambda _{3/2^{}}|^2{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{m_{3/2^{}}}{k_0}(1)P_{\mu \nu }^{(3/2)}(k)\left[e^{\tau k_0}e^{i\stackrel{}{k}\stackrel{}{x}}\left(\frac{m_{3/2^{}}+\overline{)}k}{2m_{3/2^{}}}\right)+e^{(\beta \tau )k_0}e^{i\stackrel{}{k}\stackrel{}{x}}\left(\frac{m_{3/2^{}}\overline{)}k}{2m_{3/2^{}}}\right)\right]}`$
$``$ $`|\lambda _{3/2^+}|^2{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{m_{3/2^+}}{k_0}(1)P_{\mu \nu }^{(3/2)}(k)\left[e^{\tau k_0}e^{i\stackrel{}{k}\stackrel{}{x}}\left(\frac{m_{3/2^+}\overline{)}k}{2m_{3/2^+}}\right)+e^{(\beta \tau )k_0}e^{i\stackrel{}{k}\stackrel{}{x}}\left(\frac{m_{3/2^+}+\overline{)}k}{2m_{3/2^+}}\right)\right]},`$
where $`k_0\sqrt{m_{3/2^{}}^2+\stackrel{}{k}^2}`$ for $`J^P=3/2^{}`$, $`k_0\sqrt{m_{3/2^+}^2+\stackrel{}{k}^2}`$ for $`J^P=3/2^+`$, and the following identities are used.
$`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}u_\mu (m;k,\alpha )\overline{u}_\nu (m;k,\alpha )`$ $`=`$ $`{\displaystyle \frac{m+\overline{)}k}{2m}}P_{\mu \nu }^{(3/2)}(k)`$ (48)
$`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}v_\mu (m;k,\alpha )\overline{v}_\nu (m;k,\alpha )`$ $`=`$ $`{\displaystyle \frac{m\overline{)}k}{2m}}P_{\mu \nu }^{(3/2)}(k),`$
where $`P_{\mu \nu }^{(3/2)}(k)`$ is the spin 3/2 projection operator defined as
$$P_{\mu \nu }^{(3/2)}(k)g_{\mu \nu }\frac{1}{3}\gamma _\mu \gamma _\nu \frac{1}{3k^2}(\overline{)}k\gamma _\mu k_\nu +k_\mu \gamma _\nu \overline{)}k).$$
(49)
By performing the integration over $`\stackrel{}{x}`$ for zero-momentum projection, and by replacing the Minkowskian gamma matrices by their Euclidean counterparts, we finally arrive at the spectral representation (Eq. (8) with Eq. (II.2)).
The derivation of the spin 1/2 parts is obtained by repeating a similar procedure using the following parameterizations instead of Eq. (41) and Eq. (42) as
$`0|\psi _\mu (0)|B_{1/2^{}}(k,\alpha )`$
$`=`$ $`\left(\lambda _{1/2^{}}\gamma _\mu +\lambda _{1/2^{}}^{}k_\mu \right)u(m_{1/2^{}};k,\alpha )`$
$`0|\overline{\psi }_\mu (0)|\overline{B}_{1/2^{}}(k,\alpha )`$
$`=`$ $`\overline{v}(m_{1/2^{}};k,\alpha )\left(\lambda _{1/2^{}}^{}\gamma _\mu \lambda _{}^{}{}_{1/2^{}}{}^{}k_\mu \right)`$
$`0|\psi _\mu (0)|B_{1/2^+}(k,\alpha )`$
$`=`$ $`\left(\lambda _{1/2^+}\gamma _\mu +\lambda _{1/2^{}}^{}k_\mu \right)\gamma _5u(m_{1/2^+};k,\alpha )`$
$`0|\overline{\psi }_\mu (0)|\overline{B}_{1/2^{}}(k,\alpha )`$
$`=`$ $`\overline{v}(m_{1/2^{}};k,\alpha )\gamma _5\left(\lambda _{1/2^{}}^{}\gamma _\mu \lambda _{}^{}{}_{1/2^{}}{}^{}k_\mu \right),`$
where $`|B_{1/2^\pm }(k,\alpha )`$ and $`|\overline{B}_{1/2^\pm }(k,\alpha )`$ denote the $`J^P=1/2^\pm `$ (anti-)baryon states with momentum $`k`$, helicity $`\alpha `$ and mass $`m_{1/2^\pm }`$. $`u(m;k,\alpha )`$ and $`v(m;k,\alpha )`$ denote the Dirac bispinors for spin 1/2 particles with mass $`m`$, momentum $`k`$ and helicity $`\alpha `$. $`\lambda _{1/2^\pm }`$ and $`\lambda _{1/2^\pm }^{}`$ represent the couplings to $`J^P=1/2^\pm `$ (anti-)baryons.
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# 1 Introduction
## 1 Introduction
That none of the terms in the Standard Model (SM) violate lepton number (L) is not due to an imposed symmetry, but merely reflects the fact that all such combinations of SM fields are ruled out by consideration of gauge invariance and renormalisability . For supersymmetric extensions of the SM this is no longer true. In the Minimal Supersymmetric Standard Model (MSSM) , lepton number violating terms (and baryon number (B) violating terms) appear naturally, giving rise to tree-level processes, proton decay for example, which are already strongly constrained by experiment. Either, bounds can be set on Lagrangian parameters, or a further discrete symmetry can be imposed on the Lagrangian, such that these processes are absent. The discrete symmetry most commonly imposed is known as R-parity $`(R_P)`$ . Under R-parity the particles of the Standard Model including the scalar Higgs fields are even, while all their superpartners are odd. Imposing this symmetry has a number of effects. Firstly, any interaction terms which violate lepton number or baryon number will not appear. Secondly, the decay of the lightest supersymmetric particle (LSP) into SM particles would violate $`R_P`$; the LSP is therefore stable. The sneutrino vacuum expectation values (VEVs) are zero; without extending the MSSM field content, spontaneous generation of $`R_P`$ violating terms is phenomenologically discounted .
If $`R_P`$ conservation is not imposed, fields with different $`R_P`$ mix . In particular, the neutrinos will mix with the neutralinos and the sneutrinos will mix with the neutral scalar Higgs fields; all five complex neutral scalar fields can acquire vacuum expectation values. Minimising this ten parameter potential in general is not straightforward, it is more convenient to simplify the system by choosing an appropriate basis in the neutral scalar sector. As one of the Higgs doublets carries the same quantum numbers as the lepton doublets (apart from the non-conserved lepton number), it is convenient to introduce the notation $`_\alpha =(H_1,L_i)`$ where $`H_1`$ and $`L_i`$ are the chiral superfields containing one Higgs doublet and the leptons, respectively ($`\alpha =0,\mathrm{},3`$ and $`i=1,\mathrm{},3`$). Furthermore, starting from the interaction basis, we are free to rotate the fields and choose the direction corresponding to that of the “Higgs” field. Assuming that the system defining the five complex vacuum expectation values of the fields was solved, four complex VEVs $`v_\alpha `$ would define a direction in the four dimensional $`(H_1,L_i)`$ space. One can then choose the basis vector which defines the Higgs fields to point in the direction defined by the vacuum expectation values. We refer to this basis, in which the “sneutrino” (as we call the fields perpendicular to the “Higgs” field) VEVs are zero, as the vanishing sneutrino VEV basis . This basis has the virtue of simplifying the mass matrices and vertices of the theory and thus is better suited for calculations.
Basis independent parameterisations can be chosen which explicitly show the amount of physical lepton number violation . Values for physical observables such as sneutrino masses and mass splitting between CP-even and CP-odd sneutrinos have been derived in the literature in terms of these combinations but usually under some approximations (for example the number of generations or CP-conservation). We find this procedure in general complicated for practical applications and we shall not adopt it here.
Instead, we present in the next section a calculable procedure for framing the most general MSSM scalar potential in the vanishing sneutrino VEV basis. An advantage of our procedure is to obtain a diagonal “slepton” mass matrix, two real non-zero vacuum expectation values and real parameters of the neutral scalar potential in the rotated basis. The latter proves that the neutral scalar sector of the most general R-parity violating MSSM exhibits neither spontaneous nor explicit CP-violation in agreement with . In Section 3, the tree-level masses and mixing of the neutral scalar sector is investigated. Using the Courant-Fischer theorem for the interlaced eigenvalues, we prove that there is always at least one neutral scalar which is lighter than the $`Z`$-gauge boson. We present approximate formulae which relate the Higgs masses, mixing and Higgs-gauge boson vertices of the R-parity conserving (RPC) case with the R-parity violating (RPV) one. In Section 4, the positiveness of the scalar mass matrices and stability of the vacuum is discussed.
## 2 Basis choice in the neutral scalar sector
In this section we develop a procedure connecting a general neutral scalar basis with the vanishing sneutrino VEV basis, the latter being more convenient for practical applications. The most general, renormalizable, gauge invariant superpotential that contains the minimal content of fields, is given by
$`𝒲`$ $`=`$ $`ϵ_{ab}\left[{\displaystyle \frac{1}{2}}\lambda _{\alpha \beta k}_\alpha ^a_\beta ^b\overline{E}_k+\lambda _{\alpha jk}^{}_\alpha ^aQ_j^{bx}\overline{D}_{kx}\mu _\alpha _\alpha ^aH_2^b+(Y_U)_{ij}Q_i^{ax}H_2^b\overline{U}_{jx}\right]`$ (2.1)
$`+`$ $`{\displaystyle \frac{1}{2}}ϵ_{xyz}\lambda _{ijk}^{\prime \prime }\overline{U}_i^x\overline{D}_j^y\overline{D}_k^z,`$
where $`Q_i^{ax},\overline{D}_i^x,\overline{U}_i^x,_i^a,\overline{E}_i,H_1^a,H_2^a`$ are the chiral superfield particle content, $`i=1,2,3`$ is a generation index, $`x=1,2,3`$ and $`a=1,2`$ are $`SU(3)`$ and $`SU(2)`$ gauge indices, respectively. The simple form of (2.1) results when combining the chiral doublet superfields with common hypercharge $`Y=\frac{1}{2}`$ into $`_{\alpha =0,\mathrm{},3}^a=(H_1^a,L_{i=1,2,3}^a)`$. $`\mu _\alpha `$ is the generalized dimensionful $`\mu `$-parameter, and $`\lambda _{\alpha \beta k},\lambda _{\alpha jk}^{},\lambda _{ijk}^{\prime \prime },(Y_U)_{ij}`$ are Yukawa matrices with $`ϵ_{ab}`$ and $`ϵ_{xyz}`$ being the totally anti-symmetric tensors, with $`ϵ_{12}=ϵ_{123}=+1`$. Then the five neutral scalar fields, $`\stackrel{~}{\nu }_{L\alpha },h_2^0`$ from the $`SU(2)`$ doublets, $`_\alpha =(\stackrel{~}{\nu }_{L\alpha },\stackrel{~}{e}_{L\alpha }^{})^T`$ and $`H_2=(h_2^+,h_2^0)^T`$, form the most general neutral scalar potential of the MSSM,
$`V_{\text{neutral}}`$ $`=`$ $`\left(m_\stackrel{~}{}^2\right)_{\alpha \beta }\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\beta }^{}+\mu _\alpha ^{}\mu _\beta \stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\beta }^{}+\mu _\alpha ^{}\mu _\alpha h_2^0h_2^0+m_{H_2}^2h_2^0h_2^0`$ (2.2)
$``$ $`b_\alpha \stackrel{~}{\nu }_{L\alpha }^{}h_2^0b_\alpha ^{}\stackrel{~}{\nu }_{L\alpha }^{}h_2^0+{\displaystyle \frac{1}{8}}(g^2+g_2^2)[h_2^0h_2^0\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\alpha }^{}]^2,`$
where general complex parameters $`b_\alpha `$, an hermitian matrix $`\left(m_\stackrel{~}{}^2\right)_{\alpha \beta }`$ and $`m_{H_2}^2`$ all arise from the supersymmetry breaking sector of the theory. The last term in (2.2) originates from the D-term contributions of the superfields $`_\alpha ,H_2`$. Defining
$`\left(_\stackrel{~}{}^2\right)_{\alpha \beta }\left(m_\stackrel{~}{}^2\right)_{\alpha \beta }+\mu _\alpha ^{}\mu _\beta ,\mathrm{and}m_2^2m_{H_2}^2+\mu _\alpha ^{}\mu _\alpha ,`$ (2.3)
one can rewrite the potential in (2.2) in a compact form as
$`V_{\text{neutral}}`$ $`=`$ $`\left(_\stackrel{~}{}^2\right)_{\alpha \beta }\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\beta }^{}+m_2^2h_2^0h_2^0(b_\alpha \stackrel{~}{\nu }_{L\alpha }^{}h_2^0+\mathrm{H}.\mathrm{c})`$ (2.4)
$`+`$ $`{\displaystyle \frac{1}{8}}(g^2+g_2^2)[h_2^0h_2^0\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\alpha }^{}]^2.`$
In order to go to the vanishing sneutrino VEV basis, we redefine the “Higgs-sneutrino” fields
$`\stackrel{~}{\nu }_{L\alpha }^{}=U_{\alpha \beta }\stackrel{~}{\nu }_{L\beta }^{^{}},`$ (2.5)
where $`𝐔`$ is a $`4\times 4`$ unitary matrix
$`𝐔=𝐕\mathrm{diag}(e^{i\varphi _\alpha })𝐙,`$ (2.6)
being composed of three other matrices which we define below, $`𝐕`$ unitary and $`𝐙`$ real orthogonal. The potential in the primed basis becomes,
$`V_{\text{neutral}}`$ $`=`$ $`\left[Z^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)Z\right]_{\alpha \beta }\stackrel{~}{\nu }_{L\alpha }^{{}_{}{}^{}}\stackrel{~}{\nu }_{L\beta }^{^{}}+m_2^2h_2^0h_2^0`$ (2.7)
$``$ $`[(b^{}Z)_\alpha \stackrel{~}{\nu }_{L\alpha }^{^{}}h_2^0+\mathrm{H}.\mathrm{c}]+{\displaystyle \frac{1}{8}}(g^2+g_2^2)(h_2^0h_2^0\stackrel{~}{\nu }_{L\alpha }^{{}_{}{}^{}}\stackrel{~}{\nu }_{L\alpha }^{^{}})^2,`$
where
$`\left(\widehat{^{}}_\stackrel{~}{}^2\right)=\mathrm{diag}(e^{i\varphi _\alpha })𝐕^{}\left(_\stackrel{~}{}^2\right)𝐕\mathrm{diag}(e^{i\varphi _\alpha }),b^{}_{}{}^{}T=b^T𝐕\mathrm{diag}(e^{i\varphi _\alpha }).`$ (2.8)
The unitary matrix, $`𝐕`$, is chosen such that $`\left(\widehat{^{}}_\stackrel{~}{}^2\right)`$ is real and diagonal - the hat $`(\widehat{})`$ is used to denote a diagonal matrix. The phases $`\varphi _\alpha `$ are chosen such that $`b_\alpha ^{}`$ is real and positive \[they are equal to the phases of $`(b^TV)_\alpha ^{}`$\]. The minimisation conditions for the scalar fields are now derived, to obtain conditions for the vacuum expectation values,
$`{\displaystyle \frac{V}{\stackrel{~}{\nu }_{L\alpha }^{{}_{}{}^{}}}}|_{\text{vac}}`$ $`=`$ $`\left[Z^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)Z\right]_{\alpha \beta }\stackrel{~}{\nu }_{L\beta }^{^{}}(b^{}Z)_\alpha h_2^0{\displaystyle \frac{1}{4}}(g^2+g_2^2)\left(h_2^0h_2^0\stackrel{~}{\nu }_{L\gamma }^{{}_{}{}^{}}\stackrel{~}{\nu }_{L\gamma }^{^{}}\right)\stackrel{~}{\nu }_{L\alpha }^{^{}}|_{\text{vac}}=0,`$
$`{\displaystyle \frac{V}{h_2^0}}|_{\text{vac}}`$ $`=`$ $`m_2^2h_2^0(b^{}Z)_\alpha \stackrel{~}{\nu }_{L\alpha }^{{}_{}{}^{}}+{\displaystyle \frac{1}{4}}(g^2+g_2^2)\left(h_2^0h_2^0\stackrel{~}{\nu }_{L\gamma }^{{}_{}{}^{}}\stackrel{~}{\nu }_{L\gamma }^{^{}}\right)h_2^0|_{\text{vac}}=0,`$ (2.9)
where “vac” indicates that the fields have to be replaced by their VEVs,
$`\stackrel{~}{\nu }_{L\alpha }^{^{}}={\displaystyle \frac{v_\alpha }{\sqrt{2}}},h_2^0={\displaystyle \frac{v_u}{\sqrt{2}}}.`$ (2.10)
The $`U(1)_Y`$ symmetry of the unbroken Lagrangian was used to set the phase of $`v_u`$ to zero, however, at this stage all other vacuum expectation values will be treated as complex variables. By combining Eqs. (2.9,2.10) we obtain
$`\left[Z^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)Z\right]_{\alpha \beta }v_\beta (b^{}Z)_\alpha v_u{\displaystyle \frac{1}{8}}(g^2+g_2^2)(v_u^2v_\gamma ^{}v_\gamma )v_\alpha `$ $`=`$ $`0,`$ (2.11)
$`m_2^2v_u(b^{}Z)_\alpha v_\alpha ^{}+{\displaystyle \frac{1}{8}}(g^2+g_2^2)(v_u^2v_\gamma ^{}v_\gamma )v_u`$ $`=`$ $`0.`$ (2.12)
In a general basis, it is difficult to solve the above system with respect to the VEVs without making some approximations, for example assuming small “sneutrino” VEVs . In order to simplify calculations we would like to find a basis where the “sneutrino” VEVs vanish, $`v_1=v_2=v_3=0`$. In other words, we are seeking an orthogonal matrix $`Z`$, such that the following equation,
$`\left[Z^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)Z\right]_{\alpha 0}v_0(b^{}Z)_\alpha v_u{\displaystyle \frac{1}{2}}M_Z^2{\displaystyle \frac{v_u^2v_0^2}{v_u^2+v_0^2}}v_0\delta _{0\alpha }`$ $`=`$ $`0,`$ (2.13)
holds. If the above system is satisfied, then a solution with zero “sneutrino” VEVs exists. The other solutions, with non-vanishing “sneutrino” VEVs will be discussed later. In Eq. (2.13),
$`M_Z^2={\displaystyle \frac{1}{4}}\left(g^2+g_2^2\right)\left(v_u^2+v_0^2\right),`$ (2.14)
is the Z-gauge boson mass squared. It is obvious that when $`v_i=0`$, $`v_0`$ is real. It is now useful to define
$`\mathrm{tan}\beta {\displaystyle \frac{v_u}{v_0}}.`$ (2.15)
To determine $`𝐙`$, multiplying (2.13) by $`Z_{\gamma \alpha }`$, summing over $`\alpha `$ and solving for $`Z_{\alpha 0}`$, yields,
$`Z_{\alpha 0}={\displaystyle \frac{b_\alpha ^{}\mathrm{tan}\beta }{\left(\widehat{^{}}_\stackrel{~}{}^2\right)_{\alpha \alpha }\frac{1}{2}M_Z^2\frac{\mathrm{tan}^2\beta 1}{\mathrm{tan}^2\beta +1}}}.`$ (2.16)
For given set of model parameters, $`Z_{\alpha 0}`$ depends only on $`\mathrm{tan}\beta `$ which we can now fix by solving the orthonormality condition,
$`{\displaystyle \underset{\alpha =0}{\overset{3}{}}}Z_{\alpha 0}Z_{\alpha 0}={\displaystyle \underset{\alpha =0}{\overset{3}{}}}{\displaystyle \frac{b_\alpha ^{}_{}{}^{}2\mathrm{tan}^2\beta }{\left[\left(\widehat{^{}}_\stackrel{~}{}^2\right)_{\alpha \alpha }\frac{1}{2}M_Z^2\frac{\mathrm{tan}^2\beta 1}{\mathrm{tan}^2\beta +1}\right]^2}}=1.`$ (2.17)
This equation can be easily be solved numerically for any given set of model parameters.
It is worth noting that when $`b_i=0`$ and using notation more typical for this case, $`b_0^{}m_{12}^2`$, $`\left(\widehat{^{}}_\stackrel{~}{}^2\right)_{00}m_1^2`$, Eq. (2.17) reduces to one of the standard RPC MSSM equations for the Higgs VEVs:
$`m_{12}^2v_d=v_u\left[m_1^2{\displaystyle \frac{1}{8}}(g^2+g_2^2)(v_u^2v_d^2)\right].`$ (2.18)
For some parameter choices Eq. (2.17) may admit multiple solutions for $`\mathrm{tan}\beta `$. Each of the possible $`\mathrm{tan}\beta `$ specify a different basis, and each of these bases has one solution of the minimisation conditions with vanishing “sneutrino” VEVs. The subtlety highlighted earlier is the following: all possible solutions of the minimisation conditions can be found in each basis, so, in general, each basis contains a number of extrema equal to the number of possible solutions for $`\mathrm{tan}\beta `$. Hence, a solution with $`v_i=0`$ in one basis, is a solution with $`v_i0`$ in another basis. The important point to note is that by considering all possible values of $`\mathrm{tan}\beta `$, and selecting the value which corresponds to the deepest minima for the solution with vanishing sneutrino VEVs, all the solutions will have been accounted for, and the vanishing sneutrino VEV basis will have been determined correctly. The value of the potential at the vacuum, in terms of $`\mathrm{tan}\beta `$ is given by
$`V(\mathrm{tan}\beta )={\displaystyle \frac{M_Z^4}{2(g^2+g_2^2)}}\left({\displaystyle \frac{\mathrm{tan}^2\beta 1}{\mathrm{tan}^2\beta +1}}\right)^2.`$ (2.19)
The obvious conclusion from the equation above is that the deepest minimum of the potential is given by the solution for $`\mathrm{tan}\beta `$ or $`\mathrm{cot}\beta `$ which is greatest.
Knowing $`\mathrm{tan}\beta `$, one should fix $`m_2^2`$ using Eqs. (2.12,2.14-2.16) (again in the analogy with RPC MSSM where $`m_2^2`$ is usually given in terms of $`M_A,\mathrm{tan}\beta `$). Namely
$`m_2^2=Z_{\alpha 0}b_\alpha ^{}\mathrm{cot}\beta {\displaystyle \frac{1}{2}}M_Z^2{\displaystyle \frac{\mathrm{tan}^2\beta 1}{\mathrm{tan}^2\beta +1}}.`$ (2.20)
In this way $`m_2^2`$ is chosen to give the correct value of the the $`Z`$-boson mass.
Only the first column of the $`Z`$ matrix, $`Z_{\alpha 0}`$, is defined by Eq. (2.16). The remaining elements of $`𝐙`$ must still be determined. Having fixed the first column of the matrix, the other three columns can be chosen to be orthogonal to the first column and to each other. This leaves us with an $`O(3)`$ invariant subspace, such that the matrix $`𝐙`$ is given by
$`𝐙=𝐎\left(\begin{array}{cc}1& 0\\ 0& 𝐗_{3\times 3}\end{array}\right),`$ (2.23)
where
$`O=\left(\begin{array}{cccc}Z_{00}& \sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}& 0& 0\\ Z_{10}& \frac{Z_{00}Z_{10}}{\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& \frac{\sqrt{Z_{20}^2+Z_{30}^3}}{\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& 0\\ Z_{20}& \frac{Z_{00}Z_{20}}{\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& \frac{Z_{10}Z_{20}}{\sqrt{Z_{20}^2+Z_{30}^3}\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& \frac{Z_{30}}{\sqrt{Z_{20}^2+Z_{30}^3}}\\ Z_{30}& \frac{Z_{00}Z_{30}}{\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& \frac{Z_{10}Z_{30}}{\sqrt{Z_{20}^2+Z_{30}^3}\sqrt{Z_{10}^2+Z_{20}^2+Z_{30}^3}}& \frac{Z_{20}}{\sqrt{Z_{20}^2+Z_{30}^3}}\end{array}\right),`$ (2.28)
and $`𝐗`$ is an, as yet, undetermined $`3\times 3`$ orthogonal matrix determined by three angles. This remaining freedom can be used to diagonalise $`\left[𝐙^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)𝐙\right]_{ij}`$, i.e. the (real symmetric) “sneutrino” part of the $`𝐙^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)𝐙`$ matrix, with entries $`(\widehat{M}_{\stackrel{~}{\mathrm{L}}}^2)_i`$. We have now accomplished our aim of finding the matrices $`𝐕`$ and $`𝐙`$ which, after inserting into potential of Eq. (2.7) and dropping the primes, reduce the scalar potential to the form
$`V_{\text{neutral}}`$ $`=`$ $`(M_{\stackrel{~}{\mathrm{L}}}^2)_{\alpha \beta }\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\beta }^{}+m_2^2h_2^0h_2^0[B_\alpha \stackrel{~}{\nu }_{L\alpha }^{}h_2^0+\mathrm{H}.\mathrm{c}]`$ (2.29)
$`+`$ $`{\displaystyle \frac{1}{8}}(g^2+g_2^2)\left(h_2^0h_2^0\stackrel{~}{\nu }_{L\alpha }^{}\stackrel{~}{\nu }_{L\alpha }^{}\right)^2,`$
where
$`(M_{\stackrel{~}{\mathrm{L}}}^2)_{\alpha \beta }\left[Z^T\left(\widehat{^{}}_\stackrel{~}{}^2\right)Z\right]_{\alpha \beta }\mathrm{and}B_\alpha (b^{}Z)_\alpha ,`$ (2.30)
with $`\left(\widehat{^{}}_\stackrel{~}{}^2\right)`$ and $`b^{}`$ given by Eq. (2.8). In this basis the matrix $`𝐌_{\stackrel{~}{\mathrm{L}}}^\mathrm{𝟐}`$ adopts a particularly simple form
$`(M_{\stackrel{~}{\mathrm{L}}}^2)_{\alpha \beta }=\left(\begin{array}{cc}B_0\mathrm{tan}\beta \frac{1}{2}M_Z^2\mathrm{cos}2\beta & B_j\mathrm{tan}\beta \\ B_i\mathrm{tan}\beta & (\widehat{M}_{\stackrel{~}{\mathrm{L}}}^2)_i\delta _{ij}\end{array}\right),`$ (2.33)
where there is no sum over $`i`$ in the down-right part of the matrix. Notice that we did not only succeed to self consistently go to a basis where the sneutrino VEVs are zero, but also we managed to have the sneutrino masses $`(\widehat{M}_{\stackrel{~}{\mathrm{L}}}^2)_i`$ diagonal and all the parameters of the scalar potential in Eq. (2.29) real.
As a byproduct of our procedure, we denote here that the potential of Eq. (2.29) exhibits neither spontaneous nor explicit CP-violation at the tree level. The latter is in agreement with the results of Ref. following a different method. Of course, the parameters $`\mu _\alpha `$ of the superpotential and the soft supersymmetry breaking couplings stay in general complex. The result that the neutral scalar potential is CP invariant can also be seen directly from Eq. (2.4). By forming the complex basis $`(\stackrel{~}{\nu }_{L\alpha }^{},h_2^0)`$ the first line of the potential can be rewritten as a matrix; a rotation can then be performed such that the matrix is real and diagonal. After the rotation, the second line, being the contribution from D-terms, contains complex parameters in general, but the rotation matrix can be chosen such that these phases are set to zero.
A question arises when we include high order corrections to the potential. Then the vanishing “sneutrino” VEVs will be shifted to non-zero values by tadpoles originating, for example, from the $`QD`$ contribution in the superpotential (2.1). The “sneutrino” VEVs maybe set back to zero by a renormalization condition such that a counterterm for these VEVs set their one particle irreducible (1PI) tadpole corrections to zero.
To conclude, it is worth making a remark about the sign of $`B_0`$. As is clear from the form of Eqs. (2.30,2.8,2.16), if $`\left(\widehat{^{}}_\stackrel{~}{}^2\right)_{\alpha \alpha }\frac{1}{2}M_Z^2\frac{\mathrm{tan}^2\beta 1}{\mathrm{tan}^2\beta +1}>0`$ for all $`\alpha `$, $`B_0`$ is always positive in the vanishing sneutrino VEV basis.
## 3 Parametrising the neutral scalar mass matrices
The neutral scalar sector of the R-parity violating MSSM is in general very complicated. This is due to the fact that the scalars mix through the lepton number violating terms proportional to $`B_i`$, $`M_\mathrm{L}^2`$ and unless all of these parameters and VEVs are real one has a $`10\times 10`$ matrix to consider. However, for any given set of model parameters, one can always perform the basis change described in the previous section and arrive to the potential defined by Eq. (2.29), with only real parameters. Consequently, the physical neutral scalars are, at the tree level, exact CP-eigenstates. This implies that the neutral scalar mass matrix decouples into two $`5\times 5`$ matrices, one for the CP-odd particles and one for CP-even. In the same manner as in the R-parity conserving MSSM, once quantum corrections are considered, the CP invariance will generically be broken .
Ultimately, one would like to parametrise the scalar sector resulting from the potential in (2.29) with as few parameters as possible in order to make contact with phenomenology. These parameters in the case of the R-parity conserving MSSM are: the physical mass of the CP-odd Higgs boson
$`M_A^2={\displaystyle \frac{2B_0}{\mathrm{sin}2\beta }},`$ (3.1)
and $`\mathrm{tan}\beta `$. An advantage of the form of potential in Eq. (2.29,2.30,2.33) is that, $`M_A`$ and $`\mathrm{tan}\beta `$ can still be used for parametrising the general Higgs sector in the R-parity violating MSSM. $`M_A^2`$ is the mass of the lightest CP-odd Higgs boson in the R-parity conserving MSSM; as such, it is used here as a parameter. $`m_A^2`$ is used to denote the physical tree-level mass of the lightest CP-odd Higgs in the R-parity violating MSSM (the convention adopted is that masses in the RPC case, parameters in this model, are denoted by $`M`$, and the masses in the RPV model are denoted by $`m`$).
### 3.1 CP-even neutral scalar masses and couplings
The Lagrangian after spontaneous gauge symmetry breaking contains the terms
$`\left(\begin{array}{ccc}\mathrm{Re}h_0^2& \mathrm{Re}\stackrel{~}{\nu }_{L0}& \mathrm{Re}\stackrel{~}{\nu }_{Li}\end{array}\right)_{\mathrm{EVEN}}^2\left(\begin{array}{c}\mathrm{Re}h_0^2\\ \mathrm{Re}\stackrel{~}{\nu }_{L0}\\ \mathrm{Re}\stackrel{~}{\nu }_{Lj}\end{array}\right).`$ (3.6)
As such, the scalar CP-even Higgs squared mass matrix becomes
$`_{\mathrm{EVEN}}^2=\left(\begin{array}{ccc}\mathrm{cos}^2\beta M_A^2+\mathrm{sin}^2\beta M_Z^2& \frac{1}{2}\mathrm{sin}2\beta (M_A^2+M_Z^2)& B_j\\ \frac{1}{2}\mathrm{sin}2\beta (M_A^2+M_Z^2)& \mathrm{sin}^2\beta M_A^2+\mathrm{cos}^2\beta M_Z^2& B_j\mathrm{tan}\beta \\ B_i& B_i\mathrm{tan}\beta & M_i^2\delta _{ij}\end{array}\right),`$ (3.10)
where
$`M_i^2(\widehat{M}_{\stackrel{~}{\mathrm{L}}}^2)_i+{\displaystyle \frac{1}{2}}\mathrm{cos}2\beta M_Z^2,`$ (3.11)
are in fact the sneutrino physical masses of the RPC case. It is important here to notice that the top-left $`2\times 2`$ sub-matrix is identical to the RPC case, for which the Higgs masses are given by
$`M_{h,H}^2={\displaystyle \frac{1}{2}}\left(M_Z^2+M_A^2\pm \sqrt{(M_Z^2+M_A^2)^24M_A^2M_Z^2\mathrm{cos}^22\beta }\right),`$ (3.12)
and will be used as parameters in the RPV model.
The matrix (3.10) always has one eigenvalue which is smaller than $`M_Z^2`$. This may be proved as follows: one first observes that the upper left $`2\times 2`$ submatrix of (3.10), call it A, has at least one eigenvalue smaller than or equal to $`M_Z^2`$. Then using the Courant-Fischer theorem of the linear matrix algebra one proves that, for one flavour, the eigenvalues of the $`3\times 3`$ matrix $`_{\mathrm{EVEN}}^2`$, are interlaced with those of A. This means that the matrix $`_{\mathrm{EVEN}}^2`$ with $`i=1`$ has at least one eigenvalue smaller or equal than $`M_Z^2`$. Repeating this procedure twice, proves our statement. Furthermore, it is interesting to notice that in the region where $`\mathrm{tan}\beta 1`$, the eigenvector $`(\mathrm{sin}\beta ,\mathrm{cos}\beta ,0,0,0)^\mathrm{T}`$ corresponds to the eigenvalue with mass approximately $`M_Z^2`$. Notice that this is the same eigenvector as in the RPC case which corresponds to the Higgs boson which couples almost maximally to the Z-gauge boson.
Lepton flavour violating processes have not been observed as yet and therefore, bearing in mind cancellations, the parameters $`B_i\mathrm{tan}\beta `$ have to be much smaller than $`min(M_A^2,M_i^2)`$. To get a rough estimate, consider the dominant contribution from neutral scalars and neutralinos in the loop ,
$`m_\nu {\displaystyle \frac{a_{ew}}{16\pi }}{\displaystyle \frac{B^2\mathrm{tan}^2\beta }{\stackrel{~}{m}^3}}1\mathrm{eV},`$ (3.13)
with $`\stackrel{~}{m}=max(M_A,M_i)`$ and $`B𝒪(B_i)`$. This shows that
$`{\displaystyle \frac{B_i\mathrm{tan}\beta }{\stackrel{~}{m}^2}}{\displaystyle \frac{1.210^3}{\sqrt{\stackrel{~}{m}}}}0.1\%.`$ (3.14)
With this approximation, it is not hard to find a matrix $`𝐙_𝐑`$ which rotates the fields into the mass basis, such that
$`𝐙_{𝐑}^{}{}_{}{}^{T}_{\mathrm{EVEN}}^2𝐙_𝐑=\mathrm{diag}[m_{h^0}^2,m_{H^0}^2,(m_{\stackrel{~}{\nu }_+}^2)_i],`$ (3.15)
with $`m_{h^0}^2`$ being the lightest neutral scalar mass and
$`Z_R=\left(\begin{array}{ccc}\mathrm{cos}\alpha & \mathrm{sin}\alpha & \frac{\mathrm{cos}(\beta \alpha )\mathrm{cos}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_h^2)}+\frac{\mathrm{sin}(\beta \alpha )\mathrm{sin}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_H^2)}\\ \mathrm{sin}\alpha & \mathrm{cos}\alpha & \frac{\mathrm{cos}(\beta \alpha )\mathrm{sin}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_h^2)}+\frac{\mathrm{sin}(\beta \alpha )\mathrm{cos}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_H^2)}\\ \frac{\mathrm{cos}\beta P_i^hB_i}{\mathrm{cos}(\beta \alpha )}& \frac{\mathrm{cos}\beta P_i^HB_i}{\mathrm{sin}(\beta \alpha )}& \delta _{ij}\end{array}\right),`$ (3.19)
where there is no sum over $`i`$ and ($`M_j^2,M_h^2,M_H^2`$) are defined in (3.11,3.12). In addition,
$`\mathrm{tan}2\alpha =\mathrm{tan}2\beta {\displaystyle \frac{M_A^2+M_Z^2}{M_A^2M_Z^2}}\mathrm{and}P_i^{h,H}={\displaystyle \frac{M_Z^2\mathrm{cos}^22\beta M_{h,H}^2}{\mathrm{cos}^2\beta (M_H^2M_h^2)(M_i^2M_{h,H}^2)}},`$ (3.20)
(the common convention is to choose $`0\beta \pi /2`$ and $`\pi /2\alpha 0`$). The mass eigenstates of the RPV model are therefore given by
$`h^0`$ $``$ $`\mathrm{cos}\alpha \mathrm{Re}h_0^2\mathrm{sin}\alpha \mathrm{Re}\stackrel{~}{\nu }_{L0}+\left({\displaystyle \frac{\mathrm{cos}\beta P_i^hB_i}{\mathrm{cos}(\beta \alpha )}}\right)\mathrm{Re}\stackrel{~}{\nu }_{Li},`$ (3.21)
$`H^0`$ $``$ $`\mathrm{sin}\alpha \mathrm{Re}h_0^2+\mathrm{cos}\alpha \mathrm{Re}\stackrel{~}{\nu }_{L0}+\left({\displaystyle \frac{\mathrm{cos}\beta P_i^HB_i}{\mathrm{sin}(\beta \alpha )}}\right)\mathrm{Re}\stackrel{~}{\nu }_{Li},`$
$`(\stackrel{~}{\nu }_+)_i`$ $``$ $`\left({\displaystyle \frac{\mathrm{cos}(\beta \alpha )\mathrm{cos}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_h^2)}}+{\displaystyle \frac{\mathrm{sin}(\beta \alpha )B_j\mathrm{sin}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_H^2)}}\right)\mathrm{Re}h_0^2`$
$`+`$ $`\left({\displaystyle \frac{\mathrm{cos}(\beta \alpha )\mathrm{sin}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_h^2)}}+{\displaystyle \frac{\mathrm{sin}(\beta \alpha )\mathrm{cos}\alpha B_j}{\mathrm{cos}\beta (M_j^2M_H^2)}}\right)\mathrm{Re}\stackrel{~}{\nu }_{L0}+\mathrm{Re}\stackrel{~}{\nu }_{Li},`$
with corresponding masses,
$`m_{h^0}^2`$ $``$ $`M_h^2{\displaystyle \frac{M_Z^2\mathrm{cos}^22\beta M_h^2}{(M_H^2M_h^2)\mathrm{cos}^2\beta }}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{B_i^2}{M_i^2M_h^2}}+𝒪({\displaystyle \frac{B^4}{M^6\mathrm{cos}^4\beta }}),`$ (3.22)
$`m_{H^0}^2`$ $``$ $`M_H^2+{\displaystyle \frac{M_Z^2\mathrm{cos}^22\beta M_H^2}{(M_H^2M_h^2)\mathrm{cos}^2\beta }}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{B_i^2}{M_i^2M_H^2}}+𝒪({\displaystyle \frac{B^4}{M^6\mathrm{cos}^4\beta }}),`$ (3.23)
$`(m_{\stackrel{~}{\nu }_+}^2)_i`$ $``$ $`(\widehat{M}_{\stackrel{~}{\nu }}^2)_i+{\displaystyle \frac{B_i^2}{\mathrm{cos}^2\beta }}{\displaystyle \frac{M_i^2M_Z^2\mathrm{cos}^22\beta }{\left[M_i^4M_i^2(M_A^2+M_Z^2)+M_A^2M_Z^2\mathrm{cos}^22\beta \right]}}`$ (3.24)
$`+`$ $`𝒪({\displaystyle \frac{B^4}{M^6\mathrm{cos}^4\beta }}).`$
The above expressions, are useful in relating the masses of the neutral scalars in the RPC and RPV case in the valid approximation $`B\mathrm{tan}\beta <<min(M_A^2,M_i^2)`$. They are presented here for the first time except the mass in (3.24) which agrees with Ref.. We note here that these formulae are not valid if some of the diagonal entries in the mass matrix are closely degenerated - in such case even small $`B_i`$ terms lead to the strong mixing of respective fields. However in many types of calculations (e.g. various loop calculations) one can still formally use such expansion - in the final result one often gets expressions of the type $`\frac{f(m_1)f(m_2)}{m_1m_2}`$ which have a well defined and correct limit also for degenerate masses, even if the expansion used in the intermediate steps was, in principle, wrong.
It is interesting to note that the rotation matrix $`𝐔`$ defined in (2.6) , although explicitly calculated in this article, does not appear to all the neutral scalar vertices. For example, the vertices of the CP-even neutral scalars with the gauge bosons read as<sup>1</sup><sup>1</sup>1Note that the matrix $`𝐙`$ defined in (2.23) has nothing to do with neither $`𝐙_𝐑`$ nor $`𝐙_𝐀`$ defined in this section.,
$`_{\mathrm{VVH}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{g_2M_Z}{\mathrm{cos}\theta _w}}\left(\mathrm{cos}\beta Z_{R2s}+\mathrm{sin}\beta Z_{R1s}\right)Z^\mu Z_\mu H_s^0`$ (3.25)
$`+`$ $`{\displaystyle \frac{1}{2}}g_2M_W\left(\mathrm{cos}\beta Z_{R2s}+\mathrm{sin}\beta Z_{R1s}\right)W^{+\mu }W_\mu ^{}H_s^0,`$
where $`H_{s=1,\mathrm{},5}^0`$ are the Higgs boson fields, $`h^0,H^0,(\stackrel{~}{\nu }_+)_1,(\stackrel{~}{\nu }_+)_2,(\stackrel{~}{\nu }_+)_3`$ respectively. From (3.19) and $`_{\mathrm{VVH}}`$ above, it is easy to see that the light Higgs boson coupling to the vector bosons ($`V=Z,W`$), is proportional to $`\mathrm{sin}(\beta \alpha )`$ as in the RPC case<sup>2</sup><sup>2</sup>2We follow the conventions of Ref... In fact, the coupling sum rule,
$`{\displaystyle \underset{s=1}{\overset{5}{}}}g_{VVH_s^0}^2=g_{VV\varphi }^2,`$ (3.26)
valid in the RPC case for $`s=1,2`$, persists also here, where $`g_{VVH_s^0}`$ are the couplings appearing in (3.25) and $`g_{VV\varphi }`$ the corresponding coupling appearing in the Standard Model.
### 3.2 CP-odd neutral scalar masses and couplings
For the CP-odd case one finds,
$`\left(\begin{array}{ccc}\mathrm{Im}h_0^2& \mathrm{Im}\stackrel{~}{\nu }_{L0}& \mathrm{Im}\stackrel{~}{\nu }_{Li}\end{array}\right)_{\mathrm{ODD}}^2\left(\begin{array}{c}\mathrm{Im}h_0^2\\ \mathrm{Im}\stackrel{~}{\nu }_{L0}\\ \mathrm{Im}\stackrel{~}{\nu }_{Lj}\end{array}\right),`$ (3.31)
where the CP-odd mass matrix reads,
$`_{\mathrm{ODD}}^2=\left(\begin{array}{ccc}\mathrm{cos}^2\beta M_A^2+\xi \mathrm{sin}^2\beta M_Z^2& \frac{1}{2}\mathrm{sin}2\beta (M_A^2\xi M_Z^2)& B_j\\ \frac{1}{2}\mathrm{sin}2\beta (M_A^2\xi M_Z^2)& \mathrm{sin}^2\beta M_A^2+\xi \mathrm{cos}^2\beta M_Z^2& B_j\mathrm{tan}\beta \\ B_i& B_i\mathrm{tan}\beta & M_i^2\delta _{ij}\end{array}\right),`$ (3.35)
and $`\xi `$ is the gauge fixing parameter in $`R_\xi `$ gauge. In fact, by using an orthogonal rotation
$`𝒱=\left(\begin{array}{ccc}\mathrm{sin}\beta & \mathrm{cos}\beta & 0\\ \mathrm{cos}& \mathrm{sin}\beta & 0\\ 0& 0& 1\end{array}\right),`$ (3.39)
we can always project out the would-be Goldstone mode, of the CP-odd scalar matrix and thus
$`𝒱^\mathrm{T}_{\mathrm{ODD}}^2𝒱=\left(\begin{array}{ccc}\xi M_Z^2& 0& 0\\ 0& M_A^2& \frac{B_j}{\mathrm{cos}\beta }\\ 0& \frac{B_i}{\mathrm{cos}\beta }& M_i^2\delta _{ij}\end{array}\right).`$ (3.43)
Under the approximation of small bilinear RPV couplings \[see Eq. (3.14)\], a solution is determined for the matrix $`𝐙_𝐀`$ which rotates the fields into the mass basis, such that
$`𝐙_{𝐀}^{}{}_{}{}^{T}_{\mathrm{ODD}}^2𝐙_𝐀=\mathrm{diag}[m_{\mathrm{G}^0}^2,m_{\mathrm{A}^0}^2,(m_{\stackrel{~}{\nu }_{}}^2)_i],`$ (3.44)
$`Z_A=\left(\begin{array}{ccc}\mathrm{sin}\beta & \mathrm{cos}\beta & \frac{B_j}{M_j^2M_A^2}\\ \mathrm{cos}\beta & \mathrm{sin}\beta & \frac{B_j\mathrm{tan}\beta }{M_j^2M_A^2}\\ 0& \frac{B_i}{\mathrm{cos}\beta (M_i^2M_A^2)}& \delta _{ij}\end{array}\right),`$ (3.48)
with the mass eigenstates given by
$`G^0`$ $``$ $`\mathrm{sin}\beta \mathrm{Im}h_0^2\mathrm{cos}\beta \mathrm{Im}\stackrel{~}{\nu }_{L0},`$
$`A^0`$ $``$ $`\mathrm{cos}\beta \mathrm{Im}h_0^2+\mathrm{sin}\beta \mathrm{Im}\stackrel{~}{\nu }_{L0}+{\displaystyle \frac{B_i}{\mathrm{cos}\beta (M_i^2M_A^2)}}\mathrm{Im}\stackrel{~}{\nu }_{Li},`$
$`(\stackrel{~}{\nu }_{})_i`$ $``$ $`{\displaystyle \frac{B_i}{M_j^2M_A^2}}\mathrm{Im}h_0^2+{\displaystyle \frac{B_j\mathrm{tan}\beta }{M_i^2M_A^2}}\mathrm{Im}\stackrel{~}{\nu }_{L0}+\mathrm{Im}\stackrel{~}{\nu }_{Li},`$ (3.49)
with corresponding masses,
$`m_A^2`$ $``$ $`M_A^2{\displaystyle \frac{1}{\mathrm{cos}^2\beta }}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{B_i^2}{M_i^2M_A^2}}+𝒪({\displaystyle \frac{B^4}{M^6\mathrm{cos}^4\beta }}),`$ (3.50)
$`(m_{\stackrel{~}{\nu }_{}}^2)_i`$ $``$ $`M_i^2{\displaystyle \frac{B_i^2}{\left(M_A^2M_i^2\right)\mathrm{cos}^2\beta }}+𝒪({\displaystyle \frac{B^4}{M^6\mathrm{cos}^4\beta }}).`$ (3.51)
The coupling of the Z gauge boson to the CP-even and CP-odd neutral scalar fields is given by
$`_{ZHA}={\displaystyle \frac{ig_2}{2c_W}}\left[(p_{H_s^0}p_{A_p^0})_\mu \left({\displaystyle \underset{\alpha =0}{\overset{3}{}}}Z_{R(2+\alpha )s}Z_{A(2+\alpha )p}Z_{R\mathrm{\hspace{0.17em}1}s}Z_{A\mathrm{\hspace{0.17em}1}p}\right)\right]Z^\mu H_s^0A_p^0,`$ (3.52)
where the four momenta $`p_{H_s^0}^\mu ,p_{A_p^0}^\mu `$ are incoming and the fields $`A_{p=1,\mathrm{}5}^0`$ correspond to $`G^0,A^0,(\stackrel{~}{\nu }_{})_1,(\stackrel{~}{\nu }_{})_2,(\stackrel{~}{\nu }_{})_3`$ respectively. One may check that the coupling $`ZG^0h^0`$ derived from (3.52) is proportional to $`\mathrm{sin}(\alpha \beta )`$ as it should be.
## 4 Positiveness and stability of the scalar potential
### 4.1 Positiveness
In general, one should inspect whether all physical masses in the CP-odd and CP-even sector are positive. For that, all diagonal square subdeterminants of mass matrices should be positive. One can easily check that both CP-odd and CP-even mass matrices in (3.10,3.35) respectively, lead, in the rotated basis, to the same set of conditions,
$`M_i^2`$ $`>`$ $`0\mathrm{with}i=1,2,3\mathrm{and}M_A^2>{\displaystyle \frac{1}{\mathrm{cos}^2\beta }}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{B_i^2}{M_i^2}}.`$ (4.1)
Using the form of $`M_A^2`$ in (3.1), the last equation can be rewritten in the form
$`B_0>\mathrm{tan}\beta {\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{B_i^2}{M_i^2}}.`$ (4.2)
Excluding some very singular mass configurations, the above conditions are rather trivially fulfilled if one takes into account the bound of Eq. (3.14).
### 4.2 Stability
The question of whether the potential is stable, i.e. bounded from below, is far more complicated. In most cases the quartic ($`D`$)term dominates and there is no problem. The only exception being when the fields follow the direction $`|h_2^0|^2=_{i=0}^4|\stackrel{~}{\nu }_{Li}|^2`$. In such a case, one should check whether the remaining part of the potential is positive along this direction.
Denoting $`R\sqrt{_{i=0}^3|\stackrel{~}{\nu }_i|^2}`$ and $`h_2^0=Re^{i\varphi }`$, where $`\varphi `$ is a free phase, and using Eqs. (2.20,2.33,3.11), one can write down the scalar potential along this direction in the vanishing snueutrino VEV basis as
$`V_{\text{neutral}}`$ $`=`$ $`{\displaystyle \frac{B_0}{\mathrm{sin}\beta \mathrm{cos}\beta }}\stackrel{~}{\nu }_{L0}^{}\stackrel{~}{\nu }_{L0}^{}+\left[M_i^2+B_0\mathrm{cot}\beta \right]\stackrel{~}{\nu }_{Li}^{}\stackrel{~}{\nu }_{Li}^{}`$ (4.3)
$`+`$ $`B_i\mathrm{tan}\beta (\stackrel{~}{\nu }_{L0}^{}\stackrel{~}{\nu }_{Li}^{}+\stackrel{~}{\nu }_{L0}^{}\stackrel{~}{\nu }_{Li}^{})B_\alpha (\stackrel{~}{\nu }_{L\alpha }^{}h_0^2+H.c.)`$
$``$ $`\stackrel{~}{\nu }_𝐋^{}𝐐\stackrel{~}{\nu }_𝐋(𝐁^T\stackrel{~}{\nu }_L^{}Re^{i\varphi }+\mathrm{H}.\mathrm{c}.).`$
where the real symmetric matrix $`𝐐`$ is
$`𝐐=\left(\begin{array}{cc}M_A^2& B_i\mathrm{tan}\beta \\ B_j\mathrm{tan}\beta & \left[M_i^2+B_0\mathrm{cot}\beta \right]\delta _{ij}\end{array}\right).`$ (4.6)
Finding the stability conditions for the potential (4.3) is difficult, it depends on nine real variables (4 moduli and five phases of the fields). To simplify the problem, we perform one more field rotation to the basis in which the matrix $`𝐐`$ is diagonal. This can be done, in general, by numerical routines (routines where already used in calculating the vanishing sneutrino VEV basis, and therefore, finding the stability conditions for the general scalar potential always has to involve some numerical analysis). We thus define the matrix $`𝐏`$, $`\stackrel{~}{\nu }_LP\stackrel{~}{\nu }_L`$, as
$`𝐏^{}𝐐𝐏=\mathrm{diag}(X_0,X_1,X_2,X_3).`$ (4.7)
In fact, $`𝐐`$ is real, so we can choose $`𝐏`$ to be real orthogonal. Also, we denote $`D_\beta B_\alpha P_{\alpha \beta }`$. Obviously, the rotation $`𝐏`$ preserves the value of $`R=|h_2^0|`$.
The potential becomes:
$`V_{\text{neutral}}`$ $`=`$ $`{\displaystyle \underset{\alpha =0}{\overset{3}{}}}[X_\alpha |\stackrel{~}{\nu }_{L\alpha }^{}|^2D_\alpha R(\stackrel{~}{\nu }_{L\alpha }^{}e^{i\varphi }+\mathrm{H}.\mathrm{c}.)],`$ (4.8)
where $`X_0`$ has to be positive, otherwise for $`\varphi =0`$ along the direction $`\stackrel{~}{\nu }_{Li}^{}=\mathrm{Im}\stackrel{~}{\nu }_{L0}^{}=0`$ the potential $`V_{\text{neutral}}=|\mathrm{Re}\stackrel{~}{\nu }_{L0}^{}|^2[X_0D_0\mathrm{sign}(\mathrm{Re}\stackrel{~}{\nu }_{L0}^{})]`$ falls to $`\mathrm{}`$ at least for one direction along the $`\mathrm{Re}\stackrel{~}{\nu }_{L0}^{}`$ axis. In fact the condition on $`X_\alpha `$ is $`X_\alpha 2|D_\alpha |`$. Thus our first conclusion is that the matrix $`𝐐`$ has to be positively defined. One can write down appropriate conditions in the same manner as for the scalar mass matrices; comparing with Eq. (4.1), it can be observed that this condition is automatically fulfilled if relation (4.1) holds.
With $`X_\alpha `$ positive, one can write down the potential as:
$`V_{\text{neutral}}`$ $`=`$ $`{\displaystyle \underset{\alpha =0}{\overset{3}{}}}\left|\sqrt{X_\alpha }\stackrel{~}{\nu }_{L\alpha }^{}{\displaystyle \frac{D_\alpha }{\sqrt{X_\alpha }}}Re^{i\varphi }\right|^2R^2{\displaystyle \underset{\alpha =0}{\overset{3}{}}}{\displaystyle \frac{D_\alpha ^2}{X_\alpha }}.`$ (4.9)
To further simplify the problem, denote $`\stackrel{~}{\nu }_{L\alpha }^{}=u_\alpha e^{i(\varphi \varphi _\alpha )}`$, where $`u_\alpha 0`$ are field moduli and $`\varphi _\alpha `$ are free phases. Then
$`V_{\text{neutral}}`$ $`=`$ $`R^2\left({\displaystyle \underset{\alpha =0}{\overset{3}{}}}\left|\sqrt{X_\alpha }{\displaystyle \frac{u_\alpha }{R}}{\displaystyle \frac{D_\alpha }{\sqrt{X_\alpha }}}e^{i\varphi _\alpha }\right|^2{\displaystyle \underset{\alpha =0}{\overset{3}{}}}{\displaystyle \frac{D_\alpha ^2}{X_\alpha }}\right),`$ (4.10)
where $`R=\sqrt{_{i=0}^3|\stackrel{~}{\nu }_i|^2}=\sqrt{_{i=0}^3u_i^2}`$. Phases $`\varphi _\alpha `$ can be adjusted independently of $`u_\alpha `$. The worst case from the point of view of potential stability, the smallest first term inside the parenthesis, occurs for $`D_\alpha e^{i\varphi _\alpha }=|D_\alpha |`$. Denoting further $`ϵ_\alpha =u_\alpha /R`$, $`0ϵ_\alpha 1`$, one can reduce our initial problem to the question whether the function
$`g(ϵ_\alpha )`$ $`=`$ $`{\displaystyle \underset{\alpha =0}{\overset{3}{}}}\left|\sqrt{X_\alpha }ϵ_\alpha {\displaystyle \frac{|D_\alpha |}{\sqrt{X_\alpha }}}\right|^2{\displaystyle \underset{\alpha =0}{\overset{3}{}}}{\displaystyle \frac{D_\alpha ^2}{X_\alpha }}={\displaystyle \underset{\alpha =0}{\overset{3}{}}}(X_\alpha ϵ_\alpha ^22|D_\alpha |ϵ_\alpha ),`$ (4.11)
depending now on four real positive parameters, is non-negative on the unit sphere $`_{\alpha =0}^3ϵ_\alpha ^2=1`$. In general such problem can be solved numerically using the method of Lagrange multipliers. For $`X_i>X_0D_0`$, the minimum occurs for
$`ϵ_\alpha ={\displaystyle \frac{|D_\alpha |}{X_\alpha +\lambda }},`$ (4.12)
where $`\lambda `$ can be found numerically as a root of the following equation:
$`{\displaystyle \underset{\alpha =0}{\overset{3}{}}}{\displaystyle \frac{D_\alpha ^2}{(X_\alpha +\lambda )^2}}=1.`$ (4.13)
For smaller $`X_i`$, the minimum is realized for $`ϵ_i=0`$ for one or more values of $`i`$ and requires analysis of various special cases. Having found the correct minimum, to prove the stability of the potential one needs to show that the function $`g`$ at the minimum is non-negative.
As shown in Eq. (3.14), $`B_i`$ terms and thus also $`D_i`$ terms are usually very small. In this case one can set approximate, sufficient conditions for the stability of the potential, without resorting to solving Eq. (4.13), numerically. Denote $`D=_{i=1}^3D_i^2`$ and $`X_{min}=min(X_1,X_2,X_3)`$. Then, using the inequality $`D_iϵ_i\sqrt{_{i=1}^3D_i^2}\sqrt{_{i=1}^3ϵ_i^2}=D\sqrt{1ϵ_0^2}`$, one has
$`g(ϵ_\alpha )`$ $``$ $`X_0ϵ_0^2+X_{min}(1ϵ_0^2)+(X_iX_{min})ϵ_i^22|D_0|ϵ_02D\sqrt{1ϵ_0^2}.`$ (4.14)
Terms $`(X_iX_{min})ϵ_i^2`$ are always non-negative. The worst case being when the vector $`(ϵ_1,ϵ_2,ϵ_3)`$ is along the minimal $`X_i`$ axis, where these terms vanish. Other terms are rotation invariant in the 3-dimensional space $`(ϵ_1,ϵ_2,ϵ_3)`$, so Eq. (4.14) is equivalent to finding parameters $`X_0,X_{min},D_0,D`$ for which the expression (4.15), depending on just one real variable, is positive:
$`g^{}(ϵ_0)=X_0ϵ_0^2+X_{min}(1ϵ_0^2)2|D_0|ϵ_02D\sqrt{1ϵ_0^2}0.`$ (4.15)
Analysis of (4.15) is further simplified by one more approximation, justified for small $`D`$:
$`g^{}(ϵ_0)`$ $``$ $`X_0ϵ_0^2+X_{min}(1ϵ_0^2)2|D_0|ϵ_02D.`$ (4.16)
The rhs of Eq. (4.16) is now trivial. Following approximate conditions for the stability of the potential can be summarized as follows:
| | $`X_{min}`$ range | | Stability requires |
| --- | --- | --- | --- |
| | $`X_{min}X_0D_0`$ | | $`X_02|D_0|+2D`$ |
| | $`0<X_{min}<X_0D_0`$ | | $`(X_0X_{min})(X_{min}2D)D_0^2`$ |
Both conditions are sufficient, but not minimal - we have made some approximations and there may be parameters which do not fall into either of the categories above, and yet still give a stable potential. For example, if $`X_0=X_1=X_2=X_3X`$, one can easily derive the exact necessary and sufficient condition for potential stability as $`X2\sqrt{D_0^2+D^2}`$, less strict than $`X2(|D_0|+|D|)`$ which would be given by the table above.
For complementary work the reader is referred to Ref..
## 5 Conclusions
In this letter we present a procedure for calculating the rotation matrix which brings the neutral scalar fields of the general R-parity violating MSSM onto the vanishing sneutrino VEV basis where they develop $`n`$ zero VEVs, with $`n`$ being the number of flavour generations. In doing so, we have made no assumption about the complexity of the parameters. We consider the case of $`n=3`$ generations, but our approach immediately applies to other cases, apart from obvious modifications of the form of $`𝐙`$ matrix defined in (2.23.2.28). As a byproduct of basis change, we prove that the tree level MSSM potential does not exhibit any form of CP-violation, neither explicit nor spontaneous. Consequently, the neutral scalar fields can be divided into CP-even and CP-odd sectors with the $`5\times 5`$ neutral scalar squared mass matrices, taking a very simple form with only small RPV masses sitting on their off diagonal elements. We can thus expand along small RPV masses and find analytic approximate formulae which relate the RPC and the RPV neutral scalar masses. Furthermore we also find, that in general there is always at least one neutral scalar field with mass lighter than $`M_Z`$ which couples maximally to the $`Z`$-gauge boson in the case of large $`\mathrm{tan}\beta `$ and large $`M_A`$. Our procedure for finding the rotation matrix $`𝐔`$ has been coded<sup>3</sup><sup>3</sup>3The code will be available from http://www.fuw.edu.pl/$``$rosiek/physics/rpv/scalar.html and is numerically stable.
In the end, we are aiming to construct the most general MSSM quantum field theory structure resorting neither to R-parity violation nor to other approximations. This will be useful for examining the phenomenology of the MSSM as a whole. The convenient choice of the basis for the neutral sector found in this paper is a first step towards this direction.
Acknowledgments
The authors thank Apostolos Pilaftsis for helpful discussions. A.D. and M.S.-S. would like to thank “The Nuffield Foundation” for financial support. J.R. would like to thank IPPP for the hospitality during his visit. His work was supported in part by the KBN grant 2 P03B 040 24 (2003-2005). S.R. acknowledges the award of a UK PPARC studentship.
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# Spin dynamics of the quasi two dimensional spin-1/2 quantum magnet Cs2CuCl4
## I Introduction
The quasi two dimensional spin-1/2 quantum magnet Cs<sub>2</sub>CuCl<sub>4</sub> has attracted much theoretical and experimental interest in recent years as a possible realization of a two dimensional quantum spin liquid. Coldea0 ; Coldea1 ; Coldea2 ; Coldea3 ; BETG ; Bocquet ; McKenzie1 ; Kim1 ; Zhang1 ; Zhang2 ; Wen1 ; Isakov This anisotropic triangular Heisenberg antiferromagnet is believed to be a promising candidate due to its small spin, quasi two dimensionality and geometrically frustrated spin interactions. Although Cs<sub>2</sub>CuCl<sub>4</sub> exhibits conventional incommensurate long range magnetic order at low temperatures, neutron scattering measurements have revealed unusual features in the spin excitation spectrum. In particular, the dynamical correlations are found to be dominated by an extended scattering continuum over a relatively large window of energies. Several workers have interpreted this observation as a signature of deconfined, fractionalized spin-1/2 (spinon) excitations, characteristic of a spin liquid phase. In this line of approach, the observed broad scattering continuum is interpreted in terms of a two-spinon scattering continuum. McKenzie1 ; Wen1 ; Kim1
However, a strong scattering continuum does not entail an underlying spin liquid phase. In fact, a conventional magnetically ordered phase with strong magnon interactions can exhibit a broad continuum due to multi magnon scattering processes. A previous examination of the inelastic neutron scattering data on Cs<sub>2</sub>CuCl<sub>4</sub> was performed in the framework of linear spin wave (LSW) theory. Coldea3 The latter predicts sharp single particle excitations and weak two magnon scattering continua, features which were argued to be in poor agreement with the data. Given that the magnetic properties derive from small $`S=1/2`$ Cu spins, one would a priori expect magnon interactions to play an important role. In order to assess the applicability of a spin wave based scenario to Cs<sub>2</sub>CuCl<sub>4</sub> it is therefore necessary to go beyond linear spin wave theory.
On a qualitative level the predictions of nonlinear spin wave theory are readily anticipated. By Goldstone’s theorem the breaking of a continuous symmetry in a magnetically ordered state enforces the presence of single particle excitations at low energies. As a result of the aforementioned interactions, these magnons acquire a finite life time, which in turn leads to a finite line width in the dynamical structure factor. Furthermore, compared to linear spin wave theory, spectral weight is transferred to higher energies via multi magnon scattering processes. In the case of Cs<sub>2</sub>CuCl<sub>4</sub> one may expect the presence of a strong scattering continuum in the ordered phase because (1) the low spin and the frustrated nature of the exchange interactions lead to a small ordered moment and strong quantum fluctuations around the ordered state; (2) the magnon interactions in non-collinear spin structures like the ones found in Cs<sub>2</sub>CuCl<sub>4</sub> induce a coupling between transverse and longitudinal spin fluctuations. This interaction provides an additional mechanism for damping the spin waves and can enhance the strength of the scattering continuum.
There is evidence of low-energy spin wave modes in the inelastic neutron scattering data. Sharp peaks are also observed at high energies near special wave vectors where a putative spin wave dispersion is at a saddle-point. It is important to note that this spin wave dispersion is dramatically “renormalized” compared to the prediction of linear spin wave theory. Coldea2 ; Coldea3
A priori it appears that nonlinear spin wave theory could have the necessary ingredients to account for the spin correlations observed in Cs<sub>2</sub>CuCl<sub>4</sub>. The issue then is whether it is possible to achieve a quantitative description of the experiments in low orders of perturbation theory in the spin wave interactions.
In the present work we go beyond linear spin wave theory and include, within the framework of a $`1/S`$ expansion, the quantum fluctuations around the classical ground state. We then apply the results to the case $`S=1/2`$, in which the formal expansion parameter becomes of order $`1`$ and is therefore not small. We are motivated by the observation that spin wave theory gives a good description of physical properties of the square-lattice spin-$`\frac{1}{2}`$ Heisenberg Hamiltonian. Canali2 ; Hamer ; Igarashi1 Indeed, higher order (in a $`1/S`$ expansion) corrections to linear spin wave theory were shown to be small in this case. Furthermore, taking these corrections into account in the calculation of static and dynamical properties leads to an improved agreement with the results of more sophisticated numerical techniques. Singh2 ; Singh3 ; Sandvik Although a corresponding analysis is not available for the frustrated triangular antiferromagnet, perturbative expansions in $`1/S`$ have shown the renormalization due to quantum effects is relatively small. Chubukov1 ; Shiba2 ; Shiba1 ; Shiba3
This paper is organized as follows. The spin Hamiltonian for Cs<sub>2</sub>CuCl<sub>4</sub> is introduced in Sec. II. In Sec. III we determine the magnon Green’s function in the framework of a large-S expansion. In Sec. IV we relate the experimentally measured dynamical correlation functions to the magnon Green’s function. The results of our analysis and comparisons to the experimental data on Cs<sub>2</sub>CuCl<sub>4</sub> are presented in Sec. V. We conclude with a summary of our results in Sec. VI.
## II Spin Model
The full spin Hamiltonian of Cs<sub>2</sub>CuCl<sub>4</sub> has been determined previously from measurements in high magnetic fields (see Ref. Coldea2, for details). For our purposes it suffices to note that the magnetic Cu<sup>2+</sup> ions form a triangular lattice with anisotropic exchange interactions. As shown in Fig. 1, the main exchange interaction $`J=0.374(5)`$ meV is along the crystallographic $`b`$ axis (“chain direction”). A weaker spin exchange $`J^{}=0.128(5)`$ meV occurs along the zig-zag bonds. Finally, a Dzyaloshinskii-Moriya (DM) interaction Dzyaloshinski58 ; Moriya60 $`D=0.020(2)`$ meV is present along the zig-zag bonds.
Denoting the spin-$`\frac{1}{2}`$ operators at the sites $`𝐑`$ by $`𝐒_𝐑`$, the quasi two dimensional Hamiltonian takes the form
$`={\displaystyle \underset{𝐑}{}}`$ $`J𝐒_𝐑𝐒_{𝐑+\delta _\mathrm{𝟏}+\delta _\mathrm{𝟐}}+J^{}\left(𝐒_𝐑𝐒_{𝐑+\delta _\mathrm{𝟏}}+𝐒_𝐑𝐒_{𝐑+\delta _\mathrm{𝟐}}\right)(1)^n𝐃𝐒_𝐑\times \left(𝐒_{𝐑+\delta _\mathrm{𝟏}}+𝐒_{𝐑+\delta _\mathrm{𝟐}}\right).`$ (1)
Here the vectors $`\delta _\mathrm{𝟏}`$ and $`\delta _\mathrm{𝟐}`$ connecting neighboring sites are shown in Fig. 1. The vector $`𝐃=(D,0,0)`$ is associated with the oriented bond between the two coupled spins connected by $`\delta _\mathrm{𝟏}`$ or $`\delta _\mathrm{𝟐}`$ and $`n`$ is a layer index. The factor $`(1)^n`$ indicates that the interaction alternates between even and odd layers, which as a result can be considered to be inverted versions of one another. A weak interlayer interaction $`J^{\prime \prime }`$ is also present between neighboring layers. However, as $`J^{\prime \prime }`$ is quite small we neglect it in the following.
Following the conventions of Coldea et al. in Ref. Coldea3, , we will discuss the dynamic response in terms of the two dimensional Brillouin zone of the triangular lattice even though the full crystal symmetry is orthorhombic. In our notation wave vectors are expressed in terms of the reciprocal lattice vectors as $`𝐤=(h,k,l)`$, which is a shorthand for $`2\pi (h/a,k/b,l/c)`$.
The Fourier transforms of the exchange and DM interactions are
$$J_𝐐=J\mathrm{cos}(2\pi k)+2J^{}\mathrm{cos}(\pi k)\mathrm{cos}(\pi l),$$
(2)
and
$$D_𝐐=2iD\mathrm{sin}(\pi k)\mathrm{cos}(\pi l).$$
(3)
It is convenient for what follows to define a quantity
$$J_𝐐^T=J_𝐐iD_𝐐.$$
(4)
Experimentally, spiral magnetic long range order is observed in Cs<sub>2</sub>CuCl<sub>4</sub> at temperatures below $`T_N=0.62(1)K`$. The ordered structure is found to lie in the $`bc`$ plane by virtue of the small easy-plane anisotropy generated by the DM interactions. The spin structure is an incommensurate cycloid with an ordering wave vector $`𝐐=(0.0,0.5+ϵ,0)`$ where $`ϵ=0.030(2)`$.
## III Large S Expansion
We now turn to a summary of our calculations. The procedure we follow is standard. We first express the fluctuations around the “classical” ground state in terms of boson operators using the Holstein-Primakoff transformation. HolsteinPrimakoff ; Nagamiya ; White01 ; Shiba2 ; Shiba1 ; Shiba3 ; Zhitomirsky02 ; Zhitomirsky98 The term quadratic in the boson operators constitutes the basis for linear spin wave theory, whereas higher order terms represent spin wave interactions. The interaction vertices of $`n`$ bosons carry a factor $`S^{2n/2}`$, where $`S`$ is the “length” of the spin. In the second step we determine the renormalized magnon Green’s function by calculating the self-energy to leading order in $`1/S`$. Finally, the experimentally observable dynamical correlation functions are expressed in terms of the Green’s function of the Holstein-Primakoff bosons.
The classical ground state is determined by treating the spins as classical vectors and then minimizing the energy. In this way one obtains a cycloidal structure with a characteristic wave vector $`𝐐`$ that is fixed by the condition that it minimizes the exchange energy per spin, i.e. $`J_𝐐^T=\mathrm{min}_𝐪J_𝐪^T`$. We find $`𝐐=(0.0,0.5+ϵ_0,0)`$ with $`ϵ_0=0.054`$. This value differs significantly from the measured incommensuration but quantum fluctuations lead to a reduction in $`ϵ_0`$ and taking them into account yields good agreement with experiments. Veillette ; McKenzie1
As we have already indicated in Eq. (1), to a good approximation the layers are decoupled. Hence we consider from now on a set of independent 2-D layers, which are subdivided into two groups, differing according to the direction of the DM vector. For the case where the layer index $`n`$ is odd (even), the DM vector is taken to point into (out of) the $`bc`$ plane.
In what follows we present the results for the even layers only. However, it is easy to see that the spin structure factor is in fact independent of the layer index and the overall result is a simple summation over all layers.
It is convenient to define a local reference frame $`(x,y,z)`$ such that the classical spin direction is aligned along the $`z`$ axis at every site
$`\left(\begin{array}{c}S_𝐑^a\\ S_𝐑^b\\ S_𝐑^c\end{array}\right)=\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}(𝐐𝐑)& \mathrm{sin}(𝐐𝐑)\\ 0& \mathrm{sin}(𝐐𝐑)& \mathrm{cos}(𝐐𝐑)\end{array}\right)\left(\begin{array}{c}S_𝐑^x\\ S_𝐑^y\\ S_𝐑^z\end{array}\right).`$ (5)
The Holstein-Primakoff transformation reads HolsteinPrimakoff
$`S_𝐑^+`$ $`=`$ $`S_𝐑^x+iS_𝐑^y=e^{i\theta }\sqrt{\left(2S\varphi _𝐑^{}\varphi _𝐑^{}\right)}\varphi _𝐑^{},`$
$`S_𝐑^{}`$ $`=`$ $`S_𝐑^xiS_𝐑^y=e^{i\theta }\varphi _𝐑^{}\sqrt{\left(2S\varphi _𝐑^{}\varphi _𝐑^{}\right)},`$
$`S_𝐑^z`$ $`=`$ $`S\varphi _𝐑^{}\varphi _𝐑^{},`$ (6)
where the boson creation and annihilation operators satisfy the canonical commutation relation $`[\varphi _𝐑^{},\varphi _𝐑^{}^{}]=\delta _{𝐑,𝐑^{}}`$. Here $`\theta `$ is an arbitrary angle which we set equal to $`\pi /2`$ in order to make contact with the notation used in Ref. Shiba1, . Introducing the Fourier transform
$$\varphi _𝐤^{}=\frac{1}{\sqrt{N}}\underset{𝐑}{}\varphi _𝐑^{}e^{i𝐤𝐑},$$
(7)
on a lattice of $`N`$ sites, the Hamiltonian of Eq. (1) takes the form
$$=_0+_2+_3+_4+\mathrm{},$$
(8)
where $`_n`$ is proportional to $`S^{2n/2}`$ and consists of normal ordered products of $`n`$ boson operators. There is no $`_1`$ term, because Eq. (8) is an expansion around a minimum of the classical energy. Linear spin wave theory takes into account only the terms $`_0`$ and $`_2`$. The higher order terms represent interactions between magnons. The leading terms in the expansion are
$`_0`$ $`=`$ $`NS^2J_𝐐^T,`$ (9)
$`_2`$ $`=`$ $`NSJ_𝐐^T+S{\displaystyle \underset{𝐤}{}}A_𝐤\left(\varphi _𝐤^{}\varphi _𝐤^{}+\varphi _𝐤^{}\varphi _𝐤^{}\right)B_𝐤\left(\varphi _𝐤^{}\varphi _𝐤^{}+\varphi _𝐤^{}\varphi _𝐤^{}\right),`$ (10)
$`_3`$ $`=`$ $`{\displaystyle \frac{i}{2}}\sqrt{{\displaystyle \frac{S}{2N}}}{\displaystyle \underset{\mathrm{𝟏},\mathrm{𝟐},\mathrm{𝟑}}{}}\delta _{\mathrm{𝟏}+\mathrm{𝟐}+\mathrm{𝟑}}\left(C_\mathrm{𝟏}+C_\mathrm{𝟐}\right)\left(\varphi _\mathrm{𝟑}^{}\varphi _\mathrm{𝟐}^{}\varphi _\mathrm{𝟏}^{}\varphi _\mathrm{𝟏}^{}\varphi _\mathrm{𝟐}^{}\varphi _\mathrm{𝟑}^{}\right),`$ (11)
$`_4`$ $`=`$ $`{\displaystyle \frac{1}{4N}}{\displaystyle \underset{\mathrm{𝟏},\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}}{}}\delta _{\mathrm{𝟏}+\mathrm{𝟐}+\mathrm{𝟑}+\mathrm{𝟒}}\{{\displaystyle \frac{2}{3}}(B_\mathrm{𝟐}+B_\mathrm{𝟑}+B_\mathrm{𝟒})(\varphi _\mathrm{𝟏}^{}\varphi _\mathrm{𝟐}^{}\varphi _\mathrm{𝟑}^{}\varphi _\mathrm{𝟒}^{}+\varphi _\mathrm{𝟒}^{}\varphi _\mathrm{𝟑}^{}\varphi _\mathrm{𝟐}^{}\varphi _\mathrm{𝟏}^{})`$ (12)
$`+`$ $`[(A_{\mathrm{𝟏}+\mathrm{𝟑}}+A_{\mathrm{𝟏}+\mathrm{𝟒}}+A_{\mathrm{𝟐}+\mathrm{𝟑}}+A_{\mathrm{𝟐}+\mathrm{𝟒}})(B_{\mathrm{𝟏}+\mathrm{𝟑}}+B_{\mathrm{𝟏}+\mathrm{𝟒}}+B_{\mathrm{𝟐}+\mathrm{𝟑}}+B_{\mathrm{𝟐}+\mathrm{𝟒}})(A_\mathrm{𝟏}+A_\mathrm{𝟐}+A_\mathrm{𝟑}+A_\mathrm{𝟒})]\varphi _\mathrm{𝟏}^{}\varphi _\mathrm{𝟐}^{}\varphi _\mathrm{𝟑}^{}\varphi _\mathrm{𝟒}^{}\}.`$
Here the sum over $`𝐤`$ is performed in the first Brillouin zone and the subscripts $`\mathrm{𝟏}\mathrm{}\mathrm{𝟒}`$ denote $`𝐤_\mathrm{𝟏}\mathrm{}𝐤_\mathrm{𝟒}`$. The quantities $`A_𝐤,B_𝐤`$ and $`C_𝐤`$ are expressed as
$`A_𝐤`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(2J_𝐤+J_{𝐐+𝐤}^T+J_{𝐐𝐤}^T\right)J_𝐐^T,`$
$`B_𝐤`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(2J_𝐤J_{𝐐+𝐤}^TJ_{𝐐𝐤}^T\right),`$
$`C_𝐤`$ $`=`$ $`J_{𝐐+𝐤}^TJ_{𝐐𝐤}^T.`$ (13)
The coefficients $`A_𝐤`$ and $`B_𝐤`$ are even functions of $`𝐤`$, whereas $`C_𝐤`$ is an odd function of $`𝐤`$. In the absence of easy-plane anisotropies, i.e when $`D`$ vanishes and inversion symmetry is present, we recover the results of Ref. Shiba1, . \[Note that our definitions in Eqs. (13) differ from those of Ref. Shiba1, by a factor of four.\] We emphasize that the cubic interaction is generated as a result of the coupling between transverse and longitudinal fluctuations and hence can only exist in non-collinear spin structures. Furthermore, we note that the vertex factor $`C_𝐤|𝐤|^3`$ for small $`𝐤`$ owing to the fact that $`J_𝐐^T`$ is at a minimum by construction.
The quadratic Hamiltonian $`_2`$ is diagonalized by a Bogoliubov transformation
$`\varphi _𝐤^{}`$ $`=`$ $`u_𝐤\gamma _𝐤^{}+v_𝐤\gamma _𝐤^{},`$
$`\varphi _𝐤^{}`$ $`=`$ $`v_𝐤\gamma _𝐤^{}+u_𝐤\gamma _𝐤^{},`$ (14)
where
$`u_𝐤^2`$ $`=`$ $`1+v_𝐤^2={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{A_𝐤}{\sqrt{A_𝐤^2B_𝐤^2}}}+1\right),`$
$`u_𝐤v_𝐤`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{B_𝐤}{\sqrt{A_𝐤^2B_𝐤^2}}}.`$ (15)
The diagonal form of the quadratic Hamiltonian is
$`_2`$ $`=`$ $`NSJ_𝐐^T+{\displaystyle \underset{𝐤}{}}\omega _𝐤\left(\gamma _𝐤^{}\gamma _𝐤^{}+{\displaystyle \frac{1}{2}}\right),`$ (16)
where $`\omega _𝐤=2S\sqrt{A_𝐤^2B_𝐤^2}`$ is the linear spin wave dispersion relation. White01 ; Nagamiya We note that $`\omega _𝐤`$ is an even function of $`𝐤`$, despite the absence of inversion symmetry in the Hamiltonian. In fact, the symmetry of $`\omega _𝐤`$ is a consequence of time-reversal symmetry, which implies the following relation between the elements of the dynamical structure factor (Eq. 30), Lovesey
$$S_{𝐤,\omega }^{\mu \nu }=S_{𝐤,\omega }^{\nu \mu }.$$
(17)
The importance of quantum fluctuations can be gauged by determining the average value of the local spin given by the standard formula
$$S_𝐑^z=S\mathrm{\Delta }S=S\frac{1}{2N}\underset{𝐤}{}u_𝐤^2+v_𝐤^2.$$
(18)
The boson Green’s function at zero temperature is expressed as
$`G_{𝐤,\omega }=i{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑te^{i\omega t}T\left[\begin{array}{c}\varphi ^{}_𝐤(t)\hfill \\ \varphi ^{}_𝐤(t)\hfill \end{array}\right]\left[\varphi _𝐤^{}(0)\varphi _𝐤^{}(0)\right],`$ (21)
where $`T`$ denotes time ordering and $`\mathrm{}`$ represents a ground state expectation value. The inverse of the unperturbed Green’s function is given by a $`2\times 2`$ matrix,
$$G_{𝐤,\omega }^{(0)1}=(2SA_𝐤+i\eta )\sigma ^0+2SB_𝐤\sigma ^x+\omega \sigma ^z.$$
(23)
Here $`\sigma ^0`$ and $`𝝈`$ denote the identity and Pauli matrices respectively and $`\eta =0^+`$.
The self-energy is defined by the Dyson equation,
$$G_{𝐤,\omega }^1=G_{𝐤,\omega }^{(0)1}\mathrm{\Sigma }_{𝐤,\omega }^{},$$
(24)
and can be parameterized as
$$\mathrm{\Sigma }_{𝐤,\omega }=O_{𝐤,\omega }\sigma ^0+X_{𝐤,\omega }\sigma ^x+Z_{𝐤,\omega }\sigma ^z.$$
(25)
The leading order (in $`1/S`$) contributions to the self-energy can be divided into two parts
$`\mathrm{\Sigma }_{𝐤,\omega }=\mathrm{\Sigma }_𝐤^{(4)}+\mathrm{\Sigma }_{𝐤,\omega }^{(3)}.`$ (26)
Here $`\mathrm{\Sigma }_𝐤^{(4)}`$ denotes the vacuum polarization contribution that arises in first order perturbation theory in $`_4`$. It is frequency independent and purely real. On the other hand, $`\mathrm{\Sigma }_{𝐤,\omega }^{(3)}`$ denotes the contribution in second order perturbation theory of the three-magnon interaction $`_3`$. It incorporates the effects of magnon decay. Using Eq. (23), the $`\mathrm{\Sigma }_𝐤^{(4)}`$ contribution to the self-energy is found to be of the form
$`O_𝐤^{(4)}`$ $`=`$ $`A_𝐤+{\displaystyle \frac{2S}{N}}{\displaystyle \underset{𝐤^{}}{}}{\displaystyle \frac{1}{\omega _𝐤^{}}}[({\displaystyle \frac{1}{2}}B_𝐤+B_𝐤^{})B_𝐤^{}`$
$`+`$ $`(A_{𝐤𝐤^{}}B_{𝐤𝐤^{}}A_𝐤^{}A_𝐤)A_𝐤^{}],`$
$`X_𝐤^{(4)}`$ $`=`$ $`B_𝐤+{\displaystyle \frac{2S}{N}}{\displaystyle \underset{𝐤^{}}{}}{\displaystyle \frac{1}{\omega _𝐤^{}}}[(B_𝐤+B_𝐤^{})A_𝐤^{}`$
$`+`$ $`(A_{𝐤𝐤^{}}B_{𝐤𝐤^{}}A_𝐤^{}{\displaystyle \frac{1}{2}}A_𝐤)B_𝐤^{}],`$
$`Z_𝐤^{(4)}`$ $`=`$ $`0.`$ (27)
The contribution $`\mathrm{\Sigma }^{(3)}`$ is most easily evaluated in the Bogoliubov basis ($`\gamma `$) and is equal to
$`O_{𝐤,\omega }^{(3)}`$ $`=`$ $`{\displaystyle \frac{S}{16N}}{\displaystyle \underset{𝐤^{}}{}}\left\{\left[\mathrm{\Phi }^{(1)}(𝐤^{},𝐤𝐤^{})\right]^2+\left[\mathrm{\Phi }^{(2)}(𝐤^{},𝐤𝐤^{})\right]^2\right\}\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}+{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right),`$
$`X_{𝐤,\omega }^{(3)}`$ $`=`$ $`{\displaystyle \frac{S}{16N}}{\displaystyle \underset{𝐤^{}}{}}\left\{\left[\mathrm{\Phi }^{(1)}(𝐤^{},𝐤𝐤^{})\right]^2\left[\mathrm{\Phi }^{(2)}(𝐤^{},𝐤𝐤^{})\right]^2\right\}\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}+{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right),`$
$`Z_{𝐤,\omega }^{(3)}`$ $`=`$ $`{\displaystyle \frac{S}{16N}}{\displaystyle \underset{𝐤^{}}{}}\left\{2\mathrm{\Phi }^{(1)}(𝐤^{},𝐤𝐤^{})\mathrm{\Phi }^{(2)}(𝐤^{},𝐤𝐤^{})\right\}\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right),`$ (28)
where
$`\mathrm{\Phi }^{(1)}(𝐤^{},𝐤𝐤^{})`$ $`=`$ $`\left(C_𝐤^{}+C_{𝐤𝐤^{}}\right)\left(u_𝐤^{}+v_𝐤^{}\right)\left(u_{𝐤𝐤^{}}+v_{𝐤𝐤^{}}\right)2C_𝐤\left(u_𝐤^{}v_{𝐤𝐤^{}}+v_𝐤^{}u_{𝐤𝐤^{}}\right),`$
$`\mathrm{\Phi }^{(2)}(𝐤^{},𝐤𝐤^{})`$ $`=`$ $`C_𝐤^{}\left(u_𝐤^{}+v_𝐤^{}\right)\left(u_{𝐤𝐤^{}}v_{𝐤𝐤^{}}\right)+C_{𝐤𝐤^{}}\left(u_{𝐤𝐤^{}}+v_{𝐤𝐤^{}}\right)\left(u_𝐤^{}v_𝐤^{}\right).`$ (29)
## IV Dynamical Correlation Function
Inelastic neutron scattering experiments probe the dynamical structure factor $`S_{𝐤,\omega }^{\mu \nu }`$. The latter is defined as the Fourier transform of the dynamical spin-spin correlation function
$$S_{𝐤,\omega }^{\mu \nu }=_{\mathrm{}}^{\mathrm{}}\frac{dt}{2\pi }e^{i\omega t}S_𝐤^\mu (0)S_𝐤^\nu (t).$$
(30)
Here $`\mu ,\nu =(a,b,c)`$ label the various crystallographic axes and the Fourier-transformed spin operators are defined by $`S_𝐤^\mu =\frac{1}{\sqrt{N}}_𝐑S_𝐑^\mu e^{i𝐤𝐑}`$.
It is convenient to introduce time-ordered spin-spin correlation functions in the rotated coordinate system
$$F_{𝐤,\omega }^{\alpha \beta }=i_{\mathrm{}}^{\mathrm{}}𝑑te^{i\omega t}TS_𝐤^\alpha (0)S_𝐤^\beta (t),$$
(31)
where $`\alpha ,\beta =(x,y,z)`$ are the rotated coordinate axes (Eq. 5). The dynamical structure factor is related to the imaginary part of the time ordered correlation function in the following way
$`S_{𝐤,\omega }^{aa}`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}F_{𝐤,\omega }^{xx},`$ (32)
$`S_{𝐤,\omega }^{bb}`$ $`=`$ $`S_{𝐤,\omega }^{cc}={\displaystyle \frac{1}{\pi }}\mathrm{Im}\left[\mathrm{\Theta }_{𝐤+𝐐,\omega }^++\mathrm{\Theta }_{𝐤𝐐,\omega }^{}\right],`$ (33)
$`S_{𝐤,\omega }^{bc}`$ $`=`$ $`S_{𝐤,\omega }^{cb}={\displaystyle \frac{i}{\pi }}\mathrm{Im}\left[\mathrm{\Theta }_{𝐤+𝐐,\omega }^+\mathrm{\Theta }_{𝐤𝐐,\omega }^{}\right],`$ (34)
where
$$\mathrm{\Theta }_{𝐤,\omega }^\pm =\frac{1}{4}\left[F_{𝐤,\omega }^{zz}+F_{𝐤,\omega }^{yy}\pm i\left(F_{𝐤,\omega }^{zy}F_{𝐤,\omega }^{yz}\right)\right].$$
(35)
To proceed further, we expand the dynamical correlation functions in inverse powers of $`S`$ to order $`𝒪(S^0)`$. The corresponding results have been derived previously by Ohyama and Shiba. Shiba1 Here we merely quote their results for the sake of completeness. The transverse correlations are
$`F_{𝐤,\omega }^{xx}`$ $`=`$ $`{\displaystyle \frac{S}{2}}c_x^2\mathrm{Tr}\left[\left(\sigma ^0\sigma ^x\right)G_{𝐤,\omega }\right],`$
$`F_{𝐤,\omega }^{yy}`$ $`=`$ $`{\displaystyle \frac{S}{2}}c_y^2\mathrm{Tr}\left[\left(\sigma ^0+\sigma ^x\right)G_{𝐤,\omega }\right],`$ (36)
where the Green’s function is given by Eq. (24) and where
$`c_x`$ $`=`$ $`1{\displaystyle \frac{1}{4SN}}{\displaystyle \underset{𝐤}{}}\left(2v_𝐤^2u_𝐤v_𝐤\right),`$
$`c_y`$ $`=`$ $`1{\displaystyle \frac{1}{4SN}}{\displaystyle \underset{𝐤}{}}\left(2v_𝐤^2+u_𝐤v_𝐤\right).`$ (37)
We note that when squaring (37) only terms to order $`𝒪(S^1)`$ must be retained. The mixing of transverse and longitudinal fluctuations manifests itself in
$`i\left(F_{𝐤,\omega }^{yz}F_{𝐤,\omega }^{zy}\right)=`$ $`c_y\{P_{𝐤,\omega }^{(1)}\mathrm{Tr}\left[(\sigma ^0+\sigma ^x)G_{𝐤,\omega }\right]`$ (38)
$`+P_{𝐤,\omega }^{(2)}\mathrm{Tr}\left[\sigma ^zG_{𝐤,\omega }\right]\}.`$
Here the functions $`P_{𝐤,\omega }^{(1,2)}`$ are defined as
$`P_{𝐤,\omega }^{(1)}=`$ $`{\displaystyle \frac{S}{4N}}`$ $`{\displaystyle \underset{𝐤^{}}{}}\mathrm{\Phi }^{(1)}(𝐤^{},𝐤𝐤^{})\left(u_𝐤^{}v_{𝐤𝐤^{}}+v_𝐤^{}u_{𝐤𝐤^{}}\right)\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}+{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right),`$
$`P_{𝐤,\omega }^{(2)}=`$ $`{\displaystyle \frac{S}{4N}}`$ $`{\displaystyle \underset{𝐤^{}}{}}\mathrm{\Phi }^{(2)}(𝐤^{},𝐤𝐤^{})\left(u_𝐤^{}v_{𝐤𝐤^{}}+v_𝐤^{}u_{𝐤𝐤^{}}\right)\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right).`$ (39)
Finally the longitudinal correlations are decomposed in inverse powers of $`S`$ as $`F_{𝐤,\omega }^{zz}=F_{𝐤,\omega }^{(0)zz}+F_{𝐤,\omega }^{(1)zz}`$, where
$`F_{𝐤,\omega }^{(0)zz}={\displaystyle \frac{1}{2N}}{\displaystyle \underset{𝐤^{}}{}}\left(u_𝐤^{}v_{𝐤𝐤^{}}+v_𝐤^{}u_{𝐤𝐤^{}}\right)^2\left({\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}\omega i\eta }}+{\displaystyle \frac{1}{\omega _𝐤^{}+\omega _{𝐤𝐤^{}}+\omega i\eta }}\right),`$ (40)
$`F_{𝐤,\omega }^{(1)zz}={\displaystyle \frac{1}{2S}}\left\{\left(P_{𝐤,\omega }^{(1)}\right)^2\mathrm{Tr}\left[\left(\sigma ^0+\sigma ^x\right)G_{𝐤,\omega }\right]+\left(P_{𝐤,\omega }^{(2)}\right)^2\mathrm{Tr}\left[\left(\sigma ^0\sigma ^x\right)G_{𝐤,\omega }\right]+2P_{𝐤,\omega }^{(1)}P_{𝐤,\omega }^{(2)}\mathrm{Tr}\left[\sigma ^zG_{𝐤,\omega }\right]\right\}.`$ (41)
We note that the $`F^{(0)zz}`$ term does not require the knowledge of the bosonic self-energy and is basically a free boson result. For this reason, it is often included in linear spin wave calculation as a source of two magnon scattering, even though it is formally a higher order contribution in $`1/S`$. In what follows, we abide by this (in some sense inconsistent) convention and consider the contribution of Eq. 40 as part of linear spin-wave theory. As a consequence we then retain the $`F^{(1)zz}`$ contribution to the dynamical structure factor, although of higher order in $`1/S`$ (i.e. $`𝒪(S^1)`$) than the other terms we take into account.
The (unpolarized) inelastic neutron scattering cross section is given by
$`{\displaystyle \frac{d^2\sigma }{d\omega d\mathrm{\Omega }}}`$ $`=`$ $`|f_𝐤|^2{\displaystyle \underset{\mu \nu }{}}\left(\delta _{\mu \nu }\widehat{𝐤}_\mu \widehat{𝐤}_\nu \right)S_{𝐤,\omega }^{\mu \nu },`$ (42)
$`=`$ $`|f_𝐤|^2\left[\left(1\widehat{𝐤}_a^2\right)S_{𝐤,\omega }^{aa}+\left(1+\widehat{𝐤}_a^2\right)S_{𝐤,\omega }^{bb}\right],`$
where $`\widehat{𝐤}_\mu `$ is the $`\mu `$-component of the unit vector in $`𝐤`$ direction. The magnetic form factor $`f_𝐤`$ is determined by the magnetic ions. For Cu<sup>2+</sup>, the isotropic form factor has a relatively weak wave vector dependence within the first Brillouin zone and will be neglected from now on. Wilson
It is well known that the $`1/S`$ expansion preserves many physical properties “order by order” in $`1/S`$. For instance, it follows from Eqs. (27, 28) that the Goldstone modes persist beyond linear spin wave theory as one expects on physical grounds. A careful examination also shows that to order $`𝒪(S^0)`$ the spectral functions are positive and that the relation (17) holds. However, due to a lack of self-consistency the $`1/S`$ expansion leads to an (unphysical) unequal treatment of the one-magnon and two magnon scattering contributions to dynamical correlation functions. Shiba3 It is worthwhile to discuss this issue in more detail. The leading order contribution to the dynamical structure factor is due to coherent single magnon excitations and is of the form $`\delta \left(\omega \omega _𝐤\right)`$. The two magnon contribution due to longitudinal fluctuations (Eq. 40) gives rise to a scattering continuum of the form $`_𝐤^{}I(𝐤,𝐤^{})\delta \left(\omega \omega _𝐤^{}\omega _{𝐤𝐤^{}}\right)`$ with some function $`I(𝐤,𝐤^{})`$. The extent of the two magnon contribution in $`𝐤\omega `$ space is determined by the lower and upper bounds of the function $`\omega _𝐤^{}+\omega _{𝐤𝐤^{}}`$ for a given $`𝐤`$.
On general grounds, we expect the lower bound of the two magnon scattering continuum to be equal to or smaller than the “true” magnon dispersion $`\overline{\omega }_𝐤`$. In fact, the existence of a zero-momentum Goldstone mode guarantees that there exists a two magnon contribution at frequencies $`\overline{\omega }_𝐤+\overline{\omega }_\mathrm{𝟎}=\overline{\omega }_𝐤`$.
It is easy to see that this property does not hold order by order in a $`1/S`$ expansion. Indeed, the first order contribution in $`1/S`$ shifts the pole of the Green’s function and leads to a renormalization of the magnon dispersion. The renormalized dispersion $`\stackrel{~}{\omega }_𝐤`$ can be determined from the Dyson equation
$$G_{𝐤,\stackrel{~}{\omega }_𝐤}^1=G_{𝐤,\stackrel{~}{\omega }_𝐤}^{(0)1}\mathrm{\Sigma }_{𝐤,\stackrel{~}{\omega }_𝐤}^{}=0.$$
(43)
However, to order $`𝒪(S^0)`$ the threshold of the two magnon contribution is still determined by the bare dispersion relation $`\omega _𝐤`$. This results in an unphysical behavior, where the two magnon scattering continuum is separated from the single magnon dispersion by a gap. In order to avoid this problem, we impose the following self-consistency condition: the linear spin wave dispersion $`\omega _𝐤`$ used in Eqs. (28, 39, 40) is to be replaced by the renormalized dispersion $`\stackrel{~}{\omega }_𝐤`$.
## V Dynamical Properties of $`\text{Cs}_2\text{CuCl}_4`$
So far our discussion of the $`1/S`$ expansion has been fairly general. In order to make contact with the experiments on Cs<sub>2</sub>CuCl<sub>4</sub> we now set the exchange constants to their appropriate values Coldea2 ; Coldea3 and fix $`S=1/2`$. We then evaluate the dynamical structure factor at a given wave vector numerically. Complex integrals such as Eqs. (28, 39) are evaluated by summing the imaginary part of the integrands over a frequency grid of 1200 points and of 1000 $`\times `$ 1000 points in wave vector space. The real parts are then determined from the Kramers-Kronig relations. The aforementioned self-consistency condition is implemented by calculating the full Green’s function iteratively on a 100 $`\times `$ 100 grid in the Brillouin zone. We observe satisfactory numerical convergence after about three iterations.
We first turn to the magnon dispersion. The linear spin wave result $`\omega _𝐤`$ vanishes at the center of the paramagnetic Brillouin zone. The corresponding Goldstone mode is associated with small fluctuations of the ordered moment within the cycloidal plane. In helimagnets, the spectrum often exhibits a second Goldstone mode at the ordering wave vector. This gapless mode is due to fluctuations of the plane of the cycloid. In the case at hand, the easy-plane anisotropy generated by the DM term forces the cycloidal structure to lie in the $`bc`$ plane and creates an excitation gap at the ordering wave vector $`𝐐`$.
The renormalization of the magnon dispersion within the framework of the $`1/S`$ expansion is obtained from the poles of the Green’s function (Eq. 43). In Fig. 4 we compare the results of the $`1/S`$ expansion with the linear spin wave theory. It is customary to quantify the effects of the “quantum” renormalization of the magnon dispersion by parametrizing the latter in terms of “effective” exchange constants $`\stackrel{~}{J},\stackrel{~}{J}^{},\stackrel{~}{D}`$ and comparing them with the “bare” parameters $`J`$, $`J^{}`$ and $`D`$.
Experimentally, the quantum renormalization is found to be rather large, namely $`\frac{\stackrel{~}{J}}{J}=1.63(5)`$ and $`\frac{\stackrel{~}{J}^{}}{J^{}}=0.84(9)`$. The renormalization of $`D`$ was not established. The $`1/S`$ expansion yields the significantly smaller renormalizations $`\frac{\stackrel{~}{J}}{J}=1.131`$, $`\frac{\stackrel{~}{J}^{}}{J^{}}=0.648`$ and $`\frac{\stackrel{~}{D}}{D}=0.72`$. The difference between the theoretical and experimental values indicates that the leading order in a $`1/S`$ expansion underestimates fluctuation effects. On the other hand one should note that the $`1/S`$ expansion gives a result of $`0.031`$ for the incommensuration, which is very close to the experimentally observed value of $`ϵ=0.030(2)`$.
Before turning to a comparison of our results for the dynamical structure factor with the experimental results, we briefly review some facts about excitations in helimagnets. Generally it is useful to distinguish between three spin wave modes. In the case at hand, the “principal” mode $`\omega _𝐤^0=\stackrel{~}{\omega }_𝐤`$ is polarized along the $`a`$ axis and is probed by the $`S_{𝐤,\omega }^{aa}`$ component of the dynamical structure factor (Eq. 32). The two “secondary” spin wave modes $`\omega _𝐤^\pm =\stackrel{~}{\omega }_{𝐤\pm 𝐐}`$ are images of the main mode but their momenta are shifted by $`\pm 𝐐`$. They are polarized in the $`bc`$ plane (Eqs. 33, 34). In linear spin wave theory, the three spin wave modes give rise to sharp $`\delta `$ functions and carry a large part of the spectral weight.
In addition to the single magnon modes there are multi magnon scattering continua. Whenever the magnon dispersion lies within a scattering continuum, the single magnon excitation gets broadened and acquires a finite line width. On the other hand, when the magnon dispersion lies at the threshold of a scattering continuum, there is no significant decay and the single magnon mode remains sharp.
The unpolarized dynamical structure factor (where the various polarizations are added according to Eq. 42) is shown in Fig. 3 as a function of energy and momentum for a particular “cut” of momentum transfers. The cut along the $`b`$ direction, i.e. from $`(000)`$ to $`(010)`$, shows large modulations of the dispersion relation due to the strong intra-chain correlations. Near the ordering wave vector $`𝐐`$ the scattering intensity increases sharply. For momentum transfers perpendicular to the chains, (i.e. along the $`(01\eta )`$ direction), the single particle modes are seen to be resolution limited. The two in-plane modes become degenerate and their dispersions are nearly featureless, whereas the out-of-plane fluctuations dip to zero energy at $`(011)`$, in accordance with Goldstone’s theorem. Along the $`(0\eta \eta )`$ direction the spectrum is symmetric across the Brillouin zone boundary. Additional structures due to two magnon scattering are clearly visible at higher energies along the $`(0\eta 0)`$ and $`(0\eta \eta )`$ directions.
In order to illustrate how the spectral weights associated with the single-particle excitations are affected by the magnon interactions, we have estimated their contributions for each polarization to the integrated spectral weights. The total integrated intensity of each polarization is given by “equal-time” correlation functions,
$`I_𝐤^0`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle 𝑑\omega F_{𝐤,\omega }^{xx}},`$
$`I_𝐤^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle 𝑑\omega \mathrm{\Theta }_{𝐤\pm 𝐐,\omega }^\pm }.`$ (44)
The one-magnon contribution to the integrated intensity of each polarization is then determined by integrating the respective correlation function in the vicinity of the single particle dispersions. In practice we find that integrating the peaks assuming a Lorentzian form is a poor prescription for strongly damped peaks. Instead we numerically integrate the intensity over an energy window of three times the width at half maximum
$`R_𝐤^0`$ $`=`$ $`{\displaystyle \frac{1}{I_𝐤^0}}{\displaystyle _{\omega _𝐤^01.5\mathrm{\Delta }\omega _𝐤^0}^{\omega _𝐤^0+1.5\mathrm{\Delta }\omega _𝐤^0}}𝑑\omega {\displaystyle \frac{F_{𝐤,\omega }^{xx}}{\pi }},`$
$`R_𝐤^\pm `$ $`=`$ $`{\displaystyle \frac{1}{I_𝐤^\pm }}{\displaystyle _{\omega _{𝐤\pm 𝐐}^\pm 1.5\mathrm{\Delta }\omega _{𝐤\pm 𝐐}^\pm }^{\omega _{𝐤\pm 𝐐}^\pm +1.5\mathrm{\Delta }\omega _{𝐤\pm 𝐐}^\pm }}𝑑\omega {\displaystyle \frac{\mathrm{\Theta }_{𝐤\pm 𝐐,\omega }^\pm }{\pi }}.`$ (45)
The results are shown in Fig. 5. We see that the integrated spectral weight is concentrated in the vicinities of the ordering wave vector $`𝐐`$ and $`(0\frac{1}{2}\frac{1}{2})`$ and is largely suppressed near the $`\mathrm{\Gamma }`$ point. The weights associated with single magnon excitations are strongly suppressed for the secondary modes. This is a consequence of the non-collinearity of the magnetic order. The in-plane modes are significantly damped as a result of the coupling between longitudinal and transverse fluctuations. Such a coupling is not present for the out-of-plane mode and therefore the principal mode carries generally more spectral weight. Nevertheless the fraction of spectral weight associated with single-particle excitations decreases significantly whenever the renormalized spin wave dispersion is pushed upwards in energy for a given momentum. For instance, near wave vector $`(0,0.8,0)`$ the principal spin wave mode lies within the two magnon continuum and as a result less than 50 % of the spectral weight is attributed to the one-magnon excitation.
The scattering intensity can also be studied by performing a wave vector average,
$$I_T(\omega )=\frac{1}{N}\underset{𝐤}{}\underset{\mu }{}S_{𝐤,\omega }^{\mu \mu }.$$
(46)
By the frequency sum rule, the scattering intensity (Eq. 46) integrated over all energies (including the elastic Bragg peaks at $`\omega =0`$) has to equal $`S(S+1)`$. However, this sum rule does not hold “order by order” in perturbation theory. For instance, the total intensity within linear spin wave theory exceeds the sum rule by $`\mathrm{\Delta }S(1+2\mathrm{\Delta }S)`$. Bearing this caveat in mind, the sum rule is a useful tool for comparing the one and two magnon contributions as well as analyzing the shift in spectral weight. In Fig. 6, we plot the scattering intensities as functions of energy within linear spin wave theory and the $`1/S`$ expansion. In linear spin wave theory the integrated intensity exhibits cusps, which are associated with van-Hove singularities in the single particle density of states. In the $`1/S`$ expansion such sharp features are absent. Above approximately $`0.5`$ meV the one magnon contribution vanishes and the scattering intensity is entirely due to multi magnon states.
To quantify the shift of the spectral weight we calculate the first moment of the normalized scattering intensity $`\omega `$. We find that the linear spin wave theory value $`\omega =0.35`$ meV is renormalized upwards to $`\omega =0.40`$ meV in the $`1/S`$ expansion. This observation is in line with the expectation that the higher orders of the $`1/S`$ expansion induce a transfer of spectral weight to higher energies via multi magnon scattering processes. In fact, as shown in Fig. 6 the two magnon contribution to the overall intensity is $`29\%`$ in linear spin wave theory but $`46\%`$ in the $`1/S`$ expansion.
### V.1 Excitation line shapes
In order to exhibit the properties of the dynamical structure factor in greater detail we have generated a series of scans in $`𝐤\omega `$ space. The inelastic neutron scattering measurements on Cs<sub>2</sub>CuCl<sub>4</sub> were not performed at constant momentum transfer but followed various trajectories in energy-wave vector space. We have generated our theoretical scans using the known parameterizations of the scans A to J of Ref. Coldea3, in $`𝐤\omega `$ space, which we summarize in Table. 1. We refer the reader to Ref. Coldea3, for further details. The various scans are shown in Fig. 8. Also shown are the regions in which significant magnetic scattering is observed experimentally and the location of the main peaks. For comparison we plot the principal and secondary spin wave dispersions obtained from the $`1/S`$ expansion. As we have already emphasized, the $`1/S`$ expansion underestimates the quantum renormalization of the exchange constants and as a result the agreement of the calculated spin wave dispersions with the main peaks observed experimentally is poor.
The experimental energy and momentum resolutions have been accounted for to make contact with experiment. We find that the effects of the finite energy resolution of $`\mathrm{\Delta }E=0.016`$ meV are generally outweighed by the effects of the finite momentum resolution. This is a consequence of the large modulation of the spin wave dispersion along the chain direction, i.e $`(0k0)`$, (whose slopes can reach $`\frac{\mathrm{\Delta }E}{\mathrm{\Delta }k}1.6`$ meV), which causes an amplification of the effects of the momentum resolution. Given that the spin waves are nearly dispersionless along the $`(00l)`$ direction, we have only taken into account the spatial resolution along the chain direction.
To illustrate this point, let us consider the results for scans B, E, G and H shown in Fig. 9. The insets of panel (4) show the results of both linear spin wave theory and the $`1/S`$ expansion for a hypothetical energy resolution of $`\mathrm{\Delta }E=0.002`$ meV which has been introduced to make the various delta function peaks visible (the momentum resolution is set to zero $`\mathrm{\Delta }k=0`$). First we consider the results for scan H (Panel 4 of Fig. 9). Linear spin wave theory predicts peaks at approximately $`0.27`$ meV and $`0.37`$ meV corresponding to the degenerate spin wave modes $`\omega ^+,\omega ^0`$ and to $`\omega ^{}`$ respectively. The $`1/S`$ correction yields a slight upward shift in the energy of these peaks. In both linear spin-wave and $`1/S`$ calculations, the two magnon scattering continuum is found to carry nearly a quarter of the integrated spectral weight. Taking into account the finite momentum resolution (the width at half maximum is $`\mathrm{\Delta }k=0.057`$) we find that the sharp peaks get broadened very significantly as is shown in panel (4). The dynamical structure factor now exhibits an extended continuum in which the single-particle excitation can no longer be resolved and merges smoothly with the two magnon continuum. This result is qualitatively similar to the experimental observations shown for comparison in panel (2) of Fig. 10.
Next we turn to scan G (Panel 3 of Fig. 9), which probes the vicinity of the wave vector $`(0,0.5,1.5)`$. Experimentally a resolution-limited peak is observed at an energy of $`0.107(10)`$ meV in this region of intense scattering, see panel (1) of Fig.10. However, about two thirds of the spectral weight is associated with a scattering continuum at higher energies. Both linear spin wave theory and the $`1/S`$ expansion predict sharp peaks in this region of the Brillouin zone. The $`1/S`$ expansion gives a spin wave peak at $`\omega ^0=0.18`$ meV carrying nearly half of the spectral weight and two further peaks at energies around $`0.25`$ meV corresponding to the two secondary spin wave modes. The two magnon scattering continuum extends up to $`0.9`$ meV and carries nearly a quarter of the spectral weight. We emphasize that, in contrast to $`\omega ^\pm `$, the principal mode $`\omega ^0`$ is close to a saddle point and therefore is nearly dispersionless. In Panel 3 the finite energy and momentum resolutions are taken into account. We see that the almost dispersionless main mode remains sharp but the secondary modes can no longer be resolved and are found to merge with the two magnon continuum. Irrespective of the discrepancies between the results of the $`1/S`$ expansions and the experimental data, our calculation suggests that the lower boundary of of the measured scattering continuum in scan G could be due to unresolved transverse magnons. Such a scenario had been previously considered and ruled out on the basis of the smallness of the ratio $`I_{\mathrm{sec}}/I_{\mathrm{pri}}`$ of spectral weights of the secondary modes to the principal mode predicted by linear spin wave theory. Coldea3 However, the results of the $`1/S`$ expansion show that spin wave interactions lead to an enhancement of this ratio for the G scan.
Next, we examine scan E (Panel 2), which probes wave vectors near $`𝐤=(0,0.25,1)`$. Linear spin wave theory predicts coherent peaks at $`\omega ^0=0.35`$ meV for the principal mode and at $`\omega ^{}=0.44`$ meV and $`\omega ^+=0.33`$ meV for the secondary modes (see the inset in Panel 2). The two magnon scattering continuum is relatively weak and carries only about 23 % of the total spectral weight. In the framework of the $`1/S`$ expansion the principal mode is pushed upwards in energy to $`\omega ^0=0.42`$ meV and occurs very close to the secondary mode $`\omega ^{}=0.45`$ meV. The other secondary mode $`\omega ^+`$ is shifted very significantly to $`0.39`$ meV, but carries only a minute fraction of the spectral weight. The two magnon continuum is also shifted upwards in energy and carries approximately a quarter of the total spectral weight. Once again the spin wave dispersion is close to a saddle point and as a result the effects of the finite momentum resolution are small. The main feature in the structure factor is a broad peak formed by the two unresolved $`\omega ^{}`$ and $`\omega ^0`$ modes. This is quite similar to what is observed experimentally (Fig. 5(E) of Ref. Coldea3, ). It is then tempting to speculate that the experimentally observed single peak is a result of the accidental near degeneracy of the $`\omega ^{}`$ and $`\omega ^0`$ modes in the vicinity of $`𝐤=(0,0.25,1)`$. This would explain both the absence of the $`\omega ^{}`$ peak in the experimental data and the anomalously large intensity of the observed peak.
In Panel 1 of Fig. 9 we plot the dynamical structure factor for scan B near the wave vector $`(2,0.25,0)`$. Here the polarization factor ($`\widehat{𝐤}_a`$) in (42) leads to a strong suppression of the out-of-plane fluctuations and the scattering is almost entirely due to the in-plane $`\omega ^\pm `$ spin wave modes. The magnon interactions renormalize $`\omega ^+`$ upwards in energy to approximately $`0.42`$ meV, whereas the $`\omega ^{}`$ mode disappears in the two magnon scattering continuum. A careful analysis shows that the narrow peak at $`0.55`$ meV is not due to a single-particle excitation but is a feature in the two magnon scattering continuum.
The dominant contribution to the dynamical structure factor in scan A in the vicinity of the wave vector $`(1.5,0.3,0)`$ comes from in-plane fluctuations because the polarization factor $`\widehat{𝐤}_a`$ suppresses out-of plane fluctuations. As can be seen in Fig. 11, the magnon interactions lead to a spectral weight transfer to higher energies. The peaks near $`0.8`$ meV and $`0.85`$ meV can be traced back to single-particle poles in the Green’s function. These poles are unphysical and are a result of the uncontrolled nature of the $`1/S`$ expansion for small values of $`S`$. It is easily seen from the Dyson equation (24) that a large self-energy at a given wave vector can lead to “extra” poles in the Green’s function at high energies above the two magnon continuum. The inclusion of higher order terms in the $`1/S`$ expansion would provide decay mechanisms at all energies and lead to a broadening of these high-energy peaks in the dynamical structure factor.
Last but not least let us consider the vicinity of $`(0.8,0.4,0)`$ (scan C). As is shown in Fig. 11 the principal spin wave mode $`\omega ^0`$ is renormalized down to a slightly lower energy of approximately $`0.42`$ meV. The $`\omega ^+`$ mode, which occurs at $`0.35`$ meV in linear spin wave theory, disappears entirely in the two magnon-continuum. The feature near $`0.60`$ meV can again be understood in terms of an enhancement of the two magnon density of states. Comparing with the neutron scattering data (Fig. 5(C) of Ref. Coldea3, ), the structure factor shows features quite similar to the experimentally observed continuum. However, the scattering continuum occurs at energies nearly $`0.10`$ meV lower than what is observed experimentally.
## VI Conclusions
In this work we have used nonlinear spin wave theory to determine the dynamical structure factor in the ordered phase of the spin-1/2 helimagnet Cs<sub>2</sub>CuCl<sub>4</sub>. We have taken into account the first subleading contribution in a $`1/S`$ expansion, which incorporates interactions between magnons and generates magnon decay processes as well as multi magnon scattering continua. Both effects are particularly pronounced in Cs<sub>2</sub>CuCl<sub>4</sub> due to the non-collinear spin ordering, the low spin value and geometrical frustration.
We found that the results of nonlinear spin wave theory explain on a qualitative level many of the features observed in neutron scattering experiments. We find a strong scattering continuum in the dynamical structure factor similar to the experimental observations. Our calculations suggest the possibility that some of the spectral weight at the low-energy boundary of the experimentally observed scattering continuum in scan G could be due to single particle excitations that are unresolved.
In the vicinity of saddle points of the spin wave dispersion relation the single-particle excitations are only weakly affected by the instrumental resolution and hence exhibit sharper peaks in the dynamical structure factor.
In spite of the qualitative agreement of the theory with experiments, crucial discrepancies remain. First and foremost, nonlinear spin wave theory fails to account for the large “quantum renormalization” of the main exchange parameter. This indicates that (to order $`𝒪(S^0)`$) the $`1/S`$ expansion still underestimates the effects of quantum fluctuations. Furthermore, there are significant quantitative differences between our calculations and the experimentally observed structure factor. One may speculate that a better agreement with experiment could be achieved by taking higher-order terms in the $`1/S`$ expansion into account.
The main lesson to be learned from our calculations is that Cs<sub>2</sub>CuCl<sub>4</sub> falls somewhere in between the two theoretical scenarios that have been proposed previously. Our analysis shows that the physics of order plays an essential part in understanding the dynamic response Cs<sub>2</sub>CuCl<sub>4</sub> at low temperatures: a large fraction of the spectral weight is carried by spin wave modes, which occur over a large range of frequencies. This is a strong indication that a putative spin-liquid ground state is plainly not a good starting point for the description of the ordered phase of Cs<sub>2</sub>CuCl<sub>4</sub>. On the other hand we have seen that (in low orders in $`1/S`$) nonlinear spin wave theory significantly underestimates the effects of quantum fluctuations and hence expansions around the ordered state also fail to account for the experimental observations.
Nonlinear spin wave theory can also be applied to investigate the effects of magnetic fields. It is known that in the presence of a field linear spin wave theory is generally a very poor approximation as it excludes all-important magnon decay processes.Zhitomirsky01 A self-consistent study of magnetic field effects in Cs<sub>2</sub>CuCl<sub>4</sub> is currently under way.EJV During completion of this work, we became aware of a parallel effort which reaches similar conclusions.OtherWork
###### Acknowledgements.
The work was supported by the EPSRC under Grant GR/R83712/01. We are grateful to John Chalker and Alan Tennant for valuable discussions. Particular thanks are due to Radu Coldea for numerous helpful discussions and suggestions as well as providing us with figures 8 and 10.
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# An Infalling Torus of Molecular Gas Around the Ultra-Compact Hii Region G28.20-0.05
## 1 Introduction
Much research has been done attempting to extend the now well-defined process of low-mass-star formation to intermediate- and high-mass-star formation. The accepted standard model for low-mass-star formation, as described in Shu et al. (1987), involves a molecular cloud core which collapses to form a protostar with an accretion disk. A bipolar outflow forms. Accretion through the disk slows. Eventually, the cloud is dispersed and the protostar quasi-statically contracts onto the main sequence, radiating away its gravitational energy. Kahn (1974), however, pointed out that, because of their enormous luminosity, the Kelvin-Helmholtz time-scale for contraction of a massive star is much shorter than for low mass stars. Massive stars may radiate away their gravitational energy so quickly that they approach the main sequence and become very luminous before the surrounding protostellar envelope has finished collapsing (Garay & Lizano, 1999). Kahn (1974) and Wolfire & Cassinelli (1987) calculated that in this case the pressure of the stellar radiation may be adequate to stop accretion under certain circumstances. Observationally, however, several examples are now known of stars or protostars with masses of order 10 M or less, i.e. early B-type, which appear to be undergoing accretion through a disk (Zhang et al., 1998b, 2002; Shepherd & Kurtz, 1999; Chini et al., 2004; Beuther et al., 2004). Most of these sources are known to also include an outflow. Thus, some aspects of the highly successful standard model of low-mass-star formation have been observed in stars up to stellar masses of 10 M, despite the potential radiation pressure problem.
At the high end of the mass scale, a very different type of accretion scenario has been proposed to explain observations of the UCHii region G10.6-0.4. The ionized gas in G10.6 has a radius of about 6000 AU (0.03 pc), and contains a grouping of massive young stars, with a luminosity of $`9.2\times 10^5`$L, a mass of 150 M, emitting $`2.2\times 10^{50}`$ Lyman continuum photons per second, including at least one early O-type star (Sollins et al., 2005). Outside the UCHii region, molecular gas can be seen to be infalling and rotating in absorption in the NH<sub>3</sub>(3,3) line (Ho & Haschick, 1986; Keto et al., 1987, 1988). This molecular accretion flow proceeds all the way up to the ionization front (Sollins et al., 2005). In the ionized gas, H66$`\alpha `$ observations show that the ionized gas is also moving inward and settling into an ionized disk in the center of the UCHii region (Keto & Wood, 2005). Keto (2002b) proposed a model in which the ionized gas is gravitationally trapped by the mass of the central stars, allowing the molecular accretion flow to pass through a stalled ionization front and continue as an ionized accretion flow. The particular combination of mass and luminosity of the stellar cluster, and angular momentum in the molecular gas puts G10.6 in a regime in which the thermal pressure of the ionized gas is confined by the gravitational force of the central stars, and the molecular gas does not settle into a disk, but rather a rotating, somewhat flattened, quasi-spherical accretion flow. This is a very different scenario than the low-mass standard model in which molecular gas accretes quasi-statically through a rotationally supported accretion disk.
In order to investigate a regime of massive-star formation in between the grouping of O-stars in G10.6 and the individual early B-type stars seen to have disks and outflows, we have observed G28.20-0.05. G28 is an UCHii region at a distance of $`5.7^{+0.5}/_{0.8}`$ kpc (Fish et al., 2003), with a radius of about 3400 AU (0.017 pc). Using the formula of Casoli et al. (1986), the far infrared luminosity based on IRAS fluxes is $`1.6\times 10^5`$ L, corresponding to one O8V star, which would have a mass of 31 M, and an ionizing flux of $`7.4\times 10^{48}`$ Lyman continuum photons per second(Vacca et al., 1996). Of course, if more than one star is present, the luminosity does not uniquely determine the mass or flux of ionizing photons. Using a broader range of infrared measurements, Walsh et al. (2003) get $`1.8\times 10^5`$ L, consistent with the Casoli et al. (1986) value when scaled to our assumed distance. When compared to G10.6, G28 has a lower luminosity, smaller radius, and presumably different central mass. These factors should put G28 into a different accretion regime than the previous work described above.
## 2 Observations
We observed the UCHii region G28.20-0.05 with the NRAO Very Large Array (VLA)<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. on four occasions in 2003, with the phase center at $`\alpha (2000)=18^\mathrm{h}42^\mathrm{m}58^\mathrm{s}.10,\delta (2000)=4^\mathrm{o}13^{}57^{\prime \prime }.87`$. We observed the (1,1), (2,2), and (3,3) inversion lines of NH<sub>3</sub> in the D-configuration, yielding a resolution of about $`3^{\prime \prime }`$ in natural weighted maps. We also observed the (2,2) and (3,3) lines in the hybrid BnA configuration yielding a resolution of about $`0^{\prime \prime }.3`$ in uniform weighted maps.
We observed 3C286 to set the absolute flux scale, and 1849+005 to calibrate the phases. The passband response was calibrated using 3C273 for the (1,1) and (2,2) D array data, and 3C454.3 for the rest of the data. After initial flux and phase calibration, all visibilities were also self-calibrated for phase and amplitude, using as the model a map made from the velocity-integrated visibilities. The resolutions and noise levels, as well as some physical parameters of the lines observed, in each of the maps are listed in Table LABEL:tab:data.
As in earlier work on the inversion lines of NH<sub>3</sub> as seen around UCHii regions (Keto et al., 1988; Sollins et al., 2005), we make this point regarding spatial resolution, and the detectability of the line in either emission or absorption. We achieve a typical sensitivity in a single channel of the line data of 4 mJy beam<sup>-1</sup>. At a resolution of $`3^{\prime \prime }`$, this flux sensitivity means that the $`3\sigma `$ detection limit in terms of temperature is about 4 K. So thermal emission of any appreciable optical depth from a hot molecular core whose actual temperature is around 100 K should be easily detectable. For the higher resolution data, the flux sensitivity is similar, but the synthesized beam is a factor of 10 smaller, resulting in a factor of 100 worse temperature sensitivity. No thermal emission, at any optical depth, will be detectable from a 100 K hot molecular core when the $`3\sigma `$ detection limit is 400 K, as in the $`0^{\prime \prime }.3`$ resolution data. The sensitivity to absorption however, is quite high in the high resolution maps. Our continuum map, based on the BnA array (2,2) data, has a peak brightness temperature of 7500 K. The $`1\sigma `$ noise level in the line channel maps is 130 K, so very optically thick absorption could be detected at signal to noise ratios exceeding 50. We thus expect to detect both line absorption and emission in the D array data, but only absorption in the BnA array data.
## 3 Results
### 3.1 The High Resolution Data
The continuum map derived from the BnA array NH<sub>3</sub>(2,2) observations is shown in Figure 1. The source is resolved and there are two peaks separated by $`0^{\prime \prime }.4`$, just over one synthesized beam. At the 50% contour, the continuum emission is $`0^{\prime \prime }.6\times 0^{\prime \prime }.8`$. The total flux in the central continuum source is 0.98 Jy. The peak is 181 mJy beam<sup>-1</sup>, corresponding to a brightness temperature of 5700 K. Assuming the electron temperature is 10000 K, the peak optical depth is $`\tau _\nu =0.84`$, and the peak emission measure is $`2.0\times 10^9`$ pc cm<sup>-6</sup>. Assuming uniform density and a recombination coefficient of $`3\times 10^{13}\mathrm{cm}^3\mathrm{s}^1`$ (Keto, 2003), the electron density is $`2.1\times 10^6`$ cm<sup>-3</sup>, the mass of ionized gas is 0.54 M, and the flux of ionizing photons necessary to balance recombinations in the UCHii region is $`1.6\times 10^{49}`$ per second. If from a single star, this ionizing flux would correspond to a O6.5 V star of mass 41 M and luminosity $`3.1\times 10^5`$ L. Since this does not match the known luminosity, we suspect that the central object is a multiple star system, with no stellar member as early as O6.5. It should also be noted that there are several assumptions in this calculation (temperature of the ionized gas, path length associated with the observed continuum flux, uniform density structure) which could change the resulting spectral type.
The key feature of our BnA array line data is that it divides neatly into two distinct components, one centered at 90 km s<sup>-1</sup> which we associate with an outflow or expansion, and one centered at 97 km s<sup>-1</sup> which we associate with infall. The two components are clearly separated in the NH<sub>3</sub>(2,2) position-velocity diagrams in Figures 2, 3 and 4. The lower resolution data fix the ambient velocity of the cloud at 95.5 km s<sup>-1</sup>, so one component is blue-shifted, while the other is red-shifted, hence the infall and outflow interpretation. Figures 5 and 6 also help to distinguish the two components, and will be discussed in detail below.
#### 3.1.1 The 90 km s<sup>-1</sup> Outflow Component
The outflow component is cool and optically thin. It is detected over the entire face of the continuum source, and its mean velocity varies little as a function of position. The first and second moments of the 90 km s<sup>-1</sup> component can be seen in the first two panels in Figure 5. The moments of the 90 km s<sup>-1</sup> component are integrated from 84 km s<sup>-1</sup> to 92 km s<sup>-1</sup>. Compared to the 97 km s<sup>-1</sup> component, both moments of the 90 km s<sup>-1</sup> component vary little with position. But as seen in Figure 6, the first moment of the 90 km s<sup>-1</sup> component does have a clear pattern. The velocity pattern shows parallel stripes monotonically increasing from 89.7 km s<sup>-1</sup> on one side to 91.4 km s<sup>-1</sup> on the other side. The optical depth of the line in the 90 km s<sup>-1</sup> component is relatively low. There is no detectable absorption in the first satellite hyperfine line corresponding to the 90 km s<sup>-1</sup> component, meaning that the optical depth in the main line is everywhere less than 2.7. Because the 90 km s<sup>-1</sup> component is detected in (2,2) but not (3,3) we can put an upper limit on its temperature of 30 K. It is important to note that the 90 km s<sup>-1</sup> component extends over the entire face of the UCHii region.
#### 3.1.2 The 97 km s<sup>-1</sup> Infall Component
The rest of the panels in Figure 5 show the kinematics of the 97 km s<sup>-1</sup> component. The most prominent feature in Figure 5 is the sharp NW-SE line in the BnA array maps of the first and second moments. This is the edge of the 97 km s<sup>-1</sup> component. Because the 90 km s<sup>-1</sup> component extends over the entire face of the continuum emission, we know the sharp line must be due to a cut off in the absorption from the 97 km s<sup>-1</sup> component. In the direction perpendicular to the edge, the beam is $`>0^{\prime \prime }.2`$, so the sharpness of the edge is striking. This edge suggests the presence of a disk-like structure, perhaps a toroid, surrounding the UCHii region. If the toroid were inclined, it would naturally produce absorption over only half of the face of the UCHii region.
The spatial pattern of the mean velocity of the 97 km s<sup>-1</sup> component shows infall and possibly weak rotation. The pattern, from northwest to southeast goes from 95 km s<sup>-1</sup> to 97 km s<sup>-1</sup>and back down to 94 km s<sup>-1</sup>. This is reminiscent of the “off-center bulls-eye” characteristic of combined infall and rotation discussed in Keto et al. (1988) and observed in G10.6-0.4 (Keto et al., 1987; Keto, 2002a; Sollins et al., 2005). If rotation is responsible for the offset of the position of fastest infall, then the axis of rotation must be northeast-southwest, with the projection into the plane of the sky of the angular momentum vector pointing northeast. Position-velocity cuts in both the direction of the axis of rotation (NE-SW) and perpendicular to the axis of rotation (NW-SE) are shown in Figure 2. The most red-shifted infalling gas can be seen at 100 km s<sup>-1</sup>. While the offset of the position of fastest infall in the first moment map is conspicuous, it is difficult to see any additional evidence of rotation in the position-velocity diagram.
The infalling gas is warm and very optically thick. Figures 3 and 4 show the same position-velocity cuts as Figure 2, but in the apparent optical depths, and the hyperfine optical depths instead of in the actual line absorption. (For an explanation of apparent and hyperfine optical depths, see the Appendix.) The separations between the main hyperfine component and first satellite hyperfine component are 16.6 km s<sup>-1</sup> and 21.5 km s<sup>-1</sup> for the (2,2) and (3,3) lines respectively, and in the (2,2) the outer satellite line is at a separation of 25.8 km s<sup>-1</sup>. Only the most optically thick gas could ever be detected in the much weaker satellite lines. It is the red-shifted, infalling gas which is detectable in the satellite lines, and is the most optically thick. Since infall models generally predict strong central condensation (Shu, 1977; Terebey et al., 1984), we expect the most red-shifted infalling gas to be the densest, and most optically thick, and it is, with peak hyperfine optical depths of 16 in the (2,2) and 47 in the the (3,3). The (2,2) and (3,3) peak optical depths correspond to a rotational temperature of 280 K, much warmer than the 30 K upper limit for the outflow component, which is detected only in (2,2). Using the 280 K temperature and an abundance of NH<sub>3</sub> relative to $`\mathrm{H}_2`$ of $`10^7`$, we calculate the mass of the gas detected in the infall component to be 9 M, using either the (3,3) or the (2,2) optical depth, i.e. the two lines are consistent. Note that this is only a fraction of the mass of the entire infall component, since only that part of the infall component which lies in front of the continuum source can be observed. Thus mass of the infall component should be at least 18 M, or perhaps more. The largest source of uncertainty here is the abundance, which may be uncertain by as much as an order of magnitude.
### 3.2 The Low Resolution Data
The molecular core surrounding the UCHii is detected in all three lines we observed in the low angular resolution mode and shows typical evidence of internal heating. Figure 7 shows 3 <sup>′′</sup> resolution, velocity-integrated maps in all three lines. The (1,1) map includes the most extended emission, while the (3,3) emission is the most compact. This is consistent with central heating of the molecular core from the UCHii region, similar to the molecular gas in the regions surrounding massive protostars, such as AFGL5142 (Zhang et al., 2002). All three line maps show two peaks immediately next to the UCHii region, the brighter to the northeast, the fainter to the southwest. We use these peaks to define an axis for a position-velocity cut running from 30 west of south, to 30 east of north. The more extended emission runs northwest to southeast, so we use that to define another position-velocity cut. Both cuts are indicated in Figure 7. The absorption toward the UCHii region can be seen in (2,2) and (3,3) as negative contours at the center of the zeroth moment maps.
The position-velocity cut in the (1,1) data establishes the velocity of the ambient cloud. Figure 8 shows the position-velocity cuts indicated in Figure 7. The first satellite hyperfine components of the NH<sub>3</sub>(1,1) line are separated from the main line by $`\pm `$ 7.8 km s<sup>-1</sup>, and the main line has an intrinsic line-strength only 3.6 times that of the inner satellites. Since the kinematics in this object involve relative velocities greater than 8 km s<sup>-1</sup>, and since the satellite lines can be seen strongly in emission, the position-velocity diagrams for the (1,1) line are confusing with strong blending of the main line and the the inner satellites. Near the UCHii region, it is impossible to tell where emission and absorption from the different hyperfine components is blending. However, even in the (1,1) position-velocity diagrams, away from the UCHii region, one can see the gas returning to its ambient velocity, around 95.5 km s<sup>-1</sup>, shifting to about 97 km s<sup>-1</sup>in the far southeast of the cloud.
In the (2,2) and (3,3) position-velocity diagrams, the kinematics are less confused, showing the same outflow component detected at high angular resolution. Off the UCHii region, most of the emission is near the ambient velocity, 95.5 km s<sup>-1</sup>. At the position of the UCHii region we see several components in the main line, separated in velocity. There is blue-shifted absorption in both (2,2) and (3,3) detectable from 88 to 93 km s<sup>-1</sup>. There is also high velocity red-shifted emission out to 105 km s<sup>-1</sup>. The absorption is the same outflow component detected in the high resolution data, and the emission is the back side of the outflow, which could not be detected at high resolution. At the position of the UCHii region, the apparent optical depths at 90 km s<sup>-1</sup> in (2,2) and (3,3) are 0.25 and 0.05 respectively, giving a rotational temperature of 20 K where we have assumed the filling factors are equal to one for both lines. This is consistent with the upper limit of 30 K for the temperature of the outflow component, derived from the high resolution data. The fact that the back side of the outflow is only detected at the position of the UCHii region implies that it is no larger than the synthesized beam. The outflow cannot be too much smaller than the beam, however, since the emission from the back side of the outflow is strong, reaching a brightness temperature of 16 K in the (3,3) line, since emission much smaller than the beam will be strongly beam diluted.
The low angular resolution data also give a hint of the existence of the infall component, but only in the satellite lines, and mainly in (3,3), emphasizing the high optical depth and temperature of the infall component. The satellite is only barely detected in absorption in (2,2), but in (3,3) it is strongly detected. Since the intrinsic line-strength of the inner satellite in the (3,3) line is only about 3% of the main line, detecting the inner satellite strongly means the optical depth of the gas must be very high. The satellite absorption is at 75 to 78 km s<sup>-1</sup>. Since the (3,3) satellite line is offset by 21.5 km s<sup>-1</sup> from the main line, the corresponding high-optical-depth absorbing material should appear in the main line from 96.5 to 99.5 km s<sup>-1</sup>, exactly the velocities at which we detect the infall component at higher resolution (see Figure 4 especially). But if the optical depth toward the UCHii region is so high from 95 to 100 km s<sup>-1</sup>, why does the position-velocity diagram show little absorption and even some emission in that velocity range at the position of the continuum source? This apparent contradiction does not mean that the satellite line is somehow wrong in predicting the presence of optically thick main-line absorption. Instead, in this low spatial resolution data where the synthesized beam is larger than the continuum source, line emission in the range 95 to 100 km s<sup>-1</sup> from the area near the UCHii region can fill in the absorption which would otherwise be dominant.
In the low resolution data there are two clumps detected in all three lines, one northeast, the other southwest of the UCHii region. The satellite lines are detectable in these two clumps, only in the (1,1) line. Following the method of Ho & Townes (1983) for determining a rotational temperature when optical depth is calculable for only one rotational state, we estimate the temperature of both clumps to be 30 K. Based on the NH<sub>3</sub>(1,1) emission, the column densities of NH<sub>3</sub> are $`1.2\times 10^{15}`$cm<sup>-2</sup> and $`1.0\times 10^{15}`$cm<sup>-2</sup> in the northeast and southwest clumps respectively. Assuming an NH<sub>3</sub> abundance relative to $`\mathrm{H}_2`$ of $`10^7`$ (van Dishoeck & Blake, 1998), the northeast clump contains 12 M of gas, and the southeast clump 10 M.
## 4 Discussion
### 4.1 The Model
There are six key observational results that any model of this source must include. First, there must be two components, an infall component and an outflow component. Second, the outflow component must be seen in absorption over the entire face of the UCHii region, while the infall component must be seen over only half, cutting off in a sharp northwest-southeast line. Third, while the velocities of the infall component show a large projection effect over the $`1^{\prime \prime }`$ continuum source, the velocities of the outflow component do not project out nearly as strongly, that is the velocities do not vary strongly with position, and do not return to the ambient velocity at the edge of the continuum source. Fourth, the infall component must have high optical depth and a warm temperature (280 K), while the outflow component must have low optical depth and a lower temperature (20 to 30 K). Fifth, the outflow component must be smaller than the D array beam since the emission from the back side of the outflow is seen in the low resolution data only at the position of the UCHii region. At the same time, the outflow component cannot be much smaller than the D array beam since the emission is not too strongly beam diluted. Sixth, the model should explain the elongated shape of the continuum source, parallel to the sharp edge of the infall component. The model that we describe here fits all of these observed results.
Figure 9 shows our preferred model for G28. The infall component and the central continuum source are an inclined, infalling, possibly slowly rotating toroid, whose central region has been ionized by the central star out to a radius of 3400 AU (0.017 pc). The toroid is similar to those proposed in Beltrán et al. (2005). The ionized gas may resemble the photo-ionized disks modeled in Hollenbach et al. (1994) except that in this case there is no rotationally supported disk, just the infalling toroid. The outflow component is a molecular shell swept up by a larger, more tenuous expanding bubble of ionized gas. The radius of the shell is roughly 8300 AU (0.04 pc).
Our proposed model fits all the key observational results. The infalling and slowly rotating torus can be seen against part of the continuum source, but not all of it. The rear side of the torus obviously cannot be seen in absorption. The torus is undetectable in emission in the D-array data because of its small size compared to the low resolution beam. The outflow component provides the 88-92 km s<sup>-1</sup> absorption, and the 100+ km s<sup>-1</sup> emission seen at low resolution. The outflow is larger, more fully filling the beam at low resolution, and thus is detectable in emission. Because the expanding shell is large compared to the continuum source, its line-of-sight velocity does not return to the ambient velocity at the edge of the continuum source. The projection which is at work on the velocities of the infall component therefore does little to change the apparent velocity of the expanding shell. Because the infall component is closer to the central source it should be warmer than the outflow, which is mostly likely an isothermally shocked molecular shell. The infall component will naturally be more optically thick, since infalling gas is naturally centrally condensed. If the continuum emission is from photo-ionization of the densest central region of the toroid, that continuum emission region ought to be intrinsically flattened. Thus the inclination needed to get the sharp edge in the first moment map of the infall component would also naturally provide the elongation of the continuum source. In this way, the proposed model fits all of our key observational results.
The infall component is consistent with the existence of a central source having a mass of 79 M, while the infalling gas itself may include an amount of gas similar in mass. The local velocity is 95.5 km s<sup>-1</sup>, and the most red-shifted part of the infalling gas is at 100 km s<sup>-1</sup>. We will assume that this fastest infalling gas is at the radius of the continuum source, which has a deconvolved radius of $`0.^{\prime \prime }6`$, or 3400 AU (0.017 pc). If the continuum source is intrinsically circular, then the elongation implies an inclination of 45. Assuming the gas is in free-fall toward the central source, and all the velocity is in the plane of the disk-like structure, our estimate of the central mass, given by
$$M=\frac{R(v_{in}v_0)^2}{2Gsin^2i}$$
(1)
is 79 M. This is much larger than, but consistent with the lower limits of 41 M, based on the ionizing flux, and 31 M, based on the far-IR luminosity. (If the velocity is not purely in the plane of the disk-like structure, then the maximum velocity seen might be entirely along the line of sight, in which case $`sini=1`$ and the derived central mass would be 40 M.) The luminosity also sets an upper limit of 31 M on the largest star in the region, so this is likely to be a multiple system. This is not surprising given the high multiplicity of early type stars (Preibisch et al., 1999). The mass seen in absorption in the infall component itself, 9 M, could be only a small fraction of the mass of the toroid depending on the geometry. At the very least, the total mass is more than 18 M, since in absorption we only detect the near side. This is reminiscent of disk-like structures seen around very young early B and late O type stars which have masses which are an appreciable fraction of the stellar mass (Beuther et al., 2004; Zhang et al., 2002, 1998a).
The outflowing molecular gas could be part of a jet-driven bipolar outflow as seen around protostars, or it could be the product of pressure driven expansion of the Hii region, or even a spherical wind. Following Keto (2002a), we calculate the Bondi radius of the ionized gas. Given the mass as derived from the infall, the radius at which gravity and thermal pressure of the ionized gas balance is 240 AU, where we have assumed an ionized gas temperature of 10000 K. The radius of the detected continuum emission is 3400 AU (0.017 pc), so one could reasonably expect that the ionized gas is in a phase of pressure driven expansion. This also justifies our earlier assumption of constant density in the ionized gas. However, the velocity of the outflow component is smaller than what one would expect for pressure driven expansion. The absorption and emission seen in the low resolution (2,2) data are separated by a maximum of 17 km s<sup>-1</sup>(88 km s<sup>-1</sup> to 105 km s<sup>-1</sup>), so the expansion velocity is 8.5 km s<sup>-1</sup>. As a sound speed, this corresponds to a temperature of 3200 Kelvin for the ionized gas. Since the ionized gas is certainly hotter than that, the expansion may be impeded by the ram pressure of an infalling envelope. Still, the gravity of the central star should be weak compared to the force of the pressure imbalance between the ionized gas and the molecular gas. The ionized gas should be undergoing dynamic pressure driven expansion. Given a diameter of $`3^{\prime \prime }`$, equal to the size of the D array synthesized beam, the expansion speed implies a dynamic age of 5200 years. But if the UCHii region had an earlier gravitationally confined stage as described in Keto (2002a), the age of the source could be much larger. Another possible explanation for the expansion velocity being smaller than the sound speed of the ionized gas is that the outflow is not undergoing purely spherical expansion, but is rather a more directed bipolar outflow. In that case, the outflowing molecular gas might not be moving just along the line of sight, but also with a velocity in the plane of the sky.
The two clumps, one to the northeast, the other to the southwest, seen in Figure 7, do not fit clearly into our model. The two clumps are separated from the UCHii region by roughly 0.2 pc (41,000 AU). Thus, on the scale of Figure 9, they are off the page. The question of how to connect the gas in those clumps to the gas seen in the high resolution data remains, however. First we should note that if there were no high resolution data available, we would probably have interpreted the clumps as evidence of a rotating disk. The NH<sub>3</sub>(3,3) emission in Figure 7 is elongated along the northeast-southwest cut, and shows a velocity shift from one side of the UCHii region to the other side. This is similar to, although not quite as compelling as, the evidence for a disk in IRAS 20126+4104 (Zhang et al., 1998b). The point is moot, however, since the high resolution data firmly establish the axis of symmetry as northeast-southwest, and the plane of the toroid as northwest-southeast. Since the toroid is seen in the high resolution data in absorption against the southwest side of the UCHii region, we would expect a bipolar outflow to be blue-shifted in the northeast, and red-shifted in the southwest. This is indeed the sense of the roughly 1 km s<sup>-1</sup> velocity shift observed in the two clumps, as seen in Figure 8. But the 22 M observed in the two clumps is a great deal of mass for the outflow to accelerate. If the low velocity of the clumps is just a projection effect, and we believe that the entire 22 M is expanding with the same three dimensional speed as the outflow component seen at high resolution, i.e. 8.5 km s<sup>-1</sup>, the total momentum of the outflow would be 190 M km s<sup>-1</sup>. Alternatively, we could associate the two clumps with the infall component instead. In that case, the blue-shifted clump in the northeast would be on the far side, and the red-shifted clump in the southwest would be infalling on the near side. But the elongation is perpendicular to the elongation we expect from the geometry of the torus seen at high resolution. There are several ways to interpret the two clumps, none without problems. It should be remembered, however, that these clumps are beyond the physical size of the model, and need not detract from the fact that the model explains all the other observed results.
We have considered and ruled out another possible model, that there are two physically distinct continuum sources, each with its own associated molecular gas, which, in projection, overlap on the sky. This model is suggested by several facts. First, the continuum source is slightly resolved into two peaks. Second, the two components of the molecular material are quite distinct. Third, it quite naturally accounts for the sharp edge in the 97 km s<sup>-1</sup> as an edge in the continuum, not in the absorbing material. In this two component model, the rear continuum source is absorbed by the 97 km s<sup>-1</sup> component, and a second, nearer continuum source is absorbed only by the 90 km s<sup>-1</sup> component. This way the 90 km s<sup>-1</sup> gas is the nearest to the observer and can be seen in absorption against all the continuum. However, we rule out this model for the following reason. The 90 km s<sup>-1</sup> component is seen in the low resolution data to be paired with the 100+ km s<sup>-1</sup> emission. The two must be physically connected. This would imply that the 100+ km s<sup>-1</sup> emission is associated with the nearer continuum source. But if that were so, it would appear in absorption against the rear continuum source, which is not observed. For this reason, we reject this two component model.
### 4.2 Comparisons to Other Objects
We contrast these observations to similar observations of G10.6 (Sollins et al., 2005). In that case, where accretion is proceeding toward a central group of massive young stars, there is some flattening associated with rotation in the molecular gas. But there is no conspicuous edge like the one we see in G28. In G10.6, the lack of a sharp edge is cited as evidence that no rotationally supported, geometrically thin, optically thick disk exists around the stellar group. While G28 does show such a sharp edge, it also does not appear to have a rotationally supported disk. The circular rotation velocity implied by a 79 M central mass at a radius of 3400 AU with an inclination of 45, would be 3.2 km s<sup>-1</sup>, implying a 6.4 km s<sup>-1</sup> shift in velocity from one side of the UCHii region to the other. There is in fact less than 1 km s<sup>-1</sup>difference in velocity across the UCHii region. So the sharp edge cannot come from the fact that the gas has settled into a geometrically thin disk. The flattening seen could come from an initially flattened geometry of the cloud, as modeled in Hartmann et al. (1996), collapse along magnetic field lines, or sculpting of the infall component via interaction with the outflow.
There are some similarities and also a key difference between G28 and other young massive stars observed to have molecular gas in disk-like structures. IRAS 20126+4104, IRAS 18089-1732, and AFGL5142 all include early B type or late O type stars surrounded by a rotating disk at roughly 5000 AU scale (Zhang et al., 1998b; Beuther et al., 2002; Zhang et al., 2002). These rotating disks are sub-keplerian, and do not appear to be rotationally supported, similar to our model of G28. In G192.16-3.82, Shepherd & Kurtz (1999) find rotation in water maser spots at radii as small as 1000 AU, also around an early B type star. In that case, at smaller radii, the velocity gradient is consistent with Keplerian rotation. In all these cases, the amount of mass associated with the disk-like structure is smaller than, but comparable to the estimated mass of the central object. While G28 has a much more massive central object than any of these other regions, it seems to have much in common. The toroid through which infall is proceeding is not rotationally supported at a radius of 3000 AU. We detect 9 M of molecular gas in absorption in the infall component, but that is just in the part of the toroid that lies directly in front of the continuum source. Certainly there must be as much mass in the part of the toroid behind the continuum source as there is in front, and depending on the geometry of the toroid, the mass might be a factor of a few larger than the measured 9 M. Thus the mass of the toroid may be the same order of magnitude as the central object, similar to the other objects. Unlike any of these other objects, however, G28 has a strong central bremsstrahlung continuum source. This difference may be attributable to the larger mass, and therefore larger ionizing flux of the central star or stars. G28 could also be a more evolved version of these other sources, with more time allowing for the evolution of an UCHii region. G28 is in some ways similar to IRAS 20216+4104 and other similar young massive stars with disks, but the central source in G28 is more massive, and the system may be somewhat more evolved.
## 5 Summary
We propose a qualitative model for the source G28.20-0.05. The model has two components, an equatorial toroid of infalling molecular gas in which a central dense region has been photo-ionized to form the UCHii region, and an expanding molecular shell of larger radius, roughly 8300 AU (0.04 pc). Observationally, the molecular gas is easily divisible into two kinematically, and physically distinct components. One is seen in blue-shifted absorption and red-shifted emission, with low temperature (20-30 K) and relatively low optical depth which we associate with the outflow. The other is seen only in red-shifted absorption, with much higher temperature (280 K) and much higher optical depth (up to more that 45 in the (3,3) line) which we associate with the infall. We calculate that the central mass responsible for the infall is 79 M, and the that the mass of the torus is at least 18 M. The lack of strong rotation shows that, although the infalling gas has a flattened geometry, no rotationally supported accretion disk exists at radii as small as the continuum source, about 3000 AU.
## 6 Appendix: Optical Depths
We calculate two different types of optical depth in the inversion lines of NH<sub>3</sub> (for an example see Figure 10). The first we call “apparent optical depth” and it is found in any given channel by solving
$$T_{line}=T_{cont}(1e^{\tau _{app}})$$
(2)
for $`\tau _{app}`$, where $`T_{cont}`$ is the temperature of the continuum emission and $`T_{line}`$ is the depth of the line absorption. This equation assumes no line emission, only absorption. For a given continuum level, the lowest calculable apparent optical depth is set by the detection of absorption as different from the continuum level and the highest calculable apparent optical depth is set by detection of flux in the channel in question. Apparent optical depth includes a filling factor, so that a very clumpy gas with locally high actual optical depth may have a low apparent optical depth if the optically thick gas does not completely fill the synthesized beam. Apparent optical depth can be calculated wherever absorption is detected against the continuum and a lower limit can be derived in any channel in which the continuum is absorbed down below the detection limit. For a continuum level of 100 mJy beam<sup>-1</sup> and noise level in one channel of a spectrum of 4 mJy beam<sup>-1</sup>, if we use $`2\sigma `$ as our detection limit, the highest and lowest detectable apparent optical depths are 2.5 and 0.08 respectively. In this case, if the apparent optical depth were higher than 2.5 it would be indistinguishable from infinity, and if the apparent optical depth were less than 0.08, the absorption would be undetectable.
The second type of optical depth we calculate we will call the “hyperfine optical depth”. This is the optical depth of the main hyperfine component as calculated from the ratio of the fluxes of a satellite hyperfine component to the main hyperfine component. Assuming LTE, and that exactly the same material emits both the main hyperfine component and the inner satellite, the only difference in the optical depths of the two lines should be due to their differing intrinsic line-strengths. Since their optical depths differ by a known multiplicative constant, $`\tau _{sat}=x\tau _{main}`$, the ratio of the fluxes of the two lines can be used to solve for the optical depth of either line using the following equation.
$$\frac{T_{sat}}{T_{main}}=\frac{1e^{\tau _{sat}}}{1e^{\tau _{main}}}=\frac{1e^{x\tau _{main}}}{1e^{\tau _{main}}}$$
(3)
By convention, we solve for $`\tau _{main}`$. The details of this calculation and the hyperfine structure of the different NH<sub>3</sub> inversion lines are laid out in Ho & Townes (1983). Because only the ratio of fluxes or temperatures is used, the hyperfine optical depth can be calculated either when the line is seen in emission, or when the line is seen in absorption. Also, the use of the ratio means that the filling factor cancels leaving the hyperfine optical depth a truer measure of the optical depth. The detectable range of hyperfine optical depths varies between lines, since the hyperfine structure is different in different rotational states. For the (2,2) line as seen in absorption against a continuum brightness of 100 mJy beam<sup>-1</sup>, a noise level of 4 mJy beam<sup>-1</sup> in one spectral channel, and using a $`2\sigma `$ detection limit, the best-case-scenario range of detectable hyperfine optical depths is 0.5 to 38. For (3,3) the range would be 2.5 to 80. It should be noted that in order to calculate the hyperfine optical depth both a satellite and the main hyperfine components must be detected. Since the satellite line has an intrinsically lower optical depth, it is sensitive to much larger column densities of gas than the main component. In some cases the satellite may be detectable in absorption, but the main line may not, for instance if the continuum source is smaller than the synthesized beam, and widespread, low-optical-depth emission in the main line fills in the absorption, as in our lower spatial resolution NH<sub>3</sub>(2,2) and NH<sub>3</sub>(3,3) data. In this case, the hyperfine optical depth is not formally calculable. But one can get an estimate by noting that the apparent optical depth in the main line should just be the apparent optical depth of the satellite times the constant $`x`$. In almost every real circumstance, the apparent and hyperfine optical depths probe disjoint parts of parameter space.
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# Anisotropic flow of strange particles at RHIC
## 1 Introduction
The transverse collective flow of particles is usually decomposed onto isotropic radial flow and anisotropic components, such as directed flow, elliptic flow, etc. Directed and elliptic flows are defined as the first and the second harmonic coefficients, $`v_1`$ and $`v_2`$, of an azimuthal Fourier expansion of the particle invariant distribution
$$E\frac{d^3N}{d^3p}=\frac{1}{\pi }\frac{d^2N}{dp_t^2dy}\left[1+2v_1\mathrm{cos}(\varphi )+2v_2\mathrm{cos}(2\varphi )+\mathrm{}\right],$$
(1)
where $`\varphi `$ is the azimuthal angle between the transverse momentum of the particle and the reaction plane, and $`p_t`$ and $`y`$ is the transverse momentum and the rapidity, respectively. The directed and elliptic flows can be presented as
$$v_1\mathrm{cos}\varphi =\frac{p_x}{p_t},v_2\mathrm{cos}2\varphi =\frac{p_x^2p_y^2}{p_t^2}.$$
(2)
with $`x`$ being the impact parameter axis. The transverse momentum of a particle is simply $`p_t=\sqrt{p_x^2+p_y^2}`$.
The overlapping area of two nuclei colliding with non-zero impact parameter $`b`$ has a characteristic almond shape in the transverse plane. The fireball tries to restore spherical shape, provided the thermalization sets in rapidly and the hydrodynamic description is appropriate . When it becomes spherical, apparently, the elliptic flow stops to develop. Therefore, $`v_2`$ can carry important information about the earlier phase of ultrarelativistic heavy-ion collisions, equation of state (EOS) of hot and dense partonic matter, and is expected to be a useful tool to probe the formation and hadronization of the quark-gluon plasma (QGP).
## 2 Directed and Elliptic Flow
Figure 2 shows the rapidity dependence of the directed flow of $`\varphi ,N`$, and $`K`$ in minimum bias Au+Au collisions at $`\sqrt{s}=130`$ AGeV. The slopes of the all distributions are negative at $`|y|2`$, i.e. the antiflow component of the $`v_1`$ dominates over its normal counterpart (see ). Similar antiflow slopes of $`v_1(y)`$ are developed by $`\pi `$ and $`\mathrm{\Lambda }`$ ; its origin is traced to nuclear shadowing. At midrapidity $`|y|0.5`$ the directed flow of all hadrons is quite weak. Figure 2 depicts the simulation results for the $`v_1(\eta )`$ of charged hadrons compared to the experimental data from the PHOBOS Collaboration for 6% to 55% central Au+Au collisions at $`\sqrt{s}=200`$ AGeV
One can see that the model reproduces the $`v_1`$ data quite well both qualitatively and quantitatively, although the maxima of the directed flow around $`|\eta |2`$ are shifted to lower pseudorapidities compared to the experimental data. Similar antiflow alignment can be obtained also within the multi module model (MMM) , which is based on fluid dynamics coupled to formation of colour ropes. Microscopic models based on FRITIOF routine, e.g. UrQMD and AMPT, show a very flat and essentially zero directed flow in a broad range $`|\eta |2.5`$. Although the data seem to indicate antiflow behaviour for the directed flow of charged particles with the possible flatness at $`|\eta |1.5`$, the measured signal is quite weak, – the magnitude of the flow is less than 1% at $`|\eta |2`$. Therefore, relatively large systematic error bars do not permit us to disentangle between the different models. The other features which should be mentioned here are broadening of the antiflow region and increase of its strength as the reaction becomes more peripheral.
Microscopic models based on string phenomenology and transport theory are able to reproduce many features of the elliptic flow at ultrarelativistic energies . However, the quantitative agreement with the data is often not so good. Particularly, magnitude of the distributions $`v_2(\eta )`$ or $`v_2(p_t1.5`$ GeV/$`c)`$ appears to be too high. Does it mean that the effective EOS of hot and dense partonic-hadronic matter in microscopic models is too soft? Then, the microscopic calculations show the absence of sharp freeze-out of particles in relativistic heavy-ion collisions. What are the consequences of the continuous freeze-out for the $`v_2`$ of these particles? To study the development of the elliptic flow ca. 20$`10^3`$ gold-gold collisions with the impact parameter $`b=8`$ fm were generated at $`\sqrt{s}=130`$ AGeV. According to previous studies the elliptic flow of charged particles is close to its maximum at this impact parameter, and the multiplicity of secondaries is still quite high. The time evolutions of the $`v_2`$ of kaons and lambdas as functions of rapidity are displayed in Fig. 4(a). Here the snapshots of the $`v_2`$ profile are taken at certain time $`t=t_i`$, when all hadronic interactions are switched off and particles are propagated freely.
To avoid ambiguities, resonances were allowed to decay according to their branching ratios. Surprisingly, at $`t=2`$ fm/$`c`$ elliptic flow of kaons is weak. The flow continuously increases and reaches its maximum value $`v_2^K(y=0)6\%`$ already at $`t=8`$ fm/$`c`$. From this time the elliptic flow does not increase anymore. Instead, it becomes broader and develops a two-hump structure with a relatively weak dip at midrapidity. The flow seems to continue development till the late stages of the system evolution. However, the contributions of survived particles to the resulting elliptic flow presented in Fig. 4(b) reveal the peculiar feature: the $`v_2`$ of kaons, which are frozen already at $`t=2`$ fm/$`c`$, is the strongest among the fractions of the flow carried by kaons decoupled from the fireball later on. The later the kaons are frozen, the weaker their flow. One can conclude that the strong elliptic anisotropy of kaons, which left the system early, is caused by the absorption of kaons in the squeeze-out direction.
For lambdas the evolution picture of the $`v_2(y)`$, shown in Fig. 4(a), is similar to that for kaons. The flow is quite weak at $`t=2`$ fm/$`c`$, then it increases and gets a full strength at midrapidity between 8 fm/$`c`$ and 10 fm/$`c`$, i.e. later than the elliptic flow of kaons. Similarly to $`v_2^K(y)`$, it develops a two-hump structure, but the humps tend to dissolve at late stages of system evolution. In contrast to this behavior, the freeze-out decomposition picture of $`v_2^\mathrm{\Lambda }(y)`$, presented in Fig. 4(b), does not show monotonic tendency within first 8 fm/$`c`$ of the reaction: The flow of $`\mathrm{\Lambda }`$ frozen at 2 fm/$`c`$ is identical to that of $`\mathrm{\Lambda }`$ frozen at 8 fm/$`c`$, whereas lambdas decoupled from the system between 2 fm/$`c`$ and 8 fm/$`c`$ almost do not contribute to the resulting elliptic flow. Lambdas, which are decoupled after 8 fm/$`c`$, have significant anisotropy in the momentum space, and the later the lambdas are frozen, the stronger their elliptic flow. This picture is similar to that obtained for the development of pionic and nucleonic elliptic flows .
## 3 Conclusions
In summary, the features of the formation and development of anisotropic flow in gold-gold collisions at RHIC in the microscopic quark-gluon string model can be stated as follows. (1) The directed flow of all hadrons exhibits antiflow alignment within the pseudorapidity range $`\eta 2`$. The signal increases as the reaction becomes more peripheral. At midrapidity $`|\eta |1`$, however, the generated flow is quite weak and consistent with zero-flow behaviour reported by the STAR and PHOBOS collaborations. (2) There is no one-to-one correspondence between the apparent elliptic flow and the contribution to the final flow coming from the “survived” fraction of particles. For instance, apparent elliptic flow of kaons at $`t=2`$ fm/$`c`$ is weak, but kaons which are already decoupled from the system at this moment have the strongest elliptic anisotropy caused by their absorption in the squeeze-out direction. Elliptic flow of hadrons is formed not only during the first few fm/$`c`$, but also during the whole evolution of the system because of continuous freeze-out of particles. (3) The time evolutions of the mesonic flow and baryonic flow are quite different. The general trend in particle flow formation in microscopic models at ultrarelativistic energies is that the earlier mesons are frozen, the weaker their elliptic flow. In contrast, baryons frozen at the end of the system evolution have stronger $`v_2`$. Therefore, development of particle collective flow should not be studied independently of the freeze-out picture.
## References
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# The minimum dilatation of pseudo-Anosov 5-braids
## 1. Introduction
Let $`f:D^2D^2`$ be an orientation preserving disk homeomorphism which is the identity map on the boundary $`D^2`$, and $`\{p_i\}int(D^2)`$ be a periodic orbit of $`f`$ (or more generally a finite set invariant under $`f`$). The points $`p_i`$ move under an isotopy from the identity map on $`D^2`$ to $`f`$. Their trajectory forms a geometric braid $`\beta `$, a collection of strands in $`D^2\times [0,1]`$ connecting $`p_i\times 1`$ to $`f(p_i)\times 0`$ (see Figure 1). The isotopy class of $`\beta `$ determines the homotopy class of $`f`$ relative to $`\{p_i\}D^2`$ and vice versa. An $`n`$-braid refers to the isotopy class of a geometric braid with $`n`$ strands. The set of $`n`$-braids forms the braid group $`B_n`$.
From now on we consider $`f`$ as a homeomorphism on a punctured sphere $`f:int(D^2)\{p_i\}int(D^2)\{p_i\}`$. In particular by forgetting the boundary $`D^2`$, we ignore Dehn twists along $`D^2`$ which do not affect the dynamics of the braid $`\beta `$. In other words we consider an $`n`$-braid $`\beta `$ as a mapping class on a $`(n+1)`$-times punctured sphere with the (so called) boundary puncture fixed.
### Topological entropy
The topological entropy $`h_T(\beta )`$ of the braid $`\beta `$ is defined to be the infimum topological entropy of the disk homeomorphisms representing $`\beta `$.
$$h_T(\beta )=\underset{gf}{inf}h_T(g)$$
The topological entropy of a braid is a conjugacy invariant measuring the dynamical complexity of the braid. It is equal to the logarithm of the growth rate of the free group automorphism induced on $`\pi _1(D^2\{p_i\})`$. When $`\beta `$ is represented by a pseudo-Anosov homeomorphism $`f`$ with dilatation factor $`\lambda _f=\lambda (f)`$, we have $`h_T(\beta )=\mathrm{log}\lambda _f`$. In this case the dilatation $`\lambda (\beta )`$ of the braid is also given by $`\lambda (\beta )=\lambda _f`$.
If $`f`$ is homotopic to a periodic homeomorphism, the braid $`\beta `$ is called periodic. If there is a collection of disjoint sub-surfaces of $`int(D^2)\{p_i\}`$ with negative Euler characteristics which is homotopically invariant under $`f`$, the braid $`\beta `$ is called reducible. As we consider the dynamics of a periodic braid to be trivial, studying the dynamics of braids reduces to the maps on non-periodic irreducible components.
### Pseudo-Anosov homeomorphism
By the Nielsen-Thurston classification of surface homeomorphisms , a non-periodic irreducible braid is represented by a pseudo-Anosov homeomorphism. A pseudo-Anosov homeomorphism has several nice extremal properties. It realizes the minimum topological entropy and the minimum quasi-conformality constant in its homotopy class. It also has the minimum number of periodic orbits for each period .
A surface homeomorphism $`f:FF`$ is called a *pseudo-Anosov homeomorphism* relative to a puncture set $`\{p_i\}F`$ when the following conditions hold. First we need a singular flat metric on $`F`$ with a finite singularity set $`\{q_j\}`$ such that $`\{p_i\}\{q_j\}`$. Each singularity $`q_j`$ has its cone-angle in $`\{k\pi k_{>0}\}`$. If a singularity has cone-angle $`\pi `$, it must be one of the puncture points $`p_i`$. Now the homeomorphism $`f`$ is required to preserve $`\{q_j\}`$ and to be locally affine (hyperbolic) on $`F\{q_j\}`$ with a constant dilatation factor $`\lambda _f>1`$. In particular at a fixed point in $`F\{q_j\}`$, the map $`f`$ is locally written as $`\left[\begin{array}{cc}\lambda _f& 0\\ 0& \lambda _f^1\end{array}\right]`$.
Thus roughly speaking, if a surface homeomorphism $`f`$ represents a non-periodic irreducible mapping class, then one can simplify $`f`$ by pulling tight it everywhere until it becomes linear almost everywhere in an appropriate sense.
The horizontal directions to which $`f`$ expands by the factor $`\lambda _f`$ integrate to form one invariant measured foliation $`^s`$. The vertical directions perpendicular to $`^s`$ form the other invariant measured foliation $`^u`$. From a singularity $`q_j`$ with cone-angle $`k\pi `$, $`k`$-many singular leaves of $`^s`$ emanate. In this case $`q_j`$ is called a $`k`$-prong singularity.
Note that in the above definition of pseudo-Anosov homeomorphism we can remove or add punctures keeping the same map $`f:FF`$. When $`\{f^j(x)\}`$ is a periodic orbit of non-punctured points, *puncturing* at $`\{f^j(x)\}`$ refers to adding them to the puncture set $`\{p_i\}`$. Conversely when $`\{f^j(p_1)\}`$ is a periodic orbit of $`k`$-prong punctured singularities for $`k>1`$, *capping-off* them refers to removing them from the puncture set. For pseudo-Anosov braids, puncturing or capping-off corresponds to adding or removing some strands.
Let $`\stackrel{~}{f}:\stackrel{~}{F}\stackrel{~}{F}`$ be a lift of $`f`$ on a finite-fold cover $`\stackrel{~}{F}`$ of $`F`$ branched at some finite set of points invariant under $`f`$. Then by pulling back the flat metric on $`F`$ to $`\stackrel{~}{F}`$, the lift $`\stackrel{~}{f}`$ is also a pseudo-Anosov homeomorphism with the same dilatation factor $`\lambda _{\stackrel{~}{f}}=\lambda _f`$.
### Train track representative
Using a Markov partition (or its associated train track representative), the flat metric and the pseudo-Anosov homeomorphism can be described quite concretely (see \[10, Exposé 9\] for the definition and see Figure 2 for an example). Let $`\{R_i\}`$ be a Markov partition for a pseudo-Anosov homeomorphism $`f`$. The transition matrix $`M_f=(m_{ij})`$ is defined by that $`f(R_i)`$ crosses over $`R_j`$ $`m_{ij}`$-many times.
The transition matrix $`M_f`$ is Perron-Frobenius: for some $`k>0`$ each entry of $`M_f^k`$ is strictly positive. In particular the largest eigenvalue of $`M_f`$ is real and has an eigenvector with strictly positive coordinates \[21, Theorem 1.1 on p.1\].
The widths $`v_i`$ and the heights $`w_i`$ of $`R_i`$ satisfy the following equations.
$$\lambda _fv_i=\underset{j}{}m_{ij}v_jw_j=\underset{i}{}m_{ij}w_i/\lambda _f$$
In particular the dilatation factor $`\lambda _f`$ appear as the eigenvalue of $`M_f`$ of which the eigenvector has strictly positive coordinates.
We use train track representatives as a notational simplification for Markov partitions. As in Figure 2, each expanding edge of the invariant train track corresponds to a rectangle in the Markov partition. Once given the transition matrix of the graph map, it is easy to recover the heights and widths of the rectangles.
### Main question
Let us consider the set $`\mathrm{\Lambda }_{g,n}`$ of the dilatation factors for pseudo-Anosov homeomorphisms on an $`n`$-times punctured genus-$`g`$ surface $`F_{g,n}`$.
$$\mathrm{\Lambda }_{g,n}=\{\lambda _ff:F_{g,n}F_{g,n}\text{ pseudo-Anosov homeomorphisms}\}$$
As we can bound the number of rectangles in Markov partitions using the Euler characteristic of the punctured surface, $`\mathrm{\Lambda }_{g,n}`$ consists of eigenvalues of Perron-Frobenius matrices with bounded dimension. In particular the set $`\mathrm{\Lambda }_{g,n}`$ is discrete and has a minimum. Our current work is motivated by the following question.
###### Question 1.1.
What is $`\mathrm{min}\mathrm{\Lambda }_{g,n}`$?
The question asks to find the simplest pseudo-Anosov homeomorphism on the surface.
A pseudo-Anosov homeomorphism $`f`$ induces an isometry on the Teichmüller space equipped with the Teichmüller metric. The pair of invariant measured foliations $`(^s,^u)`$ determines a geodesic axis in the Teichmüller space on which $`f`$ acts as a translation by $`\mathrm{log}\lambda _f`$. The axis projects down to a closed geodesic in the moduli space, which is the quotient of the Teichmüller space by the action of the mapping class group. Conversely any closed geodesic in the moduli space represents the conjugacy class of some pseudo-Anosov mapping class. Therefore Question 1.1 can be rephrased as asking to find the shortest closed geodesic in the moduli space.
The hyperbolic volume of the mapping torus is another natural complexity measure for a pseudo-Anosov homeomorphism. We notice that a pseudo-Anosov homeomorphism with small dilatation tends to have the mapping torus with small hyperbolic volume and vice versa.
The question for the minimum volume of the hyperbolic mapping tori on a given surface seems to be much more difficult than Question 1.1. In the minimum volume for orientable cusped hyperbolic 3-manifolds is computed. An extensive use of computer programs is involved in its proof. In this paper we also use a computer program for the proof of the main theorem, but the algorithm and the actual code is much simpler than those of .
### Related results
The question asking for the minimum dilatation of pseudo-Anosov homeomorphisms on a given surface still remains largely unanswered since after the Nielsen-Thurston classification of surface homeomorphisms. The existence of Markov partition \[10, Exposé 10\] for a pseudo-Anosov homeomorphism implies that a dilatation should appear as the largest eigenvalue of a Perron-Frobenius matrix of bounded dimension, hence in particular should be an algebraic integer. However, it is not clear how the restriction that the symbolic dynamical system dictated by a Perron-Frobenius matrix is from a homeomorphism on a given surface, actually affects the possible values of entropy (the logarithm of the dilatation).
There are several known results relevant to this question of the minimum dilatation of pseudo-Anosov homeomorphisms. Penner gives a lower bound $`2^{1/(12g12+4n)}`$ for the dilatations on $`F_{g,n}`$ a genus-$`g`$ surface with $`n`$ punctures. In pseudo-Anosov homeomorphisms on $`F_{g,0}`$ with small dilatations are constructed showing that the minimum dilatation on $`F_{g,0}`$ converges to $`1`$ as the genus $`g`$ increases. Fehrenbach and Los compute a lower bound $`(1+\sqrt{2})^{1/n}`$ for the dilatations of pseudo-Anosov disk homeomorphisms (braids) which permute the punctures in one cycle. In a lower bound $`2+\sqrt{5}`$ for the dilatations of pseudo-Anosov pure braids is given. A pseudo-Anosov disk homeomorphism is represented by a transitive Markov tree map preserving the end point set of the tree with the same topological entropy. Baldwin gives a lower bound $`\mathrm{log}3`$ for the topological entropy of transitive Markov tree maps fixing each end point.
The exact values of the minimum dilatations are known only for few simple cases. Zhirov shows that if a pseudo-Anosov homeomorphism on $`F_{2,0}`$ has an orientable invariant foliation, then its dilatation is not less than the largest zero $`\lambda _5`$ of $`x^4x^3x^2x+1`$, and gives an example of a pseudo-Anosov homeomorphism realizing the dilatation $`\lambda _5`$.
The pseudo-Anosov 3-braid $`\sigma _2\sigma _1^1`$ is shown to be the minimum in the forcing partial order among pseudo-Anosov 3-braids by Matsuoka and Handel , hence it attains the minimum dilatation. The pseudo-Anosov 4-braid $`\sigma _3\sigma _2\sigma _1^1`$ is claimed in to have the minimum dilatation, but the proof given in unfortunately contains an error.
### Outline
In this paper we prove the following theorem giving at the same time a corrected proof of the minimality of the dilatation of $`\sigma _3\sigma _2\sigma _1^1B_4`$.
###### Theorem 3.9.
The 5-braid $`\sigma _1\sigma _2\sigma _3\sigma _4\sigma _1\sigma _2`$ attains the minimum dilatation of pseudo-Anosov 5-braids.
The dilatation of a pseudo-Anosov braid is invariant under several operations such as conjugation, composing with a full twist, taking inverse, and taking reverse. It turns out that for braid indices 3 to 5, the pseudo-Anosov braids realizing the minimum dilatations are essentially unique, modulo the aforementioned operations. This could be just a coincidence. It would be a good surprise if some uniqueness property can be proved for the minimum-dilatation pseudo-Anosov braids.
The two main ingredients of the proof of Theorem 3.9 are the construction of folding automata for generating candidate pseudo-Anosov braids for the minimum dilatation, and the following lemma for bounding the word lengths of the candidate braids.
###### Lemma 3.1.
If $`M`$ is a Perron-Frobenius matrix of dimension $`n`$ with $`\lambda >1`$ its largest eigenvalue, then
$$\lambda ^n|M|n+1$$
where $`|M|`$ denotes the sum of entries of $`M`$.
This lemma improves on \[19, Theorem 6\] and replaces erroneous Lemma 3,4 of .
Given a pseudo-Anosov homeomorphism $`f:(F,\{p_i\})(F,\{p_i\})`$ on a surface $`F`$ with punctures $`p_i`$ with negative Euler characteristic $`\chi (F\{p_i\})<0`$, there exists a train track representative of $`f`$. There exists an invariant train track $`\tau F\{p_i\}`$ which carries $`f(\tau )`$. In particular there is a splitting sequence
$$\tau =\tau _0\tau _1\mathrm{}\tau _k=f(\tau )$$
from $`\tau `$ to $`f(\tau )`$, where $`\tau _j\tau _{j+1}`$ is an elementary splitting move.
By observing that there are only finitely many diffeomorphism types of the pair $`(F\{p_i\},\tau _j)`$, one can effectively construct a *splitting automaton*, which is a finite graph with train tracks as its vertices and with splitting moves as its arrows.
The existence of the train track representative, in particular of the splitting sequence, implies that every pseudo-Anosov homeomorphism appears, up to conjugacy, as a closed path in some splitting automata (see ). Papadopoulos and Penner \[19, Theorem 6\] also gave a lower bound for the dilatation in terms of word length in automata.
In this paper we use folding automata as in , which are finite graphs with embedded train tracks as vertices and with elementary folding maps as arrows. An elementary folding map is an inverse of a splitting move.
If we are given an upper bound for the word length in terms of the dilatation, then on a fixed folding automaton, the search for the minimum dilatation in the automaton reduces to checking for finitely many closed paths.
Lemma 3.1, which is an improvement of \[19, Theorem 6\], not only gives an upper bound of the word lengths of mapping classes with dilatation bounded by a fixed number, but also trims out many branches which appear in the course of search in a big tree, namely the set of paths with bounded length. In fact Lemma 3.1 implies that it suffices to consider only such paths whose any subpath has a transition matrix with bounded norm.
For the minimum dilatation of 5-braids, the previously mentioned restriction on paths by transition matrix norm and another restriction by Lemma 3.3 significantly reduce the number of candidate braids making the computation feasible.
We think that the same method for computing the minimum dilatation would still work for a few more simple cases like on a genus-2 closed surface, although it would involve more complicated computer aided search.
## 2. Folding automata
Given a pseudo-Anosov homeomorphism $`f:(F,\{p_i\})(F,\{p_i\})`$ on a closed surface $`F`$ with punctures $`\{p_i\}`$, there exists an invariant train track $`\tau F\{p_i\}`$ and $`f`$ is represented by a train track map $`f_\tau :\tau \tau `$ .
A train track $`\tau `$ is a smooth branched 1-manifold embedded in the surface $`F\{p_i\}`$ such that each component of the complement $`F\{p_i\}\tau `$ is either a once punctured $`k`$-gon for $`k1`$ or a non-punctured $`k`$-gon for $`k3`$. The train track $`\tau `$ is called *invariant* under $`f`$ if $`f(\tau )`$ smoothly collapses onto $`\tau `$ in $`F\{p_i\}`$ inducing a smooth map $`f_\tau :\tau \tau `$, which maps branch points to branch points. In this case one may repeatedly fold (or zip) $`f(\tau )`$ nearby cusps to obtain a train track isotopic to $`\tau `$ in $`F\{p_i\}`$ (see Figure 3 and \[15, Fig. 4, 5\]).
Let $`f_\tau :\tau \tau `$ be a train track representative of a pseudo-Anosov homeomorphism $`f`$. An edge $`e`$ of $`\tau `$ is called *infinitesimal* if it is eventually periodic under $`f_\tau `$, that is, $`f_\tau ^{N+k}(e)=f_\tau ^N(e)`$ for some $`N,k>0`$. An edge of $`\tau `$ is called *expanding* if it is not infinitesimal. An expanding edge $`e`$ actually has a positive length in the sense that $`lim_N\mathrm{}|f_\tau ^N(e)|/\lambda _f^N`$ is positive where $`||`$ denotes the word length of a path and $`\lambda _f=\lambda (f)`$ denotes the dilatation factor for $`f`$.
A graph map is called *Markov* if it maps vertices to vertices, and is locally injective at points that do not map into vertices. Given a Markov map $`g:\tau \tau ^{}`$, the transition matrix $`M_g=(m_{ij})`$ is defined by that the $`j`$-th edge $`(e_j^{})^{\pm 1}`$ of $`\tau ^{}`$ occurs $`m_{ij}`$-many times in the path $`g(e_i)`$, the image of the $`i`$-th edge of $`\tau `$. When $`\tau ^{}=\tau `$, the transition matrix is square and considering its spectral radius makes sense.
The spectral radius of $`M_{f_\tau }`$ equals the dilatations factor $`\lambda (f)`$ for the pseudo-Anosov homeomorphism $`f`$. Coordinates of the corresponding eigenvectors of $`M_{f_\tau }`$ and its transpose $`M_{f_\tau }^T`$ are tangential and transverse measures of edges of $`\tau `$, which are projectively invariant under $`f`$.
An elementary folding map $`\pi :\tau \tau ^{}`$ is a smooth Markov map between two train tracks $`\tau `$ and $`\tau ^{}`$ such that for only one edge $`e`$ of $`\tau `$, the image $`\pi (e)`$ has word length $`2`$, and the other edges map to paths of length $`1`$. In other words the transition matrix $`M_\pi `$ is of the form $`P+B`$ for some permutation matrix $`P`$ and for some elementary matrix $`B`$.
When the train tracks are embedded in a surface as in our case of concern, the pairs of edges which are folded should be adjacent in the surface: the two segments of $`\tau `$ which are identified by the elementary folding map are two sides of an open triangle in $`F\{p_i\}\tau `$ (see Figure 3).
###### Proposition 2.1.
A train track representative $`f_\tau :\tau \tau `$ of a surface homeomorphism $`f:(F,\{p_i\})(F,\{p_i\})`$ admits a folding decomposition as follows:
$$f_\tau =\rho \pi _k\mathrm{}\pi _1$$
where $`\pi _j:\tau _j\tau _{j+1}`$ are elementary folding maps, $`\tau _1=\tau _{k+1}=\tau `$, and $`\rho :\tau \tau `$ is an isomorphism induced by a periodic surface homeomorphism $`(F\{p_i\},\tau )(F\{p_i\},\tau )`$ preserving $`\tau `$.
###### Proof.
It follows from . See for more details. ∎
By observing that there are only finitely many possible diffeomorphism types for the pairs $`(F\{p_i\},\tau _j)`$ appearing in the folding decomposition, we can construct folding automata. A *folding automaton* is a connected directed graph with diffeomorphism types of train tracks as vertices, with elementary folding maps and isomorphisms as arrows. See Figure 7 for a simplified version of a folding automaton. The train tracks in Figure 7 admit no non-trivial isomorphisms, that is, if $`h:(D_5,\tau )(D_5,\tau )`$ is an orientation preserving diffeomorphism fixing $`\tau `$ in the automaton, then $`h`$ is isotopic to the identity map. So in Figure 7 there are no arrows corresponding to isomorphisms.
###### Corollary 2.2.
All train track representatives of pseudo-Anosov homeomorphisms are represented by closed oriented paths in folding automata.
To each closed path based at a train track $`\tau `$ in a folding automaton, associated is a train track representative $`f_\tau :\tau \tau `$ of some homeomorphism $`f:(F,\{p_i\})(F,\{p_i\})`$. The disk homeomorphism $`f`$ is pseudo-Anosov if and only if the transition matrix $`M_{f_\tau }`$ is Perron-Frobenius (also called primitive) modulo infinitesimal edges: for some $`N>0`$, the power $`M_{f_\tau }^N`$ is strictly positive in the block of expanding edges. To find out whether $`M`$ is Perron-Frobenius, it suffices by , \[21, Theorem 2.8 on p.52\] to check if $`M^{n^22n+2}`$ has all non-zero entries where $`n`$ is the dimension of the matrix $`M`$.
Now we discuss simplifying the train track maps so that we can restrict to simplified folding automata. If the pseudo-Anosov homeomorphism $`f`$ fixes a distinguished puncture $`p_0`$: $`f(p_0)=p_0`$ (for instance when $`f`$ is from a disk homeomorphism and $`p_0`$ is the boundary puncture), we can give a restriction to the train track map $`f_\tau :\tau \tau `$, thereby reducing the size of the folding automata needed in our computation.
We first assume that only the component of $`F\tau `$ containing $`p_0`$ has expanding edges on its sides: the other components of $`F\tau `$ not containing $`p_0`$ are bounded only by infinitesimal edges. If one is given a train track representative $`f_\tau :\tau \tau `$ not satisfying this assumption, he can apply a splitting operation \[5, Section 5\] nearby $`p_0`$ (when $`p_0`$ is enclosed only by infinitesimal edges) then apply a sequence of folding operations \[4, p.15\] \[16, Section 2.2\] nearby other punctures $`p_i`$, $`i0`$, until all the components of the train track complement not containing $`p_0`$ shrink to be infinitesimal, to obtain a new train track representative satisfying the assumption \[4, Proposition 3.3\].
Applying some more folding operations (see Figure 4), we can also remove any cusp occurring between an expanding edge and an infinitesimal edge.
We assume that cusps occur only at corners of infinitesimal multigons. If one is given a train track representative with a cusp incident only to expanding edges, not satisfying this assumption, he can apply a splitting operation at the cusp until the cusp hits an infinitesimal multigon (see Figure 5).
Therefore a pseudo-Anosov braid has an invariant train track which is locally modeled by infinitesimal $`k`$-gons to which expanding edges are joined (possibly) forming cusps only between expanding edges (see Figure 6).
In this paper we use simplified versions of folding automata, of which train tracks satisfy the previously given conditions, and each arrow is either an isomorphism or a composite of two elementary folding maps by which one expanding edge and one infinitesimal edge is absorbed into another expanding edge. It is not hard to see that simplified folding automata also generate all the conjugacy classes of pseudo-Anosov homeomorphisms.
In this paper our subject of interest is pseudo-Anosov homeomorphisms on a 5-times punctured disk $`D_5`$, or equivalently on a 6-times punctured sphere $`F_{0,6}`$ with a distinguished boundary puncture.
We explain how to read Figure 7, which depicts a simplified version of a folding automaton. Each train track is embedded in a 5-times punctured disk, with each puncture enclosed by an infinitesimal monogon. Each embedding is chosen arbitrarily, and only the orientation preserving diffeomorphism types of embedded train tracks count.
An arrow is a composite of two elementary folding maps, one involving an infinitesimal edge and another involving only expanding edges. We ignore the infinitesimal edges in computing the transition matrix because the occurrences of infinitesimal edges do not affect the resulting dilatation factor.
An arrow is drawn dashed if it induces a homeomorphism isotopic to identity, and it is drawn by a solid line otherwise. Note that a folding map $`\pi :\tau \tau ^{}`$ determines a disk homeomorphism $`f:D_5D_5`$ up to isotopy when the embeddings $`\tau D_5`$ and $`\tau ^{}D_5`$ of the two train tracks are fixed. In particular $`\tau ^{}f(\tau )`$, that is, $`f(\tau )`$ folds to be $`\tau ^{}`$ inducing the folding map $`\pi `$. To each solid arrow, a braid word is assigned representing the associated disk homeomorphism.
Edges of a train track are numbered by $`\{1,2,\mathrm{},6\}`$ in such a way that in the peripheral word running clockwise from a cusp, new edges appear in an increasing order. This naming of edges amounts to fixing a groupoid homomorphism from paths in the automaton to transition matrices, that is, for two paths $`\gamma `$ and $`\delta `$, $`M(\gamma \delta )=M(\gamma )M(\delta )`$ if $`\gamma `$ ends at the starting vertex of $`\delta `$, where $`M(\gamma )`$ denotes the transition matrix for $`\gamma `$.
Each arrow is associated with a permutation $`i_1i_2i_3i_4i_5i_6`$ and a rule $`mn`$, meaning that under the elementary folding map, the edge $`j`$ maps to $`i_j`$ for $`jm`$, and $`m`$ maps to $`i_mn`$. (Here we concern only the transition matrix so that the direction of edges and the order of concatenation are irrelevant.)
Given an adjacent pair $`(e_1,e_2)`$ of edges with a cusp between them, there are two possible folding maps: one under which the image of $`e_1`$ overpasses that of $`e_2`$ and another vice versa. Therefore from each train track in Figure 7, two arrows of elementary folding maps emanate. Likewise two arrows are headed for each train track, because at each cusp there are two different elementary splittings possible.
## 3. Search for the minimum dilatation
In this section we prove that the largest zero $`\lambda _5`$ of $`x^4x^3x^2x+1`$ is indeed the minimum dilatation for pseudo-Anosov 5-braids.
The problem for the minimum dilatation reduces to a search in a finite set of closed paths in folding automata because by \[19, Theorem 6\] or by Lemma 3.1 the dilatation grows as the norm of the transition matrix grows, and there are only finitely many closed paths whose transition matrices have norm bounded by a given number. For instance if a closed path in folding automata has length $`N`$, then its associated transition matrix has norm at least $`N`$.
We first restate and prove the lemma in Section 1.
###### Lemma 3.1.
If $`M`$ is a Perron-Frobenius matrix of dimension $`n`$ with $`\lambda >1`$ its largest eigenvalue, then
$$\lambda ^n|M|n+1$$
where $`|M|`$ denotes the sum of entries of $`M`$.
###### Proof.
Let $`M=(m_{ij})`$ and let $`(v_i)`$ be the eigenvector given by the equation
$$\lambda v_i=\underset{j=1}{\overset{n}{}}m_{ij}v_j$$
for $`v_i>0`$, $`1in`$.
The matrix $`M`$ is the transition matrix of a graph $`G`$ with vertex set $`V(G)=\{1,2,\mathrm{},n\}`$ such that the number of oriented edges from $`i`$ to $`j`$ equals $`m_{ij}`$. Let $`M^n=(k_{ij})`$. The number of paths with length $`n`$ from $`i`$ to $`j`$, equals $`k_{ij}`$. For each pair $`(i,j)`$ of vertices there exists an oriented path from $`i`$ to $`j`$ since $`M`$ is Perron-Frobenius.
Note that $`(v_i)`$ is also the eigenvector of $`M^n`$ with eigenvalue $`\lambda ^n`$. Choose $`v_p=\mathrm{min}_iv_i`$ the smallest coordinate of $`(v_i)`$.
$`\lambda ^nv_p`$ $`={\displaystyle \underset{j}{}}k_{pj}v_j`$
$`({\displaystyle \underset{j}{}}k_{pj})v_p\text{since }v_jv_p`$
The inequality $`\lambda ^n_jk_{pj}`$ reads that $`\lambda ^n`$ is not less than the number of length-$`n`$ paths from the vertex $`p`$ of $`G`$.
Take a maximal positive tree $`TG`$ rooted at $`p`$ : each vertex of $`G`$ is connected to $`p`$ by a unique oriented path in $`T`$. As $`T`$ is maximal, $`|V(T)|=|V(G)|=n`$ so that the number of edges of $`T`$ is $`|E(T)|=n1`$.
If an oriented path $`\gamma `$ from $`p`$ has length $`n`$, it must digress from $`T`$ at some point. Define $`a(\gamma )=eE(G)E(T)`$ to be the first edge *not* in $`E(T)`$ of $`\gamma `$. This defines a function
$$a:\{\text{length-}n\text{ paths from }p\text{ in }G\}E(G)E(T).$$
Clearly $`a`$ is a surjection since the tail of each edge in $`E(G)E(T)`$ is connected to $`p`$ by a path in $`T`$ with length at most $`n1`$. Therefore
$$\lambda ^n\underset{j}{}k_{pj}|E(G)E(T)|=|M|(n1).$$
###### Theorem 3.2 ().
If a pseudo-Anosov 4-braid has an invariant foliation with one non-punctured 3-prong singularity, then its dilatation is not less than $`\lambda _42.29663`$, the largest zero of $`x^42x^32x+1`$.
###### Proof.
Let $`\beta B_4`$ be a pseudo-Anosov 4-braid in the theorem. Up to conjugacy and multiplication by central elements, $`\beta `$ appears as a closed path $`\gamma `$ in the folding automaton in Figure 8.
If $`\lambda (\beta )<2.3`$, then by Lemma 3.1 we have a bound $`31>2.3^4+41|M_\gamma |`$ for the norm of its transition matrix $`M_\gamma `$.
By a computer aided search in the finite set of paths $`\gamma `$ with $`|M_\gamma |<31`$, we conclude that up to conjugacy, multiplication by central elements, and taking inverse, the braid $`\sigma _3\sigma _2\sigma _1^1`$ is the only pseudo-Anosov 4-braid with dilatation less than $`2.3`$. It can be easily checked that $`\lambda (\sigma _3\sigma _2\sigma _1^1)=\lambda _4`$. ∎
We say two matrices have the same *pattern* if they have zero entries and positive entries in the same positions. We write $`MM^{}`$ for $`M=(m_{ij})`$ and $`M^{}=(m_{ij}^{})`$ if $`m_{ij}m_{ij}^{}`$ for all $`i,j`$.
In the following lemma we ignore the parts of transition matrices arising from infinitesimal edges, so that for a closed path to represent a pseudo-Anosov homeomorphism implies for its transition matrix to be Perron-Frobenius.
###### Lemma 3.3.
Let $`\gamma `$ be a closed path in a folding automaton. Let $`N,k>0`$ be numbers such that the transition matrices $`M(\gamma ^{N+i+k})`$ and $`M(\gamma ^{N+i})`$ have the same pattern and $`M(\gamma ^{N+i+k})M(\gamma ^{N+i})`$ for any $`i0`$. Then a closed path of the form $`\alpha \gamma ^{N+i+k}\delta `$ in the folding automaton represents a pseudo-Anosov homeomorphism if and only if $`\alpha \gamma ^{N+i}\delta `$ does. Furthermore in this case we have an inequality
$$\lambda (\alpha \gamma ^{N+i}\delta )\lambda (\alpha \gamma ^{N+i+k}\delta )$$
between their dilatation factors.
###### Proof.
It suffices to prove the lemma for $`\delta \alpha \gamma ^{N+i+k}`$ and $`\delta \alpha \gamma ^{N+i}`$ since conjugation does not affect dilatation factor or being pseudo-Anosov.
Since $`M(\gamma ^{N+i+k})`$ and $`M(\gamma ^{N+i})`$ have the same pattern, $`M(\delta \alpha )M(\gamma ^{N+i+k})`$ and $`M(\delta \alpha )M(\gamma ^{N+i})`$ also have the same pattern. In particular one is Perron-Frobenius if and only if so is the other, which proves the first claim of the lemma.
From $`M(\gamma ^{N+i})M(\gamma ^{N+i+k})`$, we have
$$M(\delta \alpha \gamma ^{N+i})M(\delta \alpha \gamma ^{N+i+k}),$$
which by \[21, Theorem 1.1 (e)\] implies the inequality of the lemma. ∎
###### Remark 3.4.
Let $`\gamma ,N,k`$ be given as in Lemma 3.3. Then the lemma implies that when we search just for the minimum dilatation factor for pseudo-Anosov homeomorphisms, it suffices to search in the set of paths that do not contain $`\gamma ^{N+k}`$ as a subpath.
In the search in the automaton in Figure 7, we exclude paths containing several closed paths, for example $`\left(\genfrac{}{}{0pt}{}{123564}{14}\genfrac{}{}{0pt}{}{123456}{14}\right)^6`$, $`\left(\genfrac{}{}{0pt}{}{123456}{41}\genfrac{}{}{0pt}{}{312456}{43}\right)^6`$, and second iterates of length 1 loops. This reduces the size of the set of candidate braids for minimum dilatation to the extent that the computation in the proof of Theorem 3.5 becomes possible on a personal computer.
###### Theorem 3.5.
If a pseudo-Anosov 5-braid has an invariant foliation with two non-punctured 3-prong singularities, then its dilatation is not less than the largest zero $`2.01536`$ of $`x^6x^54x^3x+1`$.
###### Proof.
It is easy to check that there are only nine different diffeomorphism types of train tracks in $`D_5`$, locally modeled by infinitesimal multigons with outgoing expanding-edge legs as in Figure 6. By computing the elementary folding maps among them (more precisely composites of two elementary folding maps, one of them involving an infinitesimal edge), we have a folding automaton depicted in Figure 7.
By a computer aided search in the set of paths $`\gamma `$ with $`|M_\gamma |2.02^6+5<73`$, we conclude that up to conjugacy and multiplication by central elements, $`\sigma _4\sigma _3\sigma _1^1\sigma _2^1`$ with dilatation $`2.01536`$ is the only such pseudo-Anosov 5-braid with dilatation less than $`2.02`$. ∎
###### Lemma 3.6.
If a pseudo-Anosov 5-braid has an invariant foliation with a non-punctured 4-prong singularity, then its dilatation is not less than the largest zero $`2.15372`$ of $`x^43x^3+3x^23x+1`$.
###### Proof.
The folding automaton for this case is similar to one in Figure 8. As in the proof of Theorem 3.2, a computer aided search in the set of closed paths up to length $`56>2.2^5+(51)`$ shows that the largest zero of $`x^43x^3+3x^23x+1`$ is the minimum dilatation factor in the automaton, and it is achieved by $`\sigma _4\sigma _3\sigma _2\sigma _1^1`$. ∎
###### Lemma 3.7.
If a pseudo-Anosov 5-braid has an invariant foliation with a punctured 3-prong singularity, then its dilatation is not less than $`\lambda _5`$, the largest zero of $`x^4x^3x^2x+1`$.
###### Proof.
Let $`f:F_{0,6}F_{0,6}`$ be a pseudo-Anosov homeomorphism with an invariant foliation $``$. If the invariant measured foliation $``$ on a 6-times punctured sphere has a punctured 3-prong singularity, then it has five other punctured 1-prong singularities and no more.
We can assume that the punctured 3-prong singularity is the boundary puncture since it should be fixed by the homeomorphism $`f`$. Now we use the folding automaton that generates such pseudo-Anosov braids with three prongs at the boundary puncture.
There are eleven diffeomorphism types of train tracks to consider for this case (see Figure 9). There are fifty arrows in the automaton, which are too many to be drawn in a figure in this paper. See for details.
By the same kind of computer aided search as before, in the set of closed paths in the folding automaton up to length $`12(\lambda _5)^4+(41)`$, we conclude that $`\lambda _5`$ is the minimum dilatation factor for pseudo-Anosov braids in this automaton.
The dilatation is achieved by $`\sigma _1\sigma _2\sigma _3\sigma _4\sigma _1\sigma _2`$. ∎
###### Remark 3.8.
In , $`\lambda _5`$ is proved to be the minimum dilatation factor for a pseudo-Anosov homeomorphism with an orientable invariant foliation on a closed genus-2 surface. The proof in seems to have a gap. To complete the proof one needs to show that the golden ratio $`(1+\sqrt{5})/21.61803`$ the largest zero of $`x^43x^2+1`$, cannot be a dilatation factor for such a pseudo-Anosov homeomorphism. In , it is proved that such a pseudo-Anosov homeomorphism with quadratic dilatation factor is a lift of an Anosov homeomorphism via a branched covering. Lemma 3.6 and Lemma 3.7 follow from by taking double covers branched at odd-prong singularities.
By collecting all the results, we conclude this paper with a proof of the main theorem.
###### Theorem 3.9.
The 5-braid $`\sigma _1\sigma _2\sigma _3\sigma _4\sigma _1\sigma _2`$ attains the minimum dilatation of pseudo-Anosov 5-braids.
###### Proof.
Let $`f:F_{0,6}F_{0,6}`$ be a pseudo-Anosov homeomorphism on a 6-times punctured sphere with a punctured point fixed by $`f`$. Let $``$ be its invariant measured foliation. Since $``$ has exactly six punctures, the formula $`2=\chi (F_{0,0})=_k(1k/2)n_k`$ where $`n_k`$ denotes the number of $`k`$-prong singularities, says no singularity of $``$ can have more than four prongs.
The list of possible types of $``$ according to its singularity type is:
1. six punctured 1-prong singularities and one non-punctured 4-prong singularity: $`n_1=6,n_4=1`$
2. six punctured 1-prong singularities and two non-punctured 3-prong singularities: $`n_1=6,n_3=2`$
3. five punctured 1-prong singularities and one punctured 3-prong singularity: $`n_1=5,n_3=1`$
4. five punctured 1-prong singularities, one punctured 2-prong singularity, and one non-punctured 3-prong singularities: $`n_1=5,n_2=1,n_3=1`$
5. four punctured 1-prong singularities and two punctured 2-prong singularities: $`n_1=4,n_2=2`$
The case (5) is, by capping-off 2-prong singularity punctures, of a pseudo-Anosov homeomorphism on a four-times punctured sphere, which lifts to an Anosov homeomorphism on a torus via branched double covering. Therefore in this case $`\lambda (f)(3+\sqrt{5})/2>\lambda _5`$.
The case (4) reduces to the case (3) by capping-off the punctured 2-prong singularity and puncturing at the 3-prong singularity.
For the cases (1) and (3) we have $`\lambda (f)\lambda _5`$ by Lemma 3.6 and Lemma 3.7 or by Remark 3.8.
Finally the case (2) is covered by Theorem 3.5, so that we have $`\lambda (f)>2.01>\lambda _5`$. In fact this is the only part of the proof which actually requires a computer aided search if one uses Zhirov’s result .
Collecting all of these we conclude that $`\lambda (f)\lambda _51.72208`$. It is easily checked that $`\beta =\sigma _1\sigma _2\sigma _3\sigma _4\sigma _1\sigma _2`$ realizes this dilatation $`\lambda (\beta )=\lambda _5`$. ∎
## 4. Implementation
To search for the minimum-dilatation pseudo-Anosov homeomorphism on a given surface, we first need to generate a collection of folding automata. For 5-braids it is possible to build the necessary folding automata manually. On surfaces with more punctures and greater genus, we also need a computer program genauto to generate the folding automata. This paper will not cover the details of its implementation. The following is a pseudo-code for genauto.
genauto genus $`g`$ and the number of punctures $`n`$ folding automata on $`F_{g,n}`$ Generate the finite set of diffeomorphism types of embedded train tracks $`\tau _iF_{g,n}`$. For each $`\tau _i`$, compute all the elementary folding maps $`f_{ij}:\tau _i\tau _{ij}^{}`$ from $`\tau _i`$ if any. Compute isomorphisms $`h_{ij}:\tau _{ij}^{}\tau _k`$ from the train track $`\tau _{ij}^{}`$ to one in the set $`\{\tau _i\}`$. For each $`\tau _i`$, compute the isomorphisms $`g_i\mathrm{}:\tau _i\tau _i`$ if any. The elementary folding maps $`h_{ij}f_{ij}`$ and the isomorphisms $`g_i\mathrm{}`$ form the arrows of the folding automata. Compute their transition matrices after labeling each edge of all the train tracks $`\tau _i`$.
Note that for step 1 one needs to solve the isomorphism problem for embedded train tracks. Once step 1 is done, implementing the other steps is more straightforward.
By running genauto, we obtain the folding automata as a collection of connected directed graph with each arrow labeled by a transition matrix. The goal is to enumerate in the folding automata all the closed paths representing pseudo-Anosov mapping classes with an upper bound for the dilatation.
In this paper we deal with 5-braids using simplified folding automata. We ignore infinitesimal edges when computing transition matrices. Therefore a closed path in a folding automaton represents a pseudo-Anosov braid if and only if its associated transition matrix is Perron-Frobenius.
The following is a pseudo-code for our program fbrmin. See for details.
fbrmin a directed graph $`𝒢`$ with arrows labeled by transition matrices,
an upper bound $`\lambda `$ for the minimum dilatation,
and a set $`𝒲`$ of sub-words which are to be avoided during the search the list of closed paths in $`𝒢`$ representing pseudo-Anosov braids with dilatation less than $`\lambda `$ Set $`\mathrm{𝚖𝚊𝚡𝚗𝚘𝚛𝚖}=\lambda ^n+n1`$ where $`n`$ is the dimension of the transition matrices,
and set $`\mathrm{𝚊𝚛𝚌𝚑𝚒𝚟𝚎}=\{\}`$, in which closed paths with small dilatation are to be stored. Set $`\mathrm{𝚝𝚖𝚙}_1`$ to be the set of length-one paths in $`𝒢`$. For each $`i`$ from 2 to $`\mathrm{𝚖𝚊𝚡𝚗𝚘𝚛𝚖}`$, compute $`\mathrm{𝚌𝚑𝚒𝚕𝚍𝚛𝚎𝚗𝚙𝚊𝚝𝚑𝚜}_i`$ by appending paths in $`\mathrm{𝚝𝚖𝚙}_1`$ to paths in $`\mathrm{𝚝𝚖𝚙}_{i1}`$, in all the ways possible in $`𝒢`$ compute $`\mathrm{𝚝𝚖𝚙}_i`$ the subset of $`\mathrm{𝚌𝚑𝚒𝚕𝚍𝚛𝚎𝚗𝚙𝚊𝚝𝚑𝚜}_i`$, consisting of paths $`\beta `$ without any subword from the avoided-word set $`𝒲`$, with transition matrix $`M_\beta `$ such that $`|M_\beta |\mathrm{𝚖𝚊𝚡𝚗𝚘𝚛𝚖}`$, and $`M_\beta `$ has at least one row and one column whose row (column) sum is less than 3. Take the subset $`\mathrm{𝚜𝚎𝚕𝚎𝚌𝚝𝚎𝚍𝚌𝚊𝚗𝚜}_i`$ of $`\mathrm{𝚝𝚖𝚙}_i`$ consisting of *closed* paths representing pseudo-Anosov braids with dilatation less than $`\lambda `$, and append it to $`\mathrm{𝚊𝚛𝚌𝚑𝚒𝚟𝚎}`$. Return $`\mathrm{𝚊𝚛𝚌𝚑𝚒𝚟𝚎}`$ $`(=_i\mathrm{𝚜𝚎𝚕𝚎𝚌𝚝𝚎𝚍𝚌𝚊𝚗𝚜}_i)`$.
In step 3 (2), we use Lemma 3.13.3 to trim out much of unnecessary computation (see Remark 3.4).
When the row sums of a transition matrix $`M_\beta `$ all exceed $`3`$, then the spectral radius of $`M_\beta `$ is greater than $`3`$. In this case the same holds for every transition matrix of the form $`M_{\beta \gamma }=M_\beta M_\gamma `$ since $`M_\gamma P`$ for some permutation matrix $`P`$. Therefore as we are looking for transition matrices with spectral radius less than 3, we can safely remove such paths $`\beta \gamma `$ from our consideration, as done in step 3 (2).
For computational aspects, the proof of Theorem 3.5 using the automaton in Figure 7 is the main part which consumes most time and memory. On a 2.40 GHz machine, it took 1000 seconds of time and 150 mega-bytes of memory. During the search roughly 85000 many matrices were actually tested for its largest eigenvalue.
We do not know how far the same kind of computation would work for more complicated surfaces. We expect that at least the case for 6-braids, hence for genus-2 closed surface, can be done on a personal computer without too much difficulty.
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# New Method of Enhancing Lepton Number Nonconservation
## Abstract
The lepton number nonconserving (LENNON) conversion of type $`e^{}e^+`$ in heavy atoms, when irradiated by intense laser beam, is considered to determine the Majorana nature and precise values of neutrino masses. When the photon energy is fine tuned , the LENNON process is greatly enhanced by both the resonance effect and a large occupancy of photons. The signal of this process would be a positron of definite energy, and a further nuclear $`\gamma `$-ray in a favorable case of transition to an excited nuclear level. By constructing a united target-detector system, it is possible to explore $`O[1meV]`$ range of the neutrino mass parameter with a good positron energy resolution.
Recent observations indicate that neutrinos have finite masses. The immediate question that arises is whether these masses are of Dirac or Majorana type. If they are of the Majorana type, the lepton number is not conserved. If this nonconservation were discovered, one should expect great impacts on other areas of physics, including the possible explanation of baryon-antibaryon imbalance in our universe th leptogenesis . The most extensively examined tool to study the Majorana nature of the neutrino mass is the neutrinoless double beta decay, and many proposals and experiments are already on-going double beta limit .
We wish to propose alternatives to the neutrinoless double beta decay in order to investigate the nature of neutrino masses, both because the important issue of the lepton number nonconservation (abbreviated as LENNON) requires an independent experimental check, and because one should examine it in many other processes. We shall focus in the present work on photon stimulated electron conversion of type $`e^{}e^+`$ in heavy atoms, with no accompanying neutrino. The experimental signature of this process is a monochromatic positron of few $`MeV`$, and furthermore nuclear $`\gamma `$-ray in a favorable case of transition to nuclear excited states. Although simultaneous emission of low energy photons occurs, it might be difficult to detect them, because their energies are either down shifted or the same laser photon energy in the case of united target-detector system later discussed, and they might easily be confused amidst the background of high intensity beam. We give rates both for this process and the background of similar process of two accompanying neutrinos.
The basic low energy effective operator lty-04 that may arise in physics beyond the standard model is $`G_F^2\stackrel{~}{m}_\nu ll\overline{q}q\overline{q}q`$ where $`l`$ is the lepton doublet and $`q`$ is the quark doublet having the quantum number of the standard model. This class of effective interaction violates the lepton number by two units, $`\mathrm{\Delta }L=2`$, a class of simplest models predicting the amplitude in direct proportion to some combination of Majorana type of neutrino masses, $`\stackrel{~}{m}_\nu `$. Thus, the rate of the non-radiative LENNON conversion of the type $`e^{}(1s)e^+`$, even for a favorable case of $`{}_{}{}^{112}Sn`$, becomes of order $`10^{29}year^1`$ for the Majorana neutrino mass of order $`0.1eV`$.
We instead consider LENNON conversion stimulated by radiative absorption, of the type depicted in Fig. 1. Here $`{}_{Z}{}^{A}X`$ is a nucleus of mass number $`A`$ and atomic number $`Z`$. An electron occupying an atomic $`n`$ state is upshifted by photon irradiation to an unoccupied $`ms`$ level, and later captured by nucleus, producing a positron by $`e^{}e^+`$ conversion. Another possibility is the positron emission first by weak interaction and a subsequent absorption of photon by atomic electron of $`{}_{Z1}{}^{A}Y`$. The idea is to enhance the $`e^{}e^+`$ conversion with the help of both the resonance effect and a high occupancy of laser photons in a quantum level. Without irradiation of the initial photon, this is a process conjugate to the neutrinoless double beta decay, and can explore the same combination of the neutrino mass parameter $`m_{ee}`$, where $`m_{\alpha \beta }=_kU_{\alpha k}U_{\beta k}m_k`$ with $`m_k`$ 3 neutrino mass eigenvalues.
In Fig.2 we depict relevant nuclear levels for three adjacent nuclei of the same mass number. The final nucleus $`{}_{Z2}{}^{A}W`$ can be in an excited nuclear level, for which triple coincidence for detection including nuclear $`\gamma `$-ray may become possible. If both $`\beta ^{}`$ decay and electron capture rates of the intermediate $`{}_{Z1}{}^{A}Y`$ nucleus are experimentally known, one may reliably estimate the rate for the radiative $`e^{}e^+`$ conversion. The best candidate with this regard is the isotope $`{}_{50}{}^{112}Sn`$ Sn .
A straightforward computation for the atomic $`nm`$ transition followed by the capture gives the following matrix element;
$$\begin{array}{c}\frac{eG_F^2m_{ee}}{4\pi }\mathrm{\hspace{0.17em}2}\pi \delta (E_\gamma +\mathrm{\Delta }E_{i,f}E_++\mathrm{\Delta }ϵ_{nm})\hfill \\ \hfill \times d^3xd^3yf|\frac{J(\stackrel{}{x})J(\stackrel{}{y})}{|\stackrel{}{x}\stackrel{}{y}|}|i\stackrel{}{p}_+|\overline{e^c}(\stackrel{}{x})(1\gamma _5)e(\stackrel{}{y})\underset{m}{}(|m\frac{m|H_\gamma |n}{E_\gamma \mathrm{\Delta }ϵ_{nm}+i\stackrel{~}{\mathrm{\Gamma }}_n/2}),\end{array}$$
(1)
where $`e(\stackrel{}{y})`$ etc. is the electron field operator. Here $`\mathrm{\Delta }E_{i,f}=E_iE_f`$, with energies $`E_{i,f}`$ referring to nuclear levels of initial and final states. Thus, $`E_+=\mathrm{\Delta }E_{i,f}+E_\gamma +\mathrm{\Delta }ϵ_{nm}`$ is the monochromatic positron energy. The important natural width $`\stackrel{~}{\mathrm{\Gamma }}_n`$ in (1) refers to that of the hole of $`n`$ atomic state. The energy difference $`\mathrm{\Delta }ϵ_{nm}=ϵ_nϵ_m`$ where $`ϵ_{n,m}`$ are energies of atomic levels.
The neutrino propagator has been replaced by the instantaneous Coulomb potential, which is allowed since for low energy electrons the energy transfer is small; $`|q_0||\stackrel{}{q}|`$. At low energies the electron wave functions can be taken out from the integral (1), and one may separate out the nuclear matrix element. Furthermore, by introducing an average inter-proton distance, $`R_n`$, one may take out the factor $`1/|\stackrel{}{x}\stackrel{}{y}|=1/R_n`$ outside the nuclear matrix element. We use $`R_n(0.82A^{1/3}+0.58)fm,`$ following Sn .
The radiative electron conversion via $`ms`$ atomic state has a cross section of the form,
$`\mathrm{\Gamma }_{0\nu }^{(mS)}{\displaystyle \frac{|mS|H_\gamma |n,\gamma |^2}{(E_\gamma \mathrm{\Delta }ϵ_{nm})^2+\stackrel{~}{\mathrm{\Gamma }}_n^2/4}},`$ (2)
where $`H_\gamma `$ is QED interaction. Let us first consider the non-radiative conversion rate $`\mathrm{\Gamma }_{0\nu }^{(ms)}=|\psi _{ms}(0)|^2\sigma _{0\nu }`$. This product is an effective luminosity $`|\psi _{ms}(0)|^2`$ of confined $`ms`$ electron times the cross section for $`e^{}(\mathrm{free})+_Z^AXe^++_{Z2}^AW`$ given by
$`\sigma _{0\nu }={\displaystyle \frac{G_F^4\stackrel{~}{m}_\pi ^2|m_{ee}|^2}{16\pi ^3}}p_+E_+,`$
$`\stackrel{~}{m}_\pi {\displaystyle d^3xd^3y\frac{f|J(\stackrel{}{x})J(\stackrel{}{y})|i}{|\stackrel{}{x}\stackrel{}{y}|}}.`$ (3)
Typically, $`\sigma _{0\nu }=O[10^{66}cm^2]|m_{ee}/1eV|^2(p_+E_+/MeV^2)`$. When a photon is irradiated, the product is further multiplied by the Breit-Wigner function, which can be very large near the resonance, $`E_\gamma =\mathrm{\Delta }ϵ_{nm}E_0`$. In other words, the cross section ratio of the photon-irradiated conversion, (2) relative to the elementary $`\sigma _{0\nu }`$, is
$`K(E)`$ $`=`$ $`{\displaystyle \frac{|mS|H_\gamma |n,\gamma |^2|\psi _{ms}(0)|^2}{(EE_0)^2+\stackrel{~}{\mathrm{\Gamma }}_n^2/4}}`$ (4)
$``$ $`{\displaystyle \frac{2\pi ^2\mathrm{\Gamma }_{mSn}|\psi _{ms}(0)|^2}{E_0^2[(EE_0)^2+\stackrel{~}{\mathrm{\Gamma }}_n^2/4]}}.`$
Here the branching fraction $`B_{msn}=\mathrm{\Gamma }_{msn}/\stackrel{~}{\mathrm{\Gamma }}_n`$ is of order unity.
The actual rate, when the photon beam of luminosity density $`I(E)`$, within unit energy bin and per unit area times unit time, is irradiated, is given by
$`R`$ $`=`$ $`{\displaystyle 𝑑EI(E)\sigma _{\gamma e^{}e^+}}`$ (5)
$`=`$ $`S{\displaystyle 𝑑E\frac{I(E)}{(EE_0)^2+\stackrel{~}{\mathrm{\Gamma }}_n^2/4}}.`$
Note the dimensionless strength factor given by
$`S`$ $`=`$ $`2\pi ^2\sigma _{0\nu }|\psi _{ms}(0)|^2{\displaystyle \frac{B_{mSn}\stackrel{~}{\mathrm{\Gamma }}_n}{E_0^2}}`$ (6)
$``$ $`2\pi \sigma _{0\nu }(Z\alpha )^6m_e^3{\displaystyle \frac{B_{msn}\stackrel{~}{\mathrm{\Gamma }}_n}{E_0^2}}.`$
The integral in (5) may readily be evaluated when the photon flux has little variation over the natural width $`\stackrel{~}{\mathrm{\Gamma }}_n`$ around $`E_0`$:
$`{\displaystyle 𝑑E\frac{I(E)}{(EE_0)^2+\stackrel{~}{\mathrm{\Gamma }}_n^2/4}}2\pi {\displaystyle \frac{I(E_0)}{\stackrel{~}{\mathrm{\Gamma }}_n}}.`$
When the photon beam is tuned and the resonance $`E_0`$ is not missed, one further has $`I(E_0)F(E_0)/\mathrm{\Delta }E,`$ where $`\mathrm{\Delta }E`$ is the energy resolution of photon beam and $`F(E_0)=𝑑EI(E)`$ is the total flux integrated over the region around $`E_0`$. This leads to the rate formula,
$$R\frac{2\pi ^2F(E_0)}{E_0^2\mathrm{\Delta }E}B_{msn}2\pi |\psi _{ms}(0)|^2\sigma _{0\nu }.$$
(7)
Each factor has a clear meaning; the first $`F(E_0)/(E_0^2\mathrm{\Delta }E(2\pi ^2)^1)`$ expressing the number of occupied photons within the relevant quantum phase space, the second $`B_{msn}`$ the branching fraction of order unity, the last $`|\psi _{ms}(0)|^2\sigma _{0\nu }`$ the $`e^{}e^+`$ conversion rate from $`ms`$ state. Thus, the enhancement or reduction factor by photon irradiation relative to the capture rate $`|\psi _{1s}(0)|^2\sigma _{0\nu }`$ is given by $`\mathrm{\hspace{0.25em}4}\pi ^3F(E_0)(|\psi _{ms}(0)|^2/E_0^3)(E_0/\mathrm{\Delta }E)B_{msn}/m^6.`$ It is useful to define a quality factor $`Q`$ of photon beam defined by $`Q=4\pi ^3F(E_0)/(E_0^2\mathrm{\Delta }E)`$, which is numerically $`2\times 10^{10}(F(E_0)/1kWmm^2)(E_0/eV)^4(10^9E_0/\mathrm{\Delta }E)`$. An advantage of a strong laser beam, which can give a large $`Q`$, is obvious, when compared to X-ray which typically gives $`Q1`$ (except contemplated X-ray laser).
One might wonder the validity of the linear rise with the laser power of the rate, eq.(7), because irradiated atoms may become transparant once electrons in the lower energy level are completely lifted to the higher energy level. However, we are considering the situation of equilibrated atoms between the two levels by constant irradiation of laser beam. In this case the rate formula (7) is still valid, with a minor multiplication of the population factor in the lower level. The formula valid in the large power limit is different and shall be presented elsewhere by one of the present authors (MY).
We shall next discuss how to estimate the rate $`\mathrm{\Gamma }_{0\nu }`$ for the non-radiative electron conversion by nucleus that appears in the fundamental formula (2). By truncating the nuclear level sum in $`_kf|J|kk|J|i`$ by a single ground state $`|k=|_{Z1}^AY,1^+`$, one may replace the nuclear matrix element $`k|J|i`$ by the beta decay rate, and $`f|J|k`$ by the electron capture rate of nucleus $`{}_{Z1}{}^{A}Y`$. In the case of the nonradiative decay, $`e^{}(1s)e^+`$ process, the rate becomes
$`\mathrm{\Gamma }_{0\nu }={\displaystyle \frac{3\pi ^3|m_{ee}|^2}{4R_n^2}}{\displaystyle \frac{p_+E_+}{p_\nu ^2\mathrm{\Delta }_\beta ^5I}}\mathrm{\Gamma }_{EC}\mathrm{\Gamma }_\beta ,`$ (8)
where $`\mathrm{\Gamma }_{EC},\mathrm{\Gamma }_\beta `$ are the electron capture and the beta decay rate of $`{}_{Z1}{}^{A}Y`$, and $`p_i`$ are respective lepton momenta. The maximum $`\beta `$ energy $`\mathrm{\Delta }_\beta =Q_\beta +m_eϵ_{1s}`$, and $`p_\nu =Q_{EC}ϵ_{1s}`$. Here $`I`$ is a dimensionless phase space factor for the beta decay and $`1/30`$ in the limit of zero electron mass. We ignored the difference of $`1s`$ electron wave function of $`{}_{Z}{}^{A}X`$ and $`{}_{Z2}{}^{A}W`$, whose error should be small, of order $`4/Z`$.
For instance, in the case of $`0^+0^+`$ nuclear transition of $`{}_{50}{}^{112}Sn_{48}^{112}Cd`$, both intermediate $`{}_{49}{}^{112}In`$ (of $`1^+`$) $`\beta ^{}`$ decay and electron capture rates are known nuclear data . Thus, the rate computed according to eq.(8) is $`(2\times 10^{29}y)^1(|m_{ee}|/0.1eV)^2`$ for $`{}_{50}{}^{112}Sn`$. It appears that nuclear matrix elements are large, as pointed out in Sn .
We numerically give the laser irradiated $`e^{}e^+`$ conversion rate for one nucleus target;
$`R`$ $``$ $`5\times 10^{35}y^1`$ (9)
$`\times {\displaystyle \frac{\mathrm{\Gamma }_{0\nu }}{10^{29}y^1}}\left({\displaystyle \frac{4}{m}}\right)^6{\displaystyle \frac{E_\gamma }{\mathrm{\Delta }E}}{\displaystyle \frac{F}{Wmm^2}}\left({\displaystyle \frac{E_\gamma }{eV}}\right)^4.`$
The crucial factor to obtain a large enhancement of the rate is the inverse of resolution, and with $`\mathrm{\Delta }E/E_\gamma =10^9`$ available commercially, the rate is enhanced by $`5\times 10^6(4/m)^6`$, using a laser beam power $`1kW/mm^2`$. Linear increase with the intensity $`F`$ of this event rate is the key check point of experimental verification of the process. A few isotopes which may give rates of $`\mathrm{\Gamma }_{0\nu }>10^{31}y^1`$ are illustrated in Table 1.
The required resonance tuning might be a great practical obstacle since laser frequencies are superposition of quantized level differences. (Alternatively, the use of laser with a continuous spectrum might help greatly.) It is however possible to avoid this problem by selecting the same lasing medium as the target nucleus, thus we arrive at the concept of a united target-detector system. A good target must then be both lasing and an efficient $`e^{}e^+`$ converter. Candidates of such targets are limited, but the element $`Kr`$ is a good example. The $`Kr^+`$ laser uses the inverted population of $`4p`$ atomic state, which falls down to $`4s`$ by stimulated emission caused by irradiated laser laser . In this basic process $`4s`$ electron may very rarely be captured by $`0^+`$ $`Kr`$ nucleus to emit a positron of energy $`1.8MeV`$.
The isotope $`{}_{36}{}^{78}Kr`$ after the $`e^{}e^+`$ conversion ends up with $`{}_{34}{}^{78}Se`$, which has a few excited energy levels of $`0^+`$ below the $`{}_{36}{}^{78}Kr`$ ground level. Thus, nuclear gamma rays are expected to give an opportunity of coincident measuremet.
Let us estimate what might occur within an ideal $`Kr^+`$ or gaseous laser device containing $`Kr`$ such as $`KrF`$ excimer laser. Suppose that the gas chamber of the device contains $`10^4cm^3`$ of 1 ATM enriched $`{}_{36}{}^{78}Kr`$, which has $`3\times 10^{23}`$ $`{}_{36}{}^{78}Kr`$ atoms. The formula (9) tells that $`e^{}e^+`$ conversion occurs with a rate, $`10^4/y(I/MWmm^2),`$ assuming a resolution $`\mathrm{\Delta }E/E_\gamma =10^9`$ and $`\mathrm{\Gamma }_{0\nu }=10^{29}/y`$. For the $`Kr^+`$ laser, $`m=4,E_\gamma =1.9eV`$ for one of the main lines. The rate scales with the neutrino mass parameter as $`\mathrm{\Gamma }_{0\nu }=0.9\times 10^{29}y^1|m_{ee}/0.1eV|^2`$, using nuclear matrix elements of Sn adopted to $`{}_{36}{}^{78}Kr`$, thus one may be able to explore $`|m_{ee}|`$ down to $`1meV`$.
The important physics background to the present process is the corresponding process of two accompanying neutrinos caused by the second order weak interaction, which itself is of interest. We may estimate this rate by using the same approximation of one-level truncation. The result is given by the ratio of two processes, $`\mathrm{\Gamma }_{0\nu }/\mathrm{\Gamma }_{2\nu }(20160\pi ^2|m_{ee}|^2m_\pi ^2)/(\mathrm{\Delta }^6R_n^2)`$ which is $`0.66\times 10^4|m_{ee}/1meV|^2(m_\pi R_n)^2(\mathrm{\Delta }/MeV)^6`$.
On the other hand, the ratio when the positron energy is limited near the end point of $`e^+`$ energy width $`\delta E`$ is $`\mathrm{\Gamma }_{0\nu }/\mathrm{\Gamma }_{2\nu }(\delta E)45\pi ^2m_\pi ^2|m_{ee}|^2/(R_n^2(\delta E)^6).`$ With $`\delta E100keV`$ and for $`{}_{36}{}^{78}Kr`$,
$$\frac{\mathrm{\Gamma }_{0\nu }}{\mathrm{\Gamma }_{2\nu }(\delta E)}O[18]\left|\frac{m_{ee}}{0.01eV}\right|^2\left(\frac{\delta E}{100keV}\right)^6.$$
(10)
The positron energy spectrum for two neutrino process is given by
$`d\mathrm{\Gamma }_{2\nu }E\sqrt{E^2m_e^2}(E\mathrm{\Delta })^5dE.`$ (11)
The signal of LENNON would be an excess of positrons near the end point of two neutrino processes. In the favorable case to decay into a nuclear excited level the background $`(2\nu )`$ process is suppressed by the phase space factor $`(E\mathrm{\Delta })^5`$ if the excited level has a relatively high energy.
In summary, it appears possible to explore the Majorana neutrino mass range down to $`1meV`$ by laser irradiated $`e^{}e^+`$ conversion, when a lasing gas medium $`>10^4cm^3`$ of high power $`>1MW`$ or more, is used, along with a good measured positron energy resolution of $`<50keV`$.
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# The role of algebraic solutions in planar polynomial differential systems. The second and third authors are partially supported by a MCYT grant number BFM 2002-04236-C02-01. The second author is also partially supported by DURSI of Government of Catalonia “Distinció de la Generalitat de Catalunya per a la promoció de la recerca universitària”.
## 1 Introduction
In this work we consider planar polynomial differential systems as:
$$\dot{x}=P(x,y),\dot{y}=Q(x,y),$$
(1)
where $`P(x,y)`$ and $`Q(x,y)`$ belong to the ring of real polynomials in two variables, $`[x,y]`$. We will always assume that $`P(x,y)`$ and $`Q(x,y)`$ are coprime polynomials. We denote by $`\mathrm{d}=\mathrm{max}\{\mathrm{deg}P,\mathrm{deg}Q\}`$ and we say that $`\mathrm{d}`$ is the degree of system (1). Equivalently to system (1), we may consider the ordinary differential equation:
$$\frac{dy}{dx}=\frac{Q(x,y)}{P(x,y)}.$$
(2)
In order to shorten some formulae, we introduce the operator $`𝒳`$ associated to (1):
$$𝒳:=P(x,y)\frac{}{x}+Q(x,y)\frac{}{y}.$$
This paper is related with the study of the properties of a certain type of particular solutions of system (1). Given a function $`g(x)`$, we say that $`y=g(x)`$ is a particular solution of equation (2) if
$$g^{}(x)=\frac{Q(x,g(x))}{P(x,g(x))},$$
(3)
where $`g^{}(x)=dg(x)/dx`$. In particular, we are concerned with algebraic functions. We say that a function $`g(x)`$ is algebraic if there exists a polynomial $`f(x,y)`$ such that $`f(x,g(x))0`$. As we prove in the next section, by using some algebraic results stated for instance in , this polynomial $`f(x,y)`$ can always be chosen irreducible and it is unique modulus multiplication by constants.
The solutions of system (1) may also be given in an implicit way. An invariant algebraic curve of system (1) is an algebraic curve $`f(x,y)=0`$ satisfying
$$P(x,y)\frac{f}{x}(x,y)+Q(x,y)\frac{f}{y}(x,y)|_{f(x,y)=0}=0.$$
(4)
By using Hilbert’s Nullstellensatz, cf. , it can be shown that $`f(x,y)=0`$ is an invariant algebraic curve of system (1) if, and only if, there exists a polynomial $`k(x,y)[x,y]`$ satisfying:
$$P(x,y)\frac{f}{x}(x,y)+Q(x,y)\frac{f}{y}(x,y)=k(x,y)f(x,y).$$
This polynomial $`k(x,y)`$ is called the cofactor of the curve given by $`f(x,y)=0`$ and it can be shown that its degree is lower than or equal to $`\mathrm{d}1`$. Invariant algebraic curves, also denoted as Darboux polynomials by some authors, have been widely studied for their relation with integrability and some qualitative properties of polynomial differential systems, see for instance . As a generalization of the notion of invariant algebraic curve, we can define an exponential factor. This concept has been firstly introduced by Christopher and it is related with the notion of multiplicity of an invariant algebraic curve of system (1). An exponential factor is a function of the form $`\mathrm{exp}\{h/f_0\}`$ where $`h(x,y)`$ and $`f_0(x,y)[x,y]`$, and $`𝒳\left(\mathrm{exp}\{h/f_0\}\right)=\stackrel{~}{k}(x,y)\mathrm{exp}\{h/f_0\}`$ with $`\stackrel{~}{k}(x,y)`$ a polynomial of degree at most $`\mathrm{d}1`$. The following lemma is given in and characterizes exponential factors.
###### Lemma 1
The function $`\mathrm{exp}\{h/f_0\}`$ is an exponential factor of system (1) with cofactor $`\stackrel{~}{k}(x,y)`$ if, and only if, $`f_0(x,y)=0`$ is an invariant algebraic curve of system (1) with cofactor $`k_0(x,y)`$ and $`𝒳(h)=k_0h+\stackrel{~}{k}f_0`$.
Invariant algebraic curves characterize the existence of first integrals for system (1) belonging to a certain functional class. In order to properly state the known results about integrability using invariant algebraic curves, we need to consider complex algebraic curves $`f(x,y)=0`$, where $`f(x,y)[x,y]`$. Since system (1) is defined by real polynomials, if $`f(x,y)=0`$ is an invariant algebraic curve with cofactor $`k(x,y)`$, then its conjugate $`\overline{f}(x,y)=0`$ is also an invariant algebraic curve with cofactor $`\overline{k}(x,y)`$. Hence, its product $`f(x,y)\overline{f}(x,y)[x,y]`$ gives rise to a real invariant algebraic curve with a real cofactor $`k(x,y)+\overline{k}(x,y)`$. In the Darboux theory of integrability, quoted in the forthcoming paragraph, we need to consider invariant algebraic curves defined by polynomials in $`[x,y]`$ since they play an essential role in the theory of integrability. We also notice that in $`^2`$ the curve given by $`f(x,y)=0`$, even if $`f(x,y)[x,y]`$, may only contain a finite number of isolated singular points or be the null set.
A function of the form $`f_1^{\lambda _1}\mathrm{}f_p^{\lambda _p}\mathrm{exp}\{h/f_0\}`$, where $`f_i`$ and $`h`$ are polynomials in $`[x,y]`$ and the $`\lambda _i`$, is called a Darboux function, see for instance . System (1) is called Darboux integrable if it has a first integral or an integrating factor which is a Darboux function (for a definition of first integral and of integrating factor, see ). The following lemma gives the relation between Darboux functions and invariant algebraic curves of a system (1).
###### Lemma 2
We consider a Darboux function $`𝒟(x,y):=f_1^{\lambda _1}\mathrm{}f_p^{\lambda _p}\mathrm{exp}\{h/f_0\}`$ such that
$$P(x,y)\frac{𝒟}{x}(x,y)+Q(x,y)\frac{𝒟}{y}(x,y)=k(x,y)𝒟(x,y),$$
where $`k(x,y)`$ is a polynomial of degree at most $`\mathrm{d}1`$. Then, each $`f_i(x,y)=0`$, $`i=1,2,\mathrm{},p`$ is an invariant algebraic curve of system (1) and $`\mathrm{exp}\{h/f_0\}`$ is an exponential factor of system (1).
The proof of this lemma is analogous to the proofs of Lemma 3 and Proposition 4 in .
We recall that $`V(x,y)`$ is an inverse integrating factor of system (1) if it is a function of class $`𝒞^1`$ in some open set $`𝒰`$ of $`^2`$ and satisfies the following partial differential equation:
$$P(x,y)\frac{V}{x}(x,y)+Q(x,y)\frac{V}{y}(x,y)=\left(\frac{P}{x}(x,y)+\frac{Q}{y}(x,y)\right)V(x,y).$$
We note that the function $`1/V(x,y)`$ is an integrating factor for system (1) in $`𝒰`$. The following result, which is a summary of some well known results, relates the existence of an inverse integrating factor in a certain functional class and the existence of a first integral in another (possibly larger) functional class. The definitions of elementary function and Liouvillian function can be found in . Proposition 3 shows that when considering the integrability problem we are also addressed to study whether an inverse integrating factor belongs to a certain given class of functions. As many authors have noted, see for instance , inverse integrating factors play a fundamental role in the integrability problem, not only because they characterize the functional class of a first integral but also because they usually belong to an easier functional class. For instance, quadratic systems of the form (1) with a center at the origin always have an inverse integrating factor which is a polynomial of degree $`3`$ or $`5`$. Hence, we have that the characterization of centers for quadratic systems can be done by means of a polynomial instead of a first integral, which will be of Darboux type in a general case.
###### Proposition 3
The following three statements hold.
* System (1) has a Darboux first integral if, and only if, it has a rational inverse integrating factor.
* If system (1) has an elementary first integral, then it has an inverse integrating factor of the form $`V(x,y)=\left(A(x,y)/B(x,y)\right)^{1/N}`$, where $`N`$ and $`A,B[x,y]`$.
* System (1) has a Liouvillian first integral if, and only if, it has a Darboux inverse integrating factor.
The first statement of this proposition is proved in and its reciprocal is proved in . Statement (b) is proved in and the last statement is proved in .
It is clear, as shown in Lemma $`2`$ of , that given an invariant algebraic curve $`f(x,y)=0`$ of system (1), all the algebraic functions defined by it in an implicit way, that is, all the functions $`g(x)`$ satisfying $`f(x,g(x))0`$, are particular solutions of equation (2). In this paper, among other results, the converse result is established, that is, given an algebraic particular solution $`y=g(x)`$ of equation (2), we show that the irreducible polynomial $`f(x,y)`$ such that $`f(x,g(x))0`$ gives rise to an invariant algebraic curve $`f(x,y)=0`$ of system (1). This fact is stated and proved in Theorem 6 of Section 3. We have noticed that a Darboux function may also contain exponential factors and this fact is necessary so as to characterize the Liouvillian integrability. Hence, exponential factors appear in a natural way when considering invariant algebraic curves. In this paper we consider algebraic particular solutions $`yg(x)`$ which come naturally from invariant algebraic curves and this relationship allows us to give an analogous to exponential factors but for algebraic particular solutions. This analogy is motivated and made clear in Subsection 3.2. In relation with this fact, we give some characteristics related to particular solutions such as the structure of their cofactor, which is given in Subsection 3.3.
In and an algorithmic method to determine, for system (1), the possible existence of first integrals or integrating factors of the form $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$ or $`I(x,y)=\mathrm{exp}\{h_2(x)_{k=1}^r(ya_k(x))\mathrm{}_{j=1}^s(y\stackrel{~}{g}_j(x))\}`$ $`h_1(x)`$ $`_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$, where $`g_i(x)`$ and $`\stackrel{~}{g}_j(x)`$ are unknown particular solutions of equation (2), $`\alpha _i`$ are unknown constants, and $`a_k(x)`$, $`h(x)`$, $`h_1(x)`$ and $`h_2(x)`$ are unknown functions, is given. In both cases, if all the particular solutions $`g_i(x)`$ and $`\stackrel{~}{g}_j(x)`$ are determined, which is expressed by the non-existence of a nonlinear superposition principle as described in , they are algebraic functions (see Proposition 7 in and Theorem 2 in ). The algorithm, in this case, gives an alternative method to determine such type of solutions. In the case where all the $`g_i(x)`$, $`a_k(x)`$ and $`\stackrel{~}{g}_j(x)`$ are algebraic functions, $`h(x)`$ and $`h_1(x)`$ have a rational logarithmic derivative and $`h_2(x)`$ is a rational function, we show (cf. Propositions 16 and 17, Subsection 3.4) that $`I(x,y)`$ is a Darboux function. This result is one of the main goals obtained in the present work. Hence, using the algorithm described in , all the systems with a Liouvillian first integral can be found, as well as an explicit expression of a non-Liouvillian first integral when there is a nonlinear superposition principle. The present work is born as a complement to the results described in . In these two works the integrability problem is studied in the particular case of the existence of a first integral of the described form $`I(x,y)`$. As we have already stated, in it is shown that if when applying the described algorithm the function $`I(x,y)`$ is completely determined then it only involves algebraic functions. In that case, in all the examples studied in , the function $`I(x,y)`$ was a Darboux function. In fact, this is the general case: if when applying the algorithm described in the function $`I(x,y)`$ is completely determined then it is a Darboux function. The proof of this assertion is one of the objectives that we have achieved in the present work.
## 2 Some preliminary results on algebraic functions and polynomials
In this section we give a summary of well known algebraic results which are needed in this paper. The definitions and proofs can be found, for instance, in the books .
We always consider polynomials in $`[x,y]`$ which is a ring with the usual addition and product operations of polynomials. Equivalently, we may consider the ring $`[x,y]`$. One of the most important equivalence relations which can be defined in the ring $`[x,y]`$ is the divisibility relation. We say that the polynomial $`f_1[x,y]`$ divides $`f_2[x,y]`$, and we write $`f_1|f_2`$, if there exists a polynomial $`k[x,y]`$ such that $`f_2=kf_1`$. It can be shown that $`[x,y]`$ is a unique factorization domain (UFD). We recall that the unit elements of the ring $`[x,y]`$ are the constants different from zero, i.e., $`\{0\}`$. We say that two elements $`f_1,f_2[x,y]`$ are associates if there exists a unit element $`e`$ such that $`f_1=ef_2`$. An irreducible polynomial is a non–constant element $`f(x,y)[x,y]`$ which is only divided by its associates.
If $`f_1,f_2[x,y]`$ are such that for any $`(a,b)\{(x,y)^2:f_1(x,y)=0\}`$ we have that $`f_2(a,b)=0`$, we will write $`f_2(x,y)|_{f_1(x,y)=0}=0`$. Given a polynomial $`f[x,y]`$, it defines a curve, denoted by $`f(x,y)=0`$, which is the set $`\{(x,y)^2:f(x,y)=0\}`$. We say that $`p^2`$ is a point of intersection of two curves $`f_1(x,y)=0`$ and $`f_2(x,y)=0`$ if $`f_1(p)=f_2(p)=0`$. An intersection point must be always counted as times as its multiplicity. The multiplicity of a point of intersection is a rather complicated notion, which can be found, for instance, in page $`60`$ of . Intuitively, the multiplicity of a point of intersection takes into account the number of tangents shared by the two curves at that point.
Bézout’s theorem takes into account the degree of two curves and their points of intersection and relates them with the fact of having a common factor.
###### Theorem 4
\[Bézout\] If $`f_1(x,y)=0`$ and $`f_2(x,y)=0`$ are two curves of degrees $`m`$ and $`n`$, respectively, with more than $`mn`$ intersection points, then there is an irreducible polynomial $`r(x,y)`$ which divides both $`f_1(x,y)`$ and $`f_2(x,y)`$.
We are going to recall some results on fractionary power series that can be found in . If $`x`$ is any free variable over $``$, we denote by $`((x))`$ the field of fractions of $`[[x]]`$, where $`[[x]]`$ is the ring of entire formal power series in $`x`$ with complex coefficients. We recall that given a ring, which needs to be an integral domain, we can define its field of fractions as the smallest field containing it. The field of fractions is therefore obtained from the integral domain by adding the least needed to make of it a field, that is, the possibility of dividing by any nonzero element. Given $`n`$, we consider entire fractionary series $`s=_{ir}a_ix^{i/n}`$ where $`r`$, $`a_i`$ and $`\mathrm{min}\{i:a_i0\}0`$. An element $`s((x^{1/n}))`$ is of the form $`s=_{ir}a_ix^{i/n}`$ where $`r`$, $`a_i`$ and the $`\mathrm{min}\{i:a_i0\}`$ can be lower than zero. The elements of the ring $`[[x^{1/n}]]`$ are the entire fractionary series such that $`\mathrm{min}\{i:a_i0\}0`$. It can be shown (see pages 17 and 18 of ), that given a fractionary series $`s=_{ir}a_ix^{i/n}((x^{1/n}))`$ we can always take an equivalent series such that $`n`$ and $`\mathrm{gcd}\{i:a_i0\}`$ have no common factor. In this case, we say that $`n`$ is the polydromy order of the fractionary series $`s`$ and we denote it by $`\nu (s)`$. Let $`M(x,y)`$ be a polynomial in $`y`$ of degree $`N`$ whose coefficients in $`y`$ are fractionary series in $`((x^{1/n_i}))`$, that is, we expand $`M(x,y)`$ in powers of $`y`$: $`M(x,y)=_{i=0}^Ns_i(x)y^i`$ and we have that $`s_i(x)((x^{1/n_i}))`$, for $`i=0,1,2,\mathrm{},N`$. The polydromy order of $`M(x,y)`$ is defined as the least common multiple of the polydromy orders of $`s_i(x)`$, $`i=0,1,2,\mathrm{},N`$, that is, $`\nu (M)=\text{lcm}\{\nu (s_i):i=0,1,2,\mathrm{},N\}`$.
The fractionary series $`s`$ is said to be convergent if $`_{ir}a_it^i`$ has non-zero convergence radius, where $`t=x^{1/n}`$.
The ring $`\{x,y\}`$ is the ring of convergent power series in two variables and complex coefficients. The following result clarifies the structure of an algebraic function and it is stated and proved in (page $`26`$).
###### Theorem 5
If $`f(x,y)[x,y]`$, then there is a unique decomposition of the form $`f=ux^r_{i=1}^{\mathrm{}}(yg_i(x)),`$ where $`r\{0\}`$, $`u`$ is an invertible power series of $`\{x,y\}`$ and $`g_i(x)`$ are convergent fractionary series.
In the book this result is stated for formal power series $`f(x,y)[[x,y]]`$. In the Corollary 1.5.6 of page $`26`$ in , it is stated that $`\mathrm{}`$ is the height of the Newton polygon associated to $`f(x,y)`$. Since we are considering a polynomial $`f(x,y)`$ instead of a formal series, we deduce that $`\mathrm{}`$ is the highest degree in $`y`$ of $`f(x,y)`$ and, therefore, $`u`$ is an invertible power series of $`\{x\}`$. We define $`h(x)=ux^r`$ and we have that given $`f(x,y)[x,y]`$ of degree $`\mathrm{}`$ in $`y`$, there is a unique decomposition of the form:
$$f=h(x)\underset{i=1}{\overset{\mathrm{}}{}}(yg_i(x)),$$
(5)
where $`h(x)\{x\}`$ and $`g_i(x)`$ are fractionary series. Theorem 1.7.2 (page $`31`$) of ensures that all the $`y`$–roots of $`f(x,y)`$ are convergent. Hence, we deduce that an algebraic function $`g(x)`$ is a convergent fractionary series.
By definition, given an algebraic function $`g(x)`$, we have that there exists a polynomial such that $`f(x,g(x))0`$. If $`f(x,y)`$ is not irreducible, then $`g(x)`$ is a fractionary series appearing in the decomposition (5) of at least one of the irreducible factors of $`f(x,y)`$. Hence, without loss of generality, we may always assume that $`f(x,y)`$ is irreducible. Moreover, given $`g(x)`$, this irreducible polynomial $`f(x,y)`$ is unique (modulus multiplication by constants). This statement is clear from the fact that if $`f_1(x,y)`$ and $`f_2(x,y)`$ are two irreducible polynomials such that $`f_i(x,g(x))0`$, $`i=1,2`$, then these two polynomials have an infinite number of points of intersection (because $`g(x)`$ is a convergent fractionary series) and, by Bézout’s Theorem 4, we have that $`f_1(x,y)`$ and $`f_2(x,y)`$ must be associates.
## 3 The Main Results
### 3.1 Particular algebraic solutions
###### Theorem 6
Let $`g(x)`$ be an algebraic particular solution of equation (2) and we call $`f(x,y)`$ the irreducible polynomial satisfying $`f(x,g(x))0`$. Then, the curve $`f(x,y)=0`$ is an invariant algebraic curve of system (1).
Proof. Let us denote by $`F(x,y)`$ the polynomial in $`[x,y]`$ defined by:
$$F(x,y):=P(x,y)\frac{f}{x}(x,y)+Q(x,y)\frac{f}{y}(x,y).$$
We have that $`F(x,y)=0`$ and $`f(x,y)=0`$ intersect in all the points of the form $`(x,g(x))`$ by virtue of (3). By Bézout’s Theorem 4, we deduce that the polynomials $`f(x,y)`$ and $`F(x,y)`$ share a common factor, because they intersect in an infinite (continuum) number of points. They intersect in all the points $`(x,g(x))`$ where the fractionary power series $`g(x)`$ is convergent. Since $`f(x,y)`$ is an irreducible polynomial, we have that $`f(x,y)`$ divides $`F(x,y)`$ in the ring of real polynomials. From this fact, we conclude that there exists a polynomial $`k(x,y)`$ such that $`F(x,y)=k(x,y)f(x,y)`$ and we get that $`f(x,y)=0`$ is an invariant algebraic curve of system (1).
We recall the definition of invariant and of quasipolynomial cofactor stated in .
###### Definition 7
An invariant of (1) is a function $`\varphi (x,y)`$ such that there exists a quasipolynomial cofactor $`M(x,y)`$, where $`M(x,y)`$ is a polynomial in one of the variables $`x`$ or $`y`$ of degree $`m1`$ with $`m`$ the degree of system (1) in that variable, satisfying $`P(\varphi /x)+Q(\varphi /y)=M\varphi `$.
In case that the set of points in $`^2`$ satisfying that $`\varphi (x,y)=0`$ is not null, we have that $`\varphi `$ is an invariant if, and only if, $`P(\varphi /x)+Q(\varphi /y)_{|\varphi =0}=0`$. We say that $`\varphi =0`$ is an invariant curve in this case.
These definitions are a generalization of the so called generalized cofactor introduced in where a generalization of the Darboux integrability theory in order to find non-Liouvillian first integrals of system (1) was presented. For the special invariant curve $`\varphi (x,y):=yg(x)=0`$ of (1), where $`g(x)`$ is a particular solution of equation (2), a quasipolynomial cofactor always exists as it was established in . Many examples of invariants with a quasipolynomial cofactor are given in as well as a method to find first integrals, which are non-Liouvillian in general, for certain families of systems.
###### Proposition 8
A particular solution $`g(x)`$ of equation (2) always has a unique associated quasipolynomial cofactor of the form $`M(x,y)=k_{m1}(x)y^{m1}+\mathrm{}+k_1(x)y+k_0(x)`$, where $`m`$ is the degree in $`y`$ of system (1).
In case that $`g(x)`$ is algebraic then each $`k_i(x)`$, $`i=0,1,2,\mathrm{},m1`$ is a rational function in $`x`$ and $`g(x)`$ with coefficients in $``$.
The second part of this proposition is deduced from the proof of the first part as given in .
In case $`g(x)`$ is an algebraic function, by Theorem 5, we have that it is a fractionary series. Since $`k_i(x)`$ is a rational function in $`x`$ and $`g(x)`$ with coefficients in $``$ and $`g(x)((x^{1/n}))`$, for a certain natural $`n`$, then its quasipolynomial cofactor $`M(x,y)((x^{1/n}))[y]`$, that is, each of the functions $`k_i(x)`$, $`i=0,1,2,\mathrm{},m1`$ is a fractionary series with a polydromy order that divides the polydromy order of $`g(x)`$. Next lemma shows that in case $`M(x,y)`$ is a polynomial, then $`g(x)`$ must be a rational function.
We say that an equation (2) is linear if it is of the form
$$\frac{dy}{dx}=m_1(x)y+m_0(x),$$
where $`m_0(x)`$ and $`m_1(x)`$ are functions of $`x`$. If equation (2) does not take this form, we say that it is non-linear.
###### Lemma 9
Let $`g(x)`$ be a particular solution of a rational non-linear equation (2). If the quasipolynomial cofactor $`M(x,y)`$ related to $`g(x)`$ is a polynomial, then $`g(x)`$ is a rational function.
Proof. We have that $`M(x,y)=k_{m1}(x)y^{m1}+\mathrm{}+k_1(x)y+k_0(x)`$ is a polynomial, so, $`k_j(x)`$ is a polynomial in $`x`$ for all $`j=0,1,2,\mathrm{},m1`$. Let us expand the polynomials $`P(x,y)`$ and $`Q(x,y)`$ in powers of $`y`$: $`P(x,y)=p_0(x)+p_1(x)y+\mathrm{}+p_m(x)y^m`$, $`Q(x,y)=q_0(x)+q_1(x)y+\mathrm{}+q_m(x)y^m`$, where $`p_j(x)`$ and $`q_j(x)`$ are the polynomials in $`x`$ corresponding to the coefficients of degree $`j`$ in $`y`$ of $`P(x,y)`$ and $`Q(x,y)`$, respectively. Since $`g(x)`$ is a particular solution of equation (2), we have that $`P(x,y)g^{}(x)+Q(x,y)=M(x,y)(yg(x))`$. Equating the coefficients of order $`j`$ in $`y`$, we deduce that:
$$p_0(x)g^{}(x)k_0(x)g(x)=q_0(x),(eq_0)$$
$$p_j(x)g^{}(x)k_j(x)g(x)=q_j(x)k_{j1}(x),(eq_j)j=1,2,\mathrm{},m1,$$
$$p_m(x)g^{}(x)=q_m(x)k_{m1}(x)(eq_m).$$
If $`k_0(x)p_j(x)p_0(x)k_j(x)0`$ for some $`j=1,2,\mathrm{},m1`$, we can equate $`g(x)`$ from the equations $`(eq_0)`$ and $`(eq_j)`$ and we get
$$g(x)=\frac{p_0(x)q_j(x)q_0(x)p_j(x)p_0(x)k_{j1}(x)}{k_0(x)p_j(x)p_0(x)k_j(x)},$$
which is a rational function. If $`k_0(x)p_j(x)p_0(x)k_j(x)0`$ for all $`j=1,2,\mathrm{},m1`$, we deduce that $`p_j(x)=k_j(x)L_1(x)`$ and $`p_0(x)=k_0(x)L_1(x)`$ for all $`j=1,2,\mathrm{},m1`$, where $`L_1(x)`$ is a rational function in $`x`$.
If $`p_m(x)k_0(x)0`$, from the first and the last equations we deduce that
$$g(x)=\frac{p_0(x)(q_m(x)k_{m1}(x))p_m(x)q_0(x)}{p_m(x)k_0(x)},$$
which is a rational function.
We can also try to equate $`g(x)`$ from the equations $`(eq_m)`$ and $`(eq_j)`$ for some $`j=1,2,\mathrm{},m1`$. If $`p_m(x)k_j(x)0`$, then
$$g(x)=\frac{q_m(x)p_j(x)p_m(x)q_j(x)+p_m(x)k_{j1}(x)k_{m1}(x)p_j(x)}{p_m(x)k_j(x)},$$
which is a rational function.
In case that $`k_0(x)p_j(x)p_0(x)k_j(x)0`$, $`p_m(x)k_0(x)0`$ and $`p_m(x)k_j(x)0`$ for all $`j=1,2,\mathrm{},m1`$, first assume that $`M(x,y)0`$. Then $`yg(x)`$ would be a first integral for system (1), which means that $`Q(x,y)P(x,y)g^{}(x)0`$. Hence, $`Q(x,y)`$ and $`P(x,y)`$ are such that equation (2) is a linear one, in contradiction with our hypothesis. So, we conclude that $`p_m(x)0`$. We also have that $`P(x,y)=L_1(x)M(x,y)`$, that is, $`p_j(x)=k_j(x)L_1(x)`$ for $`j=0,1,2,\mathrm{},m1`$. From equation $`(eq_m)`$ we have that $`q_m(x)=k_{m1}(x)`$ and the other equations read for:
$$k_0(x)(L_1(x)g^{}(x)g(x))=q_0(x),(eq_0^{}),$$
$$k_j(x)(L_1(x)g^{}(x)g(x))=q_j(x)k_{j1}(x),(eq_j^{}),$$
for $`j=1,2,\mathrm{},m1`$. Equating the factor $`(L_1(x)g^{}(x)g(x))`$ from equation $`(eq_0^{})`$ and $`(eq_j^{})`$ we deduce that $`k_j(x)q_0(x)k_0(x)q_j(x)+k_0(x)k_{j1}(x)0`$ for all $`j=1,2,\mathrm{},m1`$. We write $`q_0(x)=L_0(x)k_0(x)`$, where $`L_0(x)`$ is a rational function, and we have that $`k_0(x)(k_j(x)L_0(x)q_j(x)+k_{j1}(x))0`$ for all $`j=1,2,\mathrm{},m1`$. If $`k_0(x)0`$, then $`P(x,y)`$ and $`Q(x,y)`$ share the common factor $`y`$. If $`k_0(x)0`$, then we have that $`q_j(x)=k_j(x)L_0(x)+k_{j1}(x)`$, from which we deduce that $`Q(x,y)=(y+L_0(x))M(x,y)`$.
Unless $`M(x,y)`$ is a real number (different from zero), we have that $`P(x,y)`$ and $`Q(x,y)`$ share a common factor, in contradiction with our hypothesis. Therefore, after a rescaling of the time, the only systems of the form (1) with a particular non rational solution $`g(x)`$ with a polynomial cofactor (in fact, the cofactor is the real number $`1`$) are of the form : $`\dot{x}=L_1(x)`$, $`\dot{y}=y+L_0(x)`$, which give rise to a linear equation.
In the same way as in Lemma 9, we have that if $`g(x)`$ is an algebraic particular solution of a nonlinear equation (2) with polydromy order $`\nu (g)`$, then its quasipolynomial cofactor must have the same polydromy order.
###### Lemma 10
Let $`g(x)`$ be an algebraic particular solution of a non–linear equation (2). Its quasipolynomial cofactor $`M(x,y)`$ has the same polydromy order as $`g(x)`$.
Proof. We have that if $`g(x)`$ has polydromy order $`\nu (g)`$, then the polydromy order of its quasipolynomial cofactor $`M(x,y)`$, $`\nu (M)`$, divides it, that is, $`\nu (M)|\nu (g)`$. This is because $`M(x,y)`$ is a polynomial in $`y`$ and a rational function in $`x`$ and $`g(x)`$ as it has been stated in Proposition 8.
The fact that both polydromy orders coincide is a corollary of the previous proof of Lemma 9. The reasonings are the same since we can equate $`g(x)`$ in terms of the coefficients of $`M(x,y)`$ unless the equation is of linear type.
###### Remark 11
If we have a linear equation of the form $`dy/dx=m_1(x)y+m_0(x)`$, we may have a nonrational particular solution with a polynomial cofactor. For instance, taking $`m_0(x)0`$ and $`m_1(x)1`$, we have that $`g(x):=e^x`$ is a particular solution with cofactor $`1`$. In the same way, the linear equation $`dy/dx=\mathrm{\hspace{0.17em}3}y/(2x)`$ has the algebraic solution $`y=x^{3/2}`$, whose polydromy order is $`2`$, with the polynomial cofactor $`3`$, whose polydromy order is $`1`$.
### 3.2 On the invariants giving rise to exponential factors
In the Darboux theory of integrability not only invariant algebraic curves are considered, but also exponential factors, as we have stated in the introduction. Exponential factors appear in a natural way from the coalescence of invariant algebraic curves, as it is explained in . This statement means that if we have a polynomial system (1) with an exponential factor of the form $`\mathrm{exp}\{h/f_0\}`$, then there is a $`1`$–parameter perturbation of system (1), given by a small $`\epsilon `$, with two invariant algebraic curves, namely $`f_0=0`$ and $`f_0+\epsilon h=0`$. Hence, when $`\epsilon =0`$, these two curves coalesce giving the exponential factor $`\mathrm{exp}\{h/f_0\}`$ for the system with $`\epsilon =0`$, as well as the invariant algebraic curve $`f_0=0`$ which does not disappear.
In this context, and taking algebraic particular solutions into account, the following question arises: which kind of function appears with the coalescence of two algebraic particular solutions? In the next Proposition 12, we show that the natural generalization of algebraic particular solutions in this framework is a function of the form: $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$ where $`A_1(x,y)=_{k=1}^r(ya_k(x))`$ and $`A_0(x,y)=_{j=1}^s(y\stackrel{~}{g}_j(x))`$, with $`a_k(x)`$ and $`\stackrel{~}{g}_j(x)`$ algebraic functions, $`k=1,2,\mathrm{},r`$ and $`j=1,2,\mathrm{},s`$, $`A_1(x,y)`$ and $`A_0(x,y)`$ do not share any common factor, and $`h_2(x)`$ is a rational function in $`x`$. For convention, if $`r=0`$ or $`s=0`$, we mean that $`A_1(x,y)`$ or $`A_0(x,y)`$ takes a constant value, respectively, which we may assume to be equal to $`1`$.
Since by Proposition 8, we always have that a particular solution given by $`\varphi (x,y)=yg(x)`$ has a quasipolynomial cofactor, the property of being an invariant for $`\mathrm{\Phi }(x,y)`$ is given by associating to it a quasipolynomial cofactor $`M(x,y)`$ which is a polynomial in $`y`$ of degree at most $`m1`$, where $`m`$ is the degree in $`y`$ of system (1). Moreover, we will assume that $`M(x,y)`$ is well defined over $`A_0(x,y)=0`$, that is, $`M(x,\stackrel{~}{g}_j(x))`$ is a real function of $`x`$ for all $`j=1,2,\mathrm{},s`$. The following proposition gives the characterization of this fact.
###### Proposition 12
We consider a function $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$ where $`A_1(x,y)=_{k=1}^r(ya_k(x))`$ and $`A_0(x,y)=_{j=1}^s(y\stackrel{~}{g}_j(x))`$, with $`a_k(x)`$ and $`\stackrel{~}{g}_j(x)`$ algebraic functions, $`k=1,2,\mathrm{},r`$ and $`j=1,2,\mathrm{},s`$ $`A_1(x,y)`$ and $`A_0(x,y)`$ do not share any common factor, and $`h_2(x)`$ is a rational function in $`x`$. Assume that there exists a function $`M(x,y)`$, which is a polynomial in $`y`$ of degree at most $`m1`$, where $`m`$ is the degree in $`y`$ of system (1), such that:
$$P(x,y)\frac{\mathrm{\Phi }(x,y)}{x}+Q(x,y)\frac{\mathrm{\Phi }(x,y)}{y}=M(x,y)\mathrm{\Phi }(x,y).$$
(6)
We assume that $`M(x,\stackrel{~}{g}_j(x))`$ is a real function of $`x`$ for all $`j=1,2,\mathrm{},s`$. We denote by $`\stackrel{~}{k}_i(x)`$ the coefficient of degree $`i`$ in $`y`$ of $`M(x,y)`$, that is, $`M(x,y)=\stackrel{~}{k}_0(x)+\stackrel{~}{k}_1(x)y+\stackrel{~}{k}_2(x)y^2+\mathrm{}+\stackrel{~}{k}_{m1}(x)y^{m1}`$. Then,
* Each one of the algebraic functions $`\stackrel{~}{g}_j(x)`$, $`j=1,2,\mathrm{},s`$ is a particular solution of equation (2). We denote by $`M_j(x,y)`$ its associated quasipolynomial cofactor.
* The following identity is satisfied:
$$𝒳\left(h_2(x)A_1(x,y)\right)=\left(\underset{j=1}{\overset{s}{}}M_j(x,y)\right)h_2(x)A_1(x,y)+M(x,y)A_0(x,y).$$
(7)
* Each one of the functions $`\stackrel{~}{k}_i(x)`$ is rational in $`x`$ and rational in $`a_k(x)`$ and $`\stackrel{~}{g}_j(x)`$, with $`k=1,2,\mathrm{},r`$ and $`j=1,2,\mathrm{},s`$.
Proof. We have that $`𝒳\left(\mathrm{\Phi }(x,y)\right)=M(x,y)\mathrm{\Phi }(x,y)`$, from which we deduce the following identity:
$$𝒳\left(h_2(x)A_1(x,y)\right)A_0(x,y)h_2(x)A_1(x,y)𝒳\left(A_0(x,y)\right)=M(x,y)A_0(x,y)^2.$$
(8)
Since $`A_1(x,y)`$ and $`A_0(x,y)`$ do not share any common factor and $`A_0(x,\stackrel{~}{g}_j(x))0`$ for $`j=1,2,\mathrm{},s`$, then, from equation (8), we have that $`𝒳\left(A_0(x,y)\right)|_{y=\stackrel{~}{g}_j(x)}0`$. We notice that here we are using that $`M(x,\stackrel{~}{g}_j(x))`$ is a real function of $`x`$ for all $`j=1,2,\mathrm{},s`$. Let us call $`A_{0_j}(x,y):=_{i=1,ij}^s(yg_i(x))`$ and we have that $`A_0(x,y)=A_{0_j}(x,y)(y\stackrel{~}{g}_j(x))`$. Then,
$$𝒳\left(A_0(x,y)\right)=𝒳\left(A_{0_j}(x,y)\right)(y\stackrel{~}{g}_j(x))+A_{0_j}(x,y)\left(\stackrel{~}{g}_j^{}(x)P(x,y)+Q(x,y)\right).$$
From $`𝒳\left(A_0(x,y)\right)|_{y=\stackrel{~}{g}_j(x)}0`$, we deduce that $`\stackrel{~}{g}_j^{}(x)P(x,\stackrel{~}{g}_j(x))+Q(x,\stackrel{~}{g}_j(x))0`$ and we conclude that $`\stackrel{~}{g}_j(x)`$ is a particular solution of system (1), as stated in (i).
In order to prove (ii), we consider the quasipolynomial cofactor $`M_j(x,y)`$ associated to $`y\stackrel{~}{g}_j(x)`$, whose existence is ensured by Proposition 8. Then, we have that $`𝒳(A_0)=\left(_{j=1}^sM_j\right)A_0`$. Hence, identity (8) reads for:
$$\begin{array}{c}𝒳\left(h_2(x)A_1(x,y)\right)A_0(x,y)h_2(x)A_1(x,y)\left(\underset{j=1}{\overset{s}{}}M_j(x,y)\right)A_0(x,y)=\hfill \\ =M(x,y)A_0(x,y)^2.\hfill \end{array}$$
This identity coincides with (7), after dividing both members by $`A_0(x,y)`$.
Finally, to prove (iii), we observe that equating the coefficients of the same degree in $`y`$ from identity (7), we deduce that each $`\stackrel{~}{k}_i(x)`$ is a rational function of $`x`$, $`\stackrel{~}{g}_j(x)`$, $`a_k(x)`$ and $`a_k^{}(x)`$. Since each $`a_k(x)`$ is an algebraic function, we can consider $`f_{a_k}(x,y)`$ as the irreducible polynomial in $`[x,y]`$ such that $`f_{a_k}(x,a_k(x))0`$. We can derive this last identity with respect to $`x`$ and we deduce that:
$$\frac{f_{a_k}(x,a_k(x))}{x}+\frac{f_{a_k}(x,a_k(x))}{y}a_k^{}(x)0.$$
Hence, we can substitute the value $`a_k^{}(x)`$ appearing in $`\stackrel{~}{k}_i(x)`$ by the rational expression in $`x`$ and $`a_k(x)`$: $`\left[f_{a_k}(x,a_k(x))/x\right]/\left[f_{a_k}(x,a_k(x))/y\right]`$. Therefore, each $`\stackrel{~}{k}_i(x)`$ is a rational function of $`x`$, $`\stackrel{~}{g}_j(x)`$ and $`a_k(x)`$, as we wanted to prove.
As Theorem 6 shows, an algebraic particular solution $`yg(x)`$ recovers the invariant algebraic curve to which it is related. That is, we have that $`f(x,y)=0`$ is an invariant algebraic curve of system (1) if, and only if, all its $`y`$–roots are algebraic particular solutions of equation (2). The fact of being particular solutions implies that each one of them has an associated quasipolynomial cofactor. We notice that all the $`y`$–roots appear in the factorization $`f(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))`$ described by (6). We would like to have an analogous to this statement but for exponential factors and this is what is given in the next Theorem 13.
In this context, if we have an exponential factor given by $`\mathrm{exp}\{R(x,y)\}`$, where $`R(x,y)`$ is a rational function, i.e. $`R(x,y)(x,y)`$, we would like to write a factorization for it analogous to (6). We notice that given any two functions $`R_1(x,y)`$, $`R_2(x,y)`$ with the property that $`R_1(x,y)+R_2(x,y)R(x,y)`$, we have that $`\mathrm{exp}\{R(x,y)\}\mathrm{exp}\{R_1(x,y)\}\mathrm{exp}\{R_2(x,y)\}`$, which is a “factorization” of $`\mathrm{exp}\{R(x,y)\}`$. However, not all these factorizations are useful to our purposes because it can be shown that $`\mathrm{exp}\{R_1(x,y)\}`$ and $`\mathrm{exp}\{R_2(x,y)\}`$ do not need to have an associated quasipolynomial cofactor. So, they are not invariants for system (1), as the exponential factor $`\mathrm{exp}\{R(x,y)\}`$ is. The following example exhibits this fact.
Example. Let us consider the following cubic system with the invariant straight line $`y=0`$ :
$$\dot{x}=(2x+y)(1+x)+2x^2y+y^3,\dot{y}=y(1+x+xy).$$
(9)
This system has the invariant $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{\sqrt{x}/y\}`$ with the quasipolynomial cofactor $`M(x,y)=(1+x+y^2)/(2\sqrt{x})`$. As it will be proved in Theorem 13, the existence of this invariant implies the existence of the following exponential factor: $`F(x,y):=\mathrm{exp}\{(x+y)/y^2\}`$ with the cofactor $`\stackrel{~}{k}_0(x,y)=yx`$. We notice that not all the functions $`R_1,R_2`$ satisfying $`R_1+R_2=(x+y)/y^2`$ give rise to an invariant. For instance, if we take $`R_1:=x/y^2`$ and $`R_2:=1/y`$, the following easy computation shows that $`\mathrm{exp}\{R_1\}`$ is not an invariant of system (9), and neither $`\mathrm{exp}\{R_2\}`$ is. We note that, taking $`𝒳`$ as the vector field associated to system (9), we have that $`𝒳(\mathrm{exp}\{R_1\})=(y+(1+x)/y)\mathrm{exp}\{R_1\}`$ and since the rational function $`(y+(1+x)/y)`$ is not well-defined over $`y=0`$, we deduce that $`\mathrm{exp}\{R_1\}`$ is not an invariant. However, if we consider $`\mathrm{\Psi }(x,y):=\mathrm{exp}\{(x+yy\sqrt{x})/y^2\}`$, we have that $`\mathrm{\Phi }(x,y)\mathrm{\Psi }(x,y)=F(x,y)`$ and both $`\mathrm{\Phi }(x,y)`$ and $`\mathrm{\Psi }(x,y)`$ are invariants of system (9).
The following theorem shows that if $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$, as described in Proposition 12, is such that it has an associated quasipolynomial cofactor, then there exists another function $`\mathrm{\Psi }(x,y):=\mathrm{exp}\{B_1(x,y)/B_0(x,y)\}`$, of the same type as $`\mathrm{\Phi }(x,y)`$, such that the product $`\mathrm{\Phi }(x,y)\mathrm{\Psi }(x,y)`$ gives rise to an exponential factor $`\mathrm{exp}\{R(x,y)\}`$ for system (1). This sentence means that a factorization of $`\mathrm{exp}\{R(x,y)\}`$ giving invariants for system (1) is obtained by the product $`\mathrm{\Phi }(x,y)\mathrm{\Psi }(x,y)`$. Therefore, we are able to recover $`\mathrm{exp}\{R(x,y)\}`$ from one of its factors $`\mathrm{\Phi }(x,y)`$ and, moreover, since $`\mathrm{exp}\{R(x,y)\}`$ appears by the coalescence of two invariant algebraic curves, we conclude that $`\mathrm{\Phi }(x,y)`$ needs to appear by the coalescence of algebraic particular solutions since it is formed by a product of them. This fact is the result that we were targeting to: we wanted to exhibit the analogy between the generalization of invariant algebraic curves to exponential factors with the generalization of algebraic particular solutions to invariants of the form $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$.
###### Theorem 13
Assume that the function $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$ has a quasipolynomial cofactor $`M(x,y)`$, that is, $`𝒳(\mathrm{\Phi }(x,y))=M(x,y)\mathrm{\Phi }(x,y)`$ and assume that $`M(x,y)`$ is well defined over $`\{(x,y):A_0(x,y)=0\}`$. Then, there exist quasipolynomial functions $`B_0(x,y)`$ and $`B_1(x,y)`$, which are polynomials in $`y`$ and algebraic in $`x`$, such that $`R:=h_2A_1/A_0+B_1/B_0`$ is a rational function in $`x`$ and $`y`$ and $`\mathrm{exp}\{R(x,y)\}`$ is an exponential factor of system (1).
Proof. We consider $`\varphi _A(x,y,\epsilon ):=A_0(x,y)+\epsilon h_2(x)A_1(x,y)`$ which is a polynomial in $`y`$ whose coefficients are algebraic functions. We can compute the sequence of powers $`\varphi _A(x,y,\epsilon )^{j+1}`$, for each natural number $`j`$ and the coefficients of all these polynomials in $`y`$ are algebraic functions, combination of those of $`\varphi _A(x,y,\epsilon )`$. Therefore, there exists a natural number $`N`$ such that $`\varphi _A(x,y,\epsilon )^{N+1}`$ is a linear combination of all the previous powers and the coefficients of this combination are powers of $`x`$, $`y`$ and $`\epsilon `$. This fact is due to the finiteness of the algebraic extensions given by the $`y`$–roots of $`\varphi _A(x,y,\epsilon )`$. In this way, we have that there exists an irreducible polynomial $`𝒫(x,y,\epsilon )`$ in $`x`$, $`y`$ and $`\epsilon `$, such that each of the $`y`$–roots of $`\varphi _A(x,y,\epsilon )`$ is an $`y`$–root of $`𝒫(x,y,\epsilon )`$. This polynomial $`𝒫(x,y,\epsilon )`$ is the minimal polynomial of the $`y`$-roots of $`\varphi _A(x,y,\epsilon )`$, and it can be always taken irreducible. We expand $`𝒫(x,y,\epsilon )`$ in powers of $`\epsilon `$ and we denote by $`R_i(x,y)`$ its coefficient of degree $`i`$ in $`\epsilon `$, which is a polynomial in $`x`$ and $`y`$. Let us consider the quotient $`𝒫(x,y,\epsilon )/\varphi _A(x,y,\epsilon )`$ which is a polynomial in $`y`$, denoted by $`\varphi _B(x,y,\epsilon )`$. We expand $`\varphi _B(x,y,\epsilon )`$ in powers of $`\epsilon `$ and we denote by $`B_i(x,y)`$ its coefficient of degree $`i`$ in $`\epsilon `$, which is a polynomial in $`y`$. Since $`𝒫=\varphi _A\varphi _B`$, we deduce that $`R_0=A_0B_0`$ and $`R_1=A_0B_1+h_2A_1B_0`$, equating the coefficients of $`\epsilon ^0`$ and $`\epsilon ^1`$. Thus, $`R:=h_2A_1/A_0+B_1/B_0=R_1/R_0`$ is a rational function in $`x`$ and $`y`$. In case $`R_10`$ all the reasoning works just taking the first $`R_i0`$ with $`i>0`$.
We only need to see that $`\mathrm{exp}\{R(x,y)\}`$ is an exponential factor of system (1). We have, from Proposition 12, that $`A_0`$ is the product of algebraic particular solutions $`y\stackrel{~}{g}_j(x)`$. Hence $`𝒳\left(A_0(x,y)\right)=M_0(x,y)A_0(x,y)`$, where $`M_0(x,y)`$ is the quasipolynomial cofactor $`_{j=1}^sM_j(x,y)`$, which is a polynomial in $`y`$ of degree at most $`m1`$, $`m`$ being the degree in $`y`$ of system (1). By the irreducibility of $`𝒫(x,y,ϵ)`$, we deduce that $`R_0(x,y)`$ is a power of the lowest degree polynomial containing as $`y`$–roots all the $`\stackrel{~}{g}_j(x)`$. Hence, by Theorem 6, we have that $`R_0(x,y)=0`$ is an invariant algebraic curve of system (1), that is, $`𝒳(R_0)=k_0R_0`$, where $`k_0(x,y)`$ is a polynomial in $`x`$ and $`y`$. Moreover, from $`R_0=A_0B_0`$ and $`𝒳\left(A_0\right)=M_0A_0`$ we deduce that $`𝒳\left(B_0\right)=(k_0M_0)B_0`$. By Lemma 1, we only need to show that there is a polynomial $`\stackrel{~}{k}_0(x,y)`$ such that $`𝒳\left(R_1\right)=k_0R_1+\stackrel{~}{k}_0R_0`$. We recall that from Proposition 12 we have that $`𝒳\left(h_2A_1\right)=M_0h_2A_1+MA_0`$. We consider the polynomial $`G(x,y):=𝒳\left(R_1(x,y)\right)k_0(x,y)R_1(x,y)`$. We have:
$`G`$ $`=`$ $`𝒳\left(R_1\right)k_0R_1`$
$`=`$ $`𝒳\left(A_0\right)B_1+A_0𝒳\left(B_1\right)+𝒳\left(h_2A_1\right)B_0+h_2A_1𝒳\left(B_0\right)k_0\left(A_0B_1+h_2A_1B_0\right)`$
$`=`$ $`M_0A_0B_1+A_0𝒳\left(B_1\right)+(M_0h_2A_1+MA_0)B_0+h_2A_1(k_0M_0)B_0`$
$`k_0(A_0B_1+h_2A_1B_0)`$
$`=`$ $`\left[M_0B_1+𝒳\left(B_1\right)+MB_0k_0B_1\right]A_0.`$
We deduce that $`G(x,\stackrel{~}{g}_j(x))0`$ for all the $`\stackrel{~}{g}_j(x)`$ appearing in $`A_0`$. We note that here we are using the hypothesis that $`M(x,y)`$ is well defined over $`\{(x,y):A_0(x,y)=0\}`$. Since $`R_0(x,y)`$ is the lowest degree polynomial with this property and $`G(x,y)`$ is a polynomial, we have, by Bézout’s Theorem, that $`R_0(x,y)`$ is a divisor of $`G(x,y)`$ in the ring of polynomials $`[x,y]`$. So, there exists a polynomial $`\stackrel{~}{k}_0(x,y)`$ such that $`G(x,y)=\stackrel{~}{k}_0(x,y)R_0(x,y)`$ and, thus, $`𝒳\left(R_1\right)=k_0R_1+\stackrel{~}{k}_0R_0`$. Therefore, $`\mathrm{exp}\{R_1/R_0\}`$ is an exponential factor of system (1).
### 3.3 On the structure of the quasipolynomial cofactor
The following proposition gives the form of the quasipolynomial cofactor associated to an invariant of the form $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$.
We define as $`m`$ the degree of system (1) in the variable $`y`$ and we expand the polynomials $`P(x,y)`$ and $`Q(x,y)`$ in this variable:
$$P(x,y)=\underset{i=1}{\overset{m}{}}p_i(x)y^i,Q(x,y)=\underset{i=1}{\overset{m}{}}q_i(x)y^i.$$
###### Proposition 14
Let $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$ be an invariant of system (1) with an associated quasipolynomial cofactor $`k(x,y):=k_0(x)+k_1(x)y+\mathrm{}+k_{m1}(x)y^{m1}`$. Then, $`p_m(x)h^{}(x)0`$ and
$$k_j(x)=p_j(x)\frac{h^{}(x)}{h(x)}+\underset{s=j+1}{\overset{m}{}}\left(\sigma _{s(j+1)}(x)q_s(x)\frac{\sigma _{sj}^{}(x)}{(sj)}p_s(x)\right),$$
(10)
for $`j=0,1,2,\mathrm{},m1`$, where
$$\sigma _\kappa (x)=\underset{\nu =1}{\overset{\mathrm{}}{}}\alpha _\nu g_\nu ^\kappa (x),\text{for }\kappa =0,1,2,\mathrm{},m.$$
We notice that the first term in (10) does not appear if $`h(x)`$ is a constant or the last term does not appear if $`p_m(x)`$ is zero.
Proof. As it has been shown in , each of the factors $`yg_i(x)`$ involved in the expression $`I(x,y)`$ is a particular solution of system (1). Moreover, a particular solution $`yg(x)`$ has a related quasipolynomial cofactor of degree at most $`m1`$ in $`y`$, therefore we have that $`k(x,y)`$ is a polynomial in $`y`$ of degree at most $`m1`$. We will deduce each expression of $`k_i(x)`$ from the identity:
$$P(x,y)\frac{I}{x}(x,y)+Q(x,y)\frac{I}{y}(x,y)=k(x,y)I(x,y).$$
(11)
We compute the partial derivatives of $`I(x,y)`$ in (11) and we divide each member of the resulting expression by $`I(x,y)`$. We obtain:
$$P(x,y)\left[\frac{h^{}(x)}{h(x)}\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{\alpha _\nu g_\nu ^{}(x)}{yg_\nu (x)}\right]+Q(x,y)\left[\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{\alpha _\nu }{yg_\nu (x)}\right]=k(x,y).$$
Then, we deduce:
$$\begin{array}{c}P(x,y)\left[\frac{h^{}(x)}{h(x)}\underset{\nu =1}{\overset{\mathrm{}}{}}(yg_\nu (x))\underset{\nu =1}{\overset{\mathrm{}}{}}\alpha _\nu g_\nu ^{}(x)\underset{\mu =1,\mu \nu }{\overset{\mathrm{}}{}}(yg_\mu (x))\right]+\hfill \\ +Q(x,y)\underset{\nu =1}{\overset{\mathrm{}}{}}\alpha _\nu \left(\underset{\mu =1,\mu \nu }{\overset{\mathrm{}}{}}(yg_\mu (x))\right)=k(x,y)\underset{\nu =1}{\overset{\mathrm{}}{}}(yg_\nu (x)).\hfill \end{array}$$
(12)
The expression (12) is an identity of polynomials in $`y`$ of degree $`m+\mathrm{}`$. Let us consider the equality of coefficients of degree $`m+\mathrm{}`$ in $`y`$: $`p_m(x)h^{}(x)/h(x)=0`$, which implies that either $`p_m(x)0`$ or $`h(x)`$ is a constant function.
We define $`eq_i`$ as the equation resulting from identifying the coefficients of $`y^i`$ in both members of equation (12), $`i=0,1,2,\mathrm{},m+\mathrm{}1`$. From $`eq_{\mathrm{}+m1}`$, we can equate the expression of $`k_{m1}(x)`$. Once we know this expression, from $`eq_{\mathrm{}+m2}`$, we can equate the expression of $`k_{m2}(x)`$. Once we know this function, from $`eq_{\mathrm{}+m3}`$ we can equate the expression of $`k_{m3}(x)`$, and so on. Hence, in a recursive way, from $`eq_{\mathrm{}+j}`$, we equate $`k_j(x)`$, where $`j=m1,m2,m3,\mathrm{},2,1,0`$. It can be shown by induction that these expressions are given by (10). We notice that
$$\frac{\sigma _\kappa ^{}(x)}{\kappa }=\underset{\nu =1}{\overset{\mathrm{}}{}}\alpha _\nu g_\nu ^{\kappa 1}(x)g_\nu ^{}(x),$$
for $`\kappa =0,1,2,\mathrm{},m`$.
We substitute the values of $`k_j(x)`$ given in (10) in equation (12) and we deduce that a polynomial of degree at most $`\mathrm{}1`$ in $`y`$ must be zero. We denote by $`Pol(y)`$ this polynomial in $`y`$, which is also a function of $`x`$, $`h(x)`$ and $`g_1(x),g_2(x),\mathrm{},g_{\mathrm{}}(x)`$. We fix an index $`i`$ such that $`1i\mathrm{}`$ and when we substitute $`y`$ by $`g_i(x)`$ in $`Pol(y)`$, we get that:
$$Pol(g_i(x))=\underset{s=1,si}{\overset{\mathrm{}}{}}(g_i(x)g_s(x))\left[Q(x,g_i(x))P(x,g_i(x))g_i^{}(x)\right].$$
Since each $`g_i(x)`$ is a particular solution of system (1), we deduce that each $`g_i(x)`$ is a different root of the polynomial $`Pol(y)`$. Then, we have $`\mathrm{}`$ different roots of a polynomial of degree at most $`\mathrm{}1`$, so $`Pol(y)`$ is the null polynomial. We conclude that once we have substituted the cofactor as defined by (10) in equation (12) we have that this equation is satisfied.
In the same way as in Proposition 14, we are going to give the form of the quasipolynomial cofactor associated to an invariant of the form
$$I(x,y)=\mathrm{exp}\left\{h_2(x)\frac{A_1(x,y)}{A_0(x,y)}\right\}h_1(x)\underset{i=1}{\overset{\mathrm{}}{}}(yg_i(x))^{\alpha _i},$$
(13)
where we have defined $`A_1(x,y):=_{k=1}^r(ya_k(x))`$ and $`A_0(x,y):=_{j=1}^s(y\stackrel{~}{g}_j(x))`$.
The first thing that we notice is that, without loss of generality, we can assume that $`r=s`$. In case that $`rs`$, we can consider the change of variable $`y=1/z`$. After this change (and a reparameterization of the time variable if necessary), we obtain another polynomial system (1) and the function $`I`$ is transformed to another one with the same structure but with $`r=s`$. This last assertion is clear from the following equality:
$$\begin{array}{ccc}h_2(x)\frac{A_1(x,1/z)}{A_0(x,1/z)}\hfill & =\hfill & h_2(x)\frac{_{k=1}^r\left(\frac{1}{z}a_k(x)\right)}{_{j=1}^s\left(\frac{1}{z}\stackrel{~}{g}_j(x)\right)}=\hfill \\ & =\hfill & h_2(x)\frac{\frac{1}{z^r}_{k=1}^r\left(1a_k(x)z\right)}{\frac{1}{z^s}_{j=1}^s\left(1\stackrel{~}{g}_j(x)z\right)}=\stackrel{~}{h}_2(x)\frac{z^{sr}_{k=1}^r\left(z\frac{1}{a_k(x)}\right)}{_{j=1}^s\left(z\frac{1}{\stackrel{~}{g}_j(x)}\right)},\hfill \end{array}$$
where $`\stackrel{~}{h}_2(x):=h_2(x)(1)^{sr}_{k=1}^ra_k(x)\mathrm{}_{j=1}^s\stackrel{~}{g}_j(x)`$. If we ensure the structure of a quasipolynomial cofactor for the system with variables $`(x,z)`$, we deduce its structure for the system with variables $`(x,y)`$ just undoing the change and the reparameterization of the time, if it has been done. So, without loss of generality, we can assume that $`r=s`$. In fact, an analogous proof can be done for the case $`rs`$ but as it involves many computations which are not far from the ones that we are showing, we have preferred to avoid the case $`rs`$ by means of the change of variable $`y=1/z`$.
###### Proposition 15
Let $`I(x,y)`$ be an invariant of system (1) of the form (13) with $`r=s`$ and with an associated quasipolynomial cofactor $`k(x,y):=k_0(x)+k_1(x)y+\mathrm{}+k_{m1}(x)y^{m1}`$. Then, $`p_m(x)[h_1^{}(x)+h_1(x)h_2^{}(x)]0`$ and
$$k_j(x)=\left(\frac{h_1^{}(x)}{h_1(x)}+h_2^{}(x)\right)p_j(x)+\underset{i=j+1}{\overset{m}{}}\left(\stackrel{~}{\sigma }_{ij1}(x)q_i(x)\frac{\stackrel{~}{\sigma }_{ij}^{}(x)}{(ij)}p_i(x)\right),$$
(14)
for $`j=0,1,2,\mathrm{},m1`$, where either the first term does not appear in case that $`h_1(x)=c\mathrm{exp}\{h_2(x)\}`$, with $`c`$ and $`c0`$, or the last term does not appear if $`p_m(x)0`$. The functions $`\stackrel{~}{\sigma }_\kappa (x)`$ are defined in the following way. Given $`\kappa `$ we consider the set of indexes
$$J_\kappa :=\{(ϵ_1,ϵ_2,\mathrm{},ϵ_r,i_1,i_2,\mathrm{},i_r):\underset{k=1}{\overset{r}{}}ϵ_k+\underset{j=1}{\overset{r}{}}i_j=\kappa ,ϵ_j\{0,1\},i_j\{0\}\}$$
and we have that:
$$\stackrel{~}{\sigma }_\kappa :=\underset{\nu =1}{\overset{\mathrm{}}{}}\alpha _\nu g_\nu ^\kappa +\kappa h_2\underset{J_\kappa }{}(1)^{ϵ_1+ϵ_2+\mathrm{}+ϵ_r+1}a_1^{ϵ_1}a_2^{ϵ_2}\mathrm{}a_r^{ϵ_r}\stackrel{~}{g}_1^{i_1}\stackrel{~}{g}_2^{i_2}\mathrm{}\stackrel{~}{g}_s^{i_r},$$
for $`\kappa =0,1,2,\mathrm{},m`$.
Proof. We have that if $`I(x,y)`$ is an invariant, then each one of the functions $`g_i(x)`$ and $`\stackrel{~}{g}_j(x)`$ are particular solutions and the identity (7) is satisfied. We are going to deduce the expressions of $`k_j(x)`$ only assuming that the following identity is satisfied:
$$P(x,y)\frac{I}{x}(x,y)+Q(x,y)\frac{I}{y}(x,y)=k(x,y)I(x,y),$$
(15)
We multiply (15) by
$$A_0(x,y)^2\mathrm{exp}\left\{\frac{h_2(x)A_1(x,y)}{A_0(x,y)}\right\}\underset{i=1}{\overset{\mathrm{}}{}}(yg_i(x))^{1\alpha _i}$$
so as to get an identity of polynomials in $`y`$ of degree $`m+\mathrm{}+2r`$, where $`m`$ is the degree of system (1) in $`y`$. The equality of the coefficients of highest degree in $`y`$, that is $`y^{m+\mathrm{}+2r}`$, gives $`p_m(x)[h_1^{}(x)+h_1(x)h_2^{}(x)]0`$. This condition gives two possibilities: either $`p_m(x)0`$ or $`h_1(x)=c\mathrm{exp}\{h_2(x)\}`$, with $`c`$ and $`c0`$. Since the proof for both cases is analogous, we are going to follow them simultaneously. It has been shown in that $`k(x,y)`$ is a polynomial in $`y`$ of degree at most $`m1`$ so we can collect it in this variable: $`k(x,y)=_{i=1}^{m1}k_i(x)y^i`$. The equality of coefficients of degree $`j+\mathrm{}+2r`$ gives us the expression of $`k_j(x)`$ in a recursive way. We first compute $`k_{m1}`$ from the equality of coefficients of $`y^{m1+\mathrm{}+2r}`$, once we have this one, we compute $`k_{m2}`$ from the equality of coefficients of $`y^{m2+\mathrm{}+2r}`$ and so on. Some tedious computations show that these expressions are the ones given in (14).
We substitute the given expressions of $`k_j(x)`$ in the equality (15) and from the equality of the coefficients of $`y^{\mathrm{}+2r1}`$ we can compute the function $`h_2(x)`$. Now we have a polynomial in $`y`$ of degree at most $`\mathrm{}+2r1`$ in $`y`$ that must be identically zero. We denote by $`\overline{Pol}(y)`$ this polynomial in $`y`$. In the same way as in the proof of Proposition 14, we substitute the $`y`$ variable by each one of the $`\mathrm{}+2r`$ functions $`g_i(x)`$, $`\stackrel{~}{g}_j(x)`$ and $`a_k(x)`$ and we deduce that $`\overline{Pol}(g_i(x))0`$ because $`g_i(x)`$ is a particular solution for $`i=1,2,\mathrm{},\mathrm{}`$, $`\overline{Pol}(\stackrel{~}{g}_j(x))0`$ because $`\stackrel{~}{g}_j(x)`$ is a particular solution for $`j=1,2,\mathrm{},r`$ and $`\overline{Pol}(a_k(x))0`$ because relation (7) is satisfied for $`k=1,2,\mathrm{},r`$. Therefore, we have $`\mathrm{}+2r`$ different roots of a polynomial of degree at most $`\mathrm{}+2r1`$, so $`\overline{Pol}(y)`$ is the null polynomial. We conclude that once we have substituted the cofactor defined by (14) and the function $`h_2(x)`$ in equation (15) we have that this equation is satisfied.
Example. We include an example of the form of the quasipolynomial cofactor of a function $`I(x,y)`$ as given by Proposition 15 deduced from the equation (14). We consider the following system:
$$\dot{x}=y+y^2+x^2+4yx^2,\dot{y}=x2x^3+2xy^2,$$
(16)
which has the exponential factor $`I(x,y)=\mathrm{exp}\{(2y1)/(x^2+y^2)\}`$ with cofactor $`k(x,y):=4x`$. We are going to deduce the value of this cofactor by using the formulas stated in (14) of Proposition 15. We notice that this function $`I(x,y)`$ has $`r=1<s=2`$, so we need to perform the change of variables $`y=1/z`$. The transformed system needs a reparameterization of the time, meaning multiplying it by $`z^2`$, in order to have a polynomial system. Then, the transformed system reads for:
$$\dot{x}=1+(1+4x^2)z+x^2z^2,\dot{z}=2xz^2+(x+2x^3)z^4,$$
(17)
which has the exponential factor $`I(x,z)=\mathrm{exp}\{z(2z)/(x^2z^2+1)\}`$ with cofactor $`k(x,z)=4xz^2`$. This cofactor coincides with the transformation of the cofactor in coordinates $`(x,y)`$ after the reparameterization of time. Using the notation described in Proposition 15 we have that, for system (17), $`m=4`$, $`r=s=2`$, $`\mathrm{}=0`$, $`a_1(x):=0`$, $`a_2(x):=2`$, $`h_2(x):=1/x^2`$, $`\stackrel{~}{g}_1(x):=i/x`$, $`\stackrel{~}{g}_2(x):=i/x`$ and $`h_1(x):=1`$, where $`i=\sqrt{1}`$. We recall that $`p_i(x)`$ corresponds to the coefficient of degree $`i`$ in $`z`$ of the polynomial defining $`\dot{x}`$ and $`q_i(x)`$ corresponds to the coefficient of degree $`i`$ in $`z`$ of the polynomial defining $`\dot{z}`$. We have that $`p_4(x)0`$, so the first assertion of Proposition 15 is satisfied. We compute $`h_1^{}(x)/h_1(x)+h_2^{}(x)`$ which gives $`2/x^3`$ and the values of the $`\stackrel{~}{\sigma }_\kappa (x)`$ are equal to: $`\stackrel{~}{\sigma }_0(x)=0`$, $`\stackrel{~}{\sigma }_1(x)=2/x^2`$, $`\stackrel{~}{\sigma }_2(x)=2/x^4`$, $`\stackrel{~}{\sigma }_3(x)=6/x^4`$, $`\stackrel{~}{\sigma }_4(x)=4/x^6`$. An easy computation shows that the formulas written in (14) give $`k_0(x)=k_1(x)=k_3(x)0`$ and $`k_2(x)=4x`$, which corresponds to the given value of the cofactor $`k(x,z)=4xz^2`$.
### 3.4 Darboux functions obtained from invariants
We consider a Darboux function $`I(x,y)`$ which is the product of invariant algebraic curves up to complex numbers. Then, by computing $`y`$-roots, it can be expressed in the form $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$, where $`\alpha _i`$, $`g_i(x)`$ are algebraic particular solutions and $`h(x)`$ is such that its logarithmic derivative is a rational function. We recall that the logarithmic derivative of a function $`h(x)`$ is the quotient $`h^{}(x)/h(x)`$.
In this subsection, we are concerned with the converse of this problem, that is, we give the conditions that a function of the form $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$ must satisfy in order to be a Darboux function.
###### Proposition 16
Assume that the function $`I(x,y)=h(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$, where $`\alpha _i`$ and $`g_i(x)`$ are algebraic functions, and $`h(x)`$ is such that $`h^{}(x)/h(x)`$ is a rational function, satisfies that
$$P(x,y)\frac{I}{x}(x,y)+Q(x,y)\frac{I}{y}(x,y)=k(x,y)I(x,y),$$
where $`k(x,y)`$ is a polynomial, then $`I(x,y)`$ is a Darboux function.
Proof. By Proposition $`7`$ in , we have that each $`g_i(x)`$ is a particular solution of equation (1), which is algebraic by hypothesis. Therefore, by Theorem 6, we deduce that the irreducible polynomial $`f(x,y)`$ such that $`f(x,g(x))0`$ gives rise to an invariant algebraic curve of system (1). For each $`g_i(x)`$ there is an invariant algebraic curve $`f(x,y)=0`$, but each invariant algebraic curve $`f(x,y)=0`$ can implicitly define several functions $`g_i(x)`$, as much as the degree $`s`$ of $`f(x,y)`$ in $`y`$. We denote by $`g_{j_i}(x)`$ all the algebraic functions defined by the same invariant algebraic curve $`f_j(x,y)=0`$. Assume that the finite set $`\{g_i(x):i=1,\mathrm{},\mathrm{}\}`$ defines a total of $`r`$ invariant algebraic curves $`f_j(x,y)=0`$, $`j=1,2,\mathrm{},r`$. Hence, the function $`I(x,y)`$ is written as $`I(x,y)=h(x)(yg_{j_i}(x))^{\alpha _{j_i}}`$ where the product is taken over all the possible subindexes. Some of the $`\alpha _{j_i}`$ can be null in this denomination.
By Proposition 8, we have that each particular solution $`y=g_{j_i}(x)`$ has an associated quasipolynomial cofactor, which will be denoted by $`M_{j_i}(x,y)`$. Each invariant algebraic curve $`f_j(x,y)=0`$ has a polynomial cofactor denoted by $`k_{f_j}(x,y)`$. We define $`M_0(x,y):=h^{}(x)P(x,y)/h(x)`$, which is a rational function by hypothesis. Each polynomial $`f_j(x,y)`$ decomposes, by Theorem 5, as $`f_j(x,y)=h_j(x)_{i=1}^{s_j}(yg_{j_i}(x))`$, where $`h_j(x)`$ is a function of $`x`$ such that $`h_j^{}(x)/h_j(x)`$ is a rational function, $`g_{j_i}(x)`$ are algebraic functions and $`s_j`$ is the degree of $`f_j(x,y)`$ in $`y`$. We define $`M_{j_0}(x,y):=h_j^{}(x)P(x,y)/h_j(x)`$ which is a rational function. We have that:
$`P(x,y){\displaystyle \frac{f_j}{x}}+Q(x,y){\displaystyle \frac{f_j}{y}}`$ $`=`$ $`P(x,y)h_j^{}(x){\displaystyle \underset{i=1}{\overset{s_j}{}}}(yg_{j_i}(x))+h_j(x)`$
$`\left[{\displaystyle \underset{i=1}{\overset{s_j}{}}}\left(P(x,y)g_{j_i}^{}(x)+Q(x,y)\right){\displaystyle \underset{\nu =1,\nu i}{\overset{s_j}{}}}(yg_{j_\nu }(x))\right]`$
$`=`$ $`\left(M_{j_0}(x,y)+{\displaystyle \underset{i=1}{\overset{s_j}{}}}M_{j_i}(x,y)\right)h_j(x){\displaystyle \underset{i=1}{\overset{s_j}{}}}(yg_{j_i}(x)),`$
$`=`$ $`\left(M_{j_0}(x,y)+{\displaystyle \underset{i=1}{\overset{s_j}{}}}M_{j_i}(x,y)\right)f_j(x,y).`$
On the other hand,
$$P(x,y)\frac{f_j}{x}(x,y)+Q(x,y)\frac{f_j}{y}(x,y)=k_{f_j}(x,y)f_j(x,y),$$
for being $`f_j(x,y)=0`$ an invariant algebraic curve with cofactor $`k_{f_j}(x,y)`$. Therefore, we deduce the following identities:
$$M_{j_0}(x,y)+\underset{i=1}{\overset{s_j}{}}M_{j_i}(x,y)=k_{f_j}(x,y),$$
(18)
for $`j=1,2,\mathrm{},r`$. We have that each $`M_{j_i}(x,y)`$ is a rational function of $`x`$, a polynomial in $`g_{j_i}(x)`$ and a polynomial in $`y`$, by Proposition $`7`$ in . We can change the powers of $`g_{j_i}(x)`$ which are equal to or higher than $`s_j`$ to a linear combination of lower powers by using the expression $`f_j(x,g_{j_i}(x))0`$. Therefore, equating the powers of $`y`$ (after changing all the powers $`g_{j_i}(x)`$ to combinations of $`g_{j_i}(x)^\nu `$, $`0\nu <s_j`$), we get that the $`r`$ identities (18) give a total of $`r(m1)`$ linear combinations of the powers $`g_{j_i}(x)^\nu `$, $`0\nu <s_j`$, where $`m`$ is the degree of system (1) in the variable $`y`$.
We consider the identity
$$P(x,y)\frac{I}{x}(x,y)+Q(x,y)\frac{I}{y}(x,y)=k(x,y)I(x,y),$$
from which we deduce that
$$M_0+\underset{i=1}{\overset{s_1}{}}\alpha _{1_i}M_{1_i}+\underset{i=1}{\overset{s_2}{}}\alpha _{2_i}M_{2_i}+\mathrm{}+\underset{i=1}{\overset{s_r}{}}\alpha _{r_i}M_{r_i}=k.$$
(19)
As before, we can change each power $`g_{j_i}(x)^\nu `$ with $`\nu s_j`$ to a linear combination of the powers $`g_{j_i}(x)^\nu `$, $`0\nu <s_j`$. And from (19) we deduce an identity which is a linear combination of the powers $`g_{j_i}(x)^\nu `$, $`i=1,2,\mathrm{},s_j`$, $`0\nu <s_j`$ and $`j=1,2,\mathrm{},r`$.
We notice that if $`\alpha _{j_i}=\beta _j`$, where $`\beta _j`$, for all $`1is_j`$, and $`1jr`$, then the relations given by (18) make the identity (19) compatible with the fact that $`k(x,y)`$ is a polynomial in both variables $`x`$ and $`y`$. Assume that this is not the case. Assume that we have $`\alpha _{j_\nu }\alpha _{j_\upsilon }`$ for certain $`\nu ,\upsilon `$. We can assume that $`j=1`$ without loss of generality. We can consider the $`s_j`$ symmetric polynomials defined by each $`f_j(x,y)`$ which are linear combinations of the powers $`g_{j_i}(x)^\nu `$, $`i=1,2,\mathrm{},s_j`$, $`0\nu <s_j`$, $`j=1,2,\mathrm{},r`$. By using the elimination theory, we can eliminate all the appearances of $`g_{j_i}(x)`$ with $`j>1`$ in (19). We obtain in this way a relation only involving $`g_{1_i}(x)`$. By using the symmetric polynomials associated to $`f_1(x,y)`$, we eliminate all the $`g_{1_i}(x)`$, except one, which may be $`g_{1_1}(x)`$. The resulting relation $`R(x,g_{1_1}(x))0`$ can only be of two forms: either $`R(x,y)`$ is a multiple of $`f_1(x,y)`$ or it is not. In the first case, we have that the relation given by (19) is a combination of the symmetric polynomials, which are symmetric with respect to $`g_{1_i}(x)`$, $`i=1,2,\mathrm{},s_1`$. This symmetry implies that $`\alpha _{j_\nu }=\alpha _{j_\upsilon }`$, for all $`\nu `$ and $`\upsilon `$. In the second case, we would get that $`R(x,y)`$ is a polynomial such that $`R(x,g_{1_1}(x))0`$ and the degree of $`R(x,y)`$ in $`y`$ is lower than $`s_1`$. We recall that we have already substituted all the appearances of powers of $`g_{1_1}(x)`$ of higher degree by the corresponding expression given by the equation $`f_1(x,g_{1_1}(x))0`$. The existence of a polynomial $`R(x,y)`$ such that $`R(x,g_{1_1}(x))0`$ and the degree of $`R(x,y)`$ in $`y`$ being lower than $`s_1`$ is a contradiction with the fact that $`f_1(x,y)`$ is the only irreducible polynomial satisfying $`f_1(x,g_{1_1}(x))0`$, modulus associates.
Hence, we conclude that the only possibility is that $`\alpha _{j_i}=\beta _j`$, where $`\beta _j`$, for all $`1is_j`$, and $`1jr`$. We have that:
$$I(x,y)=h(x)\frac{_{j=1}^rf_j(x,y)^{\beta _j}}{_{j=1}^rh_j(x)^{\beta _j}}$$
and we define $`\stackrel{~}{h}(x)=h(x)_{j=1}^rh_j(x)^{\beta _j}`$. We notice that since $`h(x)`$ and $`h_j(x)`$, $`j=1,2,\mathrm{},r`$ satisfy that its logarithmic derivative is a rational function, $`\stackrel{~}{h}(x)`$ also satisfies that $`\stackrel{~}{h}^{}(x)/\stackrel{~}{h}(x)`$ is a rational function. By integration, we obtain that $`\stackrel{~}{h}(x)`$ is a Darboux function. We deduce that $`I(x,y)`$ is equal to $`\stackrel{~}{h}(x)_{j=1}^rf_j(x,y)^{\beta _j}`$ which is a Darboux function, as we wanted to show.
Example. We are going to describe an example of Proposition 16 so as to make the proof clearer. Let us consider the following planar polynomial differential system:
$$\begin{array}{ccc}\dot{x}\hfill & =\hfill & 55x+15y^26x^2y+14xy^29xy^4,\hfill \\ \dot{y}\hfill & =\hfill & 5+2x3y2xy^2+6y^33y^5.\hfill \end{array}$$
(20)
This system exhibits two invariant algebraic curves of degree $`3`$: $`f_1(x,y)=0`$ with $`f_1(x,y):=y^3yx`$ and $`f_2(x,y)=0`$ with $`f_2(x,y):=xy^2x1`$. Their cofactors are, respectively, $`k_{f_1}(x,y)=3(1+2xy4y^2+3y^4)`$ and $`k_{f_2}(x,y)=5k_1(x,y)/3`$. We factorize the polynomial $`f_1(x,y)`$ as $`h_1(x)(yg_{1_1}(x))(yg_{1_2}(x))(yg_{1_3}(x))`$ where $`h_1(x):=1`$ and $`g_{1_i}(x)`$, $`i=1,2,3`$, are the corresponding $`y`$-roots of $`f_1(x,y)`$. It is easy to see that the polydromy order of $`g_{1_i}(x)`$ is $`1`$. We perform the same computations for $`f_2(x,y)`$ and we have that $`f_2(x,y)=h_2(x)(yg_{2_1}(x))(yg_{2_2}(x))`$ where $`h_2(x):=x`$, $`g_{2_1}(x):=\sqrt{1+1/x}`$ and $`g_{2_2}(x):=\sqrt{1+1/x}`$. It is easy to see that the polydromy order of $`g_{2_i}(x)`$, $`i=1,2`$, is $`2`$. We have that $`yg_{j_i}(x)`$ are algebraic particular solutions and we can compute their corresponding quasipolynomial cofactors $`M_{j_i}(x,y)`$ which are:
$`M_{1_i}(x,y)`$ $`:=`$ $`{\displaystyle \frac{3}{13g_{1_i}^2(x)}}(1+5y+2y^2y^4+(5+2x2y+2y^3)g_{1_i}(x)+`$
$`+(1+2xy4y^2+3y^4)g_{1_i}^2(x)),`$
$`M_{2_i}(x,y)`$ $`:=`$ $`{\displaystyle \frac{1}{2xg_{2_i}(x)}}(5y10x4x^2+6xy15y^36xy^3+`$
$`+(515y^24x^2y+6xy^26xy^4)g_{2_i}(x)).`$
We have that $`g_{1_i}(x)`$, $`i=1,2,3`$ are the $`y`$-roots of $`f_1(x,y)`$ and, hence, they satisfy the following relationships, given by the symmetric polynomials on the $`y`$–roots:
$$\begin{array}{cccc}V_{1_1}\hfill & :=\hfill & g_{1_1}(x)+g_{1_2}(x)+g_{1_3}(x)\hfill & =0,\hfill \\ V_{1_2}\hfill & :=\hfill & g_{1_1}(x)g_{1_2}(x)+g_{1_1}(x)g_{1_3}(x)+g_{1_2}(x)g_{1_3}(x)\hfill & =1,\hfill \\ V_{1_3}\hfill & :=\hfill & g_{1_1}(x)g_{1_2}(x)g_{1_3}(x)\hfill & =x.\hfill \end{array}$$
In the same way we have that $`V_{2_1}:=g_{2_1}(x)+g_{2_2}(x)=0`$ and $`V_{2_2}:=g_{2_1}(x)g_{2_2}(x)=11/x`$. These relationships give that $`M_{1_1}+M_{1_2}+M_{1_3}=k_1`$ and $`M_{2_0}+M_{2_1}+M_{2_2}=k_2`$ where $`M_{2_0}=h_2^{}(x)\dot{x}/h_2(x)=(55x+15y^26x^2y+14xy^29xy^4)/x`$. We consider a function of the form:
$$I(x,y)=h(x)(yg_{1_1}(x))^{\alpha _{1_1}}(yg_{1_2}(x))^{\alpha _{1_2}}(yg_{1_3}(x))^{\alpha _{1_3}}(yg_{2_1}(x))^{\alpha _{2_1}}(yg_{2_2}(x))^{\alpha _{2_2}},$$
where $`h(x)`$ is such that $`h^{}(x)/h(x)`$ is a rational function (we define $`M_0(x,y):=h^{}(x)\dot{x}/h(x)`$), $`\alpha _{j_i}`$ and $`g_{j_i}(x)`$ are the functions defined above. We also assume that $`𝒳\left(I(x,y)\right)=k(x,y)I(x,y)`$ where $`k(x,y)`$ is a polynomial in $`x`$ and $`y`$ of degree at most $`4`$. We have that:
$$M_0+\alpha _{1_1}M_{1_1}+\alpha _{1_2}M_{1_2}+\alpha _{1_3}M_{1_3}+\alpha _{2_1}M_{2_1}+\alpha _{2_2}M_{2_2}=k.$$
This identity gives rise to five equations relating the $`g_{j_i}(x)`$ and $`\alpha _{j_i}`$ when equating the coefficients of different degrees of $`y`$. The equation corresponding to the degree $`4`$ in $`y`$ is equal to:
$$3\left(\alpha _{1_1}+\alpha _{1_2}+\alpha _{1_3}+\alpha _{2_1}+\alpha _{2_2}+3x\frac{h^{}(x)}{h(x)}\right)=k_4(x),$$
(21)
where $`k_j(x)`$ is the coefficient of degree $`j`$ in $`y`$ of $`k(x,y)`$. We note that $`k_4(x)`$ is a real number since $`k(x,y)`$ is a polynomial in $`x`$ and $`y`$ of degree at most $`4`$. Then, equation (21) implies that $`h(x)=x^b`$ with $`b=(\alpha _{1_1}+\alpha _{1_2}+\alpha _{1_3}+\alpha _{2_1}+\alpha _{2_2})/3k_4/9`$. The equation corresponding to the degree $`3`$ in $`y`$ gives:
$$\underset{i=1}{\overset{3}{}}\left(\frac{6\alpha _{1_i}g_{1_i}(x)}{13g_{1_i}(x)^2}\right)\frac{3(5+2x)}{2x}\left(\frac{\alpha _{2_1}}{g_{2_1}(x)}+\frac{\alpha _{2_2}}{g_{2_2}(x)}\right)=k_3(x).$$
The following equations correspond to the degrees $`2`$, $`1`$ and $`0`$ in $`y`$, respectively:
$$(15+14x)\frac{h^{}(x)}{h(x)}+6\underset{i=1}{\overset{3}{}}\left(\frac{\alpha _{1_i}(2g_{1_i}(x)^21)}{3g_{1_i}(x)^21}\right)+\frac{3(2x5)}{2x}(\alpha _{2_1}+\alpha _{2_2})=k_2(x),$$
$$\begin{array}{c}3\underset{i=1}{\overset{3}{}}\left(\alpha _{1_i}\frac{52g_{1_i}(x)+2g_{1_i}(x)^2}{13g_{1_i}(x)^2}\right)+\underset{i=1}{\overset{2}{}}\left(\alpha _{2_i}\frac{5+6x4x^2g_{2_i}(x)}{2xg_{2_i}(x)}\right)=\hfill \\ =6x^2\frac{h^{}(x)}{h(x)}+k_1(x),\hfill \end{array}$$
$$\begin{array}{c}3\underset{i=1}{\overset{3}{}}\left(\alpha _{1_i}\frac{1+(2x+5)g_{1_i}(x)+g_{1_i}(x)^2}{13g_{1_i}(x)^2}\right)+\underset{i=1}{\overset{2}{}}\left(\alpha _{2_i}\frac{5g_{2_i}(x)10x4x^2}{2xg_{2_i}}\right)=\hfill \\ =5(1+x)\frac{h^{}(x)}{h(x)}+k_0(x).\hfill \end{array}$$
We consider a common denominator in each one of these equations and by using elimination theory among the numerators of these equations and the polynomials $`V_{j_i}`$, we deduce that the only possibilities are $`\alpha _{1_1}=\alpha _{1_2}=\alpha _{1_3}=:\beta _1`$, $`\alpha _{2_1}=\alpha _{2_2}=:\beta _2`$, $`h(x)=x^{\beta _2}`$ and $`k(x,y)=(3\beta _1+5\beta _2)(1+2xy4y^2+3y^4)`$. Hence, $`I(x,y)=f_1(x,y)^{\beta _1}f_2(x,y)^{\beta _2}`$, which is a Darboux function.
We know that taking $`y`$–roots any Darboux function can be expressed in the form $`I(x,y)=\mathrm{exp}\left\{h_2(x)_{k=1}^r(ya_k(x))\mathrm{}_{j=1}^s(y\stackrel{~}{g}_j(x))\right\}h_1(x)_{i=1}^{\mathrm{}}(yg_i(x))^{\alpha _i}`$ where $`\alpha _i`$, $`g_i(x)`$, $`\stackrel{~}{g}_j(x)`$ and $`a_k(x)`$ are algebraic functions, $`h_1(x)`$ is such that its logarithmic derivative is a rational function and $`h_2(x)`$ is a rational function. Next Proposition gives the converse of this assertion, that is, we give the conditions that a function of the form $`I(x,y)`$ must satisfy in order to be a Darboux function.
###### Proposition 17
Assume that the function
$$I(x,y)=\mathrm{exp}\left\{h_2(x)\frac{_{k=1}^r(ya_k(x))}{_{j=1}^s(y\stackrel{~}{g}_j(x))}\right\}h_1(x)\underset{i=1}{\overset{\mathrm{}}{}}(yg_i(x))^{\alpha _i},$$
where $`\alpha _i`$, $`g_i(x)`$, $`\stackrel{~}{g}_j(x)`$ and $`a_k(x)`$ are algebraic functions, $`h_1(x)`$ is such that its logarithmic derivative is a rational function and $`h_2(x)`$ is a rational function, satisfies that:
$$P(x,y)\frac{I}{x}(x,y)+Q(x,y)\frac{I}{y}(x,y)=k(x,y)I(x,y),$$
where $`k(x,y)`$ is a polynomial, then $`I(x,y)`$ is a Darboux function.
Proof. From Theorem 2 in , we have that each $`g_i(x)`$ is a particular solution of (2), so it has an associated quasipolynomial cofactor (see Proposition 8) which we denote by $`M_i(x,y)`$. We define $`\mathrm{\Phi }(x,y):=\mathrm{exp}\{h_2(x)A_1(x,y)/A_0(x,y)\}`$ where $`A_1(x,y)=_{k=1}^r(ya_k(x))`$ and $`A_0(x,y)=_{j=1}^s(y\stackrel{~}{g}_j(x))`$. Since $`𝒳(I)=kI`$, we deduce that
$$𝒳\left(\mathrm{\Phi }(x,y)\right)=\left(k(x,y)\frac{h_1^{}(x)}{h_1(x)}P(x,y)\underset{i=1}{\overset{\mathrm{}}{}}\alpha _iM_i(x,y)\right)\mathrm{\Phi }(x,y).$$
Therefore, we are under the hypothesis of Proposition 12 and Theorem 13. Hence, we realize that the proof of this assertion goes exactly as the proof of Proposition 16.
As a consequence of Theorem 6 and Propositions 14, 15, 16 and 17 we can establish the following result.
###### Theorem 18
Assume that system (1) has a first integral or an integrating factor of the form
$$I(x,y)=\mathrm{exp}\left\{h_2(x)\frac{_{k=1}^r(ya_k(x))}{_{j=1}^s(y\stackrel{~}{g}_j(x))}\right\}h_1(x)\underset{i=1}{\overset{\mathrm{}}{}}(yg_i(x))^{\alpha _i},$$
where $`\alpha _i`$, $`g_i(x)`$, $`\stackrel{~}{g}_j(x)`$ and $`a_k(x)`$ are algebraic functions, $`h(x)`$ and $`h_1(x)`$ have a rational logarithmic derivative and $`h_2(x)`$ is a rational function. Then, $`I(x,y)`$ is a Darboux function.
Proof. We are under the hypothesis of Proposition 16 or 17. In this case, $`k(x,y)`$ is identically zero or $`k(x,y)`$ is minus the divergence of the system. We deduce that $`I(x,y)`$ must be a Darboux function.
We notice that when we apply the method of constructing first integrals given in and all the $`h(x)`$ and $`g_i(x)`$ are completely determined functions, by Theorem 18 we have that this first integral is a Darbouxian function. However, the reciprocal is not true. As some examples in show, we may have a system with a Darbouxian first integral, but when we apply the method described in we get a nonlinear superposition principle. That is, not all the $`g_i(x)`$ introduced in the ansatz $`I(x,y)`$ are determined. If we choose those undetermined $`g_i(x)`$ as algebraic particular solutions, we will have that the first integral given by the superposition principle becomes a Darboux function.
Addresses and e-mails:
$`^{(1)}`$ Lab. de Mathématiques et Physique Théorique. CNRS UMR 6083.
Faculté des Sciences et Techniques. Université de Tours.
Parc de Grandmont, 37200 Tours, FRANCE.
E-mail: Hector.Giacomini@lmpt.univ-tours.fr
$`^{(2)}`$ Departament de Matemàtica. Universitat de Lleida.
Avda. Jaume II, 69. 25001 Lleida, SPAIN.
E–mails: gine@eps.udl.es, mtgrau@matematica.udl.es
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# Chandra Observations of Shock Kinematics in Supernova Remnant 1987A
## 1 Introduction
With the rapidly developing impact of the debris of Supernova 1987A with its inner circumstellar ring, we have an unprecedented opportunity to witness the birth of a supernova remnant, SNR1987A (McCray 2005). This event, the first hint of which was the 1995 appearance of Spot 1, a rapidly brightening optical ‘hotspot’ on the ring, has now evolved to the stage that the ring is encircled by hotspots (Sugerman et al. 2002). Additionally, an annular source of X-ray emission correlated with the locations of the hotspots has brightened at an accelerating rate (Park et al. 2004). Evidently, the optical hotspots appear where the supernova blast wave encounters fingers of relatively dense gas protruding inwards from the ring. In such a situation, a complex hydrodynamic interaction ensues (cf. Borkowski, Blondin & McCray 1997a,b). The optical emission from the spots comes from relatively slow ($`V_S200`$ km s<sup>-1</sup>) shocks which have had time to undergo radiative cooling (Pun et al. 2002). The X-ray emission must come from gas heated by faster ($`V_S1000`$ km s<sup>-1</sup>) shocks, either transmitted shocks entering the protrusions or shocks reflected from the protrusions.
Michael et al. (2002) reported the first observation of a dispersed X-ray spectrum of SNR 1987A, taken in October 1999 with HETG on the Chandra observatory. The X-ray emission was dominated by shock-heated gas having electron temperature $`kT_e2.6`$ keV. Due to the poor photon statistics, only a composite line profile was constructed by stacking the profiles of the individual observed lines. From the measured FWHM ($`5000`$ km s<sup>-1</sup>), they inferred that the X-ray emitting gas was moving with radial velocity $`3500`$ km s<sup>-1</sup>, roughly the same as that inferred from the observed proper motion of the non-thermal radio source (Manchester et al. 2002) and the X-ray source (Park et al. 2004).
From October 1999 to September 2004, the X-ray source SNR1987A has brightened by a factor $`10`$ (Park et al. 2005). As a result, it has become possible with Chandra to obtain dispersed X-ray spectra with very high counting statistics. In this Letter, we report the first results from the analysis of such observations.
## 2 Observations and Data Reduction
SNR 1987A was observed with Chandra in the configuration LETG-ACIS-S in five consecutive runs during Aug 26 – Sep 5, 2004, providing a total effective exposure of 289 ksec. We extracted<sup>1</sup><sup>1</sup>1For CIAO 3.1 and ATOMDB see http://cxc.harvard.edu/ciao/ and http://cxc.harvard.edu/atomdb/ positive and negative first-order LETG spectra for each of the five observations. Then, we merged the resultant spectra into one spectrum each for the positive and negative LETG arms with respective total counts of 9,241 and 6,057 in the energy range 0.4 - 7 keV. The difference in photon statistics is a result of the different sensitivities of the respective CCD detectors. We also extracted the pulse-height spectrum from the zeroth-order image with a total number of 16,557 counts in the 0.4 - 7 keV range. Figure 1 demonstrates the enormous scientific advantage of the dispersed spectrum over the pulse-height spectrum.
We fitted the strong emission line triplets of various He-like ions by a sum of three Gaussians and constant local continuum. The ratio of line centers of the triplet components were held fixed according to the values given by the Chandra atomic database<sup>1</sup> and all components shared the same full width at half maximum (FWHM). Thus, free parameters were the intensities of their components, the line center wavelength of one of the components, FWHM, and the local continuum. Likewise, we fitted the strong emission doublets of the H-like species but the component intensities were fixed through their atomic data values. We found that the centroid shifts for the strong emission lines were consistent with the red-shift of the Large Magellanic Cloud, and that the line fluxes derived from the positive and negative first-order spectra agreed within expected statistical uncertainties. Therefore, we assumed that all the line centers had the same Doppler shift, $`V_D=286`$ km s<sup>-1</sup> , and fitted the positive and negative first-order spectra simultaneously.
## 3 Line Fluxes and Ratios
Table 1 lists the fluxes of all the emission lines and multiplets that could be measured with acceptable photon statistics. Given the excellent spectral resolution and photon statistics, we can for the first time derive reliable intensities and ratios of various X-ray emission lines from SNR 1987A. As we shall discuss in § 5, we expect the actual conditions in the plasma to be sufficiently complex that we can only use the line strengths and ratios to infer typical conditions of the plasma responsible for the emission by a given ion. The ratio of He-like(K<sub>α</sub>)/H-like(Ly<sub>α</sub>) from a given element is sensitive to both temperature and ionization state of the gas. The same is true for the ‘G-ratio’ of the He-like triplet lines ($`G=\frac{f+i}{r}`$ where ‘f’, ‘i’ and ‘r’ stand for forbidden, intercombination and recombination lines, respectively) since these lines are produced not only by electron impact excitation of the He-like ions but also by K-shell ionization of the Li-like ions (Mewe 1999; Liedahl 1999).
Since we expect that shocks are responsible for heating and ionizing the X-ray emitting gas, we can compare our measured line ratios to those resulting from the XSPEC code, which provides models of the time-dependent ionization and X-ray emission from plane-parallel shocks (Borkowski, Lyerly, & Reynolds 2001). In such models, the line ratios are functions of the post-shock temperature and the ionization age, $`n_et`$, defined as the time since the gas first entered the shock times the postshock electron density. The X-ray emitting plasma in SNR 1987A is in NEI (nonequilibrium ionization), and so inner-shell ionization and excitation contribute importantly to emission line ratios. Indeed, if these processes are not included, the theoretical G-ratios cannot match the observed ones.
In Figure 2, the curves define those regions in the parameter space of electron temperature and ionization age, $`n_et`$, for which the observed line ratios agree with the ratios that would result from a plane parallel shock. Except for a very narrow range of $`n_et`$, the inferred temperatures are likely to be in the range 0.1–2 keV. We see immediately that no single combination of electron temperature and ionization age is consistent with all the observed ratios. Instead, for any given ionization age, the inferred electron temperature increases tentatively with ionization potential. Evidently, the X-ray emission comes from a distribution of shocks having a range of ionization ages and post-shock temperatures.
## 4 Line Profiles
The excellent spatial and spectral resolution of Chandra allows us for the first time to observe the kinematics of the X-ray emitting gas through the line profiles. The method is illustrated by Figure 3.<sup>2</sup><sup>2</sup>2For the general properties of LETG see §9 of the Chandra Proposer’s Observatory Guide, Version 7.0, p.195; Complexity of the spatial-spectral effects is discussed in § 8.5.3, pp.187-189; § 9.3.3, pp.209-215. We assume that the X-ray source lies on a circular ring having angular radius $`\theta _R`$ and that this ring is expanding with constant radial velocity, $`V_R`$. We also assume that the X-ray source has, like the optical inner circumstellar ring, inclination angle $`i=4445^{}`$, with the near side to the north and minor axis at P.A. $`354^{}`$ (Sugerman et al. 2002). The roll angle was chosen so that the negative ($`m=1`$) arm of the dispersion axis was aligned at P.A. $`345^{}`$, thus, the north side of the ring will be blue-shifted and the south side will be red-shifted. The dispersed images of the ring will be distorted by these Doppler shifts.<sup>2</sup> In the $`m=1`$ image, the N side of the ring will be displaced to the left, and the S side will be displaced toward the right. Thus, the minor axis of the $`m=1`$ image will be compressed. Likewise, the minor axis of the $`m=+1`$ image will be stretched.
This behavior is exactly what we see in the measured line profiles. We fitted Gaussian profiles (as explained in § 2) to each emission line or multiplet independently in the $`m=+1`$ and $`m=1`$ first-order spectra. In every case, the width of a given emission line in the $`m=+1`$ arm is greater than the corresponding width in the $`m=1`$ arm (Fig. 1). The actual line profile will likely differ from a Gaussian due to the complex spatial-spectral effects. However, the Gaussian approximation is sufficient for the analysis presented here (resulting in a reduced $`\chi ^21`$ for each of the line-profile fits).
There are two major sources of line broadening in the dispersed images. One is the spatial extent of the image itself, which can be expressed as an equivalent line broadening, independent of wavelength. The second is the broadening due to the bulk motion of the shocked gas. We express the net line width (FWHM) broadening as:
$$\mathrm{\Delta }\lambda _{tot}=2\mathrm{\Delta }\lambda _0\pm 2z_0(\lambda /\lambda _0)^\alpha \lambda ,$$
(1)
where the plus (minus) sign refers to the $`m=+1`$ ($`m=1`$) spectrum, respectively. The two sources of broadening are represented respectively by the first and second term on the right hand side of equation (1). The power-law function of wavelength allows for the possibility that the mean bulk velocity of shocked gas emitting a given line may depend on the excitation or ionization stage of the emitting ion. The parameter $`z_0`$ determines the line broadening at some fiducial wavelength, $`\lambda _0`$, and the power-law index $`\alpha `$ is to be determined.
We then fit equation (1) to the data in Figure 1. We consider two models. In the first (constant velocity) model, we assume that all emitting ions have the same mean bulk velocity ($`\alpha =0`$). In that case, we find a best-fit value $`z_0=0.0006`$. In the second (stratified) model, we allow the values of $`\alpha `$ and $`z_0`$ to vary. In that case, we find a best-fit value $`\alpha =1.3`$. Note that such a negative value of $`\alpha `$ implies that emission lines with shorter wavelength are produced by gas having greater radial velocity, as expected. The stratified model provides a slightly better fit to the data than the constant velocity model ($`\chi ^2/dof=22/19`$ vs. $`25/20`$), but the data are not good enough to distinguish between the models with confidence. We may translate the fitting parameters into an equivalent mean radial bulk velocity, $`V_r`$, of an expanding circular ring by adopting an average value for the azimuthal $`\varphi `$ ($`\overline{\mathrm{sin}\varphi }=2/\pi ,0\varphi \pi /2`$) and a value ($`i=45^{}`$) for the inclination of the inner ring. This yields a value $`V_r=397`$ km s<sup>-1</sup>. In the case of a strong radial shock with adiabatic index $`\gamma =5/3`$, the corresponding shock velocity must be $`V_S=4V_r/3=530`$ km s<sup>-1</sup> . The corresponding shock velocity in the stratified case is given by $`V_S(\lambda )=920(\lambda /10`$Å)<sup>-1.3</sup> km s<sup>-1</sup> . Finally, we note that in the both models the derived source half-size of $`\mathrm{\Delta }\lambda _0=0.047`$ Å or $`\theta _R=0\stackrel{}{\mathrm{.}}84`$ is consistent with the SNR 1987A size from an X-ray image-deconvolution technique (Burrows et al. 2000). Thus, we are confident that the simplified treatment presented here gives reliable results.
## 5 Discussion and Conclusions
The most surprising result of this observation is the relatively low velocity of the X-ray emitting gas as determined from the line profiles. When we proposed to do this observation, we expected to see kinematic velocities in the range $`25003000`$ km s<sup>-1</sup> , as we saw in the composite line profile measured with the Chandra HETG in October 1999 (Michael et al. 2002) and in the radial expansion rate of the X-ray image (Park et al 2004). At the very least, we would expect to see velocities comparable to the velocity of gas behind a shock moving sufficiently fast to heat the electrons to temperatures inferred from the X-ray line ratios. For a shock of high Mach number entering stationary gas with velocity $`V_S`$, the maximum electron temperature is: $`kT_e=\frac{3}{16}\mu V_S^2=1.4(V_S/1000\mathrm{km}\mathrm{s}^1)^2`$ keV. The value $`V_S=530`$ km s<sup>-1</sup> inferred for the constant velocity model then implies a post-shock temperature $`kT_e0.39`$ keV. This value is inconsistent with the temperatures required to account for the $`K_\alpha /L_\alpha `$ line ratios observed for Mg and Si (Fig. 2). On the other hand, the electron temperatures inferred from the line profiles in the stratified model may be consistent with the observed line ratios ($`kT_e=0.154.0`$ keV for $`V_S=3401,700`$ km s<sup>-1</sup> ).
How do we reconcile the relatively low gas velocities inferred from the line profiles with the much greater velocities inferred from the radial expansion rate of the X-ray image? We have in mind a picture in which a blast wave of velocity $`V_B3000`$ km s<sup>-1</sup> propagates through the low atomic density ($`n_0100`$ cm<sup>-3</sup>) gas inside the circumstellar ring. A zone of enhanced X-ray emission appears when the blast wave strikes a finger of relatively dense ($`n_13\times 10^33\times 10^4`$ cm<sup>-3</sup>) gas protruding inward from the ring. A transmitted shock propagates into the protrusion at velocity reduced by a factor $`(n_1/n_0)^{1/2}`$, while a reflected shock propagates backwards. The latter increases the density and temperature of the circumstellar matter to values substantially greater than those caused by the original blast wave, and it also reduces the radial velocity of the doubly-shocked circumstellar gas. The enhanced X-ray emission comes from the gas behind the transmitted shock and from the gas behind the reflected shock, with the proportional contributions depending on the details of the hydrodynamics. But, the lower velocity of the transmitted shock and the high density behind it ensure a short cooling timescale and eventually increased optical emission. At each point of impact with a protrusion, a new zone of enhanced X-ray emission appears. Thus, the mean radius of the enhanced X-ray emission appears to move outward at a substantial fraction of the blast wave velocity, while the gas responsible for most of the X-ray emission may be moving much more slowly. This picture assumes that only a small fraction of the blast wave area is covered by protrusions at the present time. It also implies that the radial expansion of the X-ray image should slow down rapidly as the blast wave overtakes the entire equatorial ring.
The scenario that we have described here accounts in a natural way for several properties observed by Chandra images and spectra of SNR 1987A: (1) the correlation of the rapid X-ray brightening with the appearance of the optical hotspots; (2) the correlation of the X-ray image with the optical hotspots; (3) the relatively rapid expansion of the X-ray image compared to the relatively slow bulk velocity of the X-ray-emitting gas; and (4) the correlation of the inferred electron temperature and expansion velocity of the shocked gas with excitation potential of X-ray emission lines.
There is much more to be learned from the Chandra spectroscopic data beyond the brief summary that we have presented here. By fitting global models for the entire X-ray spectrum, we can infer element abundances and the distribution of emission measure with temperature. By simulating the actual 2-dimensional images in the dispersed spectra (e.g., with the MARX software), we can refine our models of the kinematics and spatial distribution of the shocked gas. By constructing simulations of the hydrodynamics of the impact of the blast wave with the protrusions, we hope to provide a more refined interpretation of the observations than the analysis presented here.
This work was supported by NASA through Chandra Awards G04-5072A (to CU, Boulder, CO) and GO4-5072B (to NCSU, Raleigh, NC). The authors appreciate the careful reading and comments by an anonymous referee.
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# Formation of nematic liquid crystalline phase of F-actin varies from continuous to biphasic transition
## Abstract
We show that the isotropic to nematic liquid crystalline phase transition of F-actin can be either continuous or discontinuous, depending critically on the filament length. For F-actin with average filament length $`3\mu m`$, we confirm that the transition is continuous in both filament alignment and local concentration. In contrast, for filament length $`2\mu m`$ the F-actin solution undergoes a first order transition. Tactoidal droplets of co-existing isotropic and nematic domains were observed. Phenomena of nucleation-and-growth and spinodal decomposition both occur, depending sensitively on the exact concentration and average filament length of F-actin. In the late stage, the tactoidal droplets continually grow and occasionally coalesce to form larger granules.
Cytoskeletal protein actin is responsible for cell morphology and motility Alberts and et al. (2002). Globular actin (G-actin) polymerizes to form long filaments, F-actin. F-actin has a diameter of 8 nm Holmes et al. (1990) and a distribution of lengths characteristic of the stochastic polymerization process Sept et al. (1999). F-actin is a semiflexible polymer with a persistence length of 15-18 $`\mu m`$ Isambert et al. (1995); Gittes and et al. (1993), which is larger than their average length in cells or in vitro. There have been extensive studies concerning many remarkable properties of F-actin, including dynamic filament assembly and dissembly Fujiwara et al. (2002), phase transitions Coppin and Leavis (1992); Furakawa and et al. (1993); Viamontes and Tang (2003), and rheology Gardel et al. (2003). Many of these properties are shared by other self-assembled protein filaments such as microtubules and collagen-based intracellular matrix Alberts and et al. (2002).
Of particular relevance to this report is that F-actin undergoes an isotropic (I) to nematic (N) liquid crystalline phase transition. The onset concentration of the transition is inversely proportional to the average filament length $`\mathrm{}`$ Coppin and Leavis (1992); Furakawa and et al. (1993); Suzuki and et al. (1991); Viamontes and Tang (2003), consistent with statistical mechanical theories Onsager (1949); Flory (1969). The experimental studies use optical birefringence methods to measure the F-actin alignment across the I-N transition region. Under certain preparation conditions, ”zebra” birefringence patterns were observed, which have been attributed to the spontaneous separation of F-actin into I and N domains Suzuki and et al. (1991). Two more recent studies show Coppin and Leavis (1992); Viamontes and Tang (2003), however, that the F-actin I-N transition appears to be continuous in both filament alignment and concentration for $`\mathrm{}3\mu m`$. It has been argued that perhaps due to the extreme filament length, polydispersity, and semi flexibility, a combined outcome of defect suppression and entanglement renders the F-actin I-N transition continuous Viamontes and Tang (2003). The phenomenon may be relevant to the theory of Lammert, Rokshar, and Toner (LRT) Lammert et al. (1993, 1995), which predicts that a high disclination line defect energy renders the I-N transition into two continuous ones.
In this paper we confirm the continuous features of the I-N transition for solutions of long F-actin ($`\mathrm{}3\mu m`$), but also show for the first time the I-N co-existence and domain separation of F-actin for solutions with $`\mathrm{}2\mu m`$. Tactoidal droplets were observed in at least three types: N tactoids in an I background, I tactoids in an N background, coexistence of N and I tactoids on a uniform backgroud of weak alignment. The droplets grow in two distinct ways: 1. nucleation of scattered tactoids and growth; 2. spinodal decomposition, followed by coasening. In the late stage for both cases, coalescence of tactoidal droplets were observed and studied. Slow axis measurements of the N tactoids suggest that the director field smoothly follows the surface contour, connecting point defects at two opposite poles, which are called boojums Prinsen and van der Schoot (2003); Drzaic (1995)
G-actin was extracted from rabbit skeletal muscle following an established method Pardee and Spudich (1982). F-actin $`\mathrm{}`$ was varied by addition of gelsolin, a filament severing and end-capping protein Janmey et al. (1986); Tang and Janmey (1996). G-actin was polymerized upon addition of KCl and MgCl<sub>2</sub> upto 50 mM and 2 mM, respectively. Rectangular capillary tubes from VitroCom Inc. (Mt. Lks., N.J) of crossectional dimensions 0.2$`\times `$2 mm were used for measurements by fluorescence and birefringence microscopy. Both ends of the capillary tube were sealed by an inert glue to eliminate flow. Birefringence measurements were performed on a Nikon E-800 microscope equipped with the CRI PolScope package Viamontes and Tang (2003). PolScope is capable of measuring the optical birefringence and the direction of slow axis at each pixel position, thus reporting local alignment Oldenbourg and Mei (1995); Shribak and Oldenbourg (2003). F-actin was labeled 1 to 1000 with TRITC-Phalloidin (Sigma, St Louis, MO) for fluorescence measurements, performed as previously described Viamontes and Tang (2003). 2D Fast Fourier Transform (2D-FFT) analysis was performed similarly to Bees and Hill Bees and Hill (1997), using the MatLab 7.0 software (The MathWorks, Inc.).
Different features are observed between samples of several $`\mathrm{}`$ of F-actin in their respective ranges of concentration over which the I-N transition occurs. F-actin with no gelsolin added were measured to be of $`\mathrm{}=11\mu m`$. Fig. 1 shows representative results of birefringence and filament alignment of F-actin in the I phase (Fig. 1A), transition region (Fig. 1B), and the N phase (Fig. 1C). Of particular note is that in the I-N transition region uniform retardance is found, suggesting that F-actin is continuous in alignment (Fig. 1B) and in concentration Viamontes and Tang (2003). Zebra patterns are occationally observed, especially near the wall of a capillary tube, or at an air liquid interface, examples of which have been shown by previous studies Coppin and Leavis (1992); Suzuki and et al. (1991); Viamontes and Tang (2003). Even at the location of a zebra pattern, the local concentration of actin remains uniform, suggesting a lack of co-existence in long F-actin samples Viamontes and Tang (2003). In contrast, upon further reduction of $`\mathrm{}`$ to $`2\mu m`$ the F-actin solution phase separates into tactoidal droplets and their surrounding medium (Fig. 1D). An increase in concentration gives rise to co-existence of I and N tactoids (Fig. 1E), and I tactoids in N background (Fig. 1F). Local concentration of F-actin is measured by quantitative fluorescence labelling, which shows that N tactoids are denser in protein than the surrounding background, but only by 20-30 %. The actin concentration within the I tactoids is lower than the sorrunding N region by a similar percentage. These results comfirm the weakly first order nature of the I-N transition.
We have measured the average alignment of F-actin at four $`\mathrm{}`$ as a function of actin concentration over the range of I-N transition (Fig. 2). $`\mathrm{}`$ was determined by fluorescence imaging or AFM (for the shortest $`\mathrm{}`$) of single F-actin for at least 500 filaments for each $`\mathrm{}`$. Below a threshold concentration F-actin solution is in the I phase, thus the retardance is zero (Region A). As the concentration increases, the solution reaches the I-N transition region, characterized by the sharp increase of specific retardance (Region B). In the high concentration region, F-actin solution is completely in the N phase (Region C). As $`\mathrm{}`$ decreases, the onset concentration of the I-N transition increases, consistent with the earlier reports Coppin and Leavis (1992); Furakawa and et al. (1993); Suzuki and et al. (1991); Viamontes and Tang (2003). D,E,F on Fig. 2 indicate the I-N transition region for $`\mathrm{}=1\mu m`$, where the specific birefringence values were measured prior to phase separation as shown in Fig. 1.
When actin is polymerized at slightly above or below the co-existing I or N concentrations, nascent tactoids are nucleated in the transition region for F-actin with $`\mathrm{}2\mu m`$. Coalescence of tactoidal droplets is an efficient form of growth. Fig. 3 shows how N tactoids coalesce, viewed under a polarization microscope. Initially the tactoids are separated from each other. Upon growth, two tactoids reach close proximity and fuse. Once two tactoids have coalesced the final shape is again a similar tactoid. The process repeats itself as the 3rd tactoid coalesces to form one final tactoid. One charactistic of such a large molecular system is its slow dynamics, as the sequence in display took nearly 3 hours.
When the actin concentration is close to midway between the co-existing I and N concentrations (both are sensitive functions of $`\mathrm{}`$), we reproducibly observed the phenomenon of spinodal decomposition. At the high concentration of several mg/ml, F-actin polymerization occurs within seconds following the addition of KCl and MgCl<sub>2</sub> Carlier et al. (1984). Therefore, the solution becomes weakly aligned by the shear flow as it is injected into the capillary tube. Fig. 4A-D show a time sequence of tactoidal growth from initiation of actin polymerization to the formation of large tactoids. Fig. 4A represents the actin solution immediately after it was prepared for observation, which took about 30 seconds. A granular structure appears throughout the capillary tube within minutes after initiation of polymerization. The characteristic domain size is determined by 2D-FFT and is $`17\mu m\times 22\mu m`$ (see inset on the bottom plot). The domain size appears to be nearly constant within the first 15 or 30 min for length and width. The dominant size of the spinodal growth can be explained by adding a density gradient term to the free energy of the system following the classical Cahn-Hilliard treatment Jones (2002). Based on the observed peak wave factor $`q_m=1/\lambda `$, we predict that the thickness of the I-N interface $`\xi \lambda /4\pi =1.4\mu m`$ Larson (1999), assuming that filaments tend to align parallel to the interface. By the 30 min time point, both I and N droplets are discernable (Fig.4B). The late stage growth lasts for many hours. For instance, Fig.4C represents the time point when coalescence has become the main form of tactoidal growth. The large tactoids in Fig. 4D are formed by both continuous growth and occasional coalescence of smaller tactoids. The plot in Fig 4 shows progression of the radial peak in the Fourier spectrum density, whereas the tactoidal length and width shown in the inset were obtained by analyzing line plots along the long and short axes separately.
It is important to note that the actin tactoids we report here are fundamentally different from what we recently reported Tang et al. (2005). Actin granules reported in our recent study were induced by an actin crosslinking protein, alpha-actinin. As a result, the density of actin in the published actin granules is over 10 times higher than in this study. The dense actin granules observed in our recent report are clearly discernable with phase contrast imaging, which is not the case for either I or N droplets reported here. Also, in the present study, we have not observed variant shapes such as triangular tactoids. The tactoids formed due to I-N separation easily disappear upon dilution, unlike the permenant tactoids due to crosslinking of alpha-actinin. Nevertheless, it is remarkable that in both cases the much similar tactoidal shape prodominates. We argue that the tactoidal shape is general among granules consisting of long and stiff filaments, which is dictated by the minimization of the surface energy while accomodating the long and stiff constituent filaments. It is thus not surprising that tactoids of similar shapes are found in two distinct types of actin granules, and in concentrated suspension of Tobacco Mosaic Virus (TMV) Freundlich (1938) and filamentous phage fd Dogic and et al. (2004), as well.
In conclusion, we have shown new features of I-N phase transition of F-actin solutions as a function of $`\mathrm{}`$. At $`\mathrm{}3\mu m`$, the I-N phase transition is continuous, consistent with previous findings Viamontes and Tang (2003); Coppin and Leavis (1992). However, biphasic behavior characteristic of a first order transition is observed for $`\mathrm{}2\mu m`$, including both phenomena of nucleation-growth and spinodal decomposition. Tactoidal droplets of either I or N domains form in the N or I background. Tactoids of both phases can also co-exist with a weakly aligned background state, suggesting slow kinetics and metastability. The process for tactoidal growth involves both constant recruitment of the surrounding filaments and coalescence of existing neighboring tactoids.
This work is supported by the National Science Foundation (NSF DMR 0405156) and the Petroleum Research Fund, administered by the American Chemical Society. We thank Professors Robert Meyer, Robert Pelcovits, Tom Powers and Jim Valles for valuable suggestions.
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# Brane World Cosmology In Jordan-Brans-Dicke Theory
## Introduction
In spite of the success of general relativity now called the standard theory of gravitation, there are many other alternative theories. Among them the scalar tensor theory is the most important one. The scalar-tensor theory was conceived originally by P. Jordan who started to embed a four dimensional curved manifold in five dimensional flat space-time. He presented a general Lagrangian for the scalar field living in four-dimensional curved space-time:
$$L_J=\sqrt{g}\left[\phi _J^\gamma \left(R\omega _J\frac{1}{\phi _J^2}g^{\mu \nu }_\mu \phi _J_\nu \phi _J\right)+L_{matter}(\phi _J,\mathrm{\Psi })\right],$$
(1)
where $`\phi _J\left(x\right)`$ is Jordan’s scalar field, $`\gamma `$ and $`\omega _J`$ are constants, and $`\mathrm{\Psi }`$ represents matter fields. $`\phi _J^\gamma R`$ is the nonminimal coupling term which marked the birth of scalar-tensor theory.
Jordan’s work was taken over particularly by C.Brans and R.H. Dicke . They assumed that decoupling of scalar field from the matter part of the Lagrangian occurs. They defined their scalar field $`\phi `$ by
$$\phi =\phi _J^\gamma $$
(2)
and then Lagrangian will be
$$L_{BD}=\sqrt{g}\left(\phi R\omega \frac{1}{\phi }g^{\mu \nu }_\mu \phi _\nu \phi +L_{matter}\left(\mathrm{\Psi }\right)\right).$$
(3)
They demanded that the matter part of the Lagrangian $`\sqrt{g}L_{matter}`$ be decoupled from $`\phi \left(x\right)`$ as their requirement that the weak equivalence principle be respected, in contrast to Jordan’s model . To remove the singularity from the second term on the right hand side we introduce a new field $`\varphi `$:
$$\phi =\frac{\varphi ^2}{8\omega }.$$
(4)
Then the Brans-Dicke (B-D) action will be
$$L_{BD}=\sqrt{g}\left(\frac{\varphi ^2}{8\omega }R\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi +L_{matter}\right).$$
(5)
where in order to get cosmic acceleration, either the parameter $`\omega `$ should be time dependent , or a potential term for the scalar field could be added to the Lagrangian.
On the other hand, string theory predicts a new type of nonlinear structure, which is called a brane, a word created from ”membrane”. This also gives a new perspective to cosmology so that our universe is confined to a four dimensional space-time subspace or 3-brane. The extra dimension may have large compact toroidal topology or be unbounded with a warp factor, depending on the distance from the brane . Additionally several works have studied higher dimensional B-D theory to combine the advantages of both the five dimensional cosmology and the B-D theory . Moreover, considering the scalar field in the five dimensional bulk with Einstein gravity was proposed by many works .
Our starting point is the paper of Bander who studied five dimensional bulk whose dynamics is governed by a scalar Liouville field coupled to gravity in the usual way . Then he derived that the effective theory on the brane has a time dependent Planck mass and cosmological constant and also found expanding scale factors with no acceleration. In this paper we investigate the properties of the five dimensional bulk in Brans-Dicke theory. The layout of our paper is as follows. In section 2 we present the general framework for our five dimensional theory and compute the five dimensional Brans-Dicke equations. In section 3 we analyze the cosmological solutions. We find false vacuum energy ($`p_B=\rho _B`$) for exponentially growing scale factors and radiation dominated universe ($`p_B=\frac{1}{3}\rho _B`$) for power law scale factors in the bulk. In section 4 we derive the effective four dimensional scalar field and obtain its time dependence. Finally, we sum up our results and conclusions in section 5.
## 1 The Action and Equations of Motion
In this work we look at the five dimensional Brans-Dicke action:
$$S=d^5x\sqrt{g}\left(\frac{\varphi ^2}{8\omega }R\frac{1}{2}_A\varphi _B\varphi g^{AB}V\left(\varphi \right)\right),$$
(6)
where $`\omega `$ is the dimensionless Brans-Dicke parameter, $`\varphi `$ is the scalar field and $`V\left(\varphi \right)`$ is the scalar potential. The variation of the action with respect to $`g^{AB}`$ gives
$$\frac{1}{8\omega }\left(\varphi ^2G_{AB}\varphi _{,A;B}^2+g_{AB}\mathrm{}\varphi ^2\right)\frac{1}{2}_A\varphi _B\varphi +\frac{g_{AB}}{4}_C\varphi ^C\varphi +\frac{1}{2}g_{AB}V\left(\varphi \right)=T_{AB}.$$
(7)
We choose a general five dimensional metric anzats which can be written in an orthonormal basis as :
$$ds^2=b\left(t\right)^2dW^2+f\left(W\right)^2\left[dt^2+a\left(t\right)^2\delta _{ij}dx^idx^j\right],$$
(8)
where $`i,j=1,2,3`$ and $`f\left(W\right)`$ is the warp factor which depends on the fifth coordinate, $`a\left(t\right)`$ is the cosmological scale factor and $`b\left(t\right)`$ is the time dependent scale factor of the fifth dimension. More generally, this metric has been studied in papers . In Mendes’ work the five dimensional brane cosmology with non-minimally coupled scalar field to gravity is interpreted in Jordan frame without a scalar potential. In our work we add a scalar potential to the action.
In the orthonormal basis $`e^0=fdt,`$ $`e^i=fadx^i`$ and $`e^5=bdW,`$ the stress-energy tensor can be considered as
$$T_B^A=T_B^A_{bulk}+T_B^A_{brane},$$
(9)
where $`T_B^A_{bulk}`$is the energy momentum tensor of the bulk matter and
$$T_B^A_{bulk}=\text{diag}(\rho _B,p_B,p_B,p_B,q_B).$$
The second term corresponds to the matter content in the brane $`\left(W=0\right),`$
$$T_B^A_{brane}=\frac{\delta \left(W\right)}{b}\text{diag}(\rho ,p,p,p,0).$$
If we substitute the Einstein tensor components in eq(7) we obtain the B-D equations. In the coordinate basis, for component $`00;`$
$`{\displaystyle \frac{1}{8\omega }}\left(3\left({\displaystyle \frac{\dot{a}^2}{a^2}}+{\displaystyle \frac{\dot{a}\dot{b}}{ab}}\right){\displaystyle \frac{3f^2}{b^2}}\left({\displaystyle \frac{\stackrel{´}{f}^2}{f^2}}+{\displaystyle \frac{\text{}}{f}}\right)+3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{_t\varphi ^2}{\varphi ^2}}+{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{_t\varphi ^2}{\varphi ^2}}{\displaystyle \frac{f^2}{b^2}}\left(3{\displaystyle \frac{\stackrel{´}{f}}{f}}{\displaystyle \frac{_W\varphi ^2}{\varphi ^2}}+{\displaystyle \frac{_W^2\varphi ^2}{\varphi ^2}}\right)\right)`$ (10)
$`{\displaystyle \frac{f^2}{4b^2}}{\displaystyle \frac{\left(_W\varphi \right)^2}{\varphi ^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\left(_t\varphi \right)^2}{\varphi ^2}}{\displaystyle \frac{f^2}{2}}{\displaystyle \frac{V\left(\varphi \right)}{\varphi ^2}}={\displaystyle \frac{T_{00}}{\varphi ^2}}.`$
For components $`ii;`$
$`{\displaystyle \frac{1}{8\omega }}\left(\left({\displaystyle \frac{2\ddot{a}}{a}}+{\displaystyle \frac{\dot{a}^2}{a^2}}+2{\displaystyle \frac{\dot{a}\dot{b}}{ab}}+{\displaystyle \frac{\ddot{b}}{b}}\right)+{\displaystyle \frac{3f^2}{b^2}}\left({\displaystyle \frac{\stackrel{´}{f}^2}{f^2}}+{\displaystyle \frac{\text{}}{f}}\right){\displaystyle \frac{_t^2\varphi ^2}{\varphi ^2}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{_t\varphi ^2}{\varphi ^2}}2{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{_t\varphi ^2}{\varphi ^2}}\right)`$ (11)
$`+{\displaystyle \frac{1}{8\omega }}{\displaystyle \frac{f^2}{b^2}}\left({\displaystyle \frac{_W^2\varphi ^2}{\varphi ^2}}+3{\displaystyle \frac{\stackrel{´}{f}}{f}}{\displaystyle \frac{_W\varphi ^2}{\varphi ^2}}\right){\displaystyle \frac{1}{4}}{\displaystyle \frac{\left(_t\varphi \right)^2}{\varphi ^2}}+{\displaystyle \frac{f^2}{4b^2}}{\displaystyle \frac{\left(_W\varphi \right)^2}{\varphi ^2}}+{\displaystyle \frac{f^2}{2}}{\displaystyle \frac{V\left(\varphi \right)}{\varphi ^2}}={\displaystyle \frac{1}{a^2}}{\displaystyle \frac{T_{ii}}{\varphi ^2}}.`$
For component $`55;`$
$`{\displaystyle \frac{1}{8\omega }}\left(3\left({\displaystyle \frac{\dot{a}^2}{a^2}}+{\displaystyle \frac{\ddot{a}}{a}}\right)+6{\displaystyle \frac{f^2}{b^2}}{\displaystyle \frac{\stackrel{´}{f}^2}{f^2}}{\displaystyle \frac{_t^2\varphi ^2}{\varphi ^2}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{_t\varphi ^2}{\varphi ^2}}+{\displaystyle \frac{4f^2}{b^2}}{\displaystyle \frac{\stackrel{´}{f}}{f}}{\displaystyle \frac{_W\varphi ^2}{\varphi ^2}}\right)`$ (12)
$`{\displaystyle \frac{f^2}{4b^2}}{\displaystyle \frac{\left(_W\varphi \right)^2}{\varphi ^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{\left(_t\varphi \right)^2}{\varphi ^2}}+{\displaystyle \frac{f^2}{2}}{\displaystyle \frac{V\left(\varphi \right)}{\varphi ^2}}={\displaystyle \frac{f^2}{b^2}}{\displaystyle \frac{T_{55}}{\varphi ^2}}.`$
For component $`05;`$
$$\frac{1}{8\omega }\left(\frac{3\dot{b}\stackrel{´}{f}}{bf}\frac{_t_W\varphi ^2}{\varphi ^2}+\frac{\stackrel{´}{f}}{f}\frac{_t\varphi ^2}{\varphi ^2}+\frac{\dot{b}}{b}\frac{_W\varphi ^2}{\varphi ^2}\right)\frac{1}{2}\frac{_t\varphi _W\varphi }{\varphi ^2}=0.$$
(13)
Assume that the $`05`$ component of the energy-momentum tensor vanishes, which means that there is no flow of matter along the fifth dimension. Therefore the nonzero elements of the 5D stress-energy tensor are
$`T_{00}`$ $`=`$ $`f^2\rho _B+f^2{\displaystyle \frac{\delta \left(w\right)}{b}}\rho `$ (14)
$`T_{ii}`$ $`=`$ $`a^2f^2\rho _B+a^2f^2{\displaystyle \frac{\delta \left(w\right)}{b}}p`$
$`T_{55}`$ $`=`$ $`b^2q_B.`$
Finally variation with respect to $`\varphi `$ gives,
$$\frac{1}{4\omega }\left(\varphi R\right)\frac{V\left(\varphi \right)}{\varphi }+\mathrm{}\varphi =0,$$
which explictly reads
$$\frac{1}{4\omega }R\frac{_t^2\varphi }{f^2\varphi }+\frac{4}{b^2}\frac{\stackrel{´}{f}}{f}\frac{_W\varphi }{\varphi }\frac{3}{f^2}\frac{\dot{a}}{a}\frac{_t\varphi }{\varphi }\frac{\dot{b}}{bf^2}\frac{_t\varphi }{\varphi }+\frac{_W^2\varphi }{b^2\varphi }\frac{1}{\varphi }\frac{V\left(\varphi \right)}{\varphi }=0,$$
(15)
where the Ricci scalar $`R`$ is:
$$R=\frac{1}{f^2}\left(\frac{6\ddot{a}}{a}+\frac{2\ddot{b}}{b}+\frac{6\dot{a}^2}{a^2}+\frac{6\dot{a}\dot{b}}{ab}\right)\frac{12\stackrel{´}{f}^2}{f^2b^2}\frac{8f^{\prime \prime }}{fb^2}.$$
The metric and the B-D field are continuous across the brane localized at $`W=0`$. However their derivatives can be discontinuous at the brane. Since we have orbifold symmetry, second derivatives of scale factor and B-D field will contain Dirac delta function in the second derivatives of the metric with respect to fifth dimension. Therefore for a function $`f,`$ we have
$$f^{\prime \prime }=\widehat{f^{\prime \prime }}+\left[f^{}\right]\delta \left(W\right),$$
where $`\widehat{f^{\prime \prime }}`$ is the non-distributional part of the double derivative of $`f`$, and $`\left[f^{}\right]`$ is the jump in the first derivative of $`f`$ across $`W=0`$, it is defined as
$$\left[f^{}\right]=f^{}\left(0^+\right)f^{}\left(0^{}\right).$$
Matching the Dirac delta functions in equations (10), (11) and (15) we obtain that
$`{\displaystyle \frac{\left[f^{}\right]_0}{f_0b_0}}`$ $`=`$ $`{\displaystyle \frac{8\omega ^2}{\left(3\omega +4\right)\varphi ^2}}\rho `$ (16)
$`{\displaystyle \frac{\left[\varphi ^{}\right]_0}{\varphi _0b_0}}`$ $`=`$ $`{\displaystyle \frac{16\omega }{\left(3\omega +4\right)\varphi ^2}}\rho `$
where the subscript ‘$`0`$’ stands for the brane at $`W=0.`$ Using eq(10) and eq(11) we get the remarkable result that the cosmological constant dominates on the brane i.e.
$$\rho =p=\frac{\varphi _0^2}{8\omega }\left(\frac{3\left[f^{}\right]_0}{f_0b_0}+\frac{2\left[\varphi ^{}\right]_0}{\varphi _0b_0}\right).$$
(17)
Here choice of the scalar factor $`a\left(t\right)`$ does not make any differences on the equation of state $`\rho =p.`$ Using eq(13) to evaluate the jump condition we get the equation for the matter on the brane
$$4\dot{\rho }+3\omega \frac{\dot{b}}{b}p+6\omega \frac{\dot{\varphi }}{\varphi }\rho =0$$
(18)
where if $`2\frac{\dot{b}}{b}=\frac{\dot{\varphi }}{\varphi }`$ or in particular time derivatives of the $`b`$ and $`\varphi `$ are zero, we obtain that $`\rho `$ and $`p=\rho `$ are constant on the brane.
## 2 Solutions
Solutions of B-D equations restrict the scalar field to be in the form $`\varphi (t,W)=B\left(t\right)C\left(W\right).`$ Starting from this, to satisfy all of the B-D equations we make two possible ansatze for $`a\left(t\right).`$ The first one is exponential growth in time and the other is power law expansion.
### 2.1 <br>Exponential Expansion, $`a\left(t\right)=a_0e^{\lambda t}`$
We see from eq(10-15) that $`b\left(t\right)`$ must be constant, $`b\left(t\right)=b_0`$ and $`B\left(t\right)`$ must be in the exponential form also $`B\left(t\right)=B_0e^{\beta t}`$ and than we can easily read that
$$f\left(W\right)=\frac{W}{W_0}.$$
(19)
For a brane at $`W=W_0`$, we introduce the coordinate $`W^{}`$ such that $`\frac{W}{W_0}=1\frac{W^{}}{W_0}.`$ The metric on both sides of brane can be written
$$ds^2=b_0^2dW^2+\left(1\frac{\left|W\right|}{W_0}\right)^2\left[dt^2+e^{2\lambda t}d\stackrel{}{x}^2\right],$$
(20)
and the brane is at $`W=0`$. Here we dropped the prime for simplicity. This warp factor is the same as in Bander’s work . The brane we live in is embedded in the five-dimensional bulk space-time and the four dimensional part in the square parenthesis is the well known de-Sitter space-time. This metric is similar to a Randall- Sundrum type of model in the same sense. Instead of the exponential warp factor we obtain the linear warp factor. However for small $`W`$ it is known that
$$e^{\left|W\right|}1\left|W\right|,$$
(21)
and two the models are similar.
From equations (10) and (11) it seems to be $`p_B=\rho _B`$, which acts as a cosmological constant. In previous works this energy has been identified as the false vacuum energy density $`\rho _f.`$ During the false vacuum phase the universe supercools. It is believed that as the universe expands it cools down and then it experiences a series of phase transitions. Since the cosmic expansion continues to drive the temperature downward, the universe enters a period of supercooling. As the universe supercools the energy density acts as an effective cosmological constant. Therefore we can consider this stage as the false phase.
For this condition we get the results
$`C\left(W\right)`$ $`=`$ $`c_0\left(1{\displaystyle \frac{\left|W\right|}{W_0}}\right)^\alpha ,`$ (22)
$`\varphi `$ $`=`$ $`B_0c_0\left[e^{\beta t}\left(1{\displaystyle \frac{\left|W\right|}{W_0}}\right)\right]^\alpha ,`$ (23)
here $`\varphi `$ depends only on the distance $`W.`$ Where on the brane we live
$$\left(16\pi G\right)^1=M_p^2=\frac{\varphi ^2}{8\omega }=\frac{\left(B_0c_0\right)^2e^{2\alpha \beta }}{8\omega },$$
(24)
where $`B_0c_0`$ is required to be within a few orders of magnitude of Planck mass . We first discuss the solution where $`T_{55}=0.`$
#### 2.1.1 $`T_{55}=0:`$
From the B-D equations this condition causes $`\rho _B=p_B=0`$ (empty universe). Then solutions are very simple
| $`\rho _B=0`$ | $`p_B=0`$ | $`V_0=0`$ | $`\beta =\lambda `$ | $`\alpha =\frac{1}{2\left(1+\omega \right)}`$ |
| --- | --- | --- | --- | --- |
| | | $`V_0=\frac{\left(3\omega +4\right)\lambda ^2}{2\omega \left(1+\omega \right)^2}\left(B_0c_0\right)^{2/\alpha }`$ | $`\beta =0`$ | $`\alpha =\frac{1}{1+\omega }`$ |
where $`b_oW_o\lambda =\pm 1.`$ Here in the first row of the table the B-D equations give a scalar field which depends not only on time but also on the fifth coordinate. On the other hand in the second row the scalar field only depends on the fifth coordinate and there is a scalar potential $`V_00`$. Therefore the scalar potential is not depend on the time
$$V\left(\varphi \right)=V_0\varphi ^{2\frac{2}{\alpha }},$$
where $`V_0`$ is a constant has dimension $`L^{2\frac{3}{\alpha }}`$. From these results as $`\omega \mathrm{},`$ $`\alpha ,V_00.`$ Therefore $`V\left(\varphi \right)0.`$ This means that at the large values of the B-D parameter, the scalar field is constant $`\frac{\varphi ^2}{8\omega }=M_p^2=\frac{\left(B_0c_0\right)^2}{8\omega }`$ with no scalar potential.
The $`q=0`$ condition has been derived in where it was found that empty and flat five dimensional universe where $`{}_{}{}^{\left(5\right)}R_{PQ}^{MN}=0`$ and $`\mathrm{\Lambda }_5=0`$ gives rise to a four dimensional expanding universe with nonzero Riemann tensor and cosmological constant. This five dimensional space is a well known Minkowski universe
$$ds^2=dx_1^2+dx_2^2+dx_3^2+dx_4^2+dx_5^2$$
(25)
transformed into
$$ds^2=b_0^2dW^2+\left(1\frac{\left|W\right|}{W_0}\right)^2\left[dt^2+e^{2\lambda t}\left(dr^2+r^2d\mathrm{\Omega }_2^2\right)\right],$$
(26)
by the following transformation
$`x_1`$ $`=`$ $`b_0\left(W_0\left|W\right|\right)\left(\mathrm{sinh}\left(\lambda t\right)+{\displaystyle \frac{\lambda ^2r^2}{2}}e^{\lambda t}\right)`$ (27)
$`x_2`$ $`=`$ $`b_0\left(W_0\left|W\right|\right)\left(\mathrm{cosh}\left(\lambda t\right){\displaystyle \frac{\lambda ^2r^2}{2}}e^{\lambda t}\right)`$
$`x_3`$ $`=`$ $`b_0\left(W_0\left|W\right|\right)\lambda re^{\lambda t}\mathrm{cos}\theta `$
$`x_4`$ $`=`$ $`b_0\left(W_0\left|W\right|\right)\lambda re^{\lambda t}\mathrm{sin}\theta \mathrm{cos}\phi `$
$`x_5`$ $`=`$ $`b_0\left(W_0\left|W\right|\right)\lambda re^{\lambda t}\mathrm{sin}\theta \mathrm{sin}\phi ,`$
after some calculations we get the factor $`b_0W_0\lambda `$ in front of the four dimensional part. This was already found as unity. Therefore the four dimensional curved space time can be embedded in the five dimensional flat space time by these coordinate transformations.
#### 2.1.2 $`T_{55}0:`$
From the B-D equations we obtain that $`p_B=\rho _B0`$ and $`\beta =0.`$ As $`\omega \mathrm{},`$
$`\rho _B`$ $`=`$ $`p_B{\displaystyle \frac{\left(B_0c_0\right)^{2/\alpha }}{2\left(b_ow_o\right)^2}}{\displaystyle \frac{\alpha ^2\left(\alpha +1\right)}{\left(\alpha 1\right)}}\varphi ^{22/\alpha }`$ (28)
$`q_B`$ $``$ $`{\displaystyle \frac{\left(B_0c_0\right)^{2/\alpha }}{\left(b_0W_o\right)^2}}{\displaystyle \frac{\alpha ^2}{\alpha 1}}\varphi ^{22/\alpha }`$
$`V_0`$ $``$ $`{\displaystyle \frac{\left(B_0c_0\right)^{2/\alpha }}{\left(b_0W_0\right)^2}}{\displaystyle \frac{\alpha ^2\left(3+\alpha \right)}{2\left(\alpha 1\right)}}.`$
for all of the results $`\beta =0`$ and $`V\left(\varphi \right)=V_0\varphi ^{2\frac{2}{\alpha }},`$ therefore the scalar potential becomes again time independent for the exponentially expanding universe for $`q_B0`$.
If we suppose this phase as the false phase, the probability of a point remaining in the false phase during the bubble nucleation process is quite small as shown in . Then the universe is dominated by the true vacuum and exits from the false vacuum.
In the true vacuum we can consider a power-law expansion.
### 2.2 Power-law Expansion:
The scale factors are:
$`a\left(t\right)`$ $`=`$ $`a_0\left(t/t_0\right)^\lambda ,`$ (29)
$`b\left(t\right)`$ $`=`$ $`b_0\left(t/t_0\right)^\gamma .`$ (30)
These power law solutions restrict us to choose $`B\left(t\right)=B_0\left(\frac{t}{t_0}\right)^\beta .`$ On the other hand B-D equations is satisfied only if $`\gamma ,\beta =1,`$ and again we get same result for the warp factor, $`f\left(W\right)=\left(1\frac{\left|W\right|}{W_0}\right).`$ Then these results causes scalar field to be
$$\varphi (t,W)=B_0c_0\left[\left(\frac{t}{t_0}\right)\left(1\frac{\left|W\right|}{W_0}\right)\right]^\alpha ,$$
(31)
where $`B_0`$ and $`c_0`$ are constants. Again to satisfy the B-D equations, we find the similar scalar potential;
$$V\left(\varphi \right)=V_0\varphi ^{\frac{2}{\alpha }\left(\alpha 1\right)},$$
(32)
and $`B_0c_0`$ has dimension $`L^{3/2}`$, therefore $`V_0`$ has dimension $`L^{3/\alpha 2}`$. Here to make B-D equations simpler we set $`\frac{b_0W_0}{t_0}=1`$.
Now we want to find a general result so we consider the equation of state as:
$$p_B=\nu \rho _B.$$
(33)
Putting all of these settings in the B-D equations we find a nice result: here the interesting thing is that there is no solution other than $`\nu =\frac{1}{3}`$ for $`p_B0`$ and $`\rho _B0`$ and solutions are valid only for $`q_B=0.`$ Different values of the variables in equations (10-15) are satisfied only for a single value of $`\nu `$ which is $`\frac{1}{3}.`$ Then this ratio between the pressure and energy density corresponds to the radiation dominated universe; and $`\omega `$ dependence of $`\lambda ,\alpha ,`$ and $`V_0`$ are:
$`\rho _B`$ $`=`$ $`3p_B`$ (34)
$`\alpha _\pm `$ $`=`$ $`{\displaystyle \frac{\pm \sqrt{3\omega +4}+1}{2\left(\omega +1\right)}},`$ (35)
$`\lambda _\pm `$ $`=`$ $`{\displaystyle \frac{\omega \sqrt{3\omega +4}}{4\left(\omega +1\right)}}`$ (36)
and finally
$$V_{0\pm }=\frac{3\left(B_0c_0\right)^{2/\alpha }}{t_0^2}\frac{\left(3\omega +4\right)\left(3\omega \pm \sqrt{3\omega +4}+5\right)}{32\omega \left(\omega +1\right)^2}.$$
(37)
All of these solutions do not give a specific value for $`\omega .`$ From the time-delay measurements, experimentally $`\omega >500`$ and more recently $`\omega >3000`$ . As $`\omega \mathrm{},`$ $`\alpha 0,`$ $`\lambda _\pm \frac{1}{4}`$ and $`V_00.`$ This means that at this limit, the scalar field becomes constant and the scalar potential vanishes.
For the power-law scale factor we obtain one more solution B-D equations give the empty universe, namely $`\rho _B,`$ $`p_B=0`$ and $`\lambda =1,V_0=0`$ and $`\alpha _\pm =\frac{\pm \sqrt{3\omega +4}+1}{2\left(\omega +1\right)}`$ which are the same as previous value of $`\alpha `$ (35).
The solutions presented here represent decelerating cosmology for the radiation dominated universe and expanding cosmology with constant velocity for the empty universe. However astronomical observations show that the universe is not only expanding but also undergoing accelerated expansion . It may be possible to obtain power law acceleration in B-D theory if scale factors for external dimensions are time dependent. In string theory some cosmologies can achieve accelerating scale factors .
The metric which we found in this part may be related with the Kasner space-time . It has a cosmological singularity at $`t=0`$ where the square of Riemann tensor diverges. On the brane we live ($`W=0`$)
$$R_{\mu \nu \rho \sigma }R^{\mu \nu \rho \sigma }=\frac{24\lambda ^4}{t^4}.$$
(38)
This is the physical singularity and it cannot be avoided by any coordinate transformation . However, since the central part of the space-time is avoided in orbifold construction this has no importance for the brane world scenario .
## 3 The Effective Four Dimensional Gravitational Constant
Finally we calculate the four-dimensional effective gravitational constant on the brane and compare with the our previous results eq(24). On the left hand side of the action in eq (6) the first term is:
$$d^5x\sqrt{g}\frac{\varphi _{\left(5\right)}^2}{8\omega }R_{\left(5\right)}=d^5x\sqrt{g}M_{\left(5\right)}^3R_{\left(5\right)}=d^5x\sqrt{g}\frac{1}{16\pi G_{\left(5\right)}}R_{\left(5\right)}.$$
(39)
We can perform the $`W`$ integral to obtain the effective gravitational constant. With the same manner in the work, this equation reduces to
$$d^5x\sqrt{g}\frac{\varphi _{\left(5\right)}^2}{8\omega }R_{\left(5\right)}=d^4x𝑑W\sqrt{g^{\left(4\right)}}\frac{\varphi _{\left(5\right)}^2}{8\omega }\left(1\frac{\left|W\right|}{W_0}\right)^2b\left(t\right)R\left(g_{ij}^{\left(4\right)}\left(x\right)\right)$$
(40)
For the exponentially increasing scalar factor we have obtained time dependent scalar field that is $`\varphi _{\left(5\right)}=B_0c_0\left[e^{\beta t}\left(1\frac{\left|W\right|}{W_o}\right)\right]^\alpha `$. Then eq(40) becomes
$$d^5x\sqrt{g}\frac{\varphi _{\left(5\right)}^2}{8\omega }R_{\left(5\right)}=d^4x2\frac{b_0W_0}{8\omega }\frac{B_0^2c_0^2e^{2\alpha \beta t}}{\left(2\alpha +3\right)}\sqrt{g^{\left(4\right)}}R\left(g_{ij}^{\left(4\right)}\left(x\right)\right)$$
(41)
then since $`\alpha 0,`$ the effective gravitational constant becomes
$$\frac{1}{16\pi G_{eff}}=M_{p\left(eff\right)}^2=\frac{\varphi _{\left(4\right)}^2}{8\omega }=\frac{b_0W_0}{12\omega }\left(B_0c_0\right)^2,$$
(42)
which is independent of time. This is similar with the what we have discussed in eq(24) and here $`B_oc_o`$ is within a few orders of Planck mass.
For the power law scalar factors the scalar field is
$`\varphi _{\left(5\right)}=B_0c_0\left(\frac{t}{t_0}\left(1\frac{\left|W\right|}{W_0}\right)\right)^\alpha `$
$$d^5x\sqrt{g}\frac{\varphi _{\left(5\right)}^2}{8\omega }R_{\left(5\right)}=d^4x\frac{b_0W_0}{4\omega }\left(\frac{t}{t_0}\right)^{2\alpha +1}\frac{B_0^2c_0^2}{\left(2\alpha +3\right)}\sqrt{g^{\left(4\right)}}R\left(g_{ij}^{\left(4\right)}\left(x\right)\right),$$
(43)
here again for $`\alpha 0,`$ the four dimensional Brans-Dicke scalar field (or the inverse of the effective gravitational constant) is
$$\frac{1}{16\pi G_{eff}}=M_{p\left(eff\right)}^2=\frac{\varphi _{\left(4\right)}^2}{8\omega }=\frac{W_0}{12\omega }\left(B_0c_0\right)^2b\left(t\right).$$
(44)
Hence the four dimensional effective gravitational constant depends on time .
## 4 Conclusion
In this work we introduced a five dimensional B-D action and studied the five dimensional metric with a warp factor. We showed that the field equations imply a linear warp factor. For an inflating scale factor we found that the energy density acts as an effective cosmological constant. For power law expansion of scale factor we obtained a radiation dominated universe. Additionally we have shown that the five dimensional scalar field is nearly but cannot be exactly constant. On the other hand the four dimensional effective scalar field is constant for exponentially growing scale factor and depends on time for the power law scale factors.
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# Frequency dependence of the drifting subpulses of PSR B0031-07
## 1 Introduction
The single pulses of a pulsar are known to be composed of several smaller units of emission called subpulses. These subpulses are often seen to drift in phase across a sequence of single pulses giving rise to the well-known phenomenon of ‘drifting subpulses’, discovered in 1968 (Drake & Craft 1968). The drift pattern is seen to repeat itself after a given time which is usually denoted by $`P_3`$. The phenomenon has since been detected in many pulsars (e.g. Rankin 1986) and the process is believed to carry information on the mechanism leading to coherent radio emission from pulsars. For example, Ruderman & Sutherland (1975) have suggested a vacuum gap model in which the subpulses correspond to beams of particles (or sparks) produced in the vacuum gap over the polar cap and are thought to rotate around the magnetic axis due to the perpendicular component of the electric field and the magnetic field ($`E\times B`$ drift). Measurements of the speed of rotation of subpulses might therefore give direct information on the electric field in the vacuum gap. Recently, Deshpande & Rankin (1999) have shown that drifting subpulses observed in PSR B0943+10 can be interpreted as 20 sparks rotating around the magnetic axis at a uniform speed.
Some pulsars show clear changes in the vertical spacing between drift bands ($`P_3`$): for example PSR B0809+74 shows a changing $`P_3`$ after it goes through a null (van Leeuwen et al. 2002). PSR B0031$``$07 is particularly interesting because it shows three distinct drift-modes with different $`P_3`$, which are all very stable. They are named mode A, B and C and correspond to a $`P_3`$ of 12, 6 and 4 seconds (or 13, 7 and 4 $`P_1`$), respectively. These values are approximations and from 40,000 pulses observed at 327 MHz Vivekanand & Joshi (1997) found that they may not be harmonically related. They also found that at 327 MHz the relative occurrence rate of these modes are 15.6%, 81.8% and 2.6%, respectively. Furthermore, the pulses occur in clusters containing 30 to 100 pulses which follow each other with delays ranging from fifty to several hundred pulse periods. These clusters are constituted in one of three ways: a series of A bands followed by B bands, only B bands, or a series of B bands followed by C bands (Huguenin et al. 1970; Wright & Fowler 1981). This pulsar also shows a clear presence of Orthogonally Polarised Modes (OPM) in the integrated position angle sweep (Manchester et al. 1975). Table 1 lists some of the known parameters of PSR B0031$``$07. The values for $`\alpha `$ and $`\beta `$ were found by fitting the single vector model from Radhakrishnan & Cooke (1969) to the position angle of the dominant polarisation mode. The $`\alpha `$ and $`\beta `$ values from the position angle of the remaining polarisation mode are the same within errors.
PSR B0031$``$07 has been thoroughly studied at low observing frequencies (Huguenin et al. 1970; Krishnamohan 1980; Wright 1981; Vivekanand 1995; Vivekanand & Joshi 1997, 1999; Joshi & Vivekanand 2000), but only rarely at an observing frequency above 1 GHz (Wright & Fowler 1981; Kuzmin et al. 1986; Izvekova et al. 1993). Wright & Fowler (1981) have observed PSR B0031$``$07 at 1.62 GHz and have found the same drift-modes as seen at lower frequencies. Kuzmin et al. (1986) have studied the integrated pulse profiles of PSR B0031$``$07 at 102.7 MHz, 4.6 GHz and 10.7 GHz. Izvekova et al. (1993) have studied the subpulse characteristics of PSR B0031$``$07 at 62, 102, 406 and 1 412 MHz. They found that the switching between the three drift-modes and the nulls occur simultaneously at all frequencies.<sup>1</sup><sup>1</sup>1However, it is not clear from their paper whether they have sufficient signal to noise at 1414 MHz to see single pulses. In this paper we study the behaviour of the different modes of drift in PSR B0031$``$07 in radio observations at both low and high observing frequencies simultaneously. In Section 2 we explain how the observations have been obtained, how the different modes of drift have been determined, and what further analyses have been carried out. In Section 3 we present our results. The discussion follows in Section 4. In this last section we present a geometrical model which describes many of the observed characteristics of this pulsar.
### 1.1 Definitions
To describe the observational drift of subpulses we use three parameters, which are defined as follows: $`P_3`$ is the spacing at the same pulse phase between drift bands in units of pulsar periods ($`P_1`$); this is the “vertical” spacing when the individual radio profiles obtained during one stellar rotation are stacked as in Fig. 1. $`P_2`$ is the interval between successive subpulses within the same pulse, given in degrees. $`\mathrm{\Delta }\varphi `$, the subpulse phase drift, is the time interval over which a subpulse drifts, given in $`\mathrm{°}`$/$`P_1`$. Note that $`P_2=P_3\times \mathrm{\Delta }\varphi `$.
## 2 Data analysis
The observations of PSR B0031$``$07, were obtained on 3 February 2002 with both the Westerbork Synthesis Radio Telescope (WSRT) and the Effelsberg Radio Telescope simultaneously. These observations were obtained as part of the MFO<sup>2</sup><sup>2</sup>2The MFO collaboration undertakes simultaneous multi-frequency observations with up to seven telescopes at any one time. program. The WSRT observations were made at a frequency of 328 MHz and a bandwidth of 10 MHz. The Effelsberg observations were made at a frequency of 4.85 GHz and a bandwidth of 500 MHz. The time resolutions are 204.8 $`\mu `$s and 500 $`\mu `$s for the 328 MHz and 4.85 MHz observations, respectively. The 328-MHz observations have been corrected for Faraday rotation, dispersion and for an instrumental polarisation effects using a procedure described in the Appendix of (Edwards & Stappers 2004). Also, a 50-Hz signal present in the 4.85-GHz observation has been removed by Fourier transforming the entire sequence, removing the 50 Hz peak and Fourier transforming back. By correlating sequences of single pulses between the 328-MHz and 4.85-GHz observations that contained prominent subpulse drift the pulses could be aligned to within an accuracy of 2$`\mathrm{°}`$ of pulse longitude, which confirms the broadband nature of the pulsar signal. This alignment is sufficient for the studies presented here.
We have also used an observation from PSR B0031$``$07, obtained on 9 August 1999 with the Effelsberg Radio Telescope at a frequency of 1.41 GHz and a bandwidth of 40 MHz. The time resolution is 250 $`\mu `$s.
### 2.1 Calculation of $`P_3`$
To search for periodicities, we considered a sequence of pulses from one of the observations. For each pulse in this sequence, we took the flux at a fixed phase, and calculated the absolute values of the Fourier transform of this flux distribution. This was done for each phase of the pulsar window. The resulting transforms were then averaged over phase, giving a phase-averaged power spectrum (PAPS) from 0 up to 0.5 cycles per rotation period (hereafter c $`P_1^1`$), with a frequency resolution given by the reciprocal of the total length of the sequence.
Initially, all pulses from the observations were divided into sequences of 100 pulses, which were searched for peaks in the PAPS. When a peak was found, the beginning and end of the sequence was adjusted to get the highest signal-to-noise ratio for the peak. The signal-to-noise ratio was calculated as the peak value of the PAPS divided by the rms of the rest of the PAPS. This result was checked by visual inspection of the sequences to see whether they did indeed match the beginning and ending of a drift band. $`P_3`$ was then calculated as the reciprocal of the centre of the peak in the PAPS. When a peak would spread over multiple bins, a cubic spline interpolation was used to determine the location of the peak. The frequency resolution, given by the number of pulses in the sequence, was taken as the error on the position of the peak.
Furthermore, we have calculated the PAPS of all pulses of the 1.41- and 4.85-GHz observations in order to find signs of 6 second periodicity. For comparison, we also calculated the PAPS of the 328-MHz observation.
### 2.2 Values for $`P_2`$
Along with values for $`P_3`$, we also calculated the phase drift for each sequence of pulses. This was done by cross-correlating consecutive pulses. The fluxes in an interval around the peak of the cross-correlation were fitted with a Gaussian, from which the mean was taken as the phase drift. The error of the fit was taken as the error on the phase drift. $`P_2`$ was then determined by multiplying $`P_3`$ by the phase drift.
### 2.3 Average profiles and polarisation properties
To further study these distinct periodicities, we looked at the average-pulse profiles for the individual sequences of pulses that show mode A and mode-B drift, respectively. For the 4.85-GHz observation we compared the sequences that showed mode-A drift in the 4.85-GHz observation with the sequences that showed mode-B drift in the 328-MHz observation. In the same way, we calculated the average linear polarisation, circular polarisation and position angle as a function of pulse phase. We also measured the widths of the average total intensity at 10% and 50% of the peak values and at a height three times the rms. These analysis were done only for the 328-MHz and 4.85-GHz observations.
## 3 Results
The values of $`P_3`$ for each sequence and for both frequencies are shown in Fig. 6. Even though mode B does not seem to occur at 4.85 GHz, it should be noted that whenever mode B is active at 328 MHz there is always radiation present at 4.85 GHz. The same is true for mode C, but there is only one case of a mode C drift. Fig. 6 suggests that the transition between modes can happen within one or a few pulses. An example of how fast the drift-rate can change is shown in Fig. 1. In this figure the pulses are plotted from bottom to top. Mode A is present in the first five drift-bands. Then, within one or two pulses, the drift switches to mode B. In this example, it appears that the transition simply involves the appearance of a drift-band in a different mode, rather than the speeding up of the current drift-band. However, it should also be mentioned that this particular transition happens to occur exactly when a new drift-band would be expected to arise. In our observations there are only a few cases when there is a transition from one drift mode into the other without a null separating the two drift-bands. None of these transitions show a clear change of mode within one drift-band.
The values for $`P_3`$ from the observation at 1.41 GHz, are shown in Fig. 7. We found various examples of mode-A drift, but no mode-B drift. In this respect the 1.41-GHz observation resembles the 4.85-GHz observation. There are no drift-bands showing mode-C drift.
The PAPS of all three observations, are shown in Fig. 2. The low frequencies contain a signal due to the nulling. To make the figure clearer we have set them to zero. We see here that the 6-second periodicity is clearly present at 328-MHz and is just visible in the 1.41-GHz observation. At 4.85 GHz the PAPS does not show the 6-second periodicity. Thus, the mode-B drift gets weaker with increasing observing frequency. However, we did find small sequences in the 4.85-GHz observation where there is a weak 6-second periodicity. Fig. 3 shows the power spectrum of the flux as a function of pulse phase as well as the PAPS of a sequence of 20 pulses from the 4.85-GHz observation containing 6-second periodicity. The PAPS peaks at 6.6 seconds. We did not classify this as a mode-B drift, because the drift-bands are not clearly visible and the cross-correlation between consecutive pulses suggests the drift to be in the opposite direction of all the other drifts. It would be most difficult to explain a mode-B drift-band at high frequency that has a drift-direction different from the drift-direction of the same mode-B drift-band at low frequency. The present sequence is not significant enough to establish that this has occurred.
Table LABEL:tab:P2 shows the average values for $`P_3`$, phase drift and $`P_2`$ for each drift-mode at three frequencies.
Fig. 8 shows the average-polarisation properties of pulses which show the same modes of drift. Each panel shows the total intensity (solid line), linear polarisation (dashed line), circular polarisation (dotted line) and position angle (lower half of each panel). The left panels show the 328-MHz profiles, the right panels show the 4.85-GHz profiles. The top panels show the average polarisation of pulses containing subpulses with mode-A drift, below that are the average polarisation of pulses containing subpulses with mode-B drift and the bottom panels show the average polarisation of all pulses. There is a clear 90$`\mathrm{°}`$ jump in the position angles of all pulses in both the 328-MHz and 4.85-GHz profiles at a longitude of 24$`\mathrm{°}`$. This jump can also be seen in the pulses that only show a mode-B drift. The widths of the average-intensity profiles are listed in Table 3.
## 4 Discussion
We have analysed periodicities in two observations with 2700 pulses of PSR B0031$``$07 which were taken simultaneous at 328 MHz and 4.85 GHz. At low frequency we found that 61.8% of the time the pulsar was in one of the three drift-modes. The occurrence rate was 17.8% for mode A, 80.1% for mode B and 2.1% for mode C. This is consistent with previous results of Vivekanand & Joshi (1997). We have shown that whenever the mode-A drift is active, it is visible at both frequencies. Also, when the mode-B drift is active, it is clearly visible at 328 MHz, but not at 1.41 or 4.85 GHz (see Figs. 6 and 7). However, there is 6 second periodicity in the pulses at 1.41 GHz and a hint of 6 second periodicity in the pulses at 4.85 GHz, the latter of which is possibly drifting in the opposite direction of the drift observed at 328 MHz. This would suggest that towards higher frequency we are seeing less of the drifting subpulses and begin to see a diffuse component that is also subject to the $`E\times B`$ drift. It is difficult to explain a change in the direction of drift towards higher frequency while $`P_3`$ remains almost constant. It might indicate that the observed drift-rate is in fact an alias of the true drift-rate. Establishing and further investigating the possibility of a mode-B drift at 4.85 GHz with a drift direction opposite to that at 328 MHz might help determine the actual drift-rate and direction of the subpulses of this pulsar. The result that only one drift-mode is visible around 1.41 GHz differs from the results in Wright & Fowler (1981) and possibly differs from Izvekova et al. (1993).
In our observations we see that the drift-rate can change within one or two pulses, however we do not see any instance of a mode change within a drift-band. It should be noted that in most cases there is at least one null between two drift-bands with different drift-rates.
Table LABEL:tab:P2 shows that the $`P_2`$ of drift-modes A and B at 328 MHz is almost the same and that drift-mode C has a slightly smaller $`P_2`$ at this frequency. Within errors, however, our values agree with a constant $`P_2`$ for each drift-mode, unlike the behaviour predicted by Vivekanand & Joshi (1997), who claim that $`P_2`$ increases monotonically with $`\mathrm{\Delta }\varphi `$. With increasing frequency, the value of $`P_2`$ for polarisation mode A decreases, which can be explained by the decrease of the opening angle towards higher frequency, assuming radius-to-frequency mapping.
Fig. 8 shows that there is a distinct difference between the average profile of pulses with a 12-second periodicity (A-profile) and with a 6-second periodicity (B-profile). At 328 MHz, the A-profile seems to have two components. This can correspond to the line of sight cutting the edge of the subpulses in the centre of the profile, thereby bifurcating the average profile. It is interesting that the right component in the A-profile at 328 MHz seems to correspond to the single component of the A-profile at 4.85 GHz. Thus it appears as though the first component in the A-profile at 328 MHz disappears towards higher frequency. The profiles at 328 MHz also show that the intensity of the pulses in drift-mode A is on average lower than the intensity of the pulses in drift-mode B. A difference in average profile of the three modes has been reported before by Wright & Fowler (1981) and Vivekanand & Joshi (1997). The former authors have observed the pulsar at 1.62 GHz and found that the A-profile is more narrow than the B-profile, which is in turn more narrow than the C-profile. We cannot directly compare this result with our observation at 4.85 GHz as we do not see a mode-B drift at this frequency. But if we define the B-profile as the pulses at 4.85 GHz that show a mode-B drift at 328 MHz, then we find that at 4.85 GHz, the A-profile is indeed more narrow than the B-profile. Wright & Fowler (1981) do not note a difference in amplitude, nor an offset which are first noted by Vivekanand & Joshi (1997), who have observed the pulsar at 326.5 MHz. They show that at this frequency pulses in drift-mode B have on average more intensity than pulses in drift-mode A and C, which are of equal intensity. They also state that drift-mode A arrives earlier than drift-mode B, which in turn arrives earlier than drift-mode C. Both findings are confirmed by our results from the 328 MHz observation. However, they do not report a double component in the A-profile, which is indeed not present in their plot. This might be due to the fact that only a single linear polarisation was used in their observation. Table 3 shows that the widths from the average intensity profiles decreases from 328 MHz to 1.41 GHz and increases again from 1.41 to 4.85 GHz. This behaviour is not consistent with radius-to-frequency mapping. However, we should note that the signal-to-noise ratio of the edges of the profiles might not be high enough to detect the entire width of the pulse, which makes it difficult to draw any conclusions from these values. Furthermore, when we compare the change in 50%-width from the A-profile with the change in $`P_2`$ between 328 MHz and 4.85 GHz, we find that the value for $`P_2`$ decreases roughly by 65% while the width of the A-profile only decreases by 50%. Both should reflect the change in the size of the radioactive region due to radius-to-frequency mapping. To explain this discrepancy we suggest that at each frequency the drift-path of the subbeams is surrounded by an area of weak radio-emission which becomes relatively smaller with decreasing frequency. To support this claim we have constructed average intensity profiles with contributions from the mode-A drift-bands only. This was achieved by adding up the spectral power in the fluctuation spectrum in the domain between 0.07 and 0.1 Hz at each pulse phase. We then compared the 50%-widths of these profiles at 328 MHz and 4.85 GHz. We found that these widths were approximately 22$`\mathrm{°}`$ and 7$`\mathrm{°}`$, respectively, which is consistent with the change of $`P_2`$ between these frequencies.
Fig. 8 also shows that the average polarisation of all pulses from the 328-MHz observation has two components and a clear minimum around a pulsar phase of 24$`\mathrm{°}`$. From the average polarisation of pulses at 328 MHz that are in drift-mode A and B, it is apparent that the pulses in drift-mode A contribute only to the component on the left and the pulses in drift-mode B contribute only to the component on the right. The average position angle of all pulses shows a 90$`\mathrm{°}`$ jump at both frequencies at a pulsar phase of 24$`\mathrm{°}`$. This can be interpreted as two orthogonally polarised modes changing dominance at this pulsar phase. A straight line fit to the position angles of both modes has shown that these polarisation modes are indeed 90$`\mathrm{°}`$ apart. This jump is also visible in the average position angle of pulses that are in mode-B drift. We have searched for this jump in the average position angle of short sequences containing 20 to 50 pulses in a particular drift-mode and did not find a clear change in polarisation mode. It only manifests itself in the average of many sequences. Furthermore, the pulses in mode-A drift do not show any sign of a jump in position angle. They only show some degree of linear polarisation in the left part of the profile, just before the mode jump occurs in the average position angle of all pulses. Thus the left part of the A-profile is dominated by one of the two orthogonal polarisation modes, while the lack of polarisation in the right part of the A-profile suggests that here both polarisation modes are of equal strength. The lack of polarisation in the left part of the B-profile suggests that here both polarisation modes are also of equal strength, while the right part of the B-profile is dominated by the other polarisation mode. This means that there is a strong relationship between the drift-modes A and B and the two orthogonal modes of polarisation.
### 4.1 Modelling the observations
Since the limited length of our observations does not enable us to study mode C, we shall only attempt to model the behaviour of modes A and B, which are the most prominent drift-modes. We have found that of these two drift-modes only mode A is visible at high frequency, while at low frequency both modes are visible and are seen to have different orthogonal polarisation. At low frequency, the pulses in mode A have less intensity than the pulses in mode B, while at high frequency the pulses show the opposite behaviour. Furthermore, even though the mode B drift is not clearly visible at 4.85 GHz, we do see a hint of 6 seconds periodicity at this frequency. In the context of the potential gap model, the different rates of drift seem to suggest that the potential gap can take on different stable values. Each stable value can be associated with a particular drift rate and a particular magnetic surface of emission if we assume that the radiation is emitted tangential to the magnetic field lines. An increase in the value of the potential gap is expected to give rise to a faster drift and emission (pair production) from magnetic field-lines closer to the magnetic axis. This leads to the picture of an emission region with two concentric radiating rings. At high frequency, the line of sight intersects with magnetic field lines which are further away from the magnetic axis than at low frequency. Therefore, the drifting component of the inner ring is not seen at high frequency. This is graphically illustrated in Fig. 4.
It is also possible to construct a model wherein the emission of both modes comes from the same magnetic flux surface, but from different heights. Such a model is given in van Leeuwen et al. (2003). The difference between the two models is graphically illustrated in Fig. 5. Here, L1, L2 and L3 are the locations where the field lines are directed towards the observer for three different lines of sight: L1 when the line of sight is closest to the magnetic axis, L2 when the line of sight just touches the emission region when observing at high frequency, and L3 shows when the line of sight just touches the emission region when observing at low frequency. In case of a dipole magnetic field, the lines L1 - L3 are straight. Note that these lines do not lie in a plane as might be suggested by Fig. 5. Therefore, one must consider these images as lying in the plane through one of the three lines and the magnetic axis. At points on L1, L2 or L3 further away from the pulsar surface the curvature of the field lines decreases and the frequency of emission is believed to be lower.
In our model (left image) a mode change is caused by a shift of the emission to inner field-lines. At high frequency, the observer can only see emission coming from a region between points P3 and P4. As the mode changes from A to B, the observer will no longer see emission. At low frequency the observer can only see emission coming from a region between the points P1 and P2, and is thus able to see both modes. Note that when the line of sight is closest to the magnetic axis, mode A is not visible, thus causing a dip in the centre of the A-profile at low frequency.
In the model of van Leeuwen et al. (2003) (right image) the emission always comes from the same field-lines. Now however, the emission altitude at a fixed frequency decreases when the mode changes from A to B. This causes the observer to see emission from a point lower in the pulsar magnetosphere (indicated by an open dot and an apostrophe). At high frequency, the observer will only see emission coming from a region between P3 and P4 if the pulsar is emitting in mode A but none in mode B since the region between P3’ and P4’ is not part of the ’active’ (radiating) flux tube. Thus the observer will only see emission in mode A. At low frequency the observer will only see emission coming from a region between P1 and P2 if the pulsar is emitting in mode A and will only see emission in a region between P1’ and P2’ if the pulsar is emitting in mode B. Thus the observer will be able to see emission from both modes. Note that also in this model mode A is not visible at low frequency when the line of sight is closest to the magnetic axis.
In the case of PSR B0031$``$07 both models can explain the observed characteristics. It is therefore difficult to distinguish between them observationally, especially if observing at only one frequency. However, if we were to observe this pulsar at a large range of frequencies (preferably simultaneous) it should be possible to make quantitative statements as to favour one of the models.
On theoretical grounds, we find a possible inconsistency in the model of van Leeuwen et al. (2003) when applied to our observations. In their model, which follows the work of Ruderman & Sutherland (1975) and Melikidze et al. (2000), the transition from mode A to B occurs by a decrease in height of the voltage gap. As a result the altitude of emission decreases. However, the electric field also decreases and with it the speed of the $`E\times B`$ drift, contrary to what is observed. Of course, a possible explanation for this inconsistency could be that we are not seeing the actual drift-speed, but an alias. It should be noted that variation of drift speed can also happen due to temperature variation on the polar cap as suggested by Gil et al. (2003).
## 5 Conclusions
From an analysis of 2700 pulses from PSR B0031$``$07 taken simultaneously at 328 MHz and 4.85 GHz we found that from the three known drift-modes A, B and C of PSR B0031$``$07 only mode A is visible at high frequencies. We have constructed a geometrical model that explains how one drift-mode can disappear at high frequency while another drift-mode remains visible. Further, we have shown that the two most prominent drift-modes A and B are associated with two orthogonal modes of polarisation, respectively. To continue the study presented here, one would require more multi-frequency single pulse observations from this pulsar containing full Stokes parameters.
###### Acknowledgements.
The authors would like to thank J. Gil for his suggestions and discussion towards interpretation of the results and A. Jessner, A. Karastergiou, B. Stappers and our referee, J. Rankin, for their helpful discussions. We also thank all the members of the MFO collaboration for the establishment of the project which led to the observations being available. This paper is based on observations with the 100-m telescope of the MPIfR (Max-Planck-Institut für Radioastronomie) at Effelsberg and the Westerbork Synthesis Radio Telescope (WSRT) and we would like to thank the technical staff and scientists who have been responsible for making these observation possible.
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# FOURIER MUKAI TRANSFORMS FOR GORENSTEIN SCHEMES
## Introduction
Since its introduction by Mukai , the theory of integral functors and Fourier-Mukai transforms have been important tools in the study of the geometry of varieties and moduli spaces. At the first moment, integral functors were used mainly in connection with moduli spaces of sheaves and bundles, and provided new insights in the theory of Picard bundles on abelian varieties and in the theory of stable sheaves on abelian or K3 surfaces . In the relative version they have been also used in mirror symmetry and to produce new instances of stable sheaves on elliptic surfaces or elliptic Calabi-Yau threefolds. The reason is that the theory of integral functors is behind the spectral data constructions ; the irruption of the derived categories in string theory caused by homological mirror symmetry brought then a new interest to derived categories and integral functors (see for recent surveys of the subject and references therein).
Aside from their interest in Physics, derived categories are important geometric invariants of algebraic varieties. Much work is being done in this direction, particularly in the characterisation of all the algebraic varieties sharing the same derived category (also known as Mukai partners).
There are classic results like the theorem of Bondal and Orlov which says that if $`X`$ is a smooth projective variety whose canonical divisor is either ample or anti-ample, then $`X`$ can be reconstructed from its derived category. Mukai proved that there exist non isomorphic abelian varieties and non isomorphic K3 surfaces having equivalent derived categories. Orlov proved that two complex K3 surfaces have equivalent derived categories if and only if the transcendental lattices of their cohomology spaces are Hodge-isometric, a result now called the derived Torelli theorem for K3 surfaces. After Mukai’s work the problem of finding Fourier-Mukai partners has been contemplated by many people. Among them, we can cite Bridgeland-Maciocia and Kawamata ; they have proved that if $`X`$ is a smooth projective surface, then there is a finite number of surfaces $`Y`$ (up to isomorphism) whose derived category is equivalent to the derived category of $`X`$. Kawamata proved that if $`X`$ and $`Y`$ are smooth projective varieties with equivalent derived categories, then $`n=dimX=dimY`$ and if moreover $`\kappa (X)=n`$ (that is, $`X`$ is of general type), then there exist birational morphisms $`f:ZX`$, $`g:ZY`$ such that $`f^{}K_Xg^{}K_Y`$ (i.e. $`D`$-equivalence implies $`K`$-equivalence) . Other important contributions are owed to Bridgeland , who proved that two crepant resolutions of a projective threefold with terminal singularities have equivalent derived categories; therefore, two birational Calabi-Yau threefolds have equivalent derived categories. The proof is based on a careful study of the behaviour of flips and flops under certain integral functors and the construction of the moduli space of perverse point sheaves.
All these results support the belief that derived categories and integral functors could be most useful in the understanding of the minimal model problem in higher dimensions. And this suggests that the knowledge of both the derived categories and the properties of integral functors for singular varieties could be of great relevance.
However, very little attention has been paid so far to singular varieties in the Fourier-Mukai literature. One of the reasons may be the fact that the fundamental results on integral functors are not easily generalised to the singular situation, because they rely deeply on properties inherent to smoothness.
We would like to mention two of the most important. One is Orlov’s representation theorem according to if $`X`$ and $`Y`$ are smooth projective varieties, any (exact) fully faithful functor between their derived categories is an integral functor. Particularly, any (exact) equivalence between their derived categories is an integral functor (integral functors that are equivalences are also known as Fourier-Mukai functors). Another is Bondal and Orlov’s characterisation of those integral functors between the derived categories of two smooth varieties that are fully faithful .
Orlov’s representation theorem has been generalised by Kawamata to the smooth stack associated to a normal projective variety with only quotient singularities. Therefore $`D`$-equivalence also implies $`K`$-equivalence for those varieties when $`\kappa (X)`$ is maximal. In Van den Bergh proves using non-commutative rings that Bridgeland’s result about flopping contractions can be extended to quasi-projective varieties with only Gorenstein terminal singularities. The same result was proved by Chen ; the underlying idea is to embed such a threefold into a smooth fourfold and then use the essential smoothness. The author himself notices that his smoothing approach will not work for most general threefold flops because quotient singularities in dimension greater or equal to 3 are very rigid. In his paper, some general properties of the Fourier-Mukai transform on singular varieties can be found as well as the computation of a spanning class of the derived category of a normal projective variety with only isolated singularities. Finally, Kawamata has obtained analogous results for some $``$-Gorenstein threefolds using algebraic stacks.
This paper is divided in two parts. In the first part, we give an extension of Bondal and Orlov’s characterisation of fully faithful integral functors to proper varieties with (arbitrary) Gorenstein singularities. This is the precise statement.
###### Theorem (Theorem 1.22).
Let $`X`$ and $`Y`$ be projective Gorenstein schemes over an algebraically closed field of characteristic zero, and let $`𝒦^{}`$ be an object in $`D_c^b(X\times Y)`$ of finite projective dimension over $`X`$ and over $`Y`$. Assume also that $`X`$ is integral. Then the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is fully faithful if and only if the kernel $`𝒦^{}`$ is strongly simple over $`X`$.
One should notice that this Theorem may fail to be true in positive characteristic even in the smooth case. A counterexample is given in Remark 1.25.
In the Gorenstein case, strong simplicity (Definition 1.19) is defined in terms of locally complete intersection zero cycles instead of the structure sheaves of the closed points, as it happens in the smooth case. In the latter situation, our result improves the characterization of fully faithfulness of Bondal and Orlov.
As in the smooth case, when $`X`$ is a Gorenstein variety the skyscraper sheaves $`𝒪_x`$ form a spanning class for the derived category $`D_c^b(X)`$. Nevertheless, due to the fact that one may has an infinite number of non-zero $`\mathrm{Ext}_X^i(𝒪_x,𝒪_x)`$ when $`x`$ is a singular point, this spanning class does not allow to give an effective criterion characterising the fully faithfulness of integral functors. However, Bridgeland’s criterion that characterises when a fully faithful integral functor is an equivalence is also valid in the Gorenstein case. Moreover, since for a Gorenstein variety one has a more natural spanning class given by the structure sheaves of locally complete intersection cycles supported on closed points, one also proves the following alternative result.
###### Theorem (Theorem 1.28).
Let $`X`$, $`Y`$ and $`𝒦^{}`$ be as in the previous theorem with $`Y`$ connected. A fully faithful integral functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is an equivalence of categories if and only if for every closed point $`xX`$ there exists a locally complete intersection cycle $`Z_x`$ supported on $`x`$ such that $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})\omega _Y`$.
We also derive in the Gorenstein case some geometric consequences of the existence of Fourier-Mukai functors (Proposition 1.30) which are analogous to certain well-known properties of smooth schemes.
The second part of the paper is devoted to relative integral functors. As already mentioned, relative Fourier-Mukai transforms have been considered mainly in connection with elliptic fibrations. And besides some standard functorial properties, like compatibility with (some) base changes, more specific results or instances of Fourier-Mukai functors (equivalences of the derived categories) are known almost only for abelian schemes or elliptic fibrations.
We prove a new result that characterises when a relative integral functor is fully faithful or an equivalence, and generalises \[19, Prop. 6.2\]:
###### Theorem (Theorem 2.4).
Let $`p:XS`$ and $`q:YS`$ be locally projective Gorenstein morphisms (the base field is algebraically closed of characteristic zero). Let $`𝒦^{}D^b(X\times _SY)`$ be a kernel of finite projective dimension over both $`X`$ and $`Y`$. The relative integral functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is fully faithful (respectively an equivalence) if and only if $`\mathrm{\Phi }_{X_sY_s}^{𝐋j_s^{}𝒦^{}}:D_c^b(X_s)D_c^b(Y_s)`$ is fully faithful (respectively an equivalence) for every closed point $`sS`$, where $`j_s`$ is the immersion of $`X_s\times Y_s`$ into $`X\times _SY`$.
Though this result is probably true in greater generality, our proof needs the Gorenstein condition in an essential way. The above theorem, together with the characterisation of fully faithful integral functors and of Fourier-Mukai functors in the absolute Gorenstein case (Theorems 1.22 and 1.28) gives a criterion to ascertain when a relative integral functor between the derived categories of the total spaces of two Gorenstein fibrations is an equivalence. We expect that this theorem could be applied to very general situations. As a first application we give here a very simple and short proof of the invertibility result for elliptic fibrations:
###### Theorem (Proposition 2.7).
Let $`S`$ be an algebraic scheme over an algebraically closed field of characteristic zero, $`XS`$ an elliptic fibration with integral fibres and a section, $`\widehat{X}S`$ the dual fibration and $`𝒫`$ the relative Poincaré sheaf on $`X\times _S\widehat{X}`$. The relative integral functor
$$\mathrm{\Phi }_{X\widehat{X}}^𝒫:D_c^b(X)D_c^b(\widehat{X})$$
is an equivalence of categories.
This result has been proved elsewhere in different ways. When the total spaces of the fibrations involved are smooth the theorem can be proved, even if the fibres are singular, by considering the relative integral functor as an absolute one (defined by the direct image of the relative Poincaré to the direct product) and then applying the known criteria in the smooth case (see also ). When the total spaces are singular, there is a proof in that follows a completely different path and is much longer than ours.
In this paper, scheme means algebraic scheme (that is, a scheme of finite type) over an algebraically closed field $`k`$. By a Gorenstein morphism, we understand a flat morphism of schemes whose fibres are Gorenstein. For any scheme $`X`$ we denote by $`D(X)`$ the derived category of complexes of $`𝒪_X`$-modules with quasi-coherent cohomology sheaves. This is the essential image of the derived category of quasi-coherent sheaves in the derived category of all $`𝒪_X`$-modules. Analogously $`D^+(X)`$, $`D^{}(X)`$ and $`D^b(X)`$ will denote the derived categories of complexes which are respectively bounded below, bounded above and bounded on both sides, and have quasi-coherent cohomology sheaves. The subscript $`c`$ will refer to the corresponding subcategories of complexes with coherent cohomology sheaves.
### Acknowledgements
We would like to thank Claudio Bartocci, Ugo Bruzzo and Tom Bridgeland for discussions, indications and ideas. We also thank Alessio Corti and Constantin Teleman for organising a discussion session about Fourier-Mukai for singular varieties during the “Coherent Sheaves and Mirror Symmetry” workshop held in Cambridge (May 2005) and to the participants in the session, specially Igor Burban, Alexander Kuznetsov, Dimitri Orlov, Tony Pantev and Miles Reid, for useful comments and suggestions. Finally we owe our gratitude to the authors of the forthcoming book for sharing with us their notes and to the anonymous referee for comments and suggestions which helped us to improve the manuscript.
## 1. Fourier-Mukai transform on Gorenstein schemes
### 1.1. Preliminary results
We first recall some basic formulas which will be used in the rest of the paper.
If $`X`$ is a scheme, there is a functorial isomorphism (in the derived category)
(1.1)
$$𝐑om_{𝒪_X}^{}(^{},𝐑om_{𝒪_X}^{}(^{},^{}))\stackrel{}{}𝐑om_{𝒪_X}^{}(^{}\stackrel{𝐋}{}^{},^{})$$
where $`^{}`$, $`^{}`$ are in $`D^{}(X)`$, $`^{}`$ is in $`D^+(X)`$, and all have coherent cohomology (). One also has a functorial isomorphism
(1.2)
$$𝐑om_{𝒪_X}^{}(^{},^{})\stackrel{𝐋}{}^{}\stackrel{}{}𝐑om_{𝒪_X}^{}(^{},^{}\stackrel{𝐋}{}^{})$$
where $`^{}`$ is a bounded complex of $`𝒪_X`$-modules with coherent cohomology and either $`^{}`$ or $`^{}`$ is of finite homological dimension (i.e. locally isomorphic to a bounded complex of locally free sheaves of finite rank). The usual proof (see or ) requires that $`^{}`$ is of finite homological dimension; however, it still works when both members are defined. If we denote by $`_{}^{}{}_{}{}^{}=𝐑om_{𝒪_X}^{}(^{},𝒪_X)`$ the dual in the derived category, (1.2) implies that
(1.3)
$$_{}^{}{}_{}{}^{}\stackrel{𝐋}{}^{}\stackrel{}{}𝐑om_{𝒪_X}^{}(^{},^{}).$$
Nevertheless this formula may fail to be true when neither $`^{}`$ nor $`^{}`$ have finite homological dimension as the following example shows.
###### Example 1.1.
Let $`X`$ be a Gorenstein scheme of dimension $`n`$ over a field $`k`$. Let $`xX`$ be a singular point and let $``$ be any $`𝒪_X`$-module. Since $`𝒪_x^{}𝒪_x[n]`$, if one had
$$𝒪_x^{}\stackrel{𝐋}{}𝐑om_{𝒪_X}^{}(𝒪_x,),$$
then one would have $`\mathrm{Tor}_{ni}(𝒪_x,)\mathrm{Ext}^i(𝒪_x,)`$ for every $`i`$. It follows that $`\mathrm{Ext}^i(𝒪_x,)=0`$ for all $`i>n`$ and every $`𝒪_X`$-module $``$ and this is impossible because $`𝒪_x`$ is not of finite homological dimension. $`\mathrm{}`$
The formula (1.3) implies that if $`f:XY`$ is a morphism, there is an isomorphism
(1.4)
$$𝐋f^{}(^{})(𝐋f^{}^{})^{}$$
if either $`^{}`$ is of finite homological dimension or $`f`$ is of finite Tor-dimension (in this paper we shall only need to consider the case when $`f`$ is flat or is a regular closed immersion).
Some other formulas will be useful. When $`X`$ is a Gorenstein scheme, every object $`^{}`$ in $`D_c^b(X)`$ is reflexive, that is, one has an isomorphism in the derived category \[37, 1.17\]:
(1.5)
$$^{}(^{})^{}.$$
Then, one has
(1.6)
$$\mathrm{Hom}_{D(X)}(^{},^{})\mathrm{Hom}_{D(X)}(^{},^{})$$
for every bounded complex $`^{}`$ in $`D_c^b(X)`$ and any complex $`^{}`$.
Moreover, if $`X`$ is a zero dimensional Gorenstein scheme, the sheaf $`𝒪_X`$ is injective so that
(1.7)
$$_{}^{}{}_{}{}^{}_{}^{}{}_{}{}^{}\text{and}^i(_{}^{}{}_{}{}^{})(^i(^{}))^{}$$
for every object $`^{}`$ in $`D_c^b(X)`$, where $`_{}^{}{}_{}{}^{}=om_{𝒪_X}^{}(^{},𝒪_X)`$ is the ordinary dual.
Let $`f:XY`$ be a proper morphism of schemes. The relative Grothendieck duality states the existence of a functorial isomorphism in the derived category
(1.8)
$$𝐑om_{𝒪_Y}^{}(𝐑f_{}^{},𝒢^{})𝐑f_{}𝐑om_{𝒪_X}^{}(^{},f^!𝒢^{}).$$
for $`𝒢^{}`$ in $`D(Y)`$ and $`^{}`$ in $`D(X)`$ (see for instance ). By applying the derived functor of the global section functor, we obtain the *global duality formula*
(1.9)
$$\mathrm{Hom}_{D(Y)}(𝐑f_{}^{},𝒢^{})\mathrm{Hom}_{D(X)}(^{},f^!𝒢^{}).$$
In other words, the direct image $`𝐑f_{}:D(X)D(Y)`$ has a right adjoint $`f^!:D(Y)D(X)`$.
There is a natural map $`f^{}𝒢^{}\stackrel{𝐋}{}f^!𝒪_Yf^!𝒢^{}`$, which is an isomorphism when either $`𝒢^{}`$ has finite homological dimension or $`𝒢^{}`$ is reflexive and $`f^!𝒪_Y`$ has finite homological dimension.
When $`f`$ is a Gorenstein morphism of relative dimension $`n`$, the object $`f^!𝒪_Y`$ reduces to an invertible sheaf $`\omega _f`$, called *the relative dualizing sheaf*, located at the place $`n`$, $`f^!𝒪_Y\omega _f[n]`$.
Grothendieck duality is compatible with base-change. We state this result for simplicity only when $`f`$ is Gorenstein. In this case, since $`f`$ is flat, base-change compatibility means that if $`g:ZY`$ is a morphism and $`f_Z:Z\times _YXZ`$ is the induced morphism, then the relative dualizing sheaf for $`f_Z`$ is $`\omega _{f_Z}=g_X^{}\omega _f`$ where $`g_X:Z\times _YXX`$ is the projection.
As it is customary, when $`f`$ is the projection onto a point, we denote the dualizing sheaf by $`\omega _X`$.
### 1.2. Complexes of relative finite projective dimension
In this subsection we shall prove a weaker version of (1.2) in some cases.
###### Lemma 1.2.
Let $`^{}`$ be an object in $`D_c^b(X)`$. The following conditions are equivalent:
1. $`^{}`$ is of finite homological dimension.
2. $`^{}\stackrel{𝐋}{}𝒢^{}`$ is an object of $`D^b(X)`$ for every $`𝒢^{}`$ in $`D^b(X)`$.
3. $`𝐑om_{𝒪_X}^{}(^{},𝒢^{})`$ is in $`D^b(X)`$ for every $`𝒢^{}`$ in $`D^b(X)`$.
###### Proof.
Since $`X`$ is noetherian, the three conditions are local so that we can assume that $`X`$ is affine. It is clear that (1) implies (2) and (3). Now let us see that (3) implies (1). Let us consider a quasi-isomorphism $`^{}^{}`$ where $`^{}`$ is a bounded above complex of finite free modules. If $`𝒦^n`$ is the kernel of the differential $`^n^{n+1}`$, then for $`n`$ small enough the truncated complex $`𝒦^n^n\mathrm{}`$ is still quasi-isomorphic to $`^{}`$ because $`^{}`$ is an object of $`D_c^b(X)`$. Let $`x`$ be a point and $`𝒪_x`$ its residual field. Since $`𝐑om_{𝒪_X}^{}(^{},𝒪_x)`$ has bounded homology, one also has that $`\mathrm{Ext}_{𝒪_X}^1(𝒦^n,𝒪_x)=0`$ for $`n`$ small enough. For such $`n`$ the module $`𝒦^n`$ is free in a neighbourhood of $`x`$ and one concludes. To prove that (2) implies (1), one proceeds analogously replacing $`\mathrm{Ext}^1`$ by $`\mathrm{Tor}_1`$. ∎
This lemma suggests the following definition.
###### Definition 1.3.
Let $`f:XY`$ be a morphism of schemes. An object $`^{}`$ in $`D(X)`$ is said to be of *finite homological dimension over $`Y`$* (resp. of *finite projective dimension over $`Y`$*), if $`^{}\stackrel{𝐋}{}𝐋f^{}𝒢^{}`$ (resp. $`𝐑om_{𝒪_X}^{}(^{},f^!𝒢^{})`$), is in $`D^b(X)`$ for any $`𝒢^{}`$ in $`D^b(Y)`$. $`\mathrm{}`$
These notions are similar (though weaker) to the notions of finite Tor-amplitude and finite Ext-amplitude considered in .
In the absolute case (i.e. when $`f`$ is the identity), finite projective dimension is equivalent to finite homological dimension by the previous lemma. To characterise complexes of finite projective dimension over $`Y`$ when $`f`$ is projective, we shall need the following result (c.f. \[36, Lem. 2.13\]).
###### Lemma 1.4.
Let $`A`$ be a noetherian ring, $`f:XY=\mathrm{Spec}A`$ a projective morphism and $`𝒪_X(1)`$ a relatively very ample line bundle.
1. Let $`^{}`$ be an object of $`D^{}(X)`$. Then $`^{}=0`$ (resp. is an object of $`D^b(X)`$) if and only if $`𝐑f_{}(^{}(r))=0`$ (resp. is an object of $`D^b(Y)`$) for every integer $`r`$.
2. Let $`g:^{}𝒩^{}`$ be a morphism in $`D^{}(X)`$. Then $`g`$ is an isomorphism if and only if the induced morphism $`𝐑f_{}(^{}(r))𝐑f_{}(𝒩^{}(r))`$ is an isomorphism in $`D^{}(Y)`$ for every integer $`r`$.
(As it is usual, we set $`^{}(r)=^{}𝒪_X(r)`$.)
###### Proof.
Let $`i:X_A^N`$ be the closed immersion of $`A`$-schemes defined by $`𝒪_X(1)`$. Since $`^{}=0`$ if and only if $`i_{}^{}=0`$ and $`^{}`$ has bounded cohomology if and only if $`i_{}^{}`$ has bounded cohomology as well, we can assume that $`X=_A^N`$. Now one has an exact sequence (Beilinson’s resolution of the diagonal)
$$0_N\mathrm{}_1_0𝒪_\mathrm{\Delta }0$$
where $`_j=\pi _1^{}𝒪_{_A^N/A}(j)\pi _2^{}\mathrm{\Omega }_{_A^N/A}^j(j)`$, $`\pi _1`$ and $`\pi _2`$ being the projections of $`_A^N\times _A^N`$ onto its factors. Then $`𝒪_\mathrm{\Delta }`$ is an object of the smallest triangulated subcategory of $`D^b(_A^N\times _A^N)`$ that contains the sheaves $`_j`$ for $`0jN`$. Since $`F(^{})=𝐑\pi _2(\pi _1^{}(^{})\stackrel{𝐋}{}^{})`$ is an exact functor $`D^b(_A^N\times _A^N)D^{}(_A^N)`$, $`^{}F(𝒪_\mathrm{\Delta })`$ is an object of the smallest triangulated category generated by the objects $`F(_j)`$ for $`0jN`$. Thus to prove (1) we have only to see that $`F(_j)=0`$ (resp. have bounded homology) for all $`0jN`$. This follows because we have
$$F(_j)𝐑\pi _2(\pi _1^{}(^{}(j)))\stackrel{𝐋}{}\mathrm{\Omega }_{_A^N/A}^j(j)f^{}𝐑f_{}(^{}(j)))\mathrm{\Omega }^j_{_A^N/A}(j)$$
by the projection formula \[35, Prop. 5.3\] and flat base-change.
By applying the first statement to the cone of $`g`$, the second statement follows. ∎
One can also easily prove that $`^{}=0`$ if and only if $`𝐑f_{}(^{}(r))=0`$ for all $`r`$ by using the spectral sequence $`R^pf_{}(^q(^{}(r)))R^{p+q}f_{}(^{}(r))`$.
###### Lemma 1.5.
Let $`f:XY`$ be a proper morphism and $`^{}`$ an object of $`D_c^b(X)`$. If $`^{}`$ is either of finite projective dimension or of finite homological dimension over $`Y`$, then $`𝐑f_{}^{}`$ is of finite homological dimension.
###### Proof.
The duality isomorphism (1.8) together with Lemma 1.2 imply that $`𝐑f_{}^{}`$ is of finite homological dimension when $`^{}`$ is of finite projective dimension over $`Y`$. If $`^{}`$ is of finite homological dimension over $`Y`$, we use the same lemma and the projection formula. ∎
###### Proposition 1.6.
Let $`f:XY`$ be a projective morphism and $`^{}`$ an object of $`D_c^b(X)`$. The following conditions are equivalent:
1. $`^{}`$ is of finite projective dimension over $`Y`$.
2. $`𝐑f_{}(^{}(r))`$ is of finite homological dimension for every integer $`r`$.
3. $`^{}`$ is of finite homological dimension over $`Y`$.
Thus, if $`f`$ is locally projective, $`^{}`$ is of finite projective dimension over $`Y`$ if and only if it is of finite homological dimension over $`Y`$.
###### Proof.
If $`^{}`$ is of finite projective dimension (resp. of finite homological dimension) over $`Y`$, so is $`^{}(r)`$ for every $`r`$, and then $`𝐑f_{}(^{}(r))`$ is of finite homological dimension by Lemma 1.5. Assume that (2) is satisfied. Then (1) is a consequence of the duality isomorphism $`𝐑f_{}(𝐑om_{𝒪_X}^{}(^{},f^!𝒢^{})(r))𝐑om_{𝒪_Y}^{}(𝐑f_{}(^{}(r)),𝒢^{})`$ and Lemma 1.2, whilst (3) follows from the same lemma and the projection formula. ∎
###### Corollary 1.7.
Let $`f:XY`$ be a projective morphism and $`^{}`$ an object of $`D_c^b(X)`$. If $`^{}`$ is of finite projective dimension over $`Y`$, then $`𝐑om_{𝒪_X}^{}(^{},f^!𝒪_Y)`$ is also of finite projective dimension over $`Y`$. In particular, if $`f`$ is Gorenstein, $`_{}^{}{}_{}{}^{}`$ is of finite projective dimension over $`Y`$.
###### Proof.
Let us write $`𝒩^{}=𝐑om_{𝒪_X}^{}(^{},f^!𝒪_Y)`$. By Proposition 1.6, it suffices to see that $`𝐑f_{}(𝒩^{}(r))`$ is of finite homological dimension for every $`r`$. This follows again by Proposition 1.6, due to the isomorphism $`𝐑f_{}(𝒩^{}(r))[𝐑f_{}(^{}(r))]^{}`$. ∎
###### Proposition 1.8.
Let $`f:XY`$ be a locally projective Gorenstein morphism of schemes and $`^{}`$ an object of $`D_c(X)`$ of finite projective dimension over $`Y`$. One has
$$^{}\stackrel{𝐋}{}f^{}𝒢^{}\omega _f[n]𝐑om_{𝒪_X}^{}(^{},f^!𝒢^{})$$
for $`𝒢^{}`$ in $`D_c^b(Y)`$. Moreover, if $`Y`$ is Gorenstein, then
$$^{}\stackrel{𝐋}{}f^{}𝒢^{}𝐑om_{𝒪_X}^{}(^{},f^{}𝒢^{}).$$
###### Proof.
One has natural morphisms
(1.10) $`𝐑om_{𝒪_X}(^{},𝒪_X)\stackrel{𝐋}{}f^{}𝒢^{}\omega _f[n]`$ $`𝐑om_{𝒪_X}(^{},f^{}𝒢^{}\omega _f[n])`$
$`𝐑om_{𝒪_X}(^{},f^!𝒢^{}).`$
We have to prove that the composition is an isomorphism. This is a local question on $`Y`$, so that we can assume that $`Y=\mathrm{Spec}A`$.
By Lemma 1.4 we have to prove that the induced morphism
(1.11)
$$𝐑f_{}(𝐑om_{𝒪_X}(^{},𝒪_X)\stackrel{𝐋}{}f^{}𝒢^{}\omega _f[n]𝒪_X(r))𝐑f_{}(𝐑om_{𝒪_X}(^{},f^!𝒢^{})𝒪_X(r))$$
is an isomorphism in $`D^{}(Y)`$ for any integer $`r`$. The first member is isomorphic to
$$𝐑om_{𝒪_Y}(𝐑f_{}(^{}(r)),𝒪_Y)\stackrel{𝐋}{}𝒢^{}$$
by the projection formula and relative duality; the second one is isomorphic to
(1.12)
$$𝐑f_{}(𝐑om_{𝒪_X}(^{},f^!𝒢^{})𝒪_X(r))𝐑om_{𝒪_Y}(𝐑f_{}(^{}(r)),𝒢^{}).$$
Thus, we have to prove that the natural morphism
(1.13)
$$𝐑om_{𝒪_Y}(𝐑f_{}^{}(r),𝒪_Y)\stackrel{𝐋}{}𝒢^{}𝐑om_{𝒪_Y}(𝐑f_{}^{}(r),𝒢^{}),$$
is an isomorphism. Since $`𝐑f_{}^{}(r)`$ is of finite homological dimension by Proposition 1.6, one concludes by (1.2). ∎
### 1.3. Depth and local properties of Cohen-Macaulay and Gorenstein schemes
Here we state some preliminary results about depth on singular schemes and local properties of Cohen-Macaulay and Gorenstein schemes. We first recall a local property of Cohen-Macaulay schemes.
###### Lemma 1.9.
\[38, Prop. 6.2.4\] Let $`A`$ be a noetherian local ring. $`A`$ is Cohen-Macaulay if and only if there is an ideal $`I`$ of $`A`$ with $`dimA/I=0`$ and such that $`A/I`$ has finite homological dimension.
###### Proof.
Let $`n`$ be the dimension of the ring $`A`$. If $`A`$ is Cohen-Macaulay, $`0pt(A)=n`$. Then there is a regular sequence $`(a_1,\mathrm{},a_n)`$ in $`A`$ and taking $`I=(a_1,\mathrm{},a_n)`$ we conclude. Conversely, if $`I`$ is an ideal satisfying $`dimA/I=0`$ and $`\mathrm{hdim}(A/I)=s<\mathrm{}`$, then the Auslander and Buchsbaum’s formula $`0pt(A/I)+\mathrm{hdim}(A/I)=0pt(A)`$ \[30, Thm. 19.1\] proves that $`s=0pt(A)`$; then $`sn`$. Moreover, if $`^{}`$ is a free resolution of $`A/I`$ of length $`s`$, one has $`sn`$ by the intersection theorem \[38, 6.2.2.\]. Thus $`A`$ is Cohen-Macaulay. ∎
Let $``$ be a coherent sheaf on a scheme $`X`$ of dimension $`n`$. We write $`n_x`$ for the dimension of the local ring $`𝒪_{X,x}`$ of $`X`$ at a point $`xX`$ and $`_x`$ for the stalk of $``$ at $`x`$. $`_x`$ is a $`𝒪_{X,x}`$-module. The integer number $`\mathrm{codepth}(_x)=n_x0pt(_x)`$ is called the codepth of $``$ at $`x`$. For any integer $`m`$, the *$`m`$-th singularity set* of $``$ is defined to be
$$S_m()=\{xX|\mathrm{codepth}(_x)nm\}.$$
Then, if $`X`$ is equidimensional, a closed point $`x`$ is in $`S_m()`$ if and only if $`0pt(_x)m`$. If $`x`$ is a point of $`X`$ (not necessarily closed) the zero cycles $`Z_x`$ of $`\mathrm{Spec}𝒪_{X,x}`$ supported on the closed point $`x`$ of $`\mathrm{Spec}𝒪_{X,x}`$ will be called *zero cycles (of $`X`$) supported on $`x`$* by abuse of language.
Since $`0pt(_x)`$ is the first integer $`i`$ such that either
* $`\mathrm{Ext}^i(𝒪_x,)0`$ or
* $`H_x^i(\mathrm{Spec}𝒪_{X,x},_x)0`$ or
* $`\mathrm{Ext}^i(𝒪_Z,_x)0`$ for some zero cycle $`Z`$ supported on $`x`$ or
* $`\mathrm{Ext}^i(𝒪_Z,_x)0`$ for every zero cycle $`Z`$ supported on $`x`$
(see for instance ), we have alternative descriptions of $`S_m()`$:
(1.14) $`S_m()`$ $`=\{xX|H_x^i(\mathrm{Spec}𝒪_{X,x},_x)0\text{ for some }im+n_xn\}`$
$`=\{xX|\mathrm{Ext}^i(𝒪_Z,_x)0\text{ for some }im+n_xn\text{ and some}`$
$`\text{zero cycle }Z\text{ supported on }x\}`$
$`=\{xX|\mathrm{Ext}^i(𝒪_Z,_x)0\text{ for some }im+n_xn\text{ and any}`$
$`\text{zero cycle }Z\text{ supported on }x\}`$
###### Lemma 1.10.
If $`X`$ is smooth, then the $`m`$-th singularity set of $``$ can be described as
$$S_m()=_{pnm}\{xX|L_pj_x^{}0\},$$
where $`j_x`$ is the immersion of the point $`x`$.
###### Proof.
Let $`xX`$ be a point and $`^{}`$ the Koszul complex associated locally to a regular sequence of generators of the maximal ideal of $`𝒪_{X,x}`$. Since $`_{}^{}{}_{}{}^{}^{}[n_x]`$, one has an isomorphism $`\mathrm{Ext}^i(𝒪_x,_x)L_{n_xi}j_x^{}`$ which proves the result. ∎
In the singular case, this characterization of $`S_m()`$ is not true. However, there is a similar interpretation for Cohen-Macaulay schemes as we shall see now. By Lemma 1.9, if $`X`$ is Cohen-Macaulay, for every point $`x`$ there exist zero cycles supported on $`x`$ defined locally by a regular sequence; we refer to them as *locally complete intersection* or *l.c.i.* cycles. If $`ZX`$ is such a l.c.i. cycle, by the Koszul complex theory the structure sheaf $`𝒪_Z`$ has *finite homological dimension* as an $`𝒪_X`$-module.
We denote by $`j_Z`$ the immersion of $`Z`$ in $`X`$. Recall that for every object $`𝒦^{}`$ in $`D^b(X)`$, $`L_ij_Z^{}𝒦^{}`$ denotes the cohomology sheaf $`^i(j_Z^{}^{})`$ where $`^{}`$ is a bounded above complex of locally free sheaves quasi-isomorphic to $`𝒦^{}`$.
###### Lemma 1.11.
If $`X`$ is Cohen-Macaulay, then the $`m`$-th singularity set $`S_m()`$ can be described as
$`S_m()`$ $`=\{xX|\text{there is an integer }inm\text{ with }L_ij_{Z_x}^{}0`$
$`\text{ for any l.c.i zero cycle }Z_x\text{ supported on }x\}.`$
###### Proof.
Let $`Z_x`$ be a l.c.i. zero cycle supported on $`x`$ and $`^{}`$ the Koszul complex associated locally to a regular sequence of generators of the ideal of $`Z_x`$. As in the smooth case, we have that $`_{}^{}{}_{}{}^{}^{}[n_x]`$ and then an isomorphism $`\mathrm{Ext}^i(𝒪_{Z_x},)L_{n_xi}j_{Z_x}^{}`$. The result follows from (1.14). ∎
###### Lemma 1.12.
If $`j:XW`$ is a closed immersion and $``$ is a coherent sheaf on $`X`$, then $`S_m()=S_m(j_{})`$.
###### Proof.
Since $`H_x^i(\mathrm{Spec}𝒪_{X,x},_x)=H_x^i(\mathrm{Spec}𝒪_{W,x},(j_{})_x)`$, the result follows from (1.14). ∎
###### Proposition 1.13.
Let $`X`$ be an equidimensional scheme of dimension $`n`$ and $``$ a coherent sheaf on $`X`$.
1. $`S_m()`$ is a closed subscheme of $`X`$ and $`\mathrm{codim}S_m()nm`$.
2. If $`Z`$ is an irreducible component of the support of $``$ and $`c`$ is the codimension of $`Z`$ in $`X`$, then $`\mathrm{codim}S_{nc}()=c`$ and $`Z`$ is also an irreducible component of $`S_{nc}()`$.
###### Proof.
All questions are local and then, by Lemma 1.12, we can assume that $`X`$ is affine and smooth. By Lemma 1.10, $`S_m()=_{pnm}X_p()`$, where $`X_p()=\{xX|L_pj_x^{}0\}`$. To prove (1), we have only to see that $`X_p()`$ is closed of codimension greater or equal than $`p`$. This can be seen by induction on $`p`$. If $`p=0`$, then $`X_0()`$ is the support of $``$ and the statement is clear. For $`p=1`$, $`X_1()`$ is the locus of points where $``$ is not locally free, which is closed of codimension greater or equal than 1, since $``$ is always free at the generic point. If $`p>1`$, let us consider an exact sequence $`0𝒩0`$ where $``$ is free and finitely generated. Then $`L_pj_x^{}L_{p1}j_x^{}𝒩`$ so that $`X_p()=X_{p1}(𝒩)`$ which is closed by induction. Moreover, if $`xX_p()`$, then $`L_pj_x^{}0`$, so that $`pdim𝒪_{X,x}`$ because $`𝒪_{X,x}`$ is a regular ring. It follows that $`\mathrm{codim}X_p()=\mathrm{max}_{xX_p()}\{dim𝒪_{X,x}\}p`$.
We finally prove (2). By \[30, Thm. 6.5\], the prime ideal of $`Z`$ is also a minimal associated prime to $``$. Thus, if $`x`$ is the generic point of $`Z`$, the maximal ideal of the local ring $`𝒪_{X,x}`$ is a prime associated to $`_x`$, and then $`\mathrm{Hom}(𝒪_x,_x)0`$. This proves that $`xS_{nc}()`$ and then $`ZS_{nc}()`$. The result follows.
###### Corollary 1.14.
Let $`X`$ be a Cohen-Macaulay scheme and let $``$ be a coherent $`𝒪_X`$-module. Let $`h:YX`$ be an irreducible component of the support of $``$ and $`c`$ the codimension of $`Y`$ in $`X`$. There is a non-empty open subset $`U`$ of $`Y`$ such that for any l.c.i. zero cycle $`Z_x`$ supported on $`xU`$ one has
$`L_cj_{Z_x}^{}`$ $`0`$
$`L_{c+i}j_{Z_x}^{}`$ $`=0,\text{ for every }i>0\text{.}`$
###### Proof.
By Lemma 1.11 the locus of the points that verify the conditions is $`U=Y(S_{nc}()S_{nc1}())`$, which is open in $`Y`$ by Proposition 1.13. Proving that $`U`$ is not empty is a local question, and we can then assume that $`Y`$ is the support of $``$. Now $`Y=S_{nc}()`$ by (2) of Proposition 1.13 and $`U=S_{nc}()S_{nc1}()`$ is non-empty because the codimension of $`S_{nc1}()`$ in $`X`$ is greater or equal than $`c+1`$ again by Proposition 1.13. ∎
The following proposition characterises objects of the derived category supported on a closed subscheme.
###### Proposition 1.15.
*\[9, Prop. 1.5\]*. Let $`j:YX`$ be a closed immersion of codimension $`d`$ of irreducible Cohen-Macaulay schemes and $`𝒦^{}`$ an object of $`D_c^b(X)`$. Assume that
1. If $`xXY`$ is a closed point, then $`𝐋j_{Z_x}^{}𝒦^{}=0`$ for some l.c.i. zero cycle $`Z_x`$ supported on $`x`$.
2. If $`xY`$ is a closed point, then $`L_ij_{Z_x}^{}𝒦^{}=0`$ for some l.c.i. zero cycle $`Z_x`$ supported on $`x`$ when either $`i<0`$ or $`i>d`$.
Then there is a sheaf $`𝒦`$ on $`X`$ whose topological support is contained in $`Y`$ and such that $`𝒦^{}𝒦`$ in $`D_c^b(X)`$. Moreover, this topological support coincides with $`Y`$ unless $`𝒦^{}=0`$.
###### Proof.
Let us write $`^q=^q(𝒦^{})`$. For every zero cycle $`Z_x`$ in $`X`$ there is a spectral sequence
$$E_2^{p,q}=L_pj_{Z_x}^{}^qE_{\mathrm{}}^{p+q}=L_{pq}j_{Z_x}^{}𝒦^{}$$
Let $`q_0`$ be the maximum of the $`q`$’s with $`^q0`$. If $`x\mathrm{supp}(^{q_0})`$, one has $`j_{Z_x}^{}^{q_0}0`$ for every l.c.i. zero cycle $`Z_x`$ supported on $`x`$. A nonzero element in $`j_{Z_x}^{}^{q_0}`$ survives up to infinity in the spectral sequence. Since there is a l.c.i. zero cycle $`Z_x`$ such that $`E_{\mathrm{}}^q=L_qj_{Z_x}^{}𝒦^{}=0`$ for every $`q>0`$ by hypothesis, one has $`q_00`$. A similar argument shows that the topological support of all the sheaves $`^q`$ is contained in $`Y`$: assume that this is not true and let us consider the maximum $`q_1`$ of the $`q`$’s such that $`j_x^{}^q0`$ for a certain point $`xXY`$; then $`j_{Z_x}^{}^{q_1}0`$ and a nonzero element in $`j_{Z_x}^{}^{q_1}`$ survives up to infinity in the spectral sequence, which is impossible since $`𝐋j_{Z_x}^{}𝒦^{}=0`$.
Let $`q_2q_0`$ be the minimum of the $`q`$’s with $`^q0`$. We know that $`^{q_2}`$ is topologically supported on a closed subset of $`Y`$. Take a component $`Y^{}Y`$ of the support. If $`cd`$ is the codimension of $`Y^{}`$, then there is a non-empty open subset $`U`$ of $`Y^{}`$ such that $`L_cj_{Z_x}^{}^{q_2}0`$ for any closed point $`xU`$ and any l.c.i. zero cycle $`Z_x`$ supported on $`x`$, by Corollary 1.14. Elements in $`L_cj_{Z_x}^{}^{q_2}`$ would be killed in the spectral sequence by $`L_pj_{Z_x}^{}^{q_2+1}`$ with $`pc+2`$. By Lemma 1.11 the set
$$\{xX|L_ij_{Z_x}^{}^{q_2+1}0\text{ for some }ic+2\text{ and any l.c.i. cycle }Z_x\}$$
is equal to $`S_{n(c+2)}(^{q_2+1})`$ and then has codimension greater or equal than $`c+2`$ by Proposition 1.13. Thus there is a point $`xY^{}`$ such that any nonzero element in $`L_cj_{Z_x}^{}^{q_2}`$ survives up to the infinity in the spectral sequence. Therefore, $`L_{cq_2}j_{Z_x}^{}𝒦^{}0`$ for any l.c.i. zero cycle $`Z_x`$ supported on $`x`$. Thus $`cq_2d`$ which leads to $`q_2cd0`$ and then $`q_2=q_0=0`$. So $`𝒦^{}=^0`$ in $`D^b(X)`$ and the topological support of $`𝒦=^0`$ is contained in $`Y`$. Actually, if $`𝒦^{}0`$, then this support is the whole of $`Y`$: if this was not true, since $`Y`$ is irreducible, the support would have a component $`Y^{}Y`$ of codimension $`c>d`$ and one could find, reasoning as above, a non-empty subset $`U`$ of $`Y^{}`$ such that $`L_cj_{Z_x}^{}𝒦^{}0`$ for all $`xU`$ and all l.c.i. zero cycle $`Z_x`$ supported on $`x`$. This would imply that $`cd`$, which is impossible. ∎
Taking into account that $`𝒪_{Z_x}^{}=𝒪_{Z_x}[n]`$ where $`n=dimX`$, Proposition 1.15 may be reformulated as follows:
###### Proposition 1.16.
Let $`j:YX`$ be a closed immersion of codimension $`d`$ of irreducible Cohen-Macaulay schemes of dimensions $`m`$ and $`n`$ respectively, and let $`𝒦^{}`$ be an object of $`D_c^b(X)`$. Assume that for any closed point $`xX`$ there is a l.c.i. zero cycle $`Z_x`$ supported on $`x`$ such that
$$\mathrm{Hom}_{D(X)}^i(𝒪_{Z_x},𝒦^{})=0,$$
unless $`xY`$ and $`min`$. Then there is a sheaf $`𝒦`$ on $`X`$ whose topological support is contained in $`Y`$ and such that $`𝒦^{}𝒦`$ in $`D_c^b(X)`$. Moreover, the topological support is $`Y`$ unless $`𝒦^{}=0`$. ∎
### 1.4. Integral functors
Let $`X`$ and $`Y`$ be proper schemes. We denote the projections of the direct product $`X\times Y`$ to $`X`$ and $`Y`$ by $`\pi _X`$ and $`\pi _Y`$.
Let $`𝒦^{}`$ be an object in $`D^b(X\times Y)`$. The integral functor defined by $`𝒦^{}`$ is the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D^{}(X)D^{}(Y)`$ given by
$$\mathrm{\Phi }_{XY}^𝒦^{}(^{})=𝐑\pi _Y(\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{}).$$
If the kernel $`𝒦^{}D_c^b(X\times Y)`$ is of finite homological dimension over $`X`$, then the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$ is defined over the whole $`D(X)`$ and maps $`D_c^b(X)`$ to $`D_c^b(Y)`$.
If $`Z`$ is a third proper scheme and $`^{}`$ is an object of $`D_c^b(Y\times Z)`$, arguing exactly as in the smooth case, we prove that there is an isomorphism of functors
$$\mathrm{\Phi }_{YZ}^{^{}}\mathrm{\Phi }_{XY}^𝒦^{}\mathrm{\Phi }_{XZ}^^{}𝒦^{}$$
where
(1.15)
$$^{}𝒦^{}=𝐑\pi _{X,Z}(\pi _{X,Y}^{}𝒦^{}\stackrel{𝐋}{}\pi _{Y,Z}^{}^{}).$$
If either $`𝒦^{}`$ or $`^{}`$ is of finite homological dimension over $`Y`$, then $`^{}𝒦^{}`$ in bounded.
### 1.5. Adjoints
We can describe nicely the adjoints to an integral functor when we work with Gorenstein schemes. In this subsection $`X`$ and $`Y`$ are projective Gorenstein schemes.
###### Proposition 1.17.
Let $`𝒦^{}`$ be an object in $`D_c^b(X\times Y)`$ of finite projective dimension over $`X`$ and $`Y`$.
1. The functor $`\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}:D_c^b(Y)D_c^b(X)`$ is a left adjoint to the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$.
2. The functor $`\mathrm{\Phi }_{YX}^{𝒦^{}\pi _X^{}\omega _X[m]}:D_c^b(Y)D_c^b(X)`$ is a right adjoint to the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$.
(here $`m=dimX`$ and $`n=dimY`$)
###### Proof.
We shall freely use (1.4) for the projections $`\pi _X`$ and $`\pi _Y`$.
(1) We first notice that one has
(1.16)
$$(\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{})^{}𝐑om_{𝒪_{X\times Y}}(𝒦^{},\pi _X^{}_{}^{}{}_{}{}^{})𝒦_{}^{}{}_{}{}^{}\stackrel{𝐋}{}\pi _X^{}_{}^{}{}_{}{}^{},$$
for $`^{}`$ in $`D_c^b(X)`$ by (1.1) and Proposition 1.8. The latter applies because $`\pi _X`$ is a projective morphism and $`𝒦^{}`$ is of finite projective dimension over $`X`$. Now, if $`𝒢^{}`$ is an object of $`D_c^b(Y)`$ there is a chain of isomorphisms
$`\mathrm{Hom}_{D(Y)}`$ $`(𝒢^{},\mathrm{\Phi }_{XY}^𝒦^{}(^{}))\mathrm{Hom}_{D(X\times Y)}(\pi _Y^{}𝒢^{},\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{})`$
$`\mathrm{Hom}_{D(X\times Y)}((\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{})^{},(\pi _Y^{}𝒢^{})^{})`$
$`\mathrm{Hom}_{D(X\times Y)}(\pi _X^{}_{}^{}{}_{}{}^{}\stackrel{𝐋}{}𝒦_{}^{}{}_{}{}^{},\pi _Y^{}𝒢_{}^{}{}_{}{}^{})`$
$`\mathrm{Hom}_{D(X\times Y)}(\pi _X^{}_{}^{}{}_{}{}^{},𝐑om_{𝒪_{X\times Y}}(𝒦_{}^{}{}_{}{}^{},\pi _Y^{}𝒢_{}^{}{}_{}{}^{}))`$
$`\mathrm{Hom}_{D(X\times Y)}(\pi _X^{}_{}^{}{}_{}{}^{},(𝒦_{}^{}{}_{}{}^{}\stackrel{𝐋}{}\pi _Y^{}𝒢^{})^{})`$
$`\mathrm{Hom}_{D(X)}(_{}^{}{}_{}{}^{},𝐑\pi _X((𝒦_{}^{}{}_{}{}^{}\stackrel{𝐋}{}\pi _Y^{}𝒢^{})^{}))`$
$`\mathrm{Hom}_{D(X)}(_{}^{}{}_{}{}^{},𝐑om_{𝒪_X}(𝐑\pi _X(𝒦_{}^{}{}_{}{}^{}\stackrel{𝐋}{}\pi _Y^{}𝒢^{}\pi _Y^{}\omega _Y[n]),𝒪_X))`$
$`\mathrm{Hom}_{D(X)}(_{}^{}{}_{}{}^{},[\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}(𝒢^{})]^{})`$
$`\mathrm{Hom}_{D(X)}(\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}(𝒢^{}),^{}).`$
where the second follows from (1.6) which applies because $`\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{}`$ is bounded, the third is (1.16), the forth and the fifth are (1.1), the seventh is relative duality and the ninth is again (1.6).
(2) The adjunction between the direct and inverse images and relative duality proves that the functor
$$H(𝒢^{})=𝐑\pi _{X,}(𝐑om_{𝒪_{X\times Y}}^{}(𝒦^{},\pi _Y^!𝒢^{}))$$
satisfies
(1.17)
$$\mathrm{Hom}_{D(Y)}(\mathrm{\Phi }_{XY}^𝒦^{}(^{}),𝒢^{})\mathrm{Hom}_{D(X)}(^{},H(𝒢^{})).$$
Then we conclude by Proposition 1.8 since $`\pi _Y`$ is a projective morphism. ∎
We shall need some basic results about adjoints and fully faithfulness which we state without proof.
###### Proposition 1.18.
Let $`\mathrm{\Phi }:𝒜`$ a functor and $`G:𝒜`$ a left adjoint (resp. $`H:𝒜`$ a right adjoint). Then the following conditions are equivalent:
1. $`\mathrm{\Phi }`$ is fully faithful.
2. $`G\mathrm{\Phi }`$ is fully faithful (resp. $`H\mathrm{\Phi }`$ is fully faithful).
3. The counit morphism $`G\mathrm{\Phi }\mathrm{Id}`$ is an isomorphism (resp. the unit morphism $`\mathrm{Id}H\mathrm{\Phi }`$ is an isomorphism).
Moreover, $`\mathrm{\Phi }`$ is an equivalence if and only if $`\mathrm{\Phi }`$ and $`G`$ (resp. $`\mathrm{\Phi }`$ and $`H`$) are fully faithful. ∎
### 1.6. Strongly simple objects
Let $`X`$ and $`Y`$ be proper Gorenstein schemes. In this situation, the notion of strong simplicity is the following.
###### Definition 1.19.
An object $`𝒦^{}`$ in $`D_c^b(X\times Y)`$ is *strongly simple* over $`X`$ if it satisfies the following conditions:
1. For every closed point $`xX`$ there is a l.c.i. zero cycle $`Z_x`$ supported on $`x`$ such that
$$\mathrm{Hom}_{D(Y)}^i(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_{x_1}}),\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_{x_2}}))=0$$
unless $`x_1=x_2`$ and $`0idimX`$.
2. $`\mathrm{Hom}_{D(Y)}^0(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x),\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x))=k`$ for every closed point $`xX`$.
$`\mathrm{}`$
The last condition can be written as $`\mathrm{Hom}_{D(Y)}^0(𝐋j_x^{}𝒦^{},𝐋j_x^{}𝒦^{})=k`$, because the restriction $`𝐋j_x^{}𝒦^{}`$ of $`𝒦^{}`$ to the fibre $`j_x:Y\{x\}\times YX\times Y`$ can also be computed as $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)`$.
In order to fix some notation, for any zero-cycle $`Z_x`$ of $`X`$ and any scheme $`S`$, we shall denote by $`j_{Z_x}`$ the immersion $`Z_x\times SX\times S`$.
###### Proposition 1.20.
Assume that $`Y`$ is projective, and let $`𝒦^{}`$ be a kernel in $`D^b(X\times Y)`$ of finite projective dimension over $`X`$. If $`𝒦^{}`$ is strongly simple over $`X`$, its dual $`𝒦^{}`$ is strongly simple over $`X`$ as well.
###### Proof.
If $`Z_x`$ is a l.c.i. zero cycle supported on $`x`$, one has that $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})=p_2𝐋j_{Z_x}^{}𝒦^{}`$, with $`p_2:Z_x\times YY`$ the second projection. Since $`\omega _{p_2}𝒪_{Z_x\times Y}`$ because $`Z_x`$ is zero dimensional and Gorenstein, one obtains
$$\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})^{}p_2(𝐋j_{Z_x}^{}𝒦^{})^{}.$$
Moreover, $`(𝐋j_{Z_x}^{}𝒦^{})^{}𝐋j_{Z_x}^{}(𝒦^{})`$ by (1.4) since $`j_{Z_x}`$ is a regular closed immersion. Then,
$$\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})^{}\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x}).$$
It follows that $`𝒦^{}`$ satisfies condition (1) of Definition 1.19. To see that it also fulfils condition (1), we have to prove that $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)^{}\mathrm{\Phi }_{XY}^𝒦_{}^{}{}_{}{}^{}(𝒪_x)`$, and this is equivalent to the base change formula $`𝐋j_x^{}(𝒦^{})(𝐋j_x^{}𝒦^{})^{}`$. Since we cannot longer use (1.4) because $`j_x`$ may fail to be of finite Tor-dimension, we proceed in a different way. To see that the natural morphism $`𝐋j_x^{}(𝒦^{})(𝐋j_x^{}𝒦^{})^{}`$ is an isomorphism, it suffices to check that $`j_x𝐋j_x^{}(𝒦^{})j_x[(𝐋j_x^{}𝒦^{})^{}]`$ since $`j_x`$ is a closed embedding. On the one hand, we have
$$j_x𝐋j_x^{}(𝒦^{})𝒦^{}\stackrel{𝐋}{}j_x𝒪_Y𝒦^{}\stackrel{𝐋}{}\pi _X^{}𝒪_x,$$
whilst on the other hand,
$`j_x[(𝐋j_x^{}𝒦^{})^{}]`$ $`=j_x𝐑om_{𝒪_Y}^{}(𝐋j_x^{}𝒦^{},𝒪_Y)𝐑om_{𝒪_{X\times Y}}^{}(𝒦^{},j_x𝒪_Y)`$
$`𝐑om_{𝒪_{X\times Y}}^{}(𝒦^{},\pi _X^{}𝒪_x).`$
We conclude by Proposition 1.8.
###### Remark 1.21.
When $`X`$ and $`Y`$ are smooth, strong simplicity is usually defined by the following conditions (see ):
1. $`\mathrm{Hom}_{D(Y)}^i(𝐋j_{x_1}^{}𝒦^{},𝐋j_{x_2}^{}𝒦^{})=0`$ unless $`x_1=x_2`$ and $`0idimX`$;
2. $`\mathrm{Hom}_{D(Y)}^0(𝐋j_x^{}𝒦^{},𝐋j_x^{}𝒦^{})=k`$ for every closed point $`x`$.
Since our definition is weaker, Theorem 1.22 improves Bondal and Orlov’s result \[9, Thm. 1.1\].
We now give the criterion for an integral functor between derived categories of Gorenstein proper schemes to be fully faithful.
###### Theorem 1.22.
Let $`X`$ and $`Y`$ be projective Gorenstein schemes over an algebraically closed field of characteristic zero, and let $`𝒦^{}`$ be an object in $`D_c^b(X\times Y)`$ of finite projective dimension over $`X`$ and over $`Y`$. Assume also that $`X`$ is integral. Then the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is fully faithful if and only if the kernel $`𝒦^{}`$ is strongly simple over $`X`$.
###### Proof.
If the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$ is fully faithful, then $`𝒦^{}`$ is strongly simple over $`X`$.
Let us prove the converse. Before starting, we fix some notation: we denote by $`\pi _i`$ the projections of $`X\times X`$ onto its factors and by $`U`$ the smooth locus of $`X`$, which is not empty because $`X`$ is integral. We also denote $`m=dimX`$, $`n=dimY`$ and $`\mathrm{\Phi }=\mathrm{\Phi }_{XY}^𝒦^{}`$.
By Proposition 1.17, $`\mathrm{\Phi }`$ has a left adjoint $`G=\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}`$ and a right adjoint $`H=\mathrm{\Phi }_{YX}^{𝒦^{}\pi _X^{}\omega _X[m]}`$. By Proposition 1.18 it suffices to show that $`G\mathrm{\Phi }`$ is fully faithful. We know that $`H\mathrm{\Phi }\mathrm{\Phi }_{XX}^{^{}}`$, and $`G\mathrm{\Phi }\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}`$, with $`^{}`$ and $`\stackrel{~}{^{}}`$ given by (1.4). Notice that since $`𝒦^{}`$ is of finite projective dimension over $`X`$ and $`Y`$, $`^{}`$ and $`\stackrel{~}{^{}}`$ are bounded.
The strategy of the proof is as follows: we first show that both $`^{}`$ and $`\stackrel{~}{^{}}`$ are single sheaves supported topologically on the image $`\mathrm{\Delta }`$ of the diagonal morphism $`\delta :XX\times X`$; then we prove that $`\stackrel{~}{^{}}`$ is actually schematically supported on the diagonal, that is, $`\stackrel{~}{^{}}=\delta _{}𝒩`$ for a coherent sheaf $`𝒩`$ on $`X`$ and finally that $`𝒩`$ is a line bundle; this will imply that $`\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}`$ is the twist by $`𝒩`$ which is an equivalence of categories, in particular fully faithful.
a) *$`^{}`$ and $`\stackrel{~}{^{}}`$ are single sheaves topologically supported on the diagonal.*
Let us fix a closed point $`(x_1,x_2)X\times X`$ and consider the l.c.i. zero cycles $`Z_{x_1}`$ and $`Z_{x_2}`$ of the first condition of the definition of strongly simple object. One has
$$\mathrm{Hom}_{D(X)}^i(𝒪_{Z_{x_1}},\mathrm{\Phi }_{XX}^{^{}}(𝒪_{Z_{x_2}}))\mathrm{Hom}_{D(Y)}^i(\mathrm{\Phi }(𝒪_{Z_{x_1}}),\mathrm{\Phi }(𝒪_{Z_{x_2}})),$$
which is zero unless $`x_1=x_2`$ and $`0im`$ because $`𝒦^{}`$ is strongly simple. Applying Proposition 1.16 to the immersion $`\{x_2\}X`$ we have that $`\mathrm{\Phi }_{XX}^{^{}}(𝒪_{Z_{x_2}})`$ reduces to a coherent sheaf topologically supported at $`x_2`$. Since $`\mathrm{\Phi }_{XX}^{^{}}(𝒪_{Z_{x_2}})p_2𝐋j_{Z_{x_2}}^{}^{}`$, where $`p_2:Z_{x_2}\times XX`$ is the second projection, the complex $`𝐋j_{Z_{x_2}}^{}^{}`$ is isomorphic to a single coherent sheaf $``$ topologically supported at $`(x_2,x_2)`$. If we denote by $`i_{Z_{x_1}}:Z_{x_2}\times Z_{x_1}Z_{x_2}\times X`$ and $`j_{Z_{x_2}\times Z_{x_1}}:Z_{x_2}\times Z_{x_1}X\times X`$ the natural immersions, we have
$$𝐋j_{Z_{x_2}\times Z_{x_1}}^{}^{}𝐋i_{Z_{x_1}}^{}𝐋j_{Z_{x_2}}^{}^{}𝐋i_{Z_{x_1}}^{}.$$
Thus, $`L_pj_{Z_{x_2}\times Z_{x_1}}^{}^{}=0`$ unless $`x_1=x_2`$ and $`0pm`$. Applying now Proposition 1.15 to $`\delta `$, we obtain that $`^{}`$ reduces to a coherent sheaf $``$ supported topologically on the diagonal as claimed.
For $`\stackrel{~}{^{}}`$, we proceed as follows. We have
$`^i(𝐋j_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{Z_{x_1}})^{})`$ $`\mathrm{Hom}_{D(Z_{x_2})}^i(𝐋j_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{Z_{x_1}}),𝒪_{Z_{x_2}})`$
$`\mathrm{Hom}_{D(X)}^i(\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{Z_{x_1}}),𝒪_{Z_{x_2}})\mathrm{Hom}_{D(Y)}^i(\mathrm{\Phi }(𝒪_{Z_{x_1}}),\mathrm{\Phi }(𝒪_{Z_{x_2}})),`$
which is zero unless $`x_1=x_2`$ and $`0im`$ because $`𝒦^{}`$ is strongly simple. Since $`Z_{x_2}`$ is a zero dimensional Gorenstein scheme, (1.7) implies that $`L_ij_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{Z_{x_1}})=0`$ again unless $`x_1=x_2`$ and $`0im`$. By Proposition 1.15 for the immersion $`\{x_1\}X`$, one has that $`\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{Z_{x_1}})`$ is a sheaf supported topologically at $`x_1`$. Now, a similar argument to the one used for $`^{}`$ proves that $`\stackrel{~}{^{}}`$ reduces to a coherent sheaf $`\stackrel{~}{}`$ supported topologically on the diagonal.
b) *$`\stackrel{~}{}`$ is schematically supported on the diagonal, that is, $`\stackrel{~}{}=\delta _{}𝒩`$ for a coherent sheaf $`𝒩`$ on $`X`$; moreover $`𝒩`$ is a line bundle.*
It might happen that the schematic support is an infinitesimal neighborhood of the diagonal; we shall see that this is not the case. Let us denote by $`\overline{\delta }:WX\times X`$ the schematic support of $`\stackrel{~}{}`$ so that $`\stackrel{~}{}=\overline{\delta }_{}𝒩`$ for a coherent sheaf $`𝒩`$ on $`W`$. The diagonal embedding $`\delta `$ factors through a closed immersion $`\tau :XW`$ which topologically is a homeomorphism.
b1) *$`\pi _2(\stackrel{~}{})`$ is locally free.*
To see this, we shall prove that $`\mathrm{Hom}_{D(X)}^1(\pi _2(\stackrel{~}{}),𝒪_x)=0`$ for every closed point $`xX`$. Since $`\stackrel{~}{}`$ is topologically supported on the diagonal, we have that $`\pi _2(\stackrel{~}{})𝐑\pi _2(\stackrel{~}{})\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_X)`$. We have
$$\mathrm{Hom}_{D(X)}^1(\pi _2(\stackrel{~}{}),𝒪_x)\mathrm{Hom}_{D(X)}^1(\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_X),𝒪_x)\mathrm{Hom}_{D(X)}^1(𝒪_X,\mathrm{\Phi }_{XX}^{^{}}(𝒪_x))$$
for every closed point $`xX`$, because $`H\mathrm{\Phi }`$ is a right adjoint to $`G\mathrm{\Phi }`$. Since $`\mathrm{\Phi }_{XX}^{^{}}(𝒪_x)𝐋j_x^{}`$ has only negative cohomology sheaves and all of them are supported at $`x`$, one has that $`\mathrm{Hom}_{D(X)}^1(\pi _2(\stackrel{~}{}),𝒪_x)=0`$ and $`\pi _2(\stackrel{~}{})`$ is locally free.
b2) *$`\pi _1(\stackrel{~}{})`$ is a line bundle on the smooth locus $`U`$ of $`X`$.*
We know that $`\mathrm{\Phi }_{XX}^{}(𝒪_{Z_{x_2}})`$ reduces to a single sheaf supported at $`x_2`$. Then, for every point $`x_1U`$ one has
$`^i(𝐋j_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{x_1})^{})`$ $`\mathrm{Hom}_{D(Z_{x_2})}^i(𝐋j_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{x_1}),𝒪_{Z_{x_2}})`$
$`\mathrm{Hom}_{D(X)}^i(\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{x_1}),𝒪_{Z_{x_2}})\mathrm{Hom}_{D(X)}^i(𝒪_{x_1},\mathrm{\Phi }_{XX}^{}(𝒪_{Z_{x_2}}))`$
which is zero unless $`x_2=x_1`$ and $`0im`$ because $`x_1`$ is a smooth point. Since $`Z_{x_2}`$ is a zero dimensional Gorenstein scheme, (1.7) implies that whenever $`x_1`$ is a smooth point, then $`L_ij_{Z_{x_2}}^{}\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{x_1})=0`$ unless $`x_2=x_1`$ and $`0im`$. By Proposition 1.15, $`\mathrm{\Phi }_{XX}^\stackrel{~}{^{}}(𝒪_{x_1})`$ reduces to a single sheaf supported at $`x_1`$. In particular $`𝐋j_x^{}(\stackrel{~}{})j_x^{}(\stackrel{~}{})`$ for every smooth point $`x`$, and thus the restriction of $`\stackrel{~}{}`$ to $`U\times X`$ is flat over $`U`$. Moreover, for every point $`xU`$, we have that
$$\mathrm{Hom}_X(j_x^{}\stackrel{~}{},𝒪_x)\mathrm{Hom}_{D(X)}^0(\mathrm{\Phi }(𝒪_x),\mathrm{\Phi }(𝒪_x))k.$$
By \[12, Lemmas 5.2 and 5.3\] there is a point $`x_0`$ in $`U`$ such that the Kodaira-Spencer map for the family $`\stackrel{~}{}_{|U\times X}`$ is injective at $`x_0`$. We now proceed as in the proof of \[12, Thm. 5.1\]: the morphism $`\mathrm{Hom}_{D(X)}^1(𝒪_{x_0},𝒪_{x_0})\mathrm{Hom}_{D(X)}^1(G\mathrm{\Phi }(𝒪_{x_0}),G\mathrm{\Phi }(𝒪_{x_0}))`$ is injective so that the morphism $`\mathrm{Hom}_{D(X)}^1(𝒪_{x_0},𝒪_{x_0})\mathrm{Hom}_{D(Y)}^1(\mathrm{\Phi }(𝒪_{x_0}),\mathrm{\Phi }(𝒪_{x_0}))`$ is injective as well and then the counit morphism $`j_{x_0}^{}\stackrel{~}{}G\mathrm{\Phi }(𝒪_{x_0})𝒪_{x_0}`$ is an isomorphism. Thus, the rank of $`\pi _1(\stackrel{~}{})`$ at the point $`x_0`$ is one, and then it is one everywhere in $`U`$.
b3) *$`\tau _{|U}:UW_U=W(U\times X)`$ is an isomorphism and $`𝒩^{}=𝒩_{|U}`$ is a line bundle.*
We proceed locally. We then write $`U=\mathrm{Spec}A`$, $`W_U=\mathrm{Spec}B`$ so that $`\tau `$ corresponds to a surjective ring morphism $`BA0`$ and the projection $`q_1=\pi _{1}^{}{}_{|U}{}^{}:W_UU`$ to an immersion $`AB`$. Now $`𝒩^{}`$ is a $`B`$-module which is isomorphic to $`A`$ as an $`A`$-module, $`𝒩^{}eA`$, because $`q_1(𝒩^{})=\pi _1(\stackrel{~}{})_{|U}`$ is a line bundle. It follows that $`𝒩^{}`$ is also generated by $`e`$ as a $`B`$-module. The kernel of $`B𝒩^{}eB0`$ is the annihilator of $`𝒩^{}`$ and then it is zero by the very definition of $`W`$. It follows that $`BA`$ as an $`A`$-module and then the morphism $`BA0`$ is an isomorphism. Hence, $`W_UU`$, $`q_1`$ is the identity map, and $`𝒩^{}q_1(𝒩^{})`$ is a line bundle.
b4) *$`\tau :XW`$ is an isomorphism and $`𝒩`$ is a line bundle.*
Since $`UW_U`$, $`\pi _2\stackrel{~}{}_{|U}𝒩_{|U}\pi _1\stackrel{~}{}_{|U}`$, which is a line bundle on $`U`$. Then, the locally free sheaf $`\pi _2\stackrel{~}{}`$ has to be a line bundle. Then the same argument used in b3) proves the remaining statement. ∎
###### Corollary 1.23.
An object $`𝒦^{}`$ in $`D_c^b(X\times Y)`$ satisfying the conditions of Theorem 1.22 is strongly simple over $`X`$ if and only if
1. $`\mathrm{Hom}_{D(Y)}^i(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_{x_1}}),\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_{x_2}}))=0`$ for any pair $`Z_{x_1}`$ and $`Z_{x_2}`$ of l.c.i. zero cycles (supported on $`x_1`$, $`x_2`$ respectively) unless $`x_1=x_2`$ and $`0idimX`$;
2. $`\mathrm{Hom}_{D(Y)}^0(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x),\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x))=k`$ for every point $`xX`$.
From Propositions 1.20 and 1.18 and Corollary 1.7, we obtain:
###### Corollary 1.24.
Let $`X`$ and $`Y`$ be projective integral Gorenstein schemes over an algebraically closed field of characteristic zero, and let $`𝒦^{}`$ be an object in $`D_c^b(X\times Y)`$ of finite projective dimension over both factors. The integral functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$ is an equivalence if and only if $`𝒦^{}`$ is strongly simple over both factors.
###### Remark 1.25.
Theorem 1.22 is false in positive characteristic even in the smooth case. Let $`X`$ be a smooth projective scheme of dimension $`m`$ over a field $`k`$ of characteristic $`p>0`$, and $`F:XX^{(p)}`$ the *relative Frobenius morphism* \[25, 3.1\], which is topologically a homeomorphism. Let $`\mathrm{\Gamma }X\times X^{(p)}`$ be the graph of $`F`$, whose associated integral functor is the direct image $`F_{}:D_c^b(X)D_c^b(X^{(p)})`$. Since $`F_{}(𝒪_x)𝒪_{F(x)}`$, one easily sees that $`\mathrm{\Gamma }`$ is strongly simple over $`X`$. However, $`F_{}(𝒪_X)`$ is a locally free $`𝒪_{X^{(p)}}`$-module of rank $`p^m`$ \[25, 3.2\], so that $`\mathrm{Hom}_{D(X^{(p)})}^0(F_{}(𝒪_X),𝒪_{F(x)})k^{p^m}`$ whereas $`\mathrm{Hom}_{D(X)}^0(𝒪_X,𝒪_x)k`$; thus $`F_{}`$ is not fully faithful. $`\mathrm{}`$
### 1.7. A criterion for equivalence
The usual Bridgeland criterion \[12, Thm. 5.1\] that characterises when an integral functor over the derived category of a smooth variety is an equivalence (or a Fourier-Mukai functor) also works in the Gorenstein case. The original proof is based on the fact that if $`X`$ is smooth, the skyscraper sheaves $`𝒪_x`$ form a spanning class for the derived category $`D_c^b(X)`$ . This is also true for Gorenstein varieties. Moreover in this case there is a more natural spanning class (see for a similar statement), that allows to give a similar criterion.
###### Lemma 1.26.
If $`X`$ is a Gorenstein scheme, then the following sets are spanning classes for $`D_c^b(X)`$:
1. $`\mathrm{\Omega }_1=\{𝒪_x\}`$ for all closed points $`xX`$.
2. $`\mathrm{\Omega }_2=\{𝒪_{Z_x}\}`$ for all closed points $`xX`$ and all l.c.i. zero cycles $`Z_x`$ supported on $`x`$.
###### Proof.
(1) Arguing as in \[12, Lemma 2.2\], one proves that if $`\mathrm{Hom}^i(^{},𝒪_x)=0`$ for every $`i`$ and every $`xX`$, then $`^{}=0`$. Suppose now that $`\mathrm{Hom}^i(𝒪_x,^{})=0`$ for every $`i`$ and every $`xX`$. By (1.6), $`\mathrm{Hom}^i(𝒪_x,^{})\mathrm{Hom}^i(^{},𝒪_x^{})`$ and since $`𝒪_x^{}𝒪_x[m]`$ where $`m=dimX`$ because $`X`$ is Gorenstein, we have that $`\mathrm{Hom}^{im}(^{},𝒪_x)=0`$ for every $`i`$ and every $`xX`$. Then $`^{}=0`$ and from (1.5), one concludes that $`^{}=0`$.
(2) By Proposition 1.16 with $`Y=\mathrm{}`$, if $`\mathrm{Hom}_{D(X)}^i(𝒪_{Z_x},^{})=0`$ for every $`i`$ and every $`Z_x`$, then $`^{}=0`$. On the other hand, since $`𝒪_{Z_x}`$ is of finite homological dimension, Serre duality can be applied to get an isomorphism
$$\mathrm{Hom}^i(^{},𝒪_{Z_x})^{}\mathrm{Hom}^i(^{},𝒪_{Z_x}\omega _X)^{}\mathrm{Hom}^{mi}(𝒪_{Z_x},^{})$$
where $`m=dimX`$. By the first part, if $`^{}`$ is a non-zero object in $`D_c^b(X)`$ the second member is non-zero for some $`i`$ and we finish. ∎
###### Theorem 1.27.
Let $`X`$ and $`Y`$ be projective Gorenstein schemes over an algebraically closed field of characteristic zero. Assume also that $`X`$ is integral and $`Y`$ is connected. If $`𝒦^{}`$ is an object in $`D_c^b(X\times Y)`$ of finite projective dimension over both $`X`$ and $`Y`$, then the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is an equivalence of categories if and only if one has
1. $`𝒦^{}`$ is strongly simple over $`X`$.
2. For every closed point $`xX`$, $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)\omega _Y`$.
###### Proof.
By Proposition 1.17, the functor $`H=\mathrm{\Phi }_{YX}^{𝒦^{}\pi _X^{}\omega _X[m]}`$ is a right adjoint to $`\mathrm{\Phi }_{XY}^𝒦^{}`$ while $`G=\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}`$ is a left adjoint to it. If $`\mathrm{\Phi }_{XY}^𝒦^{}`$ is an equivalence, there is an isomorphism of functors $`HG`$ and then the left adjoints are also isomorphic, that is $`\mathrm{\Phi }_{XY}^𝒦^{}\mathrm{\Phi }_{XY}^{𝒦^{}\pi _Y^{}\omega _Y^1\pi _X^{}\omega _X}`$. Applying this to $`𝒪_x`$ we get $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)\omega _Y\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x)`$.
For the converse, notice first that the derived category $`D^b(Y)`$ is indecomposable because $`Y`$ is connected \[12, Ex. 3.2\]. Then we have to prove that for any object $`^{}`$ in $`D_c^b(Y)`$ the condition $`H(^{})=0`$ implies that $`G(^{})=0`$ \[12, Thm. 3.3\]. Since for every object $`^{}`$ in $`D_c^b(Y)`$ one has a functorial isomorphism
(1.18)
$$G(^{})H(^{}\omega _Y[n])\omega _X^1[m],$$
it is enough to prove that $`H(^{}\omega _Y[n])=0`$. We have
$`\mathrm{Hom}^i(𝒪_x,H(^{}\omega _Y[n]))`$ $`\mathrm{Hom}^i(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x),^{}\omega _Y[n])`$
$`\mathrm{Hom}^{n+i}(\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_x),^{})`$
$`\mathrm{Hom}^{n+i}(𝒪_x,H(^{}))=0`$
and one concludes by Lemma 1.26. ∎
Using now the second part of Lemma 1.26, we prove analogously the following:
###### Theorem 1.28.
Let $`X`$ and $`Y`$ be projective Gorenstein schemes over an algebraically closed field of characteristic zero. Assume also that $`X`$ is integral and $`Y`$ is connected. If $`𝒦^{}`$ is an object in $`D_c^b(X\times Y)`$ of finite projective dimension over both $`X`$ and $`Y`$, then the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is an equivalence of categories if and only if one has
1. $`𝒦^{}`$ is strongly simple over $`X`$.
2. For every closed point $`xX`$ there is a l.c.i. cycle $`Z_x`$ such that $`\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})\mathrm{\Phi }_{XY}^𝒦^{}(𝒪_{Z_x})\omega _Y`$.
###### Remark 1.29.
The second condition in the above lemma can be also written in either the form $`p_2(𝐋j_{Z_x}^{}𝒦^{})p_2(𝐋j_{Z_x}^{}𝒦^{})\omega _Y`$ or the form $`𝐋j_{Z_x}^{}𝒦^{}𝐋j_{Z_x}^{}𝒦^{}p_2^{}\omega _Y`$, where $`p_2:Z_x\times YY`$ is the projection. $`\mathrm{}`$
### 1.8. Geometric applications of Fourier-Mukai functors
As in the smooth case, the existence of a Fourier-Mukai functor between the derived categories of two Gorenstein schemes has important geometrical consequences. In the following proposition, we list some of them.
###### Proposition 1.30.
Let $`X`$ and $`Y`$ be projective Gorenstein schemes and let $`𝒦^{}`$ be an object in $`D_c^b(X\times Y)`$ of finite projective dimension over both $`X`$ and $`Y`$. If the integral functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is a Fourier-Mukai functor, the following statements hold:
1. The right and the left adjoints to $`\mathrm{\Phi }_{XY}^𝒦^{}`$ are functorially isomorphic
$$\mathrm{\Phi }_{YX}^{𝒦^{}\pi _X^{}\omega _X[m]}\mathrm{\Phi }_{YX}^{𝒦^{}\pi _Y^{}\omega _Y[n]}$$
and they are quasi-inverses to $`\mathrm{\Phi }_{XY}^𝒦^{}`$.
2. $`X`$ and $`Y`$ have the same dimension, that is, $`m=n`$.
3. $`\omega _X^r`$ is trivial for an integer $`r`$ if and only if $`\omega _Y^r`$ is trivial. Particularly, $`\omega _X`$ is trivial if and only if $`\omega _Y`$ is trivial. In this case, the functor $`\mathrm{\Phi }_{YX}^𝒦^{}`$ is a quasi-inverse to $`\mathrm{\Phi }_{XY}^𝒦^{}`$.
###### Proof.
(1) Since $`\mathrm{\Phi }_{XY}^𝒦^{}`$ is an equivalence, its quasi-inverse is a right and a left adjoint. The statement follows from Proposition 1.17 using the uniqueness of the adjoints.
(2) Applying the above isomorphism to $`𝒪_{Z_y}`$ where $`Z_y`$ is a l.c.i. zero cycle supported on $`y`$, one obtains $`p_X(𝐋j_{Z_y}^{}𝒦^{})[n]p_X(𝐋j_{Z_y}^{}𝒦^{})\omega _X[m]`$ where $`p_X:X\times Z_yX`$ is the projection. Since the two functors are equivalences, these are non-zero objects in $`D_c^b(X)`$. Let $`q_0`$ be the minimum (resp. maximum) of the $`q`$’s with $`^q(p_X(𝐋j_{Z_y}^{}𝒦^{}))0`$. Since $`^{q_0}(p_X(𝐋j_{Z_y}^{}𝒦^{}))^{q_0+mn}(p_X(𝐋j_{Z_y}^{}𝒦^{}))\omega _X`$ one has $`^{q_0+mn}(p_X(𝐋j_{Z_y}^{}𝒦^{}))\omega _X0`$ which contradicts the minimality (resp. maximality) if $`mn<0`$ (resp. $`>0`$). Thus, $`n=m`$.
(3) If we denote by $`H`$ the right adjoint to $`\mathrm{\Phi }_{XY}^𝒦^{}`$, thanks to (1) and (1.18) we have that $`H(^{})\omega _X^rH(^{}\omega _Y^r)`$ for every $`^{}D_c^b(Y)`$ and every integer $`r`$. If $`\omega _X^r𝒪_X`$, taking $`^{}=𝒪_Y`$ we have $`H(𝒪_Y)H(\omega _Y^r)`$ and applying the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$ to this isomorphism we get $`\omega _Y^r𝒪_Y`$. The converse is similar. ∎
## 2. Relative Fourier-Mukai transforms for Gorenstein morphisms
### 2.1. Generalities and base change properties
Let $`S`$ be a scheme and let $`p:XS`$ and $`q:YS`$ be proper morphisms. We denote by $`\pi _X`$ and $`\pi _Y`$ the projections of the fibre product $`X\times _SY`$ onto its factors and by $`\rho =p\pi _X=q\pi _Y`$ the projection of $`X\times _SY`$ onto the base scheme $`S`$ so that we have the following cartesian diagram
Let $`𝒦^{}`$ be an object in $`D^b(X\times _SY)`$. The relative integral functor defined by $`𝒦^{}`$ is the functor $`\mathrm{\Phi }_{XY}^𝒦^{}:D^{}(X)D^{}(Y)`$ given by
$$\mathrm{\Phi }_{XY}^𝒦^{}(^{})=𝐑\pi _Y(𝐋\pi _X^{}^{}\stackrel{𝐋}{}𝒦^{}).$$
We shall denote this functor by $`\mathrm{\Phi }`$ from now on.
Let $`sS`$ be a closed point. Let us denote $`X_s=p^1(s)`$, $`Y_s=q^1(s)`$, and $`\mathrm{\Phi }_s:D^{}(X_s)D^{}(Y_s)`$ the integral functor defined by $`𝒦_{}^{}{}_{s}{}^{}=𝐋j_s^{}𝒦^{}`$, with $`j_s:X_s\times Y_sX\times _SY`$ the natural embedding.
When the kernel $`𝒦^{}D_c^b(X\times _SY)`$ is of finite homological dimension over $`X`$, the functor $`\mathrm{\Phi }`$ is defined over the whole $`D(X)`$ and it maps $`D_c^b(X)`$ into $`D_c^b(Y)`$. If moreover $`q:YS`$ is flat, then $`𝒦_{}^{}{}_{s}{}^{}`$ is of finite homological dimension over $`X_s`$ for any $`sS`$.
If $`p:XS`$ and $`q:YS`$ are flat morphisms, from the base-change formula we obtain that
(2.1)
$$𝐋j_s^{}\mathrm{\Phi }(^{})\mathrm{\Phi }_s(𝐋j_s^{}^{})$$
for every $`^{}D(X)`$, where $`j_s:X_sX`$ and $`j_s:Y_sY`$ are the natural embeddings. In this situation, base change formula also gives that
(2.2)
$$j_s\mathrm{\Phi }_s(𝒢^{})\mathrm{\Phi }(j_s𝒢^{})$$
for every $`𝒢^{}D(X_s)`$.
Proposition 1.8 allows us to obtain the following result.
###### Lemma 2.1.
Let $`p:XS`$ and $`q:YS`$ be locally projective Gorenstein morphisms, and let $`𝒦^{}`$ be an object in $`D_c^b(X\times _SY)`$ of finite projective dimension over both $`X`$ and $`Y`$. Then the functor
$$H=\mathrm{\Phi }_{YX}^{𝒦^{}\stackrel{𝐋}{}\pi _X^{}\omega _{X/S}[m]}:D_c^b(Y)D_c^b(X)$$
is a right adjoint to the functor $`\mathrm{\Phi }_{XY}^𝒦^{}`$.
### 2.2. Criteria for fully faithfulness and equivalence in the relative setting
In this subsection we work over an algebraically closed field of characteristic zero.
In the relative situation the notion of strongly simple object is the following.
###### Definition 2.2.
Let $`p:XS`$ and $`q:YS`$ be proper Gorenstein morphisms. An object $`𝒦^{}D_c^b(X\times _SY)`$ is *relatively strongly simple* over $`X`$ if $`𝒦_{}^{}{}_{s}{}^{}`$ is bounded and strongly simple over $`X_s`$ for every closed point $`sS`$. $`\mathrm{}`$
###### Lemma 2.3.
Let $`ZS`$ be a proper morphism and $`^{}`$ be an object of $`D_c^b(Z)`$ such that $`𝐋j_s^{}^{}=0`$ in $`D_c^b(Z_s)`$ for every closed point $`s`$ in $`S`$, where $`j_s:Z_sZ`$ is the immersion of the fibre. Then $`^{}=0`$.
###### Proof.
For every closed point $`s`$ in $`S`$ there is a spectral sequence $`E_2^{p,q}=L_pj_s^{}^q(^{})`$ converging to $`E_{\mathrm{}}^{p+q}=^{p+q}(𝐋j_s^{}^{})=0`$. Assume that $`^{}0`$ and let $`q_0`$ be the maximum of the integers $`q`$ such that $`^q(^{})0`$. If $`s`$ is a point in the image of the support of $`^{q_0}(^{})`$, one has that $`j_s^{}^{q_0}(^{})0`$ and every non-zero element in $`E_2^{0,q_0}=j_s^{}^{q_0}(^{})`$ survives to infinity. Then $`E_{\mathrm{}}^{q_0}0`$ and this is impossible. ∎
###### Theorem 2.4.
Let $`p:XS`$ and $`q:YS`$ be locally projective Gorenstein morphisms. Let $`𝒦^{}`$ be an object in $`D_c^b(X\times _SY)`$ of finite projective dimension over both $`X`$ and $`Y`$. The relative integral functor $`\mathrm{\Phi }=\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is fully faithful (resp. an equivalence) if and only if $`\mathrm{\Phi }_s:D_c^b(X_s)D_c^b(Y_s)`$ is fully faithful (resp. an equivalence) for every closed point $`sS`$.
###### Proof.
By Proposition 1.18, if $`\mathrm{\Phi }`$ is fully faithful the unit morphism
$$\mathrm{Id}H\mathrm{\Phi }$$
is an isomorphism (where $`H`$ is the right adjoint given at Lemma 2.1). Then, given a closed point $`sS`$ and $`𝒢^{}D_c^b(X_s)`$, one has an isomorphism $`j_s𝒢^{}(H\mathrm{\Phi })(j_s𝒢^{})`$. Since $`(H\mathrm{\Phi })(j_s𝒢^{})j_s(H_s\mathrm{\Phi }_s)(𝒢^{})`$ by (2.2) and $`j_s`$ is a closed immersion, the unit morphism $`𝒢^{}(H_s\mathrm{\Phi }_s)(𝒢^{})`$ is an isomorphism; this proves that $`\mathrm{\Phi }_s`$ is fully faithful.
Now assume that $`\mathrm{\Phi }_s`$ is fully faithful for any closed point $`sS`$. Let us see that the unit morphism $`\eta :\mathrm{Id}H\mathrm{\Phi }`$ is an isomorphism. For each $`^{}D_c^b(X)`$we have an exact triangle
$$^{}\stackrel{\eta (^{})}{}(H\mathrm{\Phi })(^{})\mathrm{Cone}(\eta (^{}))^{}[1].$$
Then, by (2.1), for every closed point $`s`$ in $`S`$ we have an exact triangle
$$𝐋j_s^{}^{}(H_s\mathrm{\Phi }_s)(𝐋j_s^{}^{})𝐋j_s^{}\mathrm{Cone}(\eta (^{}))𝐋j_s^{}^{}[1].$$
so that $`𝐋j_s^{}[\mathrm{Cone}(\eta (^{}))]\mathrm{Cone}(\eta _s(𝐋j_s^{}^{}))0`$ because $`\eta _s:\mathrm{Id}H_s\mathrm{\Phi }_s`$. We finish by Lemma 2.3.
A similar argument gives the statement about equivalence. ∎
As a corollary of the previous theorem and Theorem 1.22, we obtain the following result.
###### Theorem 2.5.
Let $`p:XS`$ and $`q:YS`$ be locally projective Gorenstein morphisms with integral fibres. Let $`𝒦^{}`$ be an object in $`D_c^b(X\times _SY)`$ of finite projective dimension over each factor. The kernel $`𝒦^{}`$ is relatively strongly simple over $`X`$ (resp. over $`X`$ and $`Y`$) if and only if the functor $`\mathrm{\Phi }=\mathrm{\Phi }_{XY}^𝒦^{}:D_c^b(X)D_c^b(Y)`$ is fully faithful (resp. an equivalence).
### 2.3. Application to Weierstrass elliptic fibrations
In this subsection we work over an algebraically closed field of characteristic zero.
Let $`p:XS`$ be a relatively integral elliptic fibration, that is, a proper flat morphism whose fibres are integral Gorenstein curves with arithmetic genus 1. Generic fibres of $`p`$ are smooth elliptic curves, and the degenerated fibers are rational curves with one node or one cusp. If $`\widehat{p}:\widehat{X}S`$ denotes the dual elliptic fibration, defined as the relative moduli space of torsion free rank 1 sheaves of relative degree 0, it is known that for every closed point $`sS`$ there is an isomorphism $`\widehat{X}_sX_s`$ between the fibers of both fibrations. If we assume that the original fibration $`p:XS`$ has a section $`\sigma :SX`$ taking values in the smooth locus of $`p`$, then $`p`$ and $`\widehat{p}`$ are globally isomorphic. Let us identify from now on $`X`$ and $`\widehat{X}`$ and consider the commutative diagram
The relative Poincaré sheaf is
$$𝒫=_\mathrm{\Delta }\pi _1^{}𝒪_X(H)\pi _2^{}𝒪_X(H)\rho ^{}\omega ^1,$$
where $`H=\sigma (S)`$ is the image of the section and $`\omega =R^1p_{}𝒪_X(p_{}\omega _{X/S})^1`$.
Relatively integral elliptic fibrations have a Weierstrass form \[31, Lemma II.4.3\]: The line bundle $`𝒪_X(3H)`$ is relatively very ample and if $`=p_{}𝒪_X(3H)𝒪_S\omega ^2\omega ^3`$ and $`\overline{p}:(^{})=\mathrm{Proj}(S^{}())S`$ is the associated projective bundle, there is a closed immersion $`j:X(^{})`$ of $`S`$-schemes such that $`j^{}𝒪_{(^{})}(1)=𝒪_X(3H)`$. In particular, $`p`$ is a projective morphism.
###### Lemma 2.6.
The relative Poincaré sheaf $`𝒫`$ is of finite projective dimension and relatively strongly simple over both factors.
###### Proof.
By the symmetry of the expression of $`𝒫`$ it is enough to prove that $`𝒫`$ is of finite projective dimension and strongly simple over the first factor. For the first claim, it is enough to prove that $`_\mathrm{\Delta }`$ has finite projective dimension. Let us consider the exact sequence
$$0_\mathrm{\Delta }𝒪_{X\times _SX}\delta _{}𝒪_X0$$
where $`\delta :XX\times _SX`$ is the diagonal morphism. It suffices to see that $`\delta _{}𝒪_X`$ has finite projective dimension. We have to prove that for any $`𝒩^{}D^b(X)`$, the complex $`𝐑om_{𝒪_{X\times _SX}}^{}(\delta _{}𝒪_X,\pi _1^!𝒩^{})`$ is bounded. This is a complex supported at the diagonal, so that it suffices to see that $`𝐑\pi _1𝐑om_{𝒪_{X\times _SX}}^{}(\delta _{}𝒪_X,\pi _1^!𝒩^{})`$ is bounded. This follows from the following formulas.
$$𝐑\pi _1𝐑om_{𝒪_{X\times _SX}}^{}(\delta _{}𝒪_X,\pi _1^!𝒩^{})𝐑om_{𝒪_X}^{}(𝒪_X,𝒩^{})𝒩^{}.$$
Let us prove that $`𝒫`$ is strongly simple over the first factor. Fix a closed point $`sS`$ and consider two different points $`x_1`$ and $`x_2`$ in the fiber $`X_s`$. If both are non-singular, then
$$\mathrm{Hom}_{D(X_s)}^i(\mathrm{\Phi }_{X_sX_s}^{𝒫_s}(𝒪_{x_1}),\mathrm{\Phi }_{X_sX_s}^{𝒫_s}(𝒪_{x_2}))H^i(X_s,𝒪_{X_s}(x_1x_2))=0\text{for every }i$$
because $`𝒪_{X_s}(x_2x_1)`$ is a non-trivial line bundle of degree zero. Assume that $`x_2`$ is singular and $`x_1`$ is not, the other case being similar. Let $`Z_{x_2}`$ be a l.c.i zero cycle supported on $`x_2`$. We have
$$\mathrm{Hom}_{D(X_s)}^i(\mathrm{\Phi }_{X_sX_s}^{𝒫_s}(𝒪_{x_1}),\mathrm{\Phi }_{X_sX_s}^{𝒫_s}(𝒪_{Z_{x_2}}))=H^i(X_s,𝒥_{Z_{x_2}}𝒪_{X_s}(x_1))$$
where $`𝒥_{Z_{x_2}}`$ denotes the direct image by the finite morphism $`Z_{x_2}\times X_sX_s`$ of the ideal sheaf of the graph $`Z_{x_2}Z_{x_2}\times X_s`$ of the immersion $`Z_{x_2}X_s`$.
Let us consider the exact sequences of $`𝒪_{X_s}`$-modules
$`0𝒥_{Z_{x_2}}𝒪_{Z_{x_2}}`$ $`_k𝒪_{X_s}𝒪_{Z_{x_2}}0`$
$`0𝒥_{Z_{x_2}}(x_1)𝒪_{Z_{x_2}}`$ $`_k𝒪_{X_s}(x_1)𝒪_{Z_{x_2}}0`$
Since $`H^0(X_s,𝒪_{X_s})k`$ the morphism $`𝒪_{Z_{x_2}}_kH^0(X_s,𝒪_{X_s})𝒪_{Z_{x_2}}`$ of global sections induced by the first sequence is an isomorphism. Moreover, $`H^0(X_s,𝒪_{X_s})H^0(X_s,𝒪_{X_s}(x_1))`$ and then we also have an isomorphism of global sections $`𝒪_{Z_{x_2}}_kH^0(X_s,𝒪_{X_s})\stackrel{}{}𝒪_{Z_{x_2}}_kH^0(X_s,𝒪_{X_s}(x_1))`$. Thus, $`𝒪_{Z_{x_2}}_kH^0(X_s,𝒪_{X_s}(x_1))\stackrel{}{}𝒪_{Z_{x_2}}`$ so that $`H^i(X_s,𝒥_{Z_{x_2}}(x_1))=0`$ for $`i=0,1`$.
Finally, since $`\mathrm{Hom}_{D(X_s)}^0(𝒫_x,𝒫_x)=k`$ for every point $`xX_s`$, we conclude that $`𝒫_s`$ is strongly simple over $`X_s`$. ∎
Now by Corollary 1.24 we have
###### Proposition 2.7.
The relative integral functor
$$\mathrm{\Phi }_{XX}^𝒫:D_c^b(X)D_c^b(X)$$
defined by the Poincaré sheaf is an equivalence of categories.
Notice that the proof of this result does not use spanning classes.
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# CHSH violation and entropy – concurrence plane
## I Introduction
Physically allowed degree of entanglement and mixture for two - qubit mixed states were investigated by Munro et al. in terms of concurrence $`C(\rho )`$ and normalized linear entropy $`S_\mathrm{L}(\rho )`$ MJWK . These authors characterize the subset $`\mathrm{\Lambda }`$ on $`(C,S_\mathrm{L})`$ plane corresponding to possible states of the system and in particular identify maximally entangled mixed states $`\rho _{\mathrm{MEMS}}`$ laying on the boundary of that subset. The states $`\rho _{\mathrm{MEMS}}`$ have maximal allowed entanglement for a given degree of mixedness. They obtained this result analytically for some class of states. Numerical results suggests correctness of the picture for general two-qubit states.
In this note, we consider the problem of violation of Bell – CHSH inequalities Bell ; CHSH for mixed states. It is well known that quantum states violating these inequalities have to be entangled E ; S , but on the other hand, CHSH violation is not necessary for mixed state entanglement Werner . In the context of the results of Ref.MJWK , we address the following question: what are the subsets of $`\mathrm{\Lambda }`$ which correspond to states violating Bell – CHSH inequalities ? In our previous publication DJ , we have studied the structure of such subsets in the case of specific class of quantum states. The results show that $`\mathrm{\Lambda }`$ is a sum of disjoint subsets $`\mathrm{\Lambda }_\mathrm{V},\mathrm{\Lambda }_{\mathrm{NV}}`$ and $`\mathrm{\Lambda }_0`$ with the following properties: states belonging to $`\mathrm{\Lambda }_\mathrm{V}`$ violate CHSH inequalities, whereas states from $`\mathrm{\Lambda }_{\mathrm{NV}}`$ fulfil all CHSH inequalities. The subset $`\mathrm{\Lambda }_0`$ has somehow unexpected property: for any pair $`(S_\mathrm{L},C)\mathrm{\Lambda }_0`$ there are two families of states with the same entropy and concurrence such that all states from one family violate CHSH and at the same time, all states from the other family fulfil all CHSH inequalities. In the present paper, we continue these investigations for general class of two-qubit states. First we consider larger class of states still admitting explicit formulas for linear entropy, concurrence and degree of CHSH violation. Unfortunately, analytical analysis of the relation between these functions is not possible. Numerical investigations lead to some modifications of the picture from Ref. DJ , but the general structure is not changed. Finally, this problem is studied using numerically generated arbitrary density matrices. The results indicate that the structure of $`\mathrm{\Lambda }`$ for general two-qubit states seems to be the same.
We consider also the problem of maximal violation of CHSH inequalities. For a class of mixed states we obtain counterpart of the result of Ref.MJWK , namely we find the form of mixed states with maximal degree of CHSH violation for given linear entropy. All that states lie on specific curve on entropy – concurrence plane. Numerical results suggest also that general two – qubits states with this property satisfy the same relation between entropy and concurrence. As we show, maximal violation of Bell inequalities with fixed linear entropy is not equivalent to the maximal entanglement under the same conditions.
## II CHSH inequalities
Let $`𝒂,𝒂^{},𝒃,𝒃^{}`$ be the unit vectors in $`^3`$ and $`𝝈=(\sigma _1,\sigma _2,\sigma _3)`$. Consider the family of operators on two - qubits Hilbert space $`_{AB}=^4`$
$$B_{CHSH}=𝒂𝝈(𝒃+𝒃^{})𝝈+𝒂^{}𝝈(𝒃𝒃^{})𝝈$$
(II.1)
Then Bell - CHSH CHSH inequalities are
$$|\mathrm{tr}(\rho B_{CHSH})|2$$
(II.2)
If the above inequality is not satisfied by the state $`\rho `$ for some choice of $`𝒂,𝒂^{},𝒃,𝒃^{}`$ , we say that $`\rho `$ violates Bell-CHSH inequalities. In the case of two-qubits, the violation of Bell - CHSH inequalities by mixed states can be studied using simple necessary and sufficient condition HHH ; H . Consider real matrix
$$T_\rho =(t_{nm}),t_{nm}=\mathrm{tr}(\rho \sigma _n\sigma _m)$$
(II.3)
and real symmetric matrix
$$U_\rho =T_\rho ^TT_\rho $$
(II.4)
where $`T_\rho ^T`$ is the transposition of $`T_\rho `$. Let
$$m(\rho )=\underset{j<k}{\mathrm{max}}(u_j+u_k)$$
(II.5)
and $`u_j,j=1,2,3`$ are the eigenvalues of $`U_\rho `$. As was shown in HHH ; H
$$\underset{B_{CHSH}}{\mathrm{max}}\mathrm{tr}(\rho B_{CHSH})=2\sqrt{m(\rho )}$$
(II.6)
Thus (II.2) is violated by some choice of $`𝒂,𝒂^{},𝒃,𝒃^{}`$ if and only if $`m(\rho )>1`$.
We need also the measures of degree of entanglement and mixture for given state. In the case of two qubits, the useful measure of degree of entanglement is concurrence $`C(\rho )`$
$$C(\rho )=\mathrm{max}(0,\mathrm{\hspace{0.17em}2}\lambda _{\mathrm{max}}(\widehat{\rho })\mathrm{tr}\widehat{\rho })$$
(II.7)
where $`\lambda _{\mathrm{max}}(\widehat{\rho })`$ is the maximal eigenvalue of $`\widehat{\rho }`$ and
$$\widehat{\rho }=\sqrt{\sqrt{\rho }\stackrel{~}{\rho }\sqrt{\rho }},\stackrel{~}{\rho }=(\sigma _2\sigma _2)\overline{\rho }(\sigma _2\sigma _2)$$
with $`\overline{\rho }`$ denoting complex conjugation of the matrix $`\rho `$. It is known that $`C(\rho )`$ can be used to obtain entanglement of formation, which is natural measure of entanglement for mixed states HW ; Woo . To measure degree of mixture, or deviation from pure state, we use linear entropy
$$S_\mathrm{L}(\rho )=\frac{4}{3}(1\mathrm{tr}\rho ^2)$$
(II.8)
which is normalized such that its maximal value equals $`1`$.
## III CHSH inequalities and entropy – concurrence plane
To study the subset of entropy – concurrence plane corresponding to violation of CHSH inequalities, consider first the following class $`_0`$ of states
$$\rho =\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& a& \frac{1}{2}ce^{i\theta }& 0\\ 0& \frac{1}{2}ce^{i\theta }& b& 0\\ 0& 0& 0& 1ab\end{array}\right)$$
(III.1)
where
$$c[0,1],a,b0,\theta [0,2\pi ]$$
and
$$ab\frac{c^2}{4},a+b1$$
Notice that for $`\rho _0`$
$$C(\rho )=c$$
Define the subset $`\mathrm{\Lambda }^2`$
$$\mathrm{\Lambda }=\{(S_L(\rho ),C(\rho )):C(\rho )>0\text{and}\rho _0\}$$
(III.2)
In the paper DJ , we analysed the set (III.2) and we have shown that it is a sum of disjoint subsets $`\mathrm{\Lambda }_\mathrm{V},\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_{\mathrm{NV}}`$ with properties:
1. If $`(s,c)\mathrm{\Lambda }_\mathrm{V}`$, then every state $`\rho _0`$ such that $`S_L(\rho )=s`$ and $`C(\rho )=c`$ satisfies $`m(\rho )>1`$.
2. If $`(s,c)\mathrm{\Lambda }_0`$, then there exist states $`\rho _1,\rho _2_0`$ such that
$$S_L(\rho _1)=S_L(\rho _2)=s,C(\rho _1)=C(\rho _2)=c$$
and $`m(\rho _1)>1`$, but $`m(\rho _2)<1`$.
3. If $`(s,c)\mathrm{\Lambda }_{\mathrm{NV}}`$, then every state $`\rho _0`$ such that $`S_L(\rho )=s`$ and $`C(\rho )=c`$ satisfies $`m(\rho )<1`$.
Detailed analytic description of regions $`\mathrm{\Lambda }_\mathrm{V},\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_{\mathrm{NV}}`$ as well as the proof of the above properties, can be found in Ref. DJ .
In the present paper, we try to extend this analysis to the larger class of two - qubit states. For general mixed two-qubit states there is a bound on concurrence that guarantees violation of CHSH inequalities irrespective of linear entropy. It follows from the result of Verstraete and Wolf VW that minimal violation of CHSH inequality for given concurrence $`C`$ is equal to
$$m_{\mathrm{min}}=\mathrm{max}(1,\mathrm{\hspace{0.17em}2}C^2)$$
(III.3)
So if
$$C(\rho )>\frac{1}{\sqrt{2}}$$
(III.4)
then minimal value of $`m`$ is greater then $`1`$, and every state satisfying (III.4) violates CHSH inequality. On the other hand, there is a bound on linear entropy that guarantees fulfilling CHSH inequalities irrespective of concurrence. It is given by the result of Santos Santos that all states with (normalized) linear entropy
$$S_\mathrm{L}(\rho )>\frac{2}{3}$$
(III.5)
satisfy all CHSH inequalities. We see that these general bounds are compatible with our previous analysis (FIG. 1) and possible modifications can occur in region $`\mathrm{\Lambda }_\mathrm{V}`$ below the line $`C=\frac{1}{\sqrt{2}}`$ and in region $`\mathrm{\Lambda }_{\mathrm{NV}}`$ on the left hand side of the line $`s=\frac{2}{3}`$.
Consider now the larger class $`_1`$ of states of the form
$$\rho =\left(\begin{array}{cccc}\rho _{11}& 0& 0& \rho _{14}\\ 0& \rho _{22}& \rho _{23}& 0\\ 0& \rho _{32}& \rho _{33}& 0\\ \rho _{41}& 0& 0& \rho _{44}\end{array}\right)$$
(III.6)
One can check that for that class
$$C(\rho )=\mathrm{max}(0,C_1,C_2)$$
(III.7)
where
$$\begin{array}{cc}& C_1=2(|\rho _{14}|\sqrt{\rho _{22}\rho _{33}})\hfill \\ & C_2=2(|\rho _{23}|\sqrt{\rho _{11}\rho _{44}})\hfill \end{array}$$
(III.8)
and
$$S_\mathrm{L}(\rho )=1\rho _{11}^2\rho _{22}^2\rho _{33}^2\rho _{44}^22|\rho _{14}|^22|\rho _{23}|^2$$
(III.9)
Moreover,
$`m(\rho )=`$ $`\mathrm{max}[\mathrm{\hspace{0.17em}4}(|\rho _{41}|+|\rho _{23}|)^2,4(|\rho _{41}||\rho _{23}|)^2,`$ (III.10)
$`(\rho _{11}\rho _{22}\rho _{33}+\rho _{44})^2]`$
For the class (III.6), analytic description of regions of $`\mathrm{\Lambda }`$ where $`m(\rho )>1`$ or $`m(\rho )<1`$ is not possible, but it can be done numerically. Results of numerical analysis of the class $`_1`$ are presented on FIG. 2. We see that the region $`\mathrm{\Lambda }_\mathrm{V}`$ where all states violate CHSH inequalities is not changed, but there are states with $`m(\rho )>1`$ for some points in $`\mathrm{\Lambda }_{\mathrm{NV}}`$, so the region $`\mathrm{\Lambda }_0`$ is slightly enlarged.
To study this problem for general two-qubit density matrices, we numerically generate $`310^6`$ randomly chosen density matrices. For such two-qubit states the structure of the set of pairs $`(S_\mathrm{L},C)`$ is very simple. There are only points corresponding to $`m(\rho )>1`$ or $`m(\rho )<1`$ (FIG. 3). Notice that by the method of random choice of states, not all points of $`\mathrm{\Lambda }`$ are achieved (the boundary of generated set correspond exactly to the class of Werner states), but the obtained structure is compatible with previous results. To have some insight into the properties of the remaining part of the set $`\mathrm{\Lambda }`$, we modify the method of generation of states and consider density matrices lying close to the boundary of the set of all states i.e. such $`\rho `$ that one of its eigenvalues is almost equal to zero. For these randomly generated states, the pairs $`(S_\mathrm{L},C)`$ cover the whole set $`\mathrm{\Lambda }`$, and its structure is the same as for the class $`_1`$ (FIG. 4). The results suggest that the picture obtained using the class $`_1`$ should be correct also for all two-qubit density matrices, although for most of randomly chosen density matrices $`\rho `$ with fixed mixedness and linear entropy, either $`m(\rho )>1`$ or $`m(\rho )<1`$. Unfortunately, we do not know analytic description of the boundary of the enlarged region $`\mathrm{\Lambda }_0`$.
## IV Maximal CHSH violation
Consider now the values of $`m(\rho )`$ for states violating Bell – CHSH inequalities. We are especially interested in maximal values of $`m(\rho )`$. For the class (III.1)
$$m(\rho )=\mathrm{max}(\mathrm{\hspace{0.17em}2}c^2,(2(a+b)1)^2+c^2)$$
(IV.1)
We see that (IV.1) is maximal iff $`a+b=1`$, and then
$$m(\rho )=1+c^2$$
(IV.2)
By general result of Verstraete and Wolf VW , for any two-qubit state, (IV.2) is the maximal degree of CHSH violation for given concurrence $`c`$. But we ask another question: what is the maximum of (IV.2) for fixed linear entropy, and which states realize that maximum? We can simply answer this question for the class of states (III.1). Since
$$S_\mathrm{L}(\rho )=\frac{4}{3}\left(1a^2(1a)^2\frac{c^2}{2}\right)$$
(IV.3)
so fixing $`S_\mathrm{L}(\rho )=s`$, we obtain
$$m(\rho )=1+4(aa^2)\frac{3}{2}s$$
(IV.4)
Maximum of (IV.4) is achieved at $`a=\frac{1}{2}`$ and equals to $`2\frac{3}{2}s`$ for $`s[0,\frac{2}{3}]`$. In this way we obtain
###### Theorem 1
In the class (III.1), states maximizing degree of violation of CHSH inequalities for fixed linear entropy, lie on the curve
$$s=\frac{2}{3}(1c^2)$$
and have the form
$$\rho _{\mathrm{MVB}}=\frac{1}{2}\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 1& \sqrt{\beta 1}e^{i\theta }& 0\\ 0& \sqrt{\beta 1}e^{i\theta }& 1& 0\\ 0& 0& 0& 0\end{array}\right)$$
(IV.5)
with
$$\beta [1,2],\theta [0,2\pi ]$$
Moreover,
$$m(\rho _{\mathrm{MVB}})=\beta $$
It is instructive to compare the value of $`m(\rho _{\mathrm{MVB}})`$ with degree of CHSH violation for some other classes of states. Let $`W`$ be the family of Werner states
$$W=(1p)\frac{𝕀_4}{4}+p|\mathrm{\Psi }^{}\mathrm{\Psi }^{}|$$
(IV.6)
where $`\mathrm{\Psi }^{}`$ is a singlet state of two-qubits. Then
$$m(W)=22s$$
(IV.7)
For maximally entangled mixed states $`\rho _{\mathrm{MEMS}}`$ introduced in Ref. MJWK , the corresponding value of $`m`$ is given by
$$m(\rho _{\mathrm{MEMS}})=1\frac{3}{4}s+\sqrt{1\frac{3}{2}s}$$
(IV.8)
We see that for a fixed linear entropy
$$m(\rho _{\mathrm{MVB}})m(\rho _{\mathrm{MEMS}})m(W)$$
(IV.9)
although
$$C(\rho _{\mathrm{MEMS}})C(W)C(\rho _{\mathrm{MVB}})$$
(IV.10)
So
###### Corollary 1
The states with maximum amount of entanglement for a given linear entropy do not maximize degree of Bell – CHSH violation.
The family of states with maximal degree of violation of Bell – CHSH inequalities has another remarkable property. It is known that fidelity of state $`\rho `$ defined as
$$F(\rho )=\mathrm{max}\psi ,\rho \psi $$
(IV.11)
where the maximum is taken over all maximally entangled pure states $`\psi `$ is bounded above by VV
$$F(\rho )\frac{1+C(\rho )}{2}$$
(IV.12)
By direct computation, one can check that
$$F(\rho _{\mathrm{MVB}})=\frac{1+C(\rho _{\mathrm{MVB}})}{2}$$
Thus
###### Corollary 2
The states $`\rho _{\mathrm{MVB}}`$ maximize fidelity for given concurrence.
For a larger class $`_1`$ we have studied $`m(\rho )`$ as a function of $`S_\mathrm{L}`$ and $`C`$ numerically. Again the results agree with those obtained analytically for the class $`_0`$ (FIG. 6).
###### Acknowledgements.
L.J. acknowledges financial support by Polish Ministry of Scientific Research and Information Technology under the grant PBZ-Min-008/PO3/2003.
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# Wiggly cosmic strings accrete dark energy
## I Introduction
Cosmic strings are known as topological defects that occur in theories with spontaneous symmetry breaking of a local U(1) gauge symmetry which were formed during phase transitions in the early universe. They are trapped infinitely long very thin tubes filled with a previous false-vacuum phase characterized by a energy density per unit length $`\mu _0=T_0\sigma ^2`$, with $`T_0`$ the string tension and $`\sigma `$ the symmetry breaking scale, immersed in the a true-vacuum phase created after the phase transition . Typically, cosmic strings have been hypothesized as the seeds for ulterior galaxy formation or as the cosmic sites where primordial inflation took place . At any event, among the different theoretical objects that are thought to have populated the universe at some previous or current periods, which also include black holes with distinct sizes, Lorentzian wormholes or ringholes, etc., cosmic strings are the sole objects whose existence has been confirmed in the laboratory . Therefore, the great interest that cosmic strings raised when they were first introduced in cosmology has remained alive all the way up to now. Two general kinds of cosmic strings have been so far considered, straight strings and string loops . In the present paper we shall investigate how distinct forms of dark energy can be accreted by cosmic strings. We shall restrict ourselves to consider only straight cosmic strings. These are usually described by a static space-time exterior metric, first derived by Vilenkin
$$ds^2=dt^2+dr^2+dz^2+(18G\mu _0)r^2d\varphi ^2.$$
(1)
By defining a new cylindrical angular coordinate $`\varphi ^{}=(18G\mu _0)\varphi `$, it can immediately be seen that this metric corresponds to a flat spacetime with a conical singularity that is associated with a deficit angle given by $`\mathrm{\Delta }=8\pi G\mu _0`$.
However, an incoming or outgoing energy flux due to dark energy accretion is no longer strictly possible for an exterior locally flat metric like that of a motionless straight string having no wiggles . In fact, a motionless string with no wiggles cannot accrete anything that is motionless and homogeneous around it -in particular it could not accrete dark energy. Thus, if we want to consider accretion of dark energy onto cosmic strings we need these cosmic string to be perturbed by wiggles. In that case the exterior string metric can no longer be given by the locally flat line element (1.1), as wiggle-induced variations of the string mass per unit length and tension would convert these quantities into space-time dependent functions, $`\mu `$ and $`T`$, with the state equation $`\mu T=\mu _0^2`$ and $`\mu >T`$, whose values can initially be considered to be very similar to each other and therefore also very similar to their unperturbed counterparts in the linear approximation . The linearized wiggly string metric reads
$`ds^2=\left[1+4G\left(\mu T\right)\mathrm{ln}{\displaystyle \frac{r}{r_0}}\right]dt^2+dr^2`$
$`+\left[14G\left(\mu T\right)\mathrm{ln}{\displaystyle \frac{r}{r_0}}\right]dz^2`$
$`+\left[14G\left(\mu +T\right)\right]r^2d\theta ^2,`$ (2)
which, contrary to metric (1.1), produces a non-vanishing Newtonian potential. In this case the deficit angle is given by $`4\pi G(\mu +T)`$.
Nowadays cosmology, on the other hand, relies mainly on the idea that the total energy of the current universe and possibly that of the early universe (that is the two cosmic periods known to show accelerating expansion) is dominated by some form of the so-called dark energy . It is therefore of interest to investigate the effects that dark energy may cause in cosmic strings. Following the recent studies performed on black holes , one can actually suppose that dark energy can also be accreted onto a cosmic string, inducing some variation in its energy density per unit length $`\mu `$. This work aims at considering the effects that the accretion of dark energy may have in the fate of wiggly straight cosmic string in an accelerating universe. We shall represent dark energy as a perfect fluid characterized by a negative parameter $`1/3>w=p/\rho `$ (with $`p`$ the pressure and $`\rho `$ the energy density) filling a Friedmann-Robertson-Walker universe whose scale factor is given by
$$a(t)=a_0\left(1+\frac{3}{2}(1+w)C^{1/2}(tt_0)\right)^{2/[3(1+w)]},$$
(3)
where $`C=8\pi G\rho _0/3`$ and we have taken for the energy density $`\rho =\rho _0a^{3(1+w)}`$, with $`\rho _0`$ an integration constant, if we adopt a general quintessence model. It can be readily seen that, whereas the universe enters a steady regime of accelerating expansion which keeps it being finite all the way up to an infinite time if $`w>1`$, in the case that $`w<1`$ (a case at which the dark energy is called phantom energy ) the universe would expand along super-accelerated patterns that drive it to a singularity at a finite time in the future at which everything -even the elementary particles - loses any independent, local behavior by its own to be ripped apart under the sole influence of the global phantom cosmological law. This singularity has been dubbed the big rip and takes place at a time
$$t_{}=t_0+\frac{2}{3(|w|1)C^{1/2}}.$$
(4)
Such a rather weird behavior takes also place when the other main contender model for dark energy, that is to say the K-essence model , is assumed to dominate. In fact, if $`w<1`$ we obtain in this case
$$a(t)(tt_{})^{2\beta /[3(1\beta )]},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}<\beta <1,$$
(5)
where $`t_{}`$ again represents the time for the big rip which is an arbitrary parameter in this case.
Of particular interest is the scenario in which we consider that the dark energy is given in terms of a generalized Chaplygin gas having an equation of state $`p=A\rho ^\alpha `$, with $`A>0`$ and $`\alpha `$ a parameter . In this case the cosmic time $`t`$ relates to the scale factor by the more complicated expression
$`tt_0={\displaystyle \frac{2\left[1+\frac{B}{A}a^{3(1+\alpha )}\right]^{\frac{1+2\alpha }{2(1+\alpha )}}}{3(1+2\alpha )\sqrt{CA^{1/(1+\alpha )}/\rho _0}}}\times `$
$`F(1,{\displaystyle \frac{1+2\alpha }{2(1+\alpha )}};{\displaystyle \frac{3+4\alpha }{2(1+\alpha )}};1+{\displaystyle \frac{B}{A}}a^{3(1+\alpha )}),`$ (6)
with $`F`$ a hypergeometric function and $`B=(\rho ^{1+\alpha }A)a_0^{3(1+\alpha )}`$. It can be seen that even in the phantom energy regime a generalized Chaplygin gas does not lead to any big rip singularity in the future , but it always drives a steady regular accelerating expansion for the universe.
In this paper we use a formalism which is able to encompass the accretion of dark energy described by any of the above models onto wiggly straight cosmic strings. We obtain that as quintessence or K-essence dark energy is accreted onto a perturbed straight cosmic string the energy density per unit length of this string either progressively increases up to a constant finite value if $`w>1`$, or steadily decreases down to the unperturbed value first and might then enter a region where quantum accretion makes it reach a minimum value, quite before the occurrence of the big rip singularity if $`w<1`$. The behavior of the strings when they accrete Chaplygin gas is similar: their energy density per unit length also progressively either increases or decreases toward a extremal value, depending on whether the dominant energy condition is satisfied or violated.
The paper can be outlined as follows. In Sec. II we present the general formalism for the accretion of dark energy onto straight wiggly cosmic strings and obtain a general rate equation for the string core energy density per unit length, $`\mu `$, in terms of the internal dark energy, and apply such a formalism to quintessence and K-essence cosmological fields, so as to the generalized Chaplygin gas model. Approximate expressions of $`\mu `$ as a function of time for the first two dark energy models are also derived, both for $`w=p/\rho >1`$ and $`w=p/\rho <1`$, analyzing the corresponding evolution of the cosmic strings. Finally we conclude and add some further comments in Sec. III.
## II Dark energy accretion onto wiggly straight cosmic strings
We shall consider next how the general accretion theory can be applied to the case in which dark energy is accreted onto wiggly cosmic strings. We shall generally follow the procedure put forward by Babichev, Dokuchaev and Eroschenko for the case of Schwarzschild black holes, generalizing it to the case of straight wiggly cosmic strings. Thus, we start by integrating the energy-momentum conservation law by using the exterior metric (1.2). Although for metric (1.1) there are only two non-vanishing components of the Christoffel symbols, $`\mathrm{\Gamma }_{\theta \theta }^r=\left(18G\mu \right)r`$ and $`\mathrm{\Gamma }_{r\theta }^\theta =1/r`$, when the string is perturbed with wiggles there will be twenty one generally non-vanishing components of the Christoffel symbols which make the calculation to follow more complicated. For a cylindrical symmetry we then have from the time-component of the conservation law of the energy-momentum tensor, $`T_{\mu ;\nu }^\nu =0`$,
$$\sqrt{\mu }ru\sqrt{1h_{00}}\sqrt{1b}(1+h_{00})\sqrt{u^21}(p+\rho )=C,$$
(7)
where
$$h_{00}=4G(\mu T)\mathrm{ln}\left(r/r_0\right)$$
(8)
$$b=4G(\mu +T),$$
(9)
with $`r_0`$ and $`C`$ integration constants and $`u=dr/ds`$.
After integrating the conservation law for the energy-momentum tensor projected onto the four-velocity, $`u_\mu T_{;\nu }^{\mu \nu }=0`$, we also obtain
$$ur\sqrt{\mu (1h_{00}^2)(1b)}e^{_\rho _{\mathrm{}}^\rho \frac{d\rho }{p+\rho }}=A,$$
(10)
where we have taken into account that $`u`$ should be positive for incoming energy flux in this case, and $`A`$ is a positive constant. From Eqs. (2.1) and (2.4) we can then get
$$\sqrt{(u^21)(1+h_{00})}(p+\rho )e^{_\rho _{\mathrm{}}^\rho \frac{d\rho }{p+\rho }}=C_2,$$
(11)
in which the constant $`C_2`$ can be expressed as $`C_2=C/A=\widehat{A}\left[\rho _{\mathrm{}}+p(\rho _{\mathrm{}})\right]`$, with $`\widehat{A}>0`$ a constant, for the cylindrical symmetry used.
By integrating now the momentum density $`T_0^r`$ over the circular length element of the cylinder we can obtain the rate of change of the energy per unit length of the wiggly cosmic string, so that
$$\dot{\mu }=_0^{2\pi }rT_0^r𝑑\varphi =_0^{2\pi }r(p+\rho )(1+h_{00})\frac{dt}{ds}\frac{dr}{ds}𝑑\varphi .$$
(12)
Using then the property $`\sqrt{1+h_{00}}dt=\sqrt{\frac{dr^2}{ds^2}1}ds`$ stemming from the cylindrical symmetry being used and Eqs. (2.5) and (2.6), we finally derive the relevant rate equation for the energy density of a wiggly cosmic string
$$\dot{\mu }=\frac{2\pi \overline{A}\left[\rho _{\mathrm{}}+p(\rho _{\mathrm{}})\right]}{\sqrt{\mu (1b)(1h_{00}^2)}},$$
(13)
with $`\overline{A}=A\widehat{A}>0`$ a constant. Therefore, one has the following integral expression that governs the evolution of the wiggled mass per unit length of the cosmic string
$$_{\mu _i}^\mu \sqrt{\mu (1b)(1h_{00}^2)}𝑑\mu =2\pi \overline{A}_{t_0}^t\left[\rho _{\mathrm{}}+p\left(\rho _{\mathrm{}}\right)\right]𝑑t.$$
(14)
It is worth noticing that the above expressions restrict by themselves the interval along which the quantity $`\mu `$ is allowed to vary on its real values. In fact, one can derive the two conditions
$$\mu _0<\frac{1}{8G}$$
(15)
$$\frac{1\sqrt{164G^2\mu _0^2}}{8G}<\mu <\frac{1+\sqrt{164G^2\mu _0^2}}{8G}.$$
(16)
Condition (2.9) expresses nothing but the impossibility for an supermassive wiggleless cosmic string to reach a linear energy density larger than nearly $`1/G`$. Even though the concepts of radius and mass per unit length for a source like the string core are not unambiguously defined , specially in the presence of an interacting dark energy fluid, at the extreme supermassive case $`\mu =1/8G`$ one would expect the string to no longer exist because it then corresponded to the situation where all the exterior broken phase is collapsed into the core, leaving a pure false-vacuum phase in which the picture of a cosmic string with a core region of trapped is lost . When the string is wiggled then condition (2.9) reflects into condition (2.10) by which it is seen that a wiggly cosmic string cannot exceed a given maximum value or be less than a given minimum nonzero value. If a cosmic string has the extreme supermassive linear mass density, then it cannot be wiggled nor accrete any kind of dark energy.
Now, the integration in the left-hand-side of Eq. (2.8) appears to be very difficult to perform and, in fact, we have been unable to obtain an integrated expression from it in closed form. Nevertheless, in the physically relevant cases that $`\mu `$ is very close to $`\mu _0`$ and/or $`r`$ is very close to $`r_0`$, that term can be integrated to approximately give
$`{\displaystyle _{\mu _i}^\mu }𝑑\mu \sqrt{\mu (1b)(1h_{00}^2)}I(\mu )=`$
$`[{\displaystyle \frac{8G\mu 1}{16G}}\sqrt{4G\mu ^2+\mu 4G\mu _0^2}`$
$`+{\displaystyle \frac{64G^2\mu _0^21}{64G^{3/2}}}\mathrm{arcsin}\left({\displaystyle \frac{18G\mu }{\sqrt{164G^2\mu _0^2}}}\right)\left]\right|_{\mu _i}^\mu .`$ (17)
The integration of the right-hand-side of Eq. (2.8) will be performed in what follows for the distinct dark energy models considered in the Introduction.
### II.1 Quintessence and K-Essence
Starting with the equation of state $`p=w\rho `$, where $`w`$ is assumed constant, we can use the conservation of cosmic energy to finally derive
$$\rho =\rho _0\left(\frac{a_0}{a}\right)^{3(1+w)},$$
(18)
with $`\rho _0`$ and $`a_0`$ constants. Hence
$`2\pi \overline{A}{\displaystyle _{t_0}^t}\left[\rho _{\mathrm{}}+p\left(\rho _{\mathrm{}}\right)\right]𝑑t=`$
$`2\pi \overline{A}(1+w)\rho _0a_0^{3(1+w)}{\displaystyle _{t_0}^t}𝑑ta^{3(1+w)}.`$ (19)
We then have for the scale factor (1.3) corresponding to a general flat quintessence universe
$$t=t_0+\frac{I(\mu )}{(1+w)\left(2\pi \overline{A}\rho _0\frac{3}{2}C^{1/2}I(\mu )\right)},$$
(20)
where $`I(\mu )`$ is defined in Eq. (2.11). This is a parametric equation from which one can obtain how the energy per unit length of a wiggled cosmic string evolves in the accelerating universe. Thus, if $`w>1`$ we see that the string energy in the core will progressively increases from its initial value $`\mu _i`$, tending to the maximum value
$$\mu _{\mathrm{max}}=\frac{1+\sqrt{164G^2\mu _0^2}}{8G}.$$
The larger $`w`$ the shorter the time required by the accretion process to make the string to reach $`\mu _{\mathrm{max}}`$. If $`w<1`$, i.e. if we are in the phantom regime, then the linear energy density in the string core will rapidly decreases from its initial value down to recover its unperturbed value at $`\mu _0`$. The smaller $`w`$ the shorter the time taken by the system to reach the value $`\mu _0`$. As the string is approaching that value the gravitational potential should be getting on smaller and smaller values to finally vanish at $`\mu _0`$, so that the classical accretion process will stop at that point. Such a behavior is also checked to occur in the case that phantom K-energy is accreted.
An interesting question is however posed in the two considered kinds of phantom energy. Even though the classical, continuous accretion process must only proceed down to $`\mu _0`$, if we assumed that phantom energy accretion would proceed by discrete steps, then the limit at $`\mu _0`$ should be overtaken and the linear energy density of the string core would continue decreasing below $`\mu _0`$ as the phantom energy was being accreted. We would reach in this way a regime where $`T>\mu `$ which would end when $`\mu `$ reached the minimum value
$$\mu _{\mathrm{min}}=\frac{1\sqrt{164G^2\mu _0^2}}{8G},$$
which would never vanish provided $`\mu _0>0`$. The spacetime metric of the cosmic string given by Eq. (1.2) would then exchange the values between the $`tt`$ and $`zz`$ components, as in this case $`\mu <T`$.
### II.2 Generalized Chaplygin gas
We shall derive now the expression for the rate $`\dot{\mu }`$ in the case of a generalized Chaplygin gas. We start with the expression for the energy density
$$\rho =\left(A_{ch}+\frac{B}{a^{3(1+\alpha )}}\right)^{1/(1+\alpha )},$$
(21)
which has been obtained by integrating the cosmic conservation law for the case of the equation of state of a generalized Chaplygin gas, that is $`p=A_{ch}/\rho ^\alpha `$. Now, from the Friedmann equation we can get
$$\dot{a}=\sqrt{\frac{8\pi G}{3}}a(t)\left(A_{ch}+\frac{B}{a^{3(1+\alpha )}}\right)^{1/[2(1+\alpha )]}.$$
(22)
Hence, from Eq. (2.11) it can be obtained
$`I(\mu )=B\overline{A}\sqrt{{\displaystyle \frac{3\pi }{2G}}}{\displaystyle _{a_0}^a}{\displaystyle \frac{\frac{1}{a^{3(1+\alpha )}}}{a\left(A_{ch}+\frac{B}{a^{3(1+\alpha )}}\right)^{(2\alpha +1)/[2(1+\alpha )]}}}𝑑a`$
$`=\overline{A}\sqrt{{\displaystyle \frac{2\pi }{3G}}}\left[\left(A_{ch}+{\displaystyle \frac{B}{a^{3(1+\alpha )}}}\right)^{1/[2(1+\alpha )]}\sqrt{\rho _0}\right].`$ (23)
It follows
$$a^{3(1+\alpha )}=\frac{B}{\left(\sqrt{\rho _0}\sqrt{\frac{3G}{2\pi \overline{A}^2}}I(\mu )\right)^{2(1+\alpha )}A_{ch}}.$$
(24)
Again in this case the setting of a constant $`B>0`$ implies a progressive increase of $`\mu `$ with $`a`$ up to a maximum given by $`\mu _{\mathrm{max}}`$, and the assumption of a constant $`B<0`$ (phantom) leads to a decrease of $`\mu `$ with $`a`$ down to $`\mu _0`$ classically or to $`\mu _{\mathrm{min}}`$ if the Chaplygin phantom energy is supposed to be accreted in discrete steps.
## III Conclusions and further comments
While cosmic strings have a long tradition and incidence in theoretical cosmology, the introduction of cosmic dark energy has taken place quite more recently though not with less incidence or surprise. Perhaps therefore their potential mutual relations and interactions have not been so far considered. This paper is a first step in the task of studying the effects that the presence of dark energy may have in the fate of cosmic string in an accelerating universe. We have restricted ourselves here to just looking at an approximate model describing how straight wiggly cosmic strings accrete dark energy during the accelerating expansion of the universe, leaving for future publications the accurate treatment for both wiggly straight strings and the similar accretion onto circular strings, so as the kinematic effects that the acceleration of the universe may have on the shape and size of any cosmic strings. A generalized description has first been thus built up and then adapted to the case of the cylindrically symmetric accretion of dark energy onto straight cosmic strings. That description is based on the integration of the conservation laws for the energy-momentum tensor and its projection on four-velocity using the exterior geometry of a wiggly cosmic string. We have considered the dark energy accretion onto straight cosmic strings using several scalar field models for the cosmological vacuum, namely quintessence and K-essence field models with equation of state $`p=w\rho `$, and a generalized model of Chaplygin gas with the unusual equation of state $`p=A/\rho ^\alpha `$. An rate equation for the energy density per unit length of the strings has been in this way derived and finally integrated for each of these dark energy models. This ultimately leads to the prediction that, whereas when the energy density of the cosmic vacuum decreases with time the linear energy density of the straight strings progressively increases as the universe grows bigger for all dark energy models, if the energy density of the universe grows with expansion, inducing a universal violation of the dominant energy condition, the stringy energy density steadily decreases. That energy density dropping makes the strings to eventually become free of wiggles to get thereafter on a quantum accreating regime where the string energy density reaches finally a minimum nonzero value, before the occurrence of any future big rip singularities.
It appears that the current value of the parameter $`w`$ in the equation of state of the universe may be less than -1. So, one could be tempted to think that the above evolution of cosmic strings leading eventually to the formation of exotic topological defects with negative-wiggles perturbations would be inescapable. However, having now $`w<1`$ (provided this turns out to be definitively the case most favoured by observations) does not guarantee at all that the phantom regime will endure in the future. In fact, most general descriptions of quintessence field are based on tracking models where the parameter $`w`$ is time dependent and, therefore, it could well be that what is now less than -1 would later turn out to be greater than -1, so making the cosmic string evolution predicted by our constant-$`w`$ models inapplicable in the far future. Nevertheless, the initial string evolution implied by our phantom models looks as being probable. That behaviour by itself would still be important enough for a variety of subjects. But even such a behaviour would not be guaranteed as phantom fields are characterized by Lagrangians containing negative kinetic terms which have very weird properties and lead to unwanted instabilities making the whole phantom scenario problematic.
###### Acknowledgements.
We thank Professors J.A.S. Lima and E. Babichev for useful explanations, discussions and correspondence. We also acknowledge A. Ferrera and M. Rodríguez for constructive discussions and criticisms. This work was supported by DGICYT under Research Projects BMF2002-03758 and BFM2002-00778. JAJM wants to acknowledge IMAFF for kind hospitality.
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# 𝑆"-wave" 𝐾⁻𝜋⁺ SYSTEM IN 𝐷⁺→𝐾⁻𝜋⁺𝜋⁺ DECAYS FROM FERMILAB E791
## 1 Introduction
Decays of heavy-quark mesons are regarded as a potential source of information on the light-quark mesons they produce. For decays to three pseudo-scalar final states, kinematics and angular momentum conservation favor production of $`S\text{-wave}`$ systems, so improving our knowledge of the particularly confusing scalar meson ($`J^P=0^+`$) spectrum may be possible when the anticipated large, clean samples of such decays of $`D`$ mesons from the $`B`$ factories and the Tevatron collider become available. Extracting this information has, however, been done in model-dependent ways that make assumptions about the scalar states observed. Such assumptions can influence the results, so new approaches are required.
In this paper, we present a model-independent approach to the $`S\text{-wave}`$ system in a study of the decays $`D^+K^{}\pi ^+\pi ^+`$ <sup>a</sup><sup>a</sup>aThroughout the paper, charge conjugate states are implied unless explicitly stated otherwise. observed in data from Fermilab experiment E791. We compare the $`S\text{-wave}`$ amplitudes so obtained with our earlier, isobar model analysis $`^\mathrm{?}`$ and also with $`K^{}\pi ^+`$ scattering.
Measurements of $`K^{}\pi ^+`$ scattering come principally from SLAC experiment E135 (LASS) $`^\mathrm{?}`$, and cover the invariant mass range only above 825 $`\text{MeV}/c^2`$. Data exist below this range, but with less precision $`^\mathrm{?}`$. More information in the low mass region is required if the possibility of the existence of a $`\kappa `$ state is to be properly evaluated.
## 2 Data Sample
The selection process for events used in this analysis is described in Ref. $`^\mathrm{?}`$. A signal consisting of 15,079 $`D^+K^{}\pi _a^+\pi _b^+`$ decays, with a purity of $`94`$%, is obtained. Fig. 1 shows the Dalitz plot with $`K^{}\pi ^+`$ squared invariant mass $`s`$ plotted vs. $`s^{}`$. Horizontal (and the symmetrized vertical) bands corresponding to the $`K^{}(892)`$ resonance are clearly seen. A striking and complex pattern of both constructive and destructive interference is seen near 2 ($`\text{GeV}/c^2`$)<sup>2</sup> due to either $`K_0^{}(1430)`$, $`K_1^{}(1410)`$ or $`K_2^{}(1430)`$. There is also evidence for $`K_1^{}(1680)`$, difficult to see due to smearing of the Dalitz plot boundary resulting from the finite resolution in the three-body $`D^+`$ mass.
The most striking effect observed is the asymmetry in the $`K^{}(892)`$ bands, most easily described by interference with a significant $`S\text{-wave}`$ contribution to the decay. In this paper we are able to extract information on the $`S\text{-wave}`$ using the $`K^{}(892)`$, and also the other well established resonances in the Dalitz plot, as an interferometer.
## 3 Method
In Ref. $`^\mathrm{?}`$, as in most earlier analyses of $`D`$ decays to three pseudo-scalar particles $`ijk`$, we use the “isobar model”. Details of this are given in $`^\mathrm{?}`$. In this model, the decay amplitude $`𝒜`$ is described by a sum of quasi two-body terms $`DR+k,Ri+j`$, in each of the three channels $`k=1,2,3`$:
$`𝒜`$ $`=`$ $`d_0e^{i\delta _0}+{\displaystyle \underset{n=1}{\overset{N}{}}}d_ne^{i\delta _n}{\displaystyle \frac{F_R(p,r_R,J)}{m_{R_n}^2s_{ij}im_{R_n}\mathrm{\Gamma }_{R_n}(s_{ij})}}\times F_D(q,r_D,J)M_J(p,q)`$ (1)
In this, $`s_{ij}`$ is the squared invariant mass of the $`ij`$ system. $`J`$ is the spin, $`m_{R_n}`$ the mass and $`\mathrm{\Gamma }_{R_n}(s_{ij})`$ the width of each of the $`N`$ resonances $`R_n`$ seen to be contributing to the decay. $`F_R`$ and $`F_D`$ are form factors, with effective radius parameters $`r_R`$ and $`r_D`$, for all $`R_n`$ and for the parent $`D`$ meson, respectively. $`p`$ and $`q`$ are momenta of $`i`$ and $`k`$, respectively, in the $`ij`$ rest frame. $`M_J(p,q)`$ is a factor introduced to describe spin conservation in the decay. The complex coefficients $`d_ne^{i\delta _n}(n=0,N)`$ are determined by the $`D`$ decay dynamics and are parameters estimated by a fit to the data. The first, non-resonant ($`NR`$) term describes direct decay to $`i+j+k`$ with no intermediate resonance, and $`d_0`$ and $`\delta _0`$ are assumed to be independent of $`s_{ij}`$. For $`D^+K^{}\pi _a^+\pi _b^+`$ decays we Bose-symmetrize $`𝒜`$ with respect to interchange of $`\pi _a^+`$ and $`\pi _b^+`$.
In Ref. $`^\mathrm{?}`$ we reported that the $`NR`$ term was smaller than previously thought, and that a further term, parametrized as a new $`J=0`$ resonance $`\kappa (800)`$ with $`m_R=(797\pm 19\pm 43)`$ $`\text{MeV}/c^2`$and $`\mathrm{\Gamma }_R=(410\pm 43\pm 87)`$ $`\text{MeV}/c^2`$, gave a much better description of the data. Here, we examine the $`K^{}\pi ^+`$ $`S\text{-wave}`$ in a model-independent way. The $`S\text{-wave}`$ part of Eq. 1 (all terms with $`J=0`$, including the $`NR`$ term) is factored
$`𝒮`$ $`=`$ $`\text{S}(s_{K\pi })\times M_0^R(p,q)F_D(q,r_D)=\text{Interp}\left(c_ke^{i\gamma _k}\right)\times M_0^R(p,q)F_D(q,r_D)`$ (2)
into a partial wave S$`(s_{K\pi })`$, describing $`K^{}\pi ^+`$ scattering, and the product $`M_0^R(p,q)F_D(q,r_D)`$ describing the $`D`$ decay. S$`(s_{K\pi })`$ is interpolated between a set of points $`c_ke^{i\gamma _k}`$ defined at 40 $`K^{}\pi ^+`$ invariant mass squared values $`s_{K\pi }^k`$ indicated by the lines in Fig. 1. Each $`c_k`$ and $`\gamma _k`$ is regarded as an independent parameter determined by the data.
We factor the $`P\text{-}`$ and $`D\text{-}`$ reference waves in the same way:
$`𝒫`$ $`=`$ $`\text{P}(s_{K\pi })\times M_1^R(p,q)F_D(q,r_D);𝒟=\text{D}(s_{K\pi })\times M_2^R(p,q)F_D(q,r_D),`$ (3)
however, we parametrize the partial waves P$`(s_{K\pi })`$ and D$`(s_{K\pi })`$ exactly as in Eq. 1.
We make an unbinned likelihood fit to the data, Using the method described in Ref. $`^\mathrm{?}`$. This incorporates an incoherent background function describing the 6% of our sample not corresponding to true $`D`$ decays. We measure 86 parameters - all $`(c_k,\gamma _k)`$ and the coefficients $`d_ke^{i\delta _k}`$ for $`K^{}(892)`$, $`K_1^{}(1680)`$ in the $`P\text{-wave}`$ and $`K_2^{}(1430)`$ in the $`D\text{-wave}`$. For the $`K^{}(892)`$, we define $`d_ke^{i\delta _k}=1`$ to provide the reference phase.
The fit results in an excellent description of the data. Comparison of the observed and predicted population of the Dalitz plot gives a $`\chi ^2`$ probability of 50% for 363 bins.
## 4 Results
The $`S\text{-}`$, $`P\text{-}`$ and $`D\text{-wave}`$s resulting from the fit are shown in Fig. 2. They are compared with the model-dependent fit from Ref. $`^\mathrm{?}`$. The main $`S\text{-wave}`$ features of both fits agree well. Resonant fractions and the total $`S\text{-wave}`$ fraction (about 75%) also agree within statistical limits.
We turn now to a comparison of the $`S\text{-wave}`$ amplitudes S$`(s_{K\pi })`$ measured here with the amplitudes $`T(s_{K\pi })`$ measured in $`K^{}\pi ^+`$ elastic scattering. We expect, for each partial wave $`J`$ (for each iso-spin $`I`$) that $`\text{S}(s_{K\pi })=\sqrt{s_{K\pi }}/p^{(J+1)}Q(s_{K\pi })\text{T}(s_{K\pi })`$ where $`Q(s_{K\pi })`$ describes the dependence of $`K^{}\pi ^+`$ production in $`D`$ decays on $`s_{K\pi }`$. The Watson theorem $`^\mathrm{?}`$ requires that, provided there is no re-scattering of the $`K^{}\pi _a^+`$ from $`\pi _b^+`$, that $`Q`$ is a real function, so that phases found in $`D`$ decay should match those in $`K^{}\pi ^+`$ elastic scattering data.
$`I=1/2`$ phases measured by LASS are plotted in Fig. 2. There is a large offset in the $`S\text{-wave}`$, about 75, not seen in $`P\text{-}`$ or $`D\text{-wave}`$s. The shapes of $`S\text{-}`$ and $`P\text{-wave}`$s are also not the same. Unless significant admixture of $`I=3/2`$ $`K^{}\pi ^+`$ production occurs, these results suggest that the conditions for the Watson theorem are not met in these data.
## Acknowledgments
We thank members of the LASS collaboration for making their data available to us. We gratefully acknowledge the assistance of the staffs of Fermilab and of all the participating institutions. This research was supported by the Brazilian Conselho Nacional de Desenvolvimento Científico e Tecnológico, CONACyT (Mexico), FAPEMIG (Brazil), the Israeli Academy of Sciences and Humanities, the U.S. Department of Energy, the U.S.-Israel Binational Science Foundation, and the U.S. National Science Foundation. Fermilab is operated by the Universities Research Association for the U.S. Department of Energy.
## References
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# Rest-frame optical and far-infrared observations of extremely bright Lyman-break galaxy candidates at z ∼ 2.5
## 1 Introduction
The very large areal coverage of the Sloan Digital Sky Survey (SDSS, York et al., 2000) provides a unique opportunity to identify rare, intrinsically luminous examples of high-redshift galaxy populations, as well as similarly rare, strongly-gravitationally magnified examples of more typical high-redshift galaxies. To exploit this opportunity, Bentz & Osmer (2004) searched the SDSS Early Data Release (Stoughton et al., 2002) for unusual quasars with anomalously low C iv 154.9nm emission, and found a luminous $`z2.5`$ starburst candidate that appeared to have been incorrectly classified as a quasar by the SDSS pipeline.
Following on from this find, Bentz, Osmer & Weinberg (2004) identified a further five sources from the SDSS First Data Release (DR1) Quasar Catalog (Schneider et al., 2003) at $`z`$$``$ 2.5–2.8 with rest-frame ultraviolet (UV) colours similar to Lyman-break galaxies (LBGs, Steidel et al., 2003) and exhibiting weak or absent high-ionisation emission lines in their rest-frame UV spectra. All six objects have $`r`$-band magnitudes of 19.8–20.5, with a median of $`r20.3`$, i.e. they are an order of magnitude brighter than the most luminous objects in existing LBG surveys. Bentz, Osmer & Weinberg (2004) showed that if their UV emission arises solely from star formation then they are intensely active star-forming galaxies with strong lower limits on their star-formation rates ranging from 300 to 1100 M yr<sup>-1</sup>, assuming negligible absorption by dust and adopting a continuous star-formation rate over $`10^8`$ yrs (Kennicutt, 1998). Thus these candidate LBGs could be rare, extreme starbursts seen at the epoch, $`z2.5`$, when both the accretion luminosity density and the star-formation rate density in the Universe are believed to peak (Miyaji, Hasinger & Schmidt, 2000; Chapman et al., 2005).
Alternatively, the brightness and apparent rarity of these systems could simply reflect the fact that they are rare, strongly-magnified examples of the normal-luminosity LBG population. The LBG candidates show no evidence of multiple components at the resolution of SDSS imaging ($`>`$1 arcsec) and there are no obvious foreground lensing structures. Nevertheless, gravitational lensing cannot yet be ruled out based on existing imaging.
The best argument against these galaxies being highly-magnified examples of normal LBGs results from their spectral properties: the underlying continuua are much redder than typical LBGs, with observed 200–600-nm spectral indices ranging from $`\alpha =1.9`$ to $`2.5`$ (where $`F_\nu \nu ^{+\alpha }`$), cf. $`\alpha >1.7`$ for most LBGs (Shapley et al., 2003), and the low-resolution SDSS spectra show strong (sometimes broad) interstellar absorption lines and Ly$`\alpha `$ emission in several cases, as well as hints of broad C iii\] 190.9nm emission in two cases (Bentz, Osmer & Weinberg, 2004). These features contrast with the narrow emission and interstellar absorption lines seen in the composite LBG spectrum of Shapley et al. (2003), but may be explained due to star-formation activity (and perhaps associated winds) an order of magnitude more vigorous than the sample considered by Shapley et al. (2003). The gross spectral properties of a source should not be affected markedly by lensing, so these differences argue that these galaxies are unlikely to be highly-magnified examples of the general LBG population. Rather, the spectral properties of these sources share key characteristics with the submm-selected galaxies identified in the recent spectroscopic survey of Chapman et al. (2003, 2005). In particular, they resemble N2 850.4, a composite starburst/AGN at $`z`$ = 2.38 (Smail et al., 2003) and the BAL–Sy2/QSO SMM J02399$``$0136 at $`z`$ = 2.80 (Ivison et al., 1998; Vernet & Cimatti, 2001) — the two-component Ly$`\alpha `$ emission, broad C iii\] and the absorption seen in C ii, Si iv, Al iii and C iv — although the candidate LBGs lack such prominent P-Cygni profiles. The LBG candidates also share some common characteristics with broad absorption-line (BAL) quasars — the presence of broad C iii\]190.9nm in two examples and some very broad absorption lines — although the line profiles are not typical of BALs. There is thus a possibility that the UV emission from these galaxies is powered by accretion rather than star formation, or that the sample is a heterogeneous mix of star-forming galaxies and active galactic nuclei (AGN). A recent paper by Hall et al. (2004) discusses an unusual object, SDSS J113658.36+024220.1, whose rest-frame UV spectrum shows a single emission line corresponding to Ly$`\alpha `$ but no obvious metal-line emission, which bears some similarities to the candidate LBGs studied here. Hall et al. (2004) interpret SDSS J1136 as an AGN based in large part on tentative optical variability and its strong radio emission, $``$1.4 mJy at 1.4 GHz.
If these galaxies are truly related to starburst LBGs then their prodigious star formation should be betrayed in the rest-frame far-infrared. Here, we exploit Submm Common-User Bolometer Array — SCUBA, Holland et al. (1999) — submillimetre (submm) photometry to search for such emission. We then use new, high-resolution, near-infrared imaging to identify the morphological signatures of strong lensing. Finally, we present near-infrared spectra of these galaxies, covering a number of key rest-frame optical emission lines falling in the $`H`$ and $`K`$ atmospheric windows, to spectroscopically classify the galaxies and to search for quasar signatures such as a tell-tale broad component to the H$`\alpha `$ line. We describe our observations in §2, present our results and discussion in §3 and give our conclusions in §4.
## 2 Submillimetre and near-infrared observations
Submm photometry observations were obtained for four of our LBG candidates in service time on 2004 January 15 and 28 with SCUBA on the James Clerk Maxwell Telescope (JCMT<sup>6</sup><sup>6</sup>6The JCMT is operated by the Joint Astronomy Centre in Hilo, Hawaii, on behalf of the Particle Physics and Astronomy Research Council (PPARC) in the UK, the National Research Council of Canada, and The Netherlands Organisation for Scientific Research.). On the first night, observations of SDSS J1147, SDSS J1340 and SDSS J1444 were made in average opacity conditions ($`\tau _{850\mu \mathrm{m}}`$$``$ 0.3–0.4), and on the second night observations of SDSS J0243 were made in better conditions ($`\tau _{850\mu \mathrm{m}}`$$``$ 0.1–0.2). We used SCUBA in two-bolometer mode, giving a $``$15 per cent improvement in signal-to-noise compared with one-bolometer mode. Each source was observed for 1.8 ks. The starlink package surf was used to reduce the data for each bolometer separately. The resulting signals were then calibrated against the JCMT secondary calibrators CRL 618 and 16293$``$2422, and co-added to give weighted 850-$`\mu `$m flux densities and errors (see Table 1). Calibration uncertainties are estimated to be $``$10 per cent.
Near-infrared imaging data in the $`J`$\- and $`K`$-bands were obtained for the five bright LBG candidates accessible to the 3.8-m UK Infrared Telescope (UKIRT<sup>7</sup><sup>7</sup>7UKIRT is operated by the Joint Astronomy Centre on behalf of PPARC.) during 2004 January–April and 2004 July–August. Flexible scheduling enabled us to utilise better-than-average seeing on Mauna Kea, 0.4–0.6 arcsec, and we employed the UKIRT Fast Track Imager, UFTI (Roche et al., 2003), a 1024<sup>2</sup> HgCdTe array with 0.091-arcsec pixels, to exploit those conditions. The total integration time in each filter, built up whilst dithering every 60 s, was 3.8 ks. Contiguous observations of nearby faint standards were used to determine zero points. The 3-$`\sigma `$ detection threshold is $`K`$$``$ 21.5 in a 4-arcsec-diameter aperture. Data were reduced using orac-dr and we report effective total magnitudes/colours (measured from 6-arcsec-diameter photometry) for the LBG candidates in Table 1. Objects in the frames were then identified and catalogued using SExtractor (Bertin & Arnouts, 1996) and colours measured in 2-arcsec-diameter apertures from the aligned $`J`$\- and $`K`$-band frames. $`K`$-band images of regions around each target are shown in Fig. 1, with the $`(JK)`$$`K`$ colour-magnitude distributions for each field displayed in Fig. 2.
Spectra were obtained during 2004 April and August with the UKIRT 1–5 $`\mu `$m Imager Spectrometer, UIST (Ramsay Howat et al., 1998), which utilises a 1024<sup>2</sup> InSb array with 0.12-arcsec pixels. UIST’s $`HK`$ grism was used to cover the 1.4–2.5 $`\mu `$m region, with measured resolutions of $`\lambda /\mathrm{\Delta }\lambda `$ = 390 and 680 (for arc lines at 1.5 and 2.3 $`\mu `$m) for our 4-pixel-wide slit (0.48 $`\times `$ 120 arcsec). Acquisition was accomplished using 20–60-s sky-subtracted images of each field. We are confident that all targets were placed within a pixel of the optimal position on the slit. Each target was observed for 6.7 ks, nodding along the slit in an A-B-B-A sequence every 240 s. An Argon arc spectrum and a flatfield frame were obtained prior to observations of each target. Nearby F5V standard stars were observed contiguously to set the flux scale, having interpolated across their Hydrogen absorption lines before ratioing. The frames were reduced using orac-dr, with optimal extraction of spectra accomplished using figaro. For comparison purposes we also obtained UIST $`HK`$-spectra of the weak-lined quasar SDSS J113658.36+024220.1 ($`z`$ = 2.4917, Hall et al., 2004) and the $`z`$ = 2.320 LoBAL quasar, SDSS J135317.80$``$000501.3 (Reichard et al., 2003; Willott et al., 2003). All the spectra are shown in Fig. 3.
A further spectrum was obtained for the LBG candidate inaccessible to UKIRT (SDSS J134026.44+634433.2) using NIRSPEC on Keck-ii<sup>8</sup><sup>8</sup>8The W. M. Keck Observatory is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. during relatively poor conditions (0.9-arcsec seeing) on 2005 March 18. The full spectral range was 2.27–2.69 $`\mu `$m with a resolution of $``$1,500, utilising a 42 $`\times `$ 0.76-arcsec (4-pixel) slit. The total integration time was 1.6 ks, split into four A-B-B-A sequences with 100 s per exposure. A flux standard was not observed; otherwise, the data reduction followed that employed for UIST.
## 3 Results and discussion
We now discuss the insights provided into the nature of the candidate LBGs by our various multi-wavelength observations. The basic properties of our sample — redshifts, near-infrared photometry and 850-$`\mu `$m flux densities — are listed in Table 1. $`K`$-band imaging of the target fields are shown in Fig. 1, $`(JK)`$$`K`$ colour magnitude diagrams of these fields in Fig. 2 and the $`HK`$ spectra of the candidate LBGs in Fig. 3. Line widths and fluxes for the stronger features in these spectra are listed in Table 2.
### 3.1 Submillimetre properties
Studies of submm-selected galaxies (SMGs) have concluded that most are too faint in the UV to be identified by the photometric selection used for $`z3`$ LBG surveys (Webb et al., 2003). This suggests little overlap between these two well-studied classes of high-redshift star-forming galaxies. However, statistical measurements of the submm emission from samples of $`z3`$ LBGs, reaching below SCUBA’s confusion limit, suggest that they may contribute substantially to the SMG population at sub-mJy levels (Peacock et al., 2000; Webb et al., 2003; Kneib et al., 2004). Only a handful of brighter examples are known (Chapman et al., 2002). There is evidence of a more significant overlap between SMGs and the UV-selected population identified at somewhat lower redshifts, $`z`$$``$ 1.5–2.5 (Steidel et al., 2004), although again many of the SMGs are too faint to be included in the photometric samples (Chapman et al., 2005).
The raw star-formation rates estimated from the observed UV luminosities of the SDSS sources, uncorrected for dust extinction, would imply $`L_{\mathrm{FIR}}>10^{13}`$ L for our sample, where SFR = $`ϵ`$ 10$`{}_{}{}^{10}L_{\mathrm{FIR}}^{}`$ M yr<sup>-1</sup> and $`ϵ`$ = 0.8–2.1 (Scoville & Young, 1983; Thronson & Telesco, 1986), and thus 850-$`\mu `$m flux densities of 3–10 mJy (Blain & Longair, 1996). Assuming a correction factor for dust extinction typical of LBGs, 6$`\times `$ (Pettini et al., 2002; Erb et al., 2003), the predicted submm fluxes would increase by a similar factor. This suggests that the candidate LBGs should be detectable in the submm waveband if they have submm/UV flux ratios similar to the more typical luminosity members of this population. This conclusion holds whether these galaxies are either intrinsically bright in the UV or are lensed (assuming that the lensing does not preferentially boost the UV-bright regions).
The submm data presented in Table 1 demonstrate that the four SDSS sources we have observed are all undetected individually at flux limits of 6–8 mJy. This implies that it is unlikely that these galaxies are simply scaled-up or strongly-lensed examples of typical-luminosity $`z3`$ LBGs. Indeed, the weighted mean for the four sources ($``$0.36 $`\pm `$ 1.31 mJy) suggests that they would not have been detected in even the deepest submm survey, although we cannot rule out the possibility that the sample is heterogeneous, with a handful of faint submm emitters.
For completeness we note that a search of the FIRST radio survey (Becker, White & Helfand, 1995) yielded only upper limits at 1.4 GHz (Table 1), and that none of the galaxies were detected by ROSAT to limits appropriate for the X-Ray All-Sky Survey.
### 3.2 Photometric and morphological properties
Our near-infrared observations indicate that the median observed colours for the LBG candidates are $`JK=1.71\pm 0.20`$, $`iK=3.86\pm 0.43`$ and $`rK=3.43\pm 0.37`$. Comparing the $`JK`$ colours with the spectra of the sources, from §3.3, it is clear that the LBG candidates with the strongest line emission in the $`K`$-band also have the reddest continuum colours, suggesting that the line emission is biasing the colours we measure. With typical observed-frame equivalent widths of $``$100 to $``$150 nm, and a $`K`$-band filter width of $``$350 nm, the fluxes in $`K`$ should be corrected by approximately 0.6–0.7$`\times `$ (or +0.5 in magnitudes). No sources are classed as extremely red objects on the basis of either their $`rK`$ or $`iK`$ colours, although the reddest source in the optical/near-infrared, SDSS J1147, is also the reddest in the rest-frame UV and is the only candidate LBG in our sample which does not show Ly$`\alpha `$ emission. The median $`JK`$ colour of our sample is comparable to that seen for $`z3`$ LBGs from Shapley et al. (2001), $`JK1.63`$, although our candidate LBGs are redder on average in $`rK`$ than the standard UV-selected populations at $`z2`$ or $`z3`$, $`rK3.25`$ and $`rK2.85`$, respectively (Shapley et al., 2001; Steidel et al., 2004). This suggests that the rest-frame UV continua may be significantly redder than normal LBGs, although their rest-frame optical continua are comparable (before correcting for emission-line contributions).
Turning to the high-resolution near-infrared imaging (Fig. 1), we find that all five candidate LBGs are unresolved at the 0.45–0.55-arcsec seeing of our $`K`$-band images (as measured from multiple stars in each frame). Given the signal to noise of our detections of the LBGs and their measured FWHM relative to stars in the fields, we can place firm limits of $`<0.1`$ arcsec fwhm on the sizes of these sources or on the separation of multiple components if they are strongly lensed. In the absence of lensing, this angular limit corresponds to an intrinsic size of $`<`$1 kpc for the physical scale of these sources in their rest-frame $`V`$-band light.
We identify none of the morphological signatures expected from strong galaxy–galaxy lensing in our deep $`K`$-band images: multiple lensed components or an identifiable foreground lens. We note that the angular size limit estimated above, if taken as the Einstein diameter, would correspond to a velocity dispersion of only 50 km s<sup>-1</sup> (for a spherical isothermal lens at $`z0.5`$). This velocity dispersion would correspond to a $``$0.1 L early-type galaxy with an expected magnitude of $`K19`$ at $`z0.5`$, detectable in our imaging out to $`z1`$ (Rusin et al., 2003). No such nearby lenses are visible in our imaging on the relevant scales (0.2–2 arcsec).
Looking at a wider region around the candidate LBGs we see a compact $`K`$ = 17.9 galaxy (with $`JK0.8`$) lying only 4 arcsec away from SDSS J1553, but this would not provide a strong amplification of the source. We also find that SDSS J1432 sits in the outskirts of a dense, compact foreground group, the brightest members of which have $`K`$$``$ 16. It is possible that the LBG candidate is thus magnified by weak lensing from the foreground structure, although the total magnification is likely to be modest. The fields surrounding SDSS J0243 and SDSS J1444 are unremarkable, but we do identify a group of 5–6 faint, resolved galaxies, $`K`$$`>`$ 19.4, within 15 arcsec of SDSS J1147. These are possibly members of a cluster, either in the foreground or (given their faintness) associated with the candidate LBG. These faint galaxies exhibit a wide range in $`JK`$ colours, 0.8–3.0, including some extremely red objects. In addition, several brighter galaxies, $`K`$$`>`$ 17.3, lie within 20 arcsec. We cannot demonstrate a significant over-density, but they could be more luminous members of the same structure.
Looking at the $`(JK)`$$`K`$ colour-magnitude plots (Fig. 2) for the five fields we see that the LBG candidates are rarely the reddest galaxy in the field, although both SDSS J1147 and SDSS J1432 have very red, close neighbours ($`JK=2.2`$–3.0). The most striking feature is the colour-magnitude sequence seen in the distribution of galaxy colours in the SDSS J1432 field. This has a characteristic colour of $`JK1.8`$ at $`K17.5`$, consistent with that expected for evolved galaxies in a group or poor cluster at $`z0.7`$–0.8 (Feulner et al., 2003). The close similarity of the colours of the candidate LBG to these galaxies may either be a coincidence or could point to contamination of our photometric measurement by a superimposed member of this group (which could thus also gravitationally magnify the background source). We believe that the similarities in the colours are merely a coincidence as it is clear from the spectrum of SDSS J1432 (Fig. 3) that there is a significant contribution to the $`K`$-band light from the background source.
We conclude that the morphological information for all five of the candidate LBGs we have imaged at $`<`$0.55-arcsec resolution provides no support for them being strongly lensed. There is possible evidence for weak lensing for one or more sources, but this would not significantly affect the apparent magnitudes of these systems.
### 3.3 Spectral properties
It is immediately apparent from the rest-frame optical spectra presented in Fig. 3 that all six of the spectroscopically-observed LBG candidates contain broad-line AGN. All have broad H$`\alpha `$ emission, with fwhm line widths ranging from 5,000 to $``$14,000 km s<sup>-1</sup> (Table 2) and broad components visible in H$`\beta `$ for several systems (despite the typically poorer signal-to-noise in the $`H`$-band). SDSS J1444 has the weakest emission and there is no spectral coverage beyond the red extreme of the line due to its high redshift, but a broad H$`\alpha `$ line is still apparent. Comparing them to the two known AGN we observed, the weak-lined AGN SDSS J1136 and the low-ionisation BALQSO SDSS J1353, we see that the candidate LBGs exhibit broader H$`\alpha `$ emission than either of these two AGN.
So, clear AGN features are visible in the rest-frame optical, whereas the UV spectra of these galaxies are characterised by a strong continuum but lack the strong emission lines typical of AGN. Does the AGN contribute significantly to the UV fluxes of these galaxies?
The redshift measurements available to us tend to follow the same pattern for all the LBG candidates: the UV-determined values, based predominantly on the prominent Ly$`\alpha `$ emission line, are slightly blueward of the \[O iii\] 500.7nm and H$`\alpha `$ lines ($`\mathrm{\Delta }z=`$ 0.0051 $`\pm `$ 0.0057 and 0.0065$`\pm `$0.0037, respectively). The only exception is SDSS J1444, where the Ly$`\alpha `$ and H$`\alpha `$ redshifts are identical, but the H$`\alpha `$ redshift is poorly determined. The H$`\alpha `$ emission we see originates close to the AGN, in the broad-line region (BLR), so it is natural to assume that the UV absorption lines are due to wind-driven material in our line of sight to the BLR, and that the UV continuum also arises close to the AGN. Indeed, there is a tight correlation between the FWHM of the H$`\alpha `$ (Table 2) and the absolute $`i`$-band magnitudes from Bentz, Osmer & Weinberg (2004): the 0.07-dex scatter suggests a close relationship between the UV continuum emission and the AGN. However, it is not clear whether this is a direct relationship, or whether it arises merely because more massive AGN reside in more luminous galaxies. Nevertheless, assuming that the candidate LBGs have intrinsic power-law continua characteristic of normal quasars, with $`\alpha =0.44`$ (Vanden Berk et al., 2001), then their observed rest-frame 200–600-nm spectral slopes ($`\alpha =1.89`$ to $`2.54`$) indicate substantial dust extinction, $`A_V1.35`$–1.95, for a Calzetti extinction law (Calzetti et al., 2000).
But where are the UV emission lines usually associated with quasar activity? If they have been quenched by dust surrounding the active nuclei, why can we still see intense UV continuum emission? The spectral characteristics of these galaxies are unusual, but the presence of strong and broad absorption lines in the UV are similar to those seen in less-reddened examples of the most extreme BAL quasars found by the SDSS (Hall et al., 2002) and in the Digitized Palomar Observatory Sky Survey (Brunner et al., 2003). Indeed, very recently Appenzeller et al. (2005) have published high-resolution échelle spectroscopy of SDSS J1553 which provides much higher-quality information about the UV spectral properties of this galaxy. Based on their analysis of the detailed properties of the absorption lines, they conclude that SDSS J1553 is a low-ionisation BAL quasar (LoBALQSO), or perhaps an even rarer FeLoBALQSO. The very strong low-ionisation absorption features found in LoBALQSOs across a wide velocity range can strongly suppress the emission-line components in these systems, leading to the absorption-dominanted UV spectra we see. The UV absorption features of SDSS J1553 are typical of those seen in the other five galaxies and so we expect that deeper and higher resolution spectroscopy of the complete sample would likely lead to the same conclusion for the other sources. Indeed, based on the existing low-resolution SDSS spectra, Appenzeller et al. (2005) suggest at least two-thirds of the sample may be LoBALQSOs. The relative weakness and narrowness of the BAL features, combined with the absence of strong UV emission features, suggests that the outflows in these galaxies may differ in terms of their velocity and spatial coverage compared to those seen in typical LoBALQSOs. Alternatively, these AGN may be similar to SDSS J1136 (Fig. 3), which Hall et al. (2004) suggest for some unknown reason has weak, broad and highly blue-shifted emission lines.
Finally, we want to highlight the properties of SDSS J1340. This candidate LBG was detected at 16 $`\mu `$m using Spitzer by Teplitz et al. (2004). They interpret this detection in terms of a massive starburst, even though the optical/mid-infrared spectral energy distribution (SED) of SDSS J1340 in Teplitz et al. (2004) is best fit by the SED for the Seyfert-1, NGC 5548. Our non-detection of the source in the submm suggests that the 16 $`\mu `$m detection most likely arises from high-temperature AGN-heated dust, rather than a bolometrically luminous starburst. Further support for the presence of a bolometrically luminous AGN in this system comes from the detection of a strong and broad H$`\alpha `$ line in our near-infrared spectrum (Fig. 3).
## 4 Conclusions
We present multi-wavelength observations of a sample of six candidate LBGs at $`z=2.5`$–2.8 identified from the SDSS DR1 QSO Catalog by Bentz, Osmer & Weinberg (2004). We suggest that these sources could be either: 1) intrinsically luminous, UV-bright starbursts; 2) strongly-lensed examples of typical-luminosity LBGs; or 3) a class of quasars with extremely weak UV emission lines.
We do not detect any of the four candidate LBGs observed in the submm, placing a strong constraint on the submm emission from the ensemble. This suggests that the sources are unlikely to be strongly-lensed examples of more typical LBGs, or intrinsically-luminous LBGs, unless the far-infrared emission from such UV starbursts declines precipitously at high luminosities. Two further pieces of evidence weigh against the lensing hypothesis: first, the UV spectral properties of the candidate LBGs do not match those of typical luminosity LBGs; second, using high-resolution near-infrared imaging of five of the candidates we find no morphological evidence of strong lensing. Taking these results together, we conclude that the sources in our sample are unlikely to be either intrinsically-luminous LBGs or rare, strongly-lensed examples of more normal LBGs. This suggests that they are most likely to be unusual AGN.
Our near-infrared spectroscopy confirms this suggestion, identifying very broad lines in the rest-frame optical spectra of all six galaxies in the sample. We therefore conclude that the six apparently extremely luminous LBGs identified by Bentz, Osmer & Weinberg (2004) are likely to be LoBALQSOs whose unusually weak UV emission lines may either be an intrinsic property of these AGN (Hall et al., 2004) or result from a complex distribution of absorption in the outflow close to the AGN.
## Acknowledgments
We thank Pat Osmer and David Weinberg for their work on the SDSS LBG survey. We also thank Alastair Edge for useful conversations, and the referee for suggestions that improved the paper markedly. We acknowledge service observations from the JCMT. IRS acknowledges support from the Royal Society. MB is supported by a Graduate Fellowship from the National Science Foundation. AWB acknowledges support from NSF grant AST-0205937, the Research Corporation and the Alfred P. Sloan Foundation.
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# 1 Introduction
## 1 Introduction
Einstein, Podolsky and Rosen (EPR) presented their famous gedankenexperiment in 1935 with the aim of showing that quantum mechanics (QM) was not a complete description of physical reality. A complete description, in their view, would require the introduction of additional variables, usually referred to as hidden variables. They outlined a program to reproduce the predictions of QM using local hidden variable (LHV) theories.
This program was challenged by Bell in 1964 when he proved that any hidden variable theory that incorporated the concepts of locality and reality would be inconsistent with certain predictions of QM. In particular, he showed that it was possible to derive from the postulates of EPR an inequality which was violated by statistical predictions of QM for a pair of particles. This violation has been observed in a number of experiments involving interference of pairs of photons produced in an entangled state \[4-7\].
A shortcoming of all these experiments is that, with a pair of particles, Bell’s theorem reveals contradictions between the predictions of QM and EPR’s postulates only in situations involving imperfect statistical correlations: no contradictions appear with perfect correlations. A way to overcome this limitation, which was proposed by Greenberger, Horne and Zeilinger (GHZ) is to use three or more particles in an entangled state. Greenberger et al. (GHSZ) have shown that, in this case, contradictions emerge even at the level of perfect correlations. They also described a gedankenexperiment using three entangled photons to illustrate this point.
Recent work suggests that it is possible to generate entangled four-photon states by parametric down-conversion of two pump photons . Based on this work, we describe a design for a four-photon interference experiment which can be realized in the laboratory. We show how, with just two measurements, one can demonstrate that quantum mechanics contradicts LHV theories, even at the level of perfect correlations.
## 2 Generation of four photon entangled states
Parametric down conversion in a crystal exhibiting a $`\chi ^{(2)}`$ nonlinearity makes it possible to convert a pump photon (frequency $`\omega _p`$) into a pair of highly correlated photons with frequencies $`\omega _d^{(1)}`$ and $`\omega _d^{(2)}`$, where $`\omega _d{}_{}{}^{(1)}+\omega _d{}_{}{}^{(2)}=\omega _p`$ \[11-13\]. These photons are generated almost simultaneously (within the correlation time $`\tau _d`$ of the down-converted photons). While the frequency of each down-converted photon may vary over an appreciable range, the sum of their frequencies is fixed to within the pump bandwidth. The down-converted photons are therefore described by an energy entangled state.
Recent work has shown that it should be possible to extend this process to generate entangled four-photon states from two pump photons by achieving the required phase-matching conditions in a non-linear crystal with two non-collinear pump beams . While the susceptibility for this two-photon down-conversion process is low, the gain depends on the second power of the pump amplitude, so that it should be possible to obtain an appreciable yield by using pulsed pump beams with high peak power. Further improvements in yield may be possible by the use of a resonant cavity (See Appendix D for further details of the experiment).
In such an arrangement,
$$\omega _d{}_{}{}^{(1)}+\omega _d{}_{}{}^{(2)}+\omega _d{}_{}{}^{(3)}+\omega _d{}_{}{}^{(4)}=2\omega _p,$$
(1)
where $`\omega _d^{(1)}`$,…, $`\omega _d^{(4)}`$ are the frequencies of the down-converted photons and $`\omega _p`$ is the frequency of the pump photons, and
$$𝐤_d{}_{}{}^{(1)}+𝐤_d{}_{}{}^{(2)}+𝐤_d{}_{}{}^{(3)}+𝐤_d{}_{}{}^{(4)}=𝐤_p{}_{}{}^{(1)}+𝐤_p{}_{}{}^{(2)},$$
(2)
where $`𝐤_d^{(1)}`$,…, $`𝐤_d^{(4)}`$ are the wave vectors of the down-converted photons, and $`𝐤_p^{(1)}`$ and $`𝐤_p^{(2)}`$ are the wave vectors of the pump photons.
## 3 Four photon interferometer
The four-photon interferometer shown in figure 1 is an extension of a two-photon interferometer described by Franson in which each of the four down-converted photons enters one of four interferometers, each with a short path ($`s_j`$) of length $`S_j`$ and a long path ($`l_j`$) of length $`L_j`$. The optical path difference $`\mathrm{\Delta }L_j=L_jS_j`$ in each of the interferometers, which can be varied by translating the right-angle prisms, is greater than the coherence length $`c\tau _d`$ of the down-converted photons, so that no second-order interference effects due to single photons are observed in the individual interferometers.
If, in a pair of interferometers, we consider the four processes leading to photon counts ($`s_is_j,s_il_j,l_is_j,l_il_j`$), and the difference of the optical path differences ($`\mathrm{\Delta }L_{ij}=\mathrm{\Delta }L_i\mathrm{\Delta }L_j`$) is less than the coherence length of the pump beam, the $`l_il_j`$ and $`s_is_j`$ processes are indistinguishable from each other. However, the other two processes ($`s_il_j`$ and $`l_is_j`$) can be distinguished from the $`l_il_j`$ and $`s_is_j`$ processes by the relative time lag of the photons . It is then possible, with fast coincidence counters, to reject counts arising from the $`s_il_j`$ and $`l_is_j`$ processes, so that we are only concerned with coincidences due to the $`l_1l_2l_3l_4`$ and $`s_1s_2s_3s_4`$ processes in the four interferometers (See Appendix D for further comments on coincidence counts).
## 4 Four photon interference
A four-photon event is recorded when photons are detected in coincidence (within the detector response time) in all the four interferometers. Since the four photons are generated (almost) simultaneously, such a coincidence could either be due to four photons which all took the short path ($`|s>_1|s>_2|s>_3|s>_4`$) or the long path $`(|l>_1|l>_2|l>_3|l>_4)`$. Following Feynman , we can compute the amplitude (see refs. and ) for the arrival of four coincident photons by summing the amplitudes for these indistinguishable alternatives. We have
$$|\mathrm{\Psi }>=|s>_1|s>_2|s>_3|s>_4+\mathrm{exp}(i\mathrm{\Phi })|l>_1|l>_2|l>_3|l>_4,$$
(3)
where the relative phase $`\mathrm{\Phi }`$ of the interfering $`s_1s_2s_3s_4`$ and $`l_1l_2l_3l_4`$ processes is the sum of the relative phases acquired by the individual photons in the four interferometers, so that
$$\mathrm{\Phi }=\varphi _1+\varphi _2+\varphi _3+\varphi _4,$$
(4)
where $`\varphi _i=(\omega _p/2c)\mathrm{\Delta }L_i,(i=1,\mathrm{},4)`$ is the phase difference between the beams traversing the two arms of the $`i`$th interferometer. The predicted coincidence count is obtained by squaring the amplitude and is therefore proportional to
$$|(1/2)(1+\mathrm{exp}(i\mathrm{\Phi }))|^2=(1/2)(1+\mathrm{cos}(\mathrm{\Phi })),$$
(5)
where the constant of proportionality includes the intensity of the source, the detector efficiency (See Appendix C) and the losses in the system. A formal field-theoretic analysis which leads to the same result is presented in Appendix A. If the detectors are as nearly alike as possible, we can assume fair sampling, so that the number of coincidences actually measured in any situation is proportional to those expected for a perfect system.
As can be seen, QM predicts that the coincidence rate $`R_c`$ depends only on $`\mathrm{\Phi }`$, the sum of the phase delays $`\varphi _i`$ in the four interferometers. The coincidences will be perfectly correlated ($`R_c=1`$) when $`\mathrm{\Phi }=0`$ and perfectly anticorrelated ($`R_c=0`$) when $`\mathrm{\Phi }=\pi `$. When $`\mathrm{\Phi }=0`$, detection of a photon in three interferometers would imply the coincident detection of a photon in the fourth interferometer. When $`\mathrm{\Phi }=\pi `$, detection of a photon in three of the interferometers would preclude the coincident detection of a photon in the fourth interferometer. Let us define a parameter (analogous to the visibility for sinusoidal fringes)
$$𝒬=\frac{R_c{}_{}{}^{(0)}R_c^{(\pi )}}{R_c{}_{}{}^{(0)}+R_c^{(\pi )}}$$
(6)
whose value quantum mechanics predicts to be unity. As we will see in the next section, LHV theories cannot explain this value of $`𝒬`$.
## 5 LHV predictions
It is convenient in discussing the four-photon interferometer to use the language of spins traditionally used in the EPR literature. Traversals of the long and short arms of an interferometer are thought of as basis states $`|s>`$ and $`|l>`$ correponding to “spin up” and “spin down” along the $`z`$ axis. The superposition of these states with a phase difference $`\varphi `$
$$|\psi >=1/(\sqrt{2})(|s>+\mathrm{exp}(i\varphi )|l>),$$
(7)
in spin language, is a state on the equator of the Poincaré sphere of states of a spin-half particle where $`|l>`$ and $`|s>`$ are the North and South poles. The choice of a phase delay $`\varphi _i`$ in the $`i`$th interferometer corresponds to the choice of a direction in the $`xy`$ plane along which one measures spin in Bohm’s version of the EPR gedankenexperiment.
An LHV description of the four photon interferometer requires the use of a space $`\mathrm{\Lambda }`$, the space of complete states whose elements are written $`\lambda `$, with a probability measure $`\rho `$. The expectation value of coincidence counts is then
$$\frac{1+E^\mathrm{\Psi }(\varphi _1,\varphi _2,\varphi _3,\varphi _4)}{2},$$
(8)
where
$`E^\mathrm{\Psi }(\varphi _1,\varphi _2,\varphi _3,\varphi _4)`$ $`=`$ $`<A(\varphi _1)B(\varphi _2)C(\varphi _3)D(\varphi _4)>`$
$`=`$ $`{\displaystyle _\mathrm{\Lambda }}A_\lambda (\varphi _1)B_\lambda (\varphi _2)C_\lambda (\varphi _3)D_\lambda (\varphi _4)𝑑\rho ,`$ (9)
and $`A_\lambda (\varphi _1),B_\lambda (\varphi _2),C_\lambda (\varphi _3),D_\lambda (\varphi _4)`$ are four functions of $`\lambda `$ which take values $`\pm 1`$. Locality is built into the theory by the fact that $`A_\lambda (\varphi _1)`$ is independent of $`\varphi _2,\varphi _3,\varphi _4`$, $`B_\lambda (\varphi _2)`$ is independent of $`\varphi _1,\varphi _3,\varphi _4`$, and so on.
### 5.1 Perfect correlations: ideal experiment
Following GHZ we can show that an LHV theory cannot reproduce the predictions of QM even at the level of perfect correlations.
Proof: Let us suppose functions $`A_\lambda (\varphi _1),B_\lambda (\varphi _2),C_\lambda (\varphi _3),D_\lambda (\varphi _4)`$ exist, satisfying the relations
$$<A(\varphi _1)B(\varphi _2)C(\varphi _3)D(\varphi _4)>=1,\mathrm{for}\mathrm{\Phi }=0,$$
(10)
and
$$<A(\varphi _1)B(\varphi _2)C(\varphi _3)D(\varphi _4)>=1,\mathrm{for}\mathrm{\Phi }=\pi .$$
(11)
Since the quantity in brackets can only take values $`\pm 1`$, it follows that everywhere in $`\mathrm{\Lambda }`$ (except possibly for a set of measure zero),
$$A_\lambda (\varphi _1)B_\lambda (\varphi _2)C_\lambda (\varphi _3)D_\lambda (\varphi _4)=1,\mathrm{for}\mathrm{\Phi }=0,$$
(12)
and
$$A_\lambda (\varphi _1)B_\lambda (\varphi _2)C_\lambda (\varphi _3)D_\lambda (\varphi _4)=1,\mathrm{for}\mathrm{\Phi }=\pi .$$
(13)
It then follows that
$$A_\lambda (\varphi )C_\lambda (\varphi )D_\lambda (0)B_\lambda (0)=1,$$
(14)
and
$$B_\lambda (0)A_\lambda (\varphi )D_\lambda (\varphi )C_\lambda (0)=1.$$
(15)
Multiplying equations (14) and (15), we get, since $`(A_\lambda (\varphi ))^2=(B_\lambda (0))^2=1`$,
$$C_\lambda (0)D_\lambda (0)C_\lambda (\varphi )D_\lambda (\varphi )=1.$$
(16)
But
$$A_\lambda (0)B_\lambda (0)C_\lambda (0)D_\lambda (0)=1.$$
(17)
Therefore
$$A_\lambda (0)B_\lambda (0)C_\lambda (\varphi )D_\lambda (\varphi )=1.$$
(18)
However, if we set $`\varphi =\pi /2`$ in equation (18), it contradicts equation (13). It follows that functions $`A_\lambda (\varphi _1),B_\lambda (\varphi _2),C_\lambda (\varphi _3),D_\lambda (\varphi _4)`$ satisfying equations (12) and (13) do not exist. Accordingly, LHV theories cannot reproduce the predictions of QM even at the level of perfect correlations.
### 5.2 Perfect correlations: real experiment
From section 4 we see that QM predicts a $`𝒬`$ of unity in an ideal experiment. However, in any real experiment, one would obtain a value for $`𝒬`$ less than unity because of imperfections in the system. However, as noted by Ryff , “if a theorem is valid whenever we have perfect correlations, it cannot be totally wrong in the case of almost perfect correlations”. We show below that with four-photon interference, a value of $`𝒬`$ greater than 0.5 is enough to rule out LHV theories. We do this by going beyond the original argument of GHZ to allow for experimental imperfections($`𝒬<1`$). Mermin has given an elegant and general analysis of the contradiction between quantum mechanics and LHV theories for $`n`$ spin-$`1/2`$ particles in an entangled state and our bound on $`𝒬`$ agrees with the restriction of Mermin’s analysis to the case of four particles.
Given two functions $`f`$ and $`g`$ on $`\mathrm{\Lambda }`$, let us define an inner product (or cross correlation)
$$<fg>=_\mathrm{\Lambda }f_\lambda g_\lambda 𝑑\rho .$$
(19)
We will only need to deal with functions which satisfy the condition
$$<ff>=1.$$
(20)
We then have the following lemma.
Lemma: Let $`f,g,h`$ be three functions on $`\mathrm{\Lambda }`$ with values $`\pm 1`$. Then,
$$<fh><fg>+<gh>1.$$
(21)
We present a proof and a geometrical interpretation of this lemma in Appendix B.
Let us then suppose that there exist functions $`A_\lambda (\varphi _1),B_\lambda (\varphi _2),C_\lambda (\varphi _3),D_\lambda (\varphi _4)`$ satisfying the relations
$$<A(\varphi _1)B(\varphi _2)C(\varphi _3)D(\varphi _4)>=𝒬,\mathrm{for}\mathrm{\Phi }=0,$$
(22)
and
$$<A(\varphi _1)B(\varphi _2)C(\varphi _3)D(\varphi _4)>=𝒬,\mathrm{for}\mathrm{\Phi }=\pi ,$$
(23)
where $`0𝒬1`$. (Note that in the limit, when $`𝒬1`$, we recover equations (10) and (11).) We can no longer argue, as we did before, that the angular brackets in equations (10) and (11) can be removed. However, the lemma can be used to determine the maximum allowed value for $`𝒬`$.
From equation (22), it follows that
$$<(A(\varphi )C(\varphi )D(0))(B(0))>=𝒬,$$
(24)
and
$$<(B(0))(A(\varphi )D(\varphi )C(0))>=𝒬.$$
(25)
If we use the lemma, with
$`f`$ $`=`$ $`A(\varphi )C(\varphi )D(0),`$ (26)
$`g`$ $`=`$ $`B(0),`$ (27)
$`h`$ $`=`$ $`A(\varphi )D(\varphi )C(0),`$ (28)
and remember that $`(A_\lambda (\varphi ))^2=(B_\lambda (0))^2=1`$, we get
$$<C(0)D(0)C(\varphi )D(\varphi )>2𝒬1.$$
(29)
However,
$$<A(0)B(0)C(0)D(0)>=𝒬,$$
(30)
so that, if we apply the lemma to these two relations, we find that
$$<A(0)B(0)C(\varphi )D(\varphi )>3𝒬2.$$
(31)
If then, we set $`\varphi =\pi /2`$ and use equation (23), we find that
$$(𝒬)3𝒬2,$$
(32)
from which it follows that
$$𝒬1/2.$$
(33)
This result proves that LHV theories cannot yield a value of $`𝒬`$ greater than 0.5.
Note that in our adaptation of the original GHZ argument , it is not necessary to set the individual phase differences $`\varphi _1,\mathrm{},\varphi _4`$ to $`0`$ (or, more correctly, $`2m\pi `$): it is only necessary to set $`\mathrm{\Phi }`$, the sum of these phase differences, to $`0`$ (or $`2m\pi `$). In practice, it is difficult (nearly impossible) to set the individual phase differences to any preassigned value, since the optical path differences in the individual interferometers are greater than the coherence lengths of the down-converted photons; our adaptation eliminates this problem and makes the experiment feasible.
In the actual experiment, one of the four interferometers is adjusted initially so that the coincidence rate is a maximum. The first measurement therefore corresponds to the condition $`\mathrm{\Phi }=2m\pi `$. A phase shift of $`\pi `$ is then introduced in any one of the interferometers and the event rate is measured at the resulting minimum. (Note that the introduction of a further phase shift of $`\pi `$ in any of the interferometers would bring the event rate back to a maximum). The results of these two measurements are inserted in equation (6) to obtain the value of the quantity $`𝒬`$. Any value greater than $`0.5`$ represents a breakdown of LHV theories under perfect correlations.
The only data used correspond effectively to values of $`\mathrm{\Phi }`$ of $`0`$ (eq. 22) and $`\pi `$ (eq. 23). This is very much in the spirit of the original GHZ argument , which relies only on perfect correlations.
### 5.3 Statistical correlations
This experiment also makes it possible to demonstrate violations of the original Bell inequality (which uses statistical correlations for two particles) in systems of four particles. In this case, it is not necessary to adjust the value of $`\mathrm{\Phi }`$ for the first measurement so that the coincidence rate is a maximum; $`\mathrm{\Phi }`$ can have any arbitrary value (say) $`\mathrm{\Phi }_0`$. Three more measurements of the event rate are then made after introducing phase shifts of $`\pi /2`$, successively, in three of the interferometers. We then have four values of the event rate corresponding to values of $`\mathrm{\Phi }`$ of $`\mathrm{\Phi }_0,\mathrm{\Phi }_0+\pi /2,\mathrm{\Phi }_0+\pi `$ and $`\mathrm{\Phi }_0+3\pi /2`$.
If one assumes that the fringe profile is sinusoidal, one can easily determine the fringe visibility from these four measurements. Whereas quantum mechanics predicts a visibility of unity, it has been shown that LHV theories cannot explain a fringe visibility greater than $`1/(2\sqrt{2})`$. In this respect, a four-photon experiment offers a more probing test than three-photon experiments for which the critical visibility is $`1/2`$. However, as explained in this lower value for the critical visibility relies on the use of statistical correlations rather than perfect correlations.
## 6 Conclusion
While most theoretical studies related to EPR have involved spin-($`1/2)`$ particles, actual experiments have used optical analogs of such systems. In particular, all interferometric tests of Bell’s inequality carried out so far have used entangled two-photon states \[4-7\]. In this case, LHV theories do not contradict QM at the level of perfect correlations. Therefore, tests of LHV theories with two-photon states require measurements of statistical correlations. EPR experiments involving more than two particles (as in section 5.3) utilize extensions of Bell’s inequality and, therefore, also involve statistical correlations. On the other hand, tests such as those described in section 5.2 are based on the GHZ analysis and, therefore, only involve perfect correlations. Since perfect correlations formed the basis of the original EPR criterion for “elements of reality”, the contradiction emerging from the GHZ analysis strikes at the heart of the EPR program. We have described a realizable experiment involving four-photon interference which demonstrates the conflict between EPR and QM even at the level of perfect correlations.
## Acknowledgements
One of the authors (P.H) thanks the International Centre for Theoretical Physics for support under their Visiting Scholar Program.
## Appendix A
The field theoretic analysis presented by Franson can be easily extended to the case of four-photon interference. A field theoretic description has the advantage that it is manifestly local. We sketch the main ideas below using his notation.
We need only deal with scalar fields since the polarization is fixed throughout. The scalar field operator $`\psi (\stackrel{}{r},t)`$ is expanded in free space modes as
$$\psi (\stackrel{}{r},t)=\underset{\stackrel{}{k}}{}\frac{a_\stackrel{}{k}}{\sqrt{V}}\mathrm{exp}(i(\stackrel{}{k}.\stackrel{}{r}\omega t)).$$
(34)
The time evolution of this operator is governed by the free Hamiltonian of the electromagnetic field and since
$$\psi (x+c\mathrm{\Delta }t,t)=\psi (x,t\mathrm{\Delta }t),$$
(35)
the particle it describes moves at the speed of light.
The field operator at the detector of the $`i`$th interferometer with the beam splitter removed is given by $`\psi _0(\stackrel{}{r}_i,t)`$. For each pair $`(i,j)`$ of the interferometers, these operators satisfy the condition
$$\psi _0(\stackrel{}{r}_i,t)\psi _0(\stackrel{}{r}_j,t\pm \mathrm{\Delta }t)|0>=0,$$
(36)
which is analogous to Franson’s equation (5). With the beam splitter inserted in the interferometer, the field operator at the $`i`$th detector becomes
$$\psi (\stackrel{}{r}_i,t)=(1/2)(\psi _0(\stackrel{}{r}_i,t)+\mathrm{exp}(i\varphi _i)\psi _0(\stackrel{}{r}_i,t\mathrm{\Delta }t)),$$
(37)
where $`i=1,\mathrm{},4`$. The coincidence rate $`R_c`$ for the four detectors $`D_1,D_2,D_3,D_4`$, with the beam splitters inserted, is then
$`R_c`$ $`=`$ $`\eta _1\eta _2\eta _3\eta _4\times `$
$`<0|\psi ^{}(r_1,t)\psi ^{}(\stackrel{}{r}_2,t)\psi ^{}(\stackrel{}{r}_3,t)\psi ^{}(\stackrel{}{r}_4,t)\psi (\stackrel{}{r}_1,t)\psi (\stackrel{}{r}_2,t)\psi (\stackrel{}{r}_3,t)\psi (\stackrel{}{r}_4,t)|0>,`$
where $`\eta _i`$ is the efficiency of the $`i`$th detector. Substituting (37) in (LABEL:a5) and using (36), we obtain the result
$$R_c=\left(\frac{R_{c0}}{2^6}\right)\left(\frac{1+\mathrm{cos}(\mathrm{\Phi })}{2}\right),$$
(39)
where
$`R_{c0}`$ $`=`$ $`\eta _1\eta _2\eta _3\eta _4\times `$
$`<0|\psi _0{}_{}{}^{}(\stackrel{}{r}_1,t)\psi _0{}_{}{}^{}(\stackrel{}{r}_2,t)\psi _0{}_{}{}^{}(\stackrel{}{r}_3,t)\psi _0{}_{}{}^{}(\stackrel{}{r}_4,t)\psi _0(\stackrel{}{r}_1,t)\psi _0(\stackrel{}{r}_2,t)\psi _0(\stackrel{}{r}_3,t)\psi _0(\stackrel{}{r}_4,t)|0>`$
and $`\mathrm{\Phi }=\varphi _1+\varphi _2+\varphi _3+\varphi _4`$.
## Appendix B
Let us consider three real-valued functions $`f,g,h`$ on $`\mathrm{\Lambda }`$. One can think of these functions as elements of a vector space. We only need to work in a three dimensional subspace containing $`f,g,h`$. If $`f,g,h`$ are unit vectors, $`<fg>=\mathrm{cos}(\theta _{fg})`$, $`<gh>=\mathrm{cos}(\theta _{gh})`$ and $`<fh>=\mathrm{cos}(\theta _{fh}),`$ where the angles $`\theta _{fg},\theta _{gh},\theta _{fh}`$ are defined to be less than $`\pi `$ and represent the angles between the unit vectors $`f,g,h`$. These angles can also be interpreted as the lengths of the shortest geodesics between the tips of the vectors $`f,g,h`$ on the unit sphere. From the triangle inequality, it then follows that
$$\theta _{fh}\theta _{fg}+\theta _{gh}.$$
(41)
Since $`\mathrm{cos}(\theta )`$ is a decreasing function of its argument for $`0\theta \pi `$, we arrive at the result
$$\mathrm{arccos}(<fh>)\mathrm{arccos}(<fg>)+\mathrm{arccos}(<gh>).$$
(42)
This inequality has a clear interpretation in terms of the triangle inequality on the unit sphere.
The inequality we use in the text is similar in spirit. We only need to deal with functions $`f,g,h`$ which take values $`\pm 1`$ and, in this case, can make a stronger statement. (Such functions and their correlations are of interest in digital signal processing, in communication theory and radio astronomy ). We then have
$$<fg>=2\mathrm{\Omega }(fg)1,$$
(43)
where $`\mathrm{\Omega }(fg)`$ is the volume of the domain $`𝒟(fg)`$ in $`\mathrm{\Lambda }`$ where $`f`$ and $`g`$ agree. Similarly, $`<gh>=2\mathrm{\Omega }(gh)1`$ and $`<fh>=2\mathrm{\Omega }(fh)1`$. Since the domain of agreement $`𝒟(fh)`$ between $`f`$ and $`h`$ includes at least $`𝒟(fg)𝒟(gh)`$, the intersection of the domains of agreement between $`f`$ and $`g`$ ($`𝒟(fg)`$) and between $`g`$ and $`h`$ ($`𝒟(gh))`$, we conclude that
$$\mathrm{\Omega }(fh)\mathrm{\Omega }(fg)+\mathrm{\Omega }(gh)1.$$
(44)
It follows immediately that
$$<fh><fg>+<gh>1$$
(45)
.
It is worth noting that if we write $`F(x)=(1x)/2`$, inequality (45) reads
$$F(<fh>)F(<fg>)+F(<gh>),$$
(46)
which is similar to inequality (42), with $`F(x)`$ replacing the function $`\mathrm{arccos}(x)`$. Both these inequalities express the idea that if $`f`$ and $`g`$ are highly correlated and $`g`$ and $`h`$ are highly correlated, then $`f`$ and $`h`$ must be correlated to some extent.
## Appendix C: Detector efficiency
In avalanche photodiodes, the only significant loss mechanism is reflection of the incident photons. The detection efficiency is therefore given by the relation
$$\eta =1R,$$
(47)
where $`R`$ is the fraction of photons reflected at the surface of the photocathode. It is, therefore, possible to reduce the loss due to this cause to negligible levels by using a number of photodiodes in a light- trapping arrangement .
A simple trap-detector which can be used for photon counting uses only two photodiodes . In this arrangement, as shown in Fig.$`2`$, the incident beam undergoes three reflections at the photodiodes before exiting. The fraction of the photons lost by reflection is then
$$R_2=R^3,$$
(48)
and summation of the outputs of the two photodiodes should yield a detection efficiency
$$\eta _2=1R^3.$$
(49)
With commercial avalanche photodiodes, for which $`R`$ is typically around $`0.3`$, it should be possible to obtain an increase in detection efficiency from $`70\%`$ to $`97\%`$.
## Appendix D: Generation of four-photon states
In the usual parametric process, yielding two down-converted photons, a 25 mm long ADP crystal pumped by a 9 mW He-Cd laser ($`\lambda =325nm`$), yields, at a 2 mm aperture placed at a distance of 1 m from the crystal, a down-converted flux of $`4\times 10^5`$ photons / second, for each beam , corresponding to a down-conversion efficiency of $`3\times 10^{11}`$. However, crystals such as beta-barium borate (BBO) are now available with a nonlinear coefficient 5 times higher than ADP. In addition, it should be possible to obtain an increase in down-conversion efficiency by placing the crystal in a short resonant cavity . If we use a 1.5 cm long BBO crystal, placed in a short cavity with mirrors whose reflectivity is chosen so that the effective length of the crystal is increased to around 7.5 cm, it should be possible to obtain a down -conversion efficiency of $`4.5\times 10^{10}`$.
The nonlinear susceptibility involved in the production of the four-photon field is, to a first approximation, the square of the nonlinear susceptibility involved in the production of two down-converted photons , so that the down-conversion efficiency, in this case, would work out to $`2\times 10^{19}`$. However, with two pump beams, the gain depends on the second power of the pump amplitude . As a result, the output with pulsed pump beams with high peak power can be several orders of magnitude greater than that obtained with continuous-wave excitation at the same average power.
With a laser generating pulses with a duration of 1 $`\mu `$s, at a repetition rate of 10 pulses/second, it should be possible to obtain a peak power that is $`10^5`$ times greater than the average power, and an improvement in down- conversion efficiency by a factor of this order. Accordingly, with an average power of 100 mW (corresponding to a peak power of 10 kW), it should be possible to obtain a total down-converted flux in the four output beams of $`4\times 10^3`$ photons/second, or $`10^3`$ photons/second in each beam. After allowing for losses, it should be possible to obtain a flux of 100 photons/second at the output from each of the four interferometers, which should permit useful measurements.
The use of a pulsed pump beam might be expected, at first sight, to create problems connected with the spectral coherence of the pump beams and the time resolution of the detectors. However, with a pulse duration of 1 $`\mu `$s, the coherence length of the pump beams, with a properly designed laser cavity, would be greater than 100 m. On the other hand, the coherence length of the down-converted beams, which would be determined by the decay time of the cavity modes (in this case, about 0.5 ns), would be less than 0.15 m. Accordingly, it would be possible to avoid second-order interference fringes by working with an optical path difference greater than this value, without a significant loss in the visibility of fourth-order interference effects.
As mentioned earlier, with a crystal placed in a resonant cavity, the light beams have an intrinsic bandwidth determined by the bandwidth of the cavity. This is consistent with the picture that the four down-converted photons are produced simultaneously, but then escape independently within a time interval equal to the decay time of the cavity , which, as mentioned earlier, is around 0.5 ns. Since this time interval is much less than the time resolution of a fast photodetector (say, 1.5 ns), the effects of such a deviation from simultaneity would not be noticeable.
Finally, we need to consider the probability of accidental coincidences. Since the output from each interferometer consists of a series of pulses with a duration of 1 $`\mu `$s, each containing about 10 photons, the probability of detecting a single photon in a time window of 1.5 ns would be 0.015. The ratio of the probability of accidental coincidences at the four outputs, due to uncorrelated photons, to that for actual coincidences would be only marginally higher, at around $`0.02`$. This proportion of accidental coincidences should not have a significant effect on the visibility of fourth-order interference effects produced by the four down-converted beams.
## Figure Captions
Fig. 1. Schematic of the four-photon interferometer.
Fig. 2. Optical configuration for a single-photon trap detector using two avalanche photodiodes.
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# Untitled Document
GENERALIZED QUATERNIONS and
INVARIANTS of VIRTUAL KNOTS and LINKS
ROGER FENN
School of Mathematical Sciences, University of Sussex
Falmer, Brighton, BN1 9RH, England
e-mail addresses: rogerf@sussex.ac.uk
ABSTRACT
In this paper we show how generalized quaternions including $`2\times 2`$ matrices can be used to find solutions of the equation
$$[B,\left(A1\right)(A,B)]=0.$$
These solutions can then be used to find polynomial invariants of virtual knots and links.
1 Introduction
Consider the algebra with the following presentation
$$=\{A,BA^1B^1ABBA^1B^1A=B^1ABA\}.$$
In this paper we will call this the fundamental algebra and the single relation will be called the fundamental relation or equation. This relation arises naturally from attempts to find representations of the braid group. Representations of the fundamental algebra as matrices can be used to define representations of the virtual braid group and invariants of virtual knots and links. In \[BuF\] we found a complete set of conditions for two classic quaternions, $`A,B`$ to be solutions of the fundamental equation. In this paper this result is generalised to give necessary and sufficient conditions for generalized quaternions to satisfy the fundamental relation, except in the case of all $`2\times 2`$ matrices where only sufficient conditions are given. Particularly, we define two 4-variable polynomials of virtual knots and links. In addition, we give conclusive proof of the fact, only hinted at in earlier papers, that invariants defined in this manner do not give any new invariants for classical knots and links.
We are grateful to Jose Montesinos for suggesting the use of generalized quaternions and to Steve Budden, Daan Krammer, Dale Rolfsen and Bruce Westfield for helpful comments.
2 The fundamental equation and its justification
Given a set $`X`$ let $`S`$ be an endomorphism of $`X^2`$. In \[FJK\] , such an $`S`$ is called a switch if
1 $`S`$ is invertible and
2 the set theoretic Yang-Baxter equation
$$(S\times id)(id\times S)(S\times id)=(id\times S)(S\times id)(id\times S)$$
is satisfied. Switches are used in \[FJK\] to define biracks and biquandles by the formula
$$S(a,b)=(b_a,a^b).$$
Given a switch $`S`$ there is a representation of the braid group $`B_n`$ into the group of permutations of $`X^n`$ defined by
$$\sigma _i(id)^{i1}\times S\times (id)^{ni1}$$
where $`\sigma _i`$ are the standard generators. Denote this representation by $`\rho =\rho (S,n)`$.
In this paper we will only be interested in linear switches. So let $`R`$ be an associative but not necessarily commutative ring and let $`X`$ be a left $`R`$-module. Suppose
$$S=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
where the matrix entries $`A,B,C,D`$ are elements of $`R`$. The 3$`\times `$3 matrices of the Yang-Baxter equation are
$$S\times id=\left(\begin{array}{ccc}A& B& 0\\ C& D& 0\\ 0& 0& 1\end{array}\right)id\times S=\left(\begin{array}{ccc}1& 0& 0\\ 0& A& B\\ 0& C& D\end{array}\right)$$
The representation $`\rho `$ is now into $`n\times n`$ matrices with entries from $`R`$.
Let us consider methods to find such switches $`S`$. It is not difficult to see that the following 7 equations are necessary and sufficient conditions for an invertible $`S`$ to be a switch,
$$\begin{array}{cc}1:A=A^2+BAC& 2:[B,A]=BAD\\ 3:[C,D]=CDA& 4:D=D^2+CDB\\ 5:[A,C]=DAC& 6:[D,B]=ADB\\ 7:[C,B]=& ADADAD\end{array}$$
where $`[X,Y]=XYYX`$.
Examples of switches are
$$\begin{array}{ccc}\hfill \text{ The identity }& & 0:\hfill \\ & & \\ \hfill S=\left(\begin{array}{cc}0& B\\ C& 1BC\end{array}\right)& \text{ or }S=\left(\begin{array}{cc}1BC& B\\ C& 0\end{array}\right)\hfill & 1:\hfill \end{array}$$
where $`B`$ and $`C`$ are arbitary commuting invertible elements. This is called the Alexander switch. A special case of this, when $`B=1`$, is called the Burau switch.
$$\begin{array}{ccc}\hfill S=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)& & 2:\hfill \end{array}$$
where $`A,A1,B`$ are invertible, $`A,B`$ do not commute and satisfy the fundamental equation
$$A^1B^1ABBA^1B^1A=B^1ABA$$
moreover
$$C=A^1B^1A(1A),D=1A^1B^1AB.$$
We will call this the non-commuting switch. A special case of this is the matrix with quaternion entries
$$S=\left(\begin{array}{cc}1+i& j\\ j& 1+i\end{array}\right)$$
called the Budapest switch.
If $`S=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ is a switch then so is $`S(t)=\left(\begin{array}{cc}A& tB\\ t^1C& D\end{array}\right)`$ where $`t`$ is a commuting variable. We say that $`S(t)`$ is $`S`$ augmented by $`t`$.
In \[BuF\] and \[BF\] the following results can be found.
Theorem 2.1 Suppose $`R`$ is a division ring. Then any switch is one of the examples above.
$`\mathrm{}`$
Of course other types are possible, see \[Cs\] in which divisors of zero are used.
The representation of the braid group induced by any non-commuting switch looks complicated but is in fact equivalent to the Burau representation. This has been pointed out previously by Dehornoy, see \[De\] . The following lemma gives an explicit proof.
Lemma 2.2 Let $`S=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ be a non-commuting switch and let $`S^{}=\left(\begin{array}{cc}0& 1\\ Q& 1Q\end{array}\right)`$ be the Burau switch where $`Q=(1A)(1D)`$. Let $`M`$ be the $`n\times n`$ matrix
$$M=\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ A& B& 0& \mathrm{}& 0\\ A& BA& B^2& \mathrm{}& 0\\ & \mathrm{}& & \mathrm{}& 0\\ A& BA& B^2A& \mathrm{}& B^{n1}\end{array}\right)$$
In words: the rows of $`M`$, after the first, start with $`A`$ and then the previous row multiplied on the left by $`B`$. Clearly $`M`$ is invertible. Then $`\rho (S,n)=M^1\rho (S^{},n)M`$.
Proof A calculation shows that $`M\rho (S,n)=\rho (S^{},n)M`$. In this calculation the fundamental relation is used. For example $`Q`$ commutes with $`B`$. So to prove that
$$B^iA^2+B^{i+1}C=QB^i+(1Q)B^iA$$
we need to show that
$$A^2+BC=Q+(1Q)A$$
which follows from the fundamental relation. $`\mathrm{}`$
However, if we extend the representation to the virtual braid group, defined below, then we get a representation which is not equivalent to the Burau.
The virtual braid group, $`VB_n`$ \[KK\] , has generators $`\sigma _i,i=1,\mathrm{},n1`$ and braid group relations
$$\begin{array}{ccc}\hfill \sigma _i\sigma _j& =\sigma _j\sigma _i,|ij|>1\hfill & i)\hfill \\ \hfill \sigma _i\sigma _{i+1}\sigma _i& =\sigma _{i+1}\sigma _i\sigma _{i+1}\hfill & \end{array}$$
In addition there are generators $`\tau _i,i=1,\mathrm{},n1`$ and permutation group relations
$$\begin{array}{ccc}\hfill \tau _{i}^{}{}_{}{}^{2}& =1\hfill & ii)\hfill \\ \hfill \tau _i\tau _j& =\tau _j\tau _i,|ij|>1\hfill & \\ \hfill \tau _i\tau _{i+1}\tau _i& =\tau _{i+1}\tau _i\tau _{i+1}\hfill & \end{array}$$
together with mixed relations
$$\begin{array}{ccc}\hfill \sigma _i\tau _j& =\tau _j\sigma _i,|ij|>1\hfill & iii)\hfill \\ \hfill \sigma _i\tau _{i+1}\tau _i& =\tau _{i+1}\tau _i\sigma _{i+1}\hfill & \end{array}$$
We can extend the representation $`\rho (S,n)`$ by sending the generator $`\tau _i`$ to $`(id)^{i1}\times T\times (id)^{ni1}`$ where $`T=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, that is, Burau with unit variable.
Consider now the element $`\beta =\sigma _2\sigma _1\tau _2\sigma _1^1\sigma _2^1\tau _1`$ in $`VB_3`$. If $`S=\left(\begin{array}{cc}0& B\\ C& 1BC\end{array}\right)`$, the Alexander switch, then
$$\rho (\beta )=\left(\begin{array}{ccc}1& 0& (C^1B)(B1)\\ 0& 1& (C^1B)(1B)\\ 0& 0& 1\end{array}\right)$$
So if $`B1`$ this representation can not be equivalent to the Burau representation which has $`B=1`$.
We can now ask the following question: Let $`KB_n`$ be the kernel of the Burau representation. Let $`\overline{K}`$ denote the normal closure of $`K`$ in $`VB_n`$. If $`\beta `$ is a virtual braid and $`\rho (\beta )=1`$ for all switches $`S`$, is it true that $`\beta `$ lies in $`\overline{K}`$?
3 Quaternion Algebras
It is clear from the previous section that it is important to find solutions to the fundamental equation. The main result in the next section is a sufficient condition for two generalised quaternions to satisfy the fundamental equation. Except for $`2\times 2`$ matrices, this condition is also necessary.
In this section we describe the necessary algebra. The results which are already in the literature are mainly presented without proof. For more details see \[L\] .
Let $`F`$ be a field of characteristic not equal to 2. Pick two non-zero elements $`\lambda ,\mu `$ in $`F`$. Let $`\left(\frac{\lambda ,\mu }{F}\right)`$ denote the algebra of dimension 4 over $`F`$ with basis $`\{1,i,j,k\}`$ and relations $`i^2=\lambda ,j^2=\mu ,ij=ji=k`$. The multiplication table is given by
$$\begin{array}{cccccc}& & i& j& k& \\ i& (\mathrm{}& \lambda & k& \lambda j& )\mathrm{}\\ j& k& \mu & \mu i\\ k& \lambda j& \mu i& \lambda \mu \end{array}.$$
Throughout the paper a general quaternion algebra will be denoted by $`𝒬`$. Elements of $`𝒬`$ are called (generalized) quaternions. The field $`F`$ is called the underlying field and the elements $`\lambda \mu `$ the parameters of the algebra. We will denote quaternions by capital roman letters such as $`A,B,\mathrm{}`$ and (if pure) by bold face lower case, $`\text{a},\text{b},\mathrm{}`$. Field elements, (scalars) will be denoted by lower case roman letters such as $`a,b,\mathrm{}`$ and lower case greek letters such as $`\alpha ,\beta ,\mathrm{}`$.
The classical quaternions are $`\left(\frac{1,1}{\text{}}\right)`$. The algebra of $`2\times 2`$ matrices with entries in $`F`$ is $`M_2(F)=\left(\frac{1,1}{F}\right)`$.
3.1 Conjugation, Norm and Trace
Let $`A=a_0+a_1i+a_2j+a_3k`$ be a quaternion where $`a_0,a_1,a_2,a_3F`$. The coordinate $`a_0`$ is called the scalar part of $`A`$ and the 3-vector $`\text{a}=a_1i+a_2j+a_3k`$ is called the pure part of $`A`$. Evidently $`A=a_0+\text{a}`$ is the sum of its scalar and pure parts and is pure if its scalar part is zero and is a scalar if its pure part is zero.
The conjugate of $`A`$ is $`\overline{A}=a_0\text{a}`$, the norm of $`A`$ is $`N(A)=A\overline{A}`$ and the trace of $`A`$ is $`\text{tr}(A)=A+\overline{A}`$.
Conjugation is an anti-isomorphism of order 2. That is it satisfies
$$\overline{A+B}=\overline{A}+\overline{B},\overline{AB}=\overline{B}\overline{A},\overline{aA}=a\overline{A},\overline{\overline{A}}=A.$$
Also $`\overline{A}=A`$ if and only if $`A`$ is a scalar and $`\overline{A}=A`$ if and only if $`A`$ is pure.
The norm is a scalar satisfying $`N(AB)=N(A)N(B)`$. We will denote the set of values of the norm function by $`𝒩`$. It is a multiplicatively closed subset of $`F`$ and $`𝒩^{}=𝒩\{0\}`$ is a multiplicative subgroup of $`F^{}`$. An element $`A`$ has an inverse if and only if $`N(A)0`$ in which case $`A^1=N(A)^1\overline{A}`$.
The trace of a quaternion is twice its scalar part.
3.2 Multiplying Quaternions
Let $`A,B`$ be two quaternions.There is a bilinear form given by
$$AB=\frac{1}{2}(A\overline{B}+B\overline{A})=\frac{1}{2}(\overline{A}B+\overline{B}A)=\frac{1}{2}\text{tr}(A\overline{B}).$$
In terms of coordinates this is
$$AB=a_0b_0\lambda a_1b_1\mu a_2b_2+\lambda \mu a_3b_3.$$
Since $`\lambda `$ and $`\mu `$ are non-zero this is a non-degenerate form. The corresponding quadratic form is
$$N(A)=a_0^2\lambda a_1^2\mu a_2^2+\lambda \mu a_3^2.$$
Let $`\text{a},\text{b}`$ be pure quaternions. Then
$$\text{a}\text{b}=\text{a}\text{b}+\text{a}\times \text{b}$$
where
$$\text{a}\text{b}=\lambda a_1b_1\mu a_2b_2+\lambda \mu a_3b_3$$
is the restriction of the bilinear form to the pure quaternions and $`\text{a}\times \text{b}`$ is the cross product defined symbolically by
$$\text{a}\times \text{b}=\left|\begin{array}{ccc}\mu i& \lambda j& k\\ a_1& a_2& a_3\\ b_1& b_2& b_3\end{array}\right|$$
The cross product has the usual rules of bilinearity and skew symmetry. The triple cross product expansion
$$\text{a}\times (\text{b}\times \text{c})=(\text{c}\text{a})\text{b}(\text{b}\text{a})\text{c}$$
is easily verified. The scalar triple product is
$$[\text{a},\text{b},\text{c}]=\text{a}(\text{b}\times \text{c})=\lambda \mu \left|\begin{array}{ccc}a_1& a_2& a_3\\ b_1& b_2& b_3\\ c_1& c_2& c_3\end{array}\right|$$
from which all the usual rules (except volume) can be deduced.
3.3 Dependancy Criteria
In this subsection we will consider conditions for sets of quaternions to be linearly dependant or otherwise. A non-zero element, $`A`$, of $`𝒬`$ is called isotropic if $`N(A)=0`$ and anisotropic otherwise. So only non-zero anisotropic elements have inverses. We note the following theorem.
Theorem 3.2 The following statements about a quaternion algebra $`𝒬`$ are equivalent.
1. $`𝒬`$ contains an isotropic element.
2. $`𝒬`$ is the sum of two hyperbolic planes.
3. $`𝒬`$ is not a division algebra.
4. $`𝒬`$ is $`M_2(F)`$.
Proof See \[L\] p 58.
We will call a quaternion algebra above hyperbolic. Otherwise it is called anisotropic. The classic quaternions are anisotropic: $`2\times 2`$ matrices are hyperbolic.
Lemma 3.3 A pair of pure quaternions $`\text{a},\text{b}`$ is linearly dependant if and only if $`\text{a}\times \text{b}=0`$.
Proof The proof is clear one way using the antisymmetry of the cross product. Conversely suppose $`\text{a}\times \text{b}=0`$. Then $`(\text{a}\times \text{b})\times \text{c}=(\text{a}\text{c})\text{b}(\text{b}\text{c})\text{a}=0`$. This can be made into a linear dependancy by a suitable choice of c, for example if $`\text{a}\text{c}0`$. $`\mathrm{}`$
As a corollary we have the following
Lemma 3.4 Two quaternions commute if and only their pure parts are linearly dependant.
$`\mathrm{}`$
Now we look for conditions for the triple of pure quaternions, $`\text{a},\text{b},\text{a}\times \text{b}`$, to be linearly dependant. The required condition is given by the following lemma.
Lemma 3.5 The pure quaternions $`\text{a},\text{b},\text{a}\times \text{b}`$, are linearly dependant if and only if
$$N(\text{a})N(\text{b})=(\text{a}\text{b})^2.$$
This is equivalent to the equations
$$N(\text{a}\times \text{b})=\mu (a_2b_3a_3b_2)^2\lambda (a_1b_3a_3b_1)^2+(a_1b_2a_2b_1)^2=0,$$
ie $`\text{a}\times \text{b}`$ is isotropic or zero.
Proof Three 3-dimensional vectors are linearly dependant if and only if the determinant they form by rows is zero. In the case of pure quaternions this means the scalar triple product is zero
$$[\text{a},\text{b},\text{c}]=\text{a}(\text{b}\times \text{c})=\lambda \mu \left|\begin{array}{ccc}a_1& a_2& a_3\\ b_1& b_2& b_3\\ c_1& c_2& c_3\end{array}\right|=0.$$
Replacing c with $`\text{a}\times \text{b}`$ and expanding out using the triple cross product formula gives the first equation. Using the expansion formulæ
$$N(\text{a}\times \text{b})=N(\text{a})N(\text{b})(\text{a}\text{b})^2$$
gives the second formula. $`\mathrm{}`$
We have the following corollary.
Lemma 3.6 If $`\text{a},\text{b}`$ are linearly independant pure quaternions and $`\text{a}\times \text{b}`$ is anisotropic, then the triple $`\text{a},\text{b},\text{a}\times \text{b}`$, is linearly independant.
3.4 $`2\times 2`$ matrices
We will interpret all the previous results in terms of $`2\times 2`$ matrices, $`M_2(F)=\left(\frac{1,1}{F}\right)`$. This is the only quaternion algebra with zero divisors.
The generators of $`\left(\frac{1,1}{F}\right)`$ are, together with the identity, the Pauli matrices
$$i=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),j=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),k=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
By an abuse of notation we will often confuse the scalar matrix $`\left(\begin{array}{cc}\nu & 0\\ 0& \nu \end{array}\right)`$ with the corresponding field element $`\nu `$.
A general matrix can be written uniquely as
$$\left(\begin{array}{cc}\alpha & \beta \\ \gamma & \delta \end{array}\right)=\frac{1}{2}\left[(\alpha +\delta )+(\beta \gamma )i+(\beta +\gamma )j+(\alpha \delta )k\right]$$
Conversely
$$A=a_0+a_1i+a_2j+a_3k=\left(\begin{array}{cc}a_0+a_3& a_2+a_1\\ a_2a_1& a_0a_3\end{array}\right)$$
Conjugation is
$$\overline{A}=\text{adj}A=\left(\begin{array}{cc}\delta & \beta \\ \gamma & \alpha \end{array}\right)=\left(\begin{array}{cc}a_0a_3& a_2a_1\\ a_1a_2& a_0+a_3\end{array}\right)$$
and norm is
$$N(A)=detA=\alpha \delta \beta \gamma =a_0^2+a_1^2a_2^2a_3^2$$
The scalar part of $`A`$ is $`a_0=\text{tr}A/2=(\alpha +\delta )/2`$ and the pure part is
$$\left(\begin{array}{cc}a_3& a_2+a_1\\ a_2a_1& a_3\end{array}\right)=\left(\begin{array}{cc}(\alpha \delta )/2& \beta \\ \gamma & (\delta \alpha )/2\end{array}\right)$$
3.5 Multiplying Matrices
Lemma 3.2 Suppose $`A,BM_2(F)=\left(\frac{1,1}{F}\right)`$. Then $`AB=AB`$.
The statement is deliberately provocative. It says that multiplying $`A,B`$ as matrices and as quaternions is the same. This can be checked directly. $`\mathrm{}`$
The above lemma allows quick checking of formulæ so if
$$A=\left(\begin{array}{cc}\alpha _1& \alpha _2\\ \alpha _3& \alpha _4\end{array}\right)\text{ and }B=\left(\begin{array}{cc}\beta _1& \beta _2\\ \beta _3& \beta _4\end{array}\right)\text{ then }AB=\frac{1}{2}(\alpha _1\beta _4\alpha _2\beta _3\alpha _3\beta _2+\alpha _4\beta _1)$$
If $`\text{a}=\left(\begin{array}{cc}\alpha _1& \alpha _2\\ \alpha _3& \alpha _1\end{array}\right)`$ and $`\text{b}=\left(\begin{array}{cc}\beta _1& \beta _2\\ \beta _3& \beta _1\end{array}\right)`$ are pure then $`\text{a}\text{b}=\alpha _1\beta _1(\alpha _2\beta _3+\alpha _3\beta _2)/2`$ and
$$\text{a}\times \text{b}=\left(\begin{array}{cc}(\alpha _2\beta _3\alpha _3\beta _2)/2& \alpha _1\beta _2\alpha _2\beta _1\\ \alpha _3\beta _1\alpha _1\beta _3& (\alpha _3\beta _2\alpha _2\beta _3)/2\end{array}\right)$$
4 Solving the Fundamental Equation
In this section we give a sufficient condition for two generalised quaternions to satisfy the fundamental equation. Except for $`2\times 2`$ matrices, this condition is also necessary.
Let $`A=a_0+𝐚`$ and $`B=b_0+𝐛`$ be quaternions. We will need the following easily checked lemma.
Lemma 4.3 Conjugation by multiplication is $`B^1AB=N(B)^1\overline{B}AB`$ where
$$\overline{B}AB=a_0(b_0^2+N(\text{b}))+(b_0^2N(\text{b}))\text{a}+2(\text{a}\text{b})\text{b}+2b_0(\text{a}\times \text{b}).$$
The two commutators are
$$[A,B]=ABBA=2\text{a}\times \text{b}$$
and $`(A,B)=A^1B^1AB=N(A)^1N(B)^1\overline{A}\overline{B}AB`$ where
$$\begin{array}{cc}\hfill \overline{A}\overline{B}AB& =a_0^2b_0^2+b_0^2N(\text{a})+a_0^2N(\text{b})N(\text{a})N(\text{b})+2(\text{a}\text{b})^2\hfill \\ & 2(b_0(\text{a}\text{b})+a_0N(\text{b}))\text{a}+2(a_0(\text{a}\text{b})+b_0N(\text{a}))\text{b}+2(a_0b_0\text{a}\text{b})\text{a}\times \text{b}\hfill \end{array}.$$
$`\mathrm{}`$
We will call a quaternion, $`A`$, balanced if $`\text{tr}(A)=N(A)0`$. A quaternion $`A=a_0+a_1i+a_2j+a_3k`$ in a quaternion algebra with parameters $`\lambda ,\mu `$ is balanced if it lies on the quadric 3-fold
$$(a_01)^2\lambda a_1^2\mu a_2^2+\lambda \mu a_3^2=1,a_00.$$
A balanced classical quaternion $`A`$ lies on the 3-sphere centre 1 and radius 1. Note that if $`A`$ is balanced then $`N(A1)=1`$. A pair of invertible non-commuting quaternions, $`A,B`$, will be called matching if $`A`$ is balanced and $`AB=0`$.
Theorem 4.4 If the quaternion algebra is anisotropic then a necessary and sufficient condition for the non-commuting, invertible quaternions $`A,B`$ to be solutions of the fundamental equation is that they are a matching pair. Otherwise the condition is only sufficient.
Proof The proof formally follows \[BuF\] . In terms of quaternions the equation is
$$\overline{A}\overline{B}ABB\overline{A}\overline{B}A=N(A)\overline{B}ABN(A)N(B)A.$$
Using the formulæ and notation developed above, the left hand side is
$$4a_0N(\text{b})\text{a}+4a_0(\text{a}\text{b})\text{b}4(\text{a}\text{b})\text{a}\times \text{b}$$
whereas the right hand side is
$$2(a_0^2+N(\text{a}))N(\text{b})\text{a}+2(a_0^2+N(\text{a}))(\text{a}\text{b})\text{b}+2b_0(a_0^2+N(\text{a}))\text{a}\times \text{b}$$
Considering half the difference of the two sides we arrive at
$$\text{c}=(\text{tr}(A)N(A))N(\text{b})\text{a}+(N(A)\text{tr}(A))(\text{a}\text{b})\text{b}+(b_0(N(A)\text{tr}(A))+2AB)\text{a}\times \text{b}$$
So the two sides are equal if $`\text{c}=0`$.
If $`A,B`$ is a matching pair then $`\text{c}=0`$. Conversely if $`\text{a},\text{b},\text{a}\times \text{b}`$ are linearly independent then their coefficients will be zero and this implies that $`A,B`$ is a matching pair. The bilinear form is definite unless the algebra is $`M_2(F)`$ and so, except for this case, $`\text{a},\text{b},\text{a}\times \text{b}`$ will be linearly independent. $`\mathrm{}`$
4.1 The Non-definite Case
The condition for linear dependancy can be satisfied for the case when the quaternions are $`2\times 2`$ matrices, ie $`\lambda =1,\mu =1`$. Let us try
$$a_2b_3a_3b_2=1,a_1b_3a_3b_1=1,a_1b_2a_2b_1=0$$
One solution is
$$\text{a}=i+tjk,\text{b}=j\text{ and }\text{a}\times \text{b}=ki$$
where $`t`$ is an arbitary field element. So
$$\text{a}t\text{b}+\text{a}\times \text{b}=0$$
We will return to this case in a later paper.
5 The general matching pair
In this section we will use the results of the previous section to describe the most general matching pair. That is $`A,B`$ are $`2\times 2`$ matrices with entries in some field and satisfying
1. $`\text{tr}(A)=det(A)`$
2. $`A\text{adj}(B)+B\text{adj}(A)=0`$
Since $`B`$ can be multiplied by any non-zero scalar we may assume temporarily that
3. $`det(B)=1`$.
We can conjugate the matrices $`A,B`$ to simplify matters. Consider the following two cases:
Case 1, $`A`$ has two distinct eigenvalues and is diagonal.
$$A=\left(\begin{array}{cc}a& 0\\ 0& a/(a1)\end{array}\right)\text{ and }B=\left(\begin{array}{cc}b& c\\ (b^2+a1)/c(1a)& b/(1a)\end{array}\right)$$
where $`a,b,c`$ are general and $`c,1a`$ must be invertible (ie non zero).
Inverses are given by
$$A^1=\frac{1}{a}\left(\begin{array}{cc}1& 0\\ 0& a1\end{array}\right)$$
so $`a`$ must also be invertible and
$$B^1=\text{adj}(B)=\left(\begin{array}{cc}b/(1a)& c\\ (b^2+a1)/c(1a)& b\end{array}\right).$$
The $`4\times 4`$ matrix $`S`$ is $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ where
$$C=A^1B^1AA^1B^1A^2=\left(\begin{array}{cc}b& c/(1a)^2\\ (1a)(b^2+a1)/c& b/(1a)\end{array}\right)$$
and
$$D=1A^1B^1AB=\left(\begin{array}{cc}(23a+ab^2+a^22b^2)/(1a)^2& (a2)bc/(1a)^2\\ (a2)b(b^2+a1)/c(1a)& (23a+ab^2+a^22b^2)/(1a)\end{array}\right)$$
Call this switch $`𝐄_\mathrm{𝟐}`$
Case 2 $`A`$ has one eigenvalue and is lower triangular
$$A=\left(\begin{array}{cc}2& 0\\ x& 2\end{array}\right).$$
Then $`B`$ has the form,
$$B=\left(\begin{array}{cc}y& z\\ (xyz2y^22)/2z& (xz2y)/2\end{array}\right)$$
where $`x,y,z`$ are general and $`2,z`$ must be invertible.
Inverses are given by
$$A^1=\frac{1}{4}\left(\begin{array}{cc}2& 0\\ x& 2\end{array}\right).$$
and
$$B^1=\text{adj}(B)=\left(\begin{array}{cc}(xz2y)/2& z\\ (xyz+2y^2+2)/2z& y\end{array}\right)$$
The $`4\times 4`$ matrix $`S`$ is $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ where
$$C=A^1B^1AA^1B^1A^2=\left(\begin{array}{cc}y+xz& z\\ (x^2z^2+3xyz+2y^2+2)/2z& xz/2y\end{array}\right)$$
and
$$D=1A^1B^1AB=\left(\begin{array}{cc}xyz/2& xz^2/2\\ x(2y^2+xyz2)/4& xz(xz+2y)/4\end{array}\right)$$
Call this switch $`𝐄_\mathrm{𝟏}`$
5.1 The variable $`t`$
Let us now multiply $`B`$ by the scalar $`t`$, ie the switch is augmented by $`t`$. This means we have two possible switches: one dependant on four variables $`a,b,c,t`$ ($`E_2`$) and one dependant on four variables $`x,y,z,t`$ ($`E_1`$).
6 Determinants over Quaternion Algebras
In order to define workable invariants we consider in this section a determinental function on the matrices in $`M_n(𝒬)`$. That is $`n\times n`$ matrices with entries in a quaternion algebra $`𝒬`$. For background reading see \[As\] . In fact the invariants defined later can also be defined for any solutions of the fundamental equation over a ring with a determinant function satisfying the rules listed below.
If $`R`$ is a commutative ring let $`det:M_n(R)R`$ denote the usual determinant. The classic quaternions, $``$, may be embedded as a subalgebra of $`M_2(\text{})`$ and determinants taken in the usual way. Our aim is to generalize this.
Suppose $`𝒬`$ has underlying field $`F`$ and parameters $`\lambda ,\mu `$. Let $`\overline{F}`$ denote the algebraic closure of $`F`$. Embed $`i,j,k`$ in $`M_2(\overline{F})`$ by
$$i=\left(\begin{array}{cc}0& \sqrt{\lambda }\\ \sqrt{\lambda }& 0\end{array}\right),j=\left(\begin{array}{cc}0& \sqrt{\mu }\\ \sqrt{\mu }& 0\end{array}\right),k=\left(\begin{array}{cc}\sqrt{\lambda \mu }& 0\\ 0& \sqrt{\lambda \mu }\end{array}\right).$$
Define $`d:M_n(𝒬)\overline{F}`$ as the composition of the embedding $`M_n(𝒬)M_{2n}(\overline{F})`$ with $`det`$.
Alternatively the determinant function may be defined by induction on the size of the matrices. The value $`d(A)=N(A)`$ starts the induction. Consider a matrix in $`M_n(𝒬)`$. This may be reduced to diagonal form, by multiplying on the left and the right by elementary matrices having unit determinant, (see below). Suppose this matrix has diagonal elements $`d_1,\mathrm{},d_n`$. Define the determinant as $`d=N(d_1)\mathrm{}N(d_n)`$. So the determinant function takes values in $`𝒩`$. For $`M_2(F)`$ this subset is the whole of $`F`$: for classic quaternions it is the non-negative reals.
The determinant function satisfies the rules
0. $`d(M)=0`$ if and only if $`M`$ is singular, moreover $`d(MN)=d(M)d(N)`$. It follows that $`d(1)=1`$.
1. $`d`$ is unaltered by a permutation of the rows (columns).
2. If a row (column) is multiplied on the left (right) by a unit then $`d`$ is multiplied by $`d`$ of that unit.
3. $`d(M)`$ is unaltered by adding a left multiple of a row to another row or a right multiple of a column to another column.
4. $`d\left(\begin{array}{cc}x& 𝐮\\ \mathrm{𝟎}& M\end{array}\right)=N(x)d(M)`$ where $`𝐮`$ is any row vector and $`\mathrm{𝟎}`$ is a zero column vector.
4’. $`d\left(\begin{array}{cc}x& \mathrm{𝟎}\\ 𝐯& M\end{array}\right)=N(x)d(M)`$ where $`𝐯`$ is any column vector and $`\mathrm{𝟎}`$ is a zero row vector.
5. $`d(M^{})=\text{d}(M)`$ where $`M^{}=\overline{M^T}`$ denotes the Hermitian conjugate.
6. if the entries in $`M`$ all commute then $`d(M)=det^2(M)`$.
An elementary matrix of type 1 is a permutation matrix. An elementary matrix of type 2 is the identity matrix with one diagonal entry replaced by a unit and an elementary matrix of type 3 is a square matrix with zero entries except for $`1`$’s down the diagonal and one other entry off diagonal.
The properties $`i.`$ above for $`i=1,2,3`$ follow from multiplying $`M`$ on the right or left by an elementary matrix of type $`i`$.
The matrix $`S`$ can be written as a product of elementary matrices
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)=\left(\begin{array}{cc}A& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ C& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1A^1\end{array}\right)\left(\begin{array}{cc}1& A^1B\\ 0& 1\end{array}\right).$$
Note that $`1A^1`$ is invertible.
Hence $`d(S)=d(A)d(1A^1)=d(A1)=1`$. Therefore the representation of $`VB_n`$, induced by such an $`S`$, is into $`SL(F,2n)`$.
7 Virtual Knots and Links
Recall that classical knot theory can be described in terms of knot and link diagrams. A diagram is a 4-regular plane graph (with extra structure at its nodes representing the crossings in the link) on a plane and implicitly on a two-dimensional sphere $`S^2`$. We say that two such diagrams are equivalent if there is a sequence of moves of the types indicated in part (A) of Figure 1 (The Reidemeister Moves) taking one diagram to the other. These moves are performed locally on the 4-regular plane graph (with extra structure) that constitutes the link diagram.
Virtual knot theory is an extension of classical knot theory, see \[K\] . In this extension one adds a virtual crossing (See Figure 1) that is neither an over-crossing nor an under-crossing. We shall refer to the usual diagrammatic crossings, that is those without circles, as real crossings to distinguish them from the virtual crossings. A virtual crossing is represented by two crossing arcs with a small circle placed around the crossing point. The arcs of the graph joining adjacent classical crossings are called the semi-arcs of the diagram.
In addition to their application as a geometric realization of the combinatorics of a Gauss code, virtual links have physical, topological and homological applications. In particular, virtual links may be taken to represent a particle in space and time which dissappears and reappears. A virtual link may be represented, up to stabilisation, by a link diagram on an orientable surface, \[Ku\] . If the surface has minimal genus then this representation is unique. Finally an element of the second homology of a rack space can be represented by a labelled virtual link, see \[FRS1\] \[FRS2\] . Since the rack spaces form classifying spaces for classical links the study of virtual links may give information about classical knots and links.
The allowed moves on virtual diagrams are a generalization of the Reidemeister moves for classical knot and link diagrams. We show the classical Reidemeister moves as part (A) of Figure 1. These classical moves are part of virtual equivalence where no changes are made to the virtual crossings. Taken by themselves, the virtual crossings behave as diagrammatic permutations. Specifically, we have the flat Reidemeister moves (B) for virtual crossings as shown in Figure 1. In Figure 1 we also illustrate a basic move (C) that interrelates real and virtual crossings. In this move an arc going through a consecutive sequence of two virtual crossings can be moved across a single real crossing. In fact, it is consequence of moves (B) and (C) for virtual crossings that an arc going through any consecutive sequence of virtual crossings can be moved anywhere in the diagram keeping the endpoints fixed and writing the places where the moved arc now crosses the diagram as new virtual crossings. This is shown schematically in Figure 2. We call the move in Figure 2 the detour, and note that the detour move is equivalent to having all the moves of type (B) and (C) of Figure 1. This extended set of moves (Reidemeister moves plus the detour move or the equivalent moves (B) and (C)) constitutes the set of moves for diagrams of virtual knots and links.
The topological interpretation for this virtual theory in terms of embeddings of links in thickened surfaces is a useful idea. See \[KK\] , \[Ku\] . Regard each virtual crossing as a shorthand for a detour of one of the arcs in the crossing through a 1-handle that has been attached to the 2-sphere of the original diagram. The two choices for the 1-handle detour are homeomorphic to each other (as abstract surfaces with boundary a circle) since there is no a priori difference between the meridian and the longitude of a torus. By interpreting each virtual crossing in this way, we obtain an embedding of a collection of circles into a thickened surface $`S_g\times \text{}`$ where $`g`$ is the number of virtual crossings in the original diagram $`L`$, $`S_g`$ is a compact oriented surface of genus $`g`$ and denotes the real line. Thus to each virtual diagram $`L`$ we obtain an embedded disjoint union of circles in $`S_{g(L)}\times \text{}`$ where $`g(L)`$ is the number of virtual crossings of $`L`$. We say that two such surface embeddings are stably equivalent if one can be obtained from another by isotopy in the thickened surfaces, homeomorphisms of the surfaces and the addition or subtraction of empty handles. Then we have the
Theorem 7.2 Two virtual link diagrams are equivalent if and only if their correspondent surface embeddings are stably equivalent, \[KK\] , \[Ku\] .
The surface embedding interpretation of virtuals is useful since it converts their equivalence to a topological question. The diagrammatic version of virtuals embodies the stabilization in the detour moves. We shall rely on the diagrammatic approach here.
8 The invariant knot modules
In this section, we shall begin to show how the previous algebra can give rise to virtual knot invariants. Given an associative ring $`R`$, a $`2\times 2`$ matrix $`S`$ with entries in $`R`$ and a virtual link diagram $`𝒟`$, we define a presentation of an $`R`$-module which depends only on the link class of the diagram provided $`S`$ is a switch. This construction also works for classical knots and links but is only the Alexander module in disguise. The generators are the semi-arcs of $`𝒟`$, that is the portion of the diagram bounded by two adjacent classical crossings. There are 2 relations for each classical crossing.
Suppose the diagram $`𝒟`$ has $`n`$ classical crossings. Then there are $`2n`$ semi-arcs labelled $`a,b,\mathrm{}`$. These will be the generators of the module. Let the edges of a positive real crossing in a diagram be arranged diagonally and called geographically NW, SW, NE and SE. Assume that initially the crossing is oriented and the edges oriented towards the crossing from left to right ie west to east. The input edges, oriented towards the crossing, are in the west and the edges oriented away from the crossing, the output edges, are in the east. Let $`a`$ and $`b`$ be the labels of the input edges with $`a`$ labelling SW and $`b`$ labelling NW. For a positive crossing, $`a`$ will be the label of the undercrossing input and $`b`$ the label of the overcrossing input. Suppose now that
$$S(a,b)^T=(c,d)^T\text{ where }S=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right).$$
Then we label the undercrossing output NE by $`d`$ and we label the overcrossing output SE by $`c`$.
For a negative crossing the direction of labelling is reversed. So $`a`$ labels SE, $`b`$ labels NE, $`c`$ labels SW and $`d`$ labels NW.
Finally for a virtual crossing the labellings carry across the strings. This corresponds to the twist function $`T(a,b)=(b,a)`$.
The following figure shows the labelling for the three kind of crossings and the corresponding relations for the 2 classical crossings.
$`c=Aa+Bbd=Ca+Db`$
The diagram therefore gives rise to a presentation of an $`R`$-module with $`2n`$ generators and $`2n`$ relations. Note that in all cases $`B,C`$ are invertible since the identity switch is uninteresting.
Theorem 8.3 The module defined above for any diagram $`D`$ is invariant under the Reidemeister moves, and hence is a knot invariant, if $`S`$ is a switch.
Proof The proof of the above can be found in the papers \[FJK, BF, BuF\] , For the convenience of the reader we show how the module is invariant under the Reidemeister moves.
Refering back to the picture of the two relations defined by a crossing, it is convenient to think of the action from left to right on a positive crossing as being the action of $`S`$ and the action from right to left as being $`S^1`$.
Consider the action from top to bottom as being $`S_+^{}`$ and the action from bottom to top as being $`S_{}^+`$. By solving the equations of the labellings we see that these matrices are
$$S_{}^+=\left(\begin{array}{cc}DB^1& CDB^1A\\ B^1& B^1A\end{array}\right)S_+^{}=\left(\begin{array}{cc}C^1D& C^1\\ BAC^1D& AC^1\end{array}\right)$$
We call $`S_{}^+`$ and $`S_+^{}`$ the sideways matrices. They are invertible since $`S`$ is. Also $`(S^1)_{}^+=(S_+^{})^1`$ and $`(S^1)_+^{}=(S_{}^+)^1`$ and
$$S_{}^+(a,a)=(\lambda a,\lambda a)\text{ and }S_+^{}(a,a)=(\lambda ^1a,\lambda ^1a)$$
where $`\lambda =B^1(1A)=(1D)^1C`$. So the sideways matrices preserve the diagonal. This has the curious consequence that a linear switch which is a birack is also a biquandle in the sense of \[FJK\] .
For a negative crossing the actions are equal but with opposite orientation.
Assume for simplicity that we are dealing with a knot. The link case is similar and details can safely be left to the reader. We have a right $`R`$-module with a finite (square) presentation.
Following the orientation of the knot, label the semi-arcs with $`R`$-variables $`x_1,x_2,\mathrm{},x_{2n}`$. By an $`R`$-variable we mean a symbol standing in for any element of $`R`$.
At each crossing there is a relation of the form
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}x_i\\ x_j\end{array}\right)=\left(\begin{array}{c}x_{j+1}\\ x_{i+1}\end{array}\right)\text{ or }\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}x_i\\ x_j\end{array}\right)=\left(\begin{array}{c}x_{j1}\\ x_{i1}\end{array}\right)$$
depending on whether the crossing is positive or negative. As is the usual custom, indices are taken modulo $`2n`$.
The relations can now be written in matrix form as $`M𝐱=\mathrm{𝟎}`$ where $`M`$ is a $`2n\times 2n`$ matrix and $`𝐱=(x_1,x_2\mathrm{},x_{2n})^T`$. The non-zero entries in each row of the matrix are $`A,B,1`$ or $`C,D,1`$.
Let $`=(S,D)`$ be the module defined by these relations. We now show that the modules defined by diagrams representing the same virtual link are isomorphic. We do this by showing that a single Reidemeister move defines an isomorphism. The proof has the same structure as the proof, say, that the Alexander module of a classical link is an invariant as in \[Alex\] but we give the details because of the care needed due to non-commutativity.
Any module defined by a presentation of the form $`M𝐱=\mathrm{𝟎}`$ is invariant under the following moves and their inverses applied to the matrix $`M`$.
1. permutations of rows and columns,
2. multiplying any row on the left or any column on the right by a unit,
3. adding a left multiple of a row to another row or a right multiple of a column to another column,
4. changing $`M`$ to $`\left(\begin{array}{cc}x& 𝐮\\ \mathrm{𝟎}& M\end{array}\right)`$ where $`x`$ is a unit, $`𝐮`$ is any row vector and $`\mathrm{𝟎}`$ is a zero column vector,
5. repeating a row.
The operations $`i.`$ above for $`i=1,2,3`$ are equivalent to multiplying $`M`$ on the right or left by an elementary matrix of type $`i`$.
Recall that the matrix $`S`$ can be written as a product of elementary matrices
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)=\left(\begin{array}{cc}A& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ C& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1A^1\end{array}\right)\left(\begin{array}{cc}1& A^1B\\ 0& 1\end{array}\right).$$
Now consider the module $``$ defined above. Clearly the presentation is unaltered by any of the basic moves which involve the virtual crossing. So we look to see the changes induced by the classical Reidemeister moves and check that the presentation matrix $`M`$ is only changed by the above 5 moves. Assume
$$M=\left(\begin{array}{cccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}\end{array}\right)$$
Firstly, consider a Reidemeister move of the first kind.
This introduces (or deletes) two new equal generators $`x_{n+1}=x_{n+2}`$. Because $`S_+^{}`$ and $`S_{}^+`$ preserve the diagonal, (the biquandle condition, see \[FJK\] ) the output ($`x_{n+3}`$) is the same as the input ($`x_n`$). The generator $`x_{n+1}`$ is equal to $`\lambda ^1x_n`$ where $`\lambda =B^1(1A)`$.
So up to reordering of the columns the relation matrix is changed by
$$M\left(\begin{array}{cccccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}& 0& 0\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}& 0& 0\\ 0& \mathrm{}& 0& 0& 1& 1\\ 0& \mathrm{}& 0& 1& \lambda & 0\end{array}\right)\left(\begin{array}{cccccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}& 0& 0\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}& 0& 0\\ 0& \mathrm{}& 0& 0& 0& 1\\ 0& \mathrm{}& 0& 1& \lambda & 0\end{array}\right)$$
Since $`\lambda `$ is a unit this does not alter the module.
There are other possible inversions and mirror images of the above which can be dealt with in a similar fashion.
Secondly, consider a Reidemeister move of the second kind.
Again the outputs are unchanged from the inputs $`x_{n1},x_n`$ because of the relation $`S^1S=1`$.
Two new generators $`x_{n+1}`$ and $`x_{n+2}`$ are introduced (or deleted). They are related by the equations
$$x_{n1}=Ax_{n+1}+Bx_{n+2}\text{ and }x_n=Cx_{n+1}+Dx_{n+2}.$$
This has the following effect on the relation matrix.
$$M\left(\begin{array}{cccccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}& 0& 0\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}& 0& 0\\ 0& \mathrm{}& 0& 1& A& B\\ 0& \mathrm{}& 1& 0& C& D\end{array}\right).$$
Since $`S`$ is a product of elementary matrices this does not alter the module.
The other possible inversions and mirror images of the above can be dealt with in a similar fashion but it is worth looking at the case where the two arcs run in opposite directions. The right outputs are unchanged from the left inputs by the relation $`S_{}^+(S^1)_{}^+=1`$.
The changes to the relation matrix are given by
$$M\left(\begin{array}{cccccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}& 0& 0\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}& 0& 0\\ 0& \mathrm{}& 1& A& B& 0\\ 0& \mathrm{}& 0& C& D& 1\end{array}\right)\left(\begin{array}{cccccc}m_{11}& \mathrm{}& m_{1n1}& m_{1n}& 0& 0\\ m_{21}& \mathrm{}& m_{2n1}& m_{2n}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ m_{n1}& \mathrm{}& m_{nn1}& m_{nn}& 0& 0\\ 0& \mathrm{}& 0& 0& 1& 0\\ 0& \mathrm{}& 0& 0& 0& 1\end{array}\right)$$
Using the fact that $`B`$ is a unit. This doesn’t change the module.
Finally, consider a Reidemeister move of the third kind.
The outputs $`x_{i+2},x_{j+2},x_{k+2}`$ are unaltered by the Reidemeister move because of the Yang-Baxter equations. The inner generators $`x_{i+1},x_{j+1},x_{k+1}`$ are related to the inputs $`x_i,x_j,x_k`$ by the following matrix
$$\left(\begin{array}{ccc}C& DA& DB\\ 0& C& D\\ 0& A& B\end{array}\right)$$
and the inner generators $`x_{i+1}^{},x_{j+1}^{},x_{k+1}^{}`$ are related to $`x_i,x_j,x_k`$ by the following matrix
$$\left(\begin{array}{ccc}C& D& 0\\ A& B& 0\\ AC& AD& B\end{array}\right).$$
Both are the product of elementary matrices and the proof is finished. $`\mathrm{}`$
9 Determinant Invariants
9.1 The Determinant $`\mathrm{\Delta }_0`$
Given a module with a square presentation the obvious invariant of the module is the determinant, if it can be defined. This will be the case if the ring is represented by matrices with commuting entries, for example the ring of generalised quaternions. In this case if $`d`$ denotes the determinant and $`M\text{x}=0`$ is the presentation let $`\mathrm{\Delta }_0=d(M)`$. Since the module depends on the switch $`S`$ we illustrate this dependency by $`\mathrm{\Delta }_0=\mathrm{\Delta }_0(S)`$.
A close look at how the presentation of the module changes under the Reidemeister moves shows that $`\mathrm{\Delta }_0`$ is invariant up to multiplication by $`d(B)`$ or $`d(C)`$. Typically $`d(B)`$ is denoted by the variable $`t`$ and $`d(C)`$ is $`t^1`$. If we take the switch to be $`E_1`$ ($`E_2`$) defined in section 5 then $`\mathrm{\Delta }_0`$ is a polynomial $`p_1`$ ($`p_2`$) in the four variables $`x,y,z,t`$ ($`a,b,c,t`$). We can normalise these polynomial so that as a polynomial in $`t`$ it has a non-zero constant term and only positive powers of $`t`$.
Let us illustrate the previous discussion by calculating invariants for the virtual trefoil as shown in the figure.
If we label as indicated then the module has a presentation with 4 generators $`a,b,c,d`$ and relations $`c=Ab+Ba,a=Ac+Bd,b=Cc+Dd,d=Cb+Da`$. Restricting to the $`E_1`$ case gives the polynomial $`p_1`$ equal to
$$\frac{64+4x^3tz^3+4x^3t^3z^3+128t^264xtzx^4t^2z^464xt^3z+8x^2t^2z^24x^2z^2+64t^44x^2t^4z^2}{16}$$
For the $`E_2`$ case we get the polynomial $`p_2`$ equal to
$$\frac{\begin{array}{c}4+16a40tba^34t^4a^4+16t^4a^324t^4a^2+16t^4a36tb^3a^2+72a^2tb+4b^2+13a^2t^4b^2\\ +40tb^3a+4t^4b^24t^424a^212b^2a+13a^2b^2+16a^312t^4b^2a+16tb6t^4b^2a^3\\ +t^4b^2a^456atb16tb^3+b^2a^46b^2a^3+8ba^4t2b^3a^4t+14b^3a^3t4a^48t^2a^4+32t^2a^3\\ 48t^2a^2+32t^2a8t^2+a^4t^2b^42a^4t^2b^28a^3t^2b^48t^2b^2+16t^2b^4+16t^3b16t^3b^3\\ 26a^2t^2b^2+12a^3t^2b^2+24a^2t^2b^4+8a^4t^3b40a^3t^3b+72a^2t^3b+24at^2b^232at^2b^4\\ 56at^3b2t^3b^3a^4+14t^3b^3a^336t^3b^3a^2+40t^3b^3a\end{array}}{t^2\left(a1\right)^4}$$
Note that the fundamental quandle (and hence group) as defined by the Wirtinger presentation is trivial.
The following virtual knot is interesting in having a trivial Jones-polynomial as well as a trivial fundamental rack.
In this case if $`S`$ is the Alexander switch then
$$\mathrm{\Delta }_0=(B1)(C^2(B+1)C(B+1)(B^1+1)B).$$
9.2 The Determinant $`\mathrm{\Delta }_1`$
For many knots and links, including the classical, the determinant $`\mathrm{\Delta }_0`$ is zero.
For example, as we have seen earlier any switch $`S`$ with entries in the ring $`R`$ defines a representation of the virtual braid group $`VB_n`$ into the group of invertible $`n\times n`$ matrices with entries in $`R`$ by sending the standard generator $`\sigma _i`$ to $`S_i=(id)^{i1}\times S\times (id)^{ni1}`$ and the generator $`\tau _i`$ to $`T_i=(id)^{i1}\times T\times (id)^{ni1}`$. This representation is denoted by $`\rho =\rho (S,n)`$. For classical braids this representation is equivalent to the Burau representation and so we would expect the closure of a classical braid to have $`\mathrm{\Delta }_0`$ zero. We now confirm this by looking at the fixed points of $`S_i`$ both on the left and right.
Lemma 9.2 Let $`P=A^1B^1A`$ and $`Q=B^1(1A)`$. Then
$$(P^{n1},\mathrm{},P,1)S_i=(P^{n1},\mathrm{},P,1)()$$
and
$$S_i(1,Q,\mathrm{},Q^{n1})^T=(1,Q,\mathrm{},Q^{n1})^T().$$
Proof We need only check that
$$(P,1)S=(P,1)\text{ and }S(1,Q)^T=(1,Q)^T.$$
$`\mathrm{}`$
Therefore the following lemma gives a necessary condition for the knot or link to be classical.
Theorem 9.3 For all classical knots and links $`\mathrm{\Delta }_0=0`$.
Proof Since $`\mathrm{\Delta }_0`$ is an invariant of the module we can assume that the diagram from which it is defined is the closure of a braid. However from $``$ (or $``$) there is a linear relationship amongst the rows (columns), so $`\mathrm{\Delta }_0=0`$. $`\mathrm{}`$
The Kishino knots $`K_1,K_2`$ and $`K_3`$ are illustrated below.
All are ways of forming the connected sum of two unknots. $`K_1`$ and $`K_2`$ are mirror images and $`K_3`$ is amphichæral. Both have trivial racks and Jones polynomial. The invariant $`\mathrm{\Delta }_0`$ is zero in all three cases.
It is clear that we need an invariant for these cases. Let $`M`$ be the $`n\times n`$ presentation matrix. Let $`M_1,M_2,\mathrm{},M_{n^2}`$ be the submatrices obtained by deleting a row and a column. Let $`d_1,d_2,\mathrm{},d_{n^2}`$ be the determinants. These all lie in a ring of polynomials with coefficients in a field. Therefore the determinants have an hcf, call it $`\mathrm{\Delta }_1`$, which is well defined up to multiplication by a unit. Now look at what happens to this construction under a Reidemeister move. The hcf, $`\mathrm{\Delta }_1`$, is multiplied by $`d`$ of a unit.
Returning to the Kishino knots, a calculation with the Alexander switch shows that for $`K_1`$, $`\mathrm{\Delta }_1`$ is $`1+BCB`$ and for $`K_2`$, $`\mathrm{\Delta }_1`$ is $`1+CCB`$. Since these are neither units nor associates in the ring, $`K_1,K_2`$ are non-trivial and non amphichæral.
On the other hand for $`K_3`$, $`\mathrm{\Delta }_1`$ is 1. We will show shortly that $`K_3`$ is non-trivial by using the Budapest switch augmented by $`t`$. Then $`\mathrm{\Delta }_1=2+5t^2+2t^4`$ see \[BuF\] .
For a classical knot or link the invariant $`\mathrm{\Delta }_1`$ is not just the Alexander polynomial in disguised form but is independant of the deleted rows or columns chosen, up to multiplication by a power of $`t`$.
Theorem 9.4 Let $`𝒟`$ be a diagram of a classical knot or link. If $`M`$ is the presentation matrix associated with $`𝒟`$. Let $`\mathrm{\Delta }_1=d(M_{ij})`$ where $`M_{ij}`$ is obtained from $`M`$ by deleting the $`i`$th row and the $`j`$th column. Then $`\mathrm{\Delta }_1`$ is independant of $`i,j`$ up to multiplication by a power of $`t`$.
Proof Assume initially that $`𝒟`$ is the closure of a braid.
Write
$$M=\left(\begin{array}{cccc}C_1& C_2& \mathrm{}& C_n\end{array}\right)$$
in terms of its columns and let
$$M_{ij}=\left(\begin{array}{cccc}C_1^0& C_2^0& \mathrm{}& C_n^0\end{array}\right)$$
where each column has its $`i`$th component removed and $`C_j^0`$ does not appear in the list. From $``$,
$$C_j^0=C_1^0Q^{(1j)}\mathrm{}C_n^0Q^{nj}$$
and $`C_j^0`$ does not appear on the right hand side of the equation.
So by column operations which do not change the value of the determinant we can change any column to $`C_j^0`$. Now note that the value of the determinant is unchanged by interchanging two columns. A similar argument works for the rows.
A general diagram is obtained from $`𝒟`$ by a sequence of Reidemeister moves. A glance at the change of $`M`$ under the Reidemeister moves shows that $`\mathrm{\Delta }_1`$ is invariant up to multiplication by a power of $`t`$ . $`\mathrm{}`$
Let us now return to $`K_3`$. This is the closure of the braid
$$\tau _2(\sigma _1\sigma _2\sigma _1)\tau _2(\sigma _1\sigma _2\sigma _1)^1$$
Suppose the representation of this as a $`3\times 3`$ matrix, using the Budapest switch augmented by $`t`$, is $`M`$. Then the representation matrix of the module is $`Mid`$. The nine codimension 1 subdeterminants are
$$p=2t^2+5+2t^2,\text{ (4 times) }q=(2+2t^2)p(t^1),\text{ (twice) }q(t^1),p^2$$
This not only shows that $`K_3`$ is non-trivial but that it cannot be classical by 9.4.
10 Epilogue
Most calculations in this paper are done with Maple. At Andy Bartholomew’s website at http://www.layer8.co.uk/maths/ it is possible to download a C-program which calculates the invariants.
It is extremely unlikely that there are no non-trivial virtual knots for which these methods fail to distinguish it from the trivial knot. For example if the braid $`\beta \overline{K}`$ (see section 2), then the closure of $`\beta `$ possibly provides an infinite set of examples. However, to prove that an infinite set exists would require different methods.
References
\[Alex\] J. W. Alexander, Topological Invariants of Knots and Links, Trans. American Math. Soc. 30(1928) pp 255-306
\[As\] Helmer Aslaksen, Quaternionic Determinants, Math. Intel. Vol 18 no. 3 (1996)
\[BF\] A. Bartholomew and Roger Fenn. Quaternionic Invariants of Virtual Knots and Links, preprint. Preprint vailable from http://www.maths.sussex.ac.uk////Staff/RAF/Maths/Current/Andy/
\[BuF\] S. Budden and Roger Fenn. The equation
$$[b,(a1)(a,b)]=0$$
and virtual knots and links, Fund Math 184 (2004) pp 19-29.
\[FJK\] R. Fenn, M. Jordan, L. Kauffman, Biquandles and Virtual Links, Topology and its Applications, 145 (2004) 157-175
\[CS\] J. S. Carter, D. Jelsovsky, Seiichi Kamada, Laurel Langford, Masahico Saito. Quandle Cohomology and State-sum Invariants of Knotted Curves and Surfaces http://arxiv.org/abs/math.GT/9903135
\[De\] P. Dehornoy, Non Commutative Versions of the Burau Representation, C.R.Acad. Roy. Sci. Canada; 17-1; (1995) pp 53-58
\[Dr\] V. Drinfeld, On some Unsolved Problems in Quantum Group Theory, Quantum Groups, Lectures Notes in Maths. 1510, Springer 1-8 (1990)
\[ESS\] P. Etingof, T. Schedler and A. Soloviev, Set-Theoretic Solutions to the Quantum Yang-Baxter Equations. Duke Math Journal 100 no. 2 169-209 (1999)
\[EGS\] P. Etingof, R. Guralnik and A. Soloviev, Indecomposable set-theoretical solutions to the Quantum Yang-Baxter Equation on a set with prime number of elements. J of Algebra 242 (2001) 709-719
\[FR\] R. Fenn, C. Rourke. Racks and Links in Codimension Two. JKTR, No. 4, 343-406 (1992).
\[FRS1\] R. Fenn, C. Rourke B. Sanderson. An Introduction to Species and the Rack Space. M. E. Bozhuyu (ed.) Topics in Knot Theory, Kluwer Academic, pp 33-55.(1993)
\[FRS2\] James Bundles (with C.Rourke and B.Sanderson). Proceedings of the LMS (3) 89 (2004) 217-240
\[J\] D. Joyce A classifying invariant of knots, the knot quandle. J. Pure Appl. Algebra 23, 37-65 (1982).
\[KK\] N. Kamada and S. Kamada, Abstract Link Diagrams and Virtual Knots, JKTR 9, 93-106 (2000).
\[KK2\] N. Kamada and S. Kamada, Braid Presentations of Virtual and Welded Knots, preprint GT/0008092
\[Ku\] G. Kuperberg, What is a Virtual Link? Algebraic and Geometric Topology 587-591 (2003)
\[K\] L.Kauffman. Virtual Knot Theory, European J. Comb. Vol 20, 663-690, (1999)
\[KS\] T. Kishino and S. Satoh, A note on classical polynomials, preprint (2001)
\[L\] T. Y. Lam. The Algebraic Theory of Quadratic Forms, Benjamin (1973)
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# Complex Relativity: Gravity and Electromagnetic Fields
## 1 Introduction
This work, related to the General Relativity and electromagnetic field, was inspired by some recent results in the asymmetric continuum theories including the spin motions; therefore, we shall, first of all, quote the numerous attempts to extend the General Relativity to include the spin motions. The first one, that in the Cartan works was influenced by work by Cosserat brothers in which a moment stress tensor is included in a generalized continuum. A gradual development of the Einstein-Cartan Theory (ECT) started by works of Sciama , Kibble and Trautman ; for a review see: . Kopczy ski has proved that in the ECT the cosmological solutions become free from the singularities leading to the modified Friedmann equation supplemented with the conservation laws for mass and spin .
In the XX century, we observed an enormous development of the continuum theories: the micropolar and micromorphic theories were developed basing on the Cosserat brothers’ work (for a review see: ); the relations joining the theories of a continuum containing defects (dislocation and disclination densities) with the Riemannian curvature and torsion were considered by Bilby, Bullough and Smith , Kondo and followed by Hollander , Ben-Abraham and many other authors; (for a review see: ); the thermal stresses were found to have the same form as that related to dislocation field and on this basis the thermal effects were included in the continuum with a Riemannian curvature by Kroner , Teisseyre and Stojanovic et al. ; for a review see: ).
Recently the continuum theories have been generalized to the asymmetic form in which an additional constitutive law for the antisymmetric part of stresses, replacing the stress moments, joins the spin motion with a new constant, rotation rigidity modulus, to account for the rotation of the point-grains and propagation of the spin elastic waves. Such waves can exist when the elastic bonds related to the rotation motions of particles are postulated. When , in such a continuum, the material bonds for the displacement neglecting originated deformations, there remain only the rotation fields of spin and twist types; the respective equations appear to have exactly the form of electromagnetic equations. On this basis, the degenerated continuum theory (in which there exist only the spin and twist axial motions) has been considered in the last papers by Teisseyre and Teisseyre, Białecki and Górski .
Influenced by these results, we introduce in this paper the natural potentials, as defined in the way suitable for constructing the first order distortions in the metric tensor. A transition from antisymmetric tensors to the symmetric ones helped us to define these natural electromagnetic potentials; their form fits a system of the Dirac matrices and this representation leads to distortions of the metric tensor.
The definitions helped us to propose a generalization of the General Relativity; the new theory - the Complex Relativity - includes, beside gravitation, the electromagnetic equations in a first order approximation.
## 2 Analogies to Asymmetric Continuum
We consider the physical rotation fields which can be related to the curvature deformations of the complex space. However, we shall mention that an inspiration for this idea has its source in considering the rotation and twist motion in the asymmetric elastic continuum; of course, such motions in an elastic continuum are bounded to some constitutive relations describing the elastic bonds. This is main difference in comparison with motion in the space.
In our approach a homogeneous elastic continuum with the rotation nuclei - of spin and twist type - is supplemented, beside the classical ideal elasticity constitutive law for the symmetric strain-stress relation, by the relation between the rotation and asymmetric stresses; such stresses appear when including in a medium the rotation nuclei. By this supplementary constitutive law for the anti-symmetric fields, we can evade an influence of the Hook law, which, when used as the unique law in the ideal elasticity, rules out an existence of rotation waves. Thus, it comes out that the rotation vibrations can propagate and are not attenuated, unlike as the elastic waves in the ideal elastic continuum.
The twist motion differs from the pure rotation; formaly, it is a motion composed of the rotation and the mirror reflection. It presents the simultaneously occurring opposite rotation motions like shear oscillations (some analogy, in a world of the linear displacements, presents a thermal expansion/compression motion differing from the simple displacements), but it can be also related to the axial motions like those of the polarization type, or, when assuming a posibility of material-space curvature, to a bending of the 3D space, (by analogy to the situation of a flat jellyfish with the bending motions (pulsating motions) leading from 2D form to 3D one).
Our approach to the asymmetry of fields follows from the antisymmetric stresses introduced by and related to the internal friction caused by the grain motions under friction forces. Note that in the asymmetric continuum also the related asymmetric incompatibility tensors split into symmetric and anti-symmetric parts.
In our former papers we have also analyzed the theory of asymmetric continuum with defect distribution (with the dislocation and disclination densities and the densities of rotation nuclei). Special consideration was paid to rotation and twist motions related to the definition of the twist-bend tensor.
The dislocation - stress relations and the equations of motion for symmetric and asymmetric parts of stresses were derived.The obtained relations for elastic fields, given by difference of the total and self-fields, can be split into the selfparts prevailing on the fracture plane and the total parts describing seismic radiation field in a surrounding space. Some applications were shortly discussed.
Finally, we shall note that a more complex deformation field, like that with the dislocation and disclination densities, leads, when applying the material coordinate system, to description of a deformed state in terms of the Riemannian geometry, or even non-Riemannian one . This remark applies also to thermal deformation field .
## 3 Natural electromagnetic potentials
We have tested many variants of definition for the EM potentials, and finally we propose the following one. We introduce the 3D vector potentials: $`\stackrel{~}{A}_s`$, $`\widehat{A}_s`$, and charge-current potentials $`\phi `$, $`\psi _s`$ instead of the standard 4D vector potential $`A_\mu `$. We call them the natural potentials and assume they fulfill the following equations:
$`B_k`$ $`=`$ $`ϵ_{kbs}\stackrel{~}{A}_{s,b}\stackrel{~}{A}_{s,s}=0`$ (1)
$`E_k+\phi _{,k}`$ $`=`$ $`ϵ_{kbs}\widehat{A}_{s,b}\widehat{A}_{s,s}=0`$ (2)
$`{\displaystyle \frac{4\pi }{c}}J_k`$ $`=`$ $`{\displaystyle \frac{1}{c}}\dot{\phi }_{,k}+{\displaystyle \frac{1}{c}}ϵ_{kbs}\psi _{s,b}`$ (3)
$`\phi _{,kk}`$ $`=`$ $`4\pi \rho ,`$ (4)
where $`k,b,s\{1,2,3\}`$. An index *after* a comma denotes differentiation and the summation convention for repeated indices is used. These new natural potentials yield
$`ϵ_{kbs}\stackrel{~}{A}_{s,b}{\displaystyle \frac{1}{c}}{\displaystyle \frac{}{t}}\widehat{A}_k`$ $`=`$ $`{\displaystyle \frac{1}{c}}\psi _k\stackrel{~}{A}_{b,b}=0`$ (5)
$`ϵ_{kbs}\widehat{A}_{s,b}+{\displaystyle \frac{1}{c}}{\displaystyle \frac{}{t}}\stackrel{~}{A}_k`$ $`=`$ $`0\widehat{A}_{b,b}=0.`$ (6)
When applying to (5)-(6) the operator $`ϵ_{ndk}\frac{}{x_d}`$, we arrive at the Maxwell equations:
$`ϵ_{ndk}B_{k,d}{\displaystyle \frac{1}{c}}{\displaystyle \frac{}{t}}E_n`$ $`=`$ $`{\displaystyle \frac{4\pi }{c}}J_kB_{k,k}=0`$ (7)
$`ϵ_{ndk}E_{k,d}+{\displaystyle \frac{1}{c}}{\displaystyle \frac{}{t}}B_n`$ $`=`$ $`0E_{k,k}=4\pi \rho `$ (8)
Now, we construct the complex antisymmetric tensor $`A_{\alpha \beta }`$ ($`\alpha ,\beta \{1,2,3,4\}`$) for potentials $`\stackrel{~}{A}`$, $`\widehat{A}`$. We define $`A_{\alpha \beta }=\stackrel{~}{A}_{\alpha \beta }+\text{i}\widehat{A}_{\alpha \beta }`$ by
$$A_{\alpha \beta }=\left[\begin{array}{cccc}0\hfill & \stackrel{~}{A}_3\hfill & \stackrel{~}{A}_2\hfill & \text{i}\widehat{A}_1\hfill \\ \stackrel{~}{A}_3\hfill & 0\hfill & \stackrel{~}{A}_1\hfill & \text{i}\widehat{A}_2\hfill \\ \stackrel{~}{A}_2\hfill & \stackrel{~}{A}_1\hfill & 0\hfill & \text{i}\widehat{A}_3\hfill \\ \text{i}\widehat{A}_1\hfill & \text{i}\widehat{A}_2\hfill & \text{i}\widehat{A}_3\hfill & 0\hfill \end{array}\right]+\text{i}\left[\begin{array}{cccc}0\hfill & \widehat{A}_3\hfill & \widehat{A}_2\hfill & \text{i}\stackrel{~}{A}_1\hfill \\ \widehat{A}_3\hfill & 0\hfill & \widehat{A}_1\hfill & \text{i}\stackrel{~}{A}_2\hfill \\ \widehat{A}_2\hfill & \widehat{A}_1\hfill & 0\hfill & \text{i}\stackrel{~}{A}_3\hfill \\ \text{i}\stackrel{~}{A}_1\hfill & \text{i}\stackrel{~}{A}_2\hfill & \text{i}\stackrel{~}{A}_3\hfill & 0\hfill \end{array}\right]$$
(9)
or
$$A_{\alpha \beta }=\left[\begin{array}{cccc}0\hfill & \overline{A}_3\hfill & \overline{A}_2\hfill & \overline{A}_1\hfill \\ \overline{A}_3\hfill & 0\hfill & \overline{A}_1\hfill & \overline{A}_2\hfill \\ \overline{A}_2\hfill & \overline{A}_1\hfill & 0\hfill & \overline{A}_3\hfill \\ \overline{A}_1\hfill & \overline{A}_2\hfill & \overline{A}_3\hfill & 0\hfill \end{array}\right]$$
(10)
where $`\overline{A}_k=\stackrel{~}{A}_k+i\widehat{A}_k`$. The form of equation (9) follows the fact, that these potentials are constructed in a similar manner as the tensor $`f_{\alpha \beta }`$ of the EM field is constructed from the EM vector fields $`B_s`$ and $`E_s`$.
According to (5)-(6) the tensor $`A_{\alpha \beta }`$ fulfils the condition
$$A_{\alpha \beta |\beta }=\frac{1}{c}\psi _\alpha \psi _\alpha =\{\psi _k\text{ , }0\}k=1,2,3.$$
(11)
Comparing this relation with equations for the potentials $`\stackrel{~}{A}_s`$, $`\widehat{A}_s`$ we obtain
$`\mathrm{}\stackrel{~}{A}_n`$ $`=`$ $`{\displaystyle \frac{1}{c}}ϵ_{ndk}\psi _{k,d}\stackrel{~}{A}_{b,b}=0`$ (12)
$`\mathrm{}\widehat{A}_n`$ $`=`$ $`{\displaystyle \frac{1}{c^2}}{\displaystyle \frac{}{t}}\psi _n\widehat{A}_{b,b}=0`$ (13)
and when defining
$$\frac{1}{c}ϵ_{ndk}\psi _{k,d}=\stackrel{~}{J}_n\text{}\frac{1}{c^2}\frac{}{t}\psi _k=\widehat{J}_k$$
(14)
we arrive at
$$\mathrm{}\stackrel{~}{A}_n=\stackrel{~}{J}_n\text{ , }\mathrm{}\widehat{A}_n=\widehat{J}_n$$
(15)
Then, we define tensor $`J_{\alpha \beta }`$ :
$$J_{\alpha \beta }=\left[\begin{array}{cccc}0\hfill & \overline{J}_3\hfill & \overline{J}_2\hfill & \overline{J}_1\hfill \\ \overline{J}_3\hfill & 0\hfill & \overline{J}_1\hfill & \overline{J}_2\hfill \\ \overline{J}_2\hfill & \overline{J}_1\hfill & 0\hfill & \overline{J}_3\hfill \\ \overline{J}_1\hfill & \overline{J}_2\hfill & \overline{J}_3\hfill & 0\hfill \end{array}\right]$$
(16)
where $`\overline{J}_n=\stackrel{~}{J}_n+`$i$`\widehat{J}_n`$ and we arrive to relation
$$\mathrm{}A_{\alpha \beta }=J_{\alpha \beta }.$$
(17)
Applying the operator $`ϵ_{sbn}\frac{}{x_b}`$ we get, with the help of (12)-(13),
$`\mathrm{}B_s`$ $`=`$ $`{\displaystyle \frac{1}{c}}\psi _{s,dd}\text{ }={\displaystyle \frac{4\pi }{c}}ϵ_{spk}J_{k,p}`$ (18)
$`\mathrm{}E_s`$ $`=`$ $`{\displaystyle \frac{1}{c^2}}ϵ_{sbn}{\displaystyle \frac{}{t}}\psi _{n,b}\text{ }\mathrm{}\phi _{,s}={\displaystyle \frac{4\pi }{c^2}}\dot{J}_s+4\pi \rho _{,s}.`$ (19)
## 4 Natural EM Potential <br>— the Symmetric Tensor
To include the EM field into the Riemannian or non-Riemannian geometry we search for a way how to build a metric tensor which in the first order approximation could describe the gravity and EM fields.
Owing to the fact that for the six potentials $`\stackrel{~}{A}_s`$ , $`\widehat{A}_s`$ we can introduce the two additional conditions, we define the other set of four potentials:
$$\stackrel{~}{N}_n=\{\stackrel{~}{N}_1\text{ , }\stackrel{~}{N}_{2\text{ }},0\}\text{and}\widehat{N}_n=\{\widehat{N}_1,\widehat{N}_2,0\}.$$
(20)
With these conditions we introduce the natural symmetric tensor of potentials:
$`N_{\alpha \beta }`$ $`=`$ $`\stackrel{~}{N}_{\alpha \beta }+\text{i}\widehat{N}_{\alpha \beta }\text{ }N_s=\text{ }\stackrel{~}{N}_s+\text{i}\widehat{N}_s`$ (21)
$`N_{\alpha \beta }`$ $`=`$ $`\left[\begin{array}{cccc}0& 0& N_2& N_1\\ 0& 0& N_1& N_2\\ N_2& N_1& 0& 0\\ N_1& N_2& 0& 0\end{array}\right]`$ (26)
### The bridge relations.
We demand that first bridge relation is satisfied:
$$ϵ_{kbs}\stackrel{~}{N}_{s,b}=ϵ_{kbs}\stackrel{~}{A}_{s,b},ϵ_{kbs}\widehat{N}_{s,b}=ϵ_{kbs}\widehat{A}_{s,b}$$
(27)
and with the conditions
$$N_s=ϵ_{skn}\mathrm{\Omega }_{n,k}\text{ , }\mathrm{\Omega }_1=\mathrm{\Theta }_{,1}\text{ , }\mathrm{\Omega }_2=\mathrm{\Theta }_{,2}$$
(28)
we have
$$\text{ }N_{1,1}+N_{2,2}=0\text{ , }N_3=0.$$
(29)
The tensor $`N_{\alpha \beta }`$ (i.e. symmetric EM potential) has changed the signs of the terms below the diagonal of the matrix when comparing to tensor $`A_{\alpha \beta }`$ (i.e. antisymmetric EM potential).
To connect the antisymmetric potentials with the symmetric ones we shall introduce also a new definition for current potential
$$\eta _\alpha =\{\eta _1,\eta _2,0,0\},\eta _{s,s}=0.$$
(30)
We demand the validity of the second bridge relations:
$$ϵ_{kbS}\eta _{S,b}=ϵ_{kbn}\psi _{n,b};S=1,2.$$
(31)
Remark. A class of these EM potentials which can be transformed into the form described above will be called *two-component* EM potentials. We claim such potentials have proper form to contribute to metric tensor as described in section 5.
### The field relations.
Having these bridge relations we return to symmetric tensors. Due to this fact we get the equations for new symmetric potentials, similar to (12)-(13),
$$\mathrm{}\stackrel{~}{N}_K=\frac{1}{c}ϵ_{KdN}\eta _{N,d}\text{ , }\widehat{N}_{S,S}=0\text{ }\text{}\text{;}\text{ }\mathrm{}\widehat{N}_K=\frac{1}{c^2}\frac{}{t}\eta _K\text{ , }\stackrel{~}{N}_{S,S}=0.$$
(32)
Instead of relation (14) for currents $`\stackrel{~}{J}_n`$
and$`\widehat{J}_k`$ we shall introduce a new definition of currents $`\stackrel{~}{Y}_N`$
and $`\widehat{Y}_N`$ and tensor $`Y_{\alpha \beta }:`$
$$Y_{\alpha \beta }=\left[\begin{array}{cccc}0& 0& \stackrel{~}{Y}_2& \stackrel{~}{Y}_1\\ 0& 0& \stackrel{~}{Y}_1& \stackrel{~}{Y}_2\\ \stackrel{~}{Y}_2& \stackrel{~}{Y}_1& 0& 0\\ \stackrel{~}{Y}_1& \stackrel{~}{Y}_2& 0& 0\end{array}\right]+\text{i}\left[\begin{array}{cccc}0& 0& \widehat{Y}_2& \widehat{Y}_1\\ 0& 0& \widehat{Y}_1& \widehat{Y}_2\\ \widehat{Y}_2& \widehat{Y}_1& 0& 0\\ \widehat{Y}_1& \widehat{Y}_2& 0& 0\end{array}\right]$$
(33)
$$\frac{1}{c}\eta _{1,3}=\stackrel{~}{Y}_2\text{ }\text{}\text{}\frac{1}{c}\eta _{2,3}=\stackrel{~}{Y}_1\text{ ;}\frac{1}{c^2}\frac{}{t}\eta _N=\widehat{Y}_N$$
(34)
where similarly to (15) we have:
$$\mathrm{}\stackrel{~}{N}_N=\stackrel{~}{Y}_N,\mathrm{}\widehat{N}_N=\widehat{Y}_N$$
(35)
The following relations for the potentials are to be noted
$$N_{\alpha \gamma ,\gamma }=\frac{1}{c}\eta _\alpha \text{ , }\mathrm{}N_{\alpha \beta }=Y_{\alpha \beta }$$
(36)
We will show, further on, that all these symmetric matrices are the tensors related to the Dirac tensors.
### Complex metric tensor and EM natural potentials.
The natural tensor of potentials (26) can be presented as follows:
$$N_{\alpha \beta }=\stackrel{~}{N}_1ϵ^1+\stackrel{~}{N}_2ϵ^2+\text{i}\widehat{N}_1ϵ^1+\text{i}\widehat{N}_2ϵ^2=N_1ϵ^1+N_2ϵ^2$$
(37)
where
$$ϵ^1=\left[\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right],ϵ^2=\left[\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right]$$
(38)
The matrices $`ϵ^\nu `$ fulfil the conditions for the Dirac’s matrices:
$$ϵ^\alpha ϵ^\beta +ϵ^\beta ϵ^\alpha =2\eta ^{\alpha \beta }$$
(39)
with
$$ϵ^3=\left[\begin{array}{cccc}0& 0& 0& \text{i}\\ 0& 0& \text{i}& 0\\ 0& \text{i}& 0& 0\\ \text{i}& 0& 0& 0\end{array}\right],ϵ^4=\left[\begin{array}{cccc}\text{i}& 0& 0& 0\\ 0& \text{i}& 0& 0\\ 0& 0& \text{i}& 0\\ 0& 0& 0& \text{i}\end{array}\right].$$
(40)
However, when disturbing these matrices in the way indicated below
$$\gamma ^1=(1+N_1)ϵ^1\text{ , }\gamma ^2=(1+N_2)ϵ^2\text{ , }\gamma ^3=ϵ^3\text{ , }\gamma ^4=ϵ^4$$
(41)
we obtain the relation for the Dirac’s matrices in non-Euclidean space:
$$\gamma ^\alpha \gamma ^\beta +\gamma ^\beta \gamma ^\alpha =2g^{\alpha \beta }$$
(42)
This relation justifies our approach in which we have assumed that the natural potentials (26) can be used as the disturbances to metric tensor in any reference system.
These complex disturbances to metric tensor can be combined with the gravity disturbances $`h_{\alpha \beta }^G`$ . However let us first discuss another possible fields constructed similarly as above.
Comment: Field V. Thus, our next question is related to what would happen when one disturbes the $`ϵ^3`$ and $`ϵ^4`$ matrices. To this aim, we will consider the complex tensor potential $`V_{\alpha \beta }`$ defined as disturbances to $`ϵ^3,ϵ^4`$ similarly as $`N_{\alpha \beta }`$ disturbs $`ϵ^1,ϵ^2`$ in (37):
$$V_{\alpha \beta }=\text{i}(V_3ϵ^3+V_4ϵ^4).$$
(43)
The disturbances related to $`V_3`$ are not symmetric hence we put $`V_3=0.`$
With the undisturbed matrices $`\overline{\gamma }^1=ϵ^1`$ , $`\overline{\gamma }^2=ϵ^2`$ , $`\overline{\gamma }^3=ϵ^3`$, we define the disturbed matrics $`\overline{\gamma }^4:`$
$$\text{ }\overline{\gamma }^4=\text{i}(1+V_4)ϵ^4$$
(44)
Further on, we obtain the related metric tensor:
$$\overline{g}_{\mu \nu }=\frac{1}{2}(\overline{\gamma }_\mu \overline{\gamma }_\nu +\overline{\gamma }_\nu \overline{\gamma }_\mu )$$
(45)
Both relations correspond to those given in (41) and (42) and, with Einstein equations, can lead us to $`V^{\alpha \beta },_\beta =Z^\alpha `$, $`V^{\alpha \beta },_\gamma ^\gamma =U^{\alpha \beta }`$ where $`Z^\alpha `$, $`U^{\alpha \beta }`$ would represent some source fields.
Further on, we neglect the field $`V^{\alpha \beta }`$.
Remark. We shall note that, instead of the 4D presentation, it was also possible to present our relations in the 2D forms with the tensor
$$N_{AB}=\left[\begin{array}{cc}N_2\hfill & N_1\hfill \\ N_1\hfill & N_2\hfill \end{array}\right]$$
and with the help of the Pauli 2D tensors.
## 5 Complex metric and curvature tensors
Our formulation of the Complex Relativity is based on the symmetric form of perturbations introduced into the metric tensor.
The respective metric tensor can be constructed when considering the coordinates $`X^\alpha `$ ($`X^4=`$i$`ct`$) for which the numerical values, ascribed to the points $`x^\alpha `$ of the Minkovski space, do not change under deformation (for example see ):
\- before deformation we can write
$$\text{d}s^2=\eta _{\alpha \beta }\text{d}x^\alpha \text{d}x^\beta $$
(46)
\- while after deformation the first order disturbances into the metric tensor are:
$$\text{d}S^2=g_{\alpha \beta }\text{d}X^\alpha \text{d}X^\beta \text{ , }h_{\alpha \beta }g_{\alpha \beta }\eta _{\alpha \beta }$$
(47)
For weak fields this formalism can lead us to equations for electromagnetic and gravitation fields, when assuming that the disturbances $`h_{\alpha \beta }`$ can be related to the following fields:
* $`h_{\alpha \beta }^G`$ — the classical disturbances related to gravity
* $`h_{\alpha \beta }^N=N_{\alpha \beta }`$ — the disturbances related to EM field.
We propose the metric tensor of the following form
$$g_{\alpha \beta }=\eta _{\alpha \beta }+h_{\alpha \beta }^N+h_{\alpha \beta }^G.$$
(48)
Introducing such first order disturbances, we can consider the complex Riemann $`R_{\alpha \beta }`$ and Einstein tensors
$$G_{\alpha \beta }=R_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }$$
(49)
and the related basic field equation:
$$G_{\alpha \beta }=0$$
(50)
Now we consider (50) up to the first order terms.
### Contributions from the non-diagonal terms $`h_{\alpha \beta }^N`$.
For non-diagonal terms we perturb the metric by the potentials $`h_{\alpha \beta }^N=N_{\alpha \beta }`$ fulfiling the relation (32).
Considering the first order contributions to the Einstein tensor we obtain (an index *after* $`|`$ denotes differentiation)
$$G_{\alpha \beta }^N(\frac{1}{2}N_{\mu \beta }|{}_{\alpha }{}^{\mu }+\frac{1}{2}N_{\mu \alpha }|{}_{\beta }{}^{\mu })+\frac{1}{2}N_{\alpha \beta }|{}_{\nu }{}^{\nu }=\frac{1}{2}\text{ }Y_{\alpha \beta }$$
(51)
where accroding to (36) we have $`\frac{1}{2}N_{\mu \beta }\left|{}_{\alpha }{}^{\mu }+\frac{1}{2}N_{\mu \alpha }\right|{}_{\beta }{}^{\mu }=\frac{1}{2c}(\eta _{\alpha ,\beta }+\eta _{\beta ,\alpha }).`$
### Contributions from the h$`{}_{}{}^{G}{}_{\alpha \beta }{}^{}`$ terms
The classical General Relativity relations
$$G_{\alpha \beta }^G=R_{\alpha \beta }1/2g_{\alpha \beta }=\frac{8\pi G}{c^4}T_{\alpha \beta }$$
(52)
with the matter-energy tensor
$$T_{\alpha \beta }=c^2\rho _0\upsilon _\alpha \upsilon _\beta ,\upsilon _\alpha =\{\upsilon _s/c,i\}$$
(53)
can be included into a new complex form, with the assumption that the elements $`h_{k4}^G`$ present the imaginary values:
$$\frac{1}{2}h_{\beta \alpha }^G\left|{}_{\nu }{}^{\nu }+\frac{1}{2}\eta ^{\nu \mu }h_{\mu \nu }^G\right|{}_{\alpha \beta }{}^{}\frac{8\pi G}{c^4}T_{\alpha \beta }$$
(54)
### First order contributions to $`G_{\alpha \beta }`$.
For the disturbances given by equation (48) we obtain
$$G_{\alpha \beta }\frac{1}{2}N_{\alpha \beta }\left|{}_{\nu }{}^{\nu }\frac{1}{2}h_{\beta \alpha }^G\right|{}_{\nu }{}^{\nu }\frac{1}{2}\eta ^{\nu \mu }h_{\mu \nu }^G|_{\alpha \beta }$$
(55)
or if we define
$$\overline{h}^{\mu \nu }=h^{\mu \nu }\frac{1}{2}\eta ^{\mu \nu }h_\alpha ^\alpha $$
(56)
we arrive at
$$G_{\alpha \beta }\frac{1}{2}N_{\alpha \beta }\left|{}_{\nu }{}^{\nu }\frac{1}{2}\overline{h}_{\beta \alpha }^G\right|{}_{\nu }{}^{\nu }\frac{1}{c}\eta _{(\alpha ,\beta )}\frac{1}{2}Y_{\alpha \beta }\frac{8\pi G}{c^4}T_{\alpha \beta }.$$
(57)
The matrics $`Y_{\alpha \beta }`$ is built also with the help of the $`ϵ`$ matrices
$$Y_{\alpha \beta }=\stackrel{~}{Y}_1ϵ^1+\stackrel{~}{Y}_2ϵ^2+\text{i}\widehat{Y}_1ϵ^1+\text{i}\widehat{Y}_2ϵ^2$$
(58)
while the tensor $`\frac{1}{2c}(\eta _{\alpha ,\beta }+\eta _{\beta ,\alpha })`$ with the condition $`\eta _1=\theta _{,1}`$and $`\eta _2=\theta _{,2}`$ ($`\eta _{1,2}+\eta _{2,1}=0`$ and $`\eta _{1,2}`$ $`\eta _{2,1}=2\theta _{,12}`$ ) becomes
$$\frac{1}{2c}(\eta _{\alpha ,\beta }+\eta _{\beta ,\alpha })=\left[\begin{array}{cccc}0& 0& \frac{1}{2}\eta _{1,3}& \frac{1}{2}\eta _{1,4}\\ 0& 0& \frac{1}{2}\eta _{2,3}& \frac{1}{2}\eta _{2,4}\\ \frac{1}{2}\eta _{1,3}& \frac{1}{2}\eta _{2,3}& 0& 0\\ \frac{1}{2}\eta _{1,4}& \frac{1}{2}\eta _{2,4}& 0& 0\end{array}\right]$$
Finaly, we obtain
$$G_{\alpha \beta }\frac{1}{c}\eta _{(\alpha ,\beta )}\frac{1}{2}Y_{\alpha \beta }\frac{8\pi G}{c^4}T_{\alpha \beta }=\left[\begin{array}{cccc}0& 0& \text{i}\frac{1}{2}\widehat{Y}^2& \frac{1}{2}\stackrel{~}{Y}_1\\ 0& 0& \text{i}\frac{1}{2}\widehat{Y}& \frac{1}{2}\stackrel{~}{Y}_2\\ \text{i}\frac{1}{2}\widehat{Y}^2& \text{i}\frac{1}{2}\widehat{Y}^1& 0& 0\\ \frac{1}{2}\stackrel{~}{Y}_1& \frac{1}{2}\stackrel{~}{Y}_2& 0& 0\end{array}\right]\frac{8\pi G}{c^4}T_{\alpha \beta }$$
Notice:
It is worth noticing a complimentary structures of the amplitude related tensors $`T_{\alpha \beta }`$ and $`E_{\alpha \beta }=(\frac{1}{2c}(\eta _{\alpha ,\beta }+\eta _{\beta ,\alpha })\frac{1}{2}Y_{\alpha \beta }):`$
$$T_{\alpha \beta }=\left[\begin{array}{cccc}\text{Re}& \text{Re}& \text{Re}& \text{Im}\\ \text{Re}& \text{Re}& \text{Re}& \text{Im}\\ \text{Re}& \text{Re}& \text{Re}& \text{Im}\\ \text{Im}& \text{Im}& \text{Im}& \text{Re}\end{array}\right]\text{ },\text{ }E_{\alpha \beta }=\left[\begin{array}{cccc}0& 0& \text{Im}& \text{Re}\\ 0& 0& \text{Im}& \text{Re}\\ \text{Im}& \text{Im}& 0& 0\\ \text{Re}& \text{Re}& 0& 0\end{array}\right]$$
(59)
which implies a possibility of separation of the EM and gravity fields on a linear level.
### Remark: Meaning of natural potentials.
Let us confine ourselves to the EM fields and let us consider the 3D curvilinear complex space (or 6D real space) with coordinates
$$\overline{X}_s=\stackrel{~}{X}_s+\text{i}\widehat{X}_s.$$
(60)
The potentials $`\overline{A}_s=\stackrel{~}{A}_s+`$i$`\widehat{A}_s`$ can be identified with such frames of the complex space $`\overline{X}_s=\stackrel{~}{X}_s+`$i$`\widehat{X}_s`$; same holds for the respecvtive tensors: $`A_{\alpha \beta }\overline{X}_{\alpha \beta }:`$
$$\overline{X}_s=\overline{A}_s=\stackrel{~}{A}_s+i\widehat{A}_s\text{ ; }\overline{X}_{\alpha \beta }\overline{A}_{\alpha \beta }$$
(61)
We assume that at each point of this 6D complex space continuum there can appear the independent spin and twist motions<sup>1</sup><sup>1</sup>1It can be realized as ”internal” motions of points in such ”grained” (micropolar) space based on the Planck length. and that complex space is combined to the point of such continuum
$$ϵ_{abs}\overline{X}_s=ϵ_{abs}\stackrel{~}{X}_s+\text{i}ϵ_{abs}\widehat{X}_s\text{ , }\overline{A}_{s,s}=\text{ }\overline{X}_{s,s}=\stackrel{~}{X}_{s,s}+\text{i}\widehat{X}_{s,s}=0.$$
(62)
## Conclusion and perspectives
In this initial work we have presented the Complex Relativity theory as inspired by recent progress in asymmeric continuum approach to material sciences. The goal of this work is to propose the common relativity framework for electromagnetic and gravity fields as close as possible to the classical General Relativity formulation. For another classical approach see .
We mention here possible ways of further development of our approach.
The relations for the symmetric tensors as expressed by the $`\gamma `$ -tensors remain valid in any reference system, but we might also return to the system with the antisymmetric potentials ($`\overline{A}^3`$ different from zero; see relations between natural potentials in both systems) in order to get notation comparable with the gravity part. We shall be also aware that in such a case we return to antisymmetric natural potential tensor; we would obtain
$$G_{\alpha \beta }^{ANTISYM}=\frac{1}{2c}\times $$
$$\times \left[\begin{array}{cccc}0\hfill & ϵ_{3dk}\psi _{k,d}\text{i}\frac{}{t}\psi _3\hfill & ϵ_{2dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _2\hfill & ϵ_{1dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _1\hfill \\ ϵ_{3dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _3\hfill & 0\hfill & ϵ_{1dk}\psi _{k,d}\text{i}\frac{}{t}\psi _1\hfill & ϵ_{2dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _2\hfill \\ ϵ_{2dk}\psi _{k,d}\text{i}\frac{}{t}\psi _2\hfill & ϵ_{1dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _1\hfill & 0\hfill & ϵ_{3dk}\psi _{k,d}+\text{i}\frac{}{t}\psi _3\hfill \\ ϵ_{1dk}\psi _{k,d}\text{i}\frac{}{t}\psi _1\hfill & ϵ_{2dk}\psi _{k,d}\text{i}\frac{}{t}\psi _2\hfill & ϵ_{3dk}\psi _{k,d}\text{i}\frac{}{t}\psi _3\hfill & 0\hfill \end{array}\right]$$
(63)
where $`G_{\alpha \beta }^{ANTISYM}`$ is no longer the complex Einstein tensor, but stands for the corresponding antisymmmetric expression.
In the above considerations we tried to preserve the symmetric property when constructing the Complex Relativity, however, it seems more natural to admit the possibility that the perturbations into metric and to Einstein tensor can be asymmetric or even antisymmetric.
These problems will be discussed in our next paper.
## Acknowledgments
The authors thank Andrzej Czechowski for useful remarks. This work is partially supported by Polish Ministry of Science and Information Society Technologies, project 2PO4D 060 28.
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# Exact results for the 1D interacting mixed Bose-Fermi gas
## Abstract
The exact solution of the $`1`$D interacting mixed Bose-Fermi gas is used to calculate ground-state properties both for finite systems and in the thermodynamic limit. The quasimomentum distribution, ground-state energy and generalized velocities are obtained as functions of the interaction strength both for polarized and non-polarized fermions. We do not observe any demixing instability of the system for repulsive interactions.
Quantum mixed Fermi gas, momentum distribution function, phase separation
The cooling and trapping of quantum Fermi gases of ultracold atoms poses a number of additional challenges to those faced for bosons BEC . A key point is that no more than one identical fermion can occupy a single state due to the Pauli exclusion principle. However, a mixed gas of fermions and bosons provides an effective means of cooling single-component fermions by thermal collisions with evaporatively cooled bosons Fermi-1 ; Fermi-2 ; Fermi-3 ; Fermi-4 , providing an avenue for investigating many-body quantum effects in degenerate Fermi gases. Another development is the use of Feshbach resonance, in which the energy of a bound state of two colliding atoms is magnetic field tuned to vary the scattering strength from $`\mathrm{}`$ to $`\mathrm{}`$, allowing the investigation of the crossover from BCS superfluidity to Bose-Einstein condensation BEC-F1 ; BEC-F2 ; BEC-F3 . These achievements in realizing quantum Fermi gases of ultracold atoms may also provide insights into other areas of physics, such as ultracold superstrings superstring .
Particular theoretical and experimental interest has been paid to the Fermi gas confined in 1D geometry Fermi-1D1 ; Fermi-1D2 ; BEC-BCS1 ; BEC-BCS2 ; BEC-BCS3 . Recent attention has turned to the 1D model of mixed bosons and polarized fermions mix-1 ; mix-2 ; mix-3 ; mix-4 , revealing quantum effects such as interaction-driven phase separation phase-s1 ; mix-3 ; mix-4 , bright solitons in degenerate Bose-Fermi mixtures kar04 and Luttinger liquid behaviour mix-1 ; mix-2 . In this model, only $`s`$-wave scattering for boson-boson and boson-fermion interactions is considered, with the boson-fermion collisions minimizing the Pauli-blocking effect in momentum space.
The theoretical studies of the mixed Bose-Fermi gas have focused on the case of bosons and single-component fermions mix-1 ; mix-2 ; mix-3 ; mix-4 . Our aim here is to investigate the $`1`$D model of a Bose-Fermi mixture on a line of length $`L`$ with periodic boundary conditions, with Hamiltonian $`H=_{i=1}^N\frac{\mathrm{}^2}{2m_i}\frac{^2}{x_i^2}+_{i<j}g_{i,j}\delta (x_ix_j)`$, where among the $`N`$ particles there are $`N_f=N_{}+N_{}`$ fermions and $`N_b`$ bosons. The pairwise $`\delta `$-interaction has strength $`g_{i,j}`$. The crucial observation from a theoretical point of view is that if all particles have equal masses and if the interaction strength between all particles is the same, the above model is exactly solvable by Bethe Ansatz (BA) L-Y . Although this restricts the applicability of the model, experiments with isotopes of atoms are expected to meet the BA conditions mix-4 . Then $`g_{i,j}=2\mathrm{}^2/(ma)`$, where $`c=2/a`$ is the inverse 1D scattering length of the confined particles. For convenience of notation, the energy is measured in units of $`\mathrm{}^2/(2m)`$, such that
$`H={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}+{\displaystyle \underset{i<j}{}}2c\delta (x_ix_j).`$ (1)
The model contains two special limiting cases: (i) the $`1`$D interacting Fermi gas with arbitrary polarization Yang-Gaudin ; Takahashi ; Wadati ; BBGN , and (ii) the $`1`$D interacting Bose gas LL63 . The mixed model was recently discussed for the case of single component (fully polarized) fermions mix-4 . Here we consider the more general case of two-component fermions with arbitrary polarization. We present analytical and numerical results for the ground state energy, quasimomentum density profile, and the spin and charge velocities. In the weak coupling limit, these quantities reveal the typical signatures of pure quantum gases with an additional weak interaction due to the mixture. In the other extreme, that is for strong repulsive interactions, the ground state properties resemble those of a single-component non-interacting Fermi gas.
The BA equations (BE), determining the quantum numbers of the $`N`$-particle system, are given by L-Y
$`\text{e}^{\mathrm{i}p_{\mathrm{}}L}`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{p_{\mathrm{}}\mathrm{\Lambda }_j+\text{i}c^{}}{p_{\mathrm{}}\mathrm{\Lambda }_j\text{i}c^{}}}`$ (2)
$`{\displaystyle \underset{\mathrm{}}{}}{\displaystyle \frac{\mathrm{\Lambda }_kp_{\mathrm{}}\text{i}c^{}}{\mathrm{\Lambda }_kp_{\mathrm{}}+\text{i}c^{}}}`$ $`=`$ $`{\displaystyle \underset{j,m}{}}{\displaystyle \frac{\mathrm{\Lambda }_k\mathrm{\Lambda }_j\text{i}c}{\mathrm{\Lambda }_k\mathrm{\Lambda }_j+\text{i}c}}{\displaystyle \frac{\mathrm{\Lambda }_kA_m+\text{i}c^{}}{\mathrm{\Lambda }_kA_m\text{i}c^{}}}`$ (3)
$`1`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{A_n\mathrm{\Lambda }_j\text{i}c^{}}{A_n\mathrm{\Lambda }_j+\text{i}c^{}}},`$ (4)
where $`c^{}=c/2`$ and $`j,k=1,\mathrm{},N_b+N_{}`$, $`\mathrm{}=1,\mathrm{},N`$, $`m,n=1,\mathrm{},N_b`$. Without loss of generality, we take $`N_{}<N_{}`$. Here the set of $`N+N_{}+2N_b`$ many quantum numbers is divided into three subsets: $`\{p_{\mathrm{}}\}_{1,\mathrm{},N}`$, $`\{\mathrm{\Lambda }_k\}_{1,\mathrm{},N_{}+N_b}`$, $`\{A_n\}_{1,\mathrm{},N_b}`$. It turns out that the energy eigenvalues $`E`$ are given by the members of the first set alone, namely $`E=_{\mathrm{}=1}^Np_{\mathrm{}}^2`$.
We begin with a finite system, for which the Bethe roots are found analytically in the weak coupling limit, thereby yielding the ground state energy. The second step is to carry out the thermodynamic limit (TL). In this limit, the Bethe roots are distributed smoothly over a certain interval of the real axis, giving rise to continuous densities. Integral equations for these densities have been obtained previously L-Y . On the one hand, these equations allow a comparison between the analytical results for weak coupling and finite systems with the weak coupling expansion in the TL. On the other hand, generalized velocities can be calculated within this framework. In the weak coupling limit, these reduce to the spin and charge velocities of a pure Fermi gas and the charge velocity of a pure bosonic system.
In considering the ground state for weak interaction, it is convenient to distinguish between unpaired $`p_j^{(\mathrm{u})}`$ ($`j=1,\mathrm{},N_{}N_{}`$), paired $`p_j^{(\mathrm{p})}`$ ($`j=1,\mathrm{},N_{}1,N_{}+2,\mathrm{},2N_{}`$) and bosonic $`p_j^{(\mathrm{b})}`$ ($`j=1,\mathrm{},N_b+2`$) quasimomenta note1 . Expanding Eqs. (2)-(4) to $`𝒪(c)`$, one obtains $`p_j^{(\mathrm{u})}=\pi (N_{}1+2j)/L+\delta _j^{(\mathrm{u})}`$, $`j=1,\mathrm{},(N_{}N_{})/2`$; $`p_j^{(\mathrm{u})}=\pi (2N_{}N_{}1+2j)/L+\delta _j^{(\mathrm{u})},j=(N_{}N_{})/2+1,\mathrm{},N_{}N_{}`$ and $`p_j^{(\mathrm{p})}=\pi (1N_{}+2j_+)+\delta _{j_+}^{(\mathrm{p})}\pm \sqrt{c/L}`$. Here $`j_+=j`$ if $`j`$ odd and $`j_+=j1`$ if $`j`$ even. The deviations $`\delta _j`$ from $`p_j^{(\mathrm{u},\mathrm{p})}`$ are linear in $`c`$, with
$`\delta _j^{(\mathrm{u})}`$ $`=`$ $`{\displaystyle \frac{c}{L}}\left[{\displaystyle \underset{k}{}}{\displaystyle \frac{1}{p_{j,0}^{(\mathrm{u})}p_{k,0}^{(\mathrm{p})}}}+{\displaystyle \frac{N_b+1}{p_{j,0}^{(\mathrm{u})}}}\right],`$ (5)
$`\delta _j^{(\mathrm{p})}`$ $`=`$ $`{\displaystyle \frac{c}{L}}\left[{\displaystyle \underset{kj}{}}{\displaystyle \frac{1}{p_{j,0}^{(\mathrm{p})}p_{k,0}^{(\mathrm{p})}}}+{\displaystyle \underset{k}{}}{\displaystyle \frac{1/2}{p_{j,0}^{(\mathrm{p})}p_{k,0}^{(\mathrm{u})}}}+{\displaystyle \frac{N_b+1}{p_{j,0}^{(\mathrm{p})}}}\right].`$ (6)
Here $`p_{j,0}^{(\mathrm{p},\mathrm{u})}`$ denotes the quasimomenta in the free particle limit as given above. The sums in (5), (6) over paired momenta $`p_{k,0}^{(\mathrm{p})}`$ count each pair only once. The bosonic momenta $`p_j^{(\mathrm{b})}`$ collapse at the origin for $`c=0`$. Their $`c`$-dependence is given by
$`p_j^{(\mathrm{b})}={\displaystyle \frac{2c}{L}}{\displaystyle \underset{kj}{}}{\displaystyle \frac{1}{p_j^{(\mathrm{b})}p_k^{(\mathrm{b})}}}.`$ (7)
According to (7), the $`p_j^{(\mathrm{b})}`$ are the roots of Hermite polynomials of degree $`N_b+2`$ gau71 ; BGM .
From the above equations we obtain the ground state energy, $`E=E^{(0)}+E^{(1)}`$, to leading order in $`c`$. The energy of the free particles, $`E^0`$, is given by the corresponding expression for a free Fermi gas with $`N_f`$ fermions BBGN . The first correction in $`c`$ is $`E^{(1)}=E_b^{(1)}+E_f^{(1)}+2\frac{c}{L}N_bN_f`$, where $`E_b^{(1)}`$ ($`E_f^{(1)}`$) is the linear order in $`c`$ for a pure Bose (Fermi) gas of $`N_b`$ bosons BGM ($`N_f`$ fermions BBGN ). The last term is due to the interactions between bosons and fermions. Thus
$`E`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{L}}\right)^2({\displaystyle \frac{1}{6}}N_{}(N_{}^21)+`$ (8)
$`{\displaystyle \frac{1}{12}}(N_{}N_{})(1+N_{}^2+N_f^2N_fN_{}))`$
$`+{\displaystyle \frac{2c}{L}}\left(N_{}N_{}+{\displaystyle \frac{1}{2}}N_b(N_b1)+N_bN_f\right).`$
This expression is valid for both weak repulsive and attractive interaction. Especially, the TL is well defined in the weakly attractive case, as opposed to the strongly attractive case Takahashi . Let us now carry out the TL, i.e., $`N_\alpha ,L\mathrm{}`$ where the densities $`n_\alpha =N_\alpha /L`$ are held constant, with $`\alpha =,,b`$. In this limit, Eqs. (5) and (6) give the distribution of quasimomenta per unit length,
$`\rho _u(p)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}+{\displaystyle \frac{c}{2\pi ^2}}\left({\displaystyle \frac{A}{p^2A^2}}+{\displaystyle \frac{B}{p^2}}\right),A<|p|<C`$ (9)
$`\rho _p(p)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{c}{2\pi ^2}}\left({\displaystyle \frac{A}{A^2p^2}}+{\displaystyle \frac{C}{C^2p^2}}{\displaystyle \frac{2B}{p^2}}\right),`$ (10)
$`\mathrm{\hspace{0.17em}2}\sqrt{cB/\pi }<|p|<C`$
$`\rho _b(p)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi c}}\left(4cB/\pi p^2\right)^{1/2},|p|<2\sqrt{cB/\pi },`$ (11)
where $`A=\pi N_{}/L`$, $`B=\pi N_b/L`$, $`C=\pi (N_{}N_{})/L`$.
The quasimomentum distribution function of the bosonic momenta (11) for $`c>0`$ is given by the semi-circular law behaviour as for a pure bosonic system LL63 ; gau71 ; BGM . For $`c<0`$, the $`p_j^{(\mathrm{b})}`$ in (7) are imaginary, so that (11) with imaginary $`k`$ yields the dark-soliton like distributions of the bosonic quasimomenta on the imaginary axis. These encode the binding energy of the bosons, a quantity potentially accessible by experiments. In the distribution of the fermionic quasimomenta, divergences occur near the cutoffs. The quasimomentum distribution functions calculated from the BE (2)-(4) are compared with the approximations (9)-(11) in Fig. 1.
In order to study the effect of arbitrary interaction in the TL, we use the linear integral equations derived by Lai and Yang L-Y which determine the densities of the quasimomenta. In the following, we restrict ourselves to two limiting cases, namely $`N_{}=N_{}`$ and $`N_{}=0`$. In the latter case, the fermions do not interact among themselves due to the Pauli principle, so that the interaction potential in (1) is only effective between fermions and bosons. The ground-state energy for this case has been calculated recently mix-4 . In the other limiting case, $`N_{}=N_{}`$, all the fermions interact with each other and with the bosons. Fig. 2 shows the ground-state energy per unit length $`E/L`$ for different densities $`n_b`$ and $`n_f`$ obtained from numerical solution of the integral equations. Also shown is a comparison between the analytic result (8) and the TL.
Up to this point, we have focussed on the density $`\rho =\rho _u+\rho _p+\rho _b`$ of the quasimomenta $`p_{\mathrm{}}^{(\mathrm{u},\mathrm{p},\mathrm{b})}`$. In an analogous fashion, one may introduce densities $`\sigma `$, $`\tau `$ of the roots $`\mathrm{\Lambda }_j,A_m`$. Using the dressed energy formalism Takahashi , one can calculate the corresponding dressed energies $`ϵ,\varphi ,\psi `$, which are the energies necessary to add or remove a root to or from the seas of $`p_{\mathrm{}},\mathrm{\Lambda }_j,A_m`$. The dressed energies give rise to generalized velocities, which determine the asymptotics of correlation functions frahm . In the case of pure bosons, there is only one species of BE numbers, associated with one dressed energy function, yielding the charge (or sound) velocity $`v_c^{(b)}`$ LL63 ; Takahashi . For pure fermions, the two sets of BE numbers are linked with the charge and spin velocities $`v_{c,s}^{(f)}`$. As has been proven by Haldane hal81a , these velocities coincide with those calculated in the harmonic-fluid (or bosonization) approach hal81b .
The situation is less clear for the mixture of bosons and fermions. In the BA approach, this corresponds to two seas of BE numbers $`p_{\mathrm{}},\mathrm{\Lambda }_j`$, and correspondingly two dressed energies and two velocities, which we call $`\stackrel{~}{v}_c^{(b,f)}`$. For $`c=0`$ we have $`\stackrel{~}{v}_c^{(f)}=2\pi n_f`$ and $`\stackrel{~}{v}_c^{(b)}=0`$. The dependence of $`\stackrel{~}{v}_c^{(b,f)}`$ on $`c`$ for different values of $`n_f,n_b`$ is shown in Fig. 3. The lowest $`c`$-order of $`\stackrel{~}{v}_c^{(b)}`$ is obtained from the Bethe-Ansatz as $`\stackrel{~}{v}_c^{(b)}=2n_b\left[c/n_b(1/2\pi )(c/n_b)^{3/2}\right]^{1/2}`$, which is the velocity of a non-interacting Bose gas for small $`c`$ LL63 .
We now compare our results to those of the hydrodynamic (HD) approach mix-1 . In this framework, the pure gases correspond to harmonic oscillators, such that the coupling between them leads to new normal modes $`v_\pm `$, with $`v_\pm (c=0)=\stackrel{~}{v}_c^{(f,b)}(c=0)`$. In the weak interaction limit, $`v_{}=2\sqrt{cn_b}`$, $`v_+=2\pi n_f+\frac{n_bc^2}{n_f^2\pi ^3}`$. For the sake of comparison, the latter result is also shown in Fig. 4. In the strongly repulsive limit, a demixing instability is predicted from the HD mix-1 and mean-field approaches mix-3 . We do not observe any instability for repulsive interaction for the integrable model, in agreement with L-Y ; mix-4 . The reason for the discrepancy is that to our present understanding, the HD approach for the mixture, especially the calculation of normal modes, is a low-energy theory. That is, it is expected to yield reliable results for small interaction strengths $`|c|1`$. Investigation of $`\stackrel{~}{v}_c^{(f,b)}`$ within the BA approach in the strongly repulsive limit yields
$`\stackrel{~}{v}_c^{(f)}`$ $`=`$ $`2\pi n\left[1{\displaystyle \frac{4n}{\pi c}}\left[{\displaystyle \frac{\pi n_b}{n}}+\mathrm{sin}\left({\displaystyle \frac{\pi n_b}{n}}\right)\right]\right],`$
$`\stackrel{~}{v}_c^{(b)}`$ $`=`$ $`{\displaystyle \frac{4\pi ^2n^2}{3c}}\mathrm{sin}\left({\displaystyle \frac{\pi n_b}{n}}\right),`$
where $`n=n_f+n_b`$.
It remains a very interesting question for future research if the normal modes of the field theory coincide with the generalized velocities from the BA for mixed Bose-Fermi systems. Furthermore, it should be carefully investigated whether or not the HD approach is applicable in the strongly interacting regime.
The generalized velocities can also be computed from the BA for interacting fermions ($`n_{}=n_{}`$). In this case, the three dressed energies $`ϵ,\varphi ,\psi `$ give rise to three velocities, which we call $`\stackrel{~}{v}_{c,s}^{(f)}`$ and $`\stackrel{~}{v}_c^{(b)}`$. As shown in Fig. 4, we observe that in the weak interaction limit, $`\stackrel{~}{v}_{c,s}^{(f)}=\pi n_f(1\pm c/(n_f\pi ^2))`$, which corresponds to the charge and spin velocities of a pure fermionic system BEC-BCS1 . In the strong interaction limit, $`\stackrel{~}{v}_c^{(b)}=\stackrel{~}{v}_s^{(f)}=0`$, whereas $`\stackrel{~}{v}_c^{(f)}=2\pi n`$. This again indicates the fermionic nature of the system in the strongly interacting limit.
In conclusion, we have investigated ground state properties of the 1D interacting Bose-Fermi model from its exact BA solution. We obtained results for the distribution of quasimomenta and the ground state energy both for weak attractive and repulsive interactions. We computed the generalized velocities and compared them to the normal modes obtained from the HD approach. In contrast with other approaches mix-1 ; mix-3 ; mori , we do not observe any instability or demixing in the system.
Acknowledgements. This work has been supported by the Australian Research Council and the German Science Foundation under grant number BO/2538.
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# The INT/WFC survey of the Monoceros Ring: Accretion origin or Galactic Anomaly?
## 1 Introduction
The formation and evolution of galaxies remains one of the big questions in astronomy. In the currently favoured $`\mathrm{\Lambda }`$CDM model, galaxies are built up over time via the accretion of smaller systems (e.g. Searle & Zinn, 1978; White & Rees, 1978; White, 1978; Abadi, Navarro, Steinmetz, & Eke, 2003a, b). The picture is not wholly satisfactory and some parts of the Milky Way may have formed in an initial collapse of baryonic material, somewhat akin to the model of galaxy formation proposed by Eggen, Lynden-Bell, & Sandage (1962).
One firm prediction of the $`\mathrm{\Lambda }`$CDM model is that this accretion of smaller systems should still be ongoing and that the Milky Way Halo should contain a large number of satellite systems. It has been suggested that, given the model, there are too few satellites actually within the Milky Way Halo (Klypin, Kravtsov, Valenzuela, & Prada, 1999), although the extent of this discrepancy has been a matter of some debate. Many attempts have been made to resolve this issue (e.g. Moore, 2001) by correcting the measured velocity dispersion of the local satellites to a velocity dispersion of the dark matter components, while more recently Susa & Umemura (2004) used numerical methods to highlight a scenario whereby these low mass systems are evaporated at the Epoch of Re-ionization. In the meantime, more and more surveys are probing the Halo of the Milky Way revealing structures which may point to the formation history of our galaxy (see Freeman & Bland-Hawthorn, 2002).
The tidal dismemberment of a dwarf galaxy as it falls through the Milky Way Halo is a slow process, with extensive streams of tidal debris existing for long periods of time (Ibata & Lewis, 1998; Johnston et al., 1999). While ancient remnants have been identified, via phase-space analysis, in our own Galactic neighbourhood (Helmi et al., 1999; Brook, Kawata, Gibson, & Flynn, 2003), more extensive surveys of the Galactic Halo, such as the Spaghetti Survey (Morrison et al., 2000) and utilizing 2MASS (Majewski, Skrutskie, Weinberg, & Ostheimer, 2003) have concluded that there is only a single, major on-going accretion event, that of the Sagittarius dwarf galaxy (Ibata, Gilmore, & Irwin, 1994). While this accretion event is adding mass to the Galactic Halo and provides an important probe of the shape of the dark matter potential (e.g. Ibata et al., 2001; Majewski et al., 2004), the lack of other other major accretion events is somewhat disconcerting given the predictions from $`\mathrm{\Lambda }`$CDM.
As will be discussed in detail in $`\mathrm{\S }`$2, the recently discovered Monoceros Ring (MRi) can be interpreted as an additional on-going accretion event within the Milky Way. Investigating the density and extent of this structure is important when trying to fully understand the impact this type of event is having on the evolution of our galaxy both in the past and into the future. If the MRi is instead the outermost edge of the Milky Way, mapping the outer reaches of the Disk will provide insight into the Milky Way’s past. To this end, we have used the Isaac Newton Telescope Wide Field Camera to continue a campaign to detect this stellar population around the Galactic plane mapping out the extent of the MRi. This paper presents the results of a wide-field camera survey of the extensive stellar population thought to represent a continuation of the previous work in this field. The layout of the paper is as follows; $`\mathrm{\S }`$2 summarises the extant knowledge of the MRi and the associated population of stars while $`\mathrm{\S }`$3 describes the observational procedure and data reduction. $`\mathrm{\S }`$4 describes the analysis procedure, and the conclusions of this study are presented in $`\mathrm{\S }`$5.
## 2 The Monoceros Ring
The first sign of a new structure in the Galaxy, the Monoceros Ring<sup>1</sup><sup>1</sup>1The stream of stars at the edge of the Milky Way has received several names, including The One Ring (Ibata et al., 2003), the Monoceros Stream (Yanny et al., 2003) and the Galactic Anti-Centre Stellar Stream – GASS – (Majewski, Skrutskie, Weinberg, & Ostheimer, 2003). For the sake of simplicity, throughout this paper this stellar population is referred to as the Monoceros Ring (MRi), although this name does not imply we currently know the exact nature of this stellar population., came from a study of colour selected F-stars drawn from commissioning data from the Sloan Digital Sky Survey \[SDSS, Newberg et al. (2002)\]. Obtained in a narrow strip around the celestial equator these revealed two major stellar overdensities in the Galactic Halo, consistent with an intersection of the streams of tidal debris torn from the Sagittarius Dwarf galaxy (Ibata, Irwin, Lewis, & Stolte, 2001). Accompanying these, however, was an additional significant overdensity in the direction of the Galactic anti-centre, interpreted as being another tidal stream about the Galaxy, just past the edge of the stellar Disk, at a Galactocentric distance of $`18`$kpc and with a thickness of $`2`$kpc, covering over $`50^{}`$ of sky within $`|b|<30^{}`$ of the Galactic equator.
Given the results of the study of the SDSS, Ibata et al. (2003) searched for the signature of this stellar population in fields obtained for the Isaac Newton Telescope Wide Field Camera (INT/WFC) survey (see McMahon et al., 2001). Identifying the stream as a distinct population in colour-magnitude diagrams (CMDs), this study found the stars to be extensively distributed; over 100 of the sky within $`|\mathrm{b}|<30^\mathrm{o}`$ of the Galactic equator, suggesting that the stream completely rings the Galaxy. Main sequence fitting reveals that the Galactocentric distance to the stream varies between $`15`$kpc and $`20`$kpc, with an apparent scale-height of $`0.75`$kpc. Ibata et al. (2003) suggested that this stream represents debris of an accreting dwarf galaxy, although pointed out that the extant data did not rule out alternatives such as an outer spiral arm or unknown warp/flare of the Galactic Disk. In fact, they favoured the interpretation of it being a perturbation of the disc, possibly the result of ancient warps.
Simultaneously, Yanny et al. (2003) presented an analysis of a larger SDSS catalog of Halo stars obtaining a number of radial velocities in several fields over the stream, finding a velocity dispersion of $`2530`$km/s, inconsistent with any known Galactic component. Furthermore, these kinematics indicate that the stream possesses a prograde orbit about the Milky Way with a circular velocity of $`215\pm 25`$km/s \[Note: this revised value was presented in Yanny et al. (2004)\]. Metallicity estimates from these spectra indicate the stars are relatively metal poor $`([\frac{Fe}{H}]1.6)`$. Yanny et al. (2003) also concluded the stream represents a cannibalized dwarf galaxy undulating about the edge of the Galactic Disk.
Several other studies have focused upon this stream population; using the 2-Micron All Sky Survey (2MASS) Rocha-Pinto et al. (2003) identified M-giant stars beyond the Disk of the Milky Way, consistent with the population detected in the optical. Again, at a Galactocentric distance (R<sub>GC</sub>) of $`18\pm 2`$kpc, this arc of stars covers $`170^{}`$, with a higher metallicity than previously determined $`([\frac{Fe}{H}]=0.4\pm 0.3)`$. Crane et al. (2003) extended this study, obtaining stellar velocities of M-giants selected from 2MASS. Confirming a velocity dispersion of $`20`$km/s, this study concluded the stream orbits the Milky Way in a prograde fashion on an orbit with very little eccentricity. While this may seem problematic for accretion models preferring elliptical orbits, numerical simulations by Helmi et al. (2003) suggest that tidal streams in the plane of the Milky Way can possess quite circular orbits; these numerical studies, however, also suggest that the accretion event must be young, less than $`1`$Gyr since its first perigalactic passage, or any coherent structure would have dissolved. Finally, Frinchaboy et al. (2004) noted five globular clusters which are apparently associated with the stellar stream, as well as $`15`$ outer, old stellar clusters, bolstering the argument that it represents an accreting dwarf galaxy.
Martin et al. (2004a) also employed the 2MASS catalogue to identify M-giants beyond the Disk of the Milky Way. By considering the projected density of these stars, this study uncovered north-south anisotropies in the density of M-stars, with arcs above and below the plane of the Galaxy. Significantly, Martin et al. (2004a) identified a strong overdensity of these stars, $`2300`$ of them in roughly 20 of the sky at $`(l,b)(240,8)^{}`$, the constellation of Canis Major (CMa). This is a similar number of M-stars to that seen in the Sagittarius dwarf galaxy and Martin et al. (2004a) similarly interpret this population of stars as a dwarf galaxy. It is probably the progenitor, with an original mass of $`10^810^9\mathrm{M}_{}`$, of the extensive stream of stars around the edge of the Disk of the Milky Way. Given its mass, if CMa does represent an equatorial accretion event, it will, when fully dissolved, increase the mass of the Thick Disk by $`10\%`$.
Additional evidence for this accretion interpretation of the CMa dwarf comes from Bellazzini et al. (2004) who found the signature of the main body of the dwarf in the background to several Galactic open clusters, with the analysis of this population suggesting that it is somewhat metal rich with an age of $`27`$ Gyrs, although a blue plume indicates a younger population. This study also identified several globular and open clusters associated with the dwarf. Recently, Forbes, Strader, & Brodie (2004) have shown that the globular clusters associated with CMa possess their own age-metallicity relationship which is distinct from that of the Galactic population. Furthermore, the clusters are smaller than expected, if drawn from the Galactic population, strengthening the interpretation that they are of non-Galactic origin.
Momany et al. (2004), by re-analysing the 2MASS data claim the overabundance of M-giants in CMa is simply a signature of the warp in the Milky Way. In response, Martin et al. (2004b) used 2dF<sup>2</sup><sup>2</sup>2The 2dF instrument on the Anglo-Australian Telescope is a multi-fibre spectrograph with the ability to obtain spectra for 400 stars in a single pointing, inside a two-degree field. data of the Canis Major region to highlight the velocity disparities between the Milky Way Disk stars and the M-giant overdensity stars. This has been repeated more recently by Martin et al. (2005) suggesting the dwarf galaxy CMa can now be tentatively accepted as a real entity despite the current debate. Connecting CMa to the MRi is more problematic due to gaps in the detection of the ring and poor kinematic knowledge. A recently completed 2dF kinematic survey (Martin et al., 2005) of the Canis Major region may yield results in this area although currently any general conclusions linking it with the MRi are speculative.
More information about Halo substructure has been recently discovered as Rocha-Pinto et al. (2004) have identified another structure in Triangulum-Andromedae (TriAnd) which extends much further out than the MRi. Currently it is not known whether the two are related, although interestingly the detection of the MRi by Ibata et al. (2003) resides at the edge of this new structure in TriAnd suggesting a connection.
The latest evidence supporting the interpretation of the Canis Major Dwarf galaxy has been presented by Martínez-Delgado et al. (2004), showing the colour-magnitude diagram of a region (0.5 x 0.5) centred on (l,b) = (240,-8), the overdensity proposed by Martin et al. (2004b). An upper limit to the Heliocentric distance of the galaxy is found to be 5.3$`\pm `$0.2 to 8.1$`\pm `$0.4 kpc; this is comparable with previous estimates. By measuring the surface brightness and total luminosity of the dwarf galaxy an estimate of the mass range is found to be 1.0$`\times 10^8<`$ M $`<`$ 5.5$`\times `$10<sup>8</sup>. The tightness of the main sequence they found contradicts the claims that the Canis Major overdensity is the signature of the Galactic warp. Peñarrubia et al. (2004) have completed extensive modelling of the MRi and while not conclusively establishing a connection between Canis Major and the MRi, their results are highly suggestive of such a link. Peñarrubia et al. (2004) using all of the current information known about the MRi, undertook thousands of simulations, prograde and retrograde, in an attempt to find a model which best fitted the data. While some retrograde models were marginally acceptable, their preferred model was a prograde orbit which includes multiple wraps of the Milky Way. This allows for the scenario that the tidal stream has both near and far components. $`\mathrm{\S }`$4.3 discusses a field in which this phenomenon appears. If the dwarf galaxy has completed more than one orbit then the stream must be a much older structure than previously assumed \[cf. Helmi et al. (2003)\], raising the possibility that the newly discovered TriAnd structure (Rocha-Pinto et al., 2004) could be the distant arm of a multiply wrapped tidal stream.
The model of Peñarrubia et al. (2004) and the continuing work of those studying this new Galactic feature are slowly piecing together this structure, although currently there is no real coverage around the entire Galactic plane, these observations hopefully extend the knowledge of the MRi and it’s potential progenitor the Canis Major Dwarf.
## 3 Observations and Reduction
The data was obtained on the Isaac Newton Telescope Wide Field Camera (INT/WFC) at Roque de Los Muchachos in La Palma, Canary Islands. Mounted at the telescope prime focus, this covers 0.29 square degrees field per pointing, imaging onto four 4k$`\times `$2k CCDs. These possess 13.5$`\mu `$m pixels, corresponding to a pixel scale of 0.33 arcsec per pixel.
Nine survey regions were chosen, roughly equally spaced between l = $`61^o150^{}`$. To aid in determining the location of the fields, the model of Martin et al. (2004a) (See Figure 1) was used to predict where we might expect the spread of debris from an equatorial accretion. For each Galactic longitude, two regions symmetrically placed above and below the Galactic plane, were imaged. Each target region is a composite of overlapping fields, the number of which depended on time available to observe in each region (details of the observations are summarised in Table 1 and presented graphically in Figure 1), but typically with a total area of $`2`$ square degrees. With the representative integration times, the limiting magnitudes are on average for these observations $`V_{}`$ 23.3, $`i_{}`$ 22.2, $`g_{}`$ 23.8 and $`r_{}`$ 22.8. Archival data of (123,-19), taken from the M31 survey (Ibata et al., 2001; Ferguson et al., 2002), is used as the opposing field to (118,+16).
De-biassing and trimming, vignetting correction, astrometry and photometry were all undertaken with the CASU data reduction pipeline (Irwin & Lewis, 2001). The flat fielding employed a master twilight flat generated over the entire observing run. Each star was individually extinction corrected using the dust\_getval.c program supplied by Schlegel<sup>3</sup><sup>3</sup>3http://www.astro.princeton.edu/$``$schlegel/dust/data/data.html which interpolates the extinction using the maps presented by Schlegel, Finkbeiner, & Davis (1998). Observing several standard fields per night allows the calibration of the photometry to be determined, as described by (Irwin & Lewis, 2001), deriving the CCD zero-points. These are consistent to within a few percent on photometric nights. The data reduction pipeline produces a catalogue of all images in each colour-band. Rejecting non-stellar images, the catalogues can be cross-correlated and the colour for individual stars can be determined. Near the limiting magnitude, however, galaxies can appear stellar and so galactic contamination is expected for the faintest sources.
## 4 Analysis
### 4.1 Detecting the Monoceros Ring
As a first step each field was investigated by eye, looking for the presence of the MRi main sequence. In previous studies (Ibata et al., 2003) the MRi structure was observed as a coherent stellar population superimposed upon an overall Galactic CMD. Hence, we compare the observed and synthetic CMDs, as well as the North and South fields, to search for such a signal.
The synthetic CMDs were generated from the online models of the Milky Way by Robin, Reylé, Derrière & Picaud (2003)<sup>4</sup><sup>4</sup>4http://www.obs-besancon.fr/www/modele. To avoid any artificial cuts in the data, the modelled fields are extended out to 50 kpc (Heliocentric distance, R<sub>HC</sub>) and include the entire magnitude range available to the model (-99,99) for each passband. The extinction parameter is set to zero, all ages/populations and luminosity classes are included. This creates a complete picture of the field of interest and allows any later cuts on magnitude to be made at our discretion. The ($`V`$,$`i`$) figures have been created using data straight from the model while the ($`g`$,$`r`$) figures converted the model ($`V`$,$`i`$) magnitudes via the INT/WFC colour corrections<sup>5</sup><sup>5</sup>5http://www.ast.cam.ac.uk/$``$wfcsur/technical/photom/colours.
To understand the CMDs, requires that we know what Galactic features are present in the data. By deconstructing the model, the influence of each component of the galaxy in the CMD is presented. This can be seen in Figure 2, the top panels show the (61,-15) field from the model being split into its various Galactic components. For $`g_{}`$ and $`r_{}`$ CMDs, this plot illustrates where the various components of the Galaxy lie on the CMD. The lower panels are from the model field (75,+15), and show how a ($`V_{}`$,$`i_{}`$) CMD can be similarly deconstructed. Analysing Figure 2, shows that each component of the Milky Way occupies different regions of the CMD. In the following subsections the equations defining the separate Galactic structures as used in the synthetic galaxy model<sup>6</sup><sup>6</sup>6 Each of the following equations is cited and presented as per the referenced papers. A slight modification has been made to the Thick Disk equation to account for a typographical error in the paper. are presented.
#### 4.1.1 Thin Disk density profile
The youngest disc<sup>7</sup><sup>7</sup>7Robin & Creze (1986): $`R`$ and $`z`$ are the Galactocentric cylindrical coordinates (pc), with $`c`$ as the axis ratio for the ellipsoidal components. A is a constant defined by Table 1 of Robin & Creze (1986).
$$a^2=R^2+\frac{z^2}{c^2}$$
$$\rho (R,z)=A\left\{exp\left(\left(\frac{a}{K_\text{+}}\right)^2\right)exp\left(\left(\frac{a}{K_\text{-}}\right)^2\right)\right\}$$
(1)
with $`K_+=5000pc`$, $`K_{}=3000pc`$ and $`c=0.014`$
Old discs:
$$\begin{array}{c}\hfill \rho (R,z)=A\{exp((0.5^2+\left(\frac{a^2}{K_\text{+}^2}\right))^{\frac{1}{2}})\\ \hfill exp((0.5^2+\left(\frac{a^2}{K_\text{-}^2}\right))^{\frac{1}{2}})\}\end{array}$$
(2)
with $`K_+=2226pc`$ and $`K_{}=494.4pc`$ The Thin Disk is quite tightly constrained in the $`g`$,$`r`$ CMD but is quite broad in the $`V`$,$`i`$ CMD, it turns redward at $``$19<sup>th</sup> magnitude in the $`g`$,$`r`$ CMD, while for the $`V`$,$`i`$ CMD it is not as well defined and blurs down to $``$20<sup>th</sup> magnitude.
#### 4.1.2 Thick Disk density profile
$$\rho (R,z)\{\begin{array}{cc}exp\left(\frac{RR_{}}{h_R}\right)\left(1\frac{1/h_z}{x_l(2+x_l/h_z)}z^2\right)\hfill & \text{if }z\text{ }\text{ }x_l\hfill \\ exp\left(\frac{RR_{}}{h_R}\right)exp(\frac{z}{h_z})\hfill & \text{if }z\text{ }>\text{ }x_l\hfill \end{array}$$
(3)
The Thick Disk<sup>8</sup><sup>8</sup>8Reylé & Robin (2001) equations have three parameters defining the density along the $`z`$ axis: $`h_z`$, the scale height, $`\rho _{}`$, the local density and $`x_l`$ the distance above the plane where the density law becomes exponential. This third parameter is fixed by continuity of $`\rho `$($`z`$) and its derivative. It varies with the choice of scale height and local density following the potential. These variations are listed below<sup>9</sup><sup>9</sup>9Robin et al. (1996),
$$\begin{array}{c}\hfill x_l=1358.61.35nh+2.335^4nh^2\\ \hfill nr(8.1775^1+5.817E(3)nh)\end{array}$$
(4)
where the scale height density is
$$nh=\frac{h_z}{1pc}$$
(5)
and the local density is
$$nr=\frac{\rho _{}}{1.22E(3)stars.pc^3}$$
(6)
From Table 1 of Robin et al. (1996) the scale height is $`h_z`$ = 760 $`\pm `$ 50 pc, the local density $`\rho _{}`$ = 5.6 $`\pm `$ 1% and the scale length $`h_R`$ = 2.8 kpc. The Thick Disk in the $`g`$,$`r`$ CMD is well constrained at the ends of the magnitude range (keeping in mind that these figures do not take into account any incompleteness) and has a redward trend around 21<sup>st</sup> magnitude. The Thick Disk in $`V`$,$`i`$ is much more diffuse and the redward trend extends to $``$22<sup>nd</sup> magnitude.
#### 4.1.3 Halo density profile
$$\rho (R,z)=\rho _{}(R^2+\frac{z^2}{\epsilon ^2})^{\frac{n}{2}}$$
(7)
In the Halo <sup>10</sup><sup>10</sup>10Robin et al. (2000) or spheroidal equations, $`\rho _{}`$ is the local density, $`n`$ is the power law index and $`\epsilon `$ is the flattening. The best fit values for these parameters (Robin et al., 2000) are $`\rho _{}`$ = 1.64E(-4) stars.pc<sup>-3</sup>, $`n`$ = 2.44 and $`\epsilon `$ = 0.76. The Halo in the $`g`$,$`r`$ CMD seems to have a very faint trend redwards at $``$24<sup>th</sup> magnitude, however in the $`V`$,$`i`$ CMD there is a distinct turn redward in the Halo around 22<sup>nd</sup> magnitude.
These figures illustrate why the Galactic component of the CMD does not introduce any strong main sequences into the CMD. Instead, they are smooth with overdensities at the extremes of the magnitude range, with each component introducing redward trends at different places in the CMD. Understanding the CMDs in this way allows for greater understanding of the data when interpreting the observations.
CMDs of the data were produced using the matched catalogues from the data-reduction pipeline, shown as Hess plots<sup>11</sup><sup>11</sup>11A Hess plot is a pixelated Colour-Magnitude diagram where the grayscale denotes the square root of the Colour-Magnitude diagram number density in Figures 3 to 7. The Galactic coordinate notation from here on will be written in the form (l, $`\pm `$b). In Figures 3 to 7, the top-left panel is the Northern field, the top-right panel is the Southern field, each containing the actual data from the observing run. The lower panels are solely synthetic CMD’s and are again North on the lower-left and South on the lower-right.
All previous detections of the MRi have shown it to be beyond the Disk of the galaxy. This means that the MRi main sequence should be in the area of CMD typically containing the Halo and Thick Disk stars. Since the density of stars here is intrinsically low, the MRi main sequence should contrast well against these weaker components. However, if the stream approaches close to the edge of the Thick Disk, there is the possibility of confusion between the two. All of the possible detections of the stream that result from this work are in this region between the Thick Disk and the Halo and so there remains some ambiguity over how to separate these two structures from the stream. Each of the detections represents large excursions from the synthetic galaxy model, typically $``$1 magnitude or greater, corresponding to a distance shift of $``$6 kpc according to our estimation. Comparisons with the model in this manner is unsuited towards minor deviations, but given the distances involved in a $``$1 magnitude difference this suggests the features represent more than just problems with the synthetic galaxy model. If this were the case then the model is indeed very inadequate in these regions. Since we do not expect this to be the case we proceed with this method. Ultimately, kinematic surveys of these regions will be needed to resolve some of the remaining uncertainties but in the meantime we have used careful comparison with the models of Robin, Reylé, Derrière & Picaud (2003) coupled with completeness and magnitude-error estimates to draw out any extra features from the Colour-Magnitude diagrams.
#### 4.1.4 Magnitude Completeness
The magnitude completeness of the data was estimated using the overlap regions of the observations. Each field consists of several overlapping subfields (see Table 1). Stars in the overlap regions will appear in more than one subfield. So by dividing the number of stars that can be matched across subfields by the total number of stars in the overlap region (for a given magnitude bin), the completeness curve can be formed. Fitting the curve with the function,
$$CF=\frac{1}{1+e^{(mm_c)/\lambda }}$$
(8)
where m is the magnitude of the star, m<sub>c</sub> is the magnitude at 50% completeness and $`\lambda `$ is the width of the rollover from 100% completeness to 0% completeness. By applying the completeness curve of the data to the model makes for a more suitable comparison. Our method for determining completeness does not account for faint stars which lie close to bright stars and thus are excluded from our completeness estimate. The values used to model each field can be found in Table 2.
With regards to the real data, the uncertainty in the magnitude increases with increasing magnitude. This explains why the structures in the real data loses coherency at the faintest extremes of the Hess plots. By plotting the magnitude against the error of the magnitude it can be fitted with the following function.
$$f(x)=A+Be^x+Cx^2$$
(9)
where f(x) is the error in the magnitude and x is the magnitude<sup>12</sup><sup>12</sup>12Parameters for these equations will be supplied upon request (bconn@physics.usyd.edu.au).
This allows a similar effect to be introduced into the model figures when making quantitative measurements. The completeness and magnitude error functions have not been applied when making initial comparisons between the model and the data in a qualitative sense. Comparing obvious structural differences with the model and finding the distance to any new structure is a qualitative approach and does not require the model to be corrected in this manner. Then, when attempting to find the signal-to-noise of our detections quantitatively, we have applied both corrections to the model adjusting it to match the data.
#### 4.1.5 Estimating the Distance
Using the method employed by Ibata et al. (2003), we too have used the colour-transformation which converted the ridge-line of the SDSS S223+20 field \[ Newberg et al. (2002), see Figure 12.\] to a main sequence type overlay. By determining the offset of any new structure from the base position of this ridge-line we can estimate the distance. To do this we need to convert the SDSS ($`g^{}`$,$`r^{}`$) system to the Vega-normalized ($`g`$,$`r`$) and ($`V`$,$`i`$), this can be done by comparing overlapping INT and SDSS fields (SDSS fields taken from Early Data Release, Stoughton et al. (2002)). In particular, these conversions use a comparison field near the Galactic South Pole. The relevant colour conversions are:
for $`g`$ and $`r`$,
$$(gr)=0.21+0.86(g^{}r^{})$$
(10)
$$g=g^{}+0.150.16(gr)$$
(11)
for $`V`$ and $`i`$,
$$V=g0.030.42(gr)$$
(12)
$$(gi)=0.09+1.51(gr)$$
(13)
The ($`V`$,$`i`$) conversion is taken from Windhorst et al. (1991)<sup>13</sup><sup>13</sup>13These conversions may have significant systematic errors and as such any distance estimate derived using this method should merely be taken as indicative in particular, the conversion from ($`g^{}`$,$`r^{}`$) to ($`V`$,$`i`$). The zero offset Heliocentric distance estimate is assumed to be 11 kpc (Newberg et al., 2002) and any deviation in magnitude required to align this main sequence style overlay is assumed to be solely due to distance variations. The Heliocentric distance is then calculated using Eqn.14 and the Galactocentric distance is found using simple trigonometry assuming the Sun is located 8.0 kpc from the Galactic centre.
$$R_{HC}=11.0\left(10^{\frac{offset}{5.0}}\right)$$
(14)
Determining a value for the error associated with such a measurement is dependent on several factors. Most predominant of these is whether the fields have been correctly calibrated with regard to their photometry taking into account the dust extinction present within the fields (Schlegel, Finkbeiner, & Davis, 1998). Given however that the accuracy of the dust maps is $``$16%, this will dominate over the few percent calibration error in the photometric zero-points as derived by the CASU pipeline (Irwin & Lewis, 2001). Having understood the errors involved in both the determination of the photometry, extinction correction and the main-sequence style overlay, manually placing this overlay at the two extremes of an acceptable fit provides a range of distances over which this structure resides. The final value then represents the average of this estimate rounded to the nearest kiloparsec; the error in the fit is typically less than 1 kpc and thus is dominated by the $``$10% error which naturally resides in the original distance calculation. Given the large errors involved these distances can only be considered indicative of the true distance.
### 4.2 Individual Fields
#### 4.2.1 Fields $`(61,\pm 15)^{}`$
Presented in the top panels of Figure 3, the data is shown as a Hess plot. The lower panels were generated by the synthetic galactic model of Robin, Reylé, Derrière & Picaud (2003) to serve as comparison figures for the data. Note that the model fields have not been completeness corrected for these comparisons. When comparing with Figure 2, we can see each of the expected Galactic components present in the data. Thus the observed field matches the synthetic field to a high degree supporting our assumption that the synthetic fields will act as good comparisons to the data. In these fields, the Halo component is seemingly showing a redward trend brighter than the model, however the strength of this feature was deemed not sufficiently significant to infer a detection. $`\mathrm{\S }`$5.2.2 re-examines this field in light of the model produced by Peñarrubia et al. (2004).
#### 4.2.2 Fields $`(75,\pm 15)^{}`$
Figure 4 shows the data and the model for the fields (75,$`\pm `$15). These fields were observed using the $`V`$ and $`i`$ filters and this has the effect of changing the layout of the CMD. Studying Figure 2, each of the Galactic components in our data can again be easily identified . The field (75,+15) also seems to have a Halo component which turns redward brighter than the model. Similarly with the fields (61,$`\pm `$15) this feature is not significant enough to convincingly label a detection. These fields too are revisited in $`\mathrm{\S }`$5.2.2 when comparing with the results of Peñarrubia et al. (2004).
#### 4.2.3 Fields $`(90,\pm 10)^{}`$
Figure 5 shows the data and the model for the fields (90,$`\pm `$10). As with Figure 3 the layout is the same and also was observed in $`g_{}`$ and $`r_{}`$, possessing similar Galactic CMD structure. Again, visual inspection reveals no extra structure in this CMD and hence no MRi signature.
#### 4.2.4 Fields $`(118,+16)^{}`$ & $`(123,19)^{}`$
Figure 6 shows the data and the model for the fields (118,+16) and (123,-19). The (123,-19) field is from the M31 survey (Ibata et al., 2001; Ferguson et al., 2002) with the confirmation of the MRi signature presented in Ibata et al. (2003). The (118,+16) field has a Halo component in which the Main Sequence (MS) turns redward brighter than that of the model field. As from $`\mathrm{\S }`$4.2.2 the bulk Halo MS is expected to turn redward in this field at $`V_{}22.5`$. The redward trend detected occurs at $`V_{}20.021.0`$ and reaches a value of $`V_{}i_{}1.0`$ at $`V_{}23.0`$. Comparing this to the expected location of the Halo population in the model reveals that the (118,+16) field is inside this population and thus is a potential detection of the MRi. Fitting a main sequence to the feature and overlaying it on the model, as can be seen in the lower left plot of Figure 6, an offset of +0.2 magnitudes or R<sub>HC</sub> $``$12 kpc is calculated which corresponds to R<sub>GC</sub> $``$17 kpc. This is further out than is predicted by the model shown in Figure 1. The field (118,+16) is intended to be a comparison field on the symmetrically opposite side of the Galaxy. Both these fields have been observed in $`V`$ and $`i`$ and thus have the Galactic features as discussed for Figure 4. The MRi signature in (123,-19) stands out clearly against the known Galactic components. The magnitude offset is -0.8 with distances to this feature being R<sub>HC</sub> $``$7.6 kpc and R<sub>GC</sub> $``$14 kpc. Note that the circular overdensity at $`\mathrm{V}_{}23.5,\mathrm{V}_{}\mathrm{i}_{}1.2`$ is M31 in the background of this field.
#### 4.2.5 Fields $`(150,\pm 15)^{}`$
Figure 7 shows the data and the model for the fields (150,$`\pm `$15). These fields have been observed in $`V`$ and $`i`$ and thus have the Galactic features as discussed for Figure 4. Using the technique described in $`\mathrm{\S }`$4.2.4 of examining when the bulk Halo Main Sequence (MS) turns redward, in the Northern field this occurs at $`V_{}22.5`$. The data MS turns redward at $`V_{}21.0`$ and reaches a value of $`V_{}i_{}1.2`$ at $`V_{}23.0`$. Fitting a main sequence to this feature and overlaying it on the model using an offset of -0.4 magnitudes, which can be seen in the lower left plot of Figure 7, we calculate R<sub>HC</sub> $``$9 kpc, this corresponds to R<sub>GC</sub> $``$17 kpc which exceeds the model prediction in Figure 1 by $``$7 kpc. Comparing with the detection in field Mono-N (149,+20), Ibata et al. (2003) shows similar findings. They also report an offset of -0.4 magnitudes corresponding to a Heliocentric distance of $``$9 kpc and a Galactocentric distance of $``$16 kpc. In all likelihood this represents two detections of the same stream in different passbands separated by a few degrees on the sky, given this, we have interpreted it as a detection of the MRi. The (150,-15) field only vaguely resembles its synthetic field counterpart, however the strong vertical feature at $`V_{}i_{}0.5`$ seems to form a main sequence that is quite close to the edge of the Thick Disk in the current Galactic model. Determining whether this is in fact the edge of the Thick Disk is difficult considering that the Halo component is poorly defined and thus not readily available for comparison. Fitting a main sequence to this feature using an offset of -1.5, R<sub>HC</sub> $``$6 kpc and R<sub>GC</sub> $``$13 kpc.
### 4.3 A re-examination of field $`(123,19)^{}`$
Detecting the stream by noting deviations from the synthetic CMDs as in the (118,+16) and (150,$`\pm `$15) fields, rather than the obviously additional feature as seen in Figure 6 with regards to the (123,-19) field, means applying this technique to the (123,-19) field may yield a new detection. The redward trend of the Halo component begins at $`V_{}`$ 22-22.5 in the model while in the data this occurs at $`V_{}`$ 21-21.5, suggesting this feature is not of Galactic origins. Fitting a main sequence overlay with an offset of +0.8 to this, as can be seen in Figure 8, distance estimates of R<sub>HC</sub> $``$16 kpc and R<sub>GC</sub>$``$21 kpc are obtained for this new structure.
If however, this is to believed, then there are two structures present in this field. Since this field is in the vicinity of the newly discovered TriAnd structure perhaps this is a detection of it in the background. The current distance estimate to the TriAnd structure is R<sub>HC</sub>$``$15-30 kpc, in the region l $``$100 to $``$170 and b $``$-15 to $``$-60. This places both detections in this field on the edge of this structure. Although association between the MRi and TriAnd is still speculative, TriAnd may represent a wrapped tidal arm of the MRi. Pursuing this thought, a wrapped tidal arm would imply that the MRi is not a recent accretion event but rather the relic of a much older accretion. The model of Peñarrubia et al. (2004), also supports this view of wrapped tidal arms and this field in particular fits with their model. $`\mathrm{\S }`$5.2.1 discusses this in further detail. Clearly though, this is a tentative detection as the correspondence between the model and the data, in this region, is poor. This detection may also simply reveal a deficiency in the synthetic galaxy model.
### 4.4 Comparison: Data minus Model
Applying the completeness corrections and the appropriate magnitude error estimates to the model brings it as close to the real data as possible. By subtracting the corrected model from the data this minimises the residuals between the two. For those fields in which we found no “by-eye” detection, namely, (61,$`\pm `$15), (75,$`\pm `$15) and (90,$`\pm `$10), the residuals formed no coherent structures and have not been presented here. The field (150,-15) is also not shown, as the differences between the real data and the model proved too great to draw any logical conclusions. The remaining fields (118,+16), (123,-19) and (150,+15) are shown in Figures 910 and 11. The same main sequence lines as drawn in the previous Hess plots have also been overlaid on these CMDs. The features as mentioned previously are present in these figures also, see $`\mathrm{\S }`$4.2.
This procedure allowed us to quantitatively determine an estimate of the significance of these features, found by dividing the number of stars in the feature by the Poisson noise due to the stars in the region. The feature in Figure 9, from (118,+16), has a signal-to-noise ratio of $``$20. For field (123,-19) (Figure 10), the estimated signal-to-noise value is $``$20, which is in accordance with Ibata et al. (2003). Finally for (150,+15) (Figure 11), the signal-to-noise was found to be $``$12. While we may not know the origins of these features, these results confirm that they are significant contributors to the CMDs in these regions.
## 5 Discussion
### 5.1 Summary of Results
Using the Isaac Newton Telescope Wide Field Camera in La Palma, Canary Islands we have undertaken a survey of the Monoceros Ring (MRi) in the region of Galactic longitudes l = (61 - 150) in symmetric pairs above and below the plane. Table 1 shows the fields observed. The MRi can be seen in many fields across the sky (Figure 12) and also in 2MASS distribution of M-giant stars (Rocha-Pinto et al., 2003). The origin of this structure is still poorly understood and a systematic survey around the Galactic plane is required to understand the size and extent of this structure.
Surveying the region of sky l= (61 - 150) and through a re-examination of field (123,-19), four tentative detections of the MRi have been found. Fitting main sequences to those CMDs which showed a deviation from the synthetic galaxy model has yielded distance estimates in those fields. The distance estimates both Heliocentric and Galactocentric can be seen in Table 3.
With detections in each of these cases being on both sides of the Galactic plane it suggests the stream is either very broad on the plane of the sky or represents more than one arm of the tidal debris wrapping around on themselves. Also, the detection at (150,+15) is on the opposite side of the Galactic plane to the newest structure in the Milky Way, namely the TriAnd feature of Rocha-Pinto et al. (2004).
Are they both related to the MRi structure? The kinematics of most of the MRi detections are unknown, however with more detections of the MRi on both sides of the plane, the MRi structure is possibly much older than previously assumed. If indeed it is an accreted dwarf galaxy, then it has seemingly made several orbits around the galaxy.
### 5.2 Re-examination of the observations based on the model of Peñarrubia et al. (2004)
The hypothesis of multiply wrapped tidal arms is investigated by Peñarrubia et al. (2004) who have created an detailed model of the Monoceros tidal stream using all the data currently available. By testing thousands of models of tidal streams, in both prograde and retrograde directions around the Milky Way and correlating them against the current observations they refined the parameters for the tidal stream. The solution they favour is one where the accreting dwarf galaxy has made two orbits of the galaxy in the prograde direction leaving two concentric tidal streams.
#### 5.2.1 Detections: Fields (118,+16),(123,-19) & (150,$`\pm `$15)
Interestingly, Figure 7 of Peñarrubia et al. (2004) shows the location of the stream in Galactic coordinates and the corresponding Heliocentric distances to the stream around the Galactic equator, allowing for a direct comparison with our data. At (118,+16) their model shows the stream around $``$17 kpc away and at the location (123,-19) shows a region where the two streams cross, the inner ring being $``$17 kpc away and the outer ring $``$10 kpc away. (150,+15) resides in a region where both rings are present but sparsely populated. The range of distances to the stream is $``$7 - 16 kpc. (150,-15) lies predominantly on their inner ring and has a distance of $``$8 kpc. The detections correspond well with this model, excluding (150,-15), which as stated previously does not resemble the expected CMD and is open to interpretation. For a full comparison of the distances found, see Table 3. It is encouraging that the model provides some support to our findings that two tidal streams are present in the (123,-19) field.
#### 5.2.2 Non-detections: Fields (61,$`\pm `$15), (75,$`\pm `$15) & (90,$`\pm `$10)
Considering those fields in which no “by-eye” detection was made, the results of Peñarrubia et al. (2004) allow some more tentative detections. For the (61,$`\pm `$15) field, the Northern field stream is located $`>`$14 kpc away (Heliocentric distance), in a region where their model shows both the near and far streams crossing in the plane of the sky \[see Figure 7. of Peñarrubia et al. (2004)\]. Fitting a main sequence to the Halo structure in this field, it should be first noted that it turns redward brighter than the model, in the same manner as the detections in the other fields and has a distance of $``$16 kpc. The Southern field, while in a region again where two streams are crossing, the distance to these structures is $`>`$21 kpc and fitting a main sequence to the Halo component here reveals a distance of $``$19 kpc. The main sequence overlays for these distance estimates can be seen in Figure 13. The (75,$`\pm `$15) fields, being close to the (61,$`\pm `$15) fields, has similar properties as described above. The stream in the model at (75,+15) has a distance of $``$8 - 17 kpc. Fitting a main sequence to the Halo component of this field returns a distance of $``$15 kpc, this can be seen in Figure 14. It should be noted that in comparison with the model the location of the main sequence overlay seems coincident with the edge of the Thick Disk (see Figure 2). However, since the distances match those in the Peñarrubia et al. (2004) model it is being included in the list of tentative new detections. No kinematics are known about these regions and thus are needed to resolve definitively this ambiguity. The distance to the stream at (75,-15) in the model is $``$20-31 kpc. Since there is no obvious feature to fit to in our data, it can be assumed that the stream is not present due to the distances involved. The remaining field in which we found no “by-eye” signature of the stream was (90,$`\pm `$10). The region (90,+10) is associated with the inner and outer ring of the model and has a distance of $``$11 - 22 kpc. Our observations do not support the location of the stream in this part of the sky and no main sequence was fitted to any part of the corresponding CMD. (90,-10) resides on the edge of the outer ring which has a corresponding distance $``$9 - 14 kpc and there is no obvious feature again to fit to, possibly due to being on the edge of the stream in this location hence a smaller number of stream stars or differences between the model and the true location of the stream. The subtraction figures for the (61,$`\pm `$15), (75,$`\pm `$15) & (90,$`\pm `$10), as seen with the (118,+16),(123,-19) & (150,$`\pm `$15) fields have not been shown. Due to the relative weakness of the structures we have identified and the slight differences between the model and the data, the signal-to-noise estimates of these tentative detections are not informative.
The Peñarrubia et al. (2004) model was constrained so that it accords with the current published knowledge of the MRi detections. The correlation with this new dataset is an encouraging sign that the tidal stream they have modelled and the tentative detections found here are indeed a real part of this new Milky Way structure.
### 5.3 Could the Monoceros Ring be a detection of the Galactic warp?
Recently, Momany et al. (2004) have claimed that the identification of the Canis Major Dwarf (CMa) galaxy is none other than a misinterpretation of the Galactic warp. Martin et al. (2004b) have refuted this claim illustrating how the CMa population differs from that of the warp. However, what certainty do we have that the MRi is also not just a misinterpretation of the Galactic warp?
Firstly, what signature would the warp have on the sky if the MRi detections were merely sampling the extent of the warp? In the synthetic models of Robin, Reylé, Derrière & Picaud (2003), the Sun is assumed to be located along a node, with the warp reaching a maximum projected angular extent above the plane at l = 90 and a maximum angular extent below the plane at l = 270, when viewed from the Solar neighbourhood. Thus all detections of the warp should be above the plane in the region l = 0 \- 180 and below the plane in the region l = 180 \- 360. The distance to the warp would be a minimum at each point of maximum height above the Disk and reach a maximum distance at l = 0, 180. Also, for any given line-of-sight the warp should only be detected once. Calculating the distance to the warp using the model from Robin, Reylé, Derrière & Picaud (2003), for the lines of sight corresponding to the INT/WFC fields, the projected distance to warp is greater than 30 kpc. Since this is greater than the expected extent of the Galactic Disk, the warp should not present itself in these fields.
In Figure 12, all of the detections of the MRi through CMD identification have been plotted in a top-down view of the galaxy, showing the coordinates of these detections. The quavers (notes with tails) are detections above the plane of the galaxy (b $`>`$ 0) and the crochets (notes without tails) are detections below the plane of the galaxy (b $`<`$ 0). The colours denote the different observers. The orange circle marks out the solar circle around the galaxy at 8 kpc and the black circle has a radius of 17.5 kpc. In this plot, Galactic longitude l is measured anti-clockwise from the line joining the Sun to the Galactic centre. Comparing our knowledge of how the warp signature should present itself with Figure 12, it is apparent that the detections of the MRi are not divided with those above the plane residing solely in the region l = 0-180 and those below at l = 180-360. Detections of the stream have occurred both above and below the plane at similar longitudes. The distances to the detections above the plane seem to lie along a great arc of radius $``$17.5 kpc, which is inconsistent with the modelled distance to the warp. This shows that the MRi is not a detection of the warp but a distinct feature of the Milky Way.
### 5.4 Conclusion
The on-going INT/WFC Survey of the Monoceros Ring has yielded several detections of the ring in the region l,b= (118,16), (150, 15) and a tentative detection at (150,-15). Galactocentric distance estimates to these structures gave $``$17,$``$17, and $``$13 kpc respectively. These are combined with a reexamination of the field presented by Ibata et al. (2003), (123,-19), showing the position of the Halo is not in accordance with the model and possibly represents another detection of the ring. The Galactocentric distance to this feature is estimated at $``$21, kpc. This provides evidence that the ring may be wrapped around the galaxy more than once. This is also supported by the model of Peñarrubia et al. (2004).
In light of claims by Momany et al. (2004), that the recently discovered Canis Major Dwarf galaxy was in fact a misinterpretation of the Galactic warp, it seemed necessary to see whether the MRi could also be confused with the warp. The strongest evidence that the MRi is not the warp is that with this survey we have shown that the MRi can be detected on both sides of the Galactic plane at similar Galactic longitudes. Considering the warp cannot have this structure the MRi can be ruled out as a result of the warp.
All the detections have been found by comparing the data with the synthetic galactic model of Robin, Reylé, Derrière & Picaud (2003). This requires the assumption that the major components of the Milky Way have smooth distributions and that major excursions from the expected location of these features must represent something new. Our detections all represent major differences with the synthetic model, typically of the order of 1 magnitude or more which corresponds to a distance variation of $``$6 kpc (Heliocentrically). The distances to the detections have been calculated using a main-sequence style overlay modeled on the ridge line of the SDSS S223+20 detection. Given an offset from this original detection we can determine a distance. While there may be significant systematic errors in the colour conversions, especially from ($`g^{}`$,$`r^{}`$) to ($`V`$,$`i`$) we have included the upper and lower limits of all the distances derived to demonstrate the narrowness of the main sequences and the intrinsic errors of this method.
Taking into account the recent paper by Peñarrubia et al. (2004), the observations have been re-examined and the distances and locations of the MRi stream favourably match with their model. Comparing the CMDs of the remaining fields reveal several more tentative detections, namely in (61,$`\pm `$15) & (75,+15). These have Heliocentric distances of $``$16, 19 & 15 kpc and Galactocentric distances of $``$14, 17 & 15 kpc respectively. Each of these features had been considered earlier in the analysis of the data, however they were deemed to weak to be significant. This model provides support for the authenticity of these detections. Importantly, the model of Peñarrubia et al. (2004) also confirms one of our non-detections, (75,-15) in which the distance to the stream is too great to be detected. The non-detections in the fields (90,$`\pm `$10), seemingly do not correspond to the predictions of the model, although a deeper survey of this region may be needed to resolve this discrepancy. It is unknown whether the low Galactic latitude of this field is a factor in the non-detection of the stream.
Both detections and non-detections support a complex picture of the MRi. In particular, those detections above the plane suggest the MRi has an extended stream tracing an arc $``$17 kpc from the Galactic centre, while the detections below the plane, reveal a tentative detection of the TriAnd region in the background of the (123, -19) region and also the presence of a foreground stream if we link the INT/WFC detections and those of Newberg et al. (2002) & Ibata et al. (2003) (Figure 12). Obviously, the MRi is a very complex structure and additional observations are needed to unravel its origins, so to this end a kinematic survey of the region surrounding the Canis Major overdensity (Martin et al., 2005) has been conducted with the Two Degree Field (2dF) spectrograph and so too a Wide- Camera survey using the Anglo-Australian Observatory Wide Field Imager (AAO/WFI) to complete the survey of the MRi. These results are still being analysed and will form the basis of a forthcoming article.
## 6 Acknowledgements
BCC would like to thank his wife, LLL, for kindly supplementing his scholarship income, The University of Sydney for the University Postgraduate Award and the Cambridge Astronomical Survey Unit at Cambridge University and Mike Irwin for their hospitality during my week there. BCC would also like to thank Jorge Peñarrubia, for access to his Monoceros Ring model and the anonymous referee for their many helpful suggestions. GFL acknowledges the support of the Discovery Project grant DP0343508. The research of AMNF has been supported by a Marie Curie Fellowship of the European Community under contract number HPMF-CT-2002-01758. GFL would also like to thank Triple J for their chillout session on Sunday mornings which drowns out his fighting children and also their Three Hours of Power in case they choose to fight between the hours of 11 and 1 at night.
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# SUSY Anomalies Break 𝒩=2 to 𝒩=1: The Supersphere and the Fuzzy Supersphere
## 1 Overview
In this section, we give an overview of fuzzy SUSY as full details can be found elsewhere . In later sections, we will explain all the requisite details to develop instanton theory.
### 1.1 The Fuzzy Sphere
We recall that the fuzzy sphere $`S_F^2(n)`$ is the $`(n+1)\times (n+1)`$ matrix algebra $`Mat(n+1)`$. It can be realized as linear operators on $`^{n+1}`$ with the orthonormal basis vectors
$$\frac{(a_1^{})^{n_1}}{\sqrt{n_1!}}\frac{(a_2^{})^{n_2}}{\sqrt{n_2!}}|0,n_1+n_2=n,$$
(1.1)
where $`a_i,a_i^{}`$ are bosonic oscillators. The vectors (1.1) span a subspace of the Fock space with fixed particle number $`n`$:
$$N:=\underset{i}{}a_i^{}a_i,N|_{^{n+1}}=n.$$
(1.2)
In this representation, the elements of $`S_F^2(n)`$ are the linear operators
$$\underset{i,j}{}c_{i,j}^m(a_i^{})^m(a_j)^m,c_{i,j}^m,$$
restricted to the subspace $`^{n+1}`$.
The group $`SU(2)`$ acts on $`^{n+1}`$ and hence on $`S_F^2(n)`$ by its spin $`\frac{n}{2}`$ unitary irreducible representation. The angular momentum generators are
$$L_i=a^{}\frac{\sigma _i}{2}a,\sigma _i\text{ are Pauli matrices.}$$
(1.3)
### 1.2 SUSY
The $`𝒩=1`$ SUSY version of $`SU(2)`$ is $`OSp(2,1)`$. It has the graded Lie algebra $`osp(2,1)`$. Its generators (basis) can be written using oscillators if we introduce one additional fermionic oscillator $`b`$ and its adjoint $`b^{}`$. They commute with $`a_i,a_j^{}`$. Then the $`osp(2,1)`$ generators are
$`\mathrm{\Lambda }_i`$ $`=`$ $`a^{}{\displaystyle \frac{\sigma _i}{2}}a,\mathrm{\Lambda }_4={\displaystyle \frac{1}{2}}(a_1^{}b+b^{}a_2),`$ (1.4)
$`\mathrm{\Lambda }_5`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a_2^{}b+b^{}a_1),\sigma _i\text{ = Pauli matrices.}`$
The $`𝒩=2`$ SUSY version of $`SU(2)`$ is $`OSp(2,2)`$. It has the graded Lie algebra $`osp(2,2)`$. Its basis consists of the $`osp(2,1)`$ generators and three additional generators
$`\mathrm{\Lambda }_4^{}`$ $``$ $`\mathrm{\Lambda }_6={\displaystyle \frac{1}{2}}(a_1^{}bb^{}a_2),\mathrm{\Lambda }_5^{}\mathrm{\Lambda }_7={\displaystyle \frac{1}{2}}(a_2^{}b+b^{}a_1),`$ (1.5)
$`\mathrm{\Lambda }_8`$ $`=`$ $`a^{}a+2b^{}b.`$
If $`\{,\}`$ denotes the anticommutator, $`osp(2,2)`$ has the defining relations
$`[\mathrm{\Lambda }_i,\mathrm{\Lambda }_j]`$ $`=`$ $`i\epsilon _{ijk}\mathrm{\Lambda }_k,[\mathrm{\Lambda }_i,\mathrm{\Lambda }_\alpha ]={\displaystyle \frac{1}{2}}\mathrm{\Lambda }_\beta (\sigma _i)_{\beta \alpha },\{\mathrm{\Lambda }_\alpha ,\mathrm{\Lambda }_\beta \}={\displaystyle \frac{1}{2}}(\epsilon \sigma _i)_{\alpha \beta }\mathrm{\Lambda }_i,`$ (1.6)
$`[\mathrm{\Lambda }_i,\mathrm{\Lambda }_8]`$ $`=`$ $`0,[\mathrm{\Lambda }_\alpha ,\mathrm{\Lambda }_8]=\mathrm{\Lambda }_\alpha ^{},\{\mathrm{\Lambda }_\alpha ,\mathrm{\Lambda }_\alpha ^{}\}={\displaystyle \frac{1}{4}}\epsilon _{\alpha \beta }\mathrm{\Lambda }_8,`$
$`\{\mathrm{\Lambda }_\alpha ^{},\mathrm{\Lambda }_\beta ^{}\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\epsilon \sigma _i)_{\alpha \beta }\mathrm{\Lambda }_i,[\mathrm{\Lambda }_\alpha ^{},\mathrm{\Lambda }_8]=\mathrm{\Lambda }_\alpha ,\epsilon =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$ (1.9)
These relations show in particular that the additional three generators form a triplet under $`osp(2,1)`$.
Conventional Lie algebras like that of $`su(2)`$ have a $``$ or an adjoint operation defined on them. For $`\mathrm{\Lambda }_i`$, it is just $`\mathrm{\Lambda }_i^{}=\mathrm{\Lambda }_i`$. This follows from the fact that $`a_i^{}`$ is the adjoint of $`a_i`$. For $`osp(2,1)`$ and $`osp(2,2)`$, is replaced by the grade adjoint . On the oscillators, is defined by
$$a_i^{}=a_i^{},(a_i^{})^{}=(a_i^{})^{}=a_i,b^{}=b^{},(b^{})^{}=b.$$
Hence $`=`$ on bosonic oscillators.
On products of operators, is defined as follows. We assign the grade 0 to $`a_i,a_j^{}`$ and their products and 1 to $`b`$ and $`b^{}`$. The grades are additive (mod 2). The grade of an operator $`L`$ with definite grade is denoted by $`|L|`$. Then if $`L,M`$ have definite grades, $`(LM)^{}(1)^{|L||M|}M^{}L^{}`$. Hence $`(b^{}b)^{}=b^{}b`$ and
$$\mathrm{\Lambda }_i^{}=\mathrm{\Lambda }_i,\mathrm{\Lambda }_\alpha ^{}=\epsilon _{\alpha \beta }\mathrm{\Lambda }_\beta ,\mathrm{\Lambda }_\alpha ^{}^{}=\epsilon _{\alpha \beta }\mathrm{\Lambda }_\beta ^{}\mathrm{\Lambda }_8^{}=\mathrm{\Lambda }_8.$$
(1.10)
### 1.3 Irreducible Representations
Let $`osp(2,0)`$ denote $`su(2)`$, the Lie algebra of $`SU(2)`$. Its IRR’s are $`\mathrm{\Gamma }_J^0`$, $`J/2`$. (Here $`=\{0,1,2,\mathrm{}\}`$.) $`J`$ has the meaning of angular momentum.
The $`osp(2,1)`$ algebra is of rank 1 just as $`osp(2,0)`$. We can take $`\mathrm{\Lambda }_3`$ to be the generator of its Cartan subalgebra. Since
$$[\mathrm{\Lambda }_3,\mathrm{\Lambda }_4]=\frac{1}{2}\mathrm{\Lambda }_4,[\mathrm{\Lambda }_3,\mathrm{\Lambda }_+=\mathrm{\Lambda }_1+i\mathrm{\Lambda }_2]=\mathrm{\Lambda }_+,$$
$`\mathrm{\Lambda }_4,\mathrm{\Lambda }_+`$ are its raising operators. They commute:
$$[\mathrm{\Lambda }_4,\mathrm{\Lambda }_+]=0.$$
In an IRR, both vanish on the highest weight vector. The eigenvalue $`J/2`$ of $`\mathrm{\Lambda }_3`$ on the highest weight vector can be used to label its IRR’s. They are denoted by $`\mathrm{\Gamma }_J^1`$ in this paper.
When restricted to its subalgebra $`osp(2,0)`$, $`\mathrm{\Gamma }_J^1`$ splits as follows:
$$\mathrm{\Gamma }_J^1|_{osp(2,0)}=\mathrm{\Gamma }_J^0\mathrm{\Gamma }_{J\frac{1}{2}}^0,J\frac{1}{2}.$$
(1.11)
$`\mathrm{\Gamma }_0^1`$ is the trivial IRR.
The dimension of $`\mathrm{\Gamma }_J^1`$ is $`4J+1`$.
The graded Lie algebra $`osp(2,2)`$ is of rank 2. A basis for its Cartan subalgebra is $`\{\mathrm{\Lambda }_3,\mathrm{\Lambda }_8\}`$. Since
$$[\mathrm{\Lambda }_3,\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{}]=\frac{1}{2}(\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{}),[\mathrm{\Lambda }_8,\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{}]=\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{},$$
$`\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{}`$ serves as the raising operator for both $`\mathrm{\Lambda }_3`$ and $`\mathrm{\Lambda }_8`$. We also have that $`\mathrm{\Lambda }_1+i\mathrm{\Lambda }_2=\mathrm{\Lambda }_+`$ is the raising operator for $`\mathrm{\Lambda }_3`$ alone:
$$[\mathrm{\Lambda }_3,\mathrm{\Lambda }_+]=\mathrm{\Lambda }_+,[\mathrm{\Lambda }_8,\mathrm{\Lambda }_+]=0.$$
The raising operators $`\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{}`$ and $`\mathrm{\Lambda }_+`$ commute:
$$[\mathrm{\Lambda }_4+\mathrm{\Lambda }_4^{},\mathrm{\Lambda }_+]=0.$$
Both vanish on the highest weight vector in an IRR while the eigenvalues $`J/2`$ and $`k`$ of $`\mathrm{\Lambda }_3`$ and $`\mathrm{\Lambda }_8`$ on the highest weight vector can be used as labels of the IRR. They are denoted in this paper by $`\mathrm{\Gamma }_J^2(k)`$.
The $`osp(2,2)`$ IRR’s fall into classes, the typical and atypical (or short) IRR’s. In the typical IRR’s, $`|k|2J`$ or $`k=J=0`$, while in the atypical IRR’s, $`|k|=2J0`$. The typical IRR with $`|k|2J`$ restricted to $`osp(2,1)`$ splits as follows:
$$\mathrm{\Gamma }_J^2(k)|_{osp(2,1)}=\mathrm{\Gamma }_J^1\mathrm{\Gamma }_{J\frac{1}{2}}^1,J\frac{1}{2}.$$
$`\mathrm{\Gamma }_0^2(0)`$ is the trivial representation.
The atypical IRR’s $`\mathrm{\Gamma }_J^2(\pm \frac{J}{2})`$ ($`J1/2`$) remain irreducible on restriction to $`osp(2,1)`$:
$$\mathrm{\Gamma }_J^2(\pm J/2)|_{osp(2,1)}=\mathrm{\Gamma }_J^1.$$
$`\mathrm{\Gamma }_J^2(\pm J/2)`$ can also be abbreviated to $`\mathrm{\Gamma }_{J\pm }^2`$:
$$\mathrm{\Gamma }_J^2(\pm J/2)\mathrm{\Gamma }_{J\pm }^2,J1/2.$$
$`osp(2,2)`$ admits the automorphism
$$\tau :\mathrm{\Lambda }_i\mathrm{\Lambda }_i,\mathrm{\Lambda }_\alpha \mathrm{\Lambda }_\alpha ,\mathrm{\Lambda }_\alpha ^{}\mathrm{\Lambda }_\alpha ^{},\mathrm{\Lambda }_8\mathrm{\Lambda }_8$$
(1.12)
which interchanges $`\mathrm{\Gamma }_J^2(\pm k)`$:
$$\tau :\mathrm{\Gamma }_J^2(k)\mathrm{\Gamma }_J^2(k).$$
### 1.4 Casimir Operators
The $`osp(2,0):=su(2)`$ Casimir operator $`K_0`$ is well-known:
$$K_0=\mathrm{\Lambda }_i^2.$$
The $`osp(2,1)`$ Casimir operator is
$$K_1=\mathrm{\Lambda }_i^2+\epsilon _{\alpha \beta }\mathrm{\Lambda }_\alpha \mathrm{\Lambda }_\beta .$$
We have that
$$K_1|_{\mathrm{\Gamma }_J^1}=J(J+\frac{1}{2})1\mathrm{l}.$$
The $`osp(2,2)`$ quadratic Casimir operator is
$$K_2=K_1\epsilon _{\alpha \beta }\mathrm{\Lambda }_\alpha ^{}\mathrm{\Lambda }_\beta ^{}\frac{1}{4}\mathrm{\Lambda }_8^2:=K_1V_0.$$
(1.13)
It has the property
$`K_2|_{\mathrm{\Gamma }_J^2(k)}`$ $`=`$ $`J^2{\displaystyle \frac{k^2}{4}},`$
$`K_2|_{\mathrm{\Gamma }_{J\pm }^2}`$ $`=`$ $`0.`$ (1.14)
As already mentioned, the IRR’s $`\mathrm{\Gamma }_{J\pm }^2`$ can be distinguished by the sign of $`\mathrm{\Lambda }_8`$ on the highest weight vector.
$`osp(2,2)`$ also has a cubic Casimir operator , but we will not have occasion to use it.
### 1.5 Tensor Products
The basic Clebsh-Gordan series we need to know is as follows:
$$\mathrm{\Gamma }_J^1\mathrm{\Gamma }_K^1=\mathrm{\Gamma }_{J+K}^1\mathrm{\Gamma }_{J+K1/2}^1\mathrm{}\mathrm{\Gamma }_{|JK|}^1.$$
### 1.6 The Supertrace and the Grade Adjoint
Because of the decomposition (1.11), the vector space $`^{4J+1}`$ on which $`\mathrm{\Gamma }_J^1`$ acts can be written as $`^{2J+1}^{2J}`$ where the first term has angular momentum $`J`$ and the second term has angular momentum $`J1/2`$. By definition, the first term is the even subspace and the second term is the odd subspace. The supertrace $`str`$ of a matrix
$$M=\left(\begin{array}{cc}P_{(2J+1)\times (2J+1)}& Q_{(2J+1)\times 2J}\\ R_{2J\times (2J+1)}& S_{2J\times 2J}\end{array}\right)$$
is accordingly
$$strM=trPtrS.$$
The grade adjoint $`M^{}`$ can be calculated using the rules of graded vector spaces . The result is
$$M^{}=\left(\begin{array}{cc}P^{}& R^{}\\ Q^{}& S^{}\end{array}\right)$$
This formula is coherent with (1.10).
If $`Q,R=0`$, we say that $`M`$ is even, while if $`P,S=0`$, we say that $`M`$ is odd. We assign a number $`|M|=0,1`$ (mod 2) to even and odd matrices $`M`$ respectively.
### 1.7 The Free Action
The space with $`N=n`$ has maximum angular momentum $`J=n/2`$. It carries the $`osp(2,1)`$ IRR $`\mathrm{\Gamma }_{n/2}^1`$ which splits under $`su(2)`$ into $`\mathrm{\Gamma }_{n/2}^0\mathrm{\Gamma }_{(n1)/2}^0`$. It carries either of the short $`osp(2,2)`$ IRR’s as well.
The dimension of the Hilbert space with $`N=n`$ is $`2n+1`$. We denote it by $`^{2n+1}`$. It is the direct sum $`^{n+1}^n`$ where $`^{n+1}`$ is the even subspace carrying the IRR $`\mathrm{\Gamma }_{n/2}^0`$ and $`^n`$ is the odd subspace carrying the representation $`\mathrm{\Gamma }_{(n1)/2}^0`$. A basis for $`^{2n+1}`$ is
$$\frac{(a_1^{})^{n_1}}{\sqrt{n_1!}}\frac{(a_2^{})^{n_2}}{\sqrt{n_2!}}(b^{})^{n_3}|0,n_i=n,n_3(0,1)\text{ , }(b^{})^0:=1\mathrm{l}.$$
(1.15)
The fuzzy SUSY $`S_F^{2,2}`$ (in the zero instanton sector) is the matrix algebra $`Mat(4J+1)=Mat(2n+1)`$. Just as $`S_F^2`$, it can be realized using oscillators. In terms of oscillators, a typical element is
$$\underset{i,j}{}c_{i,j}^m(a_i^{})^m(a_j)^m+\underset{i,j}{}d_{i,j}^{m1}(a_i^{})^{m1}(a_j)^{m1}b^{}b,c_{i,j}^m,d_{i,j}^{m1}.$$
It is to be restricted to the space $`^{2n+1}`$.
The left- and right-actions
$$\mathrm{\Lambda }_\rho ^LM=\mathrm{\Lambda }_\rho M,\mathrm{\Lambda }_\rho ^RM=(1)^{|\mathrm{\Lambda }_\rho ||M|}M\mathrm{\Lambda }_\rho $$
of $`osp(2,𝒩)`$ on $`Mat(2n+1)`$ give two commuting IRR’s of $`osp(2,𝒩)`$. Here, $`\mathrm{\Lambda }_\rho osp(2,𝒩),𝒩=1,2,MMat(2n+1)`$ and both $`\mathrm{\Lambda }_\rho `$ and $`M`$ are of definite grade $`|\mathrm{\Lambda }_\rho |,|M|`$ (mod 2).
Combining the left- and right- representations, we get the grade adjoint representation
$$\text{gad}:\mathrm{\Lambda }_\rho \text{gad}\mathrm{\Lambda }_\rho =\mathrm{\Lambda }_\rho ^L\mathrm{\Lambda }_\rho ^R,\rho (i,\alpha ,\alpha ^{},8)$$
of $`osp(2,𝒩)`$.
With regard to gad, $`Mat(4J+1)`$ transforms as
$$\mathrm{\Gamma }_J^1\mathrm{\Gamma }_J^1=\mathrm{\Gamma }_{2J}^1\mathrm{\Gamma }_{2J1/2}^1\mathrm{\Gamma }_{2J1}^1\mathrm{}\mathrm{\Gamma }_0^1.$$
(1.16)
$`osp(2,2)`$ acts on $`Mat(4J+1)`$ by $`L,R`$ and gad representations as well. The $`L`$ and $`R`$ are the short representations $`\mathrm{\Gamma }_{J\pm }^2`$ so that under gad, $`Mat(4J+1)`$ transforms as $`\mathrm{\Gamma }_{J+}^2\mathrm{\Gamma }_J^2`$. Its reduction can be inferred from (1.16) once we know that $`\mathrm{\Gamma }_j^2(0)|_{osp(2,1)}=\mathrm{\Gamma }_j^1\mathrm{\Gamma }_{j1/2}^1`$. We will see this later. Hence
$$\mathrm{\Gamma }_{J+}^2\mathrm{\Gamma }_J^2=\mathrm{\Gamma }_{2J}^2(0)\mathrm{\Gamma }_{2J1}^2(0)\mathrm{}\mathrm{\Gamma }_0^2(0).$$
The fuzzy field $`\mathrm{\Phi }`$ is an element of fuzzy SUSY. The free action for $`\mathrm{\Phi }`$ is
$$S_0=\frac{f^2}{2}str\mathrm{\Phi }^{}V_0\mathrm{\Phi },$$
where $`f`$ is a real constant and $`V_0`$ is an $`osp(2,1)`$-invariant operator. When restricted to the odd subspace, it should become the Dirac operator of .
The limit of this operator for $`J=\mathrm{}`$ was found by Fronsdal and later used effectively by Grosse et al. For $`J=\mathrm{}`$, it is the difference $`K_1K_2`$ of the Casimir operators $`K_1`$ and $`K_2`$ written as graded differential operators. This operator, for finite $`J`$, becomes
$$V_0=\epsilon _{\alpha \beta }(\mathrm{\Lambda }_\alpha ^{})(\mathrm{\Lambda }_\beta ^{})+\frac{1}{4}(\mathrm{\Lambda }_8)^2.$$
(1.17)
The simplifications of $`S_0`$ for this choice of $`V_0`$ is given elsewhere .
It is evident that $`V_0`$ is $`osp(2,1)`$-invariant. But it is less obvious that $`\text{gad}\mathrm{\Lambda }_\alpha ^{}`$, $`\text{gad}\mathrm{\Lambda }_8`$ anti-commute with $`V_0`$:
$$\{\text{gad}\mathrm{\Lambda }_\alpha ^{},V_0\}=\{\text{gad}\mathrm{\Lambda }_8,V_0\}=0.$$
(1.18)
This means that these generators are realized as chiral symmetries. Of these, $`\text{gad}\mathrm{\Lambda }_8`$, restricted to the odd sector, is just standard chirality. Thus, these generators, associated with $`osp(2,2)/osp(2,1)`$ are SUSY generalizations of conventional chirality.
We now show these results.
## 2 SUSY Chirality
Let us first exhibit the highest weight vectors of the $`su(2)`$ IRR’s which occur in $`\mathrm{\Gamma }_j^2(0)`$. Here $`j`$ is an integer. Referring to (1.11), we have that $`\mathrm{\Gamma }_j^1|_{su(2)}=\mathrm{\Gamma }_j^0\mathrm{\Gamma }_{j1/2}^0`$ for $`j1`$. The highest weight vector of $`\mathrm{\Gamma }_j^0`$ is $`(a_1^{}a_2)^j`$ as it commutes with $`\mathrm{\Lambda }_4`$ and carries the eigenvalue $`j`$ of $`\text{gad}\mathrm{\Lambda }_3`$. $`\text{gad}\mathrm{\Lambda }_5`$ maps it to $`j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_4`$, the highest weight vector of $`\mathrm{\Gamma }_{j1/2}^0\mathrm{\Gamma }_j^1`$. Thus
$$\begin{array}{cccc}\mathrm{\Gamma }_j^1|_{su(2)}=& \mathrm{\Gamma }_j^0& & \mathrm{\Gamma }_{j1/2}^0\\ \begin{array}{c}\text{Highest}\text{weight}\\ \text{vectors}\end{array}\}& (a_1^{}a_2)^j& \stackrel{\text{gad}\mathrm{\Lambda }_5}{}& j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_4,j1.\end{array}$$
(2.19)
The equation also indicates the operator mapping one highest weight vector of $`su(2)`$ to the other.
Next consider $`\mathrm{\Gamma }_{j1/2}^1\mathrm{\Gamma }_{j1/2}^0\mathrm{\Gamma }_{j1}^0`$ for $`j1`$. To distinguish the $`su(2)`$ IRR’s here from those in $`\mathrm{\Gamma }_j^1`$, we put a prime on them:
$$\mathrm{\Gamma }_{j1/2}^1|_{su(2)}=\mathrm{\Gamma }_{j1/2}^0^{}\mathrm{\Gamma }_{j1}^0^{}.$$
The highest weight state of $`\mathrm{\Gamma }_{j1/2}^1`$, commuting with $`\mathrm{\Lambda }_4`$ and with eigenvalue $`j1/2`$ for $`\text{gad}\mathrm{\Lambda }_3`$ is $`j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_6`$. And $`\mathrm{\Lambda }_5`$ maps it to the highest weight vector $`X_{j1}`$ of $`\mathrm{\Gamma }_{j1}^0^{}`$. We show $`X_{j1}`$ below. Thus
$`\begin{array}{cccc}\mathrm{\Gamma }_{j1/2}^1|_{su(2)}=& \mathrm{\Gamma }_{j1/2}^0^{}& & \mathrm{\Gamma }_{j1}^0^{}\\ \begin{array}{c}\text{Highest}\text{weight}\\ \text{vectors}\end{array}\}& j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_6& \stackrel{\text{gad}\mathrm{\Lambda }_5}{}& X_{j1},\end{array}`$ (2.24)
$`X_{j1}={\displaystyle \frac{j2J1}{4}}(a_1^{}a_2)^{j1}+{\displaystyle \frac{12j}{4}}(a_1^{}a_2)^{j1}b^{}b,j1.`$
In calculating $`X_{j1}`$, we use
$$\mathrm{\Lambda }_4\mathrm{\Lambda }_6=\frac{1}{4}(a_1^{}a_2)(2b^{}b1),a^{}a+b^{}b=2J.$$
Now $`\text{gad}\mathrm{\Lambda }_7`$, $`\text{gad}\mathrm{\Lambda }_8`$ map the vectors in (2.19) to the vectors in (2.24). The full table is
$$\begin{array}{ccccccc}& & \mathrm{\Gamma }_j^1& & (a_1^{}a_2)^j& \stackrel{\text{gad}\mathrm{\Lambda }_5}{}& j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_4\hfill \\ & & & & \mathrm{\Gamma }_j^0& & \mathrm{\Gamma }_{j1/2}^0\hfill \\ \mathrm{\Gamma }_j^2(0)& & & & \text{gad}\mathrm{\Lambda }_7& \text{gad}\mathrm{\Lambda }_8& \text{gad}\mathrm{\Lambda }_7\hfill \\ & & & & & & \\ & & \mathrm{\Gamma }_{j1/2}^1& & j(a_1^{}a_2)^{j1}\mathrm{\Lambda }_6& \stackrel{\text{gad}\mathrm{\Lambda }_5}{}& X_{j1},j1.\hfill \\ & & & & \mathrm{\Gamma }_{j1/2}^0^{}& & \mathrm{\Gamma }_{j1}^0^{}\hfill \end{array}$$
(2.25)
For $`j=0`$, we get the trivial IRR of $`osp(2,𝒩)`$’s.
Eq. (LABEL:bigdiag) shows that $`\text{gad}\mathrm{\Lambda }_\alpha ^{}`$, $`\text{gad}\mathrm{\Lambda }_8`$ map the vectors of $`\mathrm{\Gamma }_j^1`$ to those of $`\mathrm{\Gamma }_{j1/2}^1`$ ($`j1`$) and vice versa. So if $`V_0`$ has opposite eigenvalues in the representations in (LABEL:bigdiag), then we can conclude that<sup>*</sup><sup>*</sup>*To show that $`\{\text{gad}\mathrm{\Lambda }_6,V_0\}=0`$ we use the fact that $`\text{gad}\mathrm{\Lambda }_6=[\text{gad}\mathrm{\Lambda }_4,\text{gad}\mathrm{\Lambda }_8]`$. The result follows from the graded Jacobi identity.
$$\{\text{gad}\mathrm{\Lambda }_\alpha ^{},V_0\}=\{\text{gad}\mathrm{\Lambda }_8,V_0\}=0$$
identically, since $`V_0|_{\mathrm{\Gamma }_0^1}=0`$. That means that these operators associated with $`osp(2,2)/osp(2,1)`$ are chirally realized symmetries.
## 3 Eigenvalues of $`V_0`$
As $`V_0`$ is an $`osp(2,1)`$ scalar, it is enough to compute its eigenvalue on the highest weight state of $`\mathrm{\Gamma }_j^1`$ to find $`V_0|_{\mathrm{\Gamma }_j^1}`$.
As $`\mathrm{\Lambda }_4^{}=\mathrm{\Lambda }_6`$ commutes with $`(a_1^{}a_2)^j`$, we have that
$$\epsilon _{\alpha \beta }\text{gad}\mathrm{\Lambda }_\alpha ^{}\text{gad}\mathrm{\Lambda }_\beta ^{}(a_1^{}a_2)^j=(\text{gad}\mathrm{\Lambda }_4^{}\text{gad}\mathrm{\Lambda }_5^{}+\text{gad}\mathrm{\Lambda }_5^{}\text{gad}\mathrm{\Lambda }_4^{})(a_1^{}a_2)^j$$
where the sign of the second term has been switched as it is zero anyway. Thus the left-hand side of the previous formula can be written as
$$\text{gad}\{\mathrm{\Lambda }_4^{},\mathrm{\Lambda }_5^{}\}(a_1^{}a_2)^j=\frac{1}{2}\text{gad}\mathrm{\Lambda }_3(a_1^{}a_2)^j=\frac{j}{2}(a_1^{}a_2)^j.$$
Also
$$\text{gad}\mathrm{\Lambda }_8(a_1^{}a_2)^j=0.$$
Hence
$$V_0(a_1^{}a_2)^j=\frac{j}{2}(a_1^{}a_2)^j.$$
(3.26)
One quick way to evaluate $`V_0|_{\mathrm{\Gamma }_{j1/2}^1}`$ is as follows. Since $`K_1|_{\mathrm{\Gamma }_j^1}=j(j+1/2)`$, we have
$$K_2|_{\mathrm{\Gamma }_j^1}=(K_1V_0)|_{\mathrm{\Gamma }_j^1}=j^2.$$
(3.27)
But $`K_2`$ is $`osp(2,2)`$-invariant. Hence
$$K_2|_{\mathrm{\Gamma }_{j1/2}^1}=j^2.$$
Since also $`K_1|_{\mathrm{\Gamma }_{j1/2}^1}=j(j1/2)`$, we have
$$V_0|_{\mathrm{\Gamma }_{j1/2}^1}=(K_1K_2)|_{\mathrm{\Gamma }_{j1/2}^1}=\frac{j}{2}1\mathrm{l}.$$
Thus $`V_0`$ has opposite eigenvalues on $`\mathrm{\Gamma }_j^1`$ and $`\mathrm{\Gamma }_{j1/2}^1`$ .
It is important to notice that
$$K_2=(2V_0)^2.$$
That is, $`2V_0`$ is a square root of $`K_2`$, a bit in the way that the Dirac operator is the square root of the Laplacian.
## 4 Fuzzy SUSY Instantons
The manifold $`S^2`$ admits twisted $`U(1)`$ bundles labelled by a topological index or Chern number $`k`$. In the algebraic language, sections of vector bundles associated with these $`U(1)`$ bundles are described by elements of projective modules .
When $`S^2`$ becomes the graded supersphere $`S^{2,2}`$, we expect these modules to persist, and become in some sense supersymmetric projective modules. That is in fact the case. We shall see that explicitly after first studying their fuzzy analogues.
The projective modules on $`S^2`$ and $`S_F^2`$ are associated with $`SU(2)S^3`$ via Hopf fibration and Lens spaces. In the same way, the supersymmetric projective modules on $`S^{2,2}`$ and $`S_F^{2,2}`$ get associated with $`osp(2,1)`$ and $`osp(2,2)`$.
The fuzzy algebra $`S_F^{2,2}`$ of previous sections is to be assigned $`k=0`$. The elements of this algebra are square matrices mapping the space with $`N=2J`$ to the same space $`N=2J`$. We emphasize the value of $`k`$ by writing $`S_F^{2,2}(0)`$ for $`S_F^{2,2}`$. $`S_F^{2,2}(0)`$ is a bimodule for $`osp(2,2)`$ as the latter can act on the left or right of $`S_F^{2,2}(0)`$ by the IRR’s $`\mathrm{\Gamma }_{J\pm }^2(0)`$.
For $`k0`$, $`S_F^{2,2}(k)`$ is not an algebra. It can be described using projectors or equally well as maps of the vector space with $`N=2J`$ to the one with $`N=2J+k`$ . (We take $`J+\frac{k}{2}0`$. If $`k<0`$, this means $`J\frac{|k|}{2}`$.) If a basis is chosen for domain and range of $`S_F^{2,2}(k)`$, their elements become rectangular matrices with $`2J+k`$ rows and $`2J`$ columns. $`S_F^{2,2}(k)`$ as well is a bimodule for $`osp(2,2)`$. The latter acts by $`\mathrm{\Gamma }_{(J+\frac{k}{2})+}^2`$ on the left of $`S_F^{2,2}(k)`$ and by $`\mathrm{\Gamma }_J^2`$ on the right of $`S_F^{2,2}(k)`$.
The invariant associated with $`S_F^{2,2}(k)`$ is just $`k`$. The meaning of $`k`$ is
$$k=\text{Dimension of range of }S_F^{2,2}(k)\text{Dimension of domain of }S_F^{2,2}(k).$$
Scalar fields $`\mathrm{\Phi }`$ are now elements of $`S_F^{2,2}(k)`$ while $`V_0`$ is replaced by a new operator $`V_k`$ which incorporates the appropriate connection and “topological” data. We now argue, using index theory and other considerations, that the $`osp(2,1)`$-invariant $`V_k`$ is fixed by the requirement
$$V_k^2=K_2$$
where $`K_2`$ is the Casimir invariant for $`\mathrm{\Gamma }_{(J+\frac{k}{2})+}^2\mathrm{\Gamma }_J^2`$.
## 5 Fuzzy SUSY Zero Modes and their Index Theory
We begin by analyzing the $`osp(2,1)`$ and $`osp(2,2)`$ representation content of $`S_F^{2,2}(k)`$.
As regards the gad representation of $`osp(2,1)`$, it transforms according to
$$\mathrm{\Gamma }_{J+\frac{k}{2}}^1\mathrm{\Gamma }_J^1=\left(\mathrm{\Gamma }_{2J+\frac{k}{2}}^1\mathrm{\Gamma }_{2J+\frac{k}{2}\frac{1}{2}}^1\right)\left(\mathrm{\Gamma }_{2J+\frac{k}{2}1}^1\mathrm{\Gamma }_{2J+\frac{k}{2}\frac{3}{2}}^1\right)\mathrm{}\left(\mathrm{\Gamma }_{\frac{|k|}{2}+1}^1\mathrm{\Gamma }_{\frac{|k|}{2}+\frac{1}{2}}^1\right)\mathrm{\Gamma }_{\frac{|k|}{2}}^1.$$
The analogue of (LABEL:bigdiag) is:
$`2J+{\displaystyle \frac{k}{2}}j{\displaystyle \frac{|k|}{2}}+1,`$
$`\begin{array}{ccccccc}& & \mathrm{\Gamma }_j^1& & \mathrm{\Gamma }_j^0& & \mathrm{\Gamma }_{j1/2}^0\\ & & & & & & \\ \mathrm{\Gamma }_j^2(k)& & & & & & \\ & & & & & & \\ & & \mathrm{\Gamma }_{j1/2}^1& & \mathrm{\Gamma }_{j1/2}^0& & \mathrm{\Gamma }_{j1}^0.\end{array}`$ (5.33)
Here $`|k|1`$. For $`j=\frac{|k|}{2}`$, we get the atypical representation of $`osp(2,2)`$:
$$\mathrm{\Gamma }_{\frac{|k|}{2}}^2(k)\mathrm{\Gamma }_{\frac{|k|}{2}}^1=\mathrm{\Gamma }_{\frac{|k|}{2}}^0\mathrm{\Gamma }_{\frac{|k|}{2}1}^0.$$
All this becomes explicit during the following calculation of the eigenvalues of $`K_2`$.
### 5.1 Spectrum of $`K_2`$
For $`k>0`$, the highest weight vector with angular momentum
$$j=m+\frac{|k|}{2},m=0,1,\mathrm{}$$
is
$$(a_1^{})^{|k|}(a_1^{}a_2)^m.$$
Since
$$\text{gad}\mathrm{\Lambda }_8(a_1^{})^{|k|}(a_1^{}a_2)^m=|k|(a_1^{})^{|k|}(a_1^{}a_2)^m,$$
it is the highest weight vector of $`\mathrm{\Gamma }_j^2(|k|)`$. Thus $`\mathrm{\Gamma }_j^2(|k|)`$ occurs in the reduction of the $`osp(2,2)`$ action on $`S_F^{2,2}(|k|)`$.
General theory tells us the branching rules of $`\mathrm{\Gamma }_j^2(|k|)`$ as in (5). This equation is thus established for $`k>0`$.
We can check as before that
$$\epsilon _{\alpha \beta }\text{gad}\mathrm{\Lambda }_\alpha ^{}\text{gad}\mathrm{\Lambda }_\beta ^{}(a_1^{})^{|k|}(a_1^{}a_2)^m=\frac{1}{2}\left(m+\frac{|k|}{2}\right)(a_1^{})^{|k|}(a_1^{}a_2)^m$$
while
$$\frac{1}{4}(\text{gad}\mathrm{\Lambda }_8)^2(a_1^{})^{|k|}(a_1^{}a_2)^m=\frac{k^2}{4}$$
and
$$K_1|_{\mathrm{\Gamma }_j^1}=j(j+1)1\mathrm{l}.$$
We thus have
$$K_2|_{\mathrm{\Gamma }_j^2(|k|)}=\left(j^2\frac{k^2}{4}\right)1\mathrm{l}.$$
For $`k<0`$,
$$(a_2)^{|k|}(a_1^{}a_2)^m$$
is the highest weight vector for angular momentum
$$j=m+\frac{|k|}{2}.$$
Since
$$\text{gad}\mathrm{\Lambda }_8(a_2)^{|k|}(a_1^{}a_2)^m=|k|(a_2)^{|k|}(a_1^{}a_2)^m,$$
it is the highest weight vector of $`\mathrm{\Gamma }_j^2(|k|)`$. Hence $`\mathrm{\Gamma }_j^2(|k|)`$ occurs in the reduction of the $`osp(2,2)`$ action on $`S_F^{2,2}(|k|)`$. We thus establish (5) for $`k<0`$ as well.
The eigenvalues of $`V_k`$, when restricted to $`\mathrm{\Gamma }_j^1`$ and $`\mathrm{\Gamma }_{j1/2}^1`$ and for $`j\frac{|k|}{2}+1`$, are $`\pm \sqrt{j^2\frac{k^2}{4}}`$. These eigenvalues are not zero. Hence the $`osp(2,2)`$ operators which intertwine these representations, mapping vectors of one representation to the other, anticommute with $`V_k`$: they are chiral symmetries for these representations. For $`j=\frac{|k|}{2}`$, $`V_k`$ vanishes while the representation space carries the atypical representation $`\mathrm{\Gamma }_{\frac{|k|}{2}}^2(k)`$ of $`osp(2,2)`$. Hence we can say that the above chiral operators all anticommute with $`V_k|_{J=0}`$. Hence these operators anticommute with $`V_k`$ (for any $`j`$, on all vectors of $`S_F^{2.2}(k)`$) just as standard chirality anticommutes with the massless Dirac operator.
For $`k=0`$, these operators were $`\mathrm{\Lambda }_\alpha ^{}`$, $`\mathrm{\Lambda }_8`$. But they change with $`k`$. They can be worked out. They do not occur in subsequent discussion and hence we do not show them here.
We now establish that $`V_k`$ is the correct choice of the action for the fuzzy SUSY action $`S_k`$:
$$S_k=\text{const}str\mathrm{\Phi }^{}V_k\mathrm{\Phi }.$$
(5.34)
This formula is valid also for $`k=0`$ as we saw earlier. We here focus on $`k0`$.
The Dirac operators $`D`$ for fuzzy spheres of instanton number $`k`$ are known . We first show that $`V_k`$ coincides with this operator on the Dirac sector.
It is enough to focus on typical $`osp(2,2)`$ IRR’s since both the Dirac operator and $`V_k`$ vanish on the grade-odd sector of $`\mathrm{\Gamma }_{\frac{|k|}{2}}^1`$. Thus consider $`\mathrm{\Gamma }_j^2(k)`$ for $`j\frac{1}{2}|k|+1`$. Angular momentum $`J`$ in the Dirac sector of $`\mathrm{\Gamma }_j^2(k)`$ is $`j1/2`$. Hence
$$V_k^2|_{\mathrm{\Gamma }_j^2(k)Diracsector}=\left(J\frac{|k|1}{2}\right)\left(J+\frac{|k|+1}{2}\right)1\mathrm{l}.$$
Substituting $`J=n+\frac{|k|1}{2}`$ and identifying $`|k|=2T`$, we get the answer of :
$$V_k^2|_{\mathrm{\Gamma }_j^2(k)Diracsector}=n(n+2T)1\mathrm{l}.$$
Hence $`V_k^2|_{\mathrm{\Gamma }_j^2(k)Diracsector}`$ is the correct Dirac operator .
This result and the $`osp(2,1)`$-invariance of $`V_k`$ are compelling reasons to identify it as the SUSY generalization of the Dirac and Laplacian operators for $`k0`$.
### 5.2 Index Theory and Zero Modes
There is also further evidence supporting the correctness of $`V_k`$: It gives the SUSY generalization of index theory.
Thus one knows that 1) the Dirac operator has $`|k|`$ zero modes for instanton number $`k`$ on $`S^2`$ and on the fuzzy sphere $`S_F^2(k)`$, and 2) they are left- (right-) chiral if $`k>0`$ ($`k<0`$), 3) charge conjugation interchanges these chiralities.
More precisely if $`n_{L,R}`$ are the number of left- and right-chiral zero modes,
$$n_Ln_R=k.$$
This number is “topologically stable”. The meaning of this statement in the fuzzy case can be found in .
If the Dirac operator is $`SU(2)`$-invariant, these zero modes organize themselves into $`SU(2)`$ multiplets with angular momentum $`\frac{|k|}{2}`$ .
Now $`V_k`$ has zero modes which form the atypical multiplet $`\mathrm{\Gamma }_{\frac{|k|}{2}}^2(k)`$ of $`osp(2,2)`$. The number of zero modes is $`2|k|+1`$. Of these, $`|k|`$ correspond to the grade odd sector and can be identified with the zero modes of $`S^2`$ and $`S_F^2(k)`$ Dirac operators. The remaining grade even ($`|k|+1`$) zero modes are their SUSY-partners.
The zero modes transform by inequivalent IRR’s of $`osp(2,2)`$ for the two signs of $`k`$. These two atypical $`osp(2,2)`$ representations are SUSY generalizations of left- and right- chiralities.
Identifying charge conjugation with the automorphism (1.12), we see that it exchanges these two IRR’s just as it exchanges chiralities in the Dirac sector.
## 6 Final Remarks
In this paper, we have extended the work of on the fuzzy SUSY model on $`S^2`$ to the instanton sector. A SUSY generalization of index theory and zero modes of the Dirac operator has also been established.
Following , we can try introducing interactions involving just $`\mathrm{\Phi }`$. For $`k0`$, $`\mathrm{\Phi }`$ can be thought of as a rectangular matrix. So $`\mathrm{\Phi }^{}\mathrm{\Phi }`$ and $`\mathrm{\Phi }\mathrm{\Phi }^{}`$ are square matrices of different sizes acting on $`osp(2,2)`$ representations with $`N=n`$ and $`N=n+k`$. A typical interaction may then be
$$\lambda str(\mathrm{\Phi }^{}\mathrm{\Phi })^2$$
(6.35)
where $`str`$ is over the space with $`N=n`$, the domain of $`\mathrm{\Phi }^{}\mathrm{\Phi }`$. (But note that (6.35) and the use of $`str`$ in interactions require further study.)
Fuzzy SUSY gauge theories remain to be formulated. The investigation of the graded commutative limit $`n\mathrm{}`$ with $`k`$ fixed has also not been done for $`k0`$.
Numerical simulations on fuzzy SUSY models are being initiated.
Acknowledgement
We are part of a collaboration on fuzzy physics and have benefited by inputs from its members. We are particulary grateful to Marco Panero for carefully reading the paper and suggesting several corrections. This work was supported in part by the DOE grant DE-FG02-85ER40231 and by NSF under contract number INT9908763.
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# Lepton flavor changing Higgs boson decays in SUSY with 𝜈_𝑅
## 1 Introduction
The observed neutrino masses do require a theoretical framework beyond the Standard Model of Particle Physics with just three massless left-handed neutrinos. Within the MSSM-seesaw context, which will be adopted here, the MSSM particle content is enlarged by three right handed neutrinos plus their corresponding supersymmetric (SUSY) partners, and the neutrino masses are generated by the seesaw mechanism. Three of the six resulting Majorana neutrinos have light masses, $`m_{\nu _i},i=1,2,3`$, and the other three have heavy masses, $`m_{N_i},i=1,2,3`$. These physical masses are related to the Dirac mass matrix $`m_D`$, the right-handed neutrino mass matrix $`m_M`$, and the unitary matrix $`U_{MNS}`$ by $`diag(m_{\nu _1},m_{\nu _2},m_{\nu _3})U_{MNS}^T(m_Dm_M^1m_D^T)U_{MNS}`$ and $`diag(m_{N_1},m_{N_2},m_{N_3})m_M`$, respectively. Here we have chosen an electroweak eigenstate basis where $`m_M`$ and the charged lepton mass matrix are flavor diagonal, and we have assumed that all elements in $`m_D=Y_\nu <H_2>`$, where $`Y_\nu `$ is the neutrino Yukawa coupling matrix and $`<H_2>=v\mathrm{sin}\beta `$ ($`v=174`$ GeV), are much smaller than those of $`m_M`$. The two previous relations can be rewritten together in a more convenient form for the work presented here as, $`m_D^T=im_N^{diag\mathrm{\hspace{0.17em}1}/2}Rm_\nu ^{diag\mathrm{\hspace{0.17em}1}/2}U_{MNS}^+`$, where $`R`$ is a general complex and orthogonal $`3\times 3`$ matrix, which will be parameterized by three complex angles $`\theta _i`$, $`i=1,2,3`$.
One of the most interesting features of the MSSM-seesaw model is the associated rich phenomenology due to the occurrence of lepton flavor violating (LFV) processes. Whereas in the standard (non-SUSY) seesaw models the ratios of LFV processes are small due to the smallness of the light neutrino masses, in the SUSY-seesaw models these can be large due to an important additional source of lepton flavor mixing in the soft-SUSY-breaking terms. Even in the scenarios with universal soft-SUSY-breaking parameters at the large energy scale associated to the SUSY breaking $`M_X`$, the running from this scale down to $`m_M`$ induces, via the neutrino Yukawa couplings, large lepton flavor mixing in the slepton soft masses, and provides the so-called slepton-lepton misalignment, which in turn generates non-diagonal lepton flavor interactions. These interactions can induce sizable ratios in several LFV processes with SM charged leptons in the external legs, which are actually being tested experimentally with high precision and therefore provide a very interesting window to look for indirect SUSY signals. They can also induce important contributions to other LFV processes that could be meassured in the next generation colliders, as it is the case of the MSSM Higgs boson decays into $`\tau \overline{\mu }`$, $`\tau \overline{e}`$ and $`\mu \overline{e}`$ which are the subject of our interest.
Here we compute the partial widths for these lepton flavor violating Higgs boson decays (LFVHD) to one-loop order and analyze numerically the corresponding branching ratios in terms of the mSUGRA and seesaw parameters, namely, $`M_0`$, $`M_{1/2}`$, $`\mathrm{tan}\beta `$, $`m_{N_i}`$ and $`R`$. For the one loop running of the parameters we use the mSUSPECT programme. We analyze in parallel the lepton flavor changing $`ł_jl_i\gamma `$ ($`ij`$) decays and explore the maximum predicted rates for LFVHD, mainly for $`H^0,A^0\tau \overline{\mu }`$ decays, by requiring compatibility with $`BR(ł_jl_i\gamma )`$ data. For these we use the present experimental upper bounds given by $`^\mathrm{?}`$ $`|BR(\mu e\gamma )|<1.2\times 10^{11}`$, $`|BR(\tau \mu \gamma )|<3.1\times 10^7`$ and $`|BR(\tau e\gamma )|<2.7\times 10^6`$.
For the numerical analysis we choose, $`M_X=2\times 10^{16}`$ GeV and $`A_0=0`$. The $`U_{MNS}`$ matrix elements and the $`m_{\nu _i}`$ are fixed to the most favored values by neutrino data $`^\mathrm{?}`$ with $`\sqrt{\mathrm{\Delta }m_{sol}^2}=0.008`$ eV, $`\sqrt{\mathrm{\Delta }m_{atm}^2}=0.05`$ eV, $`\theta _{12}=\theta _{sol}=30^o`$, $`\theta _{23}=\theta _{atm}=45^o`$, $`\theta _{13}=0^o`$ and $`\delta =\alpha =\beta =0`$. We consider two plaussible scenarios, one with quasi-degenerate light and degenerate heavy neutrinos and with $`m_{\nu _1}=0.2eV,m_{\nu _2}=m_{\nu _1}+\frac{\mathrm{\Delta }m_{sol}^2}{2m_{\nu _1}},m_{\nu _3}=m_{\nu _1}+\frac{\mathrm{\Delta }m_{atm}^2}{2m_{\nu _1}}`$ and $`m_{N_1}=m_{N_2}=m_{N_3}=m_N`$; and the other one with hierarchical light and hierarchical heavy neutrinos, and with $`m_{\nu _1}0eV,m_{\nu _2}=\sqrt{\mathrm{\Delta }m_{sol}^2},m_{\nu _3}=\sqrt{\mathrm{\Delta }m_{atm}^2}`$ and $`m_{N_1}m_{N_2}<m_{N_3}`$.
This is a reduced version of our more complete work $`^\mathrm{?}`$ to which we address the reader for more details.
## 2 Numerical results and conclusions
We show in figs. (1) through (4) the numerical results for the branching ratios of the LFVHD together with the branching ratios for the relevant $`l_jl_i\gamma `$ decays. The results of $`BR(H_0\tau \overline{\mu })`$ as a function of $`m_N`$, for degenerate heavy neutrinos and real R, are illustrated in fig. (1a), for several $`\mathrm{tan}\beta `$ values, $`\mathrm{tan}\beta =3,10,30,50`$. Notice that in this case, the rates do not depend on $`R`$. The explored range in $`m_N`$ is from $`10^8`$ GeV up to $`10^{14}`$ GeV which is favorable for baryogenesis. We also show in this figure, the corresponding predicted rates for the most relevant lepton decay, which in this case is $`\mu e\gamma `$, and include its upper experimental bound. We have checked that the other lepton decay channels are well within their experimental allowed range. The ratios for $`A_0`$ decays, not shown here for brevity, are very similar to those for $`H_0`$ decays in all the studied scenarios in this work. We have also found that the ratios for the light Higgs boson, $`h_0`$, behave very similarly with $`m_N`$ and $`\mathrm{tan}\beta `$ but are smaller than the heavy Higgs ones in about two orders of magnitude. From our results we learn about the high sensitivity to $`\mathrm{tan}\beta `$ of the LFVHD rates for all Higgs bosons which, at large $`\mathrm{tan}\beta `$, scale roughly as $`(\mathrm{tan}\beta )^4`$, in comparison with the lepton decay rates which scale as $`(\mathrm{tan}\beta )^2`$. The dependence of both rates on $`m_N`$ is that expected from the mass insertion approximation, where $`BR(H_xl_j\overline{l}_i)`$, $`BR(l_jl_i\gamma )|m_N\mathrm{log}(m_N)|^2`$. We find that the largest ratios, which are for $`H_0`$ and $`A_0`$, are in any case very small, at most $`10^{10}`$ in the region of high $`\mathrm{tan}\beta `$ and high $`m_N`$. Besides, the rates for $`\mu e\gamma `$ decays are below the upper experimental bound for all explored $`\mathrm{tan}\beta `$ and $`m_N`$ values. The branching ratios for the Higgs boson decays into $`\tau \overline{e}`$ and $`\mu \overline{e}`$ are much smaller than the $`\tau \overline{\mu }`$ ones, as expected, and we do not show plots for them. For instance, for $`m_N=10^{14}`$ GeV, and $`\mathrm{tan}\beta =50`$ we find $`BR(H^{(x)}\tau \overline{\mu })/BR(H^{(x)}\tau \overline{e})=4\times 10^3`$ and $`BR(H^{(x)}\tau \overline{\mu })/BR(H^{(x)}\mu \overline{e})=1.2\times 10^6`$ for the three Higgs bosons.
The case of hierarchical neutrinos gives clearly larger LFV rates than the degenerate case, as can be seen in figs. (2), (3) and (4). However, we will get restrictions on the maximum allowed Higgs decay rates coming from the experimental lepton decay bounds. For instance, the case of real $`\theta _1`$, that is illustrated in fig. (2) shows that compatibility with $`\mu e\gamma `$ data occurs only in the very narrow deeps at around $`\theta _1=0`$, $`1.9`$ and $`\pi `$. Notice that it is precisely at the points $`\theta _1=0,\pi `$ where the $`BR(H_0,A_0\tau \overline{\mu })`$ rates reach their maximum values, although these are not large, just about $`10^8`$. We have checked that for lower $`\mathrm{tan}\beta `$ values, the allowed regions in $`\theta _1`$ widen and are placed at the same points, but the corresponding maximum values of the LFVHD rates get considerably reduced. For the alternative case, not shown here, of real $`\theta _20`$, with $`\theta _1=\theta _3=0`$ we get a similar behaviour of $`BR(H_x\tau \overline{\mu })`$ with $`\theta _2`$ than with $`\theta _1`$, and the maximum values of about $`10^8`$ are now placed at $`\theta _2=0`$, $`\pi `$. In contrast, $`BR(\mu e\gamma )`$ is constant with $`\theta _2`$ and reach very small values, well below the experimental bound. In particular, for $`\mathrm{tan}\beta =50`$, $`M_0=400`$ GeV and $`M_{1/2}=300`$ GeV it is $`10^{19}`$. Regarding the dependence with $`\theta _3`$, not shown here either, a reverse situation is found, where $`BR(H_x\tau \overline{\mu })`$ is approximately constant and, for the heavy Higgs bosons, it is around $`10^8`$. On the contrary, $`BR(\mu e\gamma )`$ varies but it is always well below the experimental upper bound. In addition, we have checked that the $`BR(\tau \mu \gamma )`$ and $`BR(\tau e\gamma )`$ rates are within the experimental allowed range in all cases. In conclusion, for real R we find that the maximum allowed LFVHD rates are at or below $`10^8`$.
The case of complex $`R`$ is certainly more promissing. The examples illustrated in figs. (3) and (4) are for the most favorable case, among the ones studied here, of complex $`\theta _20`$ with $`\theta _1=\theta _3=0`$ and show that considerably larger $`BR(H_x\tau \overline{\mu })`$ rates than in the real $`R`$ case are found. Regarding the dependence with $`\theta _2`$, we find that for the explored values with $`(|\theta _2|,Arg(\theta _2))(3.5,1)`$, the Higgs rates grow with both $`|\theta _2|`$ and $`Arg(\theta _2)`$ and, for the selected values of the MSSM-seesaw parameters in fig. (3), they reach values up to around $`5\times 10^5`$. We have checked that the predicted rates for $`BR(\mu e\gamma )`$ are well below the experimental upper bound, being nearly constant with $`\theta _2`$ and around $`10^{19}`$. Similarly, for the $`\tau e\gamma `$ decay. Notice that the smallness of these two decays, in the case under study of $`\theta _20`$, is not maintained if our hypothesis on $`\theta _{13}=0`$ is changed. For instance, for $`\theta _{13}=5^o`$, which is also allowed by neutrino data, we get $`BR(\mu e\gamma )1.8\times 10^8`$ well above the experimental upper bound. Therefore, in this case of complex $`\theta _20`$, the relevant lepton decay is $`\tau \mu \gamma `$ which is illustrated in figs. (3) and (4) together with its experimental bound. For the set of parameters chosen in fig. (3), we get that the allowed region by $`\tau \mu \gamma `$ data of the $`(|\theta _2|,Arg(\theta _2))`$ parameter space implies a reduction in the Higgs rates, leading to a maximum allowed value of just $`5\times 10^8`$.
The dependence of the LFV ratios with $`M_0`$ and $`M_{1/2}`$ for hierarchical neutrinos are shown in fig. (4). We see clearly the different behaviour of the LFVHD and the lepton decays with these parameters, showing the first ones a milder dependence. This implies, that for large enough values of $`M_0`$ or $`M_{1/2}`$ or both the $`BR(\tau \mu \gamma )`$ rates get considerably suppresed, due to the decoupling of the heavy SUSY particles in the loops, and enter into the allowed region by data, whereas the $`BR(H_0\tau \overline{\mu })`$ rates are not much reduced. In fact, we see in figs. (4e) and (4f) that for the choice $`M_0=M_{1/2}`$ the $`\tau `$ decay ratio crosses down the upper experimental bound at around $`M_0=1100`$ GeV whereas the Higgs decay ratio is still quite large $`4\times 10^6`$ in the high $`M_0`$ region, around $`M_02000`$ GeV. This behaviour is a clear indication that the heavy SUSY particles in the loops do not decouple in the LFVHD. Notice also that it can be reformulated as non-decoupling in the effective $`H^{(x)}\tau \mu `$ couplings and these in turn can induce large contributions to other LFV processes that are mediated by Higgs exchange as, for instance, $`\tau \mu \mu \mu `$. However, we have checked that for the explored values in this work of $`M_0`$, $`M_{1/2}`$, $`\mathrm{tan}\beta `$, R and $`m_{N_i}`$ that lead to the anounced LFVHD ratios of about $`4\times 10^6`$, the corresponding $`BR(\tau \mu \mu \mu )`$ rates are below the present experimental upper bound.
In summary, after exploring the dependence of the LFVHD rates with all the involved MSSM-seesaw parameters, and by requiring compatibility with data of the correlated predictions for $`\mu e\gamma `$, $`\tau e\gamma `$ and $`\tau \mu \gamma `$ decays, we find that $`BR(H_0,A_0\tau \overline{\mu })`$ as large as $`10^5`$, for hierarchical neutrinos and large $`M_{SUSY}`$ in the TeV range can be reached. These rates are close but still below the expected future experimental reach of about $`10^4`$ at the LHC and next generation linear colliders.
## Acknowledgments
A. M. Curiel wishes to thank the organizers of this conference for a very fruitful, interesting and enjoyable meeting.
## References
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# Regular homotopy classes of singular maps
## 1. Introduction
This paper is a sequel to in which locally generic maps of closed surfaces into $`^3`$ were classified up to regular homotopy. It turned out that for maps with at least one cross-cap point the number of singular points was the only regular homotopy invariant. The obstruction to constructing a regular homotopy was destroyed by pushing $`1`$-cells of the surface through a singular point (see Figure 1).
In the present work we extend this result for locally generic maps of closed $`n`$-manifolds into $`^{2n1}`$ in case $`n3`$ (Theorem 1.7). However, this is not simply an adaptation of the ideas of . The idea of pushing $`(n1)`$-simplices through a singular point originates from the previous work, but a lot of technical problems have to be dealt with. For example, if $`n`$ is odd the Whitney-umbrella points have signs and the $`(n1)`$-dimensional obstruction is $``$-valued. Moreover, a new type of obstruction appears if $`n>3`$ that was not present in the case $`n=2`$. This is an $`n`$-dimensional obstruction and is related to the double point set of the map. In eliminating this obstruction we make essential use of a result of T. Ekholm .
It seems that this second type of obstruction cannot be eliminated if $`n=3`$. In fact, in addition to the number of singular points a new $``$-valued invariant comes into the picture. I intend to deal with this in a separate paper.
### 1.1. Preliminary definitions
###### Definition 1.1.
Let $`M^n`$ be a closed $`n`$-manifold and $`N^{2n1}`$ an arbitrary $`(2n1)`$-manifold . A map $`f:M^nN^{2n1}`$ is called *locally generic* if it is an immersion except for cross-cap (or Whitney-umbrella) singularities. Thus for any singular point $`p`$ there exist coordinate systems $`(x_1,\mathrm{},x_n)`$ about $`p`$ and $`(y_1,\mathrm{},y_{2n1})`$ about $`f(p)`$ such that $`f`$ is given near $`p`$ by
(1.1)
$$y_1=x_1^2,y_i=x_i,y_{n+i1}=x_1x_i(i=2,\mathrm{},n).$$
###### Notation 1.2.
Let $`L(M^n,N^{2n1})`$ denote the subspace of locally generic maps, $`G(M^n,N^{2n1})`$ the subspace of generic maps and $`\text{Imm}(M^n,N^{2n1})`$ the subspace of immersions in $`C^{\mathrm{}}(M^n,N^{2n1})`$ endowed with the $`C^{\mathrm{}}`$ topology. We introduce the notation $`L^{}(M^n,N^{2n1})`$ for the space
$$L(M^n,N^{2n1})\text{Imm}(M^n,N^{2n1}).$$
For $`fL(M^n,N^{2n1})`$ the set of singular points of the map $`f`$ is denoted by $`S(f)`$.
###### Remark 1.3.
Since $`M^n`$ is compact and the singularities of $`f`$ are isolated we have $`\left|S(f)\right|<\mathrm{}`$. It is well known that for $`M^n`$ closed $`\left|S(f)\right|`$ is an even number.
For $`n>2`$ a map $`f:M^nN^{2n1}`$ is generic iff
* it is an immersion with normal crossings except in a finite set of points
* the singular points are non-multiple cross-cap points
* $`f`$ has at most double crossings.
For $`n=2`$ a generic $`f`$ might also have triple points.
###### Definition 1.4.
A *regular homotopy* is a path in the space $`L(M^n,N^{2n1})`$. In other words, it is a smooth map $`H:M^n\times [0,1]N^{2n1}`$ such that $`H_tL(M^n,N^{2n1})`$ for every $`t[0,1]`$. Here $`H_t(x)=H(x,t)`$ for $`xM^n`$ and $`t[0,1]`$. The regular homotopy $`H`$ is called *singularity fixing* if $`S(H_t)=S(H_0)`$ for every $`t[0,1]`$.
###### Definition 1.5.
Two maps $`f,gL(M^n,N^{2n1})`$ are called *regularly homotopic* if there exists a regular homotopy $`H`$ such that $`H_0=f`$ and $`H_1=g`$. If $`H`$ can be chosen to be singularity fixing then we say that $`f`$ and $`g`$ are *regularly homotopic through a singularity fixing homotopy*.
###### Notation 1.6.
The fact that $`f`$ and $`g`$ are regularly homotopic is denoted by $`fg`$. Furthermore, $`f_sg`$ denotes that $`f`$ and $`g`$ are regularly homotopic through a singularity fixing homotopy.
### 1.2. The classification theorem
###### Theorem 1.7.
Suppose that $`n>3`$ or $`n=2`$. Let $`M^n`$ be a closed $`n`$-manifold and $`f,gL^{}(M^n,^{2n1})`$. Then
$$fg\left|S(f)\right|=\left|S(g)\right|.$$
## 2. Known results used in the proof of Theorem 1.7
### 2.1. Smale’s lemma and $`M^n`$-regular homotopies
We will use extensively the following result of Smale (this is Theorem 1.1 in , the original paper of Smale is ). First we need a few definitions.
###### Definition 2.1.
Let $`=_{k,m}`$ be the space of all $`C^{\mathrm{}}`$ immersions of $`D^k`$ in $`^m`$ in the $`C^1`$ topology. Let $`=_{k,m}`$ be the set of pairs $`(g,g^{})`$ where $`g:S^{k1}^m`$ is a $`C^{\mathrm{}}`$ immersion and $`g^{}:S^{k1}T^m`$ is a $`C^{\mathrm{}}`$ transversal field of $`g`$. $``$ is topologized as a subspace of Cartesian product of the space of immersions $`S^{k1}^m`$, in the $`C^1`$ topology, with the space of continuous maps $`S^{k1}T^m`$, in the compact-open topology.
If $`h`$ then let $`h^{}:S^{k1}T^m`$ be defined by $`h^{}(x)=`$ derivative of $`h`$ along the radius at $`xS^{k1}`$. I.e., if $`r(x)`$ is the unit tangent vector of $`D`$ that is normal to $`S^{k1}`$ at $`x`$ and which points away from the origin then $`h^{}(x)=h_{}r(x)`$. We define the map $`\pi :`$ by the formula $`\pi (h)=(h|S^{k1},h^{})`$. It is clear that $`\pi `$ is continuous. The following theorem is Theorem 1.1 in .
###### Theorem 2.2.
If $`k<m`$, then $`\pi :_{k,m}_{k,m}`$ has the covering homotopy property.
The intuitive content of this theorem is as follows: If we are given an immersed disk $`D^k`$ in $`^m`$ such that $`k<m`$ and we deform the boundary of the disk and the normal derivatives along the boundary, then we can deform the whole disk at the same time so as to induce the given deformation on the boundary and normal derivatives.
The following definitions were introduced by Hirsch in .
###### Definition 2.3.
Let $`A`$ be an arbitrary subset of the manifold $`M^n`$. Let $`h:AQ^q`$ and $`h^{}:TM^n|ATQ^q`$ ($`Q^q`$ is a manifold) be continuous maps such that $`h^{}`$ covers $`h`$. The pair $`(h,h^{})`$ is called an *$`M^n`$-regular map*, or *$`M^n`$-immersion*, *of $`A`$ in $`Q^q`$* if the following condition is satisfied: there is a neighborhood $`V`$ of $`A`$ in $`M^n`$ and an immersion $`l:VQ^q`$ such that $`dl|(TM|A)=h^{}`$. It follows that $`l|A=h`$. We say that $`(h,h^{})`$ is $`C^k`$ if $`l`$ can be chosen to be $`C^k`$.
###### Definition 2.4.
Let $`BAM^n`$ be subsets. If $`(r,r^{}),(s,s^{}):AQ^q`$ are $`M^n`$-immersions such that $`r|B=s|B`$ and $`r^{}|(TM^n|B)=s^{}|(TM^n|B)`$, we say that *$`(r,r^{})`$ and $`(s,s^{})`$ are tangent on $`B`$*, and write this as $`(r,r^{})|B=(s,s^{})|B`$.
###### Definition 2.5.
Let $`(r,r^{})`$ and $`(s,s^{})`$ be $`M^n`$-immersions of $`A`$ in $`Q^q`$ such that for a certain (possibly empty) subset $`B`$ of $`A`$, $`(r,r^{})`$ and $`(s,s^{})`$ are tangent on $`B`$. We say that $`(r,r^{})`$ and $`(s,s^{})`$ are *$`M^n`$-regularly homotopic (rel B)* if there is a path $`(h_t,h_t^{})`$ in the space of all $`M^n`$-immersions of $`A`$ in $`Q^q`$ joining $`(r,r^{})`$ to $`(s,s^{})`$, such that for each $`t`$, $`(h_t,h_t^{})|B=(r,r^{})|B`$. Such a path is called an *$`M^n`$-regular homotopy (rel B)*, and it is $`C^k`$ if every $`(h_t,h_t^{})`$ is $`C^k`$.
###### Notation 2.6.
The space of all $`C^{\mathrm{}}`$ $`^n`$-immersions of $`D^k`$ in $`^q`$ is denoted by $`(k,q;n)`$; the space of all $`C^{\mathrm{}}`$ $`^n`$-immersions of $`S^{k1}`$ in $`^q`$ is denoted by $`^{}(k,q;n)`$. Now let $`\pi _n:(k,q;n)^{}(k,q;n)`$ be defined by
$$\pi _n(f,f^{})=(f|S^{k1},f^{}|(T^n|S^{k1})).$$
If $`n=k`$, this is the map $`\pi :_{k,q}_{k,q}`$ defined in Definition 2.1.
For $`(f,f^{})(k,q;n)`$ put $`\mathrm{\Gamma }_n(f,f^{})=\pi _n^1(\pi _n(f,f^{}))`$.
The following statement is Theorem 3.5 in . It is a generalization of Smale’s lemma to $`M^n`$-immersions. To avoid confusion we will refer to it also as Smale’s lemma (despite the fact that it was proved by Hirsch).
###### Theorem 2.7.
$`\pi _n`$ has the covering homotopy property if $`k<q`$.
Theorem 3.2 in gives an alternative description of $`M^n`$-immersions:
###### Notation 2.8.
$`\{e_1(x),\mathrm{},e_n(x)\}`$ denotes the standard basis of $`T_x^n`$.
###### Lemma 2.9.
There is a homeomorphism $`\zeta `$ between the space of $`^n`$-immersions $`(h,h^{}):D^k^q`$ and the space of pairs $`(l,\mathrm{\Psi })`$, where $`l\text{Imm}(D^k,^q)`$ and $`\mathrm{\Psi }`$ is a transversal $`(nk)`$-field along $`l`$. The homeomorphism is given by $`\zeta (h,h^{})=(h,\mathrm{\Psi })`$ where $`\mathrm{\Psi }(x)=h^{}\{e_{k+1}(x),\mathrm{},e_n(x)\}`$. Moreover, $`(f,f^{})`$ is $`C^k`$ if and only if $`f`$ and $`\mathrm{\Psi }`$ are $`C^k`$.
An analogous result holds for $`S^{k1}`$ (see , Theorem 3.3):
###### Notation 2.10.
Denote the space of pairs $`(l,\mathrm{\Psi })`$, where $`l\text{Imm}(S^{k1},^q)`$ and $`\mathrm{\Psi }`$ is a transversal $`(nk+1)`$-field along $`l`$ by
$$\text{Imm}_{nk+1}(S^{k1},^q).$$
###### Lemma 2.11.
There is a homeomorphism
$$\chi :^{}(k,q;n)\text{Imm}_{nk+1}(S^{k1},^q).$$
$`\chi `$ is given as follows: Let $`\mathrm{\Phi }`$ be the normal $`(nk+1)`$-field on $`S^{k1}`$ given by
$$\mathrm{\Phi }(x)=\{r(x),e_{k+1}(x),\mathrm{},e_n(x)\},$$
where $`r(x)`$ is the outward unit normal to $`S^{k1}`$ in $`^k`$. Then $`\chi (h,h^{})=(h,\mathrm{\Psi })`$, where $`\mathrm{\Psi }(x)=f^{}(\mathrm{\Phi }(x))`$.
### 2.2. The obstructions $`\tau `$ and $`\mathrm{\Omega }`$
Hirsch defined in an invariant $`\tau (g^{})`$ for each $`(g,g^{})^{}(k,q;n)`$. The vanishing of $`\tau (g^{})`$ implies that $`(g,g^{})`$ comes from $`(k,q;n)`$.
###### Definition 2.12.
$`(f,f^{})^{}(k,q;n)`$ is said to be *extendible* if there is a $`(g,g^{})(k,q;n)`$ such that $`\pi _q(g,g^{})=(f,f^{})`$.
###### Definition 2.13.
Let $`(f,f^{}):S^{k1}^n`$ be a $`C^{\mathrm{}}`$ $`^n`$-immersion, i.e., $`(f,f^{})^{}(k,q;n)`$. The *obstruction to extending $`(f,f^{})`$*, denoted by $`\tau (f^{})\pi _{k1}(V_{q,n})`$ is the homotopy class of the map $`S^{k1}V_{q,n}`$ defined by
$$xf^{}\{e_1(x),\mathrm{},e_n(x)\}.$$
The following lemma is Theorem 3.9 in :
###### Lemma 2.14.
If $`k<q`$ and $`\tau (f^{})=0`$ then $`(f,f^{})`$ is extendible.
###### Definition 2.15.
Let $`\mathrm{\Phi }_n:(k,q;n)C^0(D^k,V_{q,n})`$ be as follows:
$$\mathrm{\Phi }_n(f,f^{})(x)=f^{}\{e_1(x),\mathrm{},e_n(x)\}.$$
Let $`(f,f^{}),(g,g^{})(k,q;n)`$, with $`\pi _n(f,f^{})=\pi _n(g,g^{})`$, so that $`(f,f^{})\mathrm{\Gamma }_n(g,g^{})`$. Then $`\mathrm{\Phi }_n(f,f^{})`$ and $`\mathrm{\Phi }_n(g,g^{})`$ are maps $`D^kV_{q,n}`$ which are tangent on $`S^{k1}`$.
###### Definition 2.16.
Let $`A`$ be a topological space, simple in dimension $`k`$. Let $`f,g:D^kA`$ and assume that $`f(x)=g(x)`$ if $`xS^{k1}`$. Then $`d(f,g)\pi _k(A)`$ is represented by mapping the ”top” hemisphere of $`S^k`$ by $`f`$ and the ”bottom” one by $`g`$, assuming that the orientation of $`S^k`$ is given by the coordinate frame $`\{e_1,\mathrm{},e_k\}`$ at the ”North” pole of $`S^k`$.
###### Definition 2.17.
$$\mathrm{\Omega }(f^{},g^{})=d(\mathrm{\Phi }_n(f,f^{}),\mathrm{\Phi }_n(g,g^{}))\pi _k(V_{q,n})$$
is called the *obstruction to an $`^n`$-regular homotopy* (rel $`S^{k1}`$) between $`(f,f^{})`$ and $`(g,g^{})`$. (This is well defined since $`V_{q,n}`$ is simple in all dimensions.)
###### Remark 2.18.
An explicit definition of $`\mathrm{\Omega }(f^{},g^{})`$ is as follows: identify the upper and lower hemispheres of $`S^k`$ with $`D^k`$. Let $`\omega :S^kV_{q,n}`$ be the map $`\omega (x)=f^{}\{e_1(x),\mathrm{},e_n(x)\}`$ if $`x`$ is in the upper hemisphere, $`\omega (x)=g^{}\{e_1(x),\mathrm{},e_n(x)\}`$ if $`x`$ is in the lower hemisphere. $`\omega (x)`$ is well defined on the equator because $`(f,f^{})`$ and $`(g,g^{})`$ agree on $`S^{k1}`$. Then $`\mathrm{\Omega }(f^{},g^{})`$ is the homotopy class of $`\omega `$.
The following lemma is Theorem 4.3 in :
###### Lemma 2.19.
Let $`(f,f^{}),(g,g^{}),(h,h^{})`$ be $`C^{\mathrm{}}`$ $`^n`$-immersions of $`\mathrm{\Delta }^k`$ in $`^q`$ which are all tangent on $`\mathrm{\Delta }^k`$. Then
(a) $`\mathrm{\Omega }(f^{},g^{})+\mathrm{\Omega }(g^{},h^{})=\mathrm{\Omega }(f^{},h^{})`$.
(b) $`\mathrm{\Omega }(f^{},f^{})=0.`$
(c) Given $`\alpha \pi _k(V_{q,n})`$ there exists $`(g,g^{})`$ such that $`\mathrm{\Omega }(f^{},g^{})=\alpha `$.
(d) Suppose that $`\mathrm{\Omega }(f^{},g^{})=0`$ and $`k<q`$. Let $`F:\mathrm{\Delta }^k\times IV_{q,n}`$ be a homotopy (rel $`\mathrm{\Delta }^k`$) between the maps $`F_0,F_1:\mathrm{\Delta }^kV_{q,n}`$ defined respectively by $`xf^{}\{e_i(x)\}`$ and $`xg^{}\{e_i(x)\}`$, $`i=1\mathrm{},n`$. Then there is a $`C^{\mathrm{}}`$ $`^n`$-regular homotopy $`(f_t,f_t^{})`$ between $`(f,f^{})`$ and $`(g,g^{})`$ such that the map $`\mathrm{\Delta }^k\times IV_{q,n}`$ defined by $`(x,t)f_t^{}(e_i(x))`$ is homotopic (rel $`(\mathrm{\Delta }^k\times I)(\mathrm{\Delta }^k\times I)`$) to $`F`$.
From the proof of Theorem 3.9 in we can deduce the following
###### Theorem 2.20.
Suppose that the $`C^{\mathrm{}}`$ $`^n`$-immersions $`(f,f^{}),(g,g^{}):\mathrm{\Delta }^n^q`$ are tangent on $`(\mathrm{\Delta }^n)\text{int}(\mathrm{\Delta }^{n1})`$, where $`\mathrm{\Delta }^{n1}`$ is an $`(n1)`$-face of the $`n`$-simplex $`\mathrm{\Delta }^n`$. Then
$$\mathrm{\Omega }(f^{}|\mathrm{\Delta }^{n1},g^{}|\mathrm{\Delta }^{n1})=\tau (f^{})\tau (g^{}).$$
### 2.3. Hirsch’s theorem and lemma
In Hirsch proves the following theorem that reduces regular homotopy classification of immersions to homotopy theory.
###### Theorem 2.21.
Suppose that $`nq`$ and let $`M^n`$ and $`Q^q`$ be manifolds such that $`M^n`$ is open if $`n=q`$. Then the natural map
$$\text{Imm}(M^n,Q^q)\text{Mono}(TM,TQ)$$
is a weak homotopy equivalence (so it induces a bijection in $`\pi _0`$).
The relative version of Hirsch’s theorem also holds:
###### Theorem 2.22.
Let $`M^n`$ and $`Q^q`$ be as in Theorem 2.21. Suppose that $`KM`$ is compact and $`V_K`$ is a neighborhood of $`K`$. Given an immersion $`h\text{Imm}(V_K,Q^q)`$ with differential
$$dh=\mathrm{\Phi }:TM|V_KTQ|V_K,$$
there exists an immersion $`l\text{Imm}(M^n,Q^q)`$ such that $`l|K=f`$ and $`dl`$ is homotopic to $`\mathrm{\Phi }`$.
As a consequence of Theorem 2.21 we obtain Hirsch’s lemma:
###### Lemma 2.23.
Let $`k1`$ and $`f\text{Imm}(M^n,^{n+k+1})`$ be an immersion with normal bundle $`\mu ^k\epsilon ^1`$, where $`\epsilon ^1`$ denotes the trivial line bundle over $`M^n`$. Let $`i`$ denote the standard embedding $`^{n+k}^{n+k+1}`$. Then there exists an immersion $`g\text{Imm}(M^n,^{n+k})`$ with normal bundle $`\mu ^k`$ such that $`fig`$. During the regular homotopy a trivialization of the $`\epsilon ^1`$ component of $`\nu (f^{n+k+1})`$ can be deformed simultaneously in the space of normal fields to be finally vertical along $`ig`$.
###### Remark 2.24.
The relative version of Hirsch’s lemma also holds, as can be seen by using the relative version of Hirsch’s theorem.
###### Corollary 2.25.
If $`k2`$ then Lemma 2.23 gives a bijection
$$b:\pi _0(\text{Imm}_1(M^n,^{n+k+1}))\pi _0(\text{Imm}(M^n,^{n+k})).$$
### 2.4. Whitney-umbrellas
In this chapter we sum up some ideas from . Fix an orientation for $`^{2n1}`$.
###### Definition 2.26.
Suppose that $`N^n`$ is an $`n`$-manifold with boundary and that $`f:N^n^{2n1}`$ is a generic map (thus $`f|N^n`$ is an embedding). Let $`\nu :N^nT^{2n1}`$ be a transversal vector field along $`f|N^n`$ such that $`\nu `$ points into $`f(N^n)`$. Then for $`\epsilon `$ sufficiently small define $`_f(N)`$ to be the linking number
$$\text{lk}(f|(N^n),f|(N^n)+ϵ\nu ).$$
###### Remark 2.27.
Notice that $`(f,\nu )`$ corresponds to the $`N^n`$-immersion $`(f,df)|N^n`$.
As a consequence of Lemma 6 in we obtain the following lemma.
###### Lemma 2.28.
Suppose that $`fL(M^n,^{2n1})`$ and $`pS(f)`$. Choose a coordinate neighborhood $`U_p`$ of $`p`$ diffeomorphic to $`D^n`$ such that $`f|U`$ is generic and has the form (1.1). Then $`_f(U_p)=\pm 1`$, moreover, the sign does not depend on the choice of $`U_p`$ if $`n`$ is odd. (Recall that for $`n`$ even $`_f`$ is defined mod $`2`$.)
###### Notation 2.29.
Denote the map $`(f,df)|U_p^{}(n,2n1;n)`$ by $`(w,w^{})`$.
###### Corollary 2.30.
For $`n`$ even $`_f(U_p)1`$ (mod $`2`$) for every neighborhood $`U_p`$.
In signs are defined for Whitney-umbrellas if $`n`$ is odd and $`^{2n1}`$ is oriented. This can be done as follows (this is not the original definition of Whitney):
###### Definition 2.31.
Suppose that $`fL(M^n,^{2n1})`$ and $`pS(f)`$. Choose $`U_p`$ as in Lemma 2.28. The sign of the Whitney-umbrella at $`p`$ is then defined as
$$\text{sgn}(p)=_f(U_p).$$
###### Notation 2.32.
If $`n`$ is odd and $`fL(M^n,^{2n1})`$ denote the set of positive (respectively negative) cross-cap points of $`f`$ by $`S_+(f)`$ (respectively $`S_{}(f)`$). Moreover let
$$\mathrm{\#}S(f)=\{\begin{array}{cc}\left|S_+(f)\right|\left|S_{}(f)\right|\hfill & \text{if }n\text{ is odd,}\hfill \\ \left|S(f)\right|\text{mod}\mathrm{\hspace{0.17em}2}\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$
where $`\left|A\right|`$ denotes the cardinality of the set $`A`$.
In Theorem 2 states the following:
###### Lemma 2.33.
Let $`N^n`$ be a compact $`n`$-manifold with boundary and
$$fG(N^n,^{2n1}).$$
Then $`_f(N^n)=\mathrm{\#}S(f)`$.
###### Remark 2.34.
If $`f`$ is only locally generic, we may perturb it slightly to obtain a generic map. From this we can generalize Lemma 2.33 to locally generic maps.
As a trivial consequence we obtain
###### Corollary 2.35.
Let $`M^n`$ be a closed $`n`$-manifold and $`fL(M^n,^{2n1})`$. Then for $`n`$ odd the algebraic number of singular points vanishes (i.e., $`\left|S_+(f)\right|=\left|S_{}(f)\right|`$), while for $`n`$ even, it vanishes mod $`2`$ (i.e., $`\left|S(f)\right|`$ is even).
We will also use Theorem $`7`$ of , which states the following:
###### Theorem 2.36.
Let $`M^n`$ be a compact and connected manifold with boundary and $`f:M^n^{2n1}`$ be a smooth map which is an embedding in a neighborhood of $`M^n`$. Suppose that $`_f(M^n)=0`$ if $`n`$ is odd and $`_f(M^n)0`$ mod $`2`$ if $`n`$ is even. Then there exists an immersion $`g:M^n^{2n1}`$ that is arbitrarily $`C^0`$-close to $`f`$ and equals $`f`$ in a neighborhood of $`M^n`$.
Finally, we present an equivalent definition of the sign of a Whitney-umbrella singularity (motivated by ). With an additional choice we can also define the sign of a singular point if $`n`$ is even.
###### Definition 2.37.
Let $`fL(M^n,^{2n1})`$ and $`pS(f)`$. Fix an orientation of $`^{2n1}`$. Then we define the sign of $`p`$ as follows.
Choose local coordinates $`(x_1,\mathrm{},x_n)`$ about $`p`$ and $`(y_1,\mathrm{},y_{2n1})`$ about $`f(p)`$ such that $`f`$ is given near $`p`$ by equation (1.1). Let $`D_\epsilon =\{y_1^2+\mathrm{}+y_{2n1}^2\epsilon \}`$ for $`\epsilon >0`$ sufficiently small. Then Lemma 2.2 in states that $`B_\epsilon =f^1(D_\epsilon )`$ is a closed disc neighborhood of $`p`$ in $`M^n`$. Choose an arbitrary orientation of $`S_\epsilon =B_\epsilon `$. Denote by $`q`$ the double point of $`f|S_\epsilon `$ and let $`f^1(q)=\{q_1,q_2\}`$. Moreover, let $`\underset{¯}{v}=(v_1,\mathrm{},v_{n1})`$ be a positive basis of $`T_{q_1}S_\epsilon `$ and $`\underset{¯}{w}=(w_1,\mathrm{},w_{n1})`$ a positive basis of $`T_{q_2}S_\epsilon `$. Orient $`D_\epsilon `$ by its outward normal vector. Then $`p`$ is called positive if $`(df(\underset{¯}{v}),df(\underset{¯}{w}))`$ is a positive basis of $`T_q(D_\epsilon )`$, and negative otherwise.
This definition does not depend on the choice of the orientation of $`S_\epsilon `$. However, if $`n`$ is even then the sign of $`p`$ does depend on the ordering of the points $`q_1`$ and $`q_2`$, i.e., if we swap the two points then the sign of $`p`$ also changes. For $`n`$ odd the sign of $`p`$ is independent of the ordering of $`q_1`$ and $`q_2`$. Thus fixing the orientation of $`^{2n1}`$ defines the sign of a Whitney-umbrella point only if $`n`$ is odd. If we also fix an ordering of the two branches of $`f`$ meeting at the open double curve ending at $`f(p)`$ then this defines a sign of $`p`$ if $`n`$ is even.
It is easy to verify that Definition 2.31 and Definition 2.37 are equivalent.
## 3. Proof of Theorem 1.7
### 3.1. Outline
The proof of the case $`n=2`$ can be found in . The generalizations of those ideas are incorporated into the present proof for the case $`n>3`$.
The implication
$$fg\left|S(f)\right|=\left|S(g)\right|$$
is Proposition 2.4 in . The proof relies on the fact that the cross-cap singularity is stable and $`M^n`$ is closed. Thus we will only prove the other implication.
For $`n>3`$ the proof is divided into two main parts. Using Hirsch-Smale theory (for a reference see ) we will see that there are two obstructions to constructing a regular homotopy between two immersions from $`M^n`$ to $`^{2n1}`$. These obstructions can be eliminated in the presence of cross-cap points.
In the first part we will construct a regular homotopy that pushes the $`(n1)`$-simplices of $`M^n`$ through a singular point to destroy the first obstruction. In the second part we will first reduce the problem to the case $`M^n=S^n`$. To remove the second obstruction we will merge the closed double curves of $`f`$ and $`g`$ with the ones connecting the singular points. This can be done using a variant of the Whitney-trick. Finally, we replace $`f`$ and $`g`$ with immersions and we use an argument of that shows that the Smale-invariant of an immersion of $`S^n^{2n1}`$ is completely determined by the geometry of its self-intersection if $`n4`$.
### 3.2. Setup
From now on $`f`$ and $`g`$ denote the maps in the statement of Theorem 1.7. Thus $`f,gL^{}(M^n,^{2n1})`$ and $`\left|S(f)\right|=\left|S(g)\right|`$. We also suppose that $`n>3`$ and $`^{2n1}`$ is oriented.
Using the lemma of homogeneity and Corollary 2.35 there exists a diffeotopy $`\{d_t:t[0,1]\}`$ of $`M^n`$ such that $`d_0=id_{M^n}`$, moreover
* $`d_1(S(g))=S(f)`$ if $`n`$ is even,
* $`d_1(S_+(g))=S_+(f)`$ and $`d_1(S_{}(g))=S_{}(f)`$ if $`n`$ is odd.
This implies that $`S(fd_1)=S(g)`$, moreover $`S_+(fd_1)=S_+(g)`$ and $`S_{}(fd_1)=S_{}(g)`$ if $`n`$ is odd. Since $`\{fd_t:t[0,1]\}`$ provides a regular homotopy connecting $`f`$ and $`fd_1`$ we might suppose that $`S(f)=S(g)`$, moreover $`S_+(f)=S_+(g)`$ and $`S_{}(f)=S_{}(g)`$ if $`n`$ is odd.
If $`n`$ is even then any two cross-caps are locally equivalent. For $`n`$ odd any two cross-caps of the same sign are locally equivalent. Thus there is a neighborhood $`U`$ of $`S(f)=S(g)`$ such that $`f`$ is regularly homotopic to $`g`$ on $`U`$. This regular homotopy can be extended to $`M^n`$ using Smale’s lemma. Thus we may suppose that $`f|U=g|U`$.
If $`U`$ is chosen small enough then $`f|U=g|U`$ is generic. We perturb $`f`$ and $`g`$ slightly outside $`U`$, using a regular homotopy, to obtain generic maps. Thus we can also suppose that $`f,gG(M^n,^{2n1})`$.
Choose a smooth simplicial decomposition of $`M^n`$ so fine that for any $`n`$-simplex $`\mathrm{\Delta }^n`$ containing a point $`p`$ of $`S(f)=S(g)`$ we have $`\mathrm{\Delta }^nU`$ and $`p\text{int}(\mathrm{\Delta }^n)`$. (The proof of the existence of a smooth simplicial decomposition can be found in .) Thus $`f|\mathrm{\Delta }^n=g|\mathrm{\Delta }^n`$ for each $`\mathrm{\Delta }^n`$ as above. For $`pS(f)`$ we also choose $`\mathrm{\Delta }^np`$ so small that $`f|\mathrm{\Delta }^n=g|\mathrm{\Delta }^n`$ has the canonical form (1.1) in a coordinate neighborhood containing $`\mathrm{\Delta }^n`$ and centered at $`p`$.
The fiber of the bundle $`\text{MONO}(TM,T^{2n1})`$ is homeomorphic to the Stiefel manifold $`V_{2n1,n}`$ of $`n`$-frames in $`^{2n1}`$. It is well known that $`\pi _i(V_{2n1,n})=0`$ for $`i<n1`$.
Thus there exists an $`M^n`$-regular homotopy connecting the $`M^n`$\- immersion
$$(f,df)|\text{sk}_{n2}(M^n)\text{with}(g,dg)|\text{sk}_{n2}(M^n).$$
Using Smale’s lemma we can extend this $`M^n`$-regular homotopy to the whole manifold $`M^n`$ keeping $`f|U`$ fixed. So we might suppose that
$$(f,df)|\text{sk}_{n2}(M^n)=(g,dg)|\text{sk}_{n2}(M^n).$$
### 3.3. The first obstruction
Our next task is to find an $`M^n`$-regular homotopy connecting $`(f,df)`$ with $`(g,dg)`$ on $`sk_{n1}(M^n)`$. The obstruction $`\mathrm{\Omega }(df,dg)`$ to finding a regular homotopy connecting $`(f,df)`$ with $`(g,dg)`$ on an $`(n1)`$-simplex $`\mathrm{\Delta }^{n1}`$, fixing the boundary $`\mathrm{\Delta }^{n1}`$, lies in
$$\pi _{n1}(V_{2n1,n})\{\begin{array}{cc}_2\hfill & \text{if }n\text{ is even,}\hfill \\ \hfill & \text{if }n\text{ is odd.}\hfill \end{array}$$
We now prove that the sum of the obstructions on the boundary of an $`n`$-simplex on which both $`f`$ and $`g`$ are immersions vanishes.
###### Lemma 3.1.
Suppose that the $`C^{\mathrm{}}`$ $`^n`$-immersions $`(f,f^{})`$ and $`(g,g^{})`$ of $`\mathrm{\Delta }^n`$ into $`^q`$ $`(q>n)`$ are tangent on $`\text{sk}_{n2}(\mathrm{\Delta }^n)`$. Denote the faces of $`\mathrm{\Delta }^n`$ by
$$\mathrm{\Delta }_0^{n1},\mathrm{},\mathrm{\Delta }_n^{n1},$$
oriented by the outward normal vectors of $`\mathrm{\Delta }^n`$ and let
$$(f_i,f_i^{})=(f,f^{})|\mathrm{\Delta }_i^{n1},(g_i,g_i^{})=(g,g^{})|\mathrm{\Delta }_i^{n1}(i=0,\mathrm{},n).$$
Then
$$\underset{i=0}{\overset{n}{}}\mathrm{\Omega }(f_i^{},g_i^{})=\tau (f^{})\tau (g^{}).$$
###### Proof.
We define a sequence of $`C^{\mathrm{}}`$ $`^n`$-immersions
$$(h_i,h_i^{}):\mathrm{\Delta }^n^{2n1}(i=0,\mathrm{},n+1)$$
such that $`(h_0,h_0^{})=(f,f^{})`$, $`(h_{n+1},h_{n+1}^{})=(g,g^{})`$ and for $`0<i<n+1`$ let
$$(h_i,h_i^{})|\mathrm{\Delta }_j^{n1}=\{\begin{array}{cc}(g_j,g_j^{})\hfill & \text{if }j<i\text{,}\hfill \\ (f_j,f_j^{})\hfill & \text{if }ijn\text{.}\hfill \end{array}$$
Then $`(h_i,h_i^{})`$ is $`C^{\mathrm{}}`$ since $`(f,f^{})`$ and $`(g,g^{})`$ are tangent on $`sk_{n2}(\mathrm{\Delta }^n)`$. Applying Theorem 2.20 to $`(h_i,h_i^{})`$ and $`(h_{i+1},h_{i+1}^{})`$ we get that
$$\mathrm{\Omega }(f_i^{},g_i^{})=\tau (h_i^{})\tau (h_{i+1}^{}).$$
Summing up these equations we obtain the required result. ∎
###### Corollary 3.2.
Let $`q>n`$ and $`f,g\text{Imm}(\mathrm{\Delta }^n,^q)`$. Suppose that $`(f,f^{})=(f,df)|\mathrm{\Delta }^n`$ and $`(g,g^{})=(g,dg)|\mathrm{\Delta }^n`$ are tangent on $`\text{sk}_{n2}(\mathrm{\Delta }^n)`$. Using the notations of Lemma 3.1 it holds that
$$\underset{i=0}{\overset{n}{}}\mathrm{\Omega }(f_i^{},g_i^{})=0.$$
###### Proof.
Since both $`(f,f^{})`$ and $`(g,g^{})`$ are extendible we have that $`\tau (f^{})=0`$ and $`\tau (g^{})=0`$. Now Lemma 3.1 gives the statement of our corollary. ∎
Lemma 2.11 implies that there is a homeomorphism
$$\chi :^{}(n,2n1;n)\text{Imm}_1(S^{n1},^{2n1}).$$
###### Notation 3.3.
Let $`\text{Emb}_1(S^k,^l)`$ denote the space of pairs $`(h,\nu )`$, where $`h:S^k^l`$ is an embedding and $`\nu `$ is a transversal $`1`$-field along $`h`$.
###### Definition 3.4.
Two elements of $`\text{Imm}_1(S^{n1},^{2n1})`$ (or $`^{}(n,2n1;n)`$) are called *regularly homotopic* if they lie in the same path-component of the corresponding space. A *regular homotopy* is a path in $`\text{Imm}_1(S^{n1},^{2n1})`$ (or $`^{}(n,2n1;n)`$).
###### Proposition 3.5.
For every immersion $`(f,\nu )\text{Imm}_1(S^{n1},^{2n1})`$ there exists an embedding $`(h,\mu )\text{Emb}_1(S^{n1},^{2n1})`$ such that $`(f,\nu )`$ is regularly homotopic to $`(h,\mu )`$.
###### Proof.
The subspace of embeddings is dense open in $`\text{Imm}(S^{n1},^{2n1})`$. Thus we can perturb $`f`$ by a regular homotopy to get an embedding $`h`$. If the perturbation is sufficiently small then $`\nu `$ can be deformed simultaneously as a transversal field along the regular homotopy. ∎
###### Definition 3.6.
Suppose that $`^{2n1}`$ is oriented. Let
$$\text{lk}:\text{Emb}_1(S^{n1},^{2n1})\{\begin{array}{cc}\hfill & \text{if }n\text{ is odd,}\hfill \\ _2\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$
be defined as follows: For $`(f,\nu )\text{Emb}_1(S^{n1},^{2n1})`$ choose $`\epsilon >0`$ sufficiently small so that $`f+\epsilon \nu `$ is an embedding. Then let $`\text{lk}(f,\nu )`$ be accordingly the $``$ or $`_2`$ linking number $`\text{lk}(f,f+\epsilon \nu )`$.
###### Proposition 3.7.
lk is a regular homotopy invariant, i.e., if the maps
$`(f,\nu ),(h,\mu )\text{Emb}_1(S^{n1},^{2n1})`$ are regularly homotopic then $`\text{lk}(f,\nu )=\text{lk}(h,\mu )`$.
###### Proof.
This is exactly Lemma 9 in . ∎
Now we can extend the definition of lk from embeddings to immersions.
###### Definition 3.8.
Suppose that $`^{2n1}`$ is oriented. Let
$$\text{lk}:\text{Imm}_1(S^{n1},^{2n1})\{\begin{array}{cc}\hfill & \text{if }n\text{ is odd,}\hfill \\ _2\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$
be defined as follows: For $`(f,\nu )\text{Imm}_1(S^{n1},^{2n1})`$ choose an embedding $`(h,\mu )`$ in $`\text{Emb}_1(S^{n1},^{2n1})`$ regularly homotopic to $`(f,\nu )`$ (such an embedding exists because of Proposition 3.5). Then let $`\text{lk}(f,\nu )=\text{lk}(h,\mu )`$. This definition does not depend on the choice of $`(h,\mu )`$ because of Proposition 3.7.
###### Proposition 3.9.
The function lk defined above is a regular homotopy invariant.
###### Proof.
This is clear from Proposition 3.7 and Definition 3.8. ∎
###### Notation 3.10.
Let $`I_n^{}=\pi _0(^{}(n,2n1;n))`$ and $`J_n^{}=\pi _0(\text{Imm}_1(S^{n1},^{2n1}))`$. Proposition 3.9 implies that lk defines a function on $`J_n^{}`$, which we denote by $`\text{lk}_{}`$.
###### Lemma 3.11.
Let $`fL(D^n,^{2n1})`$ and
$$(f,f^{})=(f,df)|S^{n1}^{}(n,2n1;n).$$
Then
$$\text{lk}(\chi (f,f^{}))=\mathrm{\#}S(f).$$
(For the definition of $`\chi `$ see Lemma 2.11.)
###### Proof.
Perturb $`f`$ using a regular homotopy to get a generic map $`h`$ (thus $`h|S^{n1}`$ is an embedding). Let $`(h,h^{})=(h,dh)|S^{n1}`$. Proposition 3.9 implies that $`\text{lk}(\chi (f,f^{}))=\text{lk}(\chi (h,h^{}))`$. Moreover, the number of singular points is also a regular homotopy invariant, thus $`\mathrm{\#}S(f)=\mathrm{\#}S(h)`$. By definition $`\text{lk}(\chi (h,h^{}))=_h(D^n)`$. On the other hand Lemma 2.33 implies that $`_h(D^n)=\mathrm{\#}S(h)`$. ∎
###### Corollary 3.12.
Let $`(f,\nu )\text{Imm}_1(S^{n1},^{2n1})`$. Then
$$\text{lk}(f,\nu )=\text{lk}(f,\nu ).$$
###### Proof.
Choose two generic maps $`f_0,f_1G(D^n,^{2n1})`$ such that
$$\chi ((f_i,df_i)|S^{n1})=(f,(1)^i\nu )$$
for $`i=0,1`$. Then $`f_0`$ and $`f_1`$ fit together to define a map $`hG(S^n,^{2n1})`$. Lemma 2.33 implies that $`\mathrm{\#}S(f_0)+\mathrm{\#}S(f_1)=\mathrm{\#}S(h)=0`$. Moreover, Lemma 3.11 implies that $`\text{lk}(f,\nu )=\mathrm{\#}S(f_0)`$ and $`\text{lk}(f,\nu )=\mathrm{\#}S(f_1)`$. Putting together these results we get that $`\text{lk}(f,\nu )+\text{lk}(f,\nu )=0`$, as required. ∎
###### Proposition 3.13.
$`\tau `$ (see Definition 2.13) is a regular homotopy invariant. Moreover, it induces a bijection $`\tau _{}:I_n^{}\pi _{n1}(V_{2n1,n})`$.
###### Proof.
A regular homotopy of a map $`(f,f^{})^{}(n,2n1;n)`$ induces a homotopy of the map $`S^{n1}V_{q,n}`$ defined by
$$xf^{}\{e_1(x),\mathrm{},e_n(x)\}.$$
Thus the homotopy class $`\tau (f^{})`$ is constant on the regular homotopy class of $`(f,f^{})`$.
Hirsch’s theorem (Theorem 2.21) implies that
$$\tau _{}:\pi _0(\text{Imm}_1(S^{n1},^{2n1}))\pi _0(\text{MONO}(TS^{n1}\epsilon ^1,^{2n1}))\pi _{n1}(V_{2n1,n})$$
is a bijection (the bundle $`TS^{n1}\epsilon ^1`$ is trivial). ∎
###### Definition 3.14.
Define the *connected sum* operation $`\mathrm{\#}`$ on elements $`x`$ and $`y`$ of $`\text{Imm}_1(S^{n1},^{2n1})`$ or $`^{}(n,2n1;n)`$ by joining them with a thin tube. This operation is also well defined on regular homotopy classes of maps: Suppose that $`x_Hx_1`$ and $`y_Ky_1`$. Smale’s lemma implies that the regular homotopies $`H`$ and $`K`$ may be kept fixed on small disks. Now a regular homotopy $`L`$ connecting $`x\mathrm{\#}x_1`$ with $`y\mathrm{\#}y_1`$ is defined by taking $`H`$ and $`K`$, and joining them with a tube attached to the disks kept fixed.
###### Proposition 3.15.
The sets $`I_n^{}`$ and $`J_n^{}`$ (introduced in Notation 3.10) endowed with the connected sum operation form abelian groups. Moreover, $`\chi _{}:I_n^{}J_n^{}`$ and $`\tau _{}`$ are group isomorphisms and $`\text{lk}_{}`$ is a group homomorphism.
###### Proof.
It is clear that $`(I_n^{},\mathrm{\#})`$ and $`(J_n^{},\mathrm{\#})`$ are abelian semigroups. Let $`i:D^n^{2n1}`$ denote the standard embedding. Then $`(i,i^{})=(i,di)|S^{n1}`$ (respectively
$`\chi ((i,di)|S^{n1})=(i,r)`$, where $`r`$ is the outward normal field of $`S^{n1}`$ in $`^n`$) represents the identity of $`I_n^{}`$ (respectively $`J_n^{}`$).
If $`x,yI_n^{}`$ then choose a representative $`(f,f^{})`$ for $`x`$ and $`(g,g^{})`$ for $`y`$. Then $`\tau (f^{}\mathrm{\#}g^{})`$ is the connected sum of the spheroids $`\tau (f^{})`$ and $`\tau (g^{})`$, which implies that $`\tau _{}(x\mathrm{\#}y)=\tau _{}(x)+\tau _{}(y)`$. Since $`(i,i^{})`$ is extendible, $`\tau _{}(i^{})=0`$. Proposition 3.13 states that $`\tau _{}`$ is a bijection. Thus $`\tau _{}`$ is a semigroup isomorphism that takes the identity of $`I_n^{}`$ to the identity of $`\pi _{n1}(V_{2n1,n})`$ or $`_2`$. This implies that $`I_n^{}`$ is a group isomorphic to $``$ or $`_2`$ (depending on the parity of $`n`$) and $`\tau _{}`$ is a group isomorphism.
Since $`\chi `$ is a homeomorphism, $`\chi _{}:I_n^{}J_n^{}`$ is a bijection. Moreover, from the definition of $`\chi `$ it is easy to see that $`\chi _{}`$ is a semigroup homomorphism. We saw above that $`\chi _{}`$ takes the identity $`(i,i^{})`$ of $`I_n^{}`$ to the identity $`(i,r)`$ of $`J_n^{}`$. Thus $`J_n^{}`$ is also a group isomorphic to $``$ or $`_2`$ and $`\chi _{}`$ is a group isomorphism.
We finally show that $`\text{lk}_{}`$ is a group homomorphism from $`J_n^{}`$ to $``$ or $`_2`$. Suppose that $`x,yJ_n^{}`$ and choose representatives $`(f_1,\nu _1)`$ and $`(f_2,\nu _2)`$ for $`x`$, respectively $`y`$, such that $`f_1`$ and $`f_2`$ can be separated by a hyperplane in $`^{2n1}`$. Also choose $`h_1,h_2G(D^n,^{2n1})`$ so that $`\chi ((h_i,dh_i)|S^{n1})=(f_i,\nu _i)`$ for $`i=1,2`$. Denote the boundary sum $`h_1\mathrm{}h_2`$ with $`h`$. Then $`\chi ((h,dh)|S^{n1})=(f_1,\nu _1)\mathrm{\#}(f_2,\nu _2)`$ and $`\mathrm{\#}S(h)=\mathrm{\#}S(h_1)+\mathrm{\#}S(h_2)`$. Using Lemma 3.11 we get that $`\text{lk}((f_1,\nu _1)\mathrm{\#}(f_2,\nu _2))=\text{lk}(f_1,\nu _1)+\text{lk}(f_2,\nu _2)`$, i.e., $`\text{lk}(x\mathrm{\#}y)=\text{lk}(x)+\text{lk}(y)`$. ∎
The following lemma is the key to connecting the works of Hirsch and Whitney.
###### Lemma 3.16.
There exists an isomorphism
$$\alpha :\pi _{n1}(V_{2n1,n})\{\begin{array}{cc}\hfill & \text{if }n\text{ is odd,}\hfill \\ _2\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$
such that
$$\text{lk}\chi =\alpha \tau .$$
I.e., the following diagram is commutative:
$$\begin{array}{ccc}^{}(n,2n1;n)& \stackrel{\chi }{}& \text{Imm}_1(S^{n1},^{2n1})\\ \tau & & \text{lk}& & \\ \pi _{n1}(V_{2n1,n})& \stackrel{\alpha }{}& \text{or}_2.\end{array}$$
###### Remark 3.17.
If $`n`$ is even then $`\alpha `$ might be omitted from the formula since $`_2`$ has only one automorphism.
###### Proof.
We saw that $`\tau `$ and lk are regular homotopy invariants, i.e., they are constant on the path-components of their domains. They define maps
$$\tau _{}:I_n^{}=\pi _0(^{}(n,2n1;n))\pi _{n1}(V_{2n1,n})$$
and
$$\text{lk}_{}:J_n^{}=\pi _0(\text{Imm}_1(S^{n1},^{2n1}))\{\begin{array}{cc}\hfill & \text{if }n\text{ is odd,}\hfill \\ _2\hfill & \text{if }n\text{ is even.}\hfill \end{array}$$
Proposition 3.15 implies that $`I_n^{}`$ and $`J_n^{}`$ are abelian groups with the connected sum operation, moreover, $`\chi _{},\tau _{}`$ are isomorphisms and $`\text{lk}_{}`$ is a homomorphism. We also know that $`I_n^{}`$ and $`J_n^{}`$ are isomorphic to $``$ or $`_2`$, according to the parity of $`n`$.
$`\text{lk}_{}`$ is an epimorphism because of Lemma 2.28:
$$\text{lk}(\chi (w,w^{}))=\pm 1.$$
So $`\text{lk}_{}`$ is a $``$ or $`_2_2`$ epimorphism, and thus it is also an isomorphism. Now define $`\alpha `$ to be $`\text{lk}_{}\chi _{}(\tau _{})^1`$. ∎
From Theorem 2.20 and Lemma 3.16 we get the following
###### Proposition 3.18.
Suppose that $`(f,\nu ),(g,\mu )\text{Imm}_1(\mathrm{\Delta }^n,^q)`$ agree on $`(\mathrm{\Delta }^n)\text{int}(\mathrm{\Delta }^{n1})`$, where $`\mathrm{\Delta }^{n1}`$ is an $`(n1)`$-face of the $`n`$-simplex $`\mathrm{\Delta }^n`$. Then for $`(f,f^{})=(\chi ^1(f,\nu ))|\mathrm{\Delta }^{n1}`$ and $`(g,g^{})=(\chi ^1(g,\mu ))|\mathrm{\Delta }^{n1}`$
$$\alpha (\mathrm{\Omega }(f^{},g^{}))=\text{lk}(f,\nu )\text{lk}(g,\mu ).$$
###### Remark 3.19.
$`\mathrm{\Delta }^{n1}`$ is co-oriented by the outward normal vector of $`\mathrm{\Delta }^n`$. Then
$$(\chi ^1(f,\nu ))|\mathrm{\Delta }^{n1}=\zeta ^1((f,\nu )|\mathrm{\Delta }^{n1}),$$
where $`\zeta `$ is the homeomorphism of Lemma 2.9. Thus $`\text{lk}(f,\nu )\text{lk}(g,\mu )`$ depends only on $`(f,\nu )|\mathrm{\Delta }^{n1}`$ and $`(g,\mu )|\mathrm{\Delta }^{n1}`$.
###### Definition 3.20.
For $`(f,f^{}),(g,g^{})(n1,2n1;n)`$ that are tangent on $`\mathrm{\Delta }^n`$ let
$$o(f^{},g^{})=\alpha (\mathrm{\Omega }(f^{},g^{})).$$
###### Corollary 3.21.
Let $`\mathrm{\Delta }^{n1}`$ be a face of the $`n`$-simplex $`\mathrm{\Delta }^n`$ and let $`(f,f^{})`$ and $`(g,g^{})`$ be $`C^{\mathrm{}}`$ $`^n`$-immersions of $`\mathrm{\Delta }^{n1}`$ into $`^{2n1}`$ tangent on $`\mathrm{\Delta }^{n1}`$. Choose extensions of $`(f,f^{})`$ and $`(g,g^{})`$ to $`\mathrm{\Delta }^n`$ that are tangent on $`\mathrm{\Delta }^n\mathrm{\Delta }^{n1}`$. Then
$$o(f^{},g^{})=\text{lk}\chi (f,f^{})\text{lk}\chi (g,g^{}).$$
###### Lemma 3.22.
Let $`(f,f^{})`$ and $`(g,g^{})`$ be as in Corollary 3.21 and $`(w,w^{})`$ be the boundary of a Whitney-umbrella of sign $`\epsilon `$ (as in Notation 2.29). Then
$$o(f^{}\mathrm{\#}w^{},g^{})=o(f^{},g^{})+\epsilon .$$
###### Proof.
Using extensions as in Corollary 3.21 we get that
$$o(f^{}\mathrm{\#}w^{},g^{})=\text{lk}\chi ((f,f^{})\mathrm{\#}(w,w^{}))\text{lk}\chi (g,g^{}).$$
Using the fact that lk and $`\chi `$ are additive we get that this equals
$$\text{lk}\chi (f,f^{})+\text{lk}\chi (w,w^{})\text{lk}\chi (g,g^{}).$$
The definition of the sign of an umbrella implies that $`\text{lk}\chi (w,w^{})=\epsilon `$. ∎
###### Proposition 3.23.
Let $`(f,f^{})`$ and $`(g,g^{})`$ be as in Corollary 3.21. Then $`(f,f^{})`$ and $`(g,g^{})`$ are regularly homotopic (rel $`\mathrm{\Delta }^{n1}`$) iff $`o(f^{},g^{})=0`$.
###### Proof.
This is a trivial consequence of Lemma 2.19. ∎
###### Remark 3.24.
Lemma 3.22 and Proposition 3.23 show us how Whitney-umbrellas destroy obstructions to moving $`f`$ to $`g`$ on an $`(n1)`$-simplex (rel boundary). We have to take the connected sum of $`f`$ with several copies of the boundary of an umbrella. This will be done by a diffeotopy of $`M^n`$ that pushes the $`(n1)`$-simplex through a singular point.
###### Lemma 3.25.
Using the notations of Lemma 3.1
$$\underset{i=0}{\overset{n}{}}o(f_i^{},g_i^{})=\text{lk}\chi (f,f^{})\text{lk}\chi (g,g^{}).$$
In particular, if $`(f,f^{})`$ and $`(g,g^{})`$ are both extendible then
$$\underset{i=0}{\overset{n}{}}o(f_i^{},g_i^{})=0.$$
###### Proof.
This is a trivial consequence of Lemma 3.1. ∎
###### Remark 3.26.
Lemma 3.25 shows us that the obstruction $`o`$ has the cocycle property, i.e., the sum of the obstructions on the boundary of an $`n`$-simplex, on which both $`f`$ and $`g`$ are immersions, vanishes.
### 3.4. Pushing (n-1)-simplices through Whitney-umbrellas
Now we are going to use the apparatus developed above for the proof of Theorem 1.7. Since until now we only worked in a standard simplex in Euclidean space, we have to globalize our results to the triangulated manifold $`M^n`$. We do *not* require $`M^n`$ to be oriented.
###### Definition 3.27.
Suppose that $`\mathrm{\Delta }^n`$ is an $`n`$-simplex of $`M^n`$. Then we can define the map $`\chi =\chi _{\mathrm{\Delta }^n}`$ taking $`M^n`$-immersions of $`\mathrm{\Delta }^n`$ to $`\text{Imm}_1(\mathrm{\Delta }^n,^{2n1})`$ using the outward normal vectors $`r(x)`$ of $`\mathrm{\Delta }^n`$ (corners rounded off) as follows: Let $`(f,f^{})`$ be an $`M^n`$-immersion of $`\mathrm{\Delta }^n`$ into $`^{2n1}`$. Then let $`\chi (f,f^{})=(f,\nu )`$, where $`\nu (x)=f^{}(r(x))`$.
###### Definition 3.28.
Suppose that $`\mathrm{\Delta }^{n1}`$ is an $`(n1)`$-simplex of $`M^n`$ that is co-oriented. Then we can define the map $`\zeta =\zeta _{\mathrm{\Delta }^{n1}}`$ taking $`M^n`$-immersions of $`\mathrm{\Delta }^{n1}`$ to $`\text{Imm}_1(\mathrm{\Delta }^{n1},^{2n1})`$ as follows: Let $`r(x)`$ be a positive normal field along $`\mathrm{\Delta }^{n1}`$. If $`(f,f^{})`$ is an $`M^n`$-immersion of $`\mathrm{\Delta }^{n1}`$ into $`^{2n1}`$ then let $`\zeta (f,f^{})=(f,\nu )`$, where $`\nu (x)=f^{}(r(x))`$.
Corollary 3.21 implies that the following definition makes sense:
###### Definition 3.29.
Suppose that $`\mathrm{\Delta }^{n1}`$ is an $`(n1)`$-simplex of $`M^n`$ that is co-oriented and let $`(f,f^{})`$ and $`(g,g^{})`$ be $`M^n`$-immersions of $`\mathrm{\Delta }^{n1}`$ into $`^{2n1}`$ tangent on $`\mathrm{\Delta }^{n1}`$. Then the obstruction $`o(f^{},g^{})`$ to finding a regular homotopy of $`(f,f^{})`$ to $`(g,g^{})`$ (rel $`\mathrm{\Delta }^{n1}`$) is given as follows: Choose an $`n`$-simplex $`\mathrm{\Delta }^n`$ of $`M^n`$ such that $`\mathrm{\Delta }^{n1}\mathrm{\Delta }^n`$. Extend $`(f,f^{})`$ and $`(g,g^{})`$ to $`\mathrm{\Delta }^n`$ to be tangent on $`\mathrm{\Delta }^n\mathrm{\Delta }^{n1}`$. Also choose a normal field $`s(x)`$ along $`\mathrm{\Delta }^n`$ that agrees with the co-orientation of $`\mathrm{\Delta }^{n1}`$. Now define $`o(f^{},g^{})`$ to be $`\text{lk}(f,f^{}s)\text{lk}(g,g^{}s)`$.
###### Remark 3.30.
If we change the co-orientation of $`\mathrm{\Delta }^{n1}`$ then $`o(f^{},g^{})`$ changes sign: To see this we choose the same extensions of $`(f,f^{})`$ and $`(g,g^{})`$ to the boundary $`\mathrm{\Delta }^n`$ of the same $`n`$-simplex $`\mathrm{\Delta }^n`$. Then we must take $`s`$ instead of $`s`$ because of the changed co-orientation of $`\mathrm{\Delta }^n`$. Now using Corollary 3.12 we get that
$$\text{lk}(f,f^{}(s))\text{lk}(g,g^{}(s))=\text{lk}(f,f^{}s))+\text{lk}(g,g^{}s)=o(f^{},g^{}).$$
Let $`f`$ and $`g`$ be the two locally generic maps of Theorem 1.7. We have supposed that $`(f,df)`$ and $`(g,dg)`$ are tangent on $`\text{sk}_{n2}M^n`$. If $`n`$ is odd choose a co-orientation $`O_{\mathrm{\Delta }^{n1}}`$ for every $`(n1)`$-simplex $`\mathrm{\Delta }^{n1}`$ of $`M^n`$ so that $`o(df,dg)0`$. This can be done because of Remark 3.30. Let $`k=\left|S(f)\right|=\left|S(g)\right|`$ and denote by $`\mathrm{\Delta }_1^n,\mathrm{},\mathrm{\Delta }_k^n`$ the $`n`$-simplices of $`M^n`$ that contain any point of $`S(f)=S(g)`$.
We are now going to construct an oriented curve $`\gamma `$ on $`M^n`$ that intersects each simplex $`\mathrm{\Delta }^{n1}`$ of $`M^n`$ in the positive direction (according to $`O_{\mathrm{\Delta }^{n1}}`$) in exactly $`\left|o(df,dg)\right|`$ points. Notice that if $`\mathrm{\Delta }^{n1}\mathrm{\Delta }_i^n`$ for some $`1ik`$ then on $`\mathrm{\Delta }^{n1}`$ we have $`o(df,dg)=0`$ since $`f|\mathrm{\Delta }_i^n=g|\mathrm{\Delta }_i^n`$. Thus $`\gamma `$ avoids the simplices $`\mathrm{\Delta }_i^n`$. This $`\gamma `$ may be thought of as the dual of the obstruction to deform $`f`$ to $`g`$ on the $`(n1)`$-skeleton of $`M^n`$.
To obtain $`\gamma `$ choose a set of $`\left|o(df,dg)\right|`$ points $`P_{\mathrm{\Delta }^{n1}}`$ on each $`(n1)`$-simplex of $`M^n`$. Now take an $`n`$-simplex $`\mathrm{\Delta }^n`$ and denote its faces by $`\mathrm{\Delta }_0^{n1},\mathrm{},\mathrm{\Delta }_n^{n1}`$, moreover let $`(f_i,f_i^{})=(f,df)|\mathrm{\Delta }_i^{n1}`$, $`(g_i,g_i^{})=(g,dg)|\mathrm{\Delta }_i^{n1}`$, $`O_i=O_{\mathrm{\Delta }_i^{n1}}`$ and $`P_i=P_{\mathrm{\Delta }_i^{n1}}`$ for $`0in`$. Lemma 3.25 implies that if each $`\mathrm{\Delta }_i^{n1}`$ is co-oriented by the outward normal vector of $`\mathrm{\Delta }^n`$ (denote this co-orientation by $`U_i`$) then
(3.1)
$$\underset{i=0}{\overset{n}{}}o(f_i^{},g_i^{})=0.$$
If we consider $`\mathrm{\Delta }_i^{n1}`$ with the co-orientation $`U_i`$ then $`o(f_i^{},g_i^{})<0`$ implies that $`U_i=O_i`$ and $`o(f_i^{},g_i^{})>0`$ implies that $`U_i=O_i`$. Give a minus sign to each point of $`P_i`$ if $`U_i=O_i`$ and a plus sign otherwise. Then equation $`(\text{3.1})`$ is equivalent to the statement that the sum of the signs in $`_{i=0}^nP_i`$ equals $`0`$. Now let $`\gamma \mathrm{\Delta }^n`$ be given as follows: Make a bijection between the $`+`$ and $``$ points of $`_{i=0}^nP_i`$ and connect each pair of points with an embedded curve segment oriented from $`+`$ to $``$. We do this so that these curve segments are pairwise disjoint. This is possible since $`n3`$.
Doing this for each $`n`$-simplex we obtain an oriented embedded curve $`\gamma `$ with the required intersection property. Now we make $`\gamma `$ connected by taking the connected sum of its components: Let for example $`\gamma _1`$ and $`\gamma _2`$ be two components of $`\gamma `$. Then choose two points $`p_1\gamma _1`$ and $`p_2\gamma _2`$. Join $`p_1`$ and $`p_2`$ with an embedded curve $`\eta `$ that avoids $`\mathrm{\Delta }_i^n`$ for $`1ik`$ and also each $`(n2)`$-simplex and intersects each $`(n1)`$-simplex transversally. Then take two parallel curves $`\eta _1`$ and $`\eta _2`$ close to $`\eta `$ and orient them according to the orientations of $`\gamma _1`$ and $`\gamma _2`$. If $`\eta `$ intersects an $`(n1)`$-simplex $`\mathrm{\Delta }^{n1}`$ at a point $`x`$ then $`\eta _1`$ and $`\eta _2`$ will intersect $`\mathrm{\Delta }^{n1}`$ in different directions near $`x`$.
Similarly, we might modify $`\gamma `$ so that it goes through exactly one point $`p`$ of $`S(f)`$, and let the Whitney-umbrella of $`f`$ at $`p`$ be positive if $`n`$ is odd. To do this, choose a small embedded curve containing $`p`$ and join it to $`\gamma `$ as above. We will also call this modified curve $`\gamma `$. Let us suppose that the $`n`$-simplex containing $`p`$ is $`\mathrm{\Delta }_1^n`$.
Thus it will hold true that the connected embedded curve $`\gamma `$ intersects *algebraically* each $`(n1)`$-simplex $`\mathrm{\Delta }^{n1}`$ in $`\left|o(df,dg)\right|`$ points if we consider $`\mathrm{\Delta }^{n1}`$ with the co-orientation $`O_{\mathrm{\Delta }^{n1}}`$. Moreover, $`\gamma `$ contains the cross-cap point $`p\mathrm{\Delta }_1^n`$ and is disjoint from $`\mathrm{\Delta }_i^n`$ for $`1<ik`$.
Let $`\nu _\gamma `$ be a thin tubular neighborhood of $`\gamma `$ in $`M^n`$. We are going to construct a diffeotopy $`\{\mathrm{\Psi }_t:0t1\}`$ of $`\nu _\gamma `$ such that $`\mathrm{\Psi }_0=id_{\nu _\gamma }`$ and for every $`t[0,1]`$ the diffeomorphism $`\mathrm{\Psi }_t`$ is the identity in a neighborhood of $`\nu _\gamma `$. Thus we can extend $`\mathrm{\Psi }_t`$ to the whole manifold $`M^n`$ to be the identity outside $`\nu _\gamma `$.
The diffeomorphism $`\mathrm{\Psi }_t`$ is constructed as follows: Denote by $`T`$ the identity of $`^{n1}`$ if $`\gamma `$ is orientation preserving, and let $`T`$ be a reflection in a hyperplane of $`^{n1}`$ if $`\gamma `$ is orientation reversing. Then $`\nu _\gamma `$ is diffeomorphic to the factor space
$$\mathrm{\Gamma }=\times D_2^{n1}/{}_{(x,y)(x+1,T(y))}{}^{},$$
where $`D_2^{n1}^{n1}`$ is the disc of radius $`2`$. Denote by $`p:\times D_2^{n1}\mathrm{\Gamma }`$ the projection. Let $`\lambda :`$ be a $`C^{\mathrm{}}`$ function such that $`\lambda (x)=0`$ for $`x<\epsilon `$ and $`\lambda (x)=1`$ for $`x>1\epsilon `$ (where $`\epsilon <1`$ is a small positive constant). First we define a diffeotopy $`\{\varphi _t:0t1\}`$ of $`\times D_2^{n1}`$ with the formula
$$\varphi _t(x,y)=\{\begin{array}{cc}(x+t,y)\hfill & \text{if }y1\text{,}\hfill \\ (x+\lambda (s)t,y)\hfill & \text{if }y=2s\text{.}\hfill \end{array}$$
$`(x,y)(x+1,T(y))`$ implies that $`(x+c,y)(x+c+1,T(y))`$ for any number $`c`$. Thus the diffeomorphism $`\varphi _t`$ factors through the projection $`p`$ to a diffeomorphism $`\mathrm{\Psi }_t`$ of $`\mathrm{\Gamma }\nu _\gamma `$. I.e., the following diagram is commutative:
$$\begin{array}{ccc}\times D_2^{n1}& \stackrel{\varphi _t}{}& \times D_2^{n1}\\ p& & p& & \\ \nu _\gamma & \stackrel{\mathrm{\Psi }_t}{}& \nu _\gamma .\end{array}$$
Since $`\mathrm{\Psi }_1`$ is the identity on $`p(\times D_1^{n1})`$, the diffeomorphism $`\mathrm{\Psi }_1`$ is the identity on a thinner tubular neighborhood of $`\gamma `$. Thus $`\mathrm{\Psi }_1(p)=p`$. Moreover, let $`F`$ denote a fiber of $`\nu _\gamma `$ (diffeomorphic to $`D^{n1}`$) with a normal framing $`v`$ (in $`M^n`$) in the direction determined by $`\gamma `$; and $`S`$ a small sphere around $`p`$ contained in $`\nu _\gamma `$ with the outward normal framing $`r`$. Then the framed submanifold $`(\mathrm{\Psi }_1(F),d\mathrm{\Psi }_1(v))`$ is equal to the connected sum $`(F,v)\mathrm{\#}(S,r)`$. On the other hand, the framed submanifold $`(\mathrm{\Psi }_1(F),d\mathrm{\Psi }_1(v))=(F,v)\mathrm{\#}(S,r)`$ (see Figure 2).
Now look at the regular homotopy $`f\mathrm{\Psi }_t`$ connecting $`f`$ with $`h=f\mathrm{\Psi }_1`$. Then $`S(h)=S(f)`$, moreover $`S_\pm (h)=S_\pm (f)`$ if $`n`$ is odd. We also get that
$$(h,dh)|F=((f,df)|F)\mathrm{\#}((f,df)|S).$$
Notice that $`(f,df)|S`$ is the $`M^n`$-immersion $`(w,w^{})`$ of Notation 2.29 and if we co-orient $`S`$ by $`r`$ then $`\text{lk}\chi (w,w^{})=1`$ since the sign of $`p`$ is $`1`$. Thus if we co-orient $`F`$ by $`v`$, as above, then Lemma 3.22 gives that $`o(dh|F,df|F)=1`$. On the other hand, if we co-orient $`F`$ by $`v`$ then $`o(dh|F,df|F)=1`$ (see Figure 2). So if we take an $`(n1)`$-simplex $`\mathrm{\Delta }^{n1}`$ of $`M^n`$ then $`o(dh|\mathrm{\Delta }^{n1},df|\mathrm{\Delta }^{n1})`$ is the algebraic intersection of $`\mathrm{\Delta }^{n1}`$ and $`\gamma `$. From the construction of $`\gamma `$ we know that
$$\mathrm{\Delta }^{n1}\gamma =\left|o(df|\mathrm{\Delta }^{n1},dg|\mathrm{\Delta }^{n1})\right|.$$
Thus on $`\mathrm{\Delta }^{n1}`$
$$o(dh,dg)=o(dh,df)+o(df,dg)=\left|o(df,dg)\right|+o(df,dg)=0,$$
since the co-orientation $`O_{\mathrm{\Delta }^{n1}}`$ was chosen so that $`o(df,dg)0`$.
Consequently, there is no obstruction to finding an $`M^n`$-regular homotopy between $`h`$ and $`g`$ on the $`(n1)`$-skeleton of $`M^n`$. Since $`f`$ is regularly homotopic to $`h`$, we suppose from now on that $`f=h`$. Therefore we can suppose that the maps $`f`$ and $`g`$ coincide on (a neighborhood of) $`\text{sk}_{n1}(M^n)`$.
### 3.5. The second obstruction
Let us recall that $`\mathrm{\Delta }_1^n,\mathrm{},\mathrm{\Delta }_k^n`$ denote those $`n`$-simplices of $`M^n`$ which contain a singular point of $`f`$, where $`k=2l`$ is even. (Each $`\text{int}(\mathrm{\Delta }_i^n)`$ contains a single Whitney umbrella point. If $`n`$ is odd then $`\mathrm{\Delta }_1^n,\mathrm{},\mathrm{\Delta }_l^n`$ contain a positive Whitney umbrella, $`\mathrm{\Delta }_{l+1}^n,\mathrm{},\mathrm{\Delta }_{2l}^n`$ a negative one for both $`f`$ and $`g`$. Moreover, $`f|\mathrm{\Delta }_i^n=g|\mathrm{\Delta }_i^n`$ for $`1ik`$.) Let
$$U=\underset{i=1}{\overset{k}{}}\text{int}(\mathrm{\Delta }_i^n).$$
Take the manifold with boundary $`N^n=M^nU`$ and let $`\widehat{f}=f|N^n`$ and $`\widehat{g}=g|N^n`$. We know from the previous section that $`\widehat{f}`$ and $`\widehat{g}`$ are regularly homotopic on $`\text{sk}_{n1}(N^n)`$. The secondary obstruction to finding a regular homotopy between $`\widehat{f}`$ and $`\widehat{g}`$, not fixing the boundary, lies in the cohomology group $`H^n(N^n;\pi _n(V_{2n1,n}))`$ with twisted coefficients. However, this group is $`0`$ since an $`n`$-manifold without closed components is homotopy equivalent to an $`(n1)`$-dimensional simplicial complex.
So there exists a regular homotopy connecting $`\widehat{f}`$ with $`\widehat{g}`$, but not necessarily fixing $`N^n`$. Using Smale’s lemma we can extend this regular homotopy to the whole manifold $`M^n`$, fixing a small neighborhood $`V`$ ($`U`$) of $`S(f)=S(g)`$. (During this regular homotopy on the ”collar” $`\overline{U}V`$ the map $`f`$ might become twisted, i.e., it may differ from $`g`$.)
Hence we might suppose that $`f|N^n=g|N^n`$ and $`f|V=g|V`$, but it might happen that $`f|(UV)g|(UV)`$. If $`n`$ is odd then $`\left|S_+(f)\right|=\left|S_+(g)\right|=l`$, let us suppose that
$$S_+(f)=S_+(g)\underset{1il}{}\mathrm{\Delta }_i^n,$$
while
$$S_{}(f)=S_{}(g)\underset{l+1ik}{}\mathrm{\Delta }_i^n.$$
Since up to this point we have not used the assumption $`n>3`$, we can formulate and prove the following result that is weaker than Theorem 1.7 and will not be used in its proof, but will be needed in dealing with the exceptional case $`n=3`$.
###### Theorem 3.31.
Suppose that $`n>2`$ and let $`M^n`$ be a closed $`n`$-manifold. Moreover, let $`f,gL^{}(M^n,^{2n1})`$ be two locally generic maps that satisfy $`\left|S(f)\right|=\left|S(g)\right|`$. Then there exists a map $`f^{}L^{}(M^n,^{2n1})`$ and a subset $`DM^n`$ diffeomorphic to the disk $`D^n`$ such that $`ff^{}`$, the maps $`f^{}|D`$ and $`g|D`$ are immersions, and $`f^{}|(M^nD)=g|(M^nD)`$.
###### Proof.
Let $`U`$ and $`V`$ be as above (where $`S(g)VU`$). I.e., $`U`$ is the union of the interiors of the simplices $`\mathrm{\Delta }_i^n`$ for $`1ik`$ (where $`k=\left|S(f)\right|=\left|S(g)\right|`$) and $`f`$ is regularly homotopic to a map $`f^{}`$ such that $`f^{}|(M^nU)=g|(M^nU)`$ and $`f^{}|V=g|V`$. We can also suppose that $`V_i=\mathrm{\Delta }_i^nV`$ is diffeomorphic to an $`n`$-disc for $`1ik`$. Let us take a small neighborhood $`U_i`$ of $`\mathrm{\Delta }_i^n`$ that is diffeomorphic to $`D^n`$. From now on we will denote by $`U`$ the union $`_{i=1}^kU_i`$. For each $`i`$ choose a $`1`$-simplex $`e_iU_iV_i`$ connecting a point of $`U_i`$ with a point of $`V_i`$ and let $`E_i`$ be a thin tubular neighborhood of $`e_i`$. Then $`D_i=U_i(V_iE_i)`$ is diffeomorphic to $`D^n`$. Since $`n>2`$ there is no obstruction to constructing a regular homotopy between $`f^{}|E_i`$ and $`g|E_i`$ fixing the ends of the tube $`E_i`$. Thus (using Smale’s lemma) we can suppose that $`f^{}|E_i=g|E_i`$ for $`1ik`$. Now we connect $`D_i`$ and $`D_{i+1}`$ for $`1ik1`$ with a tube $`T_iM^nU`$ diffeomorphic to $`I\times D^{n1}`$. (Thus each $`T_i`$ is an $`n`$-dimensional $`1`$-handle attached to $`D_i`$ and $`D_{i+1}`$.) Let
$$D=\left(\underset{i=1}{\overset{k}{}}D_i\right)\left(\underset{j=1}{\overset{k1}{}}T_j\right).$$
Then $`D`$ is diffeomorphic to $`D^n`$ and clearly $`f^{}|D`$ and $`g|D`$ are immersions, finally $`f^{}|(M^nD)=g|(M^nD)`$. ∎
Let $`F_1,\mathrm{},F_l`$ be pairwise disjoint subsets of $`M^n`$ such that $`F_i`$ is diffeomorphic to the closed $`n`$-disc $`D^n`$, moreover $`F_i`$ contains $`\mathrm{\Delta }_i^n`$ and $`\mathrm{\Delta }_{l+i}^n`$ in its interior for $`1il`$. We will show that for $`1il`$ the locally generic maps $`f_i=f|F_i`$ and $`g_i=g|F_i`$ are regularly homotopic keeping a neighborhood of $`F_i`$ fixed. This would finish the proof of Theorem 1.7. From now on let $`1il`$ be fixed.
It is apparent from the choice of $`F_i`$ that $`S(f_i)=S(g_i)`$ and $`\left|S(f_i)\right|=\left|S(g_i)\right|=2`$. Moreover, if $`n`$ is odd we have that $`S_+(f_i)=S_+(g_i)`$ and $`S_{}(f_i)=S_{}(g_i)`$ are both $`1`$-element sets. Thus $`\mathrm{\#}S(f_i)=0`$ and $`\mathrm{\#}S(g_i)=0`$. Lemma 2.33 implies that $`_{f_i}(F_i)=0`$ and $`_{g_i}(F_i)=0`$. Using Lemma 3.16 we get that
$$0=_{f_i}(F_i)=\text{lk}\chi ((f_i,df_i)|F_i)=\alpha \tau ((f_i,df_i)|F_i).$$
Since $`\alpha `$ is an isomorphism, $`\tau ((f_i,df_i)|F_i)=0`$ and from Lemma 2.14 we can see that the map $`(f_i,df_i)|F_i`$ is extendible. Similarly, $`(g_i,dg_i)|F_i`$ is also extendible. (This also follows from Theorem 2.36.)
Let $`S_+^n`$ (respectively $`S_{}^n`$) denote the northern (respectively southern) hemisphere of $`S^n`$ and we identify $`S_+^n`$ with $`F_i`$. Then we get from the previous paragraph that there exist two locally generic maps $`f_i^{},g_i^{}L(S^n,^{2n1})`$ such that $`f_i^{}|S_+^n=f_i`$, $`g_i^{}|S_+^n=g_i`$, moreover $`f_i^{}|S_{}^n=g_i^{}|S_{}^n`$ is an immersion. Thus it is sufficient to prove the following theorem:
###### Theorem 3.32.
Let $`f,gL(S^n,^{2n1})`$ be locally generic maps and suppose that $`f|S_{}^n=g|S_{}^n`$ are immersions, $`\left|S(f)\right|=\left|S(g)\right|=2`$ and $`S(f)=S(g)`$, moreover $`S_+(f)=S_+(g)`$ if $`n`$ is odd. Then there exists a singularity fixing regular homotopy connecting $`f`$ and $`g`$ that is fixed on $`S_{}^n`$, i.e., $`f_sg`$ (rel $`S_{}^n`$).
We will use the following lemma in the proof of Theorem 3.32:
###### Lemma 3.33.
Let $`f,g:S^n^{2n1}`$ be locally generic maps such that $`S(f)=S(g)\text{int}(S_+^n)`$ and $`f|S_{}^n=g|S_{}^n`$ are immersions. If $`f_sg`$ then $`f_sg`$ (rel $`S_{}^n`$).
###### Proof.
Choose a point $`xS_{}^n`$ and denote the differential $`(df)_x=(dg)_x`$ by $`j`$. Let $`S=S(f)=S(g)`$ and
$$=\{h\text{Imm}(S_{}^n,^{2n1}):(dh)_x=j\},$$
moreover
$$=\{hL(S^n,^{2n1}):S(h)=S\text{and}(dh)_x=j\}.$$
We define a map $`p:`$ by the formula $`p(h)=h|S_{}^n`$ for $`h`$.
We will now show that $`p`$ has the covering homotopy property. For this end choose a subset $`BS_+^n`$ such that $`B`$ is diffeomorphic to an open $`n`$-disc centered at the North pole $`N`$ of $`S^n`$, moreover $`SB`$. Denote by $`R`$ the intersection of $`S_+^nB`$ with the geodesic arc connecting $`N`$ and $`x`$. Let $`K=S_+^n(BR)`$, then $`\text{int}(K)`$ is diffeomorphic to an open $`n`$-disc. Let $`P`$ be a polyhedron, $`\psi :P\times I`$ and $`\varphi _0:P\times \{0\}`$, such that $`\psi |(P\times \{0\})=p\varphi _0`$. We must extend $`\varphi _0`$ to $`\varphi :P\times I`$ so that $`\psi =p\varphi `$, i.e., the following diagram is commutative:
$$\begin{array}{ccc}P\times I& \stackrel{\varphi }{}& \\ \text{id}& & p& & \\ P\times I& \stackrel{\psi }{}& .\end{array}$$
For $`rP`$ and $`tI`$ we define $`\varphi (r,t)|(\overline{B}R)`$ to be equal to $`\varphi _0(r)|(\overline{B}R)`$. Moreover, let $`(d\varphi (r,t))|R=(d\varphi _0(r))|R`$. This is possible since $`(\overline{B}R)S_{}^n=\{x\}`$ and $`(d\psi (r,t))_x=j`$ for every $`rP`$ and $`tI`$. The map $`\varphi _0(r)|K`$ is an immersion for every $`rP`$ and $`\varphi (r,t)`$ is already defined (together with normal derivatives) on $`K`$ for every $`rP`$ and $`tI`$. Thus using Smale’s lemma (Theorem 2.2) $`\varphi (r,t)`$ can be extended to $`K`$ as an immersion for every $`rP`$ and $`tI`$.
It is well known that the space $``$ is contractible, thus its homotopy groups are all trivial. Take the fiber of $`p`$ defined by $`=p^1(f|S_{}^n)`$, then by assumption $`f,g`$. Denote by $`i`$ the inclusion of $``$ into $``$. Let us look at the following part of the homotopy exact sequence of the fibration $`p:\stackrel{}{}:`$
$$\pi _1()\pi _0()\stackrel{i_{}}{}\pi _0()\stackrel{p_{}}{}\pi _0().$$
Since $`\pi _1()=0`$ and $`\pi _0()=0`$, the morphism $`i_{}`$ is an isomorphism. By assumption $`f_sg`$, moreover, a singularity fixing regular homotopy $`H`$ connecting $`f`$ and $`g`$ can be chosen so that $`(dH_t)_x=j`$ for every $`t[0,1]`$. Thus $`i(f)`$ and $`i(g)`$ lie in the same path-component of $``$, i.e., they represent the same element of $`\pi _0()`$. Using the fact that $`i_{}`$ is a monomorphism we get that $`f`$ and $`g`$ lie in the same path-component of $``$. This means that there exists a singularity fixing regular homotopy connecting $`f`$ and $`g`$ that is fixed on $`S_{}^n`$. ∎
### 3.6. Merging double curves
The idea of the procedure of eliminating the second obstruction (and achieving the coincidence of $`f`$ and $`g`$ – after a regular homotopy – also on the last part of $`M^n`$, i.e., on $`U`$) is as follows: In Ekholm showed that the regular homotopy class of an immersion $`S^n^{2n1}`$ is completely determined by the behavior of the map in a small neighborhood of the double curves if $`n4`$. We are going to reduce the problem to this result of Ekholm. First we make $`f`$ and $`g`$ coincide in a neighborhood of their singular sets. Then using Theorem 2.36 we will get to the case of the immersions of Ekholm’s theorem. To manipulate the double curves of $`f`$ and $`g`$ we will us the following construction.
In the presence of Whitney umbrella singular points we can merge the closed double curves with the double curves connecting the singular points. The merging is done by using the following construction motivated by the Whitney-trick (for a reference see or ). (Let us recall that in the Whitney trick one defines a standard model and then one shows that it can be embedded into the manifold under consideration.)
###### The construction.
In our case the standard model $`S`$ is chosen to be the standard model of the Whitney-trick (two arcs in $`D^2`$ that intersect in two points) times the interval $`[1,1]`$ (see the left side of Figure 3). Thus we have two surfaces $`P_1`$ and $`P_2`$ in $`D^2\times [1,1]`$, both diffeomorphic to the square, that intersect in two line segments $`l_1`$ and $`l_2.`$ The regular homotopy in this model is the identity near $`(D^2\times [1,1])`$ and is the isotopy of the standard model of the Whitney trick ($`l_2`$ is pushed above $`l_1`$) near $`D^2\times \{0\}`$. As a result of this deformation the double points are removed near $`D^2\times \{0\}`$, and thus the double curves of the final arrangement coincide with the connected sum $`l_1\mathrm{\#}l_2`$ (see the right side of Figure 3).
Let $`p`$ and $`q`$ be two double points of the generic map $`f:N^n^{2n1}`$, where $`N^n`$ is any compact connected manifold with boundary. All we have to do is to embed $`D^2\times [1,1]\times ^{n2}\times ^{n2}`$ into $`^{2n1}`$ so that $`P_1\times ^{n2}\times \{0\}`$ and $`P_2\times \{0\}\times ^{n2}`$ map to $`\text{Im}(f)`$, moreover $`l_1\times \{0\}`$ maps on the double curve through $`p`$ and $`l_2\times \{0\}`$ maps on the double curve through $`q`$.
Let $`f^1(p)=\{p_1,p_2\}`$ and $`f^1(q)=\{q_1,q_2\}`$. Choose an embedded curve $`c_i`$ in $`N^n`$ connecting $`p_i`$ and $`q_i`$ and let $`C_i=f(c_i)`$ for $`i=1,2`$. We might suppose that $`C_1C_2=\{p,q\}`$ and no other double or singular point of $`f`$ lies on $`c_1`$ and $`c_2`$. Thus $`f|c_1`$ and $`f|c_2`$ are embeddings. Let $`D_i`$ be an embedded (tubular) neighborhood of $`C_i`$ in $`f(N^n)`$ diffeomorphic to $`D^n`$ for $`i=1,2`$. The intersection $`D_1D_2=L_1L_2`$, where $`pL_1`$ and $`qL_2`$.
Next we choose a Whitney disk $`W`$ containing $`C_1`$ and $`C_2`$ as in . Since $`dim(\text{Im}(f))+dim(W)=n+2<2n1`$ (because $`n4`$), if $`W`$ is in general position then $`W\text{Im}(f)=C_1C_2`$. We might also suppose that $`W`$ is transversal to $`L_1`$ and $`L_2`$. Since from now on we will work only in a neighborhood of $`W`$, we might restrict our attention to $`D_1D_2`$ and forget about the rest of $`\text{Im}(f)`$.
Construct a Riemannian metric on $`^{2n1}`$ for which $`D_1D_2`$. Choose an orientation of $`D_1`$ and a co-orientation of $`D_2`$ in $`^{2n1}`$. Let
$$\alpha _1,\mathrm{},\alpha _{n1}\nu _p(D_2^{2n1})T_pD_1$$
be a positive normal frame of $`D_2`$. Let $`\stackrel{~}{\nu }(p)T_pL_1`$ be taken so that $`\stackrel{~}{\nu }(p),\alpha _1,\mathrm{},\alpha _{n1}`$ is a positive basis of $`T_pD_1`$. Similarly, let
$$\beta _1,\mathrm{},\beta _{n1}\nu _q(D_2^{2n1})T_qD_1$$
be a positive normal frame of $`D_2`$ and $`\stackrel{~}{\nu }(q)T_qL_2`$ be chosen so that $`\stackrel{~}{\nu }(q),\beta _1,\mathrm{},\beta _{n1}`$ is a *negative* basis of $`T_qD_1`$.
We can extend $`\stackrel{~}{\nu }`$ to $`C_1C_2`$ so that $`\stackrel{~}{\nu }|C_i`$ is tangent to $`D_i`$ since $`n3`$. Next we extend $`\stackrel{~}{\nu }`$ from $`C_1C_2`$ to a normal field of $`W`$: Since $`\nu (W^{2n1})`$ is the trivial bundle and $`C_1C_2`$ is homotopy equivalent to $`S^1`$, the obstruction to extend $`\stackrel{~}{\nu }`$ is an element of $`\pi _1(V_{2n3,1})=\pi _1(S^{2n4})`$. But this group is $`0`$ since $`2n4>1`$ if $`n3`$. Taking the exponential map along $`\stackrel{~}{\nu }`$ on $`W`$ we can embed the standard model $`S=D^2\times [1,1]`$ into $`^{2n1}`$. Denote the image of $`S`$ by $`U`$.
Orient $`C_1`$ from $`p`$ to $`q`$ and denote by $`\tau (x)`$ the positive unit tangent vector of $`C_1`$ at $`x`$. Moreover, let $`\nu (p)=\tau (p)`$ and $`\nu (q)=\tau (q)`$. Since $`n4`$ we can extend $`\nu `$ along $`C_2`$ as a normal field of $`D_2`$.
Choose vectors $`\xi _1(p),\mathrm{},\xi _{n2}(p)`$ in $`T_pD_1`$ so that $`\stackrel{~}{\nu }(p),\nu (p),\xi _1(p),\mathrm{},\xi _{n2}(p)`$ is a positive basis of $`T_pD_1`$. Using parallel translation extend $`\xi _1,\mathrm{},\xi _{n2}`$ to $`C_1`$. Then
$$\stackrel{~}{\nu }(q),\nu (q),\xi _1(q),\mathrm{},\xi _{n2}(q)T_qD_1$$
is a negative basis since $`\nu (q)=\tau (q)`$. Thus from the definition of $`\stackrel{~}{\nu }(q)`$ we get that $`\nu (q),\xi _1(q),\mathrm{},\xi _{n2}(q)`$ is a positive basis of $`\nu _q(D_2^{2n1})`$. Moreover, $`\nu (p),\xi _1(p),\mathrm{},\xi _{n2}(p)`$ is a positive basis of $`\nu _p(D_2^{2n1})`$ and $`\nu |C_2`$ is a continuous vector field. All this together implies that the $`(n2)`$-frames $`\underset{¯}{\xi }(p)`$ and $`\underset{¯}{\xi }(q)`$ can be extended onto $`C_2`$ to be transversal to $`D_2`$ and $`\nu `$.
We can extend the frame $`\underset{¯}{\xi }`$ to $`U`$: The pair $`(U,C_1C_2)`$ is homotopy equivalent (deformation retracts onto) $`(D^2,S^1)`$. Thus the obstruction lies in $`\pi _1(V_{2n4,n2})=0`$ since $`n4`$ implies that $`1<(2n4)(n2)`$.
Finally we construct an $`(n2)`$-frame $`\underset{¯}{\eta }`$ on $`U`$ that is tangent to $`D_2`$ on $`C_2`$. Look at the bundle of orthogonal $`(n2)`$-frames on $`U`$ that are perpendicular to $`U`$ and $`\underset{¯}{\xi }`$. Since $`U`$ is contractible, this bundle is trivial; let $`\underset{¯}{\eta }`$ be any section. Notice that $`\nu (UD_2U)`$ is spanned by the vector field $`\nu |C_2`$ (the subset $`UD_2`$ corresponds to $`P_2`$ in the standard model $`S`$). Because $`\underset{¯}{\eta }|C_2\underset{¯}{\xi },U`$, we get that
$$\underset{¯}{\eta }|C_2\underset{¯}{\xi },\nu |C_2=\nu (D_2^{2n1})|C_2.$$
Thus $`\underset{¯}{\eta }|C_2`$ is tangent to $`D_2`$, as required.
Using the exponential map on the fields $`\underset{¯}{\xi }`$ and $`\underset{¯}{\eta }`$ along $`U`$ we can embed $`S\times ^{n2}\times ^{n2}`$ into $`^{2n1}`$ as described above. ∎
The regular homotopy constructed above is clearly singularity fixing.
### 3.7. The proof of Theorem 3.32
Using Lemma 3.33 we only have to prove that $`f_sg`$. With a singularity fixing regular homotopy we can perturb $`f`$ and $`g`$ to obtain generic maps. Let us recall (see the beginning of the previous section) that the main idea of the proof is to reduce the problem to the case of immersions and then apply Ekholm’s theorem.
Let $`l_f`$ (respectively $`l_g`$) denote the closure of the double curve of $`f`$ (respectively $`g`$) that connects the two singular points of $`f`$ (respectively $`g`$). With the aid of the above construction we merge using a singularity fixing regular homotopy all closed double curves of $`f`$ (respectively $`g`$) with $`l_f`$ (respectively $`l_g`$). From now on we will suppose that $`l_f`$ is the only double curve of $`f`$ and $`l_g`$ is the only double curve of $`g`$. Let $`m_f=f^1(l_f)`$ and $`m_g=g^1(l_g)`$. Since $`n4`$ there exists a diffeotopy $`\{d_t:t[0,1]\}`$ of $`S^n`$ that fixes $`S(f)=S(g)`$ and takes $`m_g`$ to $`m_f`$. Composing $`f`$ with $`\{d_t\}`$ we can suppose that $`m_f=m_g`$. Let $`m=m_f=m_g`$ and $`S(f)=S(g)=\{p,q\}`$. If $`n`$ is even we have to be careful when choosing the diffeotopy $`\{d_t\}`$: If we order the two components of $`m\{p,q\}`$ then this defines signs of the singular points of $`f`$ and $`g`$ (see Definition 2.37), and it might happen that $`S_+(f)S_+(g)`$ (the equality $`\mathrm{\#}S(f)=\mathrm{\#}S(g)=0`$ still holds). In this case we compose $`f`$ with a diffeotopy of $`S^n`$ that fixes $`p`$ and $`q`$ and swaps the two components of $`m\{p,q\}`$ in order to swap $`S_+(f)`$ and $`S_{}(f)`$. This is possible if $`n4`$. Thus if $`n`$ is even we also fix an ordering of the components of $`m\{p,q\}`$ and suppose that for the induced signs $`S_\pm (f)=S_\pm (g)`$.
Next we compose $`f`$ with a diffeotopy of $`^{2n1}`$ that moves $`l_f`$ to $`l_g`$. Thus we can suppose that $`f|m=g|m`$ and let $`l=l_f=l_g`$. Using Lemma 2.14 of we can straighten $`f`$ and $`g`$ close to their self intersection $`\text{int}(l)`$, i.e., we can achieve that they agree with their normal derivatives on a neighborhood of $`m`$.
We now sketch the idea of the construction that follows: The singular points of $`f`$ and $`g`$ coincide and have the same signs if $`n`$ is odd. The sign is the only local isotopy invariant a Whitney-umbrella singularity can have. Moreover, any two bundles over an interval are equivalent. Since $`m`$ is an orientation preserving curve in $`S^n`$ the two cross-caps at the end of $`l_f`$ and $`l_g`$ are in the same ”relative position”. So $`f`$ and $`g`$ are isotopic if restricted to a small neighborhood of $`m`$, fixing $`S(f)=S(g)`$. This isotopy can be extended to an ambient isotopy of $`^{2n1}`$. So we might suppose that $`f`$ and $`g`$ agree in a small open tubular neighborhood $`V`$ of their common double locus $`m`$. The exact details are as follows:
Since $`S^n`$ is orientable, the normal bundle $`\nu (mS^n)`$ is trivial. Thus we can choose a normal framing $`\underset{¯}{\nu }`$ of $`m`$ in $`S^n`$. Notice that $`\mathrm{ker}(df_p)=\mathrm{ker}(dg_p)=T_pm`$ and $`\mathrm{ker}(df_q)=\mathrm{ker}(dg_q)=T_qm`$, so $`df(\underset{¯}{\nu }_p)`$ and $`dg(\underset{¯}{\nu }_p)`$ are $`(n1)`$-frames in $`T_{f(p)}^{2n1}=T_{g(p)}^{2n1}`$. Similarly, $`df(\underset{¯}{\nu }_q)`$ and $`dg(\underset{¯}{\nu }_q)`$ are $`(n1)`$-frames in $`T_{f(q)}^{2n1}=T_{g(q)}^{2n1}`$. Moreover, $`\underset{¯}{\nu }_f=df(\underset{¯}{\nu }|(m\{p,q\}))`$ and $`\underset{¯}{\nu }_g=dg(\underset{¯}{\nu }|(m\{p,q\}))`$ provide $`(2n2)`$-frames along $`\text{int}(l)`$.
Suppose that $`S_+(f)=S_+(g)=\{p\}`$ and $`S_{}(f)=S_{}(g)=\{q\}`$. Let $`a=f(p)=g(p)`$ and $`b=f(q)=g(q)`$, then $`a`$ and $`b`$ are the endpoints of the curve $`l`$. Orient $`l`$ from $`a`$ to $`b`$, this induces an orientation of $`\nu (l^{2n1})`$. (Recall that we fixed an orientation of $`^{2n1}`$ to define the signs of the Whitney-umbrellas.) Then from Definition 2.37 of the signs of the singular points we can see that $`\{\underset{¯}{\nu }_f,\underset{¯}{\nu }_g\}`$ provides a negative basis of $`\mu =\nu \left(\text{int}(l)^{2n1}\right)`$. Thus there exists a homotopy $`T:\text{int}(l)\times IGL_+(2n2,)`$ such that for every $`x\text{int}(l)`$ the transformation $`T(x,0)=\text{id}_{^{2n2}}`$ and $`(\underset{¯}{\nu }_f)_xT(x,1)=(\underset{¯}{\nu }_g)_x`$. The homotopy $`T`$ can be extended to an ambient isotopy of $`f`$ that takes $`f`$ to $`g`$ on a small open tubular neighborhood $`V`$ of $`m`$, and this isotopy is identical outside a neighborhood of $`l`$. Thus we can suppose that $`f|V=g|V`$.
Let $`DV`$ be diffeomorphic to the closed $`n`$-disc $`D^n`$, so that $`p,q\text{int}(D)`$ and $`f|D`$ ($`=g|D`$) is an embedding. Since the only double locus of $`f`$ and $`g`$ is $`mV`$, the distances $`d_1=d(f(D),f(S^nV))`$ and $`d_2=d(g(D),g(S^nV))`$ are positive. Let $`0<\epsilon <\mathrm{min}(d_1,d_2)`$. Using Theorem 2.36 choose a generic immersion $`h:D^{2n1}`$ that agrees with $`f|D`$ in a neighborhood of $`D`$, moreover the $`C^0`$-distance of $`h`$ and $`f|D`$ is $`<\epsilon `$. (This is possible since $`\mathrm{\#}S(f|D)=0`$ and $`f|D`$ is an embedding.) Denote by $`f_1`$ (respectively $`g_1`$) the immersion of $`S^n`$ in $`^{2n1}`$ that agrees with $`f`$ (respectively $`g`$) on $`S^nD`$ and with $`h`$ on $`D`$. From the choice of $`h`$ we can see that $`f(S^nV)h(D)=\mathrm{}`$ and $`g(S^nV)h(D)=\mathrm{}`$. Moreover, since $`f|(S^nV)`$ and $`g|(S^nV)`$ are embeddings that do not intersect $`f(VD)=g(VD)`$, the double loci of $`f_1`$ and $`g_1`$ are contained in $`V`$. But $`f_1|V=g_1|V`$, thus $`f_1`$ and $`g_1`$ agree in a neighborhood of their common double locus. Theorem 1 in implies that the regular homotopy class of a generic immersion $`S^n^{2n1}`$ can be expressed in terms of the geometry of the self intersection if $`n4`$ . Thus $`f_1g_1`$.
Since $`f_1|D=g_1|D`$, the regular homotopy connecting $`f_1`$ and $`g_1`$ can be chosen to be constant on $`D`$, i.e., $`f_1g_1`$ (rel $`D`$). This gives us a regular homotopy $`fg`$ (rel $`D`$), thus $`f_sg`$.
### 3.8. An application
In part 4 of the $`n=2`$ case of Theorem 1.7 was applied to projections of immersions of surfaces into $`^4`$. Knowing Theorem 1.7 the same results can be generalized without modification of the proof.
Let us first recall some definitions.
###### Definition 3.34.
Let
$$\text{Imm}_\pi (M^n,^{2n})=\{F\text{Imm}(M^n,^{2n}):\pi FL(M^n,^{2n1})\}$$
be the subspace of $`\text{Imm}(M^n,^{2n})`$ formed by those immersions whose projection in $`^{2n1}`$ is locally generic. Two immersions $`F,G\text{Imm}_\pi (M^n,^{2n})`$ are called *$`\pi `$-homotopic* (denoted by $`F_\pi G`$) if they are in the same path-component of $`\text{Imm}_\pi (M^n,^{2n})`$.
The generalization of Theorem 4.7 of is then the following.
###### Theorem 3.35.
If $`n=2`$ or $`n>3`$ and $`F,G\text{Imm}_\pi (M^n,^{2n})`$ then
$$F_\pi G[FG\text{and}\pi F\pi G].$$
## Acknowledgement
I would like to take this opportunity to thank Professor Tobias Ekholm for valuable advice and Professor András Szűcs for reading earlier versions of this paper and making several useful remarks.
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# A Study of Edge-On Galaxies with the Hubble Space Telescope’s Advanced Camera for Surveys. II. Vertical distribution of the resolved stellar population
## 1. Introduction
Galactic structure has long been recognized as a key constraint on theories of galaxy formation. The presence of two distinct components – disks and spheroids – suggests that at least two separate physical mechanisms were active during the formation of spiral galaxies. Within the Milky Way, however, detailed analyses of the kinematics and chemistry of nearby stars have revealed additional components: a thin disk with scale heights dependent on age, a thick disk and a stellar halo. Each of these components place unique constraints on the Galaxy’s evolution. In particular, the ages and metallicities of the old thin disk, thick disk and halo indicate that they are remnants of the initial stages in the assembly of the galaxy and thus constrain the early evolution of galaxies.
Several scenarios for the creation of the old thin and thick disks are consistent with current observations in the Milky Way. These scenarios can be categorized into three main types: (1) the creation of a thick disk from a thin disk by heating by molecular clouds, spiral arms, star formation or accretion events (e.g. Spitzer & Schwarzschild, 1951; Barbanis & Woltjer, 1967; Lacey, 1991; Kroupa, 2002; Quinn et al., 1993; Gnedin, 2003). (2) the slow collapse of the proto-Galaxy forming the thick and thin disk in succession (e.g. Eggen et al., 1962; Gilmore, 1984), and (3) the formation of a thick disk from mergers, either by direct accretion of stars or by in situ formation from accreted gas (e.g. Bekki & Chiba, 2001; Gilmore et al., 2002; Abadi et al., 2003; Brook et al., 2004). In this last scenario the thin disk would likely form by a settling of gas from the merger events. It is widely assumed that the stellar halo forms from accreted satellites (see review by Freeman & Bland-Hawthorn, 2002). Determination of which processes are important in galaxy formation requires the study of stellar components in other galaxies.
Unfortunately, detailed analyses of older stellar components are difficult outside the Milky Way. Some thick disks and stellar halos have been identified using broad-band surface photometry (e.g. Burstein, 1979; Tsikoudi, 1979; Dalcanton & Bernstein, 2002; Pohlen et al., 2004; Zibetti et al., 2004), but their low surface brightness precludes all but a most cursory study of their structure and stellar content. Recently, HST imaging has begun to allow richer analyses of the resolved stars in the thick disk and halo (Brown et al., 2003; Tikhonov et al., 2005; Mould, 2005). The resulting color-magnitude diagrams can provide robust constraints on the metallicities and ages of the stars which dominate these old components in other galaxies.
In this paper we take advantage of a large database of ACS imaging of nearby edge-on galaxies to analyze the stellar populations of the galaxies as a function of height above the galaxy midplane (“disk height”). The wide field of view of ACS allows us to track changes in the color-magnitude diagram to more than 3 kpc above the plane. Moreover, our sample galaxies have much lower masses than the Milky Way, and thus constrain how the extraplanar stellar content varies with galaxy mass. The properties of the sample and their global color-magnitude diagrams are given in Seth et al. (2005) (Paper I). In this paper, we study the spatially resolved stellar populations in eight fields of six galaxies that were close enough to determine accurate Tip of the Red Giant Branch (TRGB) distances. In particular, we show that there is a significant stellar population high above the plane of these galaxies and that these stars are old and metal-poor, comparable to high-latitude stellar populations in the Milky Way.
In §2, we review the results of Paper I for these six galaxies and describe the completeness corrections used to correct the stellar number counts in the remainder of the paper. In §3 we show that we trace stars in the host galaxies out to many scale heights and that these stars have a larger scale height than expected from fits to 2MASS $`K_s`$ band images. In §4 we select three populations of stars from color-magnitude diagrams (CMDs) and show that the scale height of a stellar population increases with age. We derive metallicity distribution functions for the extended stellar component in §5 and then conclude by discussing our results in a broader context in §6.
## 2. Review of the Galaxy Properties
Table 1 shows the position, type, maximum circular velocity ($`V_{\mathrm{max}}`$), distance modulus ($`mM_{\mathrm{TRGB}}`$), scale length ($`h_r`$) and the $`K_s`$ band half-light height ($`z_{1/2}`$) of the six galaxies we will discuss in this paper. The latter two parameters were determined from 2-D model fits to 2MASS $`K_s`$ band data, presented in Paper I. The vertical component of these models is defined using the distribution of an isothermal population of stars (van der Kruit & Searle, 1981):
$$\mathrm{\Sigma }(z)\mathrm{sech}^2\left(\frac{z}{z_0}\right)$$
(1)
where $`\mathrm{\Sigma }(z)`$ is the surface brightness or density at a position $`z`$ above the midplane, and $`z_0`$ is the scale height. We will use this functional form to fit vertical distributions throughout this paper. Note, however, that this is one of many equations commonly used to describe the vertical distribution of stars in galaxies, a good overview of which can be found in (Pohlen et al., 2000). Most of these functional forms vary near the midplane, but have similar exponential declines at large disk heights. When comparing galaxies in this paper, the disk heights will be normalized by the $`z_{1/2}`$ parameter, which gives the height containing 50% of the $`K_s`$ band light in the model fits. It is related to $`z_0`$, $`z_{1/2}`$=0.549 $`z_0`$, and is similar to the exponential scale height $`h_z`$, which at large scale heights is equal to $`\frac{1}{2}z_0`$ (van der Kruit & Searle, 1981). The values of $`z_{1/2}`$ for the six galaxies range from 160-280 pc. For comparison, the Milky Way thin disk has exponential scale heights ranging from $``$100 pc at young ages (Schmidt, 1963) to 330 pc (Chen et al., 2001) at older ages.
All six galaxies are within 8 Mpc and are type Sc or later. The maximum circular velocities are all below 135 km/sec suggesting that these objects are closer in mass to the LMC than to the Milky Way. We note that all the galaxies except NGC 4631 have circular velocities well below 120 km/sec, which appears to mark a break in the dust properties (Dalcanton et al., 2004) and current metallicity (Garnett, 2002) in spiral galaxies. The scale lengths of our sample galaxies range from 0.9 to 1.6 kpc, $``$2-3 times smaller than the Milky Way scale length<sup>2</sup><sup>2</sup>2The Milky Way thin & thick disk scale lengths are rather uncertain with recent values ranging between 2-4 kpc (e.g. Ng et al., 1996; Mendez & van Altena, 1998; Ojha, 2001), we will use a scale length of 3 kpc for comparisons to the Milky Way made in this work.. None of the galaxies has an apparent bulge component, although NGC 4244 does have a prominent central stellar cluster, which is clearly visible in the ACS and 2MASS $`K_s`$ band images.
Our observations include eight HST/ACS fields in the six galaxies shown in Table 1 and are described fully in Paper I (Seth et al., 2005). These observations were obtained to study the dust-lane properties of these galaxies, and hence are centered on the galaxy midplane. The dimensions of these eight fields are given in Table 2, which shows the minimum and maximum disk radius in terms of the scale length and the minimum and maximum disk height in terms of $`z_{1/2}`$. In general, each field is located close to the center of the galaxy. However, two of the galaxies, NGC 55 and NGC 4631 have additional fields located further out in the disk that are given a ’-DISK’ suffix. Note that many of the galaxies lie diagonally across the chip meaning that the extremities of the ranges given in Table 2 are not well sampled. In this paper, we focus on the vertical distributions of the stars, and where not otherwise noted, analyze all data at the same disk height together.
The approach proposed above is valid as long as the scale height of disk components does not vary substantially with radius, an assumption that has been verified through observations of edge-on galaxies (e.g. van der Kruit & Searle, 1981; Pohlen et al., 2004). We note however, that there are some analyses which indicate that the scale height of galaxies flares with increasing radius, both in our own Galaxy (López-Corredoira et al., 2002) and in edge-on galaxies (de Grijs & Peletier, 1997; Narayan & Jog, 2002). However, de Grijs & Peletier (1997) show that in their sample of edge-on galaxies, late-type galaxies such as those observed here have little or no flaring. In one of our galaxies, NGC 4244, flaring of the HI gas by a factor of $``$3 is seen between radii of 8-13 kpc (Olling, 1996). Even if similar flaring occurs in the stellar distribution, the star counts at the disk heights we consider here will still be dominated by stars located (radially) near the center of the galaxy. We note that we do see some evidence for modest flaring in our two ’-DISK’ fields which lie at large radii (see §4.2.1).
In Paper I we presented the color-magnitude diagrams (CMDs) for each of the galaxy fields discussed here. These are reproduced in Figure 1 which shows the F606W-F814W color vs. the F814W magnitude. All magnitudes in this paper are given in the VEGAmag system (see Paper I for details). Each of the CMDs show plumes at F606W-F814W colors of $``$0.1 and $``$1.3 which are young main sequence (MS) and helium burning (HeB) stars respectively. An old RGB can be seen in each galaxy at colors of $``$1 and an intermediate-age asymptotic giant branch (AGB) extends from the tip of the RGB to redder colors. Further details on the components can be found in §4.1 and in Paper I. The boxes shown in Figure 1 isolating these components will be discussed in greater detail in §4.
### 2.1. Completeness Corrections
As described in Paper I, we conducted extensive artificial star tests to characterize the completeness in each field as a function of magnitude, color and location. The galaxies in our sample have high surface brightnesses near their midplane making the completeness there much lower (see Fig. 2 from Paper I). The goal of this paper is to analyze the vertical distribution of stellar populations, and therefore correcting for this varying completeness is critical.
For each field, over 5 million artificial stars were inserted at random positions, with magnitudes between F606W of 18 and 29, and with F606W-F814W colors between -1 and 3. These stars were then run through the same pipeline used to determine the stellar photometry. Artificial stars that coincided with actual sources were considered detected only if the input magnitude of the artificial star was brighter than the actual source in both bands.
To enable completeness corrections for individual stars, we determined our completeness for the artificial stars in bins of magnitude, color and local surface brightness. For the magnitude and color, we used 0.15 magnitude wide bins. At its steepest, the completeness as a function of magnitude varies by at most 6% over 0.15 magnitudes, so any error introduced by the binning should be smaller than that. We then determined the size of the surface brightness bins such that there would be $``$50 stars in each final bin. Determining the completeness from 50 stars gives a random error in the completeness of $``$6%. An aperture around each star from 11 to 14 pixels was used to determine the local surface brightness. We determined the surface brightness level from the F606W image so as to include the effects of the HII region emission visible in a number of galaxies (most notably NGC 55, which therefore has the brightest completeness level in Fig. 2). This emission is not seen at all in the F814W images, due to the lack of strong emission lines in F814W bandpass. Using this binned completeness function, we are able to determine the completeness of any individual star based on its color, magnitude and local surface brightness level to within $``$10%. In the stellar density profiles presented in §3 & §4, the completeness corrections are up to 200% near the midplane, but fall to $`30\%`$ at $`z/z_{1/2}>3`$.
In addition to correcting for the completeness, magnitude limits must also be set to insure that we are not using stars fainter than we can detect in the higher surface brightness regions of the image (i.e. the midplane). We therefore choose to limit our analysis to regions of the CMD that fall above a conservative 20% completeness limit in regions of high surface brightness. Figure 2 shows the 20% completeness limits for the brightest regions in each field. As can be seen, the completeness limit rises towards redder colors, and steepens at colors redder than F606W-F814W of 1. To make a smooth boundary that is easily applied to our data, we fit the 20% completeness curves to two lines intersecting at F606W-F814W=1. Table 2 shows the results of these fits, by giving the completeness limit at F606W-F814W of -1 (F814W<sub>lim,-1</sub>), 1 (F814W<sub>lim,1</sub>), and 3 (F814W<sub>lim,3</sub>). We use these limits throughout this paper to insure that comparisons made between stellar populations at different disk heights are valid.
We note here that although we can correct for incompleteness due to crowding, we cannot correct for the attenuation of stars by dust. We will show in §4 that all the galaxies in our sample are optically thick near their midplanes. Therefore at low galactic latitudes, our completeness corrected stellar census will fall short of the true number of stars.
## 3. Vertical Distribution of Stars
We demonstrate in this section that there are significant numbers of stars well above the planes of all our disks and that the profiles of these stars do not fit the profiles expected from the ground-based $`K_s`$ band galaxy fits.
In Figure 3 we present the completeness corrected surface density profiles of all the detected stars ($`\mathrm{\Sigma }_{\mathrm{all}}`$) above the completeness limits given in Table 3. Two lines are shown for each galaxy, giving the profile on both sides of the disk. To determine the surface density, stars were binned as a function of scale height. After binning, the completeness-corrected number of stars was divided by the area of the bin to obtain the surface density. Only bins with an area of more than 300 arcsec<sup>2</sup> are plotted.
Figure 3 shows that we trace the stellar component of each galaxy out to large disk heights, with several galaxies being traced out beyond 10$`z_{1/2}`$. The profiles in general are fairly symmetric, the most notable exception being NGC 4631 and NGC 4631-DISK. These fields are contaminated on one side (shown using dot-dashed lines) by the presence of the companion galaxy NGC 4627 (see Fig. 3 from Paper I). The decrease of the profiles with increasing scale height out to the edge of the fields strongly suggests that the profiles remain above the surface density of foreground Galactic stars and background unresolved dwarf galaxies in our magnitude range. Only above 10$`z_{1/2}`$ in IC 5052 and NGC 5023 is there some evidence for the leveling off that would be expected as we reach the foreground/background level. Note that the galaxies are all at galactic latitudes above 35. We therefore can safely assume that a vast majority of stars in our images are located in the host galaxies and we make no corrections for foreground/background sources.
The dashed lines in Figure 3 show sech<sup>2</sup> profiles with scale heights one and two times the measured $`K_s`$ band scale height (note that because the plot is scaled by the $`K_s`$ band scale height these profiles are the same for each galaxy). As described in detail in Paper I, models were fit using a Levenberg-Marquardt algorithm with uniform weighting on all unmasked pixels to $`K_s`$ band 2MASS data. Because the 2MASS surface photometry has a limiting isophote of $`K_s20.0`$ mag/arcsec<sup>2</sup>, only relatively high surface brightness features could be fit. The typical $`K_s`$ band peak surface brightness of our galaxies is $``$18 (Paper I, Table 2), which means that galaxies are only detected in the $`K_s`$ images out to a few $`z_{1/2}`$ from the midplane. Although the $`K_s`$ band light is often thought to be dominated by the RGB stars that trace an old stellar population, we will show in §4 that in these low mass, late-type galaxies, it more closely traces the young and intermediate age populations and is thus dominated by red supergiant and AGB stars.
Figure 3 shows that the outer portions of the stellar density profiles of the galaxies appear to be broader than indicated by the $`K_s`$ band scale height. To quantify this, we fit the stellar density profiles between 5$`z_{1/2}`$ and 10$`z_{1/2}`$ on each side of the midplane in each galaxy. Only profiles with data beyond 8$`z_{1/2}`$ were fit (thus excluding NGC 55, NGC 55-DISK and one side of NGC 4244). In the Milky Way, the stellar profile deviates from the thin disk profile beyond $``$1 kpc above the plane (Gilmore & Reid, 1983). For our galaxies, 5$`z_{1/2}`$ is $``$1 kpc. Therefore we would expect our fitting range to be sensitive to a possible thick disk component in these galaxies. The scale heights of these fits above and below the midplane are shown in the third and fourth columns (respectively) of Table 4. All of the fitted stellar profiles are significantly broader ($`\times `$1.5-2.4) than the $`K_s`$ band scale height. This observation strongly suggests that these galaxies contain a more broadly distributed stellar component not traced by the $`K_s`$ band 2MASS images.
We note that the fitted scale heights, both in the $`K_s`$ band and using stellar density profiles, can differ from the true scale height due to a number of factors. First, dust attenuation can obscure light near the midplane of the galaxy. In the $`K_s`$ band this attenuation should be small, and we ameliorate this problem in the fits presented here by avoiding the midplane. Second, the galaxies may not be perfectly edge-on. Based on previous observations, the least inclined of the galaxies is NGC 55, which has a $``$80 inclination (Hummel et al., 1986). This would give a fitted scale height $``$30% greater than the intrinsic disk scale height. We also note that NGC 55 and NGC 4631 are fairly irregularly shaped making the fits to these galaxies less reliable than for the other four galaxies.
## 4. Variation in distribution with stellar population
We now turn our analysis to stars selected in regions of our CMDs that isolate stellar populations with different ages. Using this method we show that the older stellar populations have an increasing scale height. In §4.3, we examine the variation of scale height in the context of disk heating models. We then present simplistic dust models in §4.4 and show that the dust extinction in these galaxies is distributed in a component that is broader than the young stellar populations.
### 4.1. Selection of CMD regions
We attempt to separate our data into young, intermediate, and old stellar populations by selecting stars from different regions in the CMD. The young stars are found in the Main Sequence (MS) component and in the red and blue Helium Burning (RHeB,BHeB) sequences (see Paper I, Fig. 1 for a schematic CMD), all of which should contain stars under a few 100 Myr in age. For the intermediate age stars we select AGB stars brighter than and redward of the RGB, resulting in ages ranging between a few 100 Myr and a few Gyr. Lastly, for the old population of stars we select RGB stars, which have ages in excess of 1 Gyr, although some AGB stars will also be found in the same region.
The actual regions used for the selection are shown in the CMDs in Fig. 1. The bottom right figure is a cartoon illustrating the selection of the MS/HeB, AGB and RGB regions. The RGB region was selected using lines with slopes of 3.3 and 6.6 and F606W-F814W colors of 1.0 and 1.6 at the tip of the RGB. The MS region was defined by taking all stars redward of -0.5 and blueward of a line with slope of 3.3 and a color of 0.7 at the TRGB magnitude. Finally, the AGB region isolates stars less than two magnitudes brighter than the TRGB magnitude and redwards of a line with slope of 3.3 and a color of 1.2 at the TRGB magnitude. These boundaries were combined with the TRGB magnitude given in Table 1 and the completeness limits from Table 3 to determine the final regions for each galaxy shown in Figure 1. The regions were chosen somewhat conservatively - e.g. we chose to put space between the MS/HeB section and the RGB so that there would be little overlap between the two due to dust extinction or large photometric errors at faint magnitudes. From here we will refer to the stars in these CMD boxes as the MS, AGB and RGB stellar populations.
#### 4.1.1 Synthetic CMDs
To determine the typical ages of stars detected in our CMD boxes and to facilitate quantitative comparisons between galaxies and their different stellar populations, we generated synthetic CMDs using the MATCH program (Dolphin, 2002) and the IAC-STAR program<sup>3</sup><sup>3</sup>3This work has made use of the IAC-STAR Synthetic CMD computation code, IAC-STAR is supported and maintained by the computer divison of the IAC. (Aparicio & Gallart, 2004), using isochrones of Bertelli et al. (1994) and Girardi et al. (2000) in both cases. The synthetic stars were generated assuming a constant star formation rate (SFR) from 13 Gyr ago to the present, and a metallicity that steadily increased from \[Fe/H\]$`=1.7`$ to -0.4 (Garnett, 2002). We used slightly different IMFs in the two CMDs. For the MATCH CMD a pure Salpeter IMF ($`\alpha =2.35`$) is assumed between 0.1 and 120 M, whereas in the IAC-STAR CMD we used the default Kroupa et al. (1993) IMF, which is steeper at the high mass end ($`\alpha =2.70`$).
To compare these CMDs to our observations, the synthetic stars were first transformed from Johnson V & I to VEGAmag F606W & F814W colors. We then mimicked observations of each galaxy as follows. First, each star was randomly assigned a surface brightness value based on the values of detected stars in each frame. Then, using the artificial star tests, we determined the chance each star was detected and a magnitude error based on the star’s initial F606W & F814W magnitudes (assuming the distance moduli shown in Table 1) and the surface brightness value. A final CMD was then made by randomly determining if each star was detected and applying the determined errors. The resulting CMDs looked qualitatively similar to our observed CMDs, with the most notable difference being an offset of the AGB stars to somewhat brighter magnitudes in both synthetic CMDs relative to the real data and a deficit of MS/HeB stars in the IAC-STAR CMDs relative to the MATCH and the real galaxy CMDs. This could indicate either that the galaxies have enhanced recent star formation or that their IMF is not as steep as the Kroupa et al. (1993) IMF on the high mass end.
Figure 4 shows the resulting age distribution of the MS, AGB, and RGB boxes in NGC 4144 using the MATCH and IAC-STAR synthetic CMDs. The age distributions for other galaxies are similar. This figure clearly demonstrates that we are separating the stars into young, intermediate-age and old populations with our chosen CMD boxes. However, the separation is not perfect. Each bin has significant overlaps with the others due to unavoidable photometric errors and to true overlap in the colors and magnitudes of stellar populations with different ages. The IAC-STAR CMDs have age distributions similar to the MATCH CMDs, but with the AGB populations weighted more towards older ages and a more significant contamination of old stars in the MS box (probably due to the relative lack of MS stars in the IAC-STAR models). Both effects likely result from the steeper IMF assumed for the Aparicio CMD. In the following sections of the paper, we use the MATCH CMD for comparisons with observations because it more closely reproduces the ratio of young MS and HeB stars relative to the number of older stars.
We note that these sunthetic CMDs assume a constant star formation rate and thus are not useful in determining true star formation histories. However, we will be able to use them to get a sense of relative star formation histories (SFHs) as a function of scale height and to get a rough sense of the ages of the stars in our CMD regions.
### 4.2. Stellar Density Profiles
Now we compare the surface density profiles of the MS, AGB and RGB stars to examine possible variations in stellar population with disk height. Figure 5 shows the completeness-corrected profiles as a function of disk height for each field. Each profile is derived using the same methodology as in §3, typically using $``$10,000 stars per field. The surface densities are then normalized to have $`\mathrm{\Sigma }𝑑z=1`$. All the fields show a similar pattern. The MS (blue) stars have the narrowest distribution while the AGB (red) and RGB (orange/yellow) stars have broader distributions and typically show a dip near the midplane. Because we have corrected for incompleteness, the dip almost certainly due to dust absorption, as we demonstrate in §4.4 with a very simple model.
Figure 5 suggests that older stellar populations become more prominent with increasing disk height. We quantify this trend in Figure 6, which shows the ratios of surface densities in our different age bins. The ratios were normalized to those expected for a constant SFR using the MATCH synthetic CMDs (see §4.1.1). A ratio of one in Fig. 6 therefore corresponds to a constant star formation rate and increasing values correspond to older stellar populations. We note that the ratio is only plotted where the signal-to-noise of the ratio is greater than 3. The small number of MS stars at large scale heights limits our ability to trace the $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{MS}}`$ and $`\mathrm{\Sigma }_{\mathrm{AGB}}`$/$`\mathrm{\Sigma }_{\mathrm{MS}}`$ as high above the midplane as the profiles shown in Figures 3, 5 and 7. Also, NGC 4631 and NGC 4631-DISK are not included in Figure 6 because the high completeness limit results in very few RGB stars (see Fig. 1) and an increased contamination of AGB stars in the RGB box.
The top and middle panels of Figure 6 show that in each of the fields, RGB stars become more numerous relative to MS and AGB stars with increasing disk height. However, this trend shows an enormous variation from galaxy to galaxy. In IC 5052 the ratio $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{MS}}`$ becomes as high as $``$100 times the midplane value, while in NGC 55 and NGC 4244 the increase is much more moderate, to $``$10 times the midplane value. This variation is most likely the result of a range of recent SFRs in our galaxies. The ratio $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{AGB}}`$ is much more consistent from galaxy-to-galaxy, however. This may result from the overlapping time range spanned by stars in the AGB and RGB boxes and/or the large time ranges these boxes span relative to the MS. We argue in §5 that the RGB population is likely to be dominated by truly old stars much older than the AGB population. Interestingly, the field showing the the flattest $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{AGB}}`$ profile was the NGC 55-DISK field located in the outer parts of the NGC 55 disk, perhaps suggesting a different star formation or dynamical history at large radii ($`5\times h_R`$). However, this galaxy is the least inclined in our sample and is somewhat irregular in shape, therefore results for this one system should not be overinterpreted.
The low values of $`\mathrm{\Sigma }_{\mathrm{AGB}}`$/$`\mathrm{\Sigma }_{\mathrm{MS}}`$ and the high values of $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{AGB}}`$ result from a lack of AGB stars compared to the constant SFR MATCH synthetic CMD. This would seem to suggest that the galaxies’ star formation histories (SFHs) are depressed at intermediate ages and enhanced at young ages. However, as we noted in §4.1, the AGB morphologies in the synthetic CMDs are not well matched to the observational CMDs, probably because of the difficulty in modeling the AGB phase of evolution (Marigo, 2001). This discrepancy combined with the differences seen between the two sets of synthetic CMDs suggests that a derivation of accurate SFHs using just the brightest stars in a galaxy is not yet possible.
To check if the scatter in the $`\mathrm{\Sigma }_{\mathrm{RGB}}`$/$`\mathrm{\Sigma }_{\mathrm{MS}}`$ ratio was in part due to varying radial coverages of the galaxies (Table 3), we remade the plots in Figure 6 using only stars within the central scale length of each galaxy. These plots were similar to those shown and showed comparable scatter. This suggests that the observed variations from galaxy to galaxy in the stellar populations ratio reflect global differences in the galaxies’ SFHs and/or vertical structure. For instance, if we assume the trend towards older populations with increasing scale height results from disk heating, then the scatter in Figure 6 suggests substantial variations between galaxies in either the mechanisms that heated the disk, or the SFH of the disk. Despite these variations, Figure 6 gives strong evidence that overall the age of the stellar populations increases with increasing scale height.
#### 4.2.1 Stellar Population Scale Heights
To further quantify the differences in the vertical distribution of the three CMD regions, we fit each surface density profile to a sech<sup>2</sup> function in which the normalization, central position, scale height ($`z_0`$), and background level were all allowed to vary. We fit each profile only at disk heights $`>`$3$`z_{1/2}`$ to avoid the dips near the midplane, except in IC 5052 and NGC 5023 where we used disk heights $`>`$1.5$`z_{1/2}`$ to allow fitting of the very narrow $`\mathrm{\Sigma }_{\mathrm{MS}}`$ profiles. Figure 7 shows the resulting fits to the $`\mathrm{\Sigma }_{\mathrm{MS}}`$, $`\mathrm{\Sigma }_{\mathrm{AGB}}`$, and $`\mathrm{\Sigma }_{\mathrm{RGB}}`$ profile of each galaxy, in the top, middle and bottom panels respectively. The observed profiles are shown as a solid line, while the best fitting sech<sup>2</sup> function is shown as a dashed line. The dotted lines at $`\pm `$3$`z_{1/2}`$ ($`\pm `$1.5$`z_{1/2}`$ in IC 5052 and NGC 5023) delineate the region excluded from the fit. The error bars on the data points are used to weight the sech<sup>2</sup> fits and reflect Poisson errors in the number counts, but do not include uncertainties in the completeness corrections. The scale height of the best-fitting sech<sup>2</sup> function is shown in the upper-left corner of each panel, and the error shown is scaled by the square root of the reduced $`\chi ^2`$ of the fit. The reduced $`\chi ^2`$ values for a majority of the fits were between 0.8 and 1.3, but were larger for NGC 55 and NGC 4631 due to their irregular structure. Scale heights and errors for all the fits are shown in Table 4.
We find that in each galaxy, the MS scale height value is the narrowest followed by the AGB and then the RGB. In all cases the RGB population is significantly broader than the AGB, MS and $`K_s`$ band $`z_0`$ values. This result strongly suggests the presence of an older component with larger scale height. An analysis of the variations in scale height of the MS, AGB and RGB populations with scale length in each of the galaxies turned up no obvious trends. We also identified no trends with galaxy rotation speed, due to the small range of masses spanned by our sample galaxies. We note that in some cases, the $`z_0`$ derived for all stars (§3) is somewhat larger than the $`z_0`$ derived for just the RGB stars. This results from the lower scale heights used in the fits to the different stellar populations - if the lower limit for the fit to all stars is reduced from 5$`z_{1/2}`$ to 3$`z_{1/2}`$, the derived $`z_0`$ is less than or equal to the RGB star $`z_0`$ in each galaxy, as expected if the total stellar density is well characterized by a combination of the MS, AGB and RGB components. The fits presented here for these different stellar population components are in general not extremely sensitive to the range of $`z`$ values used. Varying the lower limit of the fit between 1.5-5$`z_{1/2}`$ typically changed the AGB and RGB $`z_0`$ values by less than 10%. The compact MS components were more dramatically affected because of the smaller number of stars at large scale heights.
From this analysis it appears that NGC 4631 has the “thickest” old component with an RGB scale height of $``$1250 pc, roughly 2.5 times larger than the MS and $`K_s`$ band fits. We note that the fits for NGC 4631 were truncated at large negative disk heights to prevent the contamination of stars from companion galaxy NGC 4627. IC 5052 and NGC 4144 have similar ratios ($``$2.5) of RGB to MS scale heights, while NGC 4244 has the smallest ratio, with an RGB scale height only 1.7 times that of the MS stars.
There is also evidence for a modest flaring of the stellar components between the central and ’-DISK’ pointings of NGC 55 and NGC 4631. The ’-DISK’ fields are centered 4.8 and 6.1 scale lengths (see Table 2) from the center of the NGC 55 and NGC 4631 respectively. An increase in $`z_0`$ values by a factor of 1.1 to 1.6 is seen for all three components in NGC 55 and for the AGB and RGB component in NGC 4631.
The profiles deviate from the fitted sech<sup>2</sup> profile significantly near the midplane. This deviation is almost certainly due to dust, we model this effect in §4.4. At larger disk heights ($`>`$2-3 kpc) there is also a slight overdensity in the RGB components of IC 5052, NGC 4144, NGC 4244 and NGC 5023. These overdensities hint at the possible presence of an even more broadly distributed old component. In three of these galaxies, the RGB fits had elevated $`\chi ^2`$ values relative to the MS and AGB fits. We estimate that these overdensities occur at a surface brightness $`\mu _{\mathrm{F606W}}`$ 28 mag arcsec<sup>-2</sup> assuming a luminosity function similar to galactic globular clusters (Buonanno et al., 1994; Kravtsov et al., 1997). However, without better knowledge of the background level, it is not possible to verify the existence of this component.
The scale heights measured in the $`K_s`$ band (Table 4) are closest to those measured for the MS and AGB components. Half of the galaxies have $`K_s`$ band scale heights closer to the MS value and the other half closer to the AGB value. This suggests that the $`K_s`$ band light in these galaxies is dominated by relatively young stellar populations, probably red supergiants and AGB stars (in agreement with the findings of Aoki et al., 1991). This result runs contrary to the common assumption that the NIR light primarily traces older stellar populations (e.g. Florido et al., 2001), and is significant in that NIR luminosity is often used as a proxy for stellar mass when comparing galaxies of different types and masses. However, we note that our $`K_s`$ band scale heights are biased towards higher surface brightness populations due to the bright limiting isophote of the 2MASS data from which they are derived.
#### 4.2.2 Comparison to Previous Observations
The results above indicate that there is a systematic increase in the vertical scale heights of older stellar populations in our sample of low mass, late-type disks. Before investigating possible origins for these structural differences in §4.3 & §5, we now compare our measurements of scale heights to previous observations of the vertical structure of disks.
The most detailed constraints on the scale heights of different stellar populations come from the solar circle of the Milky Way. Studies have revealed a complicated disk structure, with a young and old thin disk embedded within a more extended thick disk. The young thin disk is the narrowest of the three, having a scale height of z$`{}_{0}{}^{}200`$ pc, as traced by stars with bright absolute magnitudes ($`M_V3`$) (Schmidt, 1963). In contrast, the scale heights of the young main sequence stars in our sample are almost all significantly larger than the Milky Way value suggesting that the low mass galaxies in our sample form stars in a thicker layer than the Milky Way, consistent with Dalcanton et al. (2004). The resulting axial ratios for our samples’ young star forming disks are also much thicker as well, with $`z_0/h_r=1.86.3`$ (see Table 4) for our sample galaxies, versus $`z_0/h_r15`$ for the young thin disk of the Milky Way.
The division of the Milky Way’s older stellar populations into a thin and thick disk was first introduced by Gilmore & Reid (1983) to explain a break in the number counts of F & G stars at $`1`$ kpc. While the need for two old disk components was long debated, recent measurements of systematic $`\alpha `$-element enhancement in thick disk stars (most recently Gratton et al., 2003; Feltzing et al., 2003; Mishenina et al., 2004; Bensby et al., 2005) strongly suggest that the thick disk is indeed distinct from the old thin disk. Recent observations (Chen et al., 2001; Siegel et al., 2002) give a scale height for the old thin disk of z$`{}_{0}{}^{}600`$ pc, similar to found in the original Gilmore & Reid (1983) study. These same studies suggest that the exponential scale height of the thick disk is $`h_z700`$ pc (corresponding to z$`{}_{0}{}^{}1400`$ pc, thinner than originally claimed). The Milky Way thick disk is therefore roughly twice the height of the old thin disk, and 7 times the height of the young thin disk.
Within our own sample, the scale height of the old RGB component is mostly intermediate between the Milky Way old thin disk and the thick disk. Our sample galaxies have much lower masses and surface densities than the Milky Way, and, lacking any firm model that predicts how the properties of the old thin disk and the thick disk should vary with galaxy mass, we are hesistant to attribute the extraplanar population to either an old thin or a thick disk on the basis of the the surface brightness profiles alone. There are no dramatic inflection points in the RGB surface density profiles plotted in Figure 7 that would assist in a unique separation of old thin disk and thick disk stars, and the possible overdensity of stars above 2-3 kpc may well be due to a stellar halo. Even if the RGB component is similar to the Milky Way thick disk, this lack of inflection is not unexpected. In the Milky Way, the inflection point in the surface density of F & G dwarves that marks the separation of the thin and thick disks is likely the result of two different populations of stars separated in age. The lack of similar inflection points in our RGB profiles can easily be explained if the RGB stars don’t have as wide a range of ages as the Milky Way dwarves. We will show in §5 that the RGB stars in our galaxies may very well be dominated by a single-age population.
We do note, however, that the axial ratios of the RGB disks range between $`h_r/z_0=1.03.3`$, with a median of 1.8 (adopting the $`K_s`$-band radial scale length, and averaging the two independent measurements for NGC 55 and NGC 4631). For comparison, the axial ratios of the old thin and the thick disks of the Milky Way are 5.0 and 2.1, respectively (assuming $`h_r=3`$ kpc for both components). Thus, in terms of their *overall* structure, the RGB component we detect is significantly more round than the Milky Way’s old thin disk, and is distributed more like the Milky Way thick disk. However, without additional information we cannot ascribe a common formation scenario to our observed RGB component and the Milky Way thick disk. We revisit this issue in the discussion (§6), after analyzing the disk heating and the vertical metallicity gradients of our sample galaxies.
Outside of the Milky Way, the most detailed information comes from studies of the vertical distribution of resolved stars in HST images, similar to the work we present in this paper. Tikhonov et al. (2005) and Tikhonov & Galazutdinova (2005) present evidence for extended components in six galaxies, which they qualitatively argue correspond to thick disks and halos. Of the galaxies that overlap our sample (NGC 55, NGC 4144, and NGC 4244), they include archival WFPC2 observations to reach greater disk heights in NGC 4244 and NGC 55 than spanned by our ACS images. In both cases they assume *a priori* that the RGB stars at lower disk heights trace a thick disk. For NGC 4244, Tikhonov & Galazutdinova (2005) show an exponential distribution of RGB stars between $``$1 and $`3`$ kpc (their Figure 8) that appears to roughly match the scale height of the profile shown in Fig. 7. Beyond $`3`$ kpc, they see a flattening in the number counts which they claim is a halo, but which may also be the background level. For NGC 55, Tikhonov et al. (2005) plot an exponential distribution of RGB stars between 2 and $`7`$ kpc – i.e. at much greater disk heights than probed by our data. However, based on inspection of their Figure 12, the extended RGB component has a z<sub>0</sub> value of $`2`$ kpc, which is 2-3 times the width of the RGB component we fit. Although they assume this component is due to a thick disk based on its exponential surface density distribution, the axial ratio of this component would in fact be prolate ($`h_r/z_00.5`$) and thus may be more analogous to the Milky Way’s stellar halo<sup>4</sup><sup>4</sup>4Note, however, that NGC 55 is the least inclined galaxy in our sample, complicating the interpretation of its projected structural parameters.. The change in slope also implies a break in the RGB distribution in NGC 55 at around $`2`$ kpc. By analogy, this may indicate that the marginal overdensities we are seeing at comparably large disk heights in our RGB profiles might be the signature of an additional broader halo component. Mould (2005) also finds the presence of old stars at large scale heights, and while these stars are automatically assumed to be a thick disk component, no detailed analysis of their spatial distribution is presented.
In addition to these recent studies of resolved stars, most previous studies of the vertical structure of disks have focused on detecting thick disks and stellar halos using unresolved surface brightness profiles of the galaxies (e.g. Pohlen et al., 2004; Fry et al., 1999; Dalcanton & Bernstein, 2002; Neeser et al., 2002). Because we only detect stars at bright magnitudes, it is difficult to accurately convert our measured surface density of stars (Figure 2) to a surface brightness. However, assuming the outer parts of our galaxies have luminosity functions similar to Galactic globular clusters (Buonanno et al., 1994; Kravtsov et al., 1997), we estimate that we reach F606W surface brightnesses of $``$28 mag arcsec<sup>-2</sup>. This is comparable to the depth reached in deep ground-based observations.
Only one of our galaxies has been analyzed for vertically extended components using ground based data. Fry et al. (1999) present $`R`$-band surface photometry of NGC 4244 and find no evidence of a second thick disk component, based on the lack of an inflection point in the surface brightness distribution above the plane. They trace the vertical profile of NGC 4244 along the minor axis to $`2`$ kpc at which point it falls below their surface brightness limit of 27.5 mag arcsec<sup>-2</sup>. We trace the RGB component out to nearly $`3`$ kpc, and find a scale height that is similar to their fitted $`R`$-band scale height (assuming $`h_z`$=$`\frac{1}{2}`$z<sub>0</sub>). Their lack of an inflection in the surface brightness profile is consistent with the Tikhonov & Galazutdinova (2005) analysis to larger scale heights. This suggests that their fit was dominated by the old stars, and not a younger population. However, the lack of an inflection in the surface density distribution does not unambiguously rule out the presence of multiple components.
The ubiquity of thick disks in galaxies has previously been proposed by Dalcanton & Bernstein (2002) based on color-gradients in edge-on disk galaxies. Our observations confirm that the color-gradients (at least at the low-mass end) are the result of true differences in stellar populations. However, whether these gradients have an analagous formation mechanism to the Milky Way thick disk is not clear.
One set of observations that reaches considerably deeper than these ground-based observations is presented by Zibetti et al. (2004), who used stacked Sloan images to show that halos are common in late-type, edge-on galaxies. Their composite galaxy has a significantly wider field of view and poorer resolution than our observations. They show that the best-fitting model to their data is a disk+halo model, with the disk component dominating out to roughly 10 exponential scale heights ($``$10z<sub>1/2</sub>). Their limited resolution and combination of a heterogenous sample of galaxies would likely prevent them from seeing the RGB components we see in our galaxies. However, the possible detection of the more extended RGB components detected in IC 5052, NGC 4144, and NGC 5023 may be halos similar to the Zibetti et al. (2004) halo.
### 4.3. Disk Heating
The increase in $`z_0`$ seen in each galaxy between the MS, AGB and RGB populations (Fig. 7, Table 4) could result from a number of mechanisms, including vertical heating of a thin disk. Such a model would naturally produce the observed trend of older stellar populations having larger scale heights. In Figure 8 we plot the increase in scale height with mean stellar age for four fields that span the observed behaviour in our sample. We plot scale heights for the RGB and AGB, normalized by the MS scale heights, with the height of the symbols indicating the 1$`\sigma `$ uncertainties on $`z_0`$. Note that when interpreting our data in the context of disk heating models we are therefore implicitly assuming that the RGB and AGB stars were originally formed in a layer with a scale height comparable to that of the present day main sequence stars. We assign characteristic ages to the RGB and AGB using the MATCH synthetic CMD tests for a constant SFR (§4.1.1). However, because the galaxies’ actual star formation histories may differ significantly from the constant SFR assumed in Figure 4, we cannot assign a single age to each stellar population. Instead, we use the resulting age distributions to identify the 25th, 50th (Median) and 75th percentile ages. The resulting age ranges are shown by the width of the individual boxes in Figure 8. However, note that the actual age of the population may lie entirely outside of the boxes, for example, if the RGB stars were all formed in a single burst 12 Gyr ago. Thus, when interpeting Figure 8, one has substantial allowance in assigning an age.
Overplotted on Figure 8 are dashed lines showing a range of power-law increases in the disk scale height $`z_0`$ with time ($`z_0t^\beta `$). For an isothermal sech<sup>2</sup> profile, the $`z_0`$ values are related to the vertical velocity dispersion ($`\sigma _z`$):
$$z_0=\frac{\sigma _z^2}{2\pi G\mathrm{\Sigma }}$$
(2)
where $`\mathrm{\Sigma }`$ is the surface density of the disk (Eq. 17 in van der Kruit, 1988). Studies of disk heating traditionally use power laws in the velocity dispersion, $`\sigma _zt^\alpha `$, and thus $`\alpha =\beta /2`$. Figure 8 therefore demonstrates that the vertical velocity dispersion of our galaxies has increased no faster than $`\alpha =0.15`$. More specifically, there are no characteristic ages that can be assigned to the AGB and RGB stars that yield heating rates greater than $`\alpha =0.15`$ (with the possible exception of the more massive, interacting galaxy NGC 4631), and thus this conclusion is robust even in light of our substantial age uncertainties.
In contrast, the disk heating that has been observed in the Milky Way is comparatively rapid. The age-velocity dispersion relation (AVR) for Milky Way disk stars suggests that the vertical velocity dispersion increases with time with values of $`\alpha `$ ranging between 0.3 and 0.6 (e.g. Wielen, 1977; Binney et al., 2000; Nordström et al., 2004, see summary in Table 1 of Hanninen & Flynn 2002). In contrast, our limit of $`\alpha 0.15`$ is significantly smaller than the Milky Way value. These data immediately suggest that any disk heating in our low mass galaxies has been far less effective than in the Milky Way. Moreover, if some fraction of the extraplanar RGB stars are not due to disk heating, and are instead due to direct accretion or *in situ* formation at large scale heights, or if the RGB stars are weighted towards old ages (as we argue below in §5), then the actual rate of disk heating is even lower than suggested by Figure 8.
There are several reasons why disk heating is expected to be low for our sample galaxies. Within the Milky Way, the increase in vertical velocity dispersion with time is thought to be due to scattering by spiral arms (Barbanis & Woltjer, 1967; Sellwood & Carlberg, 1984; Carlberg & Sellwood, 1985), by molecular clouds (Spitzer & Schwarzschild, 1951), or both (Carlberg, 1987; Jenkins & Binney, 1990; Jenkins, 1992; Shapiro et al., 2003, see also the review by Lacey 1991). However, our galaxies have sufficiently low masses and surface densities that they are unlikely to be globally gravitationally unstable (Dalcanton et al., 2004; Verde et al., 2002) and thus would not host strong spiral arms. Given that scattering by spiral arms seems to be the dominant heating mechanism in the Milky Way (e.g. most recently De Simone et al., 2004), the absence of spiral arms alone should cause a drastic drop in heating rate down to $`\alpha 0.20.25`$, the expected value for heating by giant molecular clouds alone (e.g. Hänninen & Flynn, 2002). Likewise, the absence of strong dust lanes in these systems and the results of §4.4 both indicate that the cold molecular ISM is in a thicker layer than in the Milky Way. This large scale height for the cold ISM, and the general lack of molecular gas in low mass galaxies (Young & Scoville, 1991; Leroy et al., 2005) should therefore further suppress the efficiency of disk heating. Finally, the young stellar disks in our sample are apparently much thicker than in the Milky Way, which could reduce the efficiency of any heating mechanism (Freeman, 1991). Shapiro et al. (2003) also argue for reduced disk heating in late-type galaxies based on the ratio of vertical to radial velocity dispersions. However, their rationale for the observed trend is opposite from what we conclude from our data.
As an aside, the low observed heating rates may provide strong constraints on cosmologically important sources of disk heating including late-time satellite accretion (Quinn et al., 1993), massive black holes (Lacey & Ostriker, 1985), or halo substructure (e.g. Hänninen & Flynn, 2002; Benson et al., 2004). However, the expected heating rates for such models have been calibrated for massive spiral disks, not the thicker, lower surface density galaxies studied here.
### 4.4. Modelling Dust Effects on the Stellar Density Profiles
Before continuing to explore the origin of extraplanar stars, we briefly examine the vertical distribution of the dust layer. At first glance, interpretation of the stellar density profiles near the midplane in Figure 5 might be somewhat confusing. If the dips in surface density are due to dust, why does the dust appear to affect the RGB and AGB stars more than the MS stars? We suggest this may occur because the dust layer is opaque near the midplane and is distributed with a scale height greater than or equal to the MS population, but less than the AGB/RGB populations. The MS stars we are seeing would then lie entirely in front of an obscuring dust screen, while the AGB/RGB populations would have a significant population at large disk heights above where the galaxy becomes optically thin. The dip in their numbers near the midplane is then explained because the optically thick dust layer obscures some fraction of the stars along the line of sight.
To test this explanation, we built a simple ’toy model’ galaxy with MS, RGB and AGB populations distributed as sech<sup>2</sup> profiles with the $`z_0`$ as shown in Table 4. All components were given identical radial distributions with the $`K_s`$ band exponential scale length. The dust component was assumed to also follow a sech<sup>2</sup> profile with a variable scale height, $`z_{0,\mathrm{dust}}`$, and a radial distribution identical to that of the stars. For simplicity we assumed that the dust has no effect at an optical depth less than one, but is completely obscuring at greater optical depths. Thus, along a line of sight, the dust is completely transparent to $`\tau =1`$, and completely opaque beyond. For each vertical position we integrated the dust component along the line-of-sight until an optical depth of one was reached, which set the depth of the dust screen at that height. Stellar density profiles for the three separate populations were then created by totaling the number of stars in front of the dust screen at each height. We then normalized the stellar density profiles as in Figure 5.
Figure 9 shows the resulting model of NGC 4144 for three values of $`z_{0,\mathrm{dust}}`$ presented for comparison to the observations shown in Figure 5. In each case the amount of dust in the midplane is the same. The underlying values for $`z_0`$ were adopted from Table 4 (374, 699, and 934 pc for the MS, AGB and RGB respectively). The left panel shows the results for a dust layer whose scale height is narrower than all three stellar components. For this case, there is a pronounced dip in the surface density profile of all three components. In the middle panel, the value of $`z_{0,\mathrm{dust}}`$ is between the $`z_0`$ values for the MS and the AGB/RGB populations. In this case there is a dip only in the AGB/RGB, because the height of the MS layer is entirely confined within the opaque dust layer, allowing only the unobscured stars on the near side of the galaxy to be detected. The right panel has a dust layer larger than the both the MS and AGB value and therefore a dip is seen only in the RGB component. The middle panel does a good job of qualitatively matching the observations for NGC 4144 in Figure 5.
Referring back to Figure 5, we can see that the main sequence profile lacks a dip near the midplane for most of the fields, while the RGB has a dip in all cases within 2-3 $`z_{1/2}`$. This suggests that the dust in the galaxy is opaque below 2-3 $`z_{1/2}`$ and that it is distributed in a layer with thickness greater than or equal to the MS stars. Because the MS stars are already distributed in a thicker distribution than in the Milky Way, this result supports the Dalcanton et al. (2004) finding that galaxies with circular velocities below 120 km s<sup>-1</sup> have large dust scale heights and do not form thin dust lanes. All of the galaxies presented here except NGC 4631 are below this circular velocity limit (Table 1).
This model also suggests that although the depth to which we see in each galaxy is different at differing scale heights, it is the same for all the stellar populations at a single scale height. This validates the comparisons made in Fig. 6 between the ages of the stellar populations at different scale heights.
Although this model matches the gross characteristics of many of the profiles shown in Fig. 5, it fails to fully explain their details. Most notably, in NGC 55 and NGC 4244, the MS profiles are significantly lower than the best-fitting sech<sup>2</sup> function, contradicting the idea that all the stars we are seeing lie in front of a screen. This is likely the result of our model’s lack of sophistication. Physically, it doesn’t take into account the possibility of flares or changing scale lengths as a function of stellar population. In addition, it treats dust extinction in a very simplistic fashion. A more sophisticated treatment of the dust (such as the one presented by Matthews & Wood, 2001) is beyond the scope of this paper. The conclusions reached in this section should be considered tentative and will be tested in a later paper, in which we will present a more realistic dust model.
## 5. Metallicity Distribution Functions
As has been shown previously (e.g. Da Costa & Armandroff, 1990; Armandroff et al., 1993; Frayn & Gilmore, 2002), the color of the red giant branch near its tip can be used to constrain the metallicity of old stellar populations. Although reddening due to dust will prevent an accurate measurement of the metallicity of stars near the midplane, we can determine a rough metallicity for stars above the midplane where the effect of dust and the contamination from AGB stars are small (as shown in Figure 6, middle panel).
Figure 10 shows a composite CMD of all high-latitude stars (above 4$`z_{1/2}`$) in IC 5052, NGC 55, NGC 55-DISK, NGC 4144, NGC 4244, and NGC 5023. This disk height limit was chosen (1) to be well above any of the dips associated with dust in Figure 5, (2) to dramatically reduce contamination by AGB and HeB stars that might interfere in the metallicity determination, and (3) to be where thick disk stars dominate in the Milky Way (Chen et al., 2001). Overlayed on the CMD are 10 Gyr old RGB isochrones (Girardi et al., 2000) at metallicities ranging from \[Fe/H\] $`=2.3`$ to 0.0, with higher metallicities being redder. Examination of Figure 10 shows that the peak of the distribution falls between the \[Fe/H\] $`=1.3`$ and -0.7 lines. Roughly 13% of the stars fall bluewards of the \[Fe/H\] $`=2.3`$ isochrone, probably due to a combination of photometric error and the presence of AGB stars. Very few are redder than the \[Fe/H\] $`=0.4`$ isochrone. Overall, Figure 10 indicates that most of the stars above 4$`z_{1/2}`$ are moderately metal-poor.
Improving on the metallicity determinations in Paper I, here we will derive metallicity distribution functions (MDFs) using untransformed magnitudes and Padova isochrones (Girardi et al., 2000). To determine metallicity distribution functions for individual galaxies we binned the stars in up to three independent 0.2 magnitude wide bins centered on M<sub>F814W</sub> of -3.5, -3.3, and -3.1. These are shown with dashed lines in Figure 10. The brightest bin was chosen to include stars as metal-rich as \[Fe/H\] $`=0.4`$, which do not get brighter than M<sub>F814W</sub> of -3.7. We considered only the magnitude bins that were above the 20% completeness cutoff for crowded regions (thus excluding the fainter bins in NGC 4144, see Figure 1). This cut eliminated NGC 4631 and NGC 4631-DISK from the analysis, because none of the bins fell above their 20% completeness limits. We corrected the colors and magnitudes of each star for foreground reddening, but made no correction for completeness, since the change in completeness across the color-range in question is similar to the error that would be introduced by that correction. We then determined the metallicity of each star by linearly interpolating between the 10 Gyr isochrones in each bin. Since some of the stars bluewards of the \[Fe/H\] $`=2.3`$ isochrone may be RGB stars scattered to bluer colors by photometric error, we attempted to include these stars in the MDFs. In each bin, stars bluer than the \[Fe/H\] $`=2.3`$ isochrone were given positive color shifts by multiplying a Gaussian random number by their error. This correction moved $``$50% of these stars redwards of the \[Fe/H\] $`=2.3`$ isochrone, and resulted in a slight increase ($``$0.01 in the normalized units) in the MDF between a \[Fe/H\] of -1.3 and -2.3. Table 5 gives the resulting number of stars used in MDF determination, the median F606W-F814W error, and the peak and mean metallicities of the MDFs.
Figure 11 shows the resulting MDFs for each field. The shaded regions show the total range in the MDFs as derived in the different magnitude bins, and indicate that our results are consistent among the different magnitude ranges. In general, the MDFs peak at metallicities \[Fe/H\] of -0.7 to -1.1 as expected from Figure 10, and have a tail of stars to low metallicities.
Before analyzing the MDFs further, we note that there are several uncertainties in the detailed shapes of the metallicity distribution functions. First, they are based on isochrone models and not empirical data (as in Sarajedini & Van Duyne, 2001). Second, the color errors on the metal-poor end (Table 5) translate into a large error in metallicity. We estimated the errors as a function of metallicity by inserting stars of a specific metallicity/color, giving them appropriate color errors, and then determining the spread in the resulting metallicity distribution. This procedure gave errors of 0.5-0.8 dex at \[Fe/H\] $`=2.3`$. However, the shape of the distributions are much more believable on the metal-rich end where the isochrones are well separated, giving errors of less than 0.1 dex at \[Fe/H\] $`=0.4`$. At the peak of the MDF (\[Fe/H\] $``$ -0.9), typical errors are 0.2 dex, suggesting that the peak metallicities derived here are fairly reliable.
Comparing the peak metallicity of the extraplanar stars to known Milky Way populations, Figure 11 indicates that the metallicities of the extraplanar stars are a factor of ten times too high to be analogs of the Milky Way’s stellar halo ($`[\mathrm{Fe}/\mathrm{H}]1.7`$, Wyse & Gilmore, 1995). This result would not change even if all the stars bluewards of the \[Fe/H\] $`=2.3`$ isochrone are low metallicity RGB stars. We note however, that the low metallicity of the Milky Way halo may not be typical. The halo of Andromeda has been found to be much more metal-rich than in the Milky Way, with a peak \[Fe/H\] $``$ -0.6 (Holland et al., 1996; Brown et al., 2003), although there is difficulty ascribing these outer stars to a halo population per se, given M31’s complicated outer structure (Ferguson et al., 2002).
Of all the Milky Way components, we find that the peak metallicities are most consistent with those of the metallicity of the Milky Way thick disk, which has \[Fe/H\] $``$ -0.8 based on F/G dwarfs (Wyse & Gilmore, 1995). The extraplanar stars studied here are somewhat more metal poor than the Milky Way’s thick disk (by up to 0.3 dex). However, this offset may not be surprising given the lower mass of our galaxy sample ($`V_c80`$ km s<sup>-1</sup> vs $`V_c220`$ km s<sup>-1</sup> for the Milky Way).
As in the Milky Way (Wyse & Gilmore, 1995; Haywood, 2001), the metallicities of the extraplanar stars appear to be more metal poor than the thin, young, main sequence population. Although dust prevents us from reliably measuring the metallicity of stars near the midplane, we can estimate their metallicity using the current gas phase metallicity. NGC 55, the only galaxy in our sample with a gas phase abundance measurement, has $`12+\mathrm{log}(\mathrm{O}/\mathrm{H})=8.32`$ at one disk scale length (Garnett, 2002). This metallicity corresponds to \[Fe/H\] $`0.6`$ (assuming \[Fe/O\] $`=0`$), which is 0.5 dex more metal rich than the extraplanar stars at a comparable radius. Other late-type disks with similar rotational velocities from the Garnett (2002) compilation have comparable gas-phase metallicities, suggesting that the offset in metallicity between the midplane and the extraplanar populations is likely to be systematic.
Although they do not explicitly examine stars as a function of scale height, studies of metallicities in other galaxies using methods similar to ours also find broad agreement with the presence of an extended \[Fe/H\] $`1`$ population of RGB stars. The LMC (which has a mass similar to the galaxies in our sample) has a peak metallicity distribution of \[Fe/H\] $`0.6`$ for RGB stars in the disk (Cole et al., 2000). Recent papers on the outer regions of M33 also find peak \[Fe/H\] values of -1.0 (Davidge, 2003; Tiede et al., 2004). Furthermore, a recent paper by Davidge (2005), derives a \[Fe/H\] of roughly -1 for NGC 55 using near-IR photometry of resolved extraplanar stars, closely matching our peak metallicity in Figure 11. Our data and others therefore suggest the pervasive presence of a significant \[Fe/H\] $`1`$ old population in late-type galaxies. If our association of this population with a thicker disk is generally true in other galaxies, then it presents an attractive solution to the “G-dwarf” problem seen in the Milky Way by providing the necessary prompt initial enrichment for stars in the thin disk (Truran & Cameron, 1971).
Overplotted on Figure 11 as dashed lines are the expected metallicity distributions for closed-box “simple” chemical evolution models (Eq. 20 of Pagel, 1997) scaled to the peaks of the MDFs. While the basic shape of these models are similar to our MDFs, there appears to be a deficit in some galaxies of stars at both low and high metallicities. A deficit at high metallicities is expected if star formation truncates before all the gas is consumed. Within the context of thick disk formation models, this truncation may occur if some of the gas reservoir that forms the extraplanar stars instead settles into the thin disk, if the extraplanar stars were heated from a previously thin but gas-rich disk, or if the extraplanar stars were directly accreted from merging satellites that suffered from tidal stripping or supernova blowout. The apparent deficit of stars at low metallicities may be another manifestation of the widespread G-dwarf problem. Thus, while the existence of a substantial population of stars at \[Fe/H\] $`1`$ may help to solve the G-dwarf problem in the thin disk, it may have simply pushed the problem into a new component. The solution to the extraplanar G-dwarf problem will likely lie among the suite of popular models previously explored for the thin disk (see Pagel, 1997). However, some of the deficit of stars at low metallicities may also result from the photometric errors and methods used to construct the MDFs, as discussed above.
Finally, we note that the peaks of the metallicity distribution functions given in Table 5 are significantly more metal-rich than our previous determination presented in Paper I. Rather than using native F606W$``$F814W colors, these earlier measurements applied the metallicity-color relation of Lee et al. (1993) to the mean color of the giant branch transformed to the Johnson-Cousins filter system. These previous values, reproduced in Table 5, range from \[Fe/H\] of -1.2 to -1.7, versus -0.7 to -1.1 in the present work. We believe that in addition to the magnitude transformation, the difference in derived metallicity results in part from the difference in binning (in \[Fe/H\] vs. F606W-F814W), such that the mean color does not correspond to the peak metallicity. This offset can be seen in the 0.1-0.2 dex offset between the mean and peak metallicity in Table 5. Furthermore, the Paper I determination also included stars at somewhat lower disk heights, thereby increasing the number of AGB contaminants. We believe that the MDFs and their peak metallicity that we present here are more reliable than the estimate given in Paper I.
### 5.1. Vertical Metallicity & Color Gradients
Models for the origin of extraplanar stars (i.e. disk heating, direct accretion, etc.) predict different degrees of variation in the stellar metallicity with height above the plane. To investigate the vertical variation of metallicity with disk height $`z`$ we have examined the median color of the RGB stars as a function of the height above the midplane. Figure 12 shows the median color of RGB stars between M<sub>F814W</sub> of -3.2 and -3.6 and redwards of F606W-F814W=0.7, binned by the scale height of the galaxies. Data are plotted where errors in the median color are $`<`$0.05 magnitudes. The hatched region at low disk heights shows where the effects of internal reddening may impact the colors of the stars. For the three fields with profiles extending beyond 2-3 kpc we bin the RGB stars at large disk heights in a single bin to reach adequate signal-to-noise in our measurement of the median color, plotted as diamonds in Figure 12. We note that these are the same stars that comprise the possible extended components discussed in §4.2.
Figure 12 demonstrates three main points. First, the color gradients in the galaxies are relatively small, indicating that the stars have nearly uniform metallicity with increasing distance above the plane, particularly at scale heights above the region potentially affected by dust. However, we note that the stars at very large radii (shown with diamonds) do tend towards bluer colors, possibly indicating the presence of a more metal-poor population at $`z10z_{1/2}`$ (2-3 kpc). Second, the color gradients show no systematic trends, and are equally likely to be rising or falling. Finally, all the galaxies have very similar RGB colors (as demonstrated already in Figure 11).
Our metallicity gradient results are consistent with previous observations of these and other low-mass spiral galaxies (Davidge, 2005; Tikhonov et al., 2005; Mould, 2005). Mould (2005) used HST archival data to study the vertical properties of disks in four low-luminosity ($`M_V16`$) edge-on galaxies comparable to those studied here. Using AGB, RGB, and red supergiant stars over a large range of magnitudes ($`8.5M_I1.5`$), he calculates the mean colors up to 2 kpc from the plane. His main results are that there are slight or no color gradients as a function of disk height and that the metallicities of the stars at disk heights between 400 and 2000 pc are between -0.8 and -1.0 in all four galaxies, in excellent agreement with what we find here. Tikhonov et al. (2005) also state that any metallicity gradient in NGC 55 is very small, in agreement with our observed lack of color-gradients in Figure 12. Recent simulations of thick disks by Brook et al. (2005) also show a lack of any metallicity gradient with disk height.
The lack of strong metallicity gradients can be explained in a number of ways. Mould (2005) suggests the lack of metallicity gradients rules out dissipative and simple accretion models for thick disk formation, and favors a model in which thick disks form during interactions. We note that the lack of metallicity gradient may also be enhanced by an “age-bias” in the metallicity measurements. The mass of stars on the RGB changes from $``$2 M at an age of 1 Gyr to $`<`$1 M at 10 Gyr. Given the steepness of the IMF in this mass range, the RGB age distribution will be weighted towards older ages, and therefore to a more uniform metallicity. However, based on the constant star formation rate models in §4.1, it still appears that the RGB stars could span a wide range of ages, even in the presence of the expected age-bias. Therefore, the simplest interpretation of the lack of gradients in our data is that many of the RGB stars at disk heights above 4$`z_{1/2}`$ ($``$1 kpc) formed at a similar time and thus have comparable enrichment histories, eliminating any metallicity gradient. This scenario could be explained either by a sudden heating of the disk by interaction, or by accretion of gas-rich satellites which resulted in the formation of a thick component. N-body simulations have shown that early merging and accretion events can produce thick disks with old ages (Abadi et al., 2003; Brook et al., 2004, 2005).
## 6. Discussion & Conclusions
The work presented here has identified a number of main observational results:
* In low mass, late-type galaxies the thickness of a stellar population increases systematically with the age of the stars being studied. This behavior has been seen not just in all six of the galaxies studied here, but in all other HST studies of edge-on late-type galaxies, and in the Milky Way as well. The larger scale heights of older stellar populations is therefore likely to be a generic property of galaxy disks.
* All of the studied galaxies show a clear intermediate age ($`15`$ Gyr old) population whose scale height is intermediate between that of the young main sequence stars and the older red giants.
* The metallicity of the dominant old stellar population has \[Fe/H\] $`1`$, but shows little or no gradient between $`3z_{1/2}`$ and $`10z_{1/2}`$ above the plane. Above this height ($`>23`$ kpc), there are tentative indications of decreasing metallicities, which may be associated with slight overdensities in the RGB surface density at similar distances above the midplane.
* In the low mass galaxies studied here ($`V_c70130`$ km sec<sup>-1</sup>), the young stellar population is systematically thicker than in the MW, and has a vertical scale height comparable to the thickness of the dust layer. This suggests that the cold ISM has a larger scale height in low mass galaxies, consistent with the lack of dust lanes observed in such systems (Dalcanton et al., 2004).
* The young and intermediate-age stellar populations dominate the integrated *near-infrared* light of late-type low mass galaxies.
We now interpret these observational facts in the context of disk formation models. First, taken at face value, the old RGB component’s $``$3:1 axial ratio and \[Fe/H\] $`1`$ metallicity suggest a close correspondence with the Milky Way’s thick disk. However, each of our galaxies’ scale heights steadily increase from the young main sequence to the intermediate age AGB and the older RGB. The uniformity of this trend strongly suggests that our disks are not simply the superposition of two components (i.e. a thick and thin disk). Instead, the data require a more complex model incorporating some disk heating to explain the systematically larger scale height of the intermediate-age population. The necessary disk heating would also have affected any older population, and thus must make some contribution to the thicker population of RGB stars. The required amount of disk heating is much smaller than is seen in the Milky Way (§4.3), and could likely be provided by molecular clouds or minor mergers. The latter scenario is slightly favored by the large variations in the apparent change of disk scale height with time (Figures 6 & 8). Heating through satellite accretion and interaction could naturally produce these stochastic variations. However, numerical simulations of heating in diffuse low mass disks would be required to definitively constrain any of the above scenarios and to assess how significant a contribution heating may have made to the thicker RGB component.
While the above argument strongly suggests that disk heating must play some role in the production of the extraplanar stars, the lack of a metallicity gradient in RGB stars at moderate disk heights ($`5002000`$ pc, as shown here and in Mould 2005) suggests that steady disk heating cannot entirely explain the thickest component of old RGB stars. If the past star formation rate has been constant, then there is a significant overlap between the stellar ages of the AGB and RGB regions of the color-magnitude diagram (Figure 4). A significant fraction of the RGB stars should therefore have smaller scale heights, younger ages, and enriched metallicities, and would thus produce a steady increase in metallicity towards the midplane for all but the most contrived scenarios. The most attractive explanation for the lack of metallicity gradient is that instead of a constant star formation rate, the majority of RGB stars at all scale heights must have formed early and with a well-mixed metallicity distribution. Such a population would dwarf any subsequent population of enriched RGB stars with lower scale heights. While steady dynamical heating could push this ancient population to larger scale heights, it could not simultaneously account for recent dynamical observations of counter-rotating disks at these disk heights in comparable galaxies (Yoachim & Dalcanton, 2005). Taken together, these observations are better explained by scenarios involving the formation of a thick disk of stars in merger events (as in Abadi et al., 2003; Brook et al., 2004, 2005). Overall, our results require that some disk heating occurs at intermediate ages (to puff up the AGB component), but that events at earlier times (interactions or mergers) created a majority of the RGB stars over a short timescale.
Finally, we present tenuous evidence for an extended old component seen only at disk heights $`>`$ 2-3 kpc. At large scale heights we see marginal overdensities of stars in the RGB profiles of Fig. 7. There also seems to be a reduction in the metallicity of RGB stars at this height (Fig. 12). In one of our galaxies, Tikhonov et al. (2005) finds strong evidence for a very extended component of RGB stars extending from $``$2-7 kpc. While this component appears to have an exponential distribution, its z<sub>0</sub> value of $``$2 kpc (compared to a radial scale length of $``$1 kpc) strongly suggests it is not a disk. These extended components are detected at about the same height where the halo becomes prominent in the stacked Sloan images of Zibetti et al. (2004). However, based on our observations, we are unable to assess the properties or frequency of these components.
The present-day structure of galaxy disks results from a complex mixing of effects and a full explanation requires detailed knowledge of the the star formation history, merging events, and disk heating. Studies like the one presented here help to disentangle these effects and determine their relative importance as function of galaxy type and mass. This study also shows the promise that HST observations of resolved stars have for enabling the detailed analysis of low surface brightness stellar components in galaxies outside the Local Group. A comparison of our data with N-body simulations of low mass disk galaxies would assist in constraining disk heating and merging scenarios. Unfortunately, current simulations of disk galaxy formation have focused on massive galaxies like the Milky Way (Abadi et al., 2003; Brook et al., 2004, 2005). Also, deeper observations that fully resolve the red and blue horizontal branches would greatly improve constraints on the star formation histories of these galaxies and improve our understanding of their structure.
Acknowledgements: The authors would like to thank Andrew Dolphin and Antonio Aparicio for their help in generating synthetic CMDs, Leo Girardi for supplying us with isochrones, our anonymous referee for their thoughtful suggestions, and Peter Yoachim, Andrew West, and Kevin Covey for helpful discussions. This work was supported by HST-GO-09765, the Sloan foundation, and NSF Grant CAREER AST-0238683.
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