id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0506/astro-ph0506732.html
|
ar5iv
|
text
|
# Is it possible to consider Dark Energy and Dark Matter as a same and unique Dark Fluid?
## I Basis of the model and definitions
If one considers a Friedmann–Lemaître Universe with different fluids: photons, neutrinos, baryons and a dark fluid, the Friedmann equations take the form:
$`\left\{{\displaystyle \frac{\dot{a}}{a}}\right\}^2=H^2={\displaystyle \frac{8\pi G}{3}}\rho {\displaystyle \frac{k}{a^2}},{\displaystyle \frac{\ddot{a}}{a}}={\displaystyle \frac{4\pi G}{3}}\{\rho +3P\},`$ (1)
where $`P`$ and $`\rho `$ denote the total pressure and the total density in the Universe, respectively. It will be assumed in the whole article that the value of $`a`$ today is 1. For the dark fluid model, pressure and density can be expanded as:
$`P`$ $`=`$ $`P_r+P_D,`$
$`\rho `$ $`=`$ $`\rho _r+\rho _b+\rho _D,`$ (2)
where $`r`$ denotes the radiation (i.e. photons + neutrinos), $`b`$ the baryonic matter and $`D`$ the dark fluid.
In the standard model of cosmology, the equivalent expressions are:
$`P`$ $`=`$ $`P_r+P_\phi ,`$
$`\rho `$ $`=`$ $`\rho _r+\rho _b+\rho _{dm}+\rho _\phi ,`$ (3)
where $`dm`$ denotes the dark matter and $`\phi `$ the dark energy.
For each fluid, one can write the conservation of the energy–momentum tensor in a homogeneous and isotropic spacetime, which reads:
$$\frac{d}{dt}(\rho _{\text{fluid}}a^3)=P_{\text{fluid}}\frac{d(a^3)}{dt}.$$
(4)
This equation is equivalent to the second Friedmann equation. As usual, one can define a cosmological parameter corresponding to each fluid by:
$$\mathrm{\Omega }_{\text{fluid}}^0=\frac{\rho _{\text{fluid}}^0}{\rho _0^c},$$
(5)
where $`\rho _0^c`$ is the critical density today. A cosmological parameter containing the curvature term can also be written:
$$\mathrm{\Omega }_K\frac{k}{a^2\mathrm{\hspace{0.33em}3}H^2/(8\pi G)},$$
(6)
so that one gets finally a simple equation:
$$\mathrm{\Omega }_K+\mathrm{\Omega }_r+\mathrm{\Omega }_b+\mathrm{\Omega }_D=1.$$
(7)
For a flat Universe, one has $`\mathrm{\Omega }_K=0`$, and the Friedmann equations are remarkably simplified.
Finally, one can define the equation of state of the fluid as:
$$\omega _{\text{fluid}}\frac{P_{\text{fluid}}}{\rho _{\text{fluid}}},$$
(8)
and one knows that for baryonic matter $`\omega _b=0`$, for radiation $`\omega _r=1/3`$ and for a real cosmological constant $`\omega _\varphi =1`$. For simplicity reasons, it is considered in the following that this fluid is perfect, i.e. the entropy variations and the shear stress can be ignored. For other assumptions, it is necessary to specify a dark fluid model.
A basis for the dark fluid model is now defined, and a comparison with the observations of supernovæ of type Ia is given in the following.
## II Constraints from the Supernovae of Type Ia
Cosmological constraints from the supernovæ of type Ia are based on the joint observations of the redshift $`z`$ and of the apparent luminosity $`l`$ of a large number of supernovæ. Supernovæ of type Ia are often considered as standard candles, i.e. the absolute luminosity $`L`$ is approximately the same for every supernova (it is not completely true and recent studies correct the value of the absolute luminosity to reflect the deviations from the standard candle behavior Tonry et al. (2003)), so that it is possible to determine for each supernova the luminosity distance
$$d_L=\left(\frac{L}{4\pi l}\right)^{1/2}.$$
(9)
This luminosity distance depends on the reddening induced by the expansion of the Universe, and thus can reveal the presence of the cosmological components, through the equation:
$$d_L(z)=\frac{c}{H_0}\frac{1+z}{\sqrt{\left|\mathrm{\Omega }_K^0\right|}}S\left\{\sqrt{\left|\mathrm{\Omega }_K^0\right|}_0^z\frac{dz^{}}{\sqrt{F(z^{})}}\right\},$$
(10)
where $`z=a^11`$ denotes the redshift, $`S`$ is given by
$$S(x)\{\begin{array}{cc}\mathrm{sin}x\hfill & \text{if }k>0,\hfill \\ x\hfill & \text{if }k=0,\hfill \\ \mathrm{sinh}x\hfill & \text{if }k<0,\hfill \end{array}$$
(11)
and $`F`$ is defined as
$$F(z)=H_0^1(1+z)^2\frac{dz}{dt}.$$
(12)
One can see that $`F`$ is directly related to the first Friedmann equation through the term $`dz/dt=a^2da/dt`$. If there is only one component – replacing the two dark components – it shall have the same influence on the luminosity distance – and then on the expansion of the Universe – as dark matter and dark energy would have. Through the Friedmann equations, it seems clear that if
$`\rho _D`$ $`=`$ $`\rho _{dm}+\rho _\phi ,`$
$`P_D`$ $`=`$ $`P_{dm}+P_\phi =P_\phi ,`$ (13)
the dark fluid would provide the same effect on the expansion of the Universe as the two components.
Observations on the supernovæ of type Ia enable to give constraints on the dark component densities and on the dark energy behavior at low redshift Riess et al. (2004). From these constraints, it should be possible to characterize the dark fluid at low redshift. At first, the cosmological parameter corresponding to the dark fluid can be written in function of those related to dark matter and to dark energy:
$$\mathrm{\Omega }_D^0=\mathrm{\Omega }_{dm}^0+\mathrm{\Omega }_\phi ^0.$$
(14)
At low redshift, one can consider that the equation of state for the dark energy is, at first order in $`z`$:
$$\omega _\phi =\omega _\phi ^0+\omega _\phi ^1z,$$
(15)
and that the equation of state for the dark fluid can have the same form:
$$\omega _D=\omega _D^0+\omega _D^1z.$$
(16)
The observations have given constraints on the values of $`\omega _\phi ^0`$ and $`\omega _\phi ^1`$, and one would like to deduce from them constraints on $`\omega _D^0`$ and $`\omega _D^1`$. The equation of state of the dark fluid writes:
$$\omega _D=\frac{P_D}{\rho _D}=\frac{P_\phi }{\rho _{dm}+\rho _\phi }=\omega _\phi \frac{\rho _\phi }{\rho _{dm}+\rho _\phi }.$$
(17)
At first order in $`z`$, the dark matter density evolves like $`\rho _{dm}=\rho _{dm}^0a^3=\rho _{dm}^0(1+3z)`$ . It would be interesting to know the behavior of $`\rho _\phi `$. Let us assume that, at first order:
$$\rho _\phi =\rho _\phi ^0+\rho _\phi ^1z.$$
(18)
The equation of conservation of the energy-momentum tensor for each fluid satisfies:
$$\frac{d}{dt}(\rho _\phi a^3)=P_\phi \frac{d(a^3)}{dt}.$$
(19)
At first order in $`z`$, this equation becomes:
$$\frac{d}{dt}(\rho _\phi ^0+z(\rho _\phi ^13\rho _\phi ^0))=(\omega _\phi ^0\rho _\phi ^0)\frac{d(13z)}{dt},$$
(20)
so that the density of dark energy reads:
$$\rho _\phi ^1=3\rho _\phi ^0(1+\omega _\phi ^0).$$
(21)
Then, the relation between the ratio pressure/density for the dark fluid becomes:
$$\omega _D=(\omega _\phi ^0+\omega _\phi ^1z)\frac{\rho _\phi ^0(1+3(1+\omega _\phi ^0)z)}{\rho _{dm}^0(1+3z)+\rho _\phi ^0(1+3(1+\omega _\phi ^0)z)},$$
(22)
and one can determine the value of the two first terms of the expansion:
$`\omega _D^0`$ $`=`$ $`{\displaystyle \frac{\omega _\phi ^0\mathrm{\Omega }_\phi ^0}{\mathrm{\Omega }_{dm}^0+\mathrm{\Omega }_\phi ^0}},`$
$`\omega _D^1`$ $`=`$ $`{\displaystyle \frac{\omega _\phi ^1\mathrm{\Omega }_\phi ^0}{\mathrm{\Omega }_{dm}^0+\mathrm{\Omega }_\phi ^0}}+{\displaystyle \frac{3\mathrm{\Omega }_{dm}^0\mathrm{\Omega }_\phi ^0(\omega _\phi ^0)^2}{(\mathrm{\Omega }_{dm}^0+\mathrm{\Omega }_\phi ^0)^2}}.`$ (23)
The favored values for the cosmological parameters of the usual standard model from the supernovæ of type Ia Riess et al. (2004) combined with the results of other observations Tegmark et al. 2004 I are:
$`h=`$ $`0.70\pm 0.04`$
$`\mathrm{\Omega }_K^0=0.012\pm 0.022`$ $`\mathrm{\Omega }_{dm}^0=0.25\pm 0.04`$ (24)
$`\mathrm{\Omega }_b^0=0.049\pm 0.012`$ $`\mathrm{\Omega }_\phi ^0=0.712\pm 0.044`$
$`\omega _\phi ^0=1.02\pm 0.19`$ $`\omega _\phi ^1=0.6\pm 0.5.`$
From these values, one can calculate the parameters of the dark fluid:
$`\mathrm{\Omega }_D^0`$ $`=`$ $`0.962\pm 0.084`$
$`\omega _D^0`$ $`=`$ $`0.76\pm 0.25`$ (25)
$`\omega _D^1`$ $`=`$ $`1.0\pm 0.6.`$
These values are of course not completely representative of the dark fluid model, because they come from data analyses based on the usual standard model. Nevertheless, one can use them as test–parameters at low redshift.
Recent supernova observations tend to show that the dark energy has a negative pressure. Moreover, $`\omega _\phi <1`$ is not at all excluded, and in that case the dark energy cannot be explained anymore thanks to the usual models (see for example Caldwell et al. (2003) for a possible answer to this problem). From the precedent constraints, one can see that this difficulty vanishes with a dark fluid. Hence, the pressure of the fluid has to be negative today at cosmological scales, and seems to increase strongly with the redshift. As $`\omega _D^01`$, it seems possible to model the dark fluid with a scalar field.
The study of supernovæ provided us properties of the equation of state of our dark fluid at low redshift independently from the specification of a dark fluid model. We will now try to extract constraints from the information concerning large scale structures.
## III Large Scale Structures
We will not consider here a complete scenario of structure formation, which would require the specification of a precise model of dark fluid. Nevertheless, one can study the necessary conditions for the fluid parameters to enable the perturbations to grow and to give birth to large scale structures.
Let us consider the case where the equation of state of our fluid does not change during the growth of perturbations, and, to simplify, that the entropy perturbations can be ignored and that the Jeans length is smaller than any other considered scale. One can define the local density contrast of the dark fluid as:
$$\delta (\stackrel{}{x},t)\frac{\rho _D(\stackrel{}{x},t))}{\overline{\rho _D}(t)}1,$$
(26)
where $`\rho _D(\stackrel{}{x},t)`$ is the local value of the density, and $`\overline{\rho _D}(t)`$ is the mean background density, i.e. the apparent cosmological density. In the fluid approximation, one can write the evolution equation of the local density contrast Gaztañaga & Lobo (2001):
$`{\displaystyle \frac{1}{H^2}}{\displaystyle \frac{d^2\delta }{dt^2}}+\left(2+{\displaystyle \frac{\dot{H}}{H^2}}\right){\displaystyle \frac{1}{H}}{\displaystyle \frac{d\delta }{dt}}{\displaystyle \frac{2}{3}}(1+\omega _D)(1+3\omega _D)\mathrm{\Omega }_D\delta `$
$`={\displaystyle \frac{4+3\omega _D}{3(1+\omega _D)}}{\displaystyle \frac{1}{1+\delta }}\left({\displaystyle \frac{1}{1+\delta }}\right)^2{\displaystyle \frac{1}{H^2}}\left({\displaystyle \frac{d\delta }{dt}}\right)^2+{\displaystyle \frac{3}{2}}(1+\omega _D)(1+3\omega _D)\mathrm{\Omega }_D\delta ^2.`$ (27)
To solve this equation, one can define a new variable reflecting the expansion:
$$\eta =\mathrm{ln}a,$$
(28)
so that, if one assumes that the dark fluid is completely dominant at the time of growth of perturbations (in that case, the Friedmann equations reveal that $`\dot{H}/H^2=3(1+\omega _D)/2`$), equation (III) becomes:
$$\frac{d^2\delta }{d\eta ^2}+\frac{13\omega _D}{2}\frac{d\delta }{d\eta }\frac{2}{3}(1+\omega _D)(1+3\omega _D)\delta =\frac{4+3\omega _D}{3(1+\omega _D)}\frac{1}{1+\delta }\left(\frac{1}{1+\delta }\right)^2\left(\frac{d\delta }{d\eta }\right)^2+\frac{3}{2}(1+\omega _D)(1+3\omega _D)\delta ^2.$$
(29)
Because the coefficients of the above equation are time–dependant only, one can separate the spatial and temporal parts so that
$$\delta _l(\stackrel{}{x},t)=\delta _0\left(\stackrel{}{x}\right)D(t),$$
(30)
where $`D`$ is called the “linear growth factor”. In the linear approximation, where $`\delta `$ is small, equation (29) becomes:
$$\frac{d^2D}{d\eta ^2}+\frac{13\omega _D}{2}\frac{dD}{d\eta }\frac{2}{3}(1+\omega _D)(1+3\omega _D)D=0.$$
(31)
Its solutions take the form $`D=D_1a^{\alpha _1}+D_2a^{\alpha _2}`$, with $`D_1`$ and $`D_2`$ being two integration constants, and
$`\alpha _1`$ $`=`$ $`1+3\omega _D,`$ (32)
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{3}{2}}(1+\omega _D).`$ (33)
Thus, in the case of a dominant dark fluid, we have only a growing mode if $`\omega _D>1/3`$ or if $`\omega _D<1`$. One can however note that the last inequality seems very difficult to achieve with standard model for dark matter or dark energy models.
Let us now consider the observations of the cosmic microwave background to get constraints at earlier times.
## IV Cosmic Microwave Background
A power spectrum of temperature fluctuations can be deduced from the observations of the cosmic microwave background (CMB) Spergel et al. (2003). Predicting this power spectrum requires a hard work, and a program like CMBFAST Seljak & Zaldarriaga (1996) is able to produce it for the cosmological standard model. In our case, we will consider only the position of the peaks to constrain the parameters of the dark fluid, and we will make some assumptions on the dark fluid properties.
First, one should note that at high redshift, in the standard model the density of dark energy is nearly negligible in comparison to that of the dark matter. As the usual model seems to be able to correctly reproduce the fluctuations of the CMB, one can assume that our dark fluid should not behave very differently from the superposition dark matter/dark energy, and so should behave at the moment of recombination nearly like matter. Therefore, one can write the density of our fluid as a sum of a matter–like term ($`m`$) and of another term of unknown behavior ($`o`$):
$$\rho _D=\rho _{Dm}^{ls}\left(\frac{a}{a_{ls}}\right)^3+\rho _{Do}.$$
(34)
One can note that this equation gives no constraint on the behavior of the dark fluid, as the second term is not restricted to any behavior yet.
We do not want to specify a model of dark fluid and we would like to be as general as possible. Nevertheless, we will consider for simplicity only the background properties of the dark fluid, and we will not try to reproduce the whole power spectrum, but only consider the position of its peaks without trying to find their amplitude. The conformal time is defined by:
$$\tau =𝑑ta^1(t).$$
(35)
The spacing between the peaks is then given, to a good approximation, by Hu & Sugiyama (1995):
$$\mathrm{\Delta }l\pi \frac{\tau _0\tau _{ls}}{\overline{c}_s\tau _{ls}},$$
(36)
where $`\overline{c}_s`$ is the average sound speed before last scattering, and $`\tau _0`$ and $`\tau _{ls}`$ the conformal time today and at last scattering. This average sound speed reads:
$$\overline{c}_s\tau _{ls}^1_0^{\tau _{ls}}𝑑\tau \left(3+\frac{9\rho _b(t)}{4\rho _r(t)}\right)^{1/2},$$
(37)
where $`\rho _b`$ is the density of baryonic matter and $`\rho _r`$ is the density of relativistic fluids (radiation and neutrinos).
Let us consider that the Universe is flat so that the Friedmann equations are simplified. In this case, the first Friedmann equation can be written:
$$H^2=\frac{8\pi G}{3}\left(\rho _b+\rho _r+\rho _{Dm}+\rho _{Do}\right).$$
(38)
Using the evolution equation of the different densities, this becomes:
$$H^2=H_0^2\left(\mathrm{\Omega }_b^0a^3+\mathrm{\Omega }_r^0a^4+\mathrm{\Omega }_{Dm}^{ls}\left(\frac{a}{a_{ls}}\right)^3\right)+\frac{8\pi G}{3}\rho _{Do}.$$
(39)
The precedent equation cannot be solved if the form of $`\rho _{Do}`$ is not given. In our case, it is possible to assume that the fraction
$$\mathrm{\Omega }_{Do}(\tau )\frac{\rho _{Do}(\tau )}{\rho (\tau )}$$
(40)
does not vary too rapidly before the moment of last scattering (denoted $`ls`$), so that an effective average can be defined:
$$\overline{\mathrm{\Omega }}_{Do}^{ls}\tau _{ls}^1_0^{\tau _{ls}}\mathrm{\Omega }_{Do}(\tau )𝑑\tau .$$
(41)
In the following, we will only consider an approximate and effective density:
$$\rho _{Do}H^2\frac{3}{8\pi G}\overline{\mathrm{\Omega }}_{Do}^{ls}.$$
(42)
Let us then replace this density in the Friedmann equation:
$$H^2(1\overline{\mathrm{\Omega }}_{Do}^{ls})=H_0^2\left((\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)a^3+\mathrm{\Omega }_r^0a^4\right).$$
(43)
This time, provided one fixes the values of the different cosmological parameters and knowing the initial conditions, this equation can be solved. While replacing usual time by conformal time, the Friedmann equation becomes:
$$\left(\frac{da}{d\tau }\right)^2=H_0^2(1\overline{\mathrm{\Omega }}_{Do}^{ls})^1\left((\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)a(\tau )+\mathrm{\Omega }_r^0\right).$$
(44)
The value of the conformal time at the moment of last scattering is given by:
$$\tau _{ls}=2H_0^1\sqrt{\frac{1\overline{\mathrm{\Omega }}_{Do}^{ls}}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}\left\{\sqrt{a_{ls}+\frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}\sqrt{\frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}\right\}.$$
(45)
One can use the same method to evaluate the conformal time today. The Friedmann equation reads, after last scattering:
$$\left(\frac{da}{d\tau }\right)^2=H_0^2\left(\mathrm{\Omega }_b^0a(\tau )+\mathrm{\Omega }_r^0+a(\tau )^4\frac{\rho _D}{\rho _0^C}\right).$$
(46)
Let us make now the further assumption that the dark fluid has an approximate equation of state:
$$\rho _D=\stackrel{~}{\rho }_D^0a^{3(1+\overline{\omega }_D)},$$
(47)
where $`\stackrel{~}{\rho }_D^0`$ is an effective value of the dark fluid density, such that
$$\stackrel{~}{\rho }_D^0=a_{ls}^{3(1+\overline{\omega }_D)}\rho _D^{ls},$$
(48)
and $`\overline{\omega }_D`$ is the average value of $`\omega _D`$ over the conformal time, weighted by:
$$\mathrm{\Omega }_D(\tau )=\frac{\rho _D(\tau )}{\rho (\tau )}$$
(49)
to reflect the fact that the equation of state of our fluid should be more significant when its density contributes more heavily to the total density of the Universe. $`\overline{\omega }_D`$ is then given by:
$$\overline{\omega }_D\frac{{\displaystyle _0^{\tau _0}}\mathrm{\Omega }_D(\tau )\omega _D(\tau )𝑑\tau }{{\displaystyle _0^{\tau _0}}\mathrm{\Omega }_D(\tau )𝑑\tau }.$$
(50)
Defining the effective cosmological parameter $`\stackrel{~}{\mathrm{\Omega }}_D^0=3\stackrel{~}{\rho }_D^0/(8\pi GH_0^2)`$, the Friedmann equation becomes:
$$\left(\frac{da}{d\tau }\right)^2=H_0^2\left(\mathrm{\Omega }_b^0a(\tau )+\mathrm{\Omega }_r^0+\stackrel{~}{\mathrm{\Omega }}_D^0a(\tau )^{(13\overline{\omega }_D)}\right).$$
(51)
One can then integrate the equation, and show that:
$$\tau _0=2H_0^1F(\overline{\omega }_D),$$
(52)
with
$$F(\overline{\omega }_D)\frac{1}{2}_0^1𝑑a\left(\mathrm{\Omega }_b^0a+\mathrm{\Omega }_r^0+\stackrel{~}{\mathrm{\Omega }}_D^0a^{(13\overline{\omega }_D)}\right)^{1/2}.$$
(53)
There is no analytical integration for this function, but in a few cases. In particular, one has for a flat Universe:
$$F(0)=\frac{1}{\mathrm{\Omega }_b^0+\stackrel{~}{\mathrm{\Omega }}_D^0}\left(\sqrt{\mathrm{\Omega }_b^0+\mathrm{\Omega }_r^0+\stackrel{~}{\mathrm{\Omega }}_D^0}\sqrt{\mathrm{\Omega }_r^0}\right),$$
(54)
and
$$F(1/3)=\frac{1}{2}(\mathrm{\Omega }_D^0)^{1/2}\mathrm{ln}\left(\frac{1\stackrel{~}{\mathrm{\Omega }}_r^0+2\sqrt{\mathrm{\Omega }_D^0}}{\mathrm{\Omega }_b^0+2\sqrt{\mathrm{\Omega }_r^0\mathrm{\Omega }_D^0}}\right).$$
(55)
The case $`\overline{\omega }_D=1/3`$ looks much more probable than $`\overline{\omega }_D=0`$ when one considers the value of $`\omega _D^0`$ which was obtained from the supernova data. One finally gets the spacing between peaks:
$$\mathrm{\Delta }l=\pi \overline{c}_s^1\left[F(\overline{\omega }_D)\sqrt{\frac{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}{1\overline{\mathrm{\Omega }}_{Do}^{ls}}}\left\{\sqrt{a_{ls}+\frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}\sqrt{\frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}\right\}^11\right].$$
(56)
The sound velocity $`\overline{c}_s`$ is then given by:
$$\overline{c}_s=\tau _{ls}^1H_0^1\sqrt{1\overline{\mathrm{\Omega }}_{Do}^{ls}}_0^{a_{ls}}𝑑a\left[\left(3+\frac{9\mathrm{\Omega }_b^0}{4\mathrm{\Omega }_r^0}a\right)\left((\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)a+\mathrm{\Omega }_r^0\right)\right]^{1/2}.$$
(57)
This equation can be integrated analytically, and one finally gets:
$`\overline{c}_s`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0}}\right)^{1/2}\{\sqrt{a_{ls}+{\displaystyle \frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3}}}\}^1\times `$
$`\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }_r^0(7\mathrm{\Omega }_b^0+4\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)+2\sqrt{3\mathrm{\Omega }_b^0(\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)(\mathrm{\Omega }_b^0a_{ls}+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^4+\mathrm{\Omega }_r^0)(3\mathrm{\Omega }_b^0a_{ls}+4\mathrm{\Omega }_r^0)}+6a_{ls}\mathrm{\Omega }_b^0(\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)}{\mathrm{\Omega }_r^0(7\omega _b^0+4\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)+4\mathrm{\Omega }_r^0\sqrt{3\mathrm{\Omega }_b^0(\mathrm{\Omega }_b^0+\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)}}}\right).`$
One can notice that $`\overline{c}_s`$ does not depend on $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$. The approximate value of $`a_{ls}`$ can be taken from Hu et al. (2001):
$$a_{ls}^11008(1+0.00124(\mathrm{\Omega }_b^0h^2)^{0.74})(1+c_1(\mathrm{\Omega }_{Dm}^{ls}a_{ls}^3)^{c_2}),$$
(59)
where
$`c_1`$ $`=`$ $`0.0783(\mathrm{\Omega }_b^0h^2)^{0.24}(1+39.5(\mathrm{\Omega }_b^0h^2)^{0.76})^1,`$
$`c_2`$ $`=`$ $`0.56(1+21.1(\mathrm{\Omega }_b^0h^2)^{1.28})^1.`$ (60)
Unfortunately, this is an implicit equation, and we cannot find $`a_{ls}`$ analytically.
Finally, one has a direct dependence between $`\mathrm{\Delta }l`$ and the parameters $`\mathrm{\Omega }_{Dm}^{ls}`$, $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$, $`\stackrel{~}{\mathrm{\Omega }}_D^0`$ and $`\overline{\omega }_D`$. Even if this is only an approximate formula, it gives the possibility to visualize directly the effect of the dark fluid on the position of the peaks, provided that its behavior does not differ much from the requirements of the approximations on $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$ and $`\overline{\omega }_D`$.
However, the calculated $`\mathrm{\Delta }l`$ cannot be directly related to the observed spacing between peaks, as shifts of peaks can be induced by other effects. In particular, the location of the i–th peaks can be approximated by:
$$l_i=\mathrm{\Delta }l(m\varphi _i)=\mathrm{\Delta }l(m\overline{\varphi }\delta \varphi _i),$$
(61)
where $`\overline{\varphi }`$ is the shift of the first peak, corresponding to an overall shift, and the $`\delta \varphi _i`$ is the specific shift of the i–th peak. We have fortunately access to the fitting formulae of Doran & Lilley (2002):
Overall phase shift $`\overline{\varphi }`$ :
$$\overline{\varphi }=(1.4660.466n)\left[a_1r_{ls}^{a_2}+0.291\overline{\mathrm{\Omega }}_{Do}^{ls}\right],$$
(62)
where $`a_1`$ and $`a_2`$ are given by:
$`a_1`$ $`=`$ $`0.286+0.626\left(\mathrm{\Omega }_bh^2\right),`$
$`a_2`$ $`=`$ $`0.17866.308\mathrm{\Omega }_bh^2+174.9\left(\mathrm{\Omega }_bh^2\right)^21168\left(\mathrm{\Omega }_bh^2\right)^3,`$ (63)
$`n`$ is the spectral index and $`r_{ls}`$ is defined by:
$$r_{ls}=\frac{\rho _r(a_{ls})}{\rho _{Dm}(a_{ls})}=\frac{\mathrm{\Omega }_r^0}{\mathrm{\Omega }_{Dm}^{ls}a_{ls}^4}.$$
(64)
Relative shift of second peak $`\delta \varphi _2`$ :
$$\delta \varphi _2=c_0c_1r_{ls}c_2r_{ls}^{c_3}+0.05(n1),$$
(65)
with
$`c_0`$ $`=`$ $`0.1+\left(0.2130.123\overline{\mathrm{\Omega }}_{Do}^{ls}\right)\mathrm{exp}\left\{\left(5263.6\overline{\mathrm{\Omega }}_{Do}^{ls}\right)\mathrm{\Omega }_bh^2\right\},`$
$`c_1`$ $`=`$ $`0.063\mathrm{exp}\left\{3500\left(\mathrm{\Omega }_bh^2\right)^2\right\}+0.015,`$ (66)
$`c_2`$ $`=`$ $`6\times 10^6+0.137\left(\mathrm{\Omega }_bh^20.07\right)^2,`$
$`c_3`$ $`=`$ $`0.8+2.3\overline{\mathrm{\Omega }}_{Do}^{ls}+\left(70126\overline{\mathrm{\Omega }}_{Do}^{ls}\right)\mathrm{\Omega }_bh^2.`$
Relative shift of third peak $`\delta \varphi _3`$ :
$$\delta \varphi _3=10d_1r_{ls}^{d_2}+0.08(n1),$$
(67)
with
$`d_1`$ $`=`$ $`9.97+\left(3.33\overline{\mathrm{\Omega }}_{Do}^{ls}\right)\mathrm{\Omega }_bh^2,`$ (68)
$`d_2`$ $`=`$ $`0.00160.0067\overline{\mathrm{\Omega }}_{Do}^{ls}+\left(0.1960.22\overline{\mathrm{\Omega }}_{Do}^{ls}\right)\mathrm{\Omega }_bh^2+{\displaystyle \frac{(2.25+2.77\overline{\mathrm{\Omega }}_{Do}^{ls})\times 10^5}{\mathrm{\Omega }_bh^2}}.`$
One can now compare our results to the data. The WMAP experiment provides the precise location of the two first peaks Spergel et al. (2003):
$`l_{p_1}`$ $`=`$ $`220.1\pm 0.8,`$
$`l_{p_2}`$ $`=`$ $`546\pm 10,`$ (69)
and BOOMERanG gives the position of the third peak Bernardis et al. (2002):
$$l_{p_3}=825\pm 13.$$
(70)
To evaluate roughly the value of the parameters of the fluid, one can fix the other parameters as follows:
$`n`$ $`=`$ $`1,`$
$`h`$ $`=`$ $`0.70,`$ (71)
$`\mathrm{\Omega }_b^0`$ $`=`$ $`0.049,`$
$`\mathrm{\Omega }_r^0`$ $`=`$ $`9.89\times 10^5.`$
For these values, one finds in Table 1 the resulting positions of the peaks in function of parameters of the dark fluid.
One can note that a large range of values is possible. It is not so strange because we have many parameters for our fluid. A more complete analysis is not needed here, because the strongest constraints would come from the specification of a model. Without specifying a model, Table 1 shows that the values of the parameters of the dark fluid are not stringently constrained. One can nevertheless see that for large values of $`\stackrel{~}{\mathrm{\Omega }}_D^0`$, the permitted values of $`\overline{\omega }_D`$ are negative, and hence in that case one can assume that our fluid behaves today like a cosmological constant whereas it could have behaved mainly like matter at last scattering. For small values of $`\stackrel{~}{\mathrm{\Omega }}_D^0`$, $`\overline{\omega }_D`$ is positive, so that the density of dark fluid should decrease more rapidly than a matter density after last scattering. In this case, $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$ is very small, and the fluid should have behaved like matter before and around last scattering. Small values of $`\stackrel{~}{\mathrm{\Omega }}_D^0`$ look therefore unrealistic, because as $`\overline{\omega }_D`$ is then positive, unless our dark fluid has an oscillating density, $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$ should be much larger and certainly dominant.
The value of $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$, which can be as much as 0.1, also shows that before recombination, the fluid may have behaved differently from matter, and perhaps like radiation. One can also note that low values of $`\stackrel{~}{\mathrm{\Omega }}_D^0`$ corresponds to low values of $`\overline{\mathrm{\Omega }}_{Do}^{ls}`$, and then in that case the fluid mostly behaves like matter.
If one combines these results with the results from the CMB, it seems that the constraints become $`1/3<\omega _D<0`$ to enable the perturbations to grow. This shows in fact that during the growth of the perturbations the behavior of the dark fluid should not be too different from that of matter. This result also confirms that the value of the effective $`\stackrel{~}{\mathrm{\Omega }}_D^0`$ which appears in equation (51) cannot be too small (it has to be at least larger than 0.2).
If one wants to perform a much more precise study of a specified model, it would be interesting to simulate the whole process of structure formation, and to compare the results with surveys like SDSS Tegmark et al. 2004 II , or 2dF Percival et al. (2001). We will not study here further the CMB power spectrum, as other features seem more model–dependant, and we want to consider here the general case. We will now consider the results of the big-bang nucleosynthesis and their influence on the establishment of a dark fluid model.
## V Big-Bang Nucleosynthesis
Recent analyses of the big–bang nucleosynthesis (BBN) Coc et al. (2004) indicate a discrepancy between the value of the baryonic density calculated from the observed Li and <sup>4</sup>He abundances, and the one calculated with the observations of deuterium. Some explanations can be found. Problems could have appeared in the measurement of Li and <sup>4</sup>He abundances, or the Li on the stellar surface could be altered during stellar evolution, or we have no accurate knowledge of the reaction rates related to <sup>7</sup>Be destruction, or the expansion rate during BBN could have been modified through an accelerated cosmological expansion. Thus, two possibilities can be considered for the equation of state of the dark fluid. First, if it is correct to consider a Universe dominated by radiation at BBN time, the main constraint is that the dark fluid density should be small in comparison to the radiation density; otherwise Friedmann equations indicates that the expansion rate of the Universe would be different from the one in the standard BBN, changing then the temperature evolution rate and so the abundance of the elements. It means that, if one assumes that the dark fluid behavior does not change violently during BBN, the equation of state of this fluid around the time of BBN has to be $`\omega _D(\text{BBN})1/3`$, or that its density was completely negligible before BBN. In the case of a real radiative behavior $`\omega _D(\text{BBN})=1/3`$, the dark fluid behaves like extra–families of neutrinos, and its density can be constrained. The effective extra–neutrinos number at the BBN time is defined by:
$$\mathrm{\Delta }N_{\text{eff}}(\text{BBN})\frac{\rho _D(\text{BBN})}{\rho _\nu (\text{BBN})},$$
(72)
where $`\rho _\nu `$ is the standard density of a single relativistic neutrino species. The usual bound on the number of neutrinos is $`\mathrm{\Delta }N_{\text{eff}}<1`$ Burles & Tyler (1998), which corresponds in our case for a temperature around 1 MeV ($`a^134\times 10^9`$) to:
$$\rho _D(\text{BBN})<\frac{7}{8}\left(\frac{4}{11}\right)^{4/3}\frac{\pi ^2}{15}T_{\text{BBN}}^43\times 10^2(\text{MeV})^4.$$
(73)
This limit is only valid in the case $`\omega _D(\text{BBN})=1/3`$.
If the abundance of the elements is as observed, a modification of the expansion rate could provide, as presented in Salati (2003), a correction to the predicted values. If our fluid is the dominant component at the time of the BBN, it can have a big influence on the expansion rate of the Universe. Evaluating the density of the dark fluid and its evolution during the time of the BBN so that the observations are retrieved would require further studies that I will not develop here.
We have seen that a dark fluid may be compatible with the cosmological observations, and could be an interesting approach to the ambivalence of dark energy and dark matter. Let us now consider possible paths to model a dark fluid.
## VI Models of Dark Fluid
In the literature, only few fluids behaving like a dark fluid are considered. Different ways to model the dark fluid are possible, and I will consider here in particular two of them: the generalized Chaplygin gas, based on D-brane theories, and another one using scalar fields.
### VI.1 Generalized Chaplygin Gas
The generalized Chaplygin gas (GCG) is an exotic fluid derived from D-brane theories Bento et al. (2002). It can be described by an equation of state:
$$P_{ch}=\frac{A}{\rho _{ch}^\alpha },$$
(74)
where $`\alpha `$ is a constant, $`0<\alpha 1`$, and $`A`$ is another positive constant. This equation of state corresponds to a density evolving like:
$$\rho _{ch}=\rho _{ch}^0\left(A_s+\frac{1A_s}{a^{3(1+\alpha )}}\right)^{1/(1+\alpha )},$$
(75)
where $`A_s=A/(\rho _{ch}^0)^{(1+\alpha )}`$ and $`\rho _{ch}^0`$ is the Chaplygin gas density today. Such a behavior could be interesting in order to model the dark fluid, because for high values of $`a`$ this density is mainly constant, and for low values of $`a`$ it evolves like matter. This behavior has to be compared with the observations.
For the comparison with the data of supernovæ of type Ia, let us derive the equation of state of the GCG at low redshift:
$$\omega _{ch}=\frac{P_{ch}}{\rho _{ch}}=\frac{A_s}{A_s+(1A_s)(1+z)^{3(1+\alpha )}}A_s\left[13(1A_s)(1+\alpha )z\right],$$
(76)
so that one can deduce:
$`\omega _{ch}^0`$ $`=`$ $`A_s,`$
$`\omega _{ch}^1`$ $`=`$ $`3A_s(1A_s)(1+\alpha ),`$ (77)
to be compared with the constraints (25), and one finally gets:
$`As`$ $`=`$ $`0.76\pm 0.25,`$
$`\alpha `$ $`=`$ $`0.8\pm 2.4.`$ (78)
We have of course no constraint on $`\alpha `$. The analysis of CMB has been done in Bento et al. (2003), and the results for $`h=0.7`$ and $`n=1`$, combined with the ones from the supernovæ of type Ia are:
$`As`$ $`=`$ $`0.75\pm 0.13,`$ (79)
$`\alpha `$ $`=`$ $`0.4\pm 0.2.`$
If one considers now the BBN, at this time the GCG behaves like matter, so that it is compatible with the standard BBN scenario. The large scale structure formation has been studied in a case where the Chaplygin gas adds only a background density to a Universe containing cold dark matter Multamäki et al. (2004), and such a scenario seems then possible.
So, the GCG seems to be in agreement with the cosmological observations. A further analysis is still needed, in particular concerning the growth of structures with a dominant Chaplygin gas density, or the local behavior of such a fluid.
### VI.2 Scalar Fields
One can also consider the idea that the dark fluid could be explained thanks to a scalar field. Indeed, scalar fields are very useful in explaining the behavior of the dark energy today Ratra & Peebles (2000); Hebecker & Wetterich (2001), and recent analyses have shown that they can behave like matter on local scales Arbey et al. (2003); Kiselev (2005) as well as on cosmological scales Arbey et al. (2002); Guzman & Urena–Lopez (2003).
Let us therefore consider a real scalar field associated with a Lagrangian density
$$=g^{\mu \nu }_\mu \phi _\nu \phi V\left(\phi \right).$$
(80)
Its pressure and its density on cosmological scales are given by
$`P_\phi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\phi }^2V\left(\phi \right),`$
$`\rho _\phi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\phi }^2+V\left(\phi \right).`$ (81)
So, the pressure is negative if the potential dominates, and negligible if the potential equilibrates the kinetic term. Thus, a scalar field can be a good candidate for the dark fluid if it respects in particular the following constraints:
– its density at the time of the BBN decreases at least as fast as the density of radiation,
– its density from the time of last scattering to the time of structure formation evolves nearly like matter, and so $`\frac{1}{2}\dot{\phi }^2V\left(\phi \right)`$,
– after the growth of perturbations, because the scalar field is dominating the Universe, its potential does not equilibrate the kinetic term anymore and begins to dominate; thus it will behave like a cosmological constant in the future.
The main parameter of such a model is the same as that of quintessence models: the potential. When one considers quintessence models with real scalar field, one looks for potentials which provide a cosmological constant–like behavior today, and decreasing potentials seem to be favored. If one considers now complex scalar fields, it was shown in Arbey et al. (2002) that such a field can behave like cosmological matter when its potential has a dominant $`m^2|\varphi |^2`$ term, which corresponds to an increasing term in the potential. Thus, a way to find a “good” potential would be to consider a superposition of a decreasing potential which would begin to dominate today, and of the increasing quadratic term which has to dominate at least until structure formation and can nevertheless lead to an attractive effect on local scales today. A more detailed study of the possible potentials will be presented in a future publication.
One could also hope that a scalar field might explain the excess of gravity on local scales. To do that, one can imagine a Universe filled with a scalar field. In the part – and time – of the Universe where the density of baryonic matter is high, the scalar field would, through gravitational interaction, have a large kinetic term which could even equilibrate the potential, so that one has an attractive net force on local scales (easier to achieve with a complex scalar field associated to an internal rotation, see Arbey et al. (2003)), whereas in the parts where no baryons are present, the field would not vary much, and the potential dominates, providing repulsion. Thus, on local scales, where the baryon density is high, the field behaves like matter. Where the baryonic density is small, i.e. away from galaxies and clusters, the gravitational interaction is not strong enough to increase the kinetic term of the scalar field, so that the potential dominates, and one can then observe the effects of a negative pressure. In that case, the scalar field will have a negative pressure on cosmological scales, providing a locally negative pressure in average. In the past, baryons were uniformly dense, so that the kinetic term was large everywhere, and one could have then a uniform matter behavior under these conditions. In that way, the local behavior can be in agreement with the cosmological one, and a complete cosmological scenario can be built. Such a scenario has of course to be studied further.
## VII Conclusion and Perspectives
Astrophysical and cosmological observations are usually interpreted in terms of dark matter and dark energy. We have seen here that they can also be analyzed differently. Thus, it is possible to develop a model of dark fluid, which could advantageously replace a model containing in fact two dark components. Of course, hard work and studies are required to test completely the dark fluid hypothesis. Nevertheless today, as it seems difficult to find a model for dark energy and as problems concerning cold dark matter remain, it is worthwhile to investigate different ideas such as an unification of dark energy and dark matter – that finally does not seem stranger than trying to determine the nature of two components at the same time – which could be achieved in particular thanks to D-brane theories (in particular through the Chaplygin gas), or thanks to the so–useful scalar fields. Of course, other models may also account for dark fluid.
An important question remains, how to interpret the dark matter problem on local scales and could the dark fluid account for the excess of gravity inside local structures? I will only provide here a qualitative analysis of whether a fluid with a negative pressure on cosmological scale can have an attractive effect on local scale, such as it is observed in galaxies (for example, with the rotation curves of spiral galaxies Persic et al. (1986); Gentile et al. (2004); Carignan & Purton (1998)).
Let us consider only the quasi–Newtonian limit of general relativity. In that case, deviations from the Minkowski metric $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ are accounted for by the perturbation $`h_{\mu \nu }`$. In the harmonic coordinate gauge, it satisfies the condition:
$$_\alpha h_\mu ^\alpha \frac{1}{2}_\mu h_\alpha ^\alpha =\mathrm{\hspace{0.33em}0}.$$
(82)
One can show that the perturbation $`h_{\mu \nu }`$ is related to the source tensor:
$$S_{\mu \nu }=T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T_\alpha ^\alpha $$
(83)
through the integral
$$h_{\mu \nu }\left(\stackrel{}{r}\right)=\mathrm{\hspace{0.17em}4}G\frac{S_{\mu \nu }\left(\stackrel{}{r}^{}\right)}{|\stackrel{}{r}^{}\stackrel{}{r}|}d^3\stackrel{}{r}^{}.$$
(84)
If the energy–momentum tensor is written as:
$$T^{\mu \nu }=(P+\rho )U^\mu U^\nu Pg^{\mu \nu },$$
(85)
one can show that, at leading order in $`h_{\mu \nu }`$,
$$S^{\mu \nu }=(P+\rho )U^\mu U^\nu \frac{1}{2}\eta ^{\mu \nu }(\rho P).$$
(86)
The gravitational potential is in fact $`\mathrm{\Phi }=h_{00}/2`$, so that:
$$\mathrm{\Phi }\left(\stackrel{}{r}\right)=\mathrm{\hspace{0.17em}2}G\frac{S_{00}\left(\stackrel{}{r}^{}\right)}{|\stackrel{}{r}^{}\stackrel{}{r}|}d^3\stackrel{}{r}^{}.$$
(87)
For a fluid at rest, $`U^\mu =(1,0,0,0)`$, with an equation of state $`P=\omega \rho `$, one has:
$$S_{00}=\frac{1}{2}\rho (1+3\omega ),$$
(88)
and consequently this fluid has an attractive effect only if $`\omega >1/3`$. From the study of supernovæ it seems that our dark fluid is not in that state today, so that its effects are mainly repulsive. Nevertheless, it is possible that the dark fluid has a different behavior on cosmological and on local scales. Indeed, the density and pressure of the fluid on cosmological scale are spacial averages of the local density and pressure, and one can assume that:
$`\rho (t,\stackrel{}{r})`$ $`=`$ $`\rho ^{\text{cosmo}}\left(t\right)+\delta \rho (t,\stackrel{}{r}),`$
$`P(t,\stackrel{}{r})`$ $`=`$ $`P^{\text{cosmo}}\left(t\right)+\delta P(t,\stackrel{}{r}),`$ (89)
where $`\rho ^{\text{cosmo}}`$ and $`P^{\text{cosmo}}`$ are the cosmological density and pressure, with the spatial averages:
$$<\delta \rho (t,\stackrel{}{r})>=<\delta P(t,\stackrel{}{r})>=0.$$
(90)
The density of dark fluid on cosmological scales today is of the order of the critical density, i.e. $`\rho _c^09\times 10^{29}\text{g.cm}^3`$. One can compare it to the estimated matter density in the Milky Way at the radius of the Sun $`\rho ^{\text{Sun}}5\times 10^{24}\text{g.cm}^3`$ Olling & Merrifield (2001). Hence, even if the dark fluid’s local density would represent 1% of this total matter local density only, its value would be much higher than the cosmological densities today. Therefore, on local scales, one can assume that $`\delta \rho (t,\stackrel{}{r})\rho ^{\text{cosmo}}\left(t\right)`$ and thus write:
$$S_{00}\frac{1}{2}(\delta \rho +3\delta P).$$
(91)
To have a net attraction, we get finally the same kind of constraint as before, $`\delta \rho >3\delta P`$, but this time we do not have to use the cosmological constraints, because the local behavior of the dark fluid can be very different from the cosmological one. Moreover, if no model is specified, one can still hope that $`\delta P`$ could be negligible on local scales, so that the local behavior of the dark fluid is matter-like, and that the usual Newtonian equation can be retrieved:
$$\mathrm{\Phi }\left(\stackrel{}{r}\right)=G\frac{\delta \rho \left(\stackrel{}{r}^{}\right)}{|\stackrel{}{r}^{}\stackrel{}{r}|}d^3\stackrel{}{r}^{}.$$
(92)
This local behavior will of course have to be verified quantitatively for each dark fluid model, but nevertheless gives hope for a unified explanation on any scale. In particular, considering scalar fields, this qualitative analysis tends to show that it would be interesting to try to find a potential which gives today a negative pressure on cosmological scale, but which also gives a matter behavior in local structures, i.e. where the density of baryons is high.
To conclude, the dark fluid appears as an interesting possibility to explain the observations. As the properties of the dark fluid are different from dark matter and dark energy, models of dark fluid are worth to be studied, and we can now use many precise observations as strong constraints on such models.
## Acknowledgements
I would like to thank Farvah Mahmoudi, Hélène Courtois, Julien Devriendt, Thierry Sousbie and Wolfgang Hillebrandt for their comments and for useful discussions.
|
warning/0506/math0506075.html
|
ar5iv
|
text
|
# Parallel transport of 𝐻𝑜𝑚-complexes and the Lovász conjecture
## 1 Introduction
Hidden in the background of the Babson and Kozlov proof of Lovász conjecture are interesting topological and combinatorial concepts and structures associated to graphs and graph homomorphisms. The proof itself runs in two phases, each phase divided in several steps, often involving a detailed case analysis. For this reason the underlying secondary structures may not be visible or immediately recognized under the layers of intricate technical details. Recall that the crux of Babson and Kozlov approach is a skilful and technically quite involved application of spectral sequences. One of the classical applications of this method is to the (co)homology of fibered spaces which by the nature are spaces which alow some form of transport from one fibre to another.
In this paper we focus on one of these secondary structures which can be, somewhat informally, described as the “parallel transport” of graph complexes over graphs.
The introduction of this structure and recognition of its role leads to a great simplification of the proof of Lovász conjecture in some cases. Another of its features, aside from offering a conceptual “explanation” for the success of one of the technical approaches of Babson and Kozlov, is its potential for generating other statements of this type.
The “parallel transport” of graph complexes introduced here seems to be a novel concept. However the groupoids (groups of projectivities) used in its definition have already appeared in geometric combinatorics in the work of M. Joswig , see also , where they have been applied to toric manifolds, branched coverings over $`S^3`$, colorings of simple polytopes etc.
It is an exciting “coincidence” that there have been other recent developments in geometric and algebraic combinatorics where groupoids and associated objects and constructions were implicitly used, . These are not isolated examples of course. In particular one should be fully aware of a rich and deep combinatorics already present in numerous categorical constructions related to groupoids and their applications in geometry and mathematical physics.
## 2 The Lovász conjecture
One of central themes in topological combinatorics, after the landmark paper of Laszlo Lovász where he proved the classical Kneser conjecture, has been the study and applications of graph complexes.
The underlying theme is to explore how the topological complexity of a graph complex $`X(G)`$ reflects in the combinatorial complexity of the graph $`G`$ itself. The results one is usually interested in come in the form of inequalities $`\alpha (X(G))\xi (G)`$, or equivalently in the form of implications
$$\alpha (X(G))p\xi (G)q,$$
where $`\alpha (X(G))`$ is a topological invariant of $`X(G)`$, while $`\xi (G)`$ is a combinatorial invariant of the graph $`G`$.
The most interesting candidate for the invariant $`\xi `$ has been the chromatic number $`\chi (G)`$ of $`G`$, while the role of the invariant $`\alpha `$ was played by the “connectedness” of $`X(G)`$, its equivariant index, the height of an associated characteristic cohomology class etc., see for recent accounts.
The famous result of Lovász quoted above is today usually formulated in the form of an implication
$$Hom(K_2,G)\text{ is }k\text{-connected }\chi (G)k+3,$$
(1)
where $`Hom(K_2,G)`$ is the so called “box complex” of $`G`$. The box complex is a special case of a general graph complex $`Hom(H,G)`$ (also introduced by L. Lovász), a cell complex which functorially depends on the input graphs $`H`$ and $`G`$.
An outstanding conjecture in this area, refereed to as “Lovász conjecture”, was that one obtains a better bound if the graph $`K_2`$ in (1) is replaced by an odd cycle $`C_{2r+1}`$. More precisely Lovász conjectured that
$$Hom(C_{2r+1},G)\text{ is }k\text{-connected }\chi (G)k+4.$$
(2)
This conjecture was confirmed by Babson and Kozlov in , see also for a more detailed exposition.
Our objective is to develop methods which both offer a simplified approach to the proof of implication (2), at least in the case when $`k`$ is odd, and providing new insight, open a possibility of proving similar results for other classes of (hyper)graphs and simplicial complexes.
An example of such a result is Theorem 22. One of its corollaries is the following implication,
$$Hom(\mathrm{\Gamma },K)\text{ is }k\text{-connected }\chi (K)k+d+3$$
(3)
which, under a suitable assumption on the complex $`\mathrm{\Gamma }`$ and the assumption that integer $`k`$ is odd, extends (2) to the case of pure $`d`$-dimensional simplicial complexes.
## 3 Parallel transport of $`Hom`$-complexes
### 3.1 Generalities about “parallel transport”
In order to avoid any ambiguities, we briefly clarify what is in this paper meant by a “parallel transport” on a “bundle” of spaces.
A “bundle” is a map $`\varphi :XS`$. We assume that $`S`$ is a set and that $`X(i):=\varphi ^1(i)`$ is a topological space, so a bundle is just a collection of spaces (fibres) $`X(i)`$ parameterized by $`S`$. If all spaces $`X(i)`$ are homeomorphic to a fixed “model” space, this space is referred to as the fiber of the bundle $`\varphi `$.
Suppose that $`𝒢`$ is a groupoid on $`S`$ as the set of objects. In other words $`𝒢=(Ob(𝒢),Mor(𝒢))`$ is a small category where $`Ob(𝒢)=S`$, such that all morphisms $`\alpha Mor(𝒢)`$ are invertible.
A “connection” or “parallel transport” on the bundle $`𝒳=\{X(i)\}_{iS}`$ is a functor $`𝒫:𝒢Top`$ such $`X(i)=𝒫(i)`$ for each $`iS`$.
Informally speaking, the groupoid $`𝒢`$ provides a “road map” on $`S`$, while the functor $`𝒫`$ defines the associated transport from one fibre to another.
Sometimes it is convenient to view the bundle $`𝒳=\{X(i)\}_{iS}`$ as a map $`𝒳:STop`$. Then to define a “connection” on this bundle is equivalent to enriching the map $`𝒳`$ to a functor $`𝒫:𝒢Top`$.
### 3.2 Natural bundles and groupoids over simplicial complexes
Suppose that $`K`$ and $`L`$ are finite simplicial complexes and let $`k`$ be an integer such that $`0k\mathrm{dim}(K)`$. Let $`S_k=S_k(K)`$ be the set of all $`k`$-dimensional simplices in $`K`$. Define a bundle $`_k^L:S_kTop`$ by the formula
$$_k^L(\sigma )=Hom(\sigma ,L)L_\mathrm{\Delta }^{k+1}$$
(4)
where $`Hom(\sigma ,L)`$ is one of the $`Hom`$-complexes introduced in Section 5.1 and $`L_\mathrm{\Delta }^{k+1}`$ is the complex well known in topological combinatorics as the deleted product of $`L`$, Chapter 6. A typical cell in $`L_\mathrm{\Delta }^{k+1}`$ is of the form $`e=\sigma _0\times \sigma _1\times \mathrm{}\times \sigma _kL^{k+1}`$ where $`\{\sigma _i\}_{i=0}^k`$ is a collection of non-empty simplices in $`L`$ such that if $`ij`$ then $`\sigma _i\sigma _j=\mathrm{}`$. The corresponding cell in $`Hom(\sigma ,L)`$ is a function $`\eta :V(\sigma )L\{\mathrm{}\}`$, where $`V(\sigma )=\sigma ^{(0)}`$ is the set of all vertices of $`\sigma `$, and if $`v_1v_2`$ then $`\eta (v_1)\eta (v_2)=\mathrm{}`$.
Example: It is well known that if $`L\sigma ^m=\mathrm{\Delta }^{[m+1]}`$ is a $`m`$-dimensional simplex, then the associated deleted square $`(\sigma ^m)_\mathrm{\Delta }^2`$ is homeomorphic to a $`(m1)`$-dimensional sphere. In other words, $`_1^{\sigma ^m}:S_1(K)Top`$ is a spherical bundle naturally associated to the simplicial complex $`K`$.
Our next goal, in the spirit of Section 3.1, is to identify a groupoid on the set $`S_k`$ which acts on the bundle $`_k^L`$, i.e. a groupoid which provides a parallel transport of fibres of the bundle $`_k^L`$. It is a pleasant coincidence that this groupoid has already appeared in geometric combinatorics . Indeed, the groups of projectivities M. Joswig introduced and studied in these papers are just the vertex or isotropy groups of a groupoid which we call the $`k`$-th groupoid of projectivities of $`K`$ and denote by $`𝒢_k^P(K)`$. In these and in subsequent papers , the groups of projectivities found interesting applications to toric manifolds, branched coverings over $`S^3`$, colorings of simple polytopes, etc. Here is a summary of this construction.
Two $`k`$-dimensional simplices $`\sigma _0`$ and $`\sigma _1`$ in $`K`$ are called adjacent if they share a common $`(k1)`$-dimensional face $`\tau `$. A perspectivity from $`\sigma _0`$ to $`\sigma _1`$ is the unique non-degenerated simplicial map $`\stackrel{}{\sigma _0\sigma _1}=\sigma _0,\sigma _1:\sigma _0\sigma _1`$ which leaves the simplex $`\tau `$ point-wise fixed. In the special case when $`\sigma _0=\sigma _1`$, the perspectivity $`\sigma _0,\sigma _0:\sigma _0\sigma _0`$ is the identity map $`I_{\sigma _0}`$.
A projectivity between two, not necessarily adjacent, simplices $`\sigma _0`$ and $`\sigma _n`$ is a composition of perspectivities
$$𝔭=\stackrel{}{\sigma _0\sigma _1}\stackrel{}{\sigma _1\sigma _2}\mathrm{}\stackrel{}{\sigma _{n1}\sigma _n}$$
where $`𝔭=(\sigma _0,\sigma _1,\mathrm{},\sigma _n)`$ is a path of $`k`$-dimensional simplices in $`K`$ such that $`\sigma _{i1}`$ and $`\sigma _i`$ share a common $`(k1)`$-dimensional face $`\tau _i`$.
Caveat: Here we adopt a useful convention that $`(x)(fg)=(gf)(x)`$ for each two composable maps $`f`$ and $`g`$. The notation $`fg`$ is often given priority over the usual $`gf`$ if we want to emphasize that the functions act on the points from the right, that is if the arrows in the associated formulas point from left to the right.
###### Definition 1
The $`k`$-th groupoid of projectivities $`𝒢_k^P(K)`$ of a simplicial complex $`K`$, or the $`P_k`$-groupoid associated to $`K`$, is the small category
$$𝒢_k^P(K)=(Ob(𝒢_k^P(K)),Mor(𝒢_k^P(K)))$$
which has the set $`S_k=Ob(𝒢_k^P(K))`$ of all $`k`$-dimensional simplices for the set of objects, and for each two simplices $`\sigma _0,\sigma _1S_k`$, the associated morphism set $`Mor_{𝒢_k^P(K)}(\sigma _0,\sigma _1)`$ is the collection of all projectivities from $`\sigma _0`$ to $`\sigma _1`$. The associated point (isotropy) groups
$$\mathrm{\Pi }_k(K,\sigma _0):=Mor_{𝒢_k^P(K)}(\sigma _0,\sigma _0)$$
are called the groups of projectivities or the combinatorial holonomy groups of $`K`$ .
###### Proposition 2
For each finite simplicial complex $`K`$ and an auxiliary “coefficient” complex $`L`$, there exists a canonical connection $`𝒫^L=𝒫_{K,k}^L`$ on the bundle $`_k^L`$. In other words the function $`_k^L:S_kTop`$ can be enriched (extended) to a functor
$$_k^L:S_kTop.$$
Proof: If $`\stackrel{}{\sigma _0\sigma _1}`$ is a perspectivity from $`\sigma _0`$ to $`\sigma _1`$ and if $`\eta :V(\sigma _1)2^{V(L)}\{\mathrm{}\}`$ is a cell in $`Hom(\sigma ,L)`$, then $`𝒫^L:^L(\sigma _1)^L(\sigma _0)`$ is the map defined by $`𝒫^L(\stackrel{}{\sigma _0\sigma _1})(\eta ):=\stackrel{}{\sigma _0\sigma _1}\eta `$. More generally, if $`𝔭=\stackrel{}{\sigma _0\sigma _1}\stackrel{}{\sigma _1\sigma _2}\mathrm{}\stackrel{}{\sigma _{n1}\sigma _n}`$ is a projectivity between $`\sigma _0`$ and $`\sigma _n`$, then
$$𝒫^L(𝔭)=𝒫^L(\stackrel{}{\sigma _0\sigma _1})𝒫^L(\stackrel{}{\sigma _1\sigma _2})\mathrm{}𝒫^L(\stackrel{}{\sigma _{n1}\sigma _n})$$
(5)
or in other words
$$𝒫^L(𝔭)(\eta )=\stackrel{}{\sigma _0\sigma _1}\stackrel{}{\sigma _1\sigma _2}\mathrm{}\stackrel{}{\sigma _{n1}\sigma _n}\eta .$$
(6)
It is clear from the construction that the map $`𝒫^L(𝔭)`$ depends only on the projectivity $`𝔭`$ and not on the associated path $`𝔭`$. $`\mathrm{}`$
### 3.3 Parallel transport of graph complexes
The main motivation for introducing the parallel transport of $`Hom`$-complexes is the Lovász conjecture and its ramifications. This is the reason why the case of graphs and the graph complexes deserves a special attention. Additional justification for emphasizing graphs comes from the fact that graph complexes $`Hom(G,H)`$ have been studied in numerous papers and today form a well established part of graph theory and topological combinatorics. The situation with simplicial complexes is quite the opposite. In order to extend the theory of $`Hom`$-complexes from graphs to the category of simplicial complexes, many concepts should be generalized and the corresponding facts established in a more general setting. One is supposed to recognize the main driving forces and to isolate the most desirable features of the theory. A result should be a dictionary/glossary of associated concepts, cf. Table 1. Consequently, Section 3.3 should be viewed as an important preliminary step, leading to the more general theory developed in Sections 5 and 6.
In order to simplify the exposition we assume, without a serious loss of generality, that all graphs $`G=(V(G),E(G))`$ are without loops and multiple edges. In short, graphs are $`1`$-dimensional simplicial complexes. Let $`G_{\overline{xy}}K_2`$ be the restriction of $`G`$ on the edge $`\overline{xy}E(G)`$.
Following the definitions from Section 3.2 the map
$$^H:E(G)Top,$$
where $`^H(\overline{xy}):=_{\overline{xy}}^H=Hom(G_{\overline{xy}},H)`$, can be thought of as a “bundle” over the graph $`G`$, with $`_{\overline{xy}}^H=Hom(G_{\overline{xy}},H)`$ in the role of the “fibre” over the edge $`\overline{xy}`$. More generally, given a class $`𝒞`$ of subgraphs of $`G`$, say the subtrees, the chains, the $`k`$-cliques etc., one can define an associated “bundle” $`_𝒞^H:𝒞Top`$ by a similar formula $`_𝒞^H(\mathrm{\Gamma }):=Hom(\mathrm{\Gamma },H)`$, where $`\mathrm{\Gamma }𝒞`$.
The parallel transport $`𝒫^H`$, for a given graph ($`1`$-dimensional, simplicial complex) $`H`$, is a specialization of the parallel transport $`𝒫^L`$ introduced in Section 3.2. For example if $`\stackrel{}{e_1e_2}`$ is the perspectivity between adjacent edges $`e_1=\overline{x_0x_1}`$ and $`e_2=\overline{x_1x_2}`$ in $`G`$, and if $`\eta :\{x_1,x_2\}2^{V(H)}\{\mathrm{}\}`$ is a cell in $`_{\overline{x_1x_2}}^H=Hom(G_{\overline{x_1x_2}},H)`$, then $`\eta ^{}:=𝒫^H(\stackrel{}{e_1e_2})(\eta ):\{x_0,x_1\}2^{V(H)}\{\mathrm{}\}`$ is defined by
$$\eta ^{}(x_0):=\eta (x_2)\text{ and }\eta ^{}(x_1):=\eta (x_1).$$
Fundamental observation: The construction of the connections $`𝒫^L`$, respectively $`𝒫^H`$, are quite natural and elementary but it is Proposition 4, respectively its more general relative Proposition 17, that serve as an actual justification for the introduction of these objects. Proposition 4 allows us to analyze the parallel transport of homotopy types of maps from the complex $`Hom(G,H)`$ to complexes $`Hom(G_e,H)`$, where $`eE(G)`$, providing a key for a resolution of the Lovász conjecture in the case when $`k`$ is an odd integer.
Implicit in the proof of Proposition 4 is the theory of folds of graphs and the analysis of natural morphisms between graph complexes $`Hom(T,H)`$, where $`T`$ is a tree, as developed in . This theory is one of essential ingredients in the Babson and Kozlov spectral sequence approach to the solution of Lovász conjecture. Some of these results are summarized in Proposition 3, in the form suitable for application to Proposition 4.
As usual $`L_m`$ is the graph-chain of vertex-length $`m`$, while $`L_{x_1\mathrm{}x_m}`$ is the graph isomorphic to $`L_m`$ defined on a linearly ordered set of vertices $`x_1,\mathrm{},x_m`$. In this context the “flip” is a generic name for the automorphism $`\sigma :L_{x_1\mathrm{}x_m}L_{x_1\mathrm{}x_m}`$ of the graph-chain such that $`\sigma (x_j)=x_{mj+1}`$ for each $`j`$.
###### Proposition 3
Suppose that $`e_1=\overline{x_0x_1}`$ and $`e_2=\overline{x_1x_2}`$ are two distinct, adjacent edges in the graph $`G`$. Let $`\sigma :L_{x_0x_1x_2}L_{x_0x_1x_2}`$ be the flip automorphism of $`L_{x_0x_1x_2}`$ and $`\widehat{\sigma }`$ the associated auto-homeomorphism of $`Hom(L_{x_0x_1x_2},H)`$. Suppose that $`\gamma _{ij}:L_{x_ix_j}L_{x_0x_1x_2}`$ is an obvious embedding and $`\widehat{\gamma }_{ij}`$ the associated maps of graph complexes. Then,
1. the induced map $`\widehat{\sigma }:Hom(L_{x_0x_1x_2},H)Hom(L_{x_0x_1x_2},H)`$ is homotopic to the identity map $`I`$, and
2. the diagram
$$\begin{array}{ccc}Hom(L_{x_0x_1x_2},H)& \stackrel{=}{}& Hom(L_{x_0x_1x_2},H)\\ \widehat{\gamma }_{01}& & \widehat{\gamma }_{12}& & \\ Hom(L_{x_0x_1},H)& \underset{𝒫^H(\stackrel{}{e_1e_2})}{}& Hom(L_{x_1x_2},H)\end{array}$$
is commutative up to homotopy.
Proof: Both statements are corollaries of Babson and Kozlov analysis of complexes $`Hom(T,H)`$, where $`T`$ is a tree, and morphisms $`\widehat{e}:Hom(T,H)Hom(T^{},H)`$, where $`T^{}`$ is a subtree of $`T`$ and $`e:T^{}T`$ the associated embedding.
Our starting point is an observation that both $`L_{x_0x_1}`$ and $`L_{x_1x_2}`$ are retracts of the graph $`L_{x_0x_1x_2}`$ in the category of graphs and graph homomorphisms. The retraction homomorphisms $`\varphi _{ij}:L_{x_0x_1x_2}L_{x_ix_j}`$, where $`\varphi _{01}(x_0)=x_0,\varphi _{01}(x_1)=x_1,\varphi _{01}(x_2)=x_0`$ and $`\varphi _{12}(x_0)=x_2,\varphi _{12}(x_1)=x_1,\varphi _{12}(x_2)=x_2`$ are examples of foldings of graphs. By the general theory , the maps $`\widehat{\gamma }_{ij}:Hom(L_{x_0x_1x_2},H)Hom(L_{x_ix_j},H)`$ and $`\widehat{\varphi }_{ij}:Hom(L_{x_ix_j},H)Hom(L_{x_0x_1x_2},H)`$ are homotopy equivalences. Actually $`\widehat{\gamma }_{ij}`$ is a deformation retraction and $`\widehat{\varphi }_{ij}`$ is the associated embedding such that $`\widehat{\gamma }_{ij}\widehat{\varphi }_{ij}=I`$ is the identity map.
The part (a) of the proposition is an immediate consequence of the fact that $`\varphi _{01}\sigma \gamma _{01}:L_{x_0x_1}L_{x_0x_1}`$ is an identity map. It follows that $`\widehat{\gamma _{01}}\widehat{\sigma }\widehat{\varphi }_{01}=I`$, and in light of the fact that $`\widehat{\gamma _{01}}`$ and $`\widehat{\varphi }_{01}`$ are homotopy inverses to each other, we conclude that $`\widehat{\sigma }I`$.
For the part (b) we begin by an observation that $`\varphi _{12}\sigma \gamma _{01}=\stackrel{}{e_1e_2}`$. Then, $`𝒫^H(\stackrel{}{e_1e_2})=\widehat{\gamma }_{01}\widehat{\sigma }\widehat{\varphi }_{12}`$, and as a consequence of $`\widehat{\sigma }I`$ and the fact that $`\widehat{\varphi }_{12}\widehat{\gamma }_{12}I`$, we conclude that
$$𝒫^H(\stackrel{}{e_1e_2})\widehat{\gamma }_{12}=\widehat{\gamma }_{01}\widehat{\sigma }\widehat{\varphi }_{12}\widehat{\gamma }_{12}\widehat{\gamma }_{01}.$$
$`\mathrm{}`$
###### Proposition 4
Suppose that $`x_0,x_1,x_2`$ are distinct vertices in $`G`$ such that $`\overline{x_0x_1},\overline{x_1x_2}E(G)`$. Let $`\alpha _{ij}:G_{x_ix_j}G`$ be the inclusion map of graphs and $`\widehat{\alpha }_{ij}`$ the associated map of $`Hom(,H)`$ complexes. Then the following diagram commutes up to a homotopy,
$$\begin{array}{ccc}Hom(G,H)& \stackrel{=}{}& Hom(G,H)\\ \widehat{\alpha }_{01}& & \widehat{\alpha }_{12}& & \\ Hom(G_{x_0x_1},H)& \underset{𝒫^H(\stackrel{}{e_1e_2})}{}& Hom(G_{x_1x_2},H)\end{array}$$
(7)
Proof: The diagram (7) can be factored as
$$\begin{array}{ccc}Hom(G,H)& \stackrel{=}{}& Hom(G,H)\\ \widehat{\beta }& & \widehat{\beta }& & \\ Hom(G_{x_0x_1x_2},H)& \stackrel{=}{}& Hom(G_{x_0x_1x_2},H)\\ \widehat{\gamma }_{01}& & \widehat{\gamma }_{12}& & \\ Hom(G_{x_0x_1},H)& \underset{𝒫^H(\stackrel{}{e_1e_2})}{}& Hom(G_{x_1x_2},H)\end{array}$$
(8)
where $`\beta `$ and $`\gamma _{ij}`$ are obvious inclusions of indicated graphs such that $`\alpha _{ij}=\beta \gamma _{ij}`$. Then the result is a direct consequence of Proposition 3, part (b). $`\mathrm{}`$
## 4 Lovász-Babson-Kozlov result for odd $k$
The proof of Lovász conjecture splits into two main branches, corresponding to the parity of a parameter $`n`$, where $`n`$ is an integer which enters the stage as the size of the vertex set of the complete graph $`K_n`$.
The first branch relies on Theorem 2.3. (loc. cit.), more precisely on part (b) of this result, while the second branch is founded on Theorem 2.6. Both theorems are about the topology of the graph complex $`Hom(C_{2r+1},K_n)`$. Theorem 2.3. (b) is a statement about the height of the first Stiefel-Whitney class, or equivalently the Conner-Floyd $`_2`$-index of the $`_2`$-space $`Hom(C_{2r+1},K_n)`$. Theorem 2.6. claims that for $`n`$ even, $`2\iota _{K_n}^{}`$ is a zero homomorphism where
$$\iota _{K_n}^{}:\stackrel{~}{H}^{}(Hom(K_2,K_n);)\stackrel{~}{H}^{}(Hom(C_{2r+1},K_n);)$$
(9)
is the homomorphism associated to the continuous map $`\iota _{K_n}:Hom(C_{2r+1},K_n)Hom(K_2,K_n)`$, which in turn comes from the inclusion $`K_2C_{2r+1}`$.
The central idea of our paper is an observation that Theorem 2.6. can be incorporated into a more general scheme, involving the “parallel transport” of graph complexes over graphs.
###### Theorem 5
Suppose that $`\alpha :K_2C_{2r+1}`$ is an inclusion map, $`\beta :K_2K_2`$ a nontrivial automorphism of $`K_2`$, and
$$\widehat{\alpha }:Hom(C_{2r+1},H)Hom(K_2,H),\widehat{\beta }:Hom(K_2,H)Hom(K_2,H)$$
the associated maps of graph complexes. Then the following diagram is commutative up to a homotopy
$$\begin{array}{ccc}Hom(C_{2r+1},H)& \stackrel{=}{}& Hom(C_{2r+1},H)\\ \widehat{\alpha }& & \widehat{\alpha }& & \\ Hom(K_2,H)& \stackrel{\widehat{\beta }}{}& Hom(K_2,H)\end{array}$$
(10)
Proof: Assume that the consecutive vertices of $`G=C_{2r+1}`$ are $`x_0,x_1,\mathrm{},x_{2r}`$ and let $`e_i=\overline{x_{i1}x_i}`$ be the associated sequence of edges where by convention $`e_{2r+1}=\overline{x_{2r}x_0}`$. Identify the graph $`K_2`$ to the subgraph $`G_{x_0x_1}`$ of $`G=C_{2r+1}`$.
By iterating Proposition 4 we observe that the diagram
$$\begin{array}{ccc}Hom(C_{2r+1},H)& \stackrel{=}{}& Hom(C_{2r+1},H)\\ \widehat{\alpha }& & \widehat{\alpha }& & \\ Hom(G_{x_0x_1},H)& \underset{𝒫^H(𝔭)}{}& Hom(G_{x_0x_1},H)\end{array}$$
(11)
is commutative up to a homotopy, where $`𝔭=\stackrel{}{e_1e_2}\mathrm{}\stackrel{}{e_{2r+1}e_1}`$. The proof is completed by the observation that $`𝔭=\beta `$ in the groupoid $`𝒢^P(G)`$. $`\mathrm{}`$
Theorem 2.6. from , the key for the proof of Lovász conjecture for odd $`k`$, is an immediate consequence of Theorem 5.
###### Corollary 6
(, T.2.6.) If $`n`$ is even then $`2\iota _{K_n}^{}`$ is a $`0`$-map where $`\iota _{K_n}^{}`$ is the map described in line (9).
Proof: It is sufficient to observe that for $`H=K_n`$, the complex $`Hom(K_2,K_n)S^{n2}`$ is an even dimensional sphere such that the automorphism $`\widehat{\beta }`$ from the diagram (10) is essentially an antipodal map. It follows that $`\widehat{\beta }`$ changes the orientation of $`Hom(K_2,K_n)`$ and as a consequence $`\iota _{K_n}^{}=\iota _{K_n}^{}`$. $`\mathrm{}`$
## 5 Generalizations and ramifications
In this section we extend the results from Section 3.3 to the case of simplicial complexes. This generalization is based on the following basic principles.
Graphs are viewed as $`1`$-dimensional simplicial complexes. Graph homomorphisms are special cases of non-degenerated simplicial maps of simplicial complexes, . The definition of $`Hom(G,H)`$ is extended to the case of $`Hom`$-complexes $`Hom(K,L)`$ of simplicial complexes $`K`$ and $`L`$. The groupoids needed for the definition of the parallel transport of $`Hom`$-complexes are already introduced by Joswig in , see Section 3.2 for a summary. Theory of folds for graph complexes is extended in Section 5.4 to the case of $`Hom`$-complexes in sufficient generality to allow “parallel transport” of homotopy types of maps between graph complexes. This development eventually leads to Theorem 22 which extends Theorem 5 to the case of $`Hom`$-complexes $`Hom(K,L)`$ and represents the currently final stage in the evolution of Theorem 2.6. from .
### 5.1 From $`Hom(G,H)`$ to $`Hom(K,L)`$
Suppose that $`K2^{V(K)}`$ and $`L2^{V(L)}`$ are two (finite) simplicial complexes, on the sets of vertices $`V(K)`$ and $`V(L)`$ respectively.
###### Definition 7
A simplicial map $`f:KL`$ is non-degenerated if it is injective on simplices. The set of all non-degenerated simplicial maps from $`K`$ to $`L`$ is denoted by $`Hom_0(K,L)`$.
###### Definition 8
$`Hom(K,L)`$ is a cell complex with the cells indexed by the functions $`\eta :V(K)2^{V(L)}\{\mathrm{}\}`$ such that
1. for each two vertices $`uv`$, if $`\{u,v\}K`$ then $`\eta (u)\eta (v)=\mathrm{}`$,
2. for each simplex $`\sigma K`$, the join $`_{vV(\sigma )}\eta (v)\mathrm{\Delta }^{V(L)}`$ of all sets ($`0`$-dimensional complexes) $`\eta (v)`$, where $`v`$ is a vertex of $`\sigma `$, is a subcomplex of $`L`$.
More precisely, each function $`\eta `$ satisfying conditions (1) and (2) defines a cell $`c_\eta :=_{vV(K)}\mathrm{\Delta }^{\eta (v)}`$ in $`Hom(K,L)_{vV(K)}\mathrm{\Delta }^{V(L)}`$ where by definition $`\mathrm{\Delta }^S`$ is an (abstract) simplex spanned by vertices in $`S`$.
We have already used in Section 3.2 the fact that if $`K=\mathrm{\Delta }^{[m]}`$ is a $`(m1)`$-dimensional simplex spanned by $`[m]`$ as the set of vertices, then $`Hom(\mathrm{\Delta }^{[m]},L)L_\mathrm{\Delta }^m`$ is the deleted product of $`L`$ . The following example shows that $`Hom(G,H)`$ is a special case of $`Hom(K,L)`$.
###### Example 9
The definition of $`Hom(K,L)`$ is a natural extension of $`Hom(G,H)`$ and reduces to it if $`K`$ and $`L`$ are $`1`$-dimensional complexes. Moreover,
$$Hom(G,H)Hom(Clique(G),Clique(H))$$
where $`Clique(\mathrm{\Gamma })`$ is the simplicial complex of all cliques in a graph $`\mathrm{\Gamma }`$.
###### Remark 10
The set $`Hom_0(K,L)`$ is easily identified as the $`0`$-dimensional skeleton of the cell-complex $`Hom(K,L)`$. Moreover, the reader familiar with can easily check that $`Hom(K,L)`$ is determined by the family $`M=Hom_0(K,L)`$ in the sense of Definition 2.2.1. from that paper.
### 5.2 Functoriality of $`Hom(K,L)`$
The construction of $`Hom(K,L)`$ is functorial in the sense that if $`f:KK^{}`$ is a non-degenerated simplicial map of complexes $`K`$ and $`K^{}`$, then there is an associated continuous map $`\widehat{f}:Hom(K^{},L)Hom(K,L)`$ of $`Hom`$-complexes. Indeed, if $`\eta :V(K^{})2^{V(L)}\{\mathrm{}\}`$ is a multi-valued function indexing a cell in $`Hom(K^{},L)`$, then it is not difficult to check that $`\eta f:V(K)2^{V(L)}\{\mathrm{}\}`$ is a cell in $`Hom(K,L)`$.
Perhaps even more important is the functoriality of $`Hom(K,L)`$ with respect to the second variable since this implies the functoriality of the bundle $`_k^L`$.
###### Proposition 11
Suppose that $`g:LL^{}`$ is a non-degenerated, simplicial map of simplicial complexes $`L`$ and $`L^{}`$. Then there exists an associated map
$$\widehat{g}:Hom(K,L)Hom(K,L^{}).$$
Proof: Assume that $`\eta :V(K)2^{V(L)\{\mathrm{}\}}`$ is a cell in $`Hom(K,L)`$. Then $`g\eta :V(K)2^{V(L^{})\{\mathrm{}\}}`$ is a cell in $`Hom(K,L^{})`$. Suppose $`u`$ and $`v`$ are distinct vertices in $`V(K)`$. By assumption $`\eta (u)\eta (v)=\mathrm{}`$. We deduce from here that $`g(\eta (u))g(\eta (v))\mathrm{}`$, otherwise $`g`$ would be a degenerated simplicial map.
The second condition from Definition 8 is checked by a similar argument. $`\mathrm{}`$
### 5.3 Chromatic number $`\chi (K)`$ and its relatives
The chromatic number $`\chi (K)`$ of a simplicial complex $`K`$ is
$$inf\{mHom_0(K,\mathrm{\Delta }^{[m]})\mathrm{}\}.$$
In other words $`\chi (K)`$ is the minimum number $`m`$ such that there exists a non-degenerated simplicial map $`f:K\mathrm{\Delta }^{[m]}`$. It is not difficult to check that $`\chi (K)=\chi (G_K)`$ where $`G_K=(K^{(0)},K^{(1)})`$ is the vertex-edge graph of the complex $`K`$. In particular $`\chi (K)`$ reduces to the usual chromatic number if $`K`$ is a graph, that is if $`K`$ is a $`1`$-dimensional simplicial complex.
Aside from the usual chromatic number $`\chi (G)`$, there are many related colorful graph invariants . Among the best known are the fractional chromatic number $`\chi _f(G)`$ and the circular chromatic number $`\chi _c(G)`$ of $`G`$. These and other related invariants are conveniently defined in terms of graph homomorphisms into graphs chosen from a suitable family $`=\{G_i\}_{iI}`$ of test graphs. Motivated by this, we offer an extension of the chromatic number $`\chi (K)`$ in hope that some genuine invariants of simplicial complexes objects arise this way.
###### Definition 12
Suppose that $`=\{T_iiI\}`$ is a family of “test” simplicial complexes and let $`\varphi :I`$ is a real-valued function. A $`T_i`$-coloring of $`K`$ is just a non-degenerated simplicial map $`f:KT_i`$ and $`\chi _{(,\varphi )}(K)`$, the $`(,\varphi )`$-chromatic number of $`K`$, is defined as the infimum of all weights $`\varphi (i)`$ over all $`T_i`$-colorings,
$$\chi _{(,\varphi )}(K):=inf\{\varphi (i)Hom_0(K,T_i)\mathrm{}\}.$$
### 5.4 Tree-like simplicial complexes
The tree-like or vertex collapsible complexes are intended to play in the theory of $`Hom(K,L)`$-complexes the role similar to the role of trees in the theory of graph complexes $`Hom(G,H)`$.
A pure, $`d`$-dimensional simplicial complex $`K`$ is shellable , if there is a linear order $`F_1,F_2,\mathrm{},F_m`$ on the set of its facets, such that for each $`j2`$, the complex $`F_j(_{i<j}F_i)`$ is a pure $`(d1)`$-dimensional subcomplex of the simplex $`F_j`$. The restriction $`R_j`$ of the facet $`F_j`$ is the minimal new face added to the complex $`_{k<j}F_k`$ by the addition of the facet $`F_j`$. Let $`r_j:=\mathrm{dim}(R_j)\{0,1,\mathrm{},d\}`$ be the type of the facet $`F_j`$. If $`r_jd`$ for each $`j`$ then the complex $`K`$ is collapsible. The collapsing process is just the shelling order read in the opposite direction. From this point of view, $`R_j`$ can be described as a free face in the complex $`_{ij}F_i`$, and the process of removing all faces $`F`$ such that $`R_jFF_j`$ is called an elementary $`r_j`$-collapse.
###### Definition 13
A pure $`d`$-dimensional simplicial complex $`K`$ is called tree-like or vertex collapsible if it is collapsible to a $`d`$-simplex with the use of elementary $`0`$-collapses alone. In other words $`K`$ is shellable and for each $`j2`$, the intersection $`F_j(_{i<j}F_i)`$ is a proper face of $`F_j`$.
In order to establish analogs of Propositions 3 and 4 for complexes $`Hom(K,L)`$, we prove a result which shows that elementary vertex collapsing provides a good substitute and a partial generalization for the concept of “foldings” of graphs used in in the theory of graph complexes $`Hom(G,H)`$.
###### Proposition 14
Suppose that the simplicial complex $`K^{}`$ is obtained from $`K`$ by an elementary vertex collapse. In other words we assume that $`K=\sigma K^{},\sigma K^{}=\sigma ^{}`$, where $`\sigma `$ is a simplex in $`K`$ and $`\sigma ^{}`$ a facet of $`\sigma `$. Assume that $`\sigma ^{}`$ is not maximal in $`K^{}`$, i.e. that for some simplex $`\sigma ^{\prime \prime }K^{}`$ and a vertex $`u\sigma ^{\prime \prime },\sigma ^{}=\sigma ^{\prime \prime }\{u\}`$. Then for any simplicial complex $`L`$, the inclusion map $`\gamma :K^{}K`$ induces a homotopy equivalence
$$\widehat{\gamma }:Hom(K,L)Hom(K^{},L).$$
Proof: Let $`\{v\}=\sigma \sigma ^{}`$. Aside from the inclusion map $`\gamma :K^{}K`$, there is a retraction (folding) map $`\rho :KK^{}`$, where $`\rho (v)=u`$ and $`\rho |_K^{}=I_K^{}`$. Since $`\rho \gamma =I_K^{}`$, we observe that $`\widehat{\gamma }\widehat{\rho }=Id_K^{}`$ is the identity map on $`Hom(K^{},L)`$, i.e. the complex $`Hom(K^{},L)`$ is a retract of the complex $`Hom(K,L)`$. It remains to be shown that $`\widehat{\rho }\widehat{\gamma }Id_K`$ is homotopic to the identity map on $`Hom(K,L)`$.
Note that if $`\eta Hom(K,L)`$ then $`\eta ^{}:=\widehat{\rho }\widehat{\gamma }(\eta )`$ is the function defined by
$$\eta ^{}(w)=\{\begin{array}{cc}\hfill \eta (w),& \text{ if }wv\hfill \\ \hfill \eta (u),& \text{ if }w=v.\hfill \end{array}$$
Let $`\omega :Hom(K,L)Hom(K,L)`$ be the map defined by
$$\omega (\eta )(w)=\{\begin{array}{cc}\eta (w),\hfill & \text{ if }wv\hfill \\ \eta (u)\eta (v),\hfill & \text{ if }w=v.\hfill \end{array}$$
Note that $`\omega `$ is well defined since if a vertex $`x`$ is adjacent to $`v`$ it is also adjacent to $`u`$, hence the condition $`\omega (\eta )(v)\eta (x)=\mathrm{}`$ is a consequence of $`\eta (u)\eta (x)=\mathrm{}=\eta (v)\eta (x)`$.
Since for each $`\eta Hom(K,L)`$ and each vertex $`xK`$,
$$\eta (x)\omega (\eta )(x)\widehat{\rho }\widehat{\gamma }(\eta )(x),$$
by the Order Homotopy Theorem all three maps $`Id_K,\omega `$ and $`\widehat{\rho }\widehat{\gamma }`$ are homotopic. This completes the proof of the proposition. $`\mathrm{}`$
###### Corollary 15
If $`T`$ is a $`d`$-dimensional, tree-like simplicial complex than $`Hom(T,L)`$ has the same homotopy type as the deleted join $`Hom(\mathrm{\Delta }^d,L)=L_\mathrm{\Delta }^{d+1}`$.
### 5.5 Parallel transport of homotopy types of maps
As in the case of graph complexes, the real justification for the introduction of the parallel transport of $`Hom`$-complexes comes from the fact that it preserves the homotopy type of the maps $`Hom(K,L)Hom(\sigma ,L)`$. As in Section 3.3, as a preliminary step we prove an analogue of Proposition 3.
###### Proposition 16
Suppose that $`\sigma _1`$ and $`\sigma _2`$ are two distinct, adjacent $`k`$-dimensional simplices in a finite simplicial complex $`K`$ which share a common $`(k1)`$-dimensional simplex $`\tau `$. Let $`\mathrm{\Sigma }=\sigma _1\sigma _2`$. Let $`\alpha :\mathrm{\Sigma }\mathrm{\Sigma }`$ be the automorphism of $`\mathrm{\Sigma }`$ which interchanges simplices $`\sigma _1`$ and $`\sigma _2`$ keeping the common face $`\tau `$ point-wise fixed.
Suppose that $`\gamma _i:\sigma _i\mathrm{\Sigma }`$ is an obvious embedding and $`\widehat{\gamma }_i`$ the associated maps of $`Hom`$-complexes. Then,
1. the induced map $`\widehat{\alpha }:Hom(\mathrm{\Sigma },L)Hom(\mathrm{\Sigma },L)`$ is homotopic to the identity map $`I_\mathrm{\Sigma }`$, and
2. the diagram
$$\begin{array}{ccc}Hom(\mathrm{\Sigma },L)& \stackrel{=}{}& Hom(\mathrm{\Sigma },L)\\ \widehat{\gamma }_1& & \widehat{\gamma }_2& & \\ Hom(\sigma _1,L)& \underset{𝒫^H(\stackrel{}{\sigma _1\sigma _2})}{}& Hom(\sigma _2,L)\end{array}$$
is commutative up to homotopy.
Proof: By Proposition 14, both maps $`\widehat{\gamma }_i:Hom(\mathrm{\Sigma },L)Hom(\sigma _i,L)`$ for $`i=1,2`$ are homotopy equivalences. Let $`\rho _1:\mathrm{\Sigma }\sigma _1`$ and $`\rho _2:\mathrm{\Sigma }\sigma _2`$ be the folding maps. Then $`\rho _i\gamma _i=I_{\sigma _i},`$ $`\widehat{\gamma _i}\widehat{\rho _i}=I`$ and we conclude that $`\widehat{\rho }_i:Hom(\sigma _i,L)Hom(\mathrm{\Sigma },L)`$ is also a homotopy equivalence.
Part (a) of the proposition follows from the fact that $`\rho _1\alpha \gamma _1=I_{\sigma _1}`$ is an identity map. Indeed, an immediate consequence is that $`\widehat{\gamma _1}\widehat{\alpha }\widehat{\rho _1}=I:Hom(\sigma _1,L)Hom(\sigma _1,L)`$ is also an identity map and, in light of the fact that $`\widehat{\gamma _1}`$ and $`\widehat{\rho _1}`$ are homotopy inverses to each other, we deduce that $`\widehat{\alpha }I`$.
For the part (b) we begin by an observation that $`\rho _2\alpha \gamma _1=\stackrel{}{\sigma _1\sigma _2}`$. Then, $`𝒫^H(\stackrel{}{\sigma _1\sigma _2})=\widehat{\gamma }_1\widehat{\alpha }\widehat{\rho }_2`$, and as a consequence of $`\widehat{\alpha }I`$ and the fact that $`\widehat{\rho }_2\widehat{\gamma }_2I`$, we conclude that
$$𝒫^H(\stackrel{}{\sigma _1\sigma _2})\widehat{\gamma }_2=\widehat{\gamma }_1\widehat{\alpha }\widehat{\rho }_2\widehat{\gamma }_2\widehat{\gamma }_1.$$
$`\mathrm{}`$
###### Proposition 17
Suppose that $`K`$ and $`L`$ are finite simplicial complexes and $`\sigma _1,\sigma _2`$ a pair of adjacent (distinct), $`k`$-dimensional simplices in $`K`$. Let $`\alpha _i:\sigma _iK`$ be the embedding of $`\sigma _i`$ in $`K`$ and $`\widehat{\alpha }_i:Hom(K,L)Hom(\sigma _i,L)`$ the associated map of $`Hom`$-complexes. Then the following diagram commutes up to a homotopy.
$$\begin{array}{ccc}Hom(K,L)& \stackrel{=}{}& Hom(K,L)\\ \widehat{\alpha }_1& & \widehat{\alpha }_2& & \\ Hom(\sigma _1,L)& \underset{𝒫^L(\stackrel{}{e_1e_2})}{}& Hom(\sigma _2,H)\end{array}$$
(12)
Proof: Let $`\mathrm{\Sigma }:=\sigma _1\sigma _2,\tau :=\sigma _1\sigma _2`$. Then $`\alpha _i=\beta \gamma _i`$ where $`\beta :\mathrm{\Sigma }K`$ and $`\gamma _i:\sigma _i\mathrm{\Sigma }`$ are natural embeddings of complexes.
The diagram (12) can be factored as
$$\begin{array}{ccc}Hom(K,L)& \stackrel{=}{}& Hom(K,L)\\ \widehat{\beta }& & \widehat{\beta }& & \\ Hom(\mathrm{\Sigma },L)& \stackrel{=}{}& Hom(\mathrm{\Sigma },L)\\ \widehat{\gamma }_1& & \widehat{\gamma }_2& & \\ Hom(\sigma _1,L)& \underset{𝒫^L(\stackrel{}{\sigma _1\sigma _2})}{}& Hom(\sigma _2,L)\end{array}$$
(13)
Then the result is a direct consequence of Proposition 16, part (b). $`\mathrm{}`$
###### Corollary 18
Suppose that $`K`$ and $`L`$ are finite simplicial complexes, $`\sigma `$ a $`k`$-dimensional simplex in $`K`$ and $`\alpha :\sigma K`$ the associated embedding. Let $`\tau \mathrm{\Pi }(K,\sigma )`$. Then the following diagram commutes up to a homotopy.
$$\begin{array}{ccc}Hom(K,L)& \stackrel{=}{}& Hom(K,L)\\ \widehat{\alpha }& & \widehat{\alpha }& & \\ Hom(\sigma ,L)& \underset{\widehat{\tau }}{}& Hom(\sigma ,L)\end{array}$$
(14)
## 6 Main results
In this section we prove the promised extension of the Lovász-Babson-Kozlov theorem. The graphs are replaced by pure $`d`$-dimensional simplicial complexes, while the role of the odd cycle $`C_{2r+1}`$ is played by a complex $`\mathrm{\Gamma }`$ which has some special symmetry properties in the sense of the following definition.
As usual, an involution $`\omega :XX`$ is the same as a $`_2`$-action on $`X`$. An involution on a simplicial complex $`\mathrm{\Gamma }`$ induces an involution on the complex $`Hom(\mathrm{\Gamma },L)`$ for each simplicial complex $`L`$. For all other standard facts and definitions related to $`_2`$-complexes, the reader is referred to .
###### Definition 19
A pure $`d`$-dimensional simplicial complex $`\mathrm{\Gamma }`$ is a $`\mathrm{\Phi }_d`$-complex if it is a $`_2`$-complex with an invariant $`d`$-simplex $`\sigma =\{v_0,v_1,\mathrm{},v_d\}`$ such that the restriction $`\tau :=\omega |_\sigma `$ of the involution $`\omega :\mathrm{\Gamma }\mathrm{\Gamma }`$ on $`\sigma `$ is a non-trivial element of the group $`\mathrm{\Pi }(\mathrm{\Gamma },\sigma )`$.
###### Remark 20
By definition, if $`\mathrm{\Gamma }`$ is a $`\mathrm{\Phi }_d`$-complex then the inclusion map $`\alpha :\sigma \mathrm{\Gamma }`$ is $`_2`$-equivariant, so the associated map $`\widehat{\alpha }:Hom(\mathrm{\Gamma },K)Hom(\sigma ,K)`$ is also $`_2`$-equivariant for each complex $`K`$.
###### Example 21
The graph $`C_{2r+1}`$ is obviously an example of a $`\mathrm{\Phi }_1`$-complex. Figure 1 displays four examples of $`\mathrm{\Phi }_2`$-complexes, initial elements of two infinite series $`_\mu `$ and $`\mathrm{\Sigma }_\nu ,\mu ,\nu `$. The complexes $`_1`$ and $`_2`$ etc. are obtained from two triangulated annuli, glued together along a common triangle $`\sigma `$. Similarly, the complexes $`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2,\mathrm{}`$, are obtained by gluing together two triangulated Möbius strips. The associated group of projectivities are $`\mathrm{\Pi }(_\mu ,\sigma )=S_3`$ and $`\mathrm{\Pi }(\mathrm{\Sigma }_\nu ,\sigma )=_2`$.
###### Theorem 22
Suppose that $`\mathrm{\Gamma }`$ is a $`\mathrm{\Phi }_d`$-complex in the sense of Definition 19, with an associated invariant simplex $`\sigma =\{v_0,v_1,\mathrm{},v_d\}`$. Suppose that $`K`$ is a pure $`d`$-dimensional simplicial complex. Than for $`m`$ even,
$$Coind__2(Hom(\mathrm{\Gamma },K))m\chi (K)m+d+2.$$
(15)
Proof: By definition $`Coind__2(Hom(\mathrm{\Gamma },K))m`$ means that there exists a $`_2`$-equivariant map $`\mu :S^mHom(\mathrm{\Gamma },K)`$. Assume that $`\chi (K)m+d+1`$ which means that there exists a non-degenerated simplicial map $`\varphi :K\mathrm{\Delta }^{[m+d+1]}`$. By functoriality of the construction of $`Hom`$-complexes, Section 5.2, there is an induced $`_2`$-equivariant map $`\widehat{\varphi }:Hom(\mathrm{\Gamma },K)Hom(\mathrm{\Gamma },\mathrm{\Delta }^{[m+d+1]})`$ and similarly a map $`\widehat{\alpha }:Hom(\mathrm{\Gamma },\mathrm{\Delta }^{[m+d+1]})Hom(\sigma ,\mathrm{\Delta }^{[m+d+1]})`$. By Theorem 3.3.3., the complex
$$Hom(\sigma ,\mathrm{\Delta }^{[m+d+1]})Hom(K_{d+1},K_{m+d+1})$$
is a wedge of $`m`$-dimensional spheres. Since $`Hom(\sigma ,\mathrm{\Delta }^{[m+d+1]})`$ is a free $`_2`$-complex, we deduce that there exists a $`_2`$-equivariant map $`Hom(\sigma ,\mathrm{\Delta }^{[m+d+1]})S^m`$. All these maps can be arranged in the following sequence of $`_2`$-equivariant maps
$$S^m\stackrel{\mu }{}Hom(\mathrm{\Gamma },K)\stackrel{\widehat{\varphi }}{}Hom(\mathrm{\Gamma },\mathrm{\Delta }^{[m+d+1]})\stackrel{\widehat{\alpha }}{}Hom(\sigma ,\mathrm{\Delta }^{[m+d+1]})\stackrel{\nu }{}S^m.$$
By Corollary 18, there is a homotopy equivalence $`\widehat{\alpha }\tau \widehat{\alpha }`$. This is in contradiction with Proposition 23, which completes the proof of the theorem. $`\mathrm{}`$
###### Proposition 23
Suppose that $`f:XY`$ is a $`_2`$-equivariant map of free $`_2`$-complexes $`X`$ and $`Y`$ where $`_2=\{1,\omega \}`$. Assume that $`Coind__2(X)mInd__2(Y)`$, where $`m`$ is an even integer. In other words our assumption is that there exist $`_2`$-equivariant maps $`\mu `$ and $`\nu `$ such that
$$S^m\stackrel{\mu }{}X\stackrel{f}{}Y\stackrel{\nu }{}S^m.$$
Then the maps $`f`$ and $`\omega f`$ are not homotopic.
Proof: If $`f\omega f:XY`$ then $`\nu f\mu \nu \omega f\mu :S^mS^m`$ and by the equivariance of $`\nu `$, $`\omega gg:S^mS^m`$ where $`g:=\nu f\mu `$. It follows that
$$\mathrm{deg}(g)=\mathrm{deg}(\omega )\mathrm{deg}(g)=\mathrm{deg}(\omega g)=\mathrm{deg}(g),$$
i.e. $`\mathrm{deg}(g)=0`$, which is in contradiction with a well known fact that a $`_2`$-equivariant map $`g:S^mS^m`$ of even dimensional spheres must have an odd degree. $`\mathrm{}`$
###### Corollary 24
Suppose that $`\mathrm{\Gamma }`$ is a $`\mathrm{\Phi }_d`$-complex with an associated invariant simplex $`\sigma =\{v_0,v_1,\mathrm{},v_d\}`$. Suppose that $`K`$ is a pure $`d`$-dimensional simplicial complex. Than for $`k`$ odd,
$$Hom(\mathrm{\Gamma },K)\text{ }\text{is }k\text{-connected}\text{ }\chi (K)k+d+3.$$
(16)
Proof: If $`Hom(\mathrm{\Gamma },K)`$ is $`k`$-connected then $`Coind__2(Hom(\mathrm{\Gamma },K))k+1`$, hence the implication (16) is an immediate consequence of Theorem 22. $`\mathrm{}`$
|
warning/0506/math0506337.html
|
ar5iv
|
text
|
# On the scaling limit of simple random walk excursion measure in the plane
## 1 Introduction
A number of mathematically simplistic lattice models, including the self-avoiding random walk , have been introduced in an attempt to better understand critical phenomena in two-dimensional statistical physics. While these models have been studied for several decades, little progress had been made until recently. The introduction of the Schramm-Loewner evolution, a new family of conformally invariant distributions on random curves, has led to a plethora of exciting results about the scaling limits of these models at criticality. For example, the scaling limits of loop-erased random walk , uniform spanning trees , and site percolation on the triangular lattice can now be described using SLE.
One of the first successes, however, of the SLE program was the determination of the intersection exponents for random walk and Brownian motion, and the establishment of Mandelbrot’s conjecture that the Hausdorff dimension of the frontier of the planar Brownian path is $`4/3`$. (See for a survey of this work.) The Brownian excursion measure, a conformally invariant infinite measure on curves which had been introduced in previous work by Lawler and Werner , figured prominently in the explicit calculations of the intersection exponents.
The goal of this present paper is to construct a discrete object, the simple random walk excursion measure, which has the Brownian excursion measure as its scaling limit. Of course, the convergence of simple random walk on $`^2`$ to Brownian motion in $``$ has been known since Donsker’s theorem of 1951. However, what had not been established was a strong version of this result which holds for random walk and Brownian motion on any simply connected domain where the errors do not depend on the smoothness of the boundary. By proving in the present paper that for any bounded, simply connected Jordan domain, the scaling limit of discrete excursion measure is Brownian excursion measure, we establish such a result.
### 1.1 Main results
We begin with a discussion of the main results, leaving some of the precise statements to later sections. Our concern will be exclusively two dimensional, so we will identify $`^2`$ in the usual way, and write any of $`w`$, $`x`$, $`y`$, or $`z`$ for points in $``$. A domain $`D`$ is an open and connected set; write $`𝔻:=\{z:|z|<1\}`$ for the open unit disk, and $`:=\{z:𝓂(𝓏)>\mathcal{0}\}`$ for the upper half plane. A standard complex Brownian motion will be denoted $`B_t`$, $`t0`$, and $`S_n`$, $`n=0,1,\mathrm{}`$, will denote two-dimensional simple random walk, both started at the origin unless otherwise noted. We will generally use $`T`$ for stopping times for Brownian motion and $`\tau `$ for stopping times for random walk, and write $`𝔼^x`$ and $`^x`$ for expectations and probabilities, respectively, assuming $`B_0=x`$ or $`S_0=x`$.
A subset $`A^2`$ is said to be simply connected if both $`A`$ and $`^2A`$ are non-empty and connected. Write the (outer) boundary of $`A`$ as $`A:=\{z^2A:\mathrm{dist}(z,A)=1\}`$. An excursion in $`A`$ is a path $`\omega :=[\omega _0,\omega _1,\mathrm{},\omega _k]`$ with $`|\omega _j\omega _{j1}|=1`$ for all $`j`$; $`\omega _0`$, $`\omega _kA`$; and $`w_1,\mathrm{},\omega _{k1}A`$. It is implicit that $`2k<\mathrm{}`$; the length of $`\omega `$ is $`|\omega |:=k`$. We can view excursions of length $`k`$ as curves $`\omega :[0,k]`$ by linear interpolation. Write $`𝒦_A`$ for the set of excursions in $`A`$, and define the simple random walk excursion measure as the measure on $`𝒦_A`$ which assigns measure $`4^k`$ to each length $`k`$ excursion in $`A`$. That is, the excursion measure of $`\omega =[\omega _0,\omega _1,\mathrm{},\omega _k]`$ is the probability that the first $`k`$ steps of a simple random walk starting at $`\omega _0`$ are the same as $`\omega `$. Let $`D`$ be a bounded simply connected domain containing the origin, and for each $`N<\mathrm{}`$, let $`D_N`$ denote the connected domain containing the origin of the set of $`z=u+iv\frac{1}{N}^2`$ such that $`\{u^{}+iv^{}:|uu^{}|(2N)^1,|vv^{}|(2N)^1\}`$ is contained in $`D`$. For each $`N`$, we get a measure on paths denoted $`\mu _{D_N}^{\mathrm{𝗋𝗐}}`$ by considering the random walk excursion measure on $`D_N`$, and scaling the excursions by Brownian scaling: $`\omega ^{(N)}(t):=N^{1/2}\omega (2Nt)`$. As $`N\mathrm{}`$, these measures converge to $`\mu _D`$, excursion measure on $`D`$, which is an infinite measure on paths. Since Brownian motion in $``$ is conformally invariant (up to a time-change), $`\mu _D`$ is also conformally invariant. (See Proposition 3.31.) If $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$ are disjoint arcs in $`D`$, then conditioning the excursion measure to have endpoints $`z\mathrm{\Gamma }`$, $`w\mathrm{{\rm Y}}`$, gives a probability measure on excursions from $`\mathrm{\Gamma }`$ to $`\mathrm{{\rm Y}}`$ in $`D`$.
The primary result of this paper is that for any bounded, simply connected Jordan domain $`D`$, simple random walk excursion measure converges to Brownian excursion measure on $`D`$.
###### Theorem 1.1.
If $`D`$ is a bounded, simply connected domain containing the origin with $`\mathrm{inrad}(D)=1`$, $`D`$ is Jordan, and $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$, then
$$\mathrm{}(\mathrm{\hspace{0.33em}4}\mu _{D_N}^{\mathrm{𝗋𝗐}},\mu _D)0$$
where $`\mathrm{}`$ denotes the Prohorov metric.
As we are discussing the convergence of infinite measures, we need to be a little careful about how we define convergence in the Prohorov metric (which is usually defined only for finite measures). As the restriction of excursion measure to disjoint boundary arcs gives a finite measure, Theorem 1.1 is to be interpreted as meaning that for *any* pair of disjoint boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$,
$$\mathrm{}(4\mu _{D_N}^{\mathrm{𝗋𝗐}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _D(\mathrm{\Gamma },\mathrm{{\rm Y}}))0$$
where $`\mathrm{\Gamma }_n`$, $`\mathrm{{\rm Y}}_N`$ are the “associated (discrete) boundary arcs in $`D_N`$.” In Section 4.5 we prove the precise formulation of Theorem 1.1.
Since a Brownian (resp., random walk) excursion can be viewed as consisting of a Brownian motion (resp., random walk) plus tails, the proof of convergence has two distinct parts—a “global part” plus a “local part.” The strong approximation of Komlós, Major, and Tusnády is used to couple random walk and Brownian motion in the interior of the domain away from the boundary. This global part does not depend on the smoothness of the boundary. The local part concerns the tails whose behaviour can be controlled using the Beurling estimates; here the structure of the boundary does come into play. The proof of convergence also employs an estimation of the discrete excursion Poisson kernel in terms of the excursion Poisson kernel derived in which was used in that paper to prove a conjecture of Fomin . Hence, by proving the weak convergence of excursion measures, we are extending the “central limit theorem” for the endpoints of the excursions proved in .
Technically, since $`\mu _{D_N}^{\mathrm{𝗋𝗐}}`$ is supported on continuous curves, we must associate to $`D_N`$ a domain in $``$ by identifying each point in $`D_N`$ with the square of side length $`1/N`$ centred at that point. It is important that these so-called “union of squares” domains $`\stackrel{~}{D}_N`$ converge to the original domain $`D`$. However, the convergence is not in the usual topological sense, but rather in the Carathéodory sense. This is captured by the following theorem which is carefully stated and proved in Section 4.3.
###### Theorem 1.2.
If $`f_N`$, $`f`$ are conformal transformations of the unit disk $`𝔻`$ onto $`\stackrel{~}{D}_N`$, $`D`$, respectively, with $`f_N(0)=f(0)=0`$ and $`f_N^{}(0)`$, $`f^{}(0)>0`$, then $`f_Nf`$ uniformly on compact subsets of $`𝔻`$. In other words, $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$.
### 1.2 Outline of the paper
In Section 2, we establish some notation, and recall some facts from complex analysis about conformal transformations. We also review the definitions and basic facts about Green’s functions on both $``$ and $`^2`$. Section 3 is devoted to a discussion of excursions and excursion measures. Included are some fundamental ideas about spaces of curves and measures on metric spaces. We also review the Prohorov topology, and prove several easy lemmas about the Prohorov metric which are needed in the sequel. The Poisson kernel and excursion Poisson kernel are then reviewed, with an emphasis on their conformal covariance properties, and a construction of excursion measure on $`D`$, differing from that in , is carried out. The final section, Section 4, is devoted to the proofs of Theorem 1.1 and Theorem 1.2. The material in Section 4.4 relies on the results obtain in . Instead of simply recopying those results as originally proved, we have translated them to statements in terms of $`D_N`$, the $`1/N`$-scale discrete approximation to $`D`$. A review of some recent strong approximation results is included in Section 4.5.1 because of their necessity in the proof of Theorem 1.1.
## 2 Background and notation
### 2.1 Simply connected subsets of $``$ and $`^2`$
A function $`f:DD^{}`$ is a conformal transformation if $`f`$ is an analytic, univalent (i.e, one-to-one) function that is onto $`D^{}`$. It follows that $`f^{}(z)0`$ for $`zD`$, and $`f^1:D^{}D`$ is also a conformal transformation. If $`D`$ is a domain in $``$, then a connected $`\mathrm{\Gamma }D`$ is an (open) analytic arc of $`D`$ if there is a domain $`E`$ that is symmetric about the real axis and a conformal transformation $`f:Ef(E)`$ such that $`f(E)=\mathrm{\Gamma }`$ and $`f(E)=f(E)D`$. We say that $`D`$ is locally analytic at $`xD`$ if there exists an analytic arc of $`D`$ containing $`x`$. For $`D`$ with $`0D`$, define the radius (with respect to the origin) of $`D`$ to be $`\mathrm{rad}(D):=sup\{|z|:zD\}`$, and the inradius (with respect to the origin) of $`D`$ to be $`\mathrm{inrad}(D):=\mathrm{dist}(0,D):=inf\{|z|:zD\}`$. The diameter of $`D`$ is $`\mathrm{diam}(D):=sup\{|xy|:x,yD\}`$. Call a bounded domain $`D`$ a Jordan domain if $`D`$ is a Jordan curve (i.e., homeomorphic to a circle). A Jordan domain is nice if the Jordan curve $`D`$ can be expressed as a finite union of analytic curves. Note that Jordan domains are necessarily simply connected. For each $`r>0`$, let $`𝒟^r`$ be the set of nice Jordan domains containing the origin of inradius $`r`$, and write $`𝒟:=_{r>0}𝒟^r`$. We also define $`𝒟^{}`$ to be the set of Jordan domains containing the origin, and note that $`𝒟𝒟^{}`$. If $`D`$, $`D^{}𝒟^{}`$, let $`𝒯(D,D^{})`$ be the set of all conformal transformations of $`D`$ onto $`D^{}`$. The Riemann mapping theorem implies that $`𝒯(D,D^{})\mathrm{}`$, and since $`D`$, $`D^{}`$ are Jordan, the Carathéodory extension theorem tells us that $`f𝒯(D,D^{})`$ can be extended to a homeomorphism of $`\overline{D}`$ onto $`\overline{D^{}}`$. The statements and details of these two theorems may be found in \[7, §1.5\].
There are three standard ways to define the boundary of a proper subset $`A`$ of $`^2`$:
* (outer) boundary: $`A:=\{y^2A:|yx|=1\text{ for some }xA\}`$;
* inner boundary: $`_iA:=(^2A)=\{xA:|yx|=1\text{ for some }y^2A\}`$;
* edge boundary: $`_eA:=\{(x,y):xA,y^2A,|xy|=1\}`$.
To each finite, connected $`A^2`$ we associate a domain $`\stackrel{~}{A}`$ in the following way. For each edge $`(x,y)_eA`$, considered as a line segment of length one, let $`\mathrm{}_{x,y}`$ be the perpendicular line segment of length one intersecting $`(x,y)`$ in the midpoint. Let $`\stackrel{~}{A}`$ denote the union of the line segments $`\mathrm{}_{x,y}`$, and let $`\stackrel{~}{A}`$ denote the domain with boundary $`\stackrel{~}{A}`$ containing $`A`$. Observe that
$$\stackrel{~}{A}\stackrel{~}{A}=\underset{xA}{}𝒮_x\text{ where }𝒮_x:=x+\left([1/2,1/2]\times [1/2,1/2]\right).$$
That is, $`𝒮_x`$ is the closed square of side length one centred at $`x`$ whose sides are parallel to the coordinate axes. Also, note that $`\stackrel{~}{A}`$ is a simply connected domain if and only if $`A`$ is a simply connected subset of $`^2`$. We say $`\stackrel{~}{A}`$ is the “union of squares” domain associated to $`A`$.
Let $`𝒜`$ denote the set of all finite simply connected subsets of $`^2`$ containing the origin. If $`A𝒜`$, let $`\mathrm{inrad}(A):=\mathrm{min}\{|z|:z^2A\}`$ and $`\mathrm{rad}(A):=\mathrm{max}\{|z|:zA\}`$ denote the inradius and radius (with respect to the origin), respectively, of $`A`$, and define $`𝒜^n`$ to be the set of $`A𝒜`$ with $`n\mathrm{inrad}(A)2n`$; thus $`𝒜:=_{n>0}𝒜^n`$. If $`A𝒜`$, $`0x_1_iA`$, and $`[x_1,x_2,\mathrm{},x_j]`$ is a nearest neighbour path in $`A\{0\}`$, then the connected component of $`A\{x_1,\mathrm{},x_j\}`$ containing the origin is simply connected.
Finally, if $`A𝒜`$ with associated domain $`\stackrel{~}{A}`$, then we write $`f_A:=f_{\stackrel{~}{A}}`$ for the conformal transformation of $`\stackrel{~}{A}`$ onto the unit disk $`𝔻`$ with $`f_A(0)=0`$, $`f_A^{}(0)>0`$.
### 2.2 Green’s functions on $``$ and $`^2`$
Let $`D`$ be a domain whose boundary contains a curve, and write $`g_D(x,y)`$ for the Green’s function for Brownian motion on $`D`$. If $`xD`$, we can define $`g_D(x,)`$ as the unique harmonic function on $`D\{x\}`$, vanishing on $`D`$, with $`g_D(x,y)=\mathrm{log}|xy|+O(1)`$ as $`|xy|0`$. Equivalently, if $`D𝒟^{}`$, then for distinct points $`x`$, $`yD`$, $`g_D(x,y)=𝔼^x[\mathrm{log}|B_{T_D}y|]\mathrm{log}|xy|`$ where $`T_D:=inf\{t:B_tD\}`$. In particular, if $`0D`$, then $`g_D(x)=𝔼^x[\mathrm{log}|B_{T_D}|]\mathrm{log}|x|`$ for $`xD`$ where $`g_D(x):=g_D(0,x)`$. For further details, consult \[17, Chapter 2\]. Since the Green’s function is a well-known example of a conformal invariant (see, e.g., \[7, §1.8\]), in order to determine $`g_D`$ for arbitrary $`D𝒟^{}`$, it is enough to find $`f_D𝒯(D,𝔻)`$. Conversely, suppose that $`D`$ is a simply connected domain containing the origin with Green’s function $`g_D`$. The unique conformal transformation of $`D`$ onto $`𝔻`$ with $`f_D(0)=0,f_D^{}(0)>0`$ can be written as
$$f_D(x)=\mathrm{exp}\{g_D(x)+i\theta _D(x)\}.$$
(1)
Note that $`g_D+i\theta _D`$ is analytic in $`D\{0\}`$. Suppose that $`A𝒜`$, and that $`g_A(x,y):=g_{\stackrel{~}{A}}(x,y)`$. As explained in , the exact form of the Green’s function gives
$$g_A(x,y)=\mathrm{log}\left|\frac{\overline{f_A(y)}f_A(x)1}{f_A(y)f_A(x)}\right|.$$
If we write $`\theta _A:=\theta _{\stackrel{~}{A}}`$, then (1) implies $`f_A(x)=\mathrm{exp}\{g_A(x)+i\theta _A(x)\}`$.
Let $`S_n`$ be a simple random walk on $`^2`$, and let $`A^2`$. If $`\tau _A:=\mathrm{min}\{j0:S_jA\}`$, then we let
$$G_A(x,y):=\underset{j=0}{\overset{\mathrm{}}{}}^x\{S_j=y,\tau _A>j\}$$
denote the Green’s function for random walk on $`A`$, and set $`G_A(x):=G_A(x,0)=G_A(0,x)`$. Write $`a`$ for the potential kernel for simple random walk defined by
$$a(x):=\underset{j=0}{\overset{\mathrm{}}{}}\left[^0\{S_j=0\}^x\{S_j=0\}\right].$$
It is known \[16, Theorem 1.6.2\] that as $`|x|\mathrm{}`$,
$$a(x)=\frac{2}{\pi }\mathrm{log}|x|+k_0+o(|x|^{3/2})$$
(2)
where $`k_0:=(2\varsigma +3\mathrm{ln}2)/\pi `$ and $`\varsigma `$ is Euler’s constant, and that $`G_A(x)=𝔼^x[a(S_{\tau _A})]a(x)`$ for $`xA`$. The error in (2) will suffice for our purposes even though stronger results are known; see . We also recall a uniform estimate for $`G_A(x)`$, and a relationship between the Green’s functions $`G_A`$ and $`g_A`$ which is proved in \[15, Theorem 1.2\].
###### Theorem 2.1.
If $`A𝒜^n`$, then $`G_A(0)=\frac{2}{\pi }\mathrm{log}f_A^{}(0)+k_0+O(n^{1/3}\mathrm{log}n)`$. Furthermore, if $`x0`$, then
$$G_A(x)=\frac{2}{\pi }g_A(x)+k_x+O(n^{1/3}\mathrm{log}n).$$
where
$$k_x:=k_0+\frac{2}{\pi }\mathrm{log}|x|a(x).$$
(3)
We conclude by defining what it means for two boundary arcs to be separated. Note that separation is always defined in terms of distance in the unit circle.
###### Definition 2.2.
Suppose that $`A𝒜`$ and $`D𝒟^{}`$. Let $`\mathrm{\Gamma }_1`$, $`\mathrm{{\rm Y}}_1_iA`$ with $`\overline{\mathrm{\Gamma }_1}\overline{\mathrm{{\rm Y}}_1}=\mathrm{}`$, let $`\mathrm{\Gamma }_2`$, $`\mathrm{{\rm Y}}_2D`$ with $`\overline{\mathrm{\Gamma }_2}\overline{\mathrm{{\rm Y}}_2}=\mathrm{}`$, and write $`\theta _1=\theta _A`$, $`\theta _2=\theta _D`$. The separation of $`\mathrm{\Gamma }_j`$ and $`\mathrm{{\rm Y}}_j`$, $`j=1,2`$, written $`\mathrm{sep}(\mathrm{\Gamma }_j,\mathrm{{\rm Y}}_j)`$, is defined to be
$$\mathrm{sep}(\mathrm{\Gamma }_j,\mathrm{{\rm Y}}_j):=inf\{|\theta _j(x)\theta _j(y)|:x\mathrm{\Gamma }_j,y\mathrm{{\rm Y}}_j\},$$
(4)
and the spread of $`\mathrm{\Gamma }_j`$ and $`\mathrm{{\rm Y}}_j`$, written $`\mathrm{spr}(\mathrm{\Gamma }_j,\mathrm{{\rm Y}}_j)`$, is defined to be
$$\mathrm{spr}(\mathrm{\Gamma }_j,\mathrm{{\rm Y}}_j):=sup\{|\theta _j(x)\theta _j(y)|:x\mathrm{\Gamma }_j,y\mathrm{{\rm Y}}_j\}.$$
(5)
If $`\mathrm{\Gamma }_1`$, $`\mathrm{{\rm Y}}_1A`$ instead, then (4) and (5) hold with $`\theta _A`$ extended to $`A`$ in the natural way.
## 3 Excursions and excursion measure
Much of this material may be found in and in the recent book . We repeat the relevant material here without proof in order to standardize our notation, and to remind the reader of the most important facts. We do, however, prove a number of useful lemmas about the Prohorov metric in Section 3.2.
### 3.1 Metric spaces of curves
A curve $`\gamma :I`$ shall always mean a continuous mapping of an interval $`I[0,\mathrm{})`$ into $``$. Let $`𝒦`$ denote the set of curves $`\gamma :[0,t_\gamma ]`$ where $`0<t_\gamma <\mathrm{}`$, and write $`\gamma [0,t_\gamma ]:=\{z:\gamma (t)=z\text{ for some }\mathrm{\hspace{0.17em}0}tt_\gamma \}`$ and similarly for $`\gamma (0,t_\gamma )`$. There are three natural metrics that we will consider on $`𝒦`$. Following , define the metric
$$d_𝒦^{}(\gamma ,\gamma ^{}):=\underset{\phi }{inf}\left[\underset{0st_\gamma }{sup}|\gamma (s)\gamma ^{}(\phi (s))|\right]$$
where the infimum is over all increasing homeomorphisms $`\phi :[0,t_\gamma ][0,t_\gamma ^{}]`$. Call $`\stackrel{~}{\gamma }`$ a reparameterization of $`\gamma 𝒦`$ with parameterization $`\phi `$ if $`\phi :[0,t_\gamma ][0,t_{\stackrel{~}{\gamma }}]`$ is an increasing homeomorphism such that $`\gamma (t)=\stackrel{~}{\gamma }(\phi (t))`$ for each $`0tt_\gamma `$. If $`\stackrel{~}{\gamma }`$ is a reparameterization of $`\gamma `$ under $`\phi `$, then $`\gamma `$ is a reparameterization of $`\stackrel{~}{\gamma }`$ under $`\phi ^1`$, and we write $`\gamma \stackrel{𝗉𝖺𝗋}{}\stackrel{~}{\gamma }`$. Finally, let $`𝒦^{}`$ be the set of equivalence classes of curves $`\gamma 𝒦`$ under the relation $`\stackrel{𝗉𝖺𝗋}{}`$, so that the metric $`d_𝒦^{}`$ identifies curves which are equal modulo time reparameterization. In fact, it can be shown that $`(𝒦^{},d_𝒦^{})`$ is a complete metric space \[2, Lemma 2.1\].
In order to account for the time parameterization, however, we let
$$d_𝒦(\gamma ,\gamma ^{}):=\underset{\phi }{inf}\left[\underset{0st_\gamma }{sup}\left\{|\gamma (s)\gamma ^{}(\phi (s))|+|s\phi (s)|\right\}\right]$$
where again the infimum is over all increasing homeomorphisms $`\phi :[0,t_\gamma ][0,t_\gamma ^{}]`$. The metric $`d_𝒦`$ does not identify curves which are equal modulo time reparameterization. A convenient choice of parameterization is $`\phi (s)=t_\gamma ^{}s/t_\gamma `$. Define
$$𝕕(\gamma ,\gamma ^{}):=\underset{0s1}{sup}|\gamma (t_\gamma s)\gamma ^{}(t_\gamma ^{}s)|+|t_\gamma t_\gamma ^{}|$$
and note that it is straightforward to verify $`𝕕`$ is also a metric on $`𝒦`$. Neither $`(𝒦,d_𝒦)`$ nor $`(𝒦,𝕕)`$ is complete as Example 3.2 combined with the next lemma will show. For the proof of this lemma, consult \[17, Lemma 5.1\].
###### Lemma 3.1.
If $`\gamma _1`$, $`\gamma _2𝒦`$, then $`d_𝒦(\gamma _1,\gamma _2)𝕕(\gamma _1,\gamma _2)d_𝒦(\gamma _1,\gamma _2)+\mathrm{osc}(\gamma _2,2d_𝒦(\gamma _1,\gamma _2))`$, where $`\mathrm{osc}(\gamma ,\delta ):=sup\{|\gamma (t)\gamma (s)|:|ts|\delta \}`$ denotes the modulus of continuity of $`\gamma `$.
To account for the incompleteness of $`(𝒦,𝕕)`$, we consider a larger complete space $`𝒳`$, and identify subspaces of $`(𝒦,𝕕)`$ with closed subspaces of $`𝒳`$. Let $`C[0,1]`$ denote the space of continuous complex-valued functions on $`[0,1]`$ under the metric $`d_{\mathrm{}}(\gamma _1^{},\gamma _2^{}):=sup_{0r1}|\gamma _1^{}(r)\gamma _2^{}(r)|`$, and denote the usual metric on $``$ by $`\mathrm{abs}`$. Consider the separable Banach space $`𝒳:=C[0,1]\times `$ with metric $`d_𝒳:=d_{\mathrm{}}+\mathrm{abs}`$. Thus, elements of $`𝒳`$ are pairs $`(\gamma ^{},t)`$ where $`\gamma ^{}C[0,1]`$, $`t`$, and $`d_𝒳((\gamma _1^{},s),(\gamma _2^{},t))=sup_{0r1}|\gamma _1^{}(r)\gamma _2^{}(r)|+|st|`$. We can embed $`𝒦`$ into $`𝒳`$ via $`\iota :𝒦𝒳`$, $`\gamma (\gamma ^{},t_\gamma )`$, where $`\gamma ^{}(r):=\gamma (t_\gamma r)`$, $`0r1`$. However, $`\iota (𝒦)=\{(\gamma ^{},t)𝒳:t>0\}=:𝒳^+`$ is not a closed subspace of $`𝒳`$. The metric spaces $`(𝒳^+,d_𝒳)`$ and $`(𝒦,d_{𝒳,𝒦})`$ are isomorphic, where $`d_{𝒳,𝒦}`$ is the induced metric in $`𝒦`$ associated to the metric $`d_𝒳`$ in $`𝒳`$. That is, if $`\gamma _1`$, $`\gamma _2𝒦`$, then $`\iota (\gamma _i)=(\gamma _i^{},t_{\gamma _i})`$, $`i=1`$, $`2`$, so that $`d_{𝒳,𝒦}(\gamma _1,\gamma _2)=d_𝒳((\gamma _1^{},t_{\gamma _1}),(\gamma _2^{},t_{\gamma _2}))`$. It follows that $`d_{𝒳,𝒦}=𝕕`$ and $`(𝒦,𝕕)(𝒳^+,d_𝒳)`$ since
$$d_𝒳((\gamma _1^{},t_{\gamma _1}),(\gamma _2^{},t_{\gamma _2}))=\underset{0r1}{sup}|\gamma _1(t_{\gamma _1}r)\gamma _2(t_{\gamma _2}r)|+|t_{\gamma _1}t_{\gamma _2}|=𝕕(\gamma _1,\gamma _2)$$
###### Example 3.2.
Suppose $`\gamma 𝒦`$ is given by $`\gamma (r)=r+ir`$, $`0r1`$, and for $`n=1,2,\mathrm{}`$, let $`\gamma _n(r)=nr+inr`$, $`0r1/n`$. Notice that $`\gamma _n^{}=\gamma ^{}=\gamma `$. Thus, $`\iota (\gamma _n)=(\gamma _n^{},t_{\gamma _n})=(\gamma ^{},1/n)`$ so clearly $`\{(\gamma _n^{},t_{\gamma _n})\}`$ is a Cauchy sequence in $`𝒳`$, and $`\{\gamma _n\}`$ is a Cauchy sequence in $`(𝒦,𝕕)`$. Since $`𝒳`$ is complete, it has a limit, namely $`(\gamma ^{},0)𝒳`$. However, $`(\gamma ^{},0)𝒳^+=\iota (𝒦)`$ so that $`(\gamma ^{},0)`$ does not have a counterpart in $`𝒦`$. This shows that $`(𝒦,𝕕)`$ is not complete, and illustrates the reason for considering $`𝒳`$.
However, if the limit does have a counterpart in $`𝒦`$ (i.e., if $`(\gamma ^{},t)𝒳^+`$ so that $`\iota ^1(\gamma ^{},t)𝒦`$), then we have the following result. See \[17, Lemma 5.2\] for the proof.
###### Lemma 3.3.
Let $`(\gamma _n^{},t_n)𝒳^+`$ for $`n=1,2,\mathrm{}`$, so that $`\gamma _n:=\iota ^1(\gamma _n^{},t_n)𝒦`$. Suppose that for some $`(\gamma ^{},t)𝒳`$, $`d_𝒳((\gamma _n^{},t_n),(\gamma ^{},t))0`$. If $`t>0`$ so that $`(\gamma ^{},t)𝒳^+`$, then $`\gamma :=\iota ^1(\gamma ^{},t)𝒦`$, and $`d_𝒳((\gamma _n^{},t_n),(\gamma ^{},t))0`$ if and only if $`d_𝒦(\gamma _n,\gamma )0`$ as $`n\mathrm{}`$, or equivalently, $`𝕕(\gamma _n,\gamma )0`$ if and only if $`d_𝒦(\gamma _n,\gamma )0`$ as $`n\mathrm{}`$.
Consequently, $`d_𝒦`$ and $`𝕕`$ generate the same topology on $`𝒦`$. Thus, when we need to discuss convergence or continuity in $`𝒦`$, it can be with respect to whichever metric is more convenient for the given problem.
If $`a>0`$, let $`𝒦_a:=\{\gamma 𝒦:t_\gamma a\}`$, and set $`\iota (𝒦_a)=\{(\gamma ^{},t)𝒳:ta\}=:𝒳_a`$. Note that $`𝒳_a`$ *is* a closed subspace of $`𝒳`$ so that $`(𝒳_a,d_𝒳)(𝒦_a,𝕕)`$ is complete. However, $`𝒦_a`$ is not complete under $`d_𝒦`$. As an example, consider $`\gamma _n(r)=r^n`$, $`0r1`$, which is a Cauchy sequence in $`(𝒦_1,d_𝒦)`$ that has no limit. By Lemma 3.1, if $`\{\gamma _n\}`$ is a Cauchy sequence in $`(𝒦_a,d_𝒦)`$ that is equicontinuous, then it is a Cauchy sequence in $`(𝒦_a,𝕕)`$ and therefore has a limit. In what follows, we will refer to spaces of curves which are primarily subspaces of $`𝒦`$. Since such spaces are isomorphic to subspaces of $`𝒳`$, we prefer to work with $`(𝒦,𝕕)`$ rather than $`(𝒳,d_𝒳)`$ unless it is necessary to explicitly mention this isomorphism.
If $`D`$ is a simply connected proper subset of $``$, and $`\gamma 𝒦`$, then we say that $`\gamma `$ is in $`D`$ if $`\gamma (0,t_\gamma )D`$. This does not require that either $`\gamma (0)D`$ or $`\gamma (t_\gamma )D`$. We define the space $`𝒦(D)`$ as $`𝒦(D):=\{\gamma 𝒦:\gamma \text{ is in }D\}`$. For $`z`$, $`w\overline{D}`$, let $`𝒦_z(D)`$ be the set of $`\gamma 𝒦(D)`$ with $`\gamma (0)=z`$, let $`𝒦^w(D)`$ be the set of $`\gamma 𝒦(D)`$ with $`\gamma (t_\gamma )=w`$, and define $`𝒦_z^w(D):=𝒦_z(D)𝒦^w(D)`$. Finally, if $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$, write $`𝒦_\mathrm{\Gamma }^\mathrm{{\rm Y}}(D):=_{z\mathrm{\Gamma },w\mathrm{{\rm Y}}}𝒦_z^w(D)`$.
###### Definition 3.4.
Suppose $`\gamma 𝒦(D)`$. We say $`\gamma `$ is an excursion in $`D`$ if $`\gamma (0)D`$ and $`\gamma (t_\gamma )D`$, and we say $`\gamma `$ is an excursion from $`z`$ to $`w`$ in $`D`$ if $`\gamma (0)=zD`$ and $`\gamma (t_\gamma )=wD`$, i.e., if $`\gamma 𝒦_z^w(D)`$ with $`z`$, $`wD`$. If $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$, then we say $`\gamma `$ is a $`(\mathrm{\Gamma },\mathrm{{\rm Y}})`$-excursion in $`D`$ if $`\gamma (0)\mathrm{\Gamma }`$ and $`\gamma (t_\gamma )\mathrm{{\rm Y}}`$, i.e., if $`\gamma 𝒦_\mathrm{\Gamma }^\mathrm{{\rm Y}}(D)`$.
Suppose that both $`D`$ and $`D^{}`$ are simply connected domains in $``$, and $`f:DD^{}`$ is a conformal transformation. For $`\gamma 𝒦(D)`$, let
$$A_s=A_{s,f,\gamma }:=_0^s|f^{}(\gamma (r))|^2dr\text{ and }\sigma _t=\sigma _{t,f,\gamma }:=inf\{s:A_st\}.$$
If $`\gamma 𝒦(D)`$ with $`A_{t_\gamma }<\mathrm{}`$, and if $`f`$ extends to the endpoints of $`\gamma `$, then we define the image of $`\gamma `$ under $`f`$, denoted $`f\gamma 𝒦(D^{})`$, by setting $`t_{f\gamma }:=A_{t_\gamma }`$ and $`f\gamma (t):=f(\gamma (\sigma _t))`$ for $`0tA_{t_\gamma }`$. Since $`sA_{s,f,\gamma }`$ is non-negative, continuous, and strictly increasing, it follows that $`t\sigma _{t,f,\gamma }`$ is well-defined. The following is a special case.
###### Example 3.5.
Let $`D`$ be a simply connected proper subsets of $``$, and for $`a\{0\}`$, let $`f_a(z)=az`$. If $`\gamma 𝒦(D)`$, then we define the Brownian scaling map $`\mathrm{\Psi }_a:𝒦(D)𝒦(f_a(D))`$ by setting $`t_{\mathrm{\Psi }_a\gamma }:=|a|^2t_\gamma `$ and $`\mathrm{\Psi }_a\gamma (t):=a\gamma \left(|a|^2t\right)`$ for $`0tt_{\mathrm{\Psi }_a\gamma }`$.
In particular, if $`D`$, $`D^{}𝒟`$, $`\gamma `$ is an excursion in $`D`$, and $`f𝒯(D,D^{})`$ so that $`f`$ *does* extend to the endpoints of $`\gamma `$, then $`f\gamma =:\gamma ^{}𝒦(D^{})`$ is an excursion in $`D^{}`$. Note that $`t_\gamma ^{}=A_{t_\gamma }`$ (i.e., $`\sigma _{t_\gamma ^{}}=t_\gamma `$) and $`\gamma ^{}(t)=f(\gamma (\sigma _t))`$ for $`0tt_\gamma ^{}`$.
###### Example 3.6.
As an application of Brownian scaling, suppose that $`f(z)=(1+\epsilon )z`$ for $`z𝔻`$, $`0<\epsilon <1`$, and let $`\gamma `$ be an excursion from $`x`$ to $`y`$ in $`𝔻`$. Then $`\gamma ^{}:=f\gamma `$ is an excursion from $`(1+\epsilon )x`$ to $`(1+\epsilon )y`$ in $`(1+\epsilon )𝔻`$ given explicitly by $`\gamma ^{}(t)=(1+\epsilon )\gamma \left((1+\epsilon )^2t\right)`$ for $`0tt_\gamma ^{}=(1+\epsilon )^2t_\gamma `$. Furthermore, it is not very difficult to verify that there exists a constant $`C=C(\gamma )`$ such that $`𝕕(\gamma ,\gamma ^{})C\epsilon `$.
If $`E`$ is a domain containing $`\overline{D}`$ and $`f`$ is a conformal mapping of $`E`$, then it follows from the Koebe growth and distortion theorems \[7, Theorems 2.4, 2.5, 2.6\] that $`|f^{}|`$, $`|f^{\prime \prime }|`$, and $`1/|f^{}|`$ are uniformly bounded on $`D`$, and the function $`\gamma f\gamma `$ from $`𝒦(D)`$ to $`𝒦(f(D))`$ is continuous. If $`D𝒟`$, then since $`D`$ is piecewise analytic, $`D=_{i=1}^n\mathrm{\Gamma }_i`$ for some finite union of analytic curves $`\mathrm{\Gamma }_i`$. Hence, any conformal mapping $`f`$ of $`D`$ can be analytically continued across each $`\mathrm{\Gamma }_i`$ so $`\gamma f\gamma :𝒦(D)𝒦(f(D))`$ is continuous; we denote this induced map by $`f`$.
###### Definition 3.7.
If $`\gamma _1`$, $`\gamma _2𝒦`$ with $`\gamma _1(t_{\gamma _1})=\gamma _2(0)`$, then we define the concatenation of $`\gamma _1`$ and $`\gamma _2`$, denoted $`\gamma _1\gamma _2`$, by setting $`t_{\gamma _1\gamma _2}:=t_{\gamma _1}+t_{\gamma _2}`$, and
$$\gamma _1\gamma _2(t):=\{\begin{array}{cc}\gamma _1(t),\hfill & \text{if }0tt_{\gamma _1},\hfill \\ \gamma _2(tt_{\gamma _1}),\hfill & \text{if }t_{\gamma _1}tt_{\gamma _1\gamma _2}.\hfill \end{array}$$
Note that $`(\gamma _1,\gamma _2)\gamma _1\gamma _2`$ is a continuous map from $`𝒦^w\times 𝒦_w`$ to $`𝒦`$ for every $`w`$.
###### Definition 3.8.
If $`0r<st_\gamma `$, then we define the truncation operator $`\mathrm{\Theta }_r^s:𝒦𝒦`$ by setting $`t_{\mathrm{\Theta }_r^s\gamma }:=sr`$ and $`\mathrm{\Theta }_r^s\gamma (t):=\gamma (r+t)`$ for $`0tt_{\mathrm{\Theta }_r^s\gamma }`$.
Observe that $`\mathrm{\Theta }_r^s\gamma [0,t_{\mathrm{\Theta }_r^s\gamma }]=\gamma [r,s]`$, and that by definition, truncation undoes concatenation. If $`\gamma _1`$, $`\gamma _2𝒦`$ with $`\gamma _1(t_{\gamma _1})=\gamma _2(0)`$, then $`\mathrm{\Theta }_0^{t_{\gamma _1}}\gamma _1\gamma _2(t)=\gamma _1(t)`$, $`0tt_{\gamma _1}`$, and $`\mathrm{\Theta }_{t_{\gamma _1}}^{t_{\gamma _1\gamma _2}}\gamma _1\gamma _2(t)=\gamma _2(t)`$, $`0tt_{\gamma _2}`$. It is easily seen that
$$d_𝒦(\mathrm{\Theta }_r^s\gamma ,\gamma )r+(t_\gamma s)+\mathrm{diam}(\gamma [0,r])+\mathrm{diam}(\gamma [s,t_\gamma ]).$$
Therefore, if $`r_n0+`$ and $`s_nt_\gamma `$, then by Lemma 3.3, $`𝕕(\mathrm{\Theta }_{r_n}^{s_n}\gamma ,\gamma )0`$.
### 3.2 General facts about measures on metric spaces
Throughout this section, suppose that $`(\mathrm{\Xi },\rho )`$ is a metric space. Let $`_\rho :=_\rho (\mathrm{\Xi })`$ denote the Borel $`\sigma `$-algebra associated to the topology induced by $`\rho `$, so that $`(\mathrm{\Xi },_\rho )`$ is a measurable space. A measure $`m`$ on $`(\mathrm{\Xi },\rho )`$ will always be a $`\sigma `$-finite measure on $`(\mathrm{\Xi },_\rho )`$. Denote the total mass of $`m`$ by $`|m|:=m(\mathrm{\Xi })`$. If $`|m|<\mathrm{}`$, then $`m`$ is a finite measure; otherwise it is an infinite measure. Denote the space of finite (resp., probability) measures on $`(\mathrm{\Xi },_\rho )`$ by $`(\mathrm{\Xi })`$ (resp., $`𝒫(\mathrm{\Xi })`$). If $`m(\mathrm{\Xi })`$ with $`|m|>0`$, we write $`m^\mathrm{\#}:=m/|m|`$ so that $`m^\mathrm{\#}𝒫(\mathrm{\Xi })`$. Recall (see ) that every finite measure $`m`$ on $`(\mathrm{\Xi },_\rho )`$ is regular; i.e., if $`V_\rho `$ and $`\epsilon >0`$, then there exist a closed set $`F`$ and an open set $`G`$ such that $`FVG`$ and $`m(GF)<\epsilon `$. If $`V_\rho `$, then we say that $`m`$ is concentrated on $`V`$ if $`m(\mathrm{\Xi }V)=0`$. Observe that $`V`$ need not be closed.
###### Definition 3.9.
If $`m_1`$, $`m_2=(\mathrm{\Xi })`$, let $`\mathrm{}:\times [0,\mathrm{})`$ denote the Prohorov metric given by $`\mathrm{}(m_1,m_2):=inf\{\epsilon >0:m_1(F)m_2(F^{(\epsilon )})+\epsilon ,m_2(F)m_1(F^{(\epsilon )})+\epsilon F_\rho \}`$ where $`F^{(\epsilon )}:=\{x\mathrm{\Xi }:\rho (x,y)<\epsilon ,\text{ some }yF\}`$.
It is easily checked that $`((\mathrm{\Xi }),\mathrm{})`$ is itself a metric space. Observe that $`F^{(\epsilon )}`$ is Borel, and that symmetry follows since $`((F^{(\epsilon )})^c)^{(\epsilon )}F^c`$. If $`m_1`$, $`m_2𝒫(\mathrm{\Xi })`$, then an equivalent definition of $`\mathrm{}`$ is given by $`\mathrm{}(m_1,m_2)=inf\{\epsilon >0:m_1(F)m_2(F^{(\epsilon )})+\epsilon \text{ for every closed }F_\rho \}`$. It is known \[5, Theorem 2.4.2\] that both metrics on $`𝒫(\mathrm{\Xi })`$ are equivalent and consistent with the Prohorov topology. Also note that $`||m_1||m_2||\mathrm{}(m_1,m_2)\mathrm{max}\{|m_1|,|m_2|\}`$. The following two theorems are standard.
###### Theorem 3.10.
Suppose that $`f`$ is a continuous mapping of the metric space $`(\mathrm{\Xi },\rho )`$ into the metric space $`(\mathrm{\Xi }^{},\rho ^{})`$. Then a measure $`m`$ on $`(\mathrm{\Xi },_\rho )`$ determines a measure $`m^{}`$ on $`(\mathrm{\Xi }^{},_\rho ^{})`$ such that $`fm(V^{})=m^{}(V^{})=m(f^1(V^{}))`$ for any $`V^{}_\rho ^{}`$. That is, $`fm(\mathrm{\Xi }^{})`$ is given explicitly by $`fm(V^{}):=m(\{x\mathrm{\Xi }:f(x)V^{}\})`$ for any $`V^{}_\rho ^{}`$. Furthermore,
$$_\mathrm{\Xi }^{}h(x^{})m^{}(\mathrm{d}x^{})=_\mathrm{\Xi }h(f(x))m(\mathrm{d}x)$$
for any bounded, continuous function $`h:\mathrm{\Xi }^{}`$.
###### Theorem 3.11.
If $`(\mathrm{\Xi },\rho )`$ is a complete, separable metric space, then the metric space $`(𝒫(\mathrm{\Xi }),\mathrm{})`$ is also complete and separable, where $`\mathrm{}`$ is the Prohorov metric as in Definition 3.9. Furthermore, if $`m_n`$, $`m𝒫(\mathrm{\Xi })`$, then as $`n\mathrm{}`$, $`\mathrm{}(m_n,m)0`$ if and only if $`m_nm`$ weakly.
###### Important Remark 3.12.
Whenever we say a sequence of measures converges, it will be with respect to the Prohorov metric.
As noted in \[5, page 29\], Strassen proved another equivalent definition of $`\mathrm{}`$ consistent with the Prohorov topology is given by
$$\mathrm{}(m_1,m_2)=\underset{𝔐}{inf}\left[inf\{\epsilon >0:\{\rho (X_1,X_2)\epsilon \}\epsilon \}\right],$$
where $`𝔐`$ is the set of all $`\mathrm{\Xi }\times \mathrm{\Xi }`$-valued random variables $`(X_1,X_2)`$ with $`(X_1)=m_1`$ and $`(X_2)=m_2`$ where $``$ denotes law. In fact, an easy calculation shows that if $`X_i`$ are $`(\mathrm{\Xi },\rho )`$-valued random variables with $`(X_i)=m_i`$, $`i=1,2`$, and if $`\{\rho (X_1,X_2)\epsilon \}\epsilon `$, then $`\mathrm{}(m_1,m_2)\epsilon `$. If $`(\mathrm{\Xi },\rho )`$ is a complete and separable metric space, then to show a sequence of non-zero finite measures $`m_n(\mathrm{\Xi })`$ converges to $`m(\mathrm{\Xi })`$, it suffices to show that both $`|m_n||m|`$ and $`\mathrm{}(m_n^\mathrm{\#},m^\mathrm{\#})0`$, as $`n\mathrm{}`$. In particular, we record the following version of these remarks.
###### Proposition 3.13.
Let $`\gamma `$, $`\gamma ^{}`$ be $`𝒦`$-valued random variables with $`(\gamma )=\mu `$ and $`(\gamma ^{})=\mu ^{}`$, respectively. If $`\{𝕕(\gamma ,\gamma ^{})\epsilon \}\epsilon `$, then $`\mathrm{}(\mu ,\mu ^{})\epsilon `$.
###### Lemma 3.14.
Suppose that $`(\mathrm{\Xi },\rho )`$ is a complete, separable metric space, and that $`m_1`$, $`m_2(\mathrm{\Xi })`$. If $`C>0`$, then $`\mathrm{}(Cm_1,Cm_2)(C1)\mathrm{}(m_1,m_2)`$.
###### Proof.
Suppose that $`\mathrm{}(m_1,m_2)=\epsilon `$. To begin, let $`C>1`$. Then since $`m_1(F)m_2(F^{(\epsilon )})+\epsilon `$ for every $`F`$ Borel, we have $`Cm_1(F)Cm_2(F^{(\epsilon )})+C\epsilon `$. Since $`C\epsilon >\epsilon `$, we have $`F^{(\epsilon )}F^{(C\epsilon )}`$. Hence, $`Cm_1(F)Cm_2(F^{(C\epsilon )})+C\epsilon `$. Interchanging $`m_1`$ and $`m_2`$ gives $`\mathrm{}(Cm_1,Cm_2)C\epsilon `$. Suppose instead that $`C<1`$. Then since $`m_1(F)m_2(F^{(\epsilon )})+\epsilon `$, and $`C\epsilon <\epsilon `$, we have $`m_1(F)m_2(F^{(\epsilon )})+\epsilon /C`$. Multiplying by $`C`$ gives $`Cm_1(F)Cm_2(F^{(\epsilon )})+\epsilon `$. Interchanging $`m_1`$ and $`m_2`$ yields $`\mathrm{}(Cm_1,Cm_2)\epsilon `$. Thus, the conclusion follows. ∎
###### Lemma 3.15.
If $`f:(\mathrm{\Xi },\rho )(\mathrm{\Xi }^{},\rho ^{})`$ is continuous, $`m(\mathrm{\Xi })`$, and $`C`$ is a constant, then
$$f(Cm)=C(fm).$$
(6)
###### Proof.
Since $`f(Cm)(V^{})=(Cm)(f^1(V^{}))=C\left[m(f^1(V^{}))\right]=C\left[fm(V^{})\right]`$ for any $`V^{}_\rho ^{}(\mathrm{\Xi }^{})`$ the result follows. ∎
###### Lemma 3.16.
Suppose that $`(\mathrm{\Xi },\rho )`$ is a complete, separable metric space, and let $`m(\mathrm{\Xi })`$. If $`f_n,f:(\mathrm{\Xi },\rho )(\mathrm{\Xi },\rho )`$ are continuous, and $`f_nf`$ uniformly, then $`\mathrm{}(f_nm,fm)0`$.
###### Proof.
Assume first that $`m𝒫(\mathrm{\Xi })`$. If $`\mu _n:=f_nm`$ and $`\mu :=fm`$, then by Theorem 3.11, it suffices to show that $`\mu _n\mu `$ weakly. Suppose that $`h:\mathrm{\Xi }`$ is a bounded, continuous function. Hence, by Theorem 3.10, we conclude that
$$_\mathrm{\Xi }h(x)\mu _n(\mathrm{d}x)=_\mathrm{\Xi }hf_n(x)m(\mathrm{d}x)_\mathrm{\Xi }hf(x)m(\mathrm{d}x)=_\mathrm{\Xi }h(x)\mu (\mathrm{d}x)$$
since $`f_nf`$ uniformly. We next consider $`m(\mathrm{\Xi })`$. If $`m`$ is the zero measure, the result is trivial. If $`|m|>0`$, then by (6) and Lemma 3.14, $`\mathrm{}(f_nm,fm)=\mathrm{}(|m|(f_nm^\mathrm{\#}),|m|(fm^\mathrm{\#}))=(|m|1)\mathrm{}(f_nm^\mathrm{\#},fm^\mathrm{\#})0`$. ∎
###### Lemma 3.17.
Under the same assumptions as Lemma 3.16, if $`m_2(\mathrm{\Xi })`$, and $`\mathrm{}(f_nm_1,fm_1)0`$ as $`n\mathrm{}`$, then $`\mathrm{}(f_nm_1,m_2)\mathrm{}(fm_1,m_2)`$.
###### Proof.
Since $`\mathrm{}(,)`$ is a metric, we have by the triangle inequality $`\mathrm{}(f_nm_1,m_2)\mathrm{}(f_nm_1,fm_1)+\mathrm{}(fm_1,m_2)`$, so that
$$\underset{n\mathrm{}}{lim\; sup}\mathrm{}(f_nm_1,m_2)\mathrm{}(fm_1,m_2).$$
(7)
However, the reverse inequality tells us that $`\mathrm{}(fm_1,m_2)\mathrm{}(fm_1,f_nm_1)+\mathrm{}(f_nm_1,m_2)`$, so that
$$\underset{n\mathrm{}}{lim\; inf}\mathrm{}(f_nm_1,m_2)\mathrm{}(fm_1,m_2).$$
(8)
By combining (7) and (8), the result follows. ∎
We conclude this section by reviewing how to define a measure by Riemann integration. Let $`\mathrm{\Lambda }`$ be an analytic arc that is parameterized by $`\xi :[0,t_\xi ]`$ with $`t_\xi <\mathrm{}`$. Consider the measures $`\{\mu (z,):z\mathrm{\Lambda }\}(\mathrm{\Xi })`$, and let $`\{\xi (0)=z_0,z_1,\mathrm{},z_n=\xi (t_\xi )\}`$ partition $`\mathrm{\Lambda }`$. Let $`z_i^{}[z_{i1},z_i]`$, $`|\mathrm{\Delta }z_i|=|z_iz_{i1}|`$, $`i=1,\mathrm{}n`$, and set
$$\mu _n():=\underset{i=1}{\overset{n}{}}\mu (z_i^{},)|\mathrm{\Delta }z_i|.$$
Note that $`\mu _n()(\mathrm{\Xi })`$ for each $`n`$. If $`lim_n\mathrm{}\mu _n()`$ exists in $`(\mathrm{\Xi })`$, then we define the Riemann integral of the measure-valued function $`z\mathrm{\Lambda }\mu (z,)(\mathrm{\Xi })`$ to be this limiting value; that is,
$$\mu ():=_\mathrm{\Lambda }\mu (z,)|\mathrm{d}z|:=\underset{n\mathrm{}}{lim}\mu _n().$$
(9)
Several conditions guarantee the existence of the Riemann integral. For instance, if $`z\mu (z,)`$ is continuous at $`z_0`$ for all $`z_0\mathrm{\Lambda }`$, or if $`(\mathrm{\Xi },\rho )`$ is a complete and separable metric space and $`\{\mu _n()\}`$ is a Cauchy sequence, then (9) exists in $`(\mathrm{\Xi })`$.
### 3.3 Excursion Poisson kernel
We now briefly review several results about the Poisson kernel and the excursion Poisson kernel. Further details including proofs can be found in . Let $`D`$ be a domain, and write $`𝐧_x=𝐧_{x,D}`$ for the unit normal at $`x`$ pointing into $`D`$. If $`xD`$ and $`D`$ is locally analytic at $`yD`$, then both harmonic measure and its density with respect to arc length, the Poisson kernel $`H_D(x,y)`$, are well-defined. The behaviour of the Poisson kernel under a conformal transformation can be easily deduced from the Riemann mapping theorem and Lévy’s theorem on the conformal invariance of Brownian motion . See also \[15, Proposition 2.10\].
###### Proposition 3.18.
If $`f:DD^{}`$ is a conformal transformation, $`xD`$, $`D`$ is locally analytic at $`yD`$, and $`D^{}`$ is locally analytic at $`f(y)`$, then
$$H_D(x,y)=|f^{}(y)|H_D^{}(f(x),f(y)).$$
(10)
Furthermore, if $`\mathrm{\Gamma }D`$ and $`f\mathrm{\Gamma }D^{}`$ are analytic, then
$$H_D(x,\mathrm{\Gamma }):=_\mathrm{\Gamma }H_D(x,y)|\mathrm{d}y|=H_D^{}(f(x),f\mathrm{\Gamma }).$$
(11)
###### Definition 3.19.
For $`x`$, $`yD`$, $`xy`$, the excursion Poisson kernel $`H_D(x,y)`$ is given by
$$H_D(x,y):=\underset{\epsilon 0+}{lim}\frac{1}{\epsilon }H_D(x+\epsilon 𝐧_x,y)=\underset{\epsilon 0+}{lim}\frac{1}{\epsilon }H_D(y+\epsilon 𝐧_y,x).$$
For a proof of the next proposition, see \[15, Proposition 2.11\].
###### Proposition 3.20.
Suppose $`f:DD^{}`$ is a conformal transformation and $`x`$, $`y`$ are distinct points on $`D`$. Suppose that $`D`$ is locally analytic at $`x`$, $`y`$, and $`D^{}`$ is locally analytic at $`f(x)`$, $`f(y)`$. Then $`H_D(x,y)=|f^{}(x)||f^{}(y)|H_D^{}(f(x),f(y))`$.
###### Corollary 3.21.
If $`x`$, $`y𝔻`$, $`xy`$, and $`f𝒯(𝔻,𝔻)`$ with $`f(x)=x`$ and $`f(y)=y`$, then $`|f^{}(x)||f^{}(y)|=1`$.
###### Proof.
If $`f(x)=x`$ and $`f(y)=y`$, then we immediately obtain from Proposition 3.20 that $`H_𝔻(x,y)=|f^{}(x)||f^{}(y)|H_𝔻(f(x),f(y))=|f^{}(x)||f^{}(y)|H_𝔻(x,y)`$. ∎
### 3.4 Brownian excursion measures on $`(𝒦,𝕕)`$
We now remind the reader of several Brownian measures on $`(𝒦,𝕕)`$, and outline the construction of the Brownian excursion measure on $`D`$. The exposition follows , although there are some noticeable differences. We begin with a general definition.
###### Definition 3.22.
A measure $`\mu `$ on $`𝒦`$ is defined to be a $`\sigma `$-finite measure on the measurable space $`(𝒳,_{d_𝒳})`$ concentrated on $`𝒳^+`$.
Suppose $`B_t`$ is a Brownian motion with $`B_0=z`$, and let $`T_D:=inf\{t:B_tD\}`$ be its exit time from $`D`$. The process $`X_t^D:=B_{tT_D}`$, $`t0`$, is Brownian motion killed on exiting $`D`$. Let $`D𝒟`$ and suppose $`wD`$ so that the Poisson kernel $`H_D(z,w)`$ is well-defined. Define the continuous, positive martingale $`M`$ by $`M_t^D:=H_D(X_t^D,w)/H_D(z,w)`$, and the probability $`_D^{z,w}`$ by $`_D^{z,w}(A):=𝔼^z[M_t^D;A]`$ for $`A_t`$. As noted in , the law of the process $`X_t^D`$ under $`_D^{z,w}`$ is that of Brownian motion conditioned to exit $`D`$ at $`w`$.
###### Definition 3.23.
Suppose that $`D𝒟`$. The interior-to-boundary excursion measure from $`z`$ to $`w`$ in $`D`$, written $`\mu _D(z,w)`$, is defined to be $`\mu _D(z,w):=H_D(z,w)_D^{z,w}`$, and the interior-to-boundary excursion measure from $`z`$ in $`D`$, written $`\mu _D(z)`$, is defined by
$$\mu _D(z):=_D\mu _D(z,w)|\mathrm{d}w|.$$
(12)
Translation invariance and Brownian scaling imply $`w\mu _D(z,w)`$ is continuous so that (12) is well-defined as in (9). The measure on paths $`\mu _D(z)`$ is what is generally called Wiener measure. Observe that $`\mu _D(z)`$ is a measure on $`𝒦`$ concentrated on $`𝒦_z(D)`$, and that by definition, $`\mu _D(z,w)`$ is a finite measure with mass $`|\mu _D(z,w)|=H_D(z,w)`$. As such we can consider the normalized probability measure
$$\mu _D^\mathrm{\#}(z,w):=\frac{\mu _D(z,w)}{|\mu _D(z,w)|}=\frac{\mu _D(z,w)}{H_D(z,w)}:=_D^{z,w}.$$
(13)
It is well-known that two-dimensional Brownian motion is conformally invariant, and consequently so too is Wiener measure. We express this as follows. If $`D`$, $`D^{}𝒟`$, $`zD`$, and $`f𝒯(D,D^{})`$, then $`f\mu _D(z)=\mu _D^{}(f(z))`$. This definition is independent of the choice of $`f𝒯(D,D^{})`$; indeed, if $`f_1`$, $`f_2𝒯(D,D^{})`$ with $`f_1(z)=f_2(z)=z^{}`$, then $`f_1\mu _D(z)=\mu _D^{}(f_1(z))=\mu _D^{}(z^{})=\mu _D^{}(f_2(z))=f_2\mu _D(z)`$. The first part of the next proposition follows from a quick change-of-variables, while the second follows immediately from the first as a result of (10). See also \[17, Proposition 5.5\].
###### Proposition 3.24.
Suppose that $`D`$, $`D^{}𝒟`$, and $`zD`$, $`wD`$ with $`D`$ locally analytic at $`w`$. If $`f𝒯(D,D^{})`$, and $`D^{}`$ is locally analytic at $`f(w)`$, then $`f\mu _D(z,w)=|f^{}(w)|\mu _D^{}(f(z),f(w))`$ and $`f\mu _D^\mathrm{\#}(z,w)=\mu _D^{}^\mathrm{\#}(f(z),f(w))`$.
Using the interior-to-boundary excursion measure, we now define boundary-to-boundary excursion measure in $`𝔻`$, and show that it exists by an explicit calculation. It is then a simple matter to define excursion measure for other simply connected $`D`$, and to derive the important conformal invariance formula.
###### Definition 3.25.
If $`x`$, $`y𝔻`$, $`xy`$, then normalized excursion measure on excursions from $`x`$ to $`y`$ in $`𝔻`$ is the measure on $`𝒦`$, concentrated on $`𝒦_x^y(𝔻)`$, defined by
$$\underset{\epsilon 0+}{lim}\mu _𝔻^\mathrm{\#}((1\epsilon )x,y)=:\overline{\mu }_𝔻(x,y),$$
(14)
where $`\mu _𝔻^\mathrm{\#}(z,y)`$ for $`z𝔻`$, $`y𝔻`$ is as in (13).
###### Lemma 3.26.
The limit in (14) exists.
###### Proof.
Let $`\gamma 𝒦(𝔻)`$ with $`\gamma (0)=0`$, $`\gamma (t_\gamma )=1`$. Let $`f_\alpha (z)=\frac{z\alpha }{1\alpha z}`$ for $`\alpha (1,1)`$ so that $`f_\alpha 𝒯(𝔻,𝔻)`$, $`f_\alpha (0)=\alpha `$, and both $`1`$ and $`1`$ are fixed points of $`f_\alpha `$. Using the exact form of the Möbius transformation $`f_\alpha `$, a straightforward computation shows that $`lim_{\alpha 1}f_\alpha \gamma `$ exists in the metric space $`(𝒦,𝕕)`$ where $`f_\alpha \gamma `$ is defined as on page 3.1. In particular, this shows $`lim_{\epsilon 0+}\mu _𝔻^\mathrm{\#}((1\epsilon ),1)=:\overline{\mu }_𝔻(1,1)`$ exists. For other $`x`$ and $`y`$, simply use a composition of Möbius transformations. ∎
###### Definition 3.27.
We define excursion measure on excursions from $`x`$ to $`y`$ in $`𝔻`$ to be the measure on $`𝒦`$, concentrated on $`𝒦_x^y(𝔻)`$, defined by $`\mu _𝔻(x,y)=H_𝔻(x,y)\overline{\mu }_𝔻(x,y)`$ where $`H_𝔻`$ denotes the excursion Poisson kernel.
Observe that by definition, the mass of excursion measure $`\mu _{}(x,y)`$ is *defined* to be $`|\mu _𝔻(x,y)|=H_𝔻(x,y)`$. Hence, $`\mu _𝔻^\mathrm{\#}(x,y)=\overline{\mu }_𝔻(x,y)`$.
###### Definition 3.28.
Suppose that $`D𝒟`$, and $`z`$, $`wD`$ with $`D`$ locally analytic at both $`z`$ and $`w`$. Let $`h𝒯(𝔻,D)`$. Excursion measure from $`z`$ to $`w`$ in $`D`$ is defined by
$$\mu _D(z,w):=\frac{1}{|h^{}(h^1(z))||h^{}(h^1(w))|}h\mu _𝔻(h^1(z),h^1(w)).$$
(15)
A straightforward exercise in the chain rule shows that the definition of $`\mu _D(z,w)`$ given by (15) does not depend on the choice of $`h𝒯(𝔻,D)`$.
###### Proposition 3.29.
Let $`D`$, $`D^{}𝒟`$, and let $`z`$, $`wD`$ with $`D`$ locally analytic at both $`z`$, $`w`$. If $`f𝒯(D,D^{})`$, and $`D^{}`$ is locally analytic at both $`f(z)`$, $`f(w)`$, then
$$f\mu _D(z,w)=|f^{}(z)||f^{}(w)|\mu _D^{}(f(z),f(w))$$
(16)
and
$$f\mu _D^\mathrm{\#}(z,w)=\mu _D^\mathrm{\#}(f(z),f(w)).$$
(17)
If $`f_1`$, $`f_2𝒯(D,D^{})`$, then $`f_1\mu _D(z,w)=f_2\mu _D(z,w)`$ so (16) and (17) are independent of the choice of map. In particular,
$$\mu _D(z,w)=\underset{\epsilon 0+}{lim}\frac{1}{\epsilon }\mu _D(z+\epsilon 𝐧_z,w).$$
###### Definition 3.30.
Suppose that $`D𝒟`$. Excursion measure in $`D`$ is defined by
$$\mu _D:=_D_D\mu _D(z,w)|\mathrm{d}w||\mathrm{d}z|.$$
The conformal invariance of excursion measure is immediate; see \[17, Proposition 5.8\].
###### Proposition 3.31 (Conformal Invariance).
If $`D`$, $`D^{}𝒟`$ and $`f𝒯(D,D^{})`$, then $`f\mu _D=\mu _D^{}`$.
In fact, it should be noted that we can define excursion measure $`\mu _D`$ for *any* simply connected subset of $``$ by conformal invariance. The only reason to restrict to $`D𝒟^{}`$ in the next definition is so that excursions $`\gamma 𝒦(D)`$ will necessarily have $`t_\gamma <\mathrm{}`$. (We will not be concerned with excursions with $`t_\gamma =\mathrm{}`$ in this paper.)
###### Definition 3.32.
Suppose that $`D𝒟^{}`$ and $`f𝒯(𝔻,D)`$. Excursion measure in $`D`$ is defined by $`\mu _D:=f\mu _𝔻`$. Furthermore, if $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$, define $`\mu _D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ to be the measure $`\mu _D`$ restricted to those excursions $`\gamma 𝒦_\mathrm{\Gamma }^\mathrm{{\rm Y}}(D)`$, and define the excursion Poisson kernel $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ to be its mass; that is, $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}}):=|\mu _D(\mathrm{\Gamma },\mathrm{{\rm Y}})|`$.
An immediate consequence of these definitions is the following.
###### Proposition 3.33 (Conformal Invariance).
If $`D`$, $`D^{}𝒟^{}`$; $`f𝒯(D,D^{})`$; $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$; and $`\mathrm{\Gamma }^{}`$, $`\mathrm{{\rm Y}}^{}`$ are the images under $`f`$ of $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, respectively, then $`f\mu _D(\mathrm{\Gamma },\mathrm{{\rm Y}})=\mu _D^{}(\mathrm{\Gamma }^{},\mathrm{{\rm Y}}^{})`$ and $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})=H_D^{}(\mathrm{\Gamma }^{},\mathrm{{\rm Y}}^{})`$.
### 3.5 Discrete excursions and discrete excursion measure
Throughout this section, suppose that $`A𝒜`$; $`w`$, $`zA`$; $`x`$, $`yA`$; and $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}A`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$. Our goal is to define a discrete excursion and formulate the discrete analogues of the previous sections. If $`S_j`$ is a simple random walk with $`S_0=w`$, denote the one-step transition probability $`p(w,z):=^w\{S_1=z\}`$, and define the discrete Poisson kernel to be $`h_A(w,y):=^w\{S_{\tau _A}=y\}`$ where $`\tau _A:=\mathrm{min}\{j>0:S_jA\}`$. As in \[8, §3.1\],
$$q(w,z;y):=^w\{S_1=z|S_{\tau _A}=y\}=p(w,z)\frac{h_A(z,y)}{h_A(w,y)}$$
(18)
defines the one-step transition probabilities of simple random walk conditioned to exit $`A`$ at $`y`$. Note that $`h_A`$ is discrete harmonic, and that (18) an example of a discrete $`h`$-transform.
###### Definition 3.34.
A discrete excursion in $`A`$ is a path $`\omega :=[\omega _0,\omega _1,\mathrm{},\omega _k]`$ where $`\omega _0A`$, $`\omega _kA`$, $`|\omega _i\omega _{i1}|=1`$ for $`i=1,\mathrm{},k`$, and $`\omega _iA`$ for $`i=1,\mathrm{},k1`$, where $`2k<\mathrm{}`$. If $`\omega =[\omega _0,\omega _1,\mathrm{},\omega _k]`$, define the length of $`\omega `$, written $`|\omega |`$, to be $`k`$. If $`\omega `$ is a discrete excursion in $`A`$ with $`\omega _0\mathrm{\Gamma }`$ and $`\omega _k\mathrm{{\rm Y}}`$, then $`w`$ is called a $`(\mathrm{\Gamma },\mathrm{{\rm Y}})`$-discrete excursion in $`A`$. In particular, if $`\omega _0=x`$ and $`\omega _k=y`$, then $`\omega `$ is called a discrete excursion from $`x`$ to $`y`$ in $`A`$.
Discrete excursions can be generated by starting a simple random walk at $`xA`$, conditioning it to take its first step into $`A`$, and stopping it at $`\tau _A`$. Let the discrete excursion Poisson kernel $`h_A(x,y)`$ be given by $`h_A(x,y):=^x\{S_{\tau _A}=y,S_1A\}`$, and define discrete excursion measure to be the measure that assigns weight $`4^{|\omega |}`$ to each discrete excursion $`\omega `$. Denote this measure by $`\mu _A^{\mathrm{𝗋𝗐}}()`$ so that $`\mu _A^{\mathrm{𝗋𝗐}}(\omega ):=4^{|\omega |}`$. Write $`\mu _A^{\mathrm{𝗋𝗐}}(x,y)`$ to denote the measure on discrete excursions from $`x`$ to $`y`$ in $`A`$, and $`\mu _A^{\mathrm{𝗋𝗐}}(\mathrm{\Gamma },\mathrm{{\rm Y}}):=_{x\mathrm{\Gamma }}_{y\mathrm{{\rm Y}}}\mu _A^{\mathrm{𝗋𝗐}}(x,y)`$ to denote the measure on $`(\mathrm{\Gamma },\mathrm{{\rm Y}})`$-discrete excursions in $`A`$. In , Lawler and Werner defined $`\mu _A^{\mathrm{𝗋𝗐}}(\omega ):=(2\pi \mathrm{\hspace{0.17em}4}^{|\omega |})^1`$; this difference only affects things up to a constant.
We want both discrete excursion measure and Brownian excursion measure to be measures on the metric space $`(𝒦,𝕕)`$. Consequently, we need to associate to each discrete excursion $`\omega `$ a curve $`\stackrel{~}{\omega }𝒦`$. Suppose that $`\omega `$ is a discrete excursion in $`A`$, and let $`\mathrm{cl}(A):=AA`$ with associated domain $`\stackrel{~}{\mathrm{cl}(A)}`$. We associate to $`\omega `$ a curve $`\stackrel{~}{\omega }𝒦\left(\stackrel{~}{\mathrm{cl}(A)}\right)`$ by setting $`t_{\stackrel{~}{\omega }}:=2|\omega |`$, and
$$\stackrel{~}{\omega }(t):=\omega _{t/2}+\frac{1}{2}(tt)(\omega _{t/2+1}\omega _{t/2}),0tt_{\stackrel{~}{\omega }}.$$
(19)
In other words, we join the lattice points in order with line segments parallel to the coordinate axes in $`^2`$, with each segment taking time 2 to traverse. Note that $`\stackrel{~}{\omega }(0)=\omega _0`$ and $`\stackrel{~}{\omega }(t_{\stackrel{~}{\omega }})=\omega _{|\omega |}`$. Using this identification, if $`\omega `$ is an excursion from $`x`$ to $`y`$ in $`A`$, then $`\mu _A^{\mathrm{𝗋𝗐}}(x,y)(𝒦)`$ and $`\mu _{A,x,y}^{\mathrm{𝗋𝗐}}(\stackrel{~}{\omega })=4^{t_{\stackrel{~}{\omega }}}`$. In order to prove discrete excursion measure converges to Brownian excursion measure in the scaling limit, we will consider scaling excursions as the mesh of the lattice becomes finer. See (25) in Section 4.4.
As a consequence of the so-called KMT approximation (see Section 4.5.1), it follows that $`|B_tS_{2t}|=O(\mathrm{log}t)`$. Complete details may be found in and . Thus, it is simply a matter of æsthetics that a random walk path of $`|\omega |`$ steps take time $`2|\omega |`$ to traverse: if $`\gamma `$ is Brownian curve and $`\stackrel{~}{\omega }`$ is as above with $`\gamma (0)=\stackrel{~}{\omega }(0)`$, then $`|\gamma (t)\stackrel{~}{\omega }(t)|=O(\mathrm{log}t)`$.
###### Definition 3.35.
Suppose that $`A𝒜`$ and $`x`$, $`yA`$. Discrete excursion measure $`\mu _A^{\mathrm{𝗋𝗐}}(x,y)`$ is defined to be the measure on $`(𝒦,𝕕)`$, concentrated on $`V=V(x,y;A):=\{\gamma 𝒦:𝕕(\gamma ,\stackrel{~}{\omega })=0`$ for some discrete excursion $`\omega `$ from $`x`$ to $`y`$ in $`A\}`$ given by $`\mu _A^{\mathrm{𝗋𝗐}}(x,y)(\gamma ):=4^{t_\gamma }`$ for $`\gamma V`$. Note that $`\mu _A^{\mathrm{𝗋𝗐}}(x,y)(V)=h_A(x,y)`$.
## 4 The main convergence arguments
### 4.1 Carathéodory convergence
###### Definition 4.1.
Fix $`r>0`$. Suppose that $`D_n`$ is a sequence of domains with $`D_n𝒟^r`$ for each $`n`$. The kernel of $`D_n`$, written $`\mathrm{ker}(\{D_n\})`$, is the largest domain $`D`$ containing the origin and having the property that each compact subset of $`D`$ lies in all but a finite number of the domains $`D_n`$. Suppose that $`\mathrm{ker}(\{D_n\})=D`$. The sequence $`D_n`$ converges in the Carathéodory sense to $`D`$, written $`D_n\stackrel{𝖼𝖺𝗋𝖺}{}D`$, if every subsequence $`D_{n_j}`$ of $`D_n`$ has $`\mathrm{ker}(\{D_{n_j}\})=D`$.
Recall that a sequence of functions $`f_n`$ on a domain $`D`$ converges to a function $`f`$ uniformly on compacta of $`D`$ if for each compact $`KD`$, $`f_nf`$ uniformly on $`K`$. The following theorem, which roughly states that convergence of domains in the Carathéodory sense is equivalent to the uniform convergence on compacta of the appropriate Riemann maps, will be our main tool. A proof may be found in \[7, Theorem 3.1\].
###### Theorem 4.2 (Carathéodory Convergence).
Suppose that $`D_n`$ is a sequence of domains with $`D_n𝒟^{}`$ for each $`n`$, and let $`f_n𝒯(𝔻,D_n)`$ with $`f_n(0)=0`$, $`f_n^{}(0)>0`$. Suppose further that $`D𝒟`$ and $`f𝒯(𝔻,D)`$ with $`f(0)=0`$, $`f^{}(0)>0`$. Then $`f_nf`$ uniformly on compacta of $`𝔻`$ if and only if $`D_n\stackrel{𝖼𝖺𝗋𝖺}{}D`$.
###### Lemma 4.3.
Suppose that $`D_n\stackrel{𝖼𝖺𝗋𝖺}{}D`$ with $`D_n`$, $`D𝒟^{}`$. Suppose further that there exists an $`E𝒟^{}`$ with $`D_nE`$ for all $`n`$, and $`DE`$. If $`F:E𝔻`$ is the conformal transformation with $`F(0)=0`$, $`F^{}(0)>0`$, then $`F(D_n)\stackrel{𝖼𝖺𝗋𝖺}{}F(D)`$.
###### Proof.
Let $`f_n:𝔻D_n`$ and let $`f:𝔻D`$ be conformal transformations mapping 0 to 0 with positive derivative at the origin. By Theorem 4.2, the convergence of $`D_n`$ to $`D`$ is equivalent to the uniform convergence of $`f_n`$ to $`f`$ on compacta of $`𝔻`$. Set $`h_n:=Ff_n`$ and $`h:=Ff`$, and let $`K`$ be a compact subset of $`𝔻`$. If $`zK`$, then $`|h_n(z)h(z)|=|F(f_n(z))F(f(z))|0`$ uniformly as $`n\mathrm{}`$. ∎
###### Lemma 4.4.
Suppose that $`F_n`$, $`F`$ are conformal mappings of $`𝔻`$. Let $`D:=F(𝔻)`$. If $`F_nF`$ uniformly on compacta of $`𝔻`$, then $`F_nF^1I`$ uniformly on compacta of $`D`$, where $`I:DD`$ is the identity map $`I(z)=z`$.
###### Proof.
Let $`K^{}D`$ be compact. Let $`\epsilon >0`$ be given. Let $`K=F^1(K^{})𝔻`$ which is clearly compact. By uniform convergence, there exists $`N=N(\epsilon ,K)`$ such that $`|F_n(x)F(x)|<\epsilon `$ for all $`n>N`$, $`xK`$. If $`yK^{}`$, then $`y=F(x)`$ for some $`xK`$. Hence, if $`n>N`$, then $`|F_nF^1(y)I(y)|=|F_n(x)F(x)|<\epsilon `$, and the proof is complete. ∎
###### Lemma 4.5.
Suppose that $`F_n`$, $`F`$ are conformal mappings of the unit disk $`𝔻`$, and that $`F_nF`$ uniformly on compacta of $`𝔻`$. If $`\gamma 𝒦(𝔻)`$ with $`|\gamma (0)|<1`$ and $`|\gamma (t_\gamma )|<1`$, then $`𝕕(F_n\gamma ,F\gamma )0`$ as $`n\mathrm{}`$.
###### Proof.
Suppose that $`\gamma 𝒦(𝔻)`$ with $`|\gamma (0)|<1`$ and $`|\gamma (t_\gamma )|<1`$. Note that $`\gamma `$ is *not* an excursion in $`𝔻`$. Therefore, there necessarily exists a compact set $`K𝔻`$ such that $`\gamma 𝒦(K)`$. Consider $`t_{F_n\gamma }=A_{t_\gamma }^n=_0^{t_\gamma }|F_n^{}(\gamma (r))|^2dr`$ and $`t_{F\gamma }=A_{t_\gamma }=_0^{t_\gamma }|F^{}(\gamma (r))|^2dr`$. Since $`F_nF`$ uniformly on compacta of $`𝔻`$, we necessarily have that $`F_nF`$ uniformly on $`K`$. Hence, it follows that $`t_{F_n\gamma }t_{F\gamma }`$. Furthermore,
$`\underset{0s1}{sup}|F\gamma (t_{F\gamma }s)F_n\gamma (t_{F_n\gamma }s)|`$
$`\underset{0s1}{sup}|F\gamma (t_{F\gamma }s)F\gamma (t_{F_n\gamma }s)|+|F\gamma (t_{F_n\gamma }s)F_n\gamma (t_{F_n\gamma }s)|0.`$
Taken together, these imply the result. ∎
### 4.2 Construction of approximate domains $`\stackrel{~}{D}_N`$
Suppose that $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and let
$$D_N^{\prime \prime }:=\{x\frac{1}{N}^2D:\frac{1}{N}𝒮_xD\},$$
where $`𝒮_x:=x+\left([1/2,1/2]\times [1/2,1/2]\right)`$ is the unit square about $`x`$. Let $`D_N^{}`$ be the connected component of $`D_N^{\prime \prime }`$ containing the origin, and set $`D_N:=D_N^{}_iD_N^{}`$. Take $`\stackrel{~}{D}_N`$ to be the interior of the union of the scaled squares centred at those $`xD_N`$. We call $`D_N`$ the $`1/N`$-scale discrete approximation to $`D`$ (with respect to the origin), and we informally refer to $`\stackrel{~}{D}_N`$ as the associated “union of squares” domain; that is,
$$\stackrel{~}{D}_N=\mathrm{int}\left(\underset{xD_N}{}\frac{1}{N}𝒮_x\right)\text{ and }\mathrm{cl}(\stackrel{~}{D}_N):=\stackrel{~}{D}_N\stackrel{~}{D}_N=\underset{xD_N}{}\frac{1}{N}𝒮_x.$$
(20)
Let $`f𝒯(𝔻,D)`$ with $`f(0)=0`$, $`f^{}(0)>0`$. Let $`\mathrm{\Gamma }_𝔻`$, $`\mathrm{{\rm Y}}_𝔻𝔻`$ be (open) boundary arcs with $`\overline{\mathrm{\Gamma }_𝔻}\overline{\mathrm{{\rm Y}}_𝔻}=\mathrm{}`$; that is, $`\mathrm{\Gamma }_𝔻:=\{e^{i\theta }:\theta _1<\theta <\theta _2\}`$ and $`\mathrm{{\rm Y}}_𝔻:=\{e^{i\theta }:\theta _3<\theta <\theta _4\}`$, for some $`0\theta _1<\theta _2<\theta _3<\theta _4<\theta _1+2\pi `$. Define $`\mathrm{\Gamma }D`$ to be the image of $`\mathrm{\Gamma }_𝔻`$ under $`f`$, and similarly, let $`\mathrm{{\rm Y}}D`$ be the image of $`\mathrm{{\rm Y}}_𝔻`$ under $`f`$. Let $`s:=\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ as in Definition 2.2, and let $`N`$ be sufficiently large so that $`s\epsilon _n:=n^{1/48}\mathrm{log}^{2/3}n`$ if $`nN`$. If $`f_N𝒯(𝔻,\stackrel{~}{D}_N)`$ with $`f_N(0)=0`$, $`f_N^{}(0)>0`$, then define $`\stackrel{~}{\mathrm{\Gamma }}_N`$ to be the image of $`\mathrm{\Gamma }_𝔻`$ under $`f_N`$, with $`\stackrel{~}{\mathrm{{\rm Y}}}_N`$ defined similarly. In Theorem 4.9, we prove $`f_Nf`$ uniformly on compacta of $`𝔻`$ showing $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$.
We now define our approximating discrete boundary arcs. If $`\stackrel{~}{\mathrm{\Gamma }}_N\stackrel{~}{D}_N`$, then associate to $`\stackrel{~}{\mathrm{\Gamma }}_N`$ the set $`\mathrm{\Gamma }_ND_N`$ as follows. Let $`\mathrm{\Gamma }_N^{}:=\{x_iD_N:\frac{1}{N}𝒮_x\stackrel{~}{\mathrm{\Gamma }}_N\mathrm{}\}`$, and then take
$$\mathrm{\Gamma }_N:=\{yD_N:(x,y)_eD_N\text{ with }x\mathrm{\Gamma }_N^{}\text{ and }\frac{1}{N}\mathrm{}_{x,y}\stackrel{~}{\mathrm{\Gamma }}_N\}.$$
Similarly, let $`\mathrm{{\rm Y}}_N`$ be the discrete boundary arc associated to $`\stackrel{~}{\mathrm{{\rm Y}}}_N`$. Our notation is summarized in the following table.
| $`𝔻`$ | $`D`$, $`D𝒟^{}`$ | $`\stackrel{~}{D}_N`$, $`\stackrel{~}{D}_N𝒟`$ | $`D_N\frac{1}{N}^2`$, $`2ND_N𝒜^N`$ |
| --- | --- | --- | --- |
| $`\mathrm{\Gamma }_𝔻,\mathrm{{\rm Y}}_𝔻𝔻`$ | $`\mathrm{\Gamma },\mathrm{{\rm Y}}D`$ | $`\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N\stackrel{~}{D}_N`$ | $`\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_ND_N`$ |
Note that by conformal invariance, it is equivalent to specify either $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$, or $`\mathrm{\Gamma }_𝔻`$, $`\mathrm{{\rm Y}}_𝔻𝔻`$. We have (arbitrarily) chosen the latter.
### 4.3 Convergence of domains $`\stackrel{~}{D}_N`$ to $`D`$
Suppose that $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$ and let $`D_N`$ be the $`1/N`$-scale discrete approximation to $`D`$ with associated “union of squares” domain $`\stackrel{~}{D}_N`$ as in Section 4.2. The following lemmas are an immediate consequence of those definitions.
###### Lemma 4.6.
For each $`N`$, $`\stackrel{~}{D}_N𝒟`$ with $`\mathrm{cl}(\stackrel{~}{D}_N)D`$. That is, $`\stackrel{~}{D}_N`$ is a simply connected proper subset of $`D`$ with piecewise analytic boundary. Furthermore, the lattice $`\mathrm{cl}(D_N):=D_ND_ND`$.
###### Lemma 4.7.
Suppose that $`x_iD_N`$, $`yD_N`$, and $`z\stackrel{~}{D}_N`$. Then $`\mathrm{dist}(x,D)c_1N^1`$, $`\mathrm{dist}(y,D)c_2N^1`$, and $`\mathrm{dist}(z,D)c_3N^1`$ where $`c_1=2\sqrt{2}+1/\sqrt{2}`$, $`c_2=\sqrt{2}+1/\sqrt{2}`$, and $`c_3=2\sqrt{2}`$.
The next proposition follows from the Beurling estimate see \[17, Proposition 3.79\].
###### Proposition 4.8.
If $`x_iD_N`$ and $`f𝒯(D,𝔻)`$ with $`f(0)=0`$, $`f^{}(0)>0`$, then there exists a constant $`C`$ such that $`\mathrm{dist}(f(x),𝔻)CN^{1/2}`$ and $`f(\stackrel{~}{D}_N)\{|z|1CN^{1/2}\}`$.
We will now establish Theorem 1.2 with the proof of the following result.
###### Theorem 4.9.
The sets $`\stackrel{~}{D}_N`$ as defined by (20) converge to $`D`$ in the Carathéodory sense.
The proof of this theorem requires two lemmas. The first is a simple power series estimate, while the second gives good bounds on the difference of the image of a point under two different maps: the identity map from $`𝔻`$ to $`𝔻`$, and a map which is “almost the identity.”
###### Lemma 4.10.
If $`0|z|1/2`$, then $`|\mathrm{log}(1+z)z||z|/2`$.
###### Proof.
Since
$$\mathrm{log}(1+z)=\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n1}\frac{1}{n}z^n,$$
we have
$$|\mathrm{log}(1+z)z|\underset{n=2}{\overset{\mathrm{}}{}}\frac{1}{n}|z|^n\frac{1}{2}|z|\underset{n=1}{\overset{\mathrm{}}{}}|z|^n\frac{1}{2}|z|$$
provided that $`0|z|1/2`$. ∎
###### Lemma 4.11.
For $`N>4C^2`$, where $`C`$ is the constant in Lemma 4.8, suppose that $`E_N`$ is a domain with $`\{|z|1CN^{1/2}\}E_N\{|z|1+CN^{1/2}\}`$. Let $`h_N:𝔻E_N`$ be the conformal transformation with $`h_N(0)=0`$ and $`h_N^{}(0)>0`$. Then, there exists a constant $`C^{}`$ such that $`|h_N(z)z|C^{}N^{1/2}\mathrm{log}N`$ for $`|z|1CN^{1/2}`$.
###### Proof.
Without loss of generality, assume that $`h_N`$ may be extended to an analytic function in a neighbourhood of $`\overline{𝔻}`$. For if this is not the case, we may approximate $`h_N`$ by $`h_{N,r}(z):=r^1h_N(rz)`$ and take the limit as $`r1`$. From the Schwarz lemma \[1, page 135\], we can immediately see that $`1CN^{1/2}h_N^{}(0)1+CN^{1/2}`$. Let $`\kappa _N(z):=\mathrm{log}[h_N(z)/z]`$ so that $`\kappa _N=u_N+iv_N`$ is analytic on $`𝔻`$ with $`|u_N(z)|(3/2)CN^{1/2}`$ for $`|z|=1`$ using the estimate from Lemma 4.10. Thus, the maximum principle for harmonic functions tells us that $`|u_N(z)|(3/2)CN^{1/2}`$ for all $`|z|1`$. We therefore conclude that the partial derivatives of $`u_N`$ at $`z`$ are bounded by an absolute constant times $`N^{1/2}\mathrm{dist}(z,𝔻)^1`$; whence $`|\kappa _N^{}(z)|C_1N^{1/2}(1|z|)^1`$. Writing
$`\left|\mathrm{log}\left[1+{\displaystyle \frac{h_N(z)z}{z}}\right]\right|=|\kappa _N(z)|=\left|\kappa _N(0)+{\displaystyle _0^z}\kappa _N^{}(w)dw\right|`$ $`{\displaystyle \frac{C_2}{\sqrt{N}}}\left[1+\mathrm{log}{\displaystyle \frac{1}{1|z|}}\right]`$
with $`C_2=\mathrm{max}\{C,C_1\}`$, we see that if $`\epsilon >0`$ is such that
$$\left|\frac{h_N(z)z}{z}\right|\frac{1}{2}\text{ for }|z|\epsilon ,$$
(21)
then
$$\left|\frac{h_N(z)z}{z}\right|2\left|\mathrm{log}\left[1+\frac{h_N(z)z}{z}\right]\right|2C_2N^{1/2}\left[1+\mathrm{log}\frac{1}{1|z|}\right].$$
(22)
Since (21) holds for some $`\epsilon >0`$, we can iterate (22) to see that (22) must hold for all $`|z|`$ such that the right side of (22) is less than 1/2. For $`N`$ sufficiently large, this includes all $`|z|1CN^{1/2}`$. ∎
###### Proof of Theorem 4.9.
Suppose that $`f:D𝔻`$ is the conformal transformation with $`f(0)=0`$, $`f^{}(0)>0`$, and let $`\stackrel{~}{f}_N:f(\stackrel{~}{D}_N)𝔻`$ be the conformal transformation with $`\stackrel{~}{f}_N(0)=0`$, $`\stackrel{~}{f}_N^{}(0)>0`$. Let $`F_N:𝔻\stackrel{~}{D}_N`$ and $`F:𝔻D`$ be the conformal transformations with $`F_N(0)=0`$, $`F_N^{}(0)>0`$, and $`F(0)=0`$, $`F^{}(0)>0`$, respectively, which are defined by setting $`F_N:=(\stackrel{~}{f}_Nf)^1`$ and $`F:=f^1=(If)^1`$ where $`I(z)=z`$ is the identity map from $`𝔻`$ to $`𝔻`$. Finally, let $`z𝔻`$, and let $`w:=\stackrel{~}{f}_N^1(z)`$ so that $`F_N(z)=F(w)`$.
We prove that $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$ by applying Theorem 4.2 which states that it is sufficient to show $`F_NF`$ uniformly on each compact subset of $`𝔻`$. Equivalently, we will show that for each $`\delta >0`$ sufficiently small, $`F_NF`$ uniformly for all $`|z|1\delta `$. Fix $`0<\delta <1/2`$ and choose $`M`$ so that $`M>(3C^{}\delta ^1)^3`$ where $`C^{}`$ is the constant in Lemma 4.11. Let $`N>M`$. Then by Lemma 4.11, we have that for $`|z|1\delta `$,
$$|wz|\frac{C^{}\mathrm{log}N}{\sqrt{N}}|z|\left(\frac{C^{}\mathrm{log}N}{\sqrt{N}}\frac{1\delta }{\delta }\right)\delta .$$
Our choice of $`M`$ guarantees that $`C^{}\delta ^1(1\delta )N^{1/2}\mathrm{log}N<1`$ for $`N>M`$. By \[17, Corollary 3.25\], if for some $`0<r<1`$, $`|wz|r\mathrm{dist}(z,𝔻)`$, then
$$|F(w)F(z)|\frac{4\mathrm{dist}(F(z),D)}{1r^2}|wz|.$$
Hence, we conclude
$$|F_N(z)F(z)|=|F(w)F(z)|\left(\frac{4RC^{}(1\delta )}{1\left(\frac{C^{}\mathrm{log}N}{\sqrt{N}}\frac{1\delta }{\delta }\right)^2}\right)\frac{\mathrm{log}N}{\sqrt{N}}$$
where $`R:=\mathrm{rad}(D)`$ so that $`F_NF`$ uniformly; whence $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$. ∎
###### Corollary 4.12.
If $`F𝒯(D,𝔻)`$ with $`F(0)=0`$, $`F^{}(0)>0`$, then $`F(\stackrel{~}{D}_N)\stackrel{𝖼𝖺𝗋𝖺}{}𝔻`$.
###### Proof.
By Lemma 4.6, $`\stackrel{~}{D}_ND`$, so Lemma 4.3 yields the result. ∎
### 4.4 Applying results for $`A𝒜^N`$ to $`D_N`$
Suppose that $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$. In this section, we combine our construction of $`D_N`$ with Theorem 2.1 and \[15, Theorem 1.1\] to restate those results for random walk on $`D_N`$. The most difficult part of this section is keeping track of the notation.
We begin by mentioning several scaling relationships that will be needed throughout. If $`S_n`$ is a random walk on $`^2`$, then for any $`r>0`$ there is an associated random walk (which we will also denote by $`S_n`$) on the lattice $`r^2`$. In other words, there is a one-to-one correspondence between paths from $`x`$ to $`y`$ in $`A`$ on $`^2`$, and paths from $`rx`$ to $`ry`$ in $`rA`$ on $`r^2`$. Hence if $`A^2`$ and $`r>0`$, then $`G_{rA}(rx,ry)=G_A(x,y)`$, where the Green’s function on the left side is for random walk on the lattice $`r^2`$, and the Green’s function on the right side is for random walk on $`^2`$. Similarly, we have $`h_{rA}(rx,ry)=h_A(x,y)`$ for the discrete Poisson kernel, and $`h_{rA}(rx,ry)=h_A(x,y)`$ for the discrete excursion Poisson kernel.
The conformal invariance of the Green’s function for Brownian motion implies that if $`D𝒟^{}`$ and $`r>0`$, then $`g_{rD}(rx,ry)=g_D(x,y)`$. However, from the conformal covariance of the Poisson kernel (Proposition 3.18) and the excursion Poisson kernel (Proposition 3.20), it follows that $`rH_{rD}(rx,ry)=H_D(x,y)`$ and $`r^2H_{rD}(rx,ry)=H_D(x,y)`$.
Note that a random walk on $`D_N`$ is taking steps of size $`1/N`$. Therefore, let $`A_N:=2ND_N`$ so that $`A_N𝒜^N`$, and $`\stackrel{~}{A}_N:=\stackrel{~}{(2ND_N)}=2N\stackrel{~}{D}_N𝒟`$. Hence, $`z^{}A_N`$ if and only if $`z:=z^{}/2ND_N`$. Suppose $`x^{}:=2NxA^N`$ with $`xD_N`$ and $`y^{}:=2NyA^N`$ with $`yD_N`$. Thus, when the above scaling is applied to $`\stackrel{~}{A}_N`$, we conclude that
$$g_{A_N}(x^{},y^{})=g_{2ND_N}(2Nx,2Ny)=g_{D_N}(x,y),$$
(23)
where $`g_{D_N}`$ denotes the Green’s function for Brownian motion in $`\stackrel{~}{D}_N`$. In particular, if $`f_{D_N}𝒯(\stackrel{~}{D}_N,𝔻)`$ with $`f_{D_N}(0)=0`$, $`f_{D_N}^{}(0)>0`$ and $`f_{A_N}𝒯(\stackrel{~}{A}_N,𝔻)`$ with $`f_{A_N}(0)=0`$, $`f_{A_N}^{}(0)>0`$, then since $`f_{A_N}(x^{})=f_{D_N}(x)`$ and $`g_{A_N}(x^{})=g_{D_N}(x)`$, and since we can write $`f_{A_N}(x^{})=\mathrm{exp}\{g_{A_N}(x^{})+i\theta _{A_N}(x^{})\}`$ and $`f_{D_N}(x)=\mathrm{exp}\{g_{D_N}(x)+i\theta _{D_N}(x)\}`$, it follows that $`\theta _{A_N}(x^{})=\theta _{2ND_N}(2Nx)=\theta _{D_N}(x)`$. Further, in the random walk case, $`G_{A_N}(x^{},y^{})=G_{2ND_N}(2Nx,2Ny)=G_{D_N}(x,y)`$, and for $`x^{}:=2NxA_N`$ with $`xD_N`$, we have $`h_{A_N}(x^{},y^{})=h_{2ND_N}(2Nx,2Ny)=h_{D_N}(x,y)`$; similarly, $`h_{A_N}(0,x^{})=h_{D_N}(0,x)`$, and $`h_{A_N}(0,y^{})=h_{D_N}(0,y)`$.
For $`A_N𝒜^N`$, let $`A_N^{}:=\{x^{}A_N:g_{A_N}(x^{})N^{1/16}\}`$ which is consistent with the usage in . If $`x^{}A_N^{}`$, $`y^{}A_N`$, then Theorem 2.1 (in particular, its corollary \[15, Corollary 3.5\]) implies that
$$G_{A_N}(x^{},y^{})=\frac{2}{\pi }g_{A_N}(x^{},y^{})+k_{y^{}x^{}}+O(N^{7/24}\mathrm{log}N).$$
(24)
With the above notation in hand, we are finally able to state the following corollary to (24).
###### Corollary 4.13.
Let $`xD_N`$ be such that $`x^{}:=2Nx(2ND_N)^{}=A_N^{}`$, and let $`yD_N`$ with $`y^{}:=2NyA_N`$. Then,
$$G_{D_N}(x,y)=\frac{2}{\pi }g_{D_N}(x,y)+k_{y^{}x^{}}+O(N^{7/24}\mathrm{log}N)$$
where $`k_z`$ is as in (3).
Note that $`k_{y^{}x^{}}cN^{3/2}|xy|^{3/2}`$. Thus, if $`|xy|N^{29/36}`$, then $`k_{y^{}x^{}}=O(N^{7/24})`$, and we have a refined version of the previous corollary.
###### Corollary 4.14.
If $`xD_N`$ with $`x^{}:=2NxA_N^{}`$, $`yD_N`$ with $`y^{}:=2NyA_N`$, and $`|xy|N^{29/36}`$, then
$$G_{D_N}(x,y)=\frac{2}{\pi }g_{D_N}(x,y)+O(N^{7/24}\mathrm{log}N).$$
We also have the following corollary to \[15, Theorem 1.1\].
###### Corollary 4.15.
If $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$, and $`x`$, $`yD_N`$ with $`|\theta _{D_N}(x)\theta _{D_N}(y)|N^{1/16}\mathrm{log}^2N`$, then
$$h_{D_N}(x,y)=\frac{(\pi /2)h_{D_N}(0,x)h_{D_N}(0,y)}{1\mathrm{cos}(\theta _{D_N}(x)\theta _{D_N}(y))}\left[1+O\left(\frac{\mathrm{log}N}{N^{1/16}|\theta _{D_N}(x)\theta _{D_N}(y)|}\right)\right].$$
We now make several observations regarding excursion measure. Suppose $`x`$, $`y\stackrel{~}{D}_N`$ so that $`x^{}:=2Nx`$, $`y^{}:=2Ny\stackrel{~}{A}_N`$ as above. If $`f(z)=2Nz`$, then $`f𝒯(\stackrel{~}{D}_N,\stackrel{~}{A}_N)`$ with $`f(0)=0`$ and $`f^{}(z)=2N`$ for all $`z`$. Since excursion measure is conformally covariant/invariant, we are able to conclude that $`\mu _{\stackrel{~}{D}_N}(x,y)=4N^2\mu _{\stackrel{~}{A}_N}(x^{},y^{})`$ and $`\mu _{\stackrel{~}{D}_N}=\mu _{\stackrel{~}{A}_N}`$. We know from Donsker’s theorem that simple random walk converges in the scaling limit to Brownian motion provided that space and time are scaled appropriately. In order to prove Theorem 1.1, we will need to apply a similar scaling. Recall from (19) that if $`\omega `$ is a discrete excursion then we can associate to it a curve $`\stackrel{~}{\omega }𝒦`$, and that the Brownian scaling map $`\mathrm{\Psi }_a`$ was defined in Example 3.5. For $`N`$, write $`\mathrm{\Phi }_N:=\mathrm{\Psi }_{1/(2N)}`$ so that
$$\mathrm{\Phi }_N\stackrel{~}{\omega }(t)=\frac{1}{2N}\stackrel{~}{\omega }(4N^2t)\text{ for }\mathrm{\hspace{0.33em}\hspace{0.33em}0}tt_{\mathrm{\Phi }_N\stackrel{~}{\omega }}=\frac{t_{\stackrel{~}{\omega }}}{4N^2}=\frac{|\omega |}{2N^2}.$$
(25)
###### Lemma 4.16.
If $`\gamma `$, $`\gamma ^{}𝒦`$, then
$$\frac{1}{4N^2}𝕕(\gamma ,\gamma ^{})𝕕(\mathrm{\Phi }_N\gamma ,\mathrm{\Phi }_N\gamma ^{})\frac{1}{2N}𝕕(\gamma ,\gamma ^{}).$$
###### Proof.
From the definitions of $`𝕕`$ and $`\mathrm{\Phi }_N`$ we conclude that
$`𝕕(\mathrm{\Phi }_N\gamma ,\mathrm{\Phi }_N\gamma ^{})`$ $`=\underset{0s1}{sup}|\mathrm{\Phi }_N\gamma (st_{\mathrm{\Phi }_N\gamma })\mathrm{\Phi }_N\gamma ^{}(st_{\mathrm{\Phi }_N\gamma ^{}})|+|t_{\mathrm{\Phi }_N\gamma }t_{\mathrm{\Phi }_N\gamma ^{}}|`$
$`=\underset{0s1}{sup}\left|{\displaystyle \frac{1}{2N}}\gamma (st_\gamma ){\displaystyle \frac{1}{2N}}\gamma ^{}(st_\gamma ^{})\right|+{\displaystyle \frac{1}{4N^2}}|t_\gamma t_\gamma ^{}|`$
so the result follows. ∎
###### Definition 4.17.
Suppose that $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$, and $`x`$, $`yD_N`$. The $`1/N`$-scale discrete excursion measure $`\mu _{D_N}^{\mathrm{𝗋𝗐}}(x,y)`$ is defined to be the measure on $`(𝒦,𝕕)`$, concentrated on $`V_N=V_N(x,y;D):=\{\gamma 𝒦:𝕕(\gamma ,\mathrm{\Phi }_N\stackrel{~}{\omega })=0`$ for some discrete excursion $`\omega `$ from $`2Nx`$ to $`2Ny`$ in $`2ND_N\}`$ given by $`\mu _{D_N}^{\mathrm{𝗋𝗐}}(x,y)(\gamma ):=4^{4N^2t_\gamma }=4^{|\omega |}`$ for $`\gamma V_N`$. Finally, if $`\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_ND_N`$ with $`\overline{\mathrm{\Gamma }_N}\overline{\mathrm{{\rm Y}}_N}=\mathrm{}`$, then
$$\mu _{D_N}^{\mathrm{𝗋𝗐}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N):=\underset{x\mathrm{\Gamma }_N}{}\underset{y\mathrm{{\rm Y}}_N}{}\mu _{D_N}^{\mathrm{𝗋𝗐}}(x,y).$$
### 4.5 Proof of Theorem 1.1
In the present section, we establish the following theorem which, as noted in the introduction, may be regarded as the precise formulation of Theorem 1.1.
###### Theorem 4.18.
Suppose $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and let $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ be open boundary arcs with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$. For every $`\epsilon >0`$, there exists an $`N`$ such that
1. $`\left|h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N){\displaystyle \frac{1}{4}}H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})\right|\epsilon `$,
2. $`\mathrm{}(\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N),\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}))\epsilon `$, and
3. $`\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N))\epsilon `$,
where $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$, $`\stackrel{~}{D}_N𝒟`$ is the “union of squares” domain associated to $`D_N`$, and $`\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_ND_N`$ are the corresponding discrete boundary arcs with associated boundary arcs $`\stackrel{~}{\mathrm{\Gamma }}_N`$, $`\stackrel{~}{\mathrm{{\rm Y}}}_N\stackrel{~}{D}_N`$, respectively. In particular, *(a)*, *(b)*, and *(c)* imply that
$$\underset{N\mathrm{}}{lim}\mathrm{}(\mathrm{\hspace{0.33em}4}\mu _{D_N}^{\mathrm{𝗋𝗐}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _D(\mathrm{\Gamma },\mathrm{{\rm Y}}))=0.$$
Each of the three parts of Theorem 4.18 will be proved in a separate section: in Section 4.5.2 we prove Theorem 4.25 establishing (a), in Section 4.5.3 we prove Theorem 4.26 establishing (b), and finally in Section 4.5.4 we prove Theorem 4.29 establishing (c).
#### 4.5.1 Review of strong approximation of Brownian motion and random walk
In order to establish Theorem 4.18, it will be necessary to use a strong approximation result which follows from the theorem of Komlós, Major, and Tusnády . Because of its central rôle in the proof, we include the statement here for the convenience of the reader. In what follows, $`S_t`$ is defined for non-integer $`t`$ by linear interpolation.
###### Theorem 4.19 (Komlós-Major-Tusnády).
There exists $`c<\mathrm{}`$ and a probability space $`(\mathrm{\Omega },,)`$ on which are defined a two-dimensional Brownian motion $`B`$ and a two-dimensional simple random walk $`S`$ with $`B_0=S_0`$, such that for all $`\lambda >0`$ and each $`n`$,
$$\left\{\underset{0tn}{\mathrm{max}}\left|\frac{1}{\sqrt{2}}B_tS_t\right|>c(\lambda +1)\mathrm{log}n\right\}<cn^\lambda .$$
The proofs of the following two results may be found in \[15, Corollary 3.2\] and \[15, Proposition 3.3\], respectively.
###### Corollary 4.20.
There exist $`C<\mathrm{}`$ and a probability space $`(\mathrm{\Omega },,)`$ on which are defined a two-dimensional Brownian motion $`B`$ and a two-dimensional simple random walk $`S`$ with $`B_0=S_0`$ such that
$$\left\{\underset{0t\sigma _n}{\mathrm{max}}\left|\frac{1}{\sqrt{2}}B_tS_t\right|>C\mathrm{log}n\right\}=O(n^{10}),$$
where $`\sigma _n^1:=inf\{t:|S_tS_0|n^8\}`$, $`\sigma _n^2:=inf\{t:|B_tB_0|n^8\}`$, and $`\sigma _n:=\sigma _n^1\sigma _n^2`$.
###### Proposition 4.21 (Strong Approximation).
There exists a constant $`c`$ such that for every $`n`$, a Brownian motion $`B`$ and a simple random walk $`S`$ can be defined on the same probability space so that if $`A𝒜^n`$, $`1<rn^{20}`$, and $`xA`$ with $`|x|n^3`$, then $`^x\{|B_{T_A}S_{\tau _A}|cr\mathrm{log}n\}cr^{1/2}`$.
By combining the strong approximation with Theorem 4.19, the following estimate is easily deduced.
###### Proposition 4.22.
There exists a decreasing sequence $`\delta _n0`$ such that if $`A𝒜^n`$ with associated “union of squares” domain $`\stackrel{~}{A}𝒟`$, and $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}A`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$ and associated boundary arcs $`\stackrel{~}{\mathrm{\Gamma }}`$, $`\stackrel{~}{\mathrm{{\rm Y}}}\stackrel{~}{A}`$, then $`h_A(0,\mathrm{\Gamma })=H_{\stackrel{~}{A}}(0,\stackrel{~}{\mathrm{\Gamma }})+O(\delta _n)`$, and $`h_A(0,\mathrm{{\rm Y}})=H_{\stackrel{~}{A}}(0,\stackrel{~}{\mathrm{{\rm Y}}})+O(\delta _n)`$. Consequently, $`h_A(0,\mathrm{\Gamma })h_A(0,\mathrm{{\rm Y}})=H_{\stackrel{~}{A}}(0,\stackrel{~}{\mathrm{\Gamma }})H_{\stackrel{~}{A}}(0,\stackrel{~}{\mathrm{{\rm Y}}})+O(\delta _n)`$ where the error term depends on both $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$.
###### Proof.
If $`c`$ is the constant in Proposition 4.21, and $`V`$ is the set $`V:=\{xA:\mathrm{dist}(x,\mathrm{\Gamma })cn^{1/8}\mathrm{log}n\}`$, then $`h_A(0,\mathrm{\Gamma })=H_{\stackrel{~}{A}}(0,\stackrel{~}{V})+O(n^{1/16})`$. However, a simple gambler’s ruin estimate for Brownian motion shows that $`H_{\stackrel{~}{A}}(0,\stackrel{~}{V})=H_{\stackrel{~}{A}}(0,\stackrel{~}{\mathrm{\Gamma }})+O(n^{7/8}\mathrm{log}n)`$, so the result follows with $`\delta _n=n^{7/8}\mathrm{log}n`$. ∎
#### 4.5.2 Convergence of $`4h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$ to $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$
The goal of the present section is to prove that if $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ are disjoint open boundary arcs, then $`4h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ using the notation from Section 4.2, therefore establishing Theorem 4.18 (a).
It follows from the exact form of the excursion Poisson kernel in $`𝔻`$ \[17, Example 5.6\] that if $`D𝒟`$ and $`x`$, $`yD`$ with $`D`$ locally analytic at $`x`$ and $`y`$, then
$$H_D(x,y)=\frac{2\pi H_D(0,x)H_D(0,y)}{1\mathrm{cos}(\theta _D(x)\theta _D(y))}.$$
(26)
For further details, see also \[15, Example 2.14\].
###### Lemma 4.23.
If $`D𝒟^{}`$ and $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$, then
$$\frac{2\pi H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})\frac{2\pi H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}$$
where $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ is as in Definition 3.32, and $`H_D(0,\mathrm{\Gamma })`$, $`H_D(0,\mathrm{{\rm Y}})`$ are as in (11).
###### Proof.
Suppose first that $`D𝒟`$, and that $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$ are analytic open boundary arcs. Then from (26), we conclude that for all $`x\mathrm{\Gamma }`$ and for all $`y\mathrm{{\rm Y}}`$,
$$\frac{2\pi H_D(0,x)H_D(0,y)}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}H_D(x,y)\frac{2\pi H_D(0,x)H_D(0,y)}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}.$$
Since $`D𝒟`$, Proposition 3.20 implies that
$$\frac{2\pi H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})\frac{2\pi H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}.$$
Now, suppose $`D^{}𝒟^{}`$, and let $`f𝒯(D,D^{})`$. Write $`\mathrm{\Gamma }^{}`$, $`\mathrm{{\rm Y}}^{}`$ for the images under $`f`$ of $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, respectively. Conformal invariance yields $`H_D(0,\mathrm{\Gamma })=H_D^{}(0,\mathrm{\Gamma }^{})`$ and $`H_D(0,\mathrm{{\rm Y}})=H_D^{}(0,\mathrm{{\rm Y}}^{})`$. (Indeed this holds for all domains $`D𝒟^{}`$ since $`D`$ is regular.) From Proposition 3.20, it follows that $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})=H_D^{}(\mathrm{\Gamma }^{},\mathrm{{\rm Y}}^{})`$; whence the proof is complete. ∎
Let $`f𝒯(𝔻,D)`$ with $`f(0)=0`$, $`f^{}(0)>0`$. Analogous to Section 4.2, by rotating<sup>2</sup><sup>2</sup>2Both the excursion Poisson kernel for $`𝔻`$ and excursion measure in $`𝔻`$ are rotationally invariant. if necessary, it is possible to find $`0\theta _1<\theta _2<\theta _3<\theta _4<2\pi `$ such that $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, are the images under $`f`$ of $`\mathrm{\Gamma }_𝔻`$, $`\mathrm{{\rm Y}}_𝔻`$, respectively, where $`\mathrm{\Gamma }_𝔻:=\{e^{i\theta }:\theta _1<\theta <\theta _2\}`$ and $`\mathrm{{\rm Y}}_𝔻:=\{e^{i\theta ^{}}:\theta _3<\theta ^{}<\theta _4\}`$. Define the length of $`\mathrm{\Gamma }`$, written $`\mathrm{}_\mathrm{\Gamma }`$, to be length of $`\mathrm{\Gamma }_𝔻`$ so that $`\mathrm{}_\mathrm{\Gamma }:=\theta _2\theta _1`$. Similarly define $`\mathrm{}_\mathrm{{\rm Y}}:=\theta _4\theta _3`$. Note that our notion of length is simply harmonic measure so that while $`\mathrm{\Gamma }`$ may not even be rectifiable, $`\mathrm{}_\mathrm{\Gamma }`$ always exists. An easy estimate shows that if $`(\theta _3\theta _2)`$, $`(\theta _4\theta _1)`$ are fixed, then
$$\frac{1\mathrm{cos}(\theta _3\theta _2)}{1\mathrm{cos}(\theta _4\theta _1)}=1+O(\theta _4\theta _3)+O(\theta _2\theta _1)$$
as $`(\theta _4\theta _3)0`$, $`(\theta _2\theta _1)0`$, and hence, as $`\mathrm{}_\mathrm{{\rm Y}}0`$, $`\mathrm{}_\mathrm{\Gamma }0`$,
$$\frac{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}=1+O(\mathrm{}_\mathrm{{\rm Y}})+O(\mathrm{}_\mathrm{\Gamma }).$$
(27)
Thus, we have proved the following lemma.
###### Lemma 4.24.
If $`D𝒟^{}`$, then for any $`\eta >0`$ there exist open boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}=\mathrm{}`$ such that
$$1\frac{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}1+\eta .$$
Note that the lower bound holds automatically by the definitions of separation and spread.
###### Theorem 4.25.
For every $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and for every pair of open boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$, if $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$, and $`\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_N`$ are the discrete approximations to $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, respectively, as in Section 4.2, then $`4h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$.
###### Proof.
Consider $`D𝒟^{}`$, and let $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ be (open) boundary arcs with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$. Find $`M`$ so that $`\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}})\epsilon _N:=N^{1/48}\mathrm{log}^{2/3}N`$ for $`NM`$. Throughout this section, let $`NM`$. Let $`D_N`$ be the $`1/N`$-scale discrete approximation to $`D`$ with associated “union of squares” domain $`\stackrel{~}{D}_N`$, and let $`\stackrel{~}{\mathrm{\Gamma }}_N`$, $`\stackrel{~}{\mathrm{{\rm Y}}}_N\stackrel{~}{D}_N`$ with associated discrete boundary arcs $`\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_ND_N`$. From the definitions of separation and spread, and from Corollary 4.15, since $`\mathrm{\Gamma }`$ and $`\mathrm{{\rm Y}}`$ are fixed so that $`\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}})=O(1)`$, it follows that
$$\frac{(\pi /2)h_{D_N}(0,x)h_{D_N}(0,y)}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}[1+O(\epsilon _N^3)]h_{D_N}(x,y)\frac{(\pi /2)h_{D_N}(0,x)h_{D_N}(0,y)}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}[1+O(\epsilon _N^3)].$$
Summing over all $`x\mathrm{\Gamma }_N`$ and all $`y\mathrm{{\rm Y}}_N`$ yields
$$\frac{h_{D_N}(0,\mathrm{\Gamma }_N)h_{D_N}(0,\mathrm{{\rm Y}}_N)}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}[1+O(\epsilon _N^3)]\frac{2}{\pi }h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)\frac{h_{D_N}(0,\mathrm{\Gamma }_N)h_{D_N}(0,\mathrm{{\rm Y}}_N)}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}[1+O(\epsilon _N^3)].$$
where we write $`h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N):=_{x\mathrm{\Gamma }_N}_{y\mathrm{{\rm Y}}_N}h_{D_N}(x,y)`$ and similarly for $`h_{D_N}(0,\mathrm{\Gamma }_N)`$ and $`h_{D_N}(0,\mathrm{{\rm Y}}_N)`$. However, from Proposition 4.22,
$$h_{D_N}(0,\mathrm{\Gamma }_N)h_{D_N}(0,\mathrm{{\rm Y}}_N)=H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{\Gamma }}_N)H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{{\rm Y}}}_N)+O(\delta _N),$$
where $`\delta _N:=N^{7/8}\mathrm{log}N`$, so that we conclude
$`\left[{\displaystyle \frac{H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{\Gamma }}_N)H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{{\rm Y}}}_N)}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}+O(\delta _N)\right][1+O(\epsilon _N^3)]`$
$`{\displaystyle \frac{2}{\pi }}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)\left[{\displaystyle \frac{H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{\Gamma }}_N)H_{\stackrel{~}{D}_N}(0,\stackrel{~}{\mathrm{{\rm Y}}}_N)}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}+O(\delta _N)\right][1+O(\epsilon _N^3)].`$
Now, as we let $`N\mathrm{}`$, it follows that
$`{\displaystyle \frac{H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}{\displaystyle \frac{2}{\pi }}\underset{N\mathrm{}}{lim\; inf}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$
$`{\displaystyle \frac{2}{\pi }}\underset{N\mathrm{}}{lim\; sup}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N){\displaystyle \frac{H_D(0,\mathrm{\Gamma })H_D(0,\mathrm{{\rm Y}})}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}.`$
However, Lemma 4.23 implies that
$`{\displaystyle \frac{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})4\underset{N\mathrm{}}{lim\; inf}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$
$`4\underset{N\mathrm{}}{lim\; sup}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N){\displaystyle \frac{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma },\mathrm{{\rm Y}}))}}H_D(\mathrm{\Gamma },\mathrm{{\rm Y}}).`$
For any $`\eta >0`$, let $`\{\mathrm{\Gamma }_i\}`$, $`\{\mathrm{{\rm Y}}_j\}`$ be finite partitions of $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, respectively, with
$$1\frac{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}1+\eta .$$
Note that such a partitioning is possible by Lemma 4.24. Hence, the equation above becomes
$`{\displaystyle \frac{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}}H_D(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j)4\underset{N\mathrm{}}{lim\; inf}h_{D_N}(\mathrm{\Gamma }_{N,i},\mathrm{{\rm Y}}_{N,j})`$
$`4\underset{N\mathrm{}}{lim\; sup}h_{D_N}(\mathrm{\Gamma }_{N,i},\mathrm{{\rm Y}}_{N,j}){\displaystyle \frac{1\mathrm{cos}(\mathrm{spr}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}{1\mathrm{cos}(\mathrm{sep}(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j))}}H_D(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j).`$
Summing over $`i`$ and $`j`$ and noting that
$$\underset{i}{}\underset{j}{}H_D(\mathrm{\Gamma }_i,\mathrm{{\rm Y}}_j)=H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})\text{ and }\underset{i}{}\underset{j}{}h_{D_N}(\mathrm{\Gamma }_{N,i},\mathrm{{\rm Y}}_{N,j})=h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)$$
since $`\{\mathrm{\Gamma }_{N,i}\}`$, $`\{\mathrm{{\rm Y}}_{N,j}\}`$ partition $`\{\mathrm{\Gamma }_N\}`$, $`\{\mathrm{{\rm Y}}_N\}`$, respectively, gives
$`(1+\eta )^1H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})4\underset{N\mathrm{}}{lim\; inf}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$
$`4\underset{N\mathrm{}}{lim\; sup}h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)(1+\eta )H_D(\mathrm{\Gamma },\mathrm{{\rm Y}}).`$
Since $`\eta >0`$ was arbitrary, we conclude $`4h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ as $`N\mathrm{}`$. ∎
#### 4.5.3 Convergence of $`\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)`$ to $`\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}})`$
We now prove Theorem 4.18 (b) via a result which basically says that an excursion in $`D`$ can be thought of as an excursion in $`\stackrel{~}{D}_N`$ with Brownian tails.
###### Theorem 4.26.
For every $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and for every pair of open boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$,
$$\underset{N\mathrm{}}{lim}\mathrm{}(\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N),\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}))=0$$
(28)
where $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$ with associated domain $`\stackrel{~}{D}_N𝒟`$, and corresponding boundary arcs $`\stackrel{~}{\mathrm{\Gamma }}_N`$, $`\stackrel{~}{\mathrm{{\rm Y}}}_N\stackrel{~}{D}_N`$ as in Section 4.2.
By conformal invariance, we can define excursion measure $`\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ to be either the measure $`f\mu _𝔻^\mathrm{\#}(\mathrm{\Gamma }_𝔻,\mathrm{{\rm Y}}_𝔻)`$ for $`f𝒯(𝔻,D)`$, or $`\mu _D`$ restricted to those excursions $`\gamma 𝒦_\mathrm{\Gamma }^\mathrm{{\rm Y}}(D)`$ (and normalized by $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$). Also using conformal invariance, we have $`\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)=f_N\mu _𝔻^\mathrm{\#}(\mathrm{\Gamma }_𝔻,\mathrm{{\rm Y}}_𝔻)`$ for $`f_N𝒯(𝔻,\stackrel{~}{D}_N)`$, so that we conclude
$$\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)=(f_Nf^1)\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}).$$
(29)
In order to show the convergence of the masses $`h_{D_N}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$ to $`H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$, the intermediate step of showing
$$\underset{N\mathrm{}}{lim}H_{\stackrel{~}{D}_N}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)=H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})$$
(30)
is unnecessary as a consequence of the conformal invariance of the excursion Poisson kernel: $`H_{\stackrel{~}{D}_N}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)=H_D(\mathrm{\Gamma },\mathrm{{\rm Y}})`$. However, in contrast to the excursion Poisson kernel, it is not simply a matter of applying the conformal invariance of excursion measure to conclude that (cf. Lemma 4.5)
$$\mathrm{}((f_Nf^1)\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}),\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}))0.$$
(31)
Suppose that $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, and associated “union of squares” domain $`\stackrel{~}{D}_N`$. As mentioned in Lemma 4.7, if $`z\stackrel{~}{D}_N`$, then $`\mathrm{dist}(z,D)2\sqrt{2}N^1`$. It follows from the Beurling estimates (see \[17, Proposition 3.79\] and \[20, Lemma 5.3\]) that Brownian motion started at $`z`$ is likely to exit $`D`$ quickly and nearby; that is,
$$^z\{\mathrm{diam}B[0,T_D]N^{1/2}\}CN^{1/4}\text{ and }^z\{T_DN^{1/2}\}CN^{3/8}.$$
(32)
Unfortunately, if $`z\stackrel{~}{\mathrm{\Gamma }}_N`$, it may be extremely unlikely that $`\{B_{T_D}\mathrm{\Gamma }\}`$. This will be the case, for example, if $`z`$ and $`\mathrm{\Gamma }`$ are on opposite sides of a “channel” (or “fjord”). However, since $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$ by Theorem 4.9, for *fixed* $`D𝒟^{}`$, fixed disjoint open boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}`$, and for every $`\epsilon >0`$, there exists an $`N`$ such that $`\mathrm{max}\{\mathrm{dist}(\stackrel{~}{\mathrm{\Gamma }}_N,\mathrm{\Gamma }),\mathrm{dist}(\stackrel{~}{\mathrm{{\rm Y}}}_N,\mathrm{{\rm Y}})\}<\epsilon `$. The following is then a consequence of (32) and easy bounds on the Poisson kernel.
###### Lemma 4.27.
For every $`\epsilon >0`$, there exists an $`N`$ such that for all $`z\stackrel{~}{\mathrm{\Gamma }}_N`$,
$$^z\{T_D\epsilon \text{ or }\mathrm{diam}B[0,T_D]\epsilon \text{ or }B_{T_D}\mathrm{\Gamma }_\epsilon \}\epsilon $$
(33)
where $`\mathrm{\Gamma }_\epsilon :=\{zD:\mathrm{dist}(z,\mathrm{\Gamma })\epsilon \}`$.
###### Proof of Theorem 4.26.
Suppose that $`\gamma :[0,t_\gamma ]`$ is a $`(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)`$-excursion in $`\stackrel{~}{D}_N`$. Let $`b_2:[0,t_{b_2}]`$ be a Brownian motion started at $`\gamma (t_\gamma )`$ and stopped at $`t_{b_2}:=inf\{t:b_2(t)D\}`$, its hitting time of $`D`$. Let $`b^{}:[0,t_b^{}]`$ be an independent Brownian motion started at $`\gamma (0)`$, stopped at $`t_b^{}:=inf\{t:b^{}(t)D\}`$, and set $`b_1(t):=b^{}(t_b^{}t)`$. If $`\zeta :=b_1\gamma b_2`$, then by construction $`\zeta :[0,t_\zeta ]`$ has $`\zeta (0)D`$, $`\zeta (t_\zeta )D`$, $`0<t_\zeta <\mathrm{}`$, and $`\zeta (0,t_\zeta )D`$. In other words, $`\zeta `$ is an excursion in $`D`$. Unfortunately, $`\zeta `$ is not necessarily a $`(\mathrm{\Gamma },\mathrm{{\rm Y}})`$-excursion in $`D`$, but with high probability is very close to one. Indeed, if we denote by $`\nu _{\stackrel{~}{D}_N}(\mathrm{\Gamma },\mathrm{{\rm Y}})`$ the probability measure on paths obtained by this *$`(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)`$-excursion in $`\stackrel{~}{D}_N`$ plus Brownian tails* procedure, then it follows from (33) and Proposition 3.13 that for every $`\epsilon >0`$ there exists an $`N`$ such that
$$\{𝕕(\zeta ,\gamma )\epsilon \}\epsilon \text{ and therefore }\mathrm{}(\nu _{\stackrel{~}{D}_N}(\mathrm{\Gamma },\mathrm{{\rm Y}}),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N))\epsilon .$$
The proof is completed by noting that $`\mathrm{}(\nu _{\stackrel{~}{D}_N}(\mathrm{\Gamma },\mathrm{{\rm Y}}),\mu _D^\mathrm{\#}(\mathrm{\Gamma },\mathrm{{\rm Y}}))0`$ as a consequence of Proposition 3.29: $`(\mathrm{\Gamma },\mathrm{{\rm Y}})`$-Brownian excursions in $`D`$ are generated by starting $`\epsilon `$ from $`\mathrm{\Gamma }`$ inside $`D`$ and conditioning the Brownian motion to exit $`D`$ at $`\mathrm{{\rm Y}}`$. ∎
As in the discussion preceding Theorem 1.1, we can use (28) and (30) to define the convergence of the infinite measures $`\mu _{\stackrel{~}{D}_N}`$ to $`\mu _D`$.
###### Theorem 4.28.
If $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, then $`\mathrm{}(\mu _{\stackrel{~}{D}_N},\mu _D)0`$ where $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$ with associated domain $`\stackrel{~}{D}_N`$.
It must be noted, however, that by Proposition 3.31 and Definition 3.32, we *define* excursion measure $`\mu _D`$ for $`DD^{}`$ by conformal invariance. Let $`f_N𝒯(𝔻,\stackrel{~}{D}_N)`$ as above, and also suppose that $`f𝒯(𝔻,D)`$. Hence, $`\mu _{\stackrel{~}{D}_N}:=f_N\mu _𝔻`$ and $`\mu _D:=f\mu _𝔻`$ so that $`\mu _{\stackrel{~}{D}_N}=(f_Nf^1)\mu _D`$ as in (29). Thus, we can rephrase the conclusion of Theorem 4.28 as $`\mathrm{}((f_Nf^1)\mu _D,\mu _D)0`$; compare this with (31).
#### 4.5.4 Estimating $`\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N))`$
In this section we establish Theorem 4.18 (c) by proving the following result.
###### Theorem 4.29.
For every $`D𝒟^{}`$ with $`\mathrm{inrad}(D)=1`$, for every pair of open boundary arcs $`\mathrm{\Gamma }`$, $`\mathrm{{\rm Y}}D`$ with $`\overline{\mathrm{\Gamma }}\overline{\mathrm{{\rm Y}}}\mathrm{}`$, and for every $`\epsilon >0`$, there exists an $`N`$ such that
$$\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N))\epsilon $$
(34)
where $`D_N`$ is the $`1/N`$-scale discrete approximation to $`D`$ with associated domain $`\stackrel{~}{D}_N𝒟`$ and corresponding boundary arcs $`\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_ND_N`$; $`\stackrel{~}{\mathrm{\Gamma }}_N`$, $`\stackrel{~}{\mathrm{{\rm Y}}}_N\stackrel{~}{D}_N`$ as in Section 4.2.
In order to prove (34), it will be necessary to use the strong approximation of Proposition 4.21. Hence, let $`A_N:=2ND_N`$ so that $`A_N𝒜^N`$, and write $`\mathrm{\Gamma }_{N,A}:=2N\mathrm{\Gamma }_N`$, $`\mathrm{{\rm Y}}_{N,A}:=2N\mathrm{{\rm Y}}_NA_N`$ for the corresponding boundary arcs. Suppose further that $`N`$ is chosen large enough so that $`\mathrm{dist}(\mathrm{\Gamma }_{N,A},\mathrm{{\rm Y}}_{N,A})N^{15/16}`$. Since $`D𝒟^{}`$, it follows that $`A_N`$ is necessarily bounded so that $`\mathrm{rad}(A_N)\mathrm{inrad}(A_N)N`$, and furthermore, $`|\mathrm{{\rm Y}}_{N,A}||\mathrm{\Gamma }_{N,A}|N`$ where all of the constants may depend on $`D`$.
###### Proof of Theorem 4.29.
Suppose that $`xA_N^{}:=\{xA_N:g_{A_N}(x)N^{1/16}\}`$, and let $`S`$ be a simple random walk with $`S_0=x`$. As in the proof of \[15, Corollary 3.5\], it follows from the Beurling estimate that $`\mathrm{dist}(x,A)CN^{7/8}`$. Hence, a straightforward gambler’s ruin estimate shows that $`^x\{S_\tau \mathrm{{\rm Y}}_{N,A}\}N^{1/16}`$ where $`\tau =\tau _{A_N}:=\mathrm{min}\{j:S_jA\}`$. The coupling of Brownian motion and random walk provided by Corollary 4.20 is so strong that even conditioning on the rare event $`\{S_\tau \mathrm{{\rm Y}}_{N,A}\}`$ does not uncouple the processes. Hence, there exists a Brownian motion $`B`$, a simple random walk $`S`$ with $`B_0=S_0=x`$, and a constant $`C`$ such that
$$^x\left\{\underset{0t\tau }{sup}\left|\frac{1}{\sqrt{2}}B_tS_t\right|C\mathrm{log}N|S_\tau \mathrm{{\rm Y}}_{N,A}\right\}CN^8.$$
(35)
The strong approximation (Proposition 4.21) allows us to conclude that conditioned on the event $`\{S_\tau \mathrm{{\rm Y}}_{N,A}\}`$, Brownian motion and simple random walk starting $`N^{7/8}`$ away from the boundary still exit near each other; that is,
$$^x\left\{|B_TS_\tau |CN^{1/4}\mathrm{log}N|S_\tau \mathrm{{\rm Y}}_{N,A}\right\}CN^{1/16}$$
(36)
where $`T=T_{A_N}:=inf\{t:B_t\stackrel{~}{A}_N\}`$. The time version of the Beurling estimate \[20, Lemma 5.3\] says that $`^x\{|T\tau |r^2\mathrm{dist}(x,\stackrel{~}{A})^2\}Cr^{1/2}`$. Hence,
$$^x\left\{|T\tau |CN^{1/2}\mathrm{log}^2N|S_\tau \mathrm{{\rm Y}}_{N,A}\right\}CN^{1/16}.$$
(37)
We can now use Proposition 3.13 to deduce statements about convergence in $`\mathrm{}`$ from statements about convergence in $`𝕕`$. In particular, let $`\gamma :[0,t_\gamma ]`$ be given by $`t_\gamma :=T`$, $`\gamma (t):=B_t`$, $`0tt_\gamma `$, and associate to the random walk $`S`$ the curve $`\stackrel{~}{\omega }:[0,t_{\stackrel{~}{\omega }}]`$ as in (19), so that from (35), (36), and (37), we conclude that $`\{𝕕(\gamma ,\stackrel{~}{\omega })CN^{1/2}\mathrm{log}^2N\}CN^{1/16}`$, and using Lemma 4.16, we can scale our results to $`D_N`$:
$$\left\{𝕕(\mathrm{\Phi }_N\gamma ,\mathrm{\Phi }_N\stackrel{~}{\omega })CN^{1/2}\mathrm{log}^2N\right\}\left\{𝕕(\gamma ,\stackrel{~}{\omega })CN^{1/2}\mathrm{log}^2N\right\}N^{1/16}$$
(38)
where $`\mathrm{\Phi }_N:=\mathrm{\Psi }_{1/(2N)}`$ is the Brownian scaling map as in (25). Let $`V_{N,A}`$ be the set $`V_{N,A}:=\{xA_N:\mathrm{dist}(x,\mathrm{{\rm Y}}_{N,A})CN^{1/4}\mathrm{log}N\}`$, let $`\stackrel{~}{V}_{N,A}`$ be the associated subset of $`\stackrel{~}{A}_N`$, and let $`2N\stackrel{~}{V}_N=\stackrel{~}{V}_{N,A}`$. It then follows that $`(\mathrm{\Phi }_N\stackrel{~}{\omega })=\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(x,\mathrm{{\rm Y}}_N)`$ and $`(\mathrm{\Phi }_N\gamma )=\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(x,\stackrel{~}{V}_N)`$. Since $`N^{1/2}\mathrm{log}NN^{1/16}`$, Proposition 3.13 and (38) yield
$$\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(x,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(x,\stackrel{~}{V}_N))CN^{1/16}.$$
(39)
As in the proof of Proposition 4.22, $`H_{\stackrel{~}{D}_N}(x,\stackrel{~}{V}_N)=H_{\stackrel{~}{D}_N}(x,\stackrel{~}{\mathrm{{\rm Y}}}_N)+O(N^{3/4}\mathrm{log}N)`$, so it follows that
$$\mathrm{}(\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(x,\stackrel{~}{\mathrm{{\rm Y}}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(x,\stackrel{~}{V}_N))CN^{3/4}\mathrm{log}N.$$
(40)
Combining (39) and (40) then yields $`\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(x,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(x,\stackrel{~}{\mathrm{{\rm Y}}}_N))CN^{1/16}`$, and, in particular, if $`y𝒜_N^{}`$ with $`|xy|C\mathrm{log}N`$, then
$$\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(x,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(y,\stackrel{~}{\mathrm{{\rm Y}}}_N))CN^{1/16}.$$
(41)
To complete the proof, suppose that $`S^{}`$ is a simple random walk on the scaled lattice $`\frac{1}{2N}^2`$, and let $`D_N^{}:=\frac{1}{2N}A_N^{}`$ so that $`D_N^{}=\{zD_N:g_{D_N}(z)N^{1/16}\}`$ by (23) where $`g_{D_N}`$ is the Green’s function for Brownian motion on $`\stackrel{~}{D}_N`$. Also recall from Theorem 4.9 that $`\stackrel{~}{D}_N\stackrel{𝖼𝖺𝗋𝖺}{}D`$. Hence, if $`\eta _N=\eta (D,N):=\mathrm{min}\{j0:S_j^{}D_N^{}D_N^c\}`$ and $`xD_ND_N^{}`$, then it follows from \[15, Lemma 3.11\] that for every $`\epsilon >0`$, there exists an $`N`$ such that
$$^x\left\{\eta _N\epsilon \right|S_{\eta _N}^{}D_N^{}\}\epsilon .$$
(42)
Furthermore, using \[15, Lemma 3.11\] again, we can find constants $`C`$, $`\alpha `$ such that
$$^x\left\{\underset{0j\eta 1}{\mathrm{max}}|f_{D_N}(S_j^{})f_{D_N}(x)|N^{1/16}\mathrm{log}N\right\}CN^\alpha ,$$
(43)
and
$$^x\left\{|f_{D_N}(S_\eta ^{})f_{D_N}(x)|N^{1/16}\mathrm{log}N|S_\eta ^{}D_N^{}\right\}CN^\alpha .$$
(44)
Suppose further that $`\stackrel{~}{B}`$ is a Brownian motion started at $`xD_ND_N^{}`$. As in Lemma 4.27, if $`\stackrel{~}{\eta }_N=\stackrel{~}{\eta }(D,N):=inf\{t0:\stackrel{~}{B}_t\stackrel{~}{D}_N^{}\stackrel{~}{D}_N^c\}`$, then for every $`\epsilon >0`$, there exists an $`N`$ such that
$$^x\left\{\stackrel{~}{\eta }_N\epsilon \text{ or }\mathrm{diam}B[0,\stackrel{~}{\eta _N}]\epsilon \right|B_{\stackrel{~}{\eta }_N}\stackrel{~}{D}_N^{}\}\epsilon .$$
(45)
If we let $`\stackrel{~}{\gamma }:[0,t_{\stackrel{~}{\gamma }}]`$ be given by $`t_{\stackrel{~}{\gamma }}:=\stackrel{~}{\eta }_N`$, $`\stackrel{~}{\gamma }(t):=\stackrel{~}{B}_t`$, $`0tt_{\stackrel{~}{\gamma }}`$, and associate to the (scaled) random walk $`S^{}`$ the (scaled) curve $`\stackrel{~}{\omega }^{}:[0,t_{\stackrel{~}{\omega }^{}}]`$ as in (25) (i.e., Brownian scaled in both time and space), then letting $`\underset{\stackrel{~}{}}{\gamma }:=\stackrel{~}{\gamma }\mathrm{\Phi }_N\gamma `$ and $`\underset{\stackrel{~}{}}{\omega }:=\stackrel{~}{\omega }^{}\mathrm{\Phi }_N\stackrel{~}{\omega }`$ we see that $`(\underset{\stackrel{~}{}}{\gamma })=\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N)`$ and $`(\underset{\stackrel{~}{}}{\omega })=\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N)`$. Hence, by combining (42), (43), (44), and (45) with (41), we conclude that for every $`\epsilon >0`$, there exists an $`N`$ with
$$\mathrm{}(\mu _{D_N}^{\mathrm{𝗋𝗐},\mathrm{\#}}(\mathrm{\Gamma }_N,\mathrm{{\rm Y}}_N),\mu _{\stackrel{~}{D}_N}^\mathrm{\#}(\stackrel{~}{\mathrm{\Gamma }}_N,\stackrel{~}{\mathrm{{\rm Y}}}_N))\epsilon .\mathit{}$$
## Acknowledgements
Much of this research was done by the author in his Ph.D. dissertation under the supervision of Greg Lawler. The author wishes to express his gratitude to Prof. Lawler for his continued guidance and support. Thanks are also due to Christophe Garban for valuable comments, to the Banff International Research Station for Mathematical Innovation and Discovery where the final writing of this paper was done, and to Christian G. Beneš and José A. Trujillo Ferreras for many fruitful discussions.
|
warning/0506/cond-mat0506164.html
|
ar5iv
|
text
|
# The cavity method for large deviations
## I Introduction
The algorithmic complexity of a problem is traditionally measured on an ensemble of possible inputs (instances) by retaining the largest time it takes for an algorithm to solve one of the instances Papadimitriou and Steiglitz (1982). Statistical physics studies have however suggested a different characterization of hardness, based on the average case rather than the worst case Mézard et al. (1987). This alternative approach is motivated by a generic phenomenon of concentration, according to which a particular instance behaves almost surely as the average case in the limit of infinite size. This self-averaging property is common to many disordered systems whose environment is specified by a quenched random variable from a prescribed ensemble : in the thermodynamical limit, the properties of a sample tend to be independent of the particular realization of the disorder. Due to this parallel, the methods first developed for physical disordered systems have been successfully applied to combinatorial optimization problems Mézard et al. (1987). The interest in typical properties is however far from limited to physics or optimization ; information Cover and Thomas (1991) and graph theories Bollobás (2001) are two other major fields where they play a key role, as illustrated by the seminal works of Shannon Shannon (1948) and Erdős-Rényi Erdős and Rényi (1960) respectively.
In any practical implementation however, optimization or coding theories face large but finite system sizes. In such situations, controlling the deviations from the typical case becomes of primarily practical interest. The scope of large deviations theory Dembo and Zeitouni (1993) is precisely to evaluate the probability of the rare events associated with such finite size effects. In addition, and despite the tremendous number of their elements, large deviations theory has also a direct relevance for physical systems : as will be explained, it underlies the thermodynamics of systems whose configuration space is constituted of different realizations of the quenched disorder. Consequently, when the disorder consists of an ensemble of random graphs, it allows to solve models with variable topologies. A large deviations analysis is thus useful to address the adaptability of constrained systems, with a physical example being covalent molecular networks subject to stress Barré et al. (2004a); Rivoire (2005a).
A particularly powerful method for computing typical properties of disordered systems is the celebrated replica method Mézard et al. (1987). The cavity method Mézard et al. (1987) provides an alternative tool that yields equivalent results but has two advantages over the replica method : it is based on assumptions formulated explicitly, and it applies to particular instances. In the present paper, we develop an approach of large deviations based on the cavity method, which we call the large deviations cavity method (LDCM). If large deviations have only recently raised an interest in the statistical physics community Derrida et al. (2002); Montanari and Zecchina (2002); Andreanov et al. (2004); Barré et al. (2004b); Engel et al. (2004); Rivoire (2004), they have a much longer history in probability theory Dembo and Zeitouni (1993), where they are notably used to rigorously solve statistical mechanics models Ellis (1985). At variance with this mathematical tradition, the method exposed here is non-rigorous and only provides a coherent heuristic framework for obtaining quantitative predictions, namely computing rate functions assuming a large deviations principle indeed holds. Nonetheless, as for the ”typical” cavity method of Mézard and Parisi Mézard and Parisi (2001, 2003) that we will recover as a particular case, the LDCM is hoped to be amenable to rigorous studies.
The paper is organized as follows. The first section is devoted to introducing some basic elements of combinatorial optimization and large deviations theories, with an emphasis on their links with statistical mechanics. The second section presents the LDCM in its simplest, ”replica symmetric”, form : we start by rederiving in a cavity-like fashion Cramér’s theorem, the most elementary result in large deviations theory, and then discuss different graph ensembles, with explicit calculations on the vertex-cover problem. The third section deals with systems having a non trivial internal structure : we notably generalize the large deviations approach to systems displaying a glassy ”replica symmetry breaking” phase, a situation that we illustrate in details with the coloring problem. A conclusion closes the paper by suggesting some possible applications.
## II Optimization problems and large deviations
### II.1 Optimization problems
The field of combinatorial optimization provides a broad class of disordered systems, which we use here to illustrate the potentialities of the LDCM. We therefore start with a brief introduction to this subject (see e.g. Papadimitriou and Steiglitz (1982) for more details). Combinatorial optimization is primarily concerned with minimizing cost functions, $`:𝒞`$, over some discrete configuration space $`𝒞`$. In view of quantifying their algorithmic complexity, optimization problems are defined over an ensemble $``$ of instances $`I`$, each associated with a cost function $`_I`$ . In particular, many combinatorial problems are defined over ensemble of graphs Garey and Johnson (1979), in which case an instance $`I`$ is specified by a graph $`G`$, that is a set of $`N`$ nodes $`i=1,\mathrm{},N`$ associated with a subset of the pairs $`\{(i,j)\}_{ij}`$, defining its edges.
Two prototypical examples will serve as illustration. The first one is the vertex-cover problem Garey and Johnson (1979), also known as independent set, which consists, given a graph $`G`$, in blackening as many of its nodes as possible while never blackening two connected nodes. The second, quite similar, example is the coloring problem Garey and Johnson (1979) which asks, given a graph $`G`$ and $`q`$ colors, whether it is possible to assign a color to each node of $`G`$ so that no two adjacent nodes have same color. The coloring problem is strictly speaking a decision problem (the answer must be yes or no) but it is directly related to the optimization problem of minimizing the number of edges having two end-nodes sharing a same color : if the minimum is zero, the graph is colorable, otherwise it is uncolorable.
From the statistical physics viewpoint, the cost $`_I[\sigma ]`$ represents the energy of a configuration $`\sigma 𝒞`$, and the minimal cost $`E_I=\mathrm{min}_\sigma _I[\sigma ]`$ corresponds to the ground-state energy of the disordered system having quenched disorder $`I`$. In this context, it is usual to introduce an inverse temperature $`\beta `$ and a free energy density $`f_I(\beta )`$ defined by $`f_I(\beta )=\mathrm{ln}[_\sigma \mathrm{exp}(\beta _I[\sigma ])]/(\beta N)`$, such that the ground-state energy density $`ϵ_I=E_I/N`$ is given by the $`\beta \mathrm{}`$ limit, $`ϵ_I=lim_\beta \mathrm{}f_I(\beta )`$. For the coloring problem, the associated finite temperature system is known as the antiferromagnetic Potts model Wu (1982), while it is called the hard-core model Weigt and Hartmann (2001) for vertex-cover (with $`\beta `$ representing a chemical potential).
For these two examples, the configuration space $`𝒞`$ is discrete : $`𝒞`$ can be taken as $`\{0,1\}^N`$ for vertex-cover, with $`\sigma _i=0`$ and $`\sigma _i=1`$ corresponding respectively to uncovered (white) and covered (black) nodes, and as $`\{1,\mathrm{},q\}^N`$ for coloring with $`\sigma _i\{1,\mathrm{},q\}`$ now representing the color assigned to $`i`$. Different ensembles $`𝒢`$ of graphs, defining different sets of instances $``$, can be introduced. The cavity approach followed here applies to any ensemble of locally ”tree-like” graphs, that is graphs whose degrees (the number of nodes to which a given node is connected) remain finite when $`N\mathrm{}`$. Three random graphs ensembles will be specifically addressed here. The first one, noted $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, is the set of graphs with $`N`$ nodes where each edge as a probability $`\gamma /N`$ to be present, and is known as the binomial model or Erdős-Rényi ensemble Bollobás (2001). The second one, noted $`\overline{𝒢}_N^{(\gamma )}`$, and called the uniform model Bollobás (2001), is the set of graphs with $`N`$ nodes and $`M=\gamma N/2`$ edges. Finally, the third one is defined through the degree distribution $`p(k)`$ of the nodes of its graphs Newman et al. (2001) : each of the $`N`$ nodes has degree $`k`$ with independent probability $`p(k)`$ and the edges are drawn at random subject to that constraint. This last class notably includes random regular graphs Bollobás (2001), for which $`p(k)=\delta _{r,k}`$, and power-law distributed graphs, for which $`p(k)k^\tau `$ (with an appropriate cut-off to insure normalization). Here we will also consider the Poissonian model noted $`\widehat{𝒢}_N^{(\gamma )}`$ and defined by $`p(k)=\gamma ^ke^\gamma /k!`$. In the $`N\mathrm{}`$ limit, the binomial and uniform models $`\stackrel{~}{𝒢}_N^{(\gamma )}`$ and $`\overline{𝒢}_N^{(\gamma )}`$ share with $`\widehat{𝒢}_N^{(\gamma )}`$ the same Poisson degree distribution. This equivalence between the three models extends to the typical properties of optimization problems defined on them but, as will be shown, does not hold for atypical features.
From the point of view of computational complexity, an important parameter is the size $`N`$ of the instances, which, in the case of diluted graphs, is taken as the number of nodes. As seen on the vertex-cover and coloring problems, the size of the configuration space $`𝒞`$ over which optimization is to be performed increases exponentially with $`N`$, precluding any naïve exhaustive search for large $`N`$ and possibly making the problem highly non-trivial. In fact, both the vertex-cover and coloring problems are known to be NP-hard in the worst case, implying that no algorithm is known that can solve all instances of these problems in a time growing polynomially with $`N`$ Papadimitriou and Steiglitz (1982).
As stressed in the introduction, the focus on typical instances advocated by statistical physics is justified by the self-averaging property : when it holds for an ensemble $``$ of instances, there is a typical value of the ground-state energy density $`\overline{ϵ}`$ such that, for any $`\epsilon >0`$, the probability $`[|E_I/N\overline{ϵ}|>\epsilon ]`$ for the optimum $`E_I`$ to deviate from $`N\overline{ϵ}`$ goes to zero, $`[|E_I/N\overline{ϵ}|>\epsilon ]0`$ as $`N\mathrm{}`$. Informally, large problems then typically all share a common optimum, which, physically, can often be traced back to the equivalence of their local properties. The ”typical” cavity method Mézard and Parisi (2001, 2003) have been developed precisely to compute the most probable value $`\overline{ϵ}`$ for problems on random graphs. The LDCM presented here is an extension of this approach that allows to evaluate the $`N`$ and $`\epsilon `$ dependencies of vanishing probabilities such as $`[|E_I/N\overline{ϵ}|>\epsilon ]`$.
### II.2 Large deviations
For finite $`N`$, an instance has always a finite probability to deviate from the typical case. The so-called large deviations Dembo and Zeitouni (1993) refer to the extensive deviations from $`N\overline{ϵ}`$, of order $`O(N)`$, as distinguished from the small, subextensive deviations from $`N\overline{ϵ}`$, of order $`o(N)`$, like for example $`O(\sqrt{N})`$ fluctuations (see however below for a relation between the two). The present method is based on an Ansatz, according to which large deviations are exponentially small in the size $`N`$ of the instances, that is, the probability $`_N[E_I]`$ for an instance $`I`$ taken out of the ensemble $``$ to have an optimal cost $`E_I`$ is supposed to satisfy
$$_N[E_I=Nϵ]e^{NL(ϵ)},$$
(1)
where the symbol $``$ stands here and in the sequel for an exponential equivalence defined as $`lim_N\mathrm{}\mathrm{ln}([E_I=Nϵ])/N=L(ϵ)`$. $`L(ϵ)`$ is called a rate function, or large deviations function, and, in the simplest cases, is strictly positive, except for the typical value $`\overline{ϵ}`$ where it achieves its zero minimum. The Ansatz (1) is known to indeed hold in the solvable case where $`_I`$ consists of a sum of independent identically distributed variables (Cramér’s theorem, see Sec. III.1), and this result is robust to the presence of weak correlations among the variables (Gätner-Ellis theorem, to be stated below) Dembo and Zeitouni (1993). More precisely, the relation (1) corresponds in the mathematical literature to the ”large deviations principle” Dembo and Zeitouni (1993), which, in its simplest form, can be stated as follows:
Large deviations principle: The sequence $`\{ϵ_N\}_N`$ of real valued random variables is said to satisfy the large deviations principle, with rate function $`L:^+\{\mathrm{}\}`$, if
$`(i)`$ $`M0,\{x:L(x)M\}\mathrm{is}\mathrm{compact}`$,
$`(ii)`$ for all closed subset $`F`$ of $``$, and all open subset $`O`$ of $``$,
$$\underset{N\mathrm{}}{lim\; sup}\frac{1}{N}_N[ϵ_NF]L(F),\underset{N\mathrm{}}{lim\; inf}\frac{1}{N}_N[ϵ_NO]L(O).$$
(2)
We point out right away that counterexamples are easily found for which the previous Ansatz does not hold. In the field of spin-glass models, they include the two most celebrated models, the random energy model Derrida (1980) and the SK model Sherrington and Kirkpatrick (1975). For the random energy model, elementary calculations Andreanov et al. (2004); Rivoire (2004) indeed give
$$_N(ϵ)\{\begin{array}{cc}e^{e^{Ns(ϵ)}}\hfill & \text{if }ϵ>\overline{ϵ},\hfill \\ e^{Ns(ϵ)}\hfill & \text{if }ϵ<\overline{ϵ},\hfill \end{array}$$
(3)
with $`s(ϵ)=\mathrm{ln}2ϵ^2`$ and $`\overline{ϵ}=\sqrt{\mathrm{ln}2}`$. For the SK model, numerical studies Andreanov et al. (2004) also suggest different scalings on both sides of the typical value, that is $`_N(ϵ)\mathrm{exp}[N^aL(ϵ)]`$ with $`a1.2`$ when $`ϵ<\overline{ϵ}`$, but $`a1.5`$ when $`ϵ>\overline{ϵ}`$. However, in a variety of other spin-glass models, notably including models on diluted random graphs, the Ansatz (1) is supported by numerical evidence Andreanov et al. (2004).
From the analytical viewpoint, rate functions in the context of optimization problems have been studied by Montanari Montanari and Zecchina (2002); Montanari (2002), using the replica method. The results he obtained for the vertex-cover problem Montanari and Zecchina (2002); Montanari (2002) are strictly equivalent to the ones to be derived here from the cavity method. Yet, as for the typical case, the cavity approach has the advantages over the replica method to offer a more transparent derivation, and to open the way to algorithmic implementations on particular systems. For a model at finite temperature $`1/\beta `$, the replica method basically consists in inferring rate functions from the knowledge of the moments $`𝔼[Z_I^n]`$ of the partition function $`Z_I(\beta )=_\sigma \mathrm{exp}(\beta _I[\sigma ])`$, with $`𝔼[]`$ referring to the average over the disorder, that is the different instances $`I`$. As far as no replica symmetry breaking is involved, this procedure is motivated by the following rigorous result Dembo and Zeitouni (1993) :
Gärtner-Ellis theorem : Let $`\{ϵ_N\}_N`$ be a sequence of real valued random variables and let $`:`$ be defined by
$$(y)=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{ln}𝔼[e^{yNϵ_N}].$$
(4)
If $`(y)`$ exists, is finite and differentiable for every $`y`$, then the sequence $`\{ϵ_N\}_N`$ satisfies the large deviations principle with rate function $`L(ϵ)`$ given by the Legendre transform of $`(y)`$,
$$L(ϵ)=\underset{y}{inf}[yϵ(y)],$$
(5)
where minus signs are introduced to match usual conventions in statistical physics. We stress that for sake of simplicity, this theorem is stated here with much stronger hypothesis than necessary ; in particular the assumptions about the finiteness and differentiability of $``$ can be relaxed Dembo and Zeitouni (1993).
To apply the replica method to optimization problems, the limit $`\beta \mathrm{}`$ has to be considered, and, to obtain non trivial results in this limit, the replica number $`n`$ must be rescaled with $`\beta `$, such that $`\beta \mathrm{}`$ and $`n0`$ with $`y=n/\beta `$ finite. In this limit, the replica potential $`(y)`$ coincides with the function introduced in Gärtner-Ellis theorem, Eq. (4),
$$e^{N(y)}=\underset{\beta \mathrm{}}{lim}𝔼[Z^{y/\beta }]=𝔼[e^{yE}]=e^{N[L(ϵ)+yϵ]}𝑑ϵ.$$
(6)
While proceeding differently, the cavity method to be presented will lead to the same rate function, again specified as the Legendre transform of the potential $`(y)`$. Although both the replica and the cavity methods, based on Legendre transformations, naturally yield convex functions, it should be stressed that convexity is not a necessary feature of rate functions. In fact, non-convex rate functions are associated with phase transitions and are therefore encountered in many models of interest from the statistical mechanics point of view Ellis (1985).
Large deviations deal with exponentially small probability and may appear as only an extreme feature of finite size effects, while a more refined description would consist in the complete probability distribution of $`E_I`$ over $``$. Interestingly, small fluctuations can be extracted from the knowledge of the rate function near its typical minimum. More precisely, the potential $`(y)`$ yields the cumulants of $`(E_N)^k_c`$,
$$(E_N)^k_c=N\frac{^k}{y^k}(y=0),$$
(7)
where, as usual, the cumulants $`X^k_c`$ of a random variable $`X`$ are defined by $`\mathrm{ln}𝔼[e^{tX}]=_{k=1}^{\mathrm{}}\frac{t^k}{k!}X^k_c`$. In particular, the Ansatz (1) predicts the variance of the small fluctuations to be generically of order $`\sqrt{N}`$, as given by the central limit theorem in the case of a sum of independent identically distributed variables.
### II.3 Statistical mechanics interpretation
On top of their own mathematical interest, large deviations are of direct relevance to statistical mechanics studies. In the context of optimization problems, rate functions can indeed be interpreted as defining an entropy on the space of the instances $``$, corresponding to a thermodynamics over the quenched disorder. This relation, formalized by Sanov’s theorem Ellis (1985), is presented here in the restricted context where $``$ is a class of graphs associated with a given optimization problem.
Viewing the ensemble of random graphs $`𝒢_N`$ as a phase space, each graph $`G𝒢_N`$ defines a configuration to which is associated the ground-state energy $`E_G`$, that is, the optimal cost for the optimization problem on $`G`$. If $`|𝒢_N|e^{Ns_0}`$ denotes the cardinality of $`𝒢_N`$, the microcanonical entropy $`s(ϵ)`$ of the system is given by
$$e^{Ns(ϵ)}=\mathrm{\#}\{G𝒢_N;E_G=Nϵ\}=\mathrm{\#}𝒢_N\times _N[E_G=Nϵ]e^{N[s_0L(ϵ)]},$$
(8)
where $`\mathrm{\#}A`$ denotes the cardinality of the set $`A`$. Thus, up to a linear transformation, the rate function $`L(ϵ)`$ is nothing but the microcanonical entropy $`s(ϵ)`$,
$$s(ϵ)=s_0L(ϵ).$$
(9)
Within this picture, the parameter $`y`$ appearing in the replica method and Gärtner-Ellis theorem represents the external inverse temperature that allows to study statistical mechanics on the configuration space spanned by the graphs, $`y_ϵL(ϵ)=_ϵs(ϵ)`$ ($`y`$ must be distinguished from the internal inverse temperature $`\beta `$ which is set to infinity in the context of optimization). By construction, this space has no more quenched disorder, and a large deviations analysis appears as the statistical mechanics analysis of a pure system at finite inverse temperature $`y`$. From the opposite viewpoint, large deviations theory thus provides a meaning for negative temperatures, $`y<0`$. Finally, the typical case is given by the infinite temperature limit, $`y=0`$, as prescribed by replica theory.
We stress however that the possibility of deriving the thermodynamics of the system at inverse temperature $`y`$ from the knowledge of its microcanonical entropy $`s(ϵ)`$ is based on the equivalence between the microcanonical and canonical ensembles in the thermodynamical limit, which can not always be taken for granted. In presence of non-convex rate functions indeed, the two ensembles become inequivalent, and a first order transition occurs, whose description requires a Maxwell construction ; such a construction in the context of large deviations for random graph has been recently described in Engel et al. (2004).
We have restricted so far to the simplest case where the measure over the quenched disorder is an uniform measure over an ensemble of graphs, but more complicated structures can be considered as well. In particular, the disorder can have different origins, as with spin-glass models Mézard et al. (1987) or $`K`$-SAT optimization problems Papadimitriou and Steiglitz (1982), where in addition to the graph structure, the quenched disorder comprises the specification of some random couplings between the variables. In this case, an instance $`I`$ of the problem is first defined by selecting a graph $`G`$ and then by choosing the couplings $`J`$. Large deviations can be taken with respect to $`J`$ at fixed $`G`$ : for typical graphs $`\overline{G}`$, the effective system still contains a quenched disorder (the graph) which can be handled with the usual techniques of disordered systems, but if atypical graphs have to be addressed as well, a second temperature needs to be introduced. The two temperatures are in such a case associated with two levels of probability distributions, in a construction formally identical to Parisi’s hierarchical scheme for handling replica symmetry breaking, as will be discussed in Sec. IV.3. The same scheme also applies when going to lower levels to describe the internal structure of a given instance. This will be exemplified in Sec. IV.1 where we discuss the implications of working with a finite temperature on the instances, or working with optimization problems displaying a replica symmetry breaking phase.
## III The large deviations cavity method
The ”typical” cavity method, as developed by Mézard and Parisi Mézard and Parisi (2001, 2003), applies to a given instance $`I`$ and addresses the structure of its phase space, that is the organization of the configurations $`\sigma 𝒞`$ as a function of their energy density $`[\sigma ]/N`$. The method can handle either a structure composed of a unique set (or a finite number of sets) of connected solutions, called a replica symmetric (RS) phase, or a structure composed of many disconnected clusters of configurations, called a one-step replica symmetry breaking (1RSB) phase Mézard et al. (1987). In the latter case, the crucial assumption is made that the number of clusters with a given energy density $`ϵ`$ is exponential in $`N`$,
$$𝒩_{\mathrm{clusters}}(ϵ)e^{N\mathrm{\Sigma }(ϵ)}.$$
(10)
The 1RSB cavity method is specifically designed to compute the function $`\mathrm{\Sigma }(ϵ)`$, called the complexity, with the particular RS case corresponding to $`\mathrm{\Sigma }=0`$ Monasson (1995). The formal analogy between the 1RSB Ansatz (10) defining the complexity $`\mathrm{\Sigma }(ϵ)`$ and the large deviations Ansatz (1) defining the rate function $`L(ϵ)`$ is at the root of the possibility to extend the typical cavity method yielding $`\mathrm{\Sigma }(ϵ)`$ to an atypical version yielding $`L(ϵ)`$. To emphasize further the parallel, we introduce the function $`(ϵ)`$ defined as $`(ϵ)L(ϵ)`$, such that $`(ϵ)`$ plays in the LDCM a role formally identical to the complexity in the typical cavity method:
$$_N(E=Nϵ)e^{N(ϵ)}.$$
(11)
The analogy between the complexity $`\mathrm{\Sigma }(ϵ)`$ and the rate function $`L(ϵ)`$ should not be taken for a coincidence: the complexity is fundamentally nothing but a rate function \[or more accurately the entropy associated to it, as in Aq. (9)\], which describes the large deviations of the energy over the different clusters of solutions. From this point of view, further elaborated in Rivoire (2005b), the 1RSB cavity method is itself a large deviations method, acting on the self-generated (glassy) ”internal disorder” of a given sample. For glassy optimization problems, being able to address such large deviations is crucial since ground-state clusters are atypical, that is, exponentially less numerous than clusters with higher energies. These atypical ground-state clusters must be obtained by correctly tuning the ”internal inverse temperature”, noted $`\mu `$ in this context. Remarkably, while the LDCM to be presented will also apply to the typical case $`y=0`$, the 1RSB cavity method is in general not able to describe the complete complexity curve $`\mathrm{\Sigma }(ϵ)`$, and notably fails to describe the most numerous, typical clusters, corresponding to $`\mu =0`$ foo (a).
Our presentation of the LDCM will follow closely the presentation of Mézard and Parisi of their typical cavity method Mézard and Parisi (2001, 2003), but major differences will show up in the way averages over the disorder are performed. To start with, we consider the simplest case where the underlying optimization problem is assumed to be itself RS i.e., with no clustering induced by its internal disorder.
### III.1 The cavity approach to Cramér’s theorem
Although its most interesting applications involve random graphs, the cavity method is not restricted to this particular geometry. As an illustration of the ideas in their simplest setting, we consider the case, with no geometry, of a system made of $`N`$ independent elements, each contributing to the total energy $`_N`$ by a random amount $`X_i`$. In other words, we consider here large deviations in the sum of independent identically distributed random variables. For such a system, the typical energy density follows from the law of large numbers, which, assuming the distribution $`\rho (X)`$ of the $`X_i`$’s to have a finite first moment, is $`\overline{ϵ}=𝔼[X]x\rho (x)𝑑x`$. Large deviations are concerned with deviations from the prediction $`ϵ_N/N=\overline{ϵ}`$ and, for a sum of independent variables, are completely specified by Cramér’s theorem, both a generalization of the law of large numbers and a corollary of Gärtner-Ellis theorem Dembo and Zeitouni (1993).
Cramér’s theorem : Let the sequence $`\{ϵ_N\}_N`$ of real random variables be given by $`ϵ_N=(_{i=1}^NX_i)/N`$ where the $`\{X_i\}_i`$ are independently identically distributed real random variables. If $`𝔼[e^{yX}]`$ is finite for all $`y`$, then $`\{ϵ_N\}_N`$ satisfies the large deviations principle with rate function $`L:`$ defined as Legendre transform of $`:`$ given by $`(y)\mathrm{ln}𝔼[e^{yX}]`$, that is
$$L(ϵ)(ϵ)=\underset{y}{inf}[yϵ(y)].$$
(12)
The basic idea behind the cavity approach is to estimate the change of the system upon addition of a new variable (or, equivalently, upon removal of a variable, hence the name ”cavity”). Let $`_N`$ be the extensive energy, $`_N=_{i=1}^NX_i`$. By virtue of the assumed independence of the $`X_i`$, the probability distribution for $`_{N+1}=_N+X_{N+1}`$ is given by a convolution of those of $`_N`$ and $`X_{N+1}`$, which, with the Ansatz (11), reads
$$_{N+1}(_{N+1}=E)=e^{(N+1)\left(\frac{E}{N+1}\right)}=𝔼_X[_N(_N=EX)]=\rho (\mathrm{\Delta }E)e^{N\left(\frac{E\mathrm{\Delta }E}{N}\right)}𝑑\mathrm{\Delta }E.$$
(13)
Assuming a smooth behavior of $``$, we write for large $`N`$,
$$\begin{array}{cc}& (N+1)\left(\frac{E}{N+1}\right)=N(ϵ)+(ϵ)_ϵ(ϵ)+O(1),\hfill \\ & N\left(\frac{E\mathrm{\Delta }E}{N}\right)=N(ϵ)\mathrm{\Delta }E_ϵ(ϵ)+O(1),\hfill \end{array}$$
(14)
where $`ϵE/N`$. Setting $`y_ϵ(ϵ)`$ thus yields
$$(y)yϵ(ϵ)=\mathrm{ln}𝔼[e^{yX}].$$
(15)
We conclude that $`(ϵ)`$ is given by the Legendre transform of the potential $`(y)`$,
$$\begin{array}{cc}& (ϵ)=ϵy(y),\hfill \\ & ϵ=_y(y),\hfill \end{array}$$
(16)
as prescribed by Cramér’s theorem.
### III.2 Poissonian random graphs
We consider now models defined on random graphs, first under the assumption that the internal structure of an instance is replica symmetric (RS). As a further simplification (to be relaxed later on, as for the RS hypothesis), we assume that the only source of quenched disorder lies in the graph structure, as it is the case for the vertex-cover and coloring problems. We consider here simultaneously the three ensembles of random graphs, $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, $`\overline{𝒢}_N^{(\gamma )}`$ and $`\widehat{𝒢}_N^{(\gamma )}`$ defined in II.1 and hereafter generically referred to as $`𝒢_N^{(\gamma )}`$. In the $`N\mathrm{}`$ limit, the degrees of graphs in $`𝒢_N^{(\gamma )}`$ have same limiting distribution $`\pi _\gamma `$, where $`\pi _\gamma `$ denotes the Poisson distribution with mean $`\gamma `$, $`\pi _\gamma (k)=\gamma ^ke^\gamma /k!`$.
Taking a graph uniformly at random in $`𝒢_N^{(\gamma )}`$ defines the measure over the quenched disorder with respect to which large deviations are evaluated. Following the basic principles of the cavity method, we consider the changes in the system when its size is increased from $`N`$ to $`N+1`$. The first idea would be to construct uniformly at random a graph $`G_{N+1}`$ in $`𝒢_{N+1}^{(\gamma )}`$ from a graph $`G_N`$ randomly chosen in $`𝒢_N^{(\gamma )}`$, by connecting a new node to $`k`$ nodes of $`G_N`$, with $`k`$ taken with the distribution $`\pi _\gamma `$. This construction is however too naïve, since, if the initial graph was in $`\overline{𝒢}_N^{(\gamma )}`$ for example, its extension is in $`\overline{𝒢}_{N+1}^{(\gamma ^{})}`$, with $`\gamma ^{}/2=(M+k)/(N+1)\gamma /2`$, where $`M=\gamma N/2`$ denotes the number of edges in $`G_N`$. However, for the three models, it appears that this construction yields a graph uniformly at random in $`𝒢_{N+1}^{(\gamma ^{})}`$, with $`\gamma ^{}=\gamma +\chi (\gamma ,k)/N`$, where we have obtained for $`\overline{𝒢}_N^{(\gamma )}`$ that $`\overline{\chi }(\gamma ,k)=2k\gamma `$. For $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, after addition of the new node the probability for an edge to be present is still $`\gamma /N`$ and not $`\gamma /(N+1)`$, so that it is described by $`\gamma ^{}`$ satisfying $`\gamma ^{}/(N+1)=\gamma /N`$, yielding $`\stackrel{~}{\chi }(\gamma ,k)=\gamma `$. Finally for $`\widehat{𝒢}_N^{(\gamma )}`$, after addition of the new site, the distribution of the degrees $`\pi _\gamma (K)`$ is modified to $`(1k/N)\pi _\gamma (K)+(k/N)\pi _\gamma (K1)`$ since $`k`$ of the nodes of $`G_N`$ receive an additional edge ; this leads to an effective distribution $`\pi _\gamma ^{}(K)`$ with $`\gamma ^{}=\gamma +k/N`$, so that $`\widehat{\chi }(k,\gamma )=k`$. To sum up, we obtained
$`\overline{\chi }(\gamma ,k)=`$ $`2k\gamma `$ $`\text{(Uniform model }\overline{𝒢}_N^{(\gamma )}),`$
$`\stackrel{~}{\chi }(\gamma ,k)=`$ $`\gamma `$ $`\text{(Binomial model }\stackrel{~}{𝒢}_N^{(\gamma )}),`$
$`\widehat{\chi }(\gamma ,k)=`$ $`k`$ $`\text{(Poissonian model }\widehat{𝒢}_N^{(\gamma )}).`$
The fact that in all cases $`\chi (\gamma ,k)=\gamma `$, with the average $``$ taken with respect to $`\pi _\gamma `$, reflects the equivalence of the typical properties between the three models.
When a node is added, the ground-state energy is shifted by an amount $`\mathrm{\Delta }E`$. Conditioned to the fact that the new node is connected to $`k`$ other nodes, this shift has distribution $`P_n^{(k)}(\mathrm{\Delta }E)`$ from graph to graph (and from node to node on a given graph). Given $`P_n^{(k)}(\mathrm{\Delta }E)`$, the argument followed in Eq. (13) can essentially be repeated,
$$\begin{array}{cc}& e^{(N+1)(E/(N+1),\gamma )}e^{N(ϵ)}e^{(ϵ,\gamma )yϵ}=e^{N(ϵ)}e^{(y,\gamma )}\hfill \\ & =\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{N[(E\mathrm{\Delta }E)/N,\gamma \chi (\gamma ,k)/N]}e^{N(ϵ)}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}e^{z\chi (\gamma ,k)},\hfill \end{array}$$
(17)
with the notations $`ϵE/N`$, $`y_ϵ(ϵ,\gamma )`$, $`z_\gamma (ϵ,\gamma )`$ and $`(y,\gamma )=yϵ(ϵ)`$. Eq. (17) gives the Legendre transform of the rate function, $`(y,\gamma )`$, as a function of $`y`$ and $`z`$. To determine $`z`$, we need consider the energy shift $`\mathrm{\Delta }E`$ due to a link addition, having distribution $`P_{\mathrm{}}(\mathrm{\Delta }E)`$. More precisely, the average value of the energy shift when $`\gamma \gamma +1/N`$ at fixed number of nodes $`N`$ is required, which is obtained by adding $`k`$ new edges with an appropriate distribution $`\sigma (k)`$. For $`\overline{𝒢}_N^{(\gamma )}`$, adding a single edge results in $`\gamma ^{}/2=(M+1)/N=\gamma /2+1/N`$ so we take formally $`\sigma =\delta _{1/2}`$, where $`\delta _\theta (k)=\delta _{k,\theta }`$ denotes the Kronecker symbol (this non-integer prescription could be avoided as in Mézard and Parisi (2003) by adding two nodes at once instead of one). For $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, we take $`\sigma =\pi _{1/2}`$ because it corresponds to the distribution of the number of added edges when each of the $`N^2/2`$ edges has a probability $`1/N^2`$ to be present in the $`\stackrel{~}{𝒢}_N^{(\gamma +1/N)}`$ graph, but absent in the $`\stackrel{~}{𝒢}_N^{(\gamma )}`$ one. Finally for $`\widehat{𝒢}_N^{(\gamma )}`$, the addition of one edge leads to $`\gamma ^{}=\gamma +2/N`$ so that formally $`\sigma =\delta _{1/2}`$ as in the uniform model. A $`1/N`$ expansion of $`\mathrm{exp}[N(ϵ,\gamma +1/N)]`$ then yields $`z=_\gamma (ϵ,\gamma )`$ as
$$e^z=\underset{k=0}{\overset{\mathrm{}}{}}\sigma (k)\left(𝑑\mathrm{\Delta }EP_{\mathrm{}}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}\right)^k,$$
(18)
with as derived just above,
$`\overline{\sigma }(k)=`$ $`\delta _{1/2}(k)`$ $`\text{(Uniform model }\overline{𝒢}_N^{(\gamma )}),`$
$`\stackrel{~}{\sigma }(k)=`$ $`\pi _{1/2}(k)`$ $`\text{(Binomial model }\stackrel{~}{𝒢}_N^{(\gamma )}),`$
$`\widehat{\sigma }(k)=`$ $`\delta _{1/2}(k)`$ $`\text{(Poissonian model }\widehat{𝒢}_N^{(\gamma )}).`$
As in the typical cavity method Mézard and Parisi (2001, 2003), the distributions $`P_n^{(k)}(\mathrm{\Delta }E)`$ and $`P_{\mathrm{}}(\mathrm{\Delta }E)`$ can be calculated by means of cavity fields. The fundamental assumption made at this stage is that the nodes to which a new node is added are independent in the absence of the added node. Under this assumption, the problem on a random graph is reduced to a problem on a tree with self-consistent boundary conditions (the so-called Bethe lattice). While the same procedure applies to other optimization problems, we restrict here for simplicity to the vertex-cover problem, for which we take the ground-state energy as the minimum of non-covered nodes under the constraint that neighboring nodes cannot be both covered. A recursion is written for rooted-trees with same degree distribution $`\pi _\gamma (k)`$ as the graphs (see Fig. 1). In general if the degree distribution is $`p(k)`$, the root must be assigned the distribution $`q(k)=(k+1)p(k+1)/k`$, which corresponds to the probability, when the edge $`(i0)`$ is chosen at random, that $`i`$ has $`k`$ neighbors in addition to $`0`$ ; the Poisson distribution has however the specificity that $`q(k)=p(k)=\pi _\gamma (k)`$. For the rooted tree with root-node $`i`$, let $`E_0^{(i0)}`$ be the minimal energy with the root constrained to be non-covered (white), and $`E_1^{(i0)}`$ the minimal energy with the root constrained to be covered (black). These quantities are related to those associated with the $`k`$ rooted trees generated by deletion of the edges originating from $`i`$ (see Fig. 1) by
$$\begin{array}{cc}& E_0^{(i0)}=1+\underset{j=1}{\overset{k}{}}\mathrm{min}(E_0^{(ji)},E_1^{(ji)}),\hfill \\ & E_1^{(i0)}=\underset{j=1}{\overset{k}{}}E_0^{(ji)}.\hfill \end{array}$$
(19)
A cavity field is defined for each oriented edge as $`h^{(i0)}E_0^{(i0)}\mathrm{min}(E_0^{(i0)},E_1^{(i0)})`$ ; it satisfies the relation
$$h^{(i0)}=\widehat{h}^{(k)}(\{h^{(ji)}\})=\mathrm{max}(0,1\underset{j=1}{\overset{k}{}}h^{(ji)}).$$
(20)
The addition of the new node $`i`$ is associated with an energy shift given by
$$\begin{array}{cc}\hfill \mathrm{\Delta }E_{\mathrm{node}}& =\mathrm{min}(E_0^{(i0)},E_1^{(i0)})\underset{j=1}{\overset{k}{}}\mathrm{min}(E_0^{(ji)},E_1^{(ji)})\hfill \\ & =\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h^{(ji)}\})=\mathrm{min}(1,\underset{j=1}{\overset{k}{}}h^{(ji)}).\hfill \end{array}$$
(21)
We will also need the energy shift corresponding to an edge addition, which reads
$$\begin{array}{cc}\hfill \mathrm{\Delta }E_{\mathrm{edge}}& =\mathrm{min}(E_0^{(12)}+E_0^{(21)},E_0^{(12)}+E_1^{(21)},E_1^{(12)}+E_0^{(21)})\mathrm{min}(E_0^{(12)},E_1^{(12)})\mathrm{min}(E_0^{(21)},E_1^{(21)})\hfill \\ & =\mathrm{\Delta }\widehat{E}_{\mathrm{}}(h^{(12)},h^{(21)})=\mathrm{min}(h^{(12)},h^{(21)}).\hfill \end{array}$$
(22)
With the help of these equations, the distributions $`P_n^{(k)}(\mathrm{\Delta }E)`$ and $`P_{\mathrm{}}(\mathrm{\Delta }E)`$ are easily obtained once the distribution for the fields $`P(h)`$ is known.
Again similarly to the typical cavity method Mézard and Parisi (2001, 2003), to derive the equation satisfied by $`P(h)`$, we introduce $`R_{N+1}^{(\gamma )}(h,E)`$, the probability to get an energy $`E`$ and cavity field $`h`$ when taking at random a graph in $`𝒢_{N+1}^{(\gamma )}`$ and choosing one of its node as a root. By definition, the marginalization over $`h`$ gives $`_{N+1}^{(\gamma )}(E/(N+1))`$, the probability to get a graph in $`𝒢_{N+1}^{(\gamma )}`$ with energy $`E`$,
$$𝑑hR_{N+1}^{(\gamma )}(h,E)e^{(N+1)(E/(N+1),\gamma )}.$$
(23)
Generalizing slightly Eq. (17), we write
$$\begin{array}{cc}\hfill R_{N+1}^{(\gamma )}(h,E)& =\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)𝑑\mathrm{\Delta }EQ^{(k)}(h,\mathrm{\Delta }E)e^{N[(E\mathrm{\Delta }E)/N,\gamma \chi (\gamma ,k)/N]},\hfill \\ & e^{N(ϵ)}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)𝑑\mathrm{\Delta }EQ^{(k)}(h,\mathrm{\Delta }E)e^{y\mathrm{\Delta }Ez\chi (\gamma ,k)},\hfill \end{array}$$
(24)
where $`Q^{(k)}(h,\mathrm{\Delta }E)`$ denotes the joint distribution of the cavity field $`h`$ and the energy shift $`\mathrm{\Delta }E`$. As in the typical cavity method, we verify that $`h`$ is independent of $`E`$, more precisely,
$$\begin{array}{cc}& R_{N+1}^{(\gamma )}(h,E)=e^{N(ϵ)}e^{(y,\gamma )}P(h),\hfill \\ & P(h_0)=\frac{1}{Z}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)\underset{i=1}{\overset{k}{}}dh_iP(h_i)\delta (h_0\widehat{h}^{(k)}(\{h_i\}))e^{y\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_i\})z\chi (\gamma ,k)},\hfill \\ & Z=\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}e^{z\chi (\gamma ,k)},\hfill \\ & (y,\gamma )=ye(ϵ,\gamma )=\mathrm{ln}Z,\hfill \end{array}$$
(25)
where $`P(h)`$ also depends on $`\gamma `$ and $`y`$. In the particular case where $`\chi (\gamma ,k)`$ does not depend on $`k`$, the relation for $`P(h)`$ formally corresponds to what is known in the literature as a 1RSB factorized equation with 1RSB parameter $`y`$ Wong and Sherrington (1988); Mézard and Parisi (2003); this is the case with $`\stackrel{~}{𝒢}_N^{(\gamma )}`$ but not with $`\overline{𝒢}_N^{(\gamma )}`$ and $`\widehat{𝒢}_N^{(\gamma )}`$.
Specializing now to the vertex-cover problem, the equations are simplified by the fact that $`h\{0,1\}`$, so that the distribution $`P(h)`$ can be parameterized by a single real $`\eta `$, with $`P(h)=\eta \delta (h1)+(1\eta )\delta (h)`$. As seen from Eqs. (17) and (18), the distributions $`P_n^{(k)}(\mathrm{\Delta }E)`$ and $`P_{\mathrm{}}(\mathrm{\Delta }E)`$ are needed only through their Laplace transforms, which are given by
$$\begin{array}{cc}& 𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}=\underset{i=1}{\overset{k}{}}dh_iP(h_i)e^{y\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_i\})}=e^y+(1e^y)(1\eta )^k,\hfill \\ & 𝑑\mathrm{\Delta }EP_{\mathrm{}}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}=\underset{i=1}{\overset{2}{}}dh_iP(h_i)e^{y\mathrm{\Delta }\widehat{E}_{\mathrm{}}(h_1,h_2)}=1+(e^y1)\eta ^2.\hfill \end{array}$$
(26)
As a first check, one verifies that for $`y=0`$, the typical RS ground-state energy density $`\overline{ϵ}^{(\gamma )}`$ is reobtained Weigt and Hartmann (2000), with same value for the three ensembles $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, $`\overline{𝒢}_N^{(\gamma )}`$ and $`\widehat{𝒢}_N^{(\gamma )}`$,
$$\begin{array}{cc}& \overline{ϵ}^{(\gamma )}=1\eta \gamma \eta ^2/2,\hfill \\ & \eta =e^{\gamma \eta }.\hfill \end{array}$$
(27)
The equations for $`y0`$ can also be written explicitly. For the ensemble $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, they read
$$\begin{array}{cc}& \eta =\frac{1}{1+e^y(e^{\gamma \eta }1)},z=\frac{1}{2}(e^y1)\eta ^2,\hfill \\ & (y,\gamma )=\mathrm{ln}[e^y+(1e^y)e^{\gamma \eta }]+\frac{\gamma }{2}(e^y1)\eta ^2.\hfill \end{array}$$
(28)
For the ensemble $`\overline{𝒢}_N^{(\gamma )}`$,
$$\begin{array}{cc}& \eta =\frac{1}{1+e^y(e^{\gamma \eta e^z}1)},z=\mathrm{ln}[1+(e^y1)\eta ^2],\hfill \\ & (y,\gamma )=\mathrm{ln}[e^y+(1e^y)e^{\gamma \eta e^z}]+\gamma (1e^z)\gamma z/2.\hfill \end{array}$$
(29)
And for the ensemble $`\widehat{𝒢}_N^{(\gamma )}`$,
$$\begin{array}{cc}& \eta =\frac{1}{1+e^y(e^{\gamma \eta e^z}1)},z=\frac{1}{2}\mathrm{ln}[1+(e^y1)\eta ^2],\hfill \\ & (y,\gamma )=\mathrm{ln}[e^y+(1e^y)e^{\gamma \eta e^z}]+\gamma (1e^z).\hfill \end{array}$$
(30)
The formulae (29) coincide with the result of the replica computation presented in Montanari and Zecchina (2002).
The three corresponding rate functions are plotted in Fig. 2 for $`\gamma =2`$.
A remarkable aspect of the vertex-cover problem is the presence, in the typical phase diagram, of a continuous phase transition at $`\gamma _c=e2.71`$, from an RS phase at $`\gamma <\gamma _c`$ to a presumably full-RSB phase at $`\gamma >\gamma _c`$ Weigt and Hartmann (2003). Due to its continuous character, the phase transition can be located by analyzing the stability analysis of the RS Ansatz. Extending the stability analysis from typical to atypical graphs thus provides, in the $`(\gamma ,y)`$ plane, a phase diagram with an ”AT line” Mézard et al. (1987) separating a RS phase from a full-RSB one. The three ensembles are not equivalent with respect to properties associated with atypical graphs, and we concentrate here on the binomial ensemble $`\stackrel{~}{𝒢}_N^{(\gamma )}`$.
RSB effects are much likely to appear first for negative values of $`y`$, corresponding to the most frustrated graphs. The failure of the RS approach can in fact be inferred from an asymptotic analysis of the $`y\mathrm{}`$ limit : it yields $`\eta e^{y/2}/\sqrt{\gamma }`$, $`ϵ(y=\mathrm{})=1/2`$ and $`L(y=\mathrm{})=(1\mathrm{ln}\gamma )/2`$. Clearly, this is inconsistent as soon as $`\gamma >e`$ since then it predicts then $`L(ϵ=1/2)<0`$. The value thus obtained coincides with the value of $`\gamma _c`$ for the failure of the RS approach to typical graphs Weigt and Hartmann (2000); Bauer and Golinelli (2001) (the reason for this correspondence is elucidated below). The negativeness of the rate function is however a sufficient but not necessary sign of RSB. A more refined way to detect it consists in studying the stability of the RS large deviations Ansatz. For the binomial model, it happens to be strictly equivalent to the stability analysis of a factorized 1RSB Ansatz Montanari and Ricci-Tersenghi (2003), and reads
$$(\gamma \eta )^2e^y<1.$$
(31)
It starts to be violated at $`y=\mathrm{}`$ for $`\gamma >1`$, while the typical graphs described with $`y=0`$ are not concerned before $`\gamma =e`$. Indeed for $`\gamma <1`$, the RS Ansatz is stable for all $`y`$: $`(\gamma \eta )^2e^y`$ is a decreasing function of $`y`$ and for $`y\mathrm{}`$ the asymptotic analysis yields $`(\gamma \eta )^2e^y\gamma `$. At $`\gamma =1`$, only the $`y=\mathrm{}`$ point, corresponding to the maximum achievable energy $`ϵ=1/2`$, is marginally unstable. Finally, for $`\gamma >1`$, there is a critical $`y_c`$ such that RS is stable for $`y>y_c`$ but unstable for $`y<y_c`$ ; $`y_c`$ increases when $`\gamma `$ increases and reaches $`y_c=0`$ for $`\gamma _c=e`$, the point where the typical problem undergoes the RSB transition. For $`\gamma >\gamma _c`$ while typical graphs are FRSB, some less frustrated graphs are still RS. The resulting phase diagram is shown in Fig. 4 in the plane $`(\gamma ,y)`$ and in Fig. 4 in the plane $`(\gamma ,ϵ)`$. The occurrence of RSB at $`\gamma =1`$ is particularly interesting because this point corresponds to the percolation threshold of a giant connected component Bollobás (2001), which appeared totally irrelevant when restricting to typical graphs Weigt and Hartmann (2000). In contrast, when atypical graphs are included into the picture, the emergence of a giant component seems to be responsible for the onset of RSB, as shown in Fig. 4 (a similar analysis of the uniform model however reveals that in this case RSB appears only above an average connectivity $`\gamma =2`$).
The opposite $`y+\mathrm{}`$ limit is also interesting since it is always correctly described by the RS Ansatz, with $`ϵ(y=\mathrm{},\gamma )=0`$ and $`(y=\mathrm{},\gamma )=\gamma /2`$. It can be interpreted as corresponding to graphs with no edge at all, which occurs with probability
$$_N^{(\gamma )}(\mathrm{non}\mathrm{frustrated})_N^{(\gamma )}(\mathrm{no}\mathrm{edge})\left(1\frac{\gamma }{N}\right)^{N^2/2}e^{N\gamma /2}.$$
(32)
Similar relations between the $`y=\mathrm{}`$ limit and the probability for non-frustrated samples have been reported in a variety of other models Rivoire (2004), providing consistent checks of the method.
### III.3 Random graphs with given degree distributions
The LDCM applies as well to graph ensembles with non-poissonian degrees, and an example is provided here with graph ensembles specified by their degree distribution $`p(k)`$, that is with each node having an independent probability $`p(k)`$ of being of degree $`k`$. The reasoning for arbitrary $`p(k)`$ basically follows the procedure used for the Poissonian model $`\widehat{𝒢}_N^{(\gamma )}`$ which was a particularly case where $`p(k)=\gamma ^ke^\gamma /k!`$.
We thus first consider how the ensemble is modified when a new node is connected to $`k`$ nodes of a graph made of $`N`$ nodes having degree distribution $`p(K)`$. The degree distribution becomes $`p^{}(K)`$, with
$$p^{}(K)=\left(1\frac{k}{N}\right)p(K)+\frac{k}{N}p(K1),$$
(33)
since a given node has a probability $`k/N`$ to get its degree increased by one unit (the probability that it is increased by more than one unit is $`O(1/N^2)`$ and is therefore neglected). Writing $`p(K)=_rp_r\delta _{K,r}`$, we explicitly have $`p_r^{}=p_r+(k/N)\delta p_r`$ with $`\delta p_r=p_{r+1}p_r`$. The set $`\{p_r\}`$ can serve as a characterization of the graph ensemble, and following the same scheme as for Poissonian graphs, we obtain
$$e^{(y,\{p_r\})}=\underset{k=0}{\overset{\mathrm{}}{}}p_k𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}e^{kz},$$
(34)
with as before $`y=_ϵ(ϵ,\{p_r\})`$ and now
$$z\underset{r}{}\delta p_r\frac{(ϵ,\{p_r\})}{p_r}.$$
(35)
To get $`z`$, we notice that Eq. (33) with $`k=2`$ describes the effect of an edge addition so that
$$e^{2z}=𝑑\mathrm{\Delta }EP_{\mathrm{}}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}.$$
(36)
The potential whose Legendre transform yields the rate function can therefore be written
$$(y,\{p_r\})=\mathrm{ln}\left[\underset{k=0}{\overset{\mathrm{}}{}}p_k𝑑\mathrm{\Delta }EP_n^{(k)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}\left(𝑑\mathrm{\Delta }E^{}P_{\mathrm{}}(\mathrm{\Delta }E^{})e^{y\mathrm{\Delta }E^{}}\right)^{k/2}\right].$$
(37)
Similarly, the cavity equation for the fields reads
$$P(h_0)=\frac{1}{Z}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(k+1)p_{k+1}}{k}\underset{j=1}{\overset{k}{}}dh_jP(h_j)\delta (h_0\widehat{h}^{(k)}(\{h_j\}))e^{y\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_j\})}e^{kz}$$
(38)
where $`Z`$ is the appropriate normalization and the presence of $`(k+1)p_{k+1}/k`$ instead of $`p_k`$ reflects the fact that an oriented edge is chosen at random, and not a node as in Eq. (37) (only for Poissonian graphs do these two probabilities happen to be the same).
A case of particular interest is when $`p(k)=\delta _{k,r}`$, corresponding to random $`r`$-regular graphs. In such a case, the factor $`e^{kz}`$ in Eq. (38) is a constant which can be absorbed into the normalization $`Z^{}=Ze^{(r1)z}`$,
$$P(h_0)=\frac{1}{Z^{}}\underset{j=1}{\overset{r1}{}}dh_jP(h_j)\delta (h_0\widehat{h}^{(r1)}(\{h_j\}))e^{y\mathrm{\Delta }\widehat{E}_n^{(r1)}(\{h_j\})}.$$
(39)
This is formally identical to what is known as a 1RSB factorized cavity equation Wong and Sherrington (1988); Mézard and Parisi (2003). The correspondence extends to the formula for the potential, Eq. (37), which becomes for regular graphs,
$$(y)=\mathrm{ln}\left[𝑑\mathrm{\Delta }EP_n^{(r)}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}\right]+\frac{r}{2}\mathrm{ln}\left[𝑑\mathrm{\Delta }EP_{\mathrm{}}(\mathrm{\Delta }E)e^{y\mathrm{\Delta }E}\right].$$
(40)
As a consequence, the 1RSB complexity of models defined on random regular graphs coincides with minus a rate function, as already noticed in Rivoire (2004). Obviously, the correspondence holds only within the RS approximation that has been assumed throughout since by nature of a 1RSB glassy phase, the complexity is positive while a rate function is necessarily non-negative ; it will be shown below how the formalism needs to be extended to include the possibility of RSB. When the factorization does not hold, the correspondence between rate functions over the external disorder and negative complexities is only approximate ; we have seen however that for vertex-cover on $`\stackrel{~}{𝒢}_N^{(\gamma )}`$ the rate function starts getting negative values precisely at the point $`\gamma =e`$ where typical graphs undergo a RSB transition, in agreement with the observation that the two quantities are approximatively related.
## IV Multi-step large deviations
The LDCM can naturally be extended beyond the simple case of zero-temperature systems in an RS phase with disorder only specified by a random graph ensemble. We consider here successively finite-temperature systems, models with RSB phases, and external disorders including random couplings. In the three cases, a second temperature is needed to describe the large deviations with respect to the additional source of randomness. In each case also, the equations display a common 2RSB-like structure Mézard et al. (1987), which would be promoted to the $`n`$RSB type with $`n>2`$ if $`n`$ different sources of disorder were present.
### IV.1 Finite temperature
The simplest extension requiring multi-step large deviations consists in generalizing the description of a model on a given graph from zero to finite temperature. Two inverse temperatures are now required : $`\beta `$, for the thermodynamics on a given graph, and $`y`$ for the large deviations in the graph ensemble. More precisely, large deviations now concern the density of free energy $`f(\beta )`$, with the limit $`\beta \mathrm{}`$ giving back to the large deviations for the ground-state density energy $`ϵ=lim_\beta \mathrm{}f(\beta )`$, as discussed so far. For any fixed value of $`\beta `$, the rate function $`L(f,\beta )(f,\beta )`$ is calculated as before through the Legendre transform of a potential $`(y,\beta )`$ satisfying
$$\begin{array}{cc}& e^{N(y,\beta )}=𝑑fe^{N[(f,\beta )yf]}=\frac{1}{\mathrm{\#}𝒢}\underset{G𝒢}{}Z_G(\beta )^{y/\beta },\hfill \\ & Z_G(\beta )e^{\beta Nf_G(\beta )}\underset{C𝒞_G}{}e^{\beta E(C)},\hfill \end{array}$$
(41)
where we introduced $`Z_G(\beta )`$ the partition function on the graph $`G`$ at temperature $`\beta `$, $`𝒞_G`$ the set of configurations on the graph $`G`$ and $`\mathrm{\#}𝒢`$ the cardinality of the ensemble of graphs $`𝒢`$. Note that the particular choice $`y=\beta `$ corresponds to the uniform measure over all configurations $`\{C𝒞_G\}_{G𝒢}`$ :
$$\underset{G𝒢}{}\underset{C𝒞_G}{}e^{\beta E(C)}=\underset{G𝒢}{}e^{\beta Nf_G(\beta )}=(\mathrm{\#}𝒢)e^{\beta N(\beta ,\beta )}.$$
(42)
From the technical point of view, we just have to replace in all formulae the functions $`\widehat{h}^{(k)}(\{h_i\})`$, $`\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_i\})`$ and $`\mathrm{\Delta }\widehat{E}_{\mathrm{}}(h_1,h_2)`$ by their finite-temperature extensions $`\widehat{h}^{(k)}(\{h_i\};\beta )`$, $`\mathrm{\Delta }\widehat{F}_n^{(k)}(\{h_i\};\beta )`$ and $`\mathrm{\Delta }\widehat{F}_{\mathrm{}}(h_1,h_2;\beta )`$. Taking the vertex-cover problem as an example, these quantities are derived by writing recursive equations for the conditional partition functions $`Z_0^{(i0)}(\beta )`$ and $`Z_1^{(i0)}(\beta )`$ instead of the conditional ground-state energies $`E_0^{(i0)}`$ and $`E_1^{(i0)}`$. More precisely, Eqs. (19) are replaced with
$$\begin{array}{cc}\hfill Z_0^{(i0)}(\beta )=& e^\beta \underset{j=1}{\overset{k}{}}\left(Z_0^{(ji)}(\beta )+Z_1^{(ji)}(\beta )\right),\hfill \\ \hfill Z_1^{(i0)}(\beta )=& \underset{j=1}{\overset{k}{}}Z_0^{(ji)}(\beta ).\hfill \end{array}$$
(43)
To get $`lim_\beta \mathrm{}h^{(i0)}(\beta )=h^{(i0)}`$ with $`h^{(i0)}`$ defined in Eq. (20), the cavity fields at finite temperature $`h^{(i0)}(\beta )`$ are defined as
$$h^{(i0)}(\beta )\frac{1}{\beta }\mathrm{ln}\left(\frac{Z_0^{(i0)}(\beta )}{Z_0^{(i0)}(\beta )+Z_1^{(i0)}(\beta )}\right).$$
(44)
With these definitions, the different functions required to compute the rate function $`L(f,\beta )`$ are
$$\begin{array}{cc}& \widehat{h}^{(k)}(\{h_j\},\beta )=\frac{1}{\beta }\mathrm{ln}\left(1+e^{\beta (1_{j=1}^kh_j)}\right),\hfill \\ & \mathrm{\Delta }\widehat{F}_n^{(k)}(\{h_j\};\beta )=\frac{1}{\beta }\mathrm{ln}\left(e^\beta +e^{\beta _{j=1}^kh_j}\right),\hfill \\ & \mathrm{\Delta }\widehat{F}_{\mathrm{}}(h_1,h_2;\beta )=\frac{1}{\beta }\mathrm{ln}\left(e^{\beta h_1}+e^{\beta h_2}e^{\beta (h_1+h_2)}\right),\hfill \end{array}$$
(45)
which all reduce as it should to $`\widehat{h}^{(k)}(\{h_j\})`$, $`\mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_j\})`$ and $`\mathrm{\Delta }\widehat{E}_{\mathrm{}}(h_1,h_2)`$ given in Eq. (20), (21) and (22) when $`\beta \mathrm{}`$. The only practical difference with the $`\beta =\mathrm{}`$ case is that the distribution $`P(h)`$ has no more reason to be peaked on integers and therefore cannot be parameterized by a single real number.
### IV.2 Replica symmetry breaking
The phase space of a glassy 1RSB instance is structured into clusters whose energy is controlled by a parameter $`\mu `$ in exactly the same way the finite inverse temperature $`\beta `$ controls the equilibrium configurations according to their energy. Therefore, the extension of the LDCM from paramagnetic (RS) systems to glassy (1RSB) systems, is formally similar to the extension from zero temperature to finite temperature. The counterpart of Eq. (41) reads
$$\begin{array}{cc}& e^{N(y,\mu )}=\frac{1}{\mathrm{\#}𝒢}\underset{G𝒢}{}e^{yN\varphi _G(\mu )}=𝑑\varphi e^{N[(\varphi ,\mu )y\varphi ]},\hfill \\ & e^{N\mu \varphi _G(\mu )}=\underset{\alpha G}{}e^{\mu Nϵ_\alpha }=𝑑ϵe^{N[\mathrm{\Sigma }_G(ϵ)\mu ϵ]}\hfill \end{array}$$
(46)
where $`\varphi _G(\mu )`$ is the 1RSB potential on graph $`G`$ and $`L(\varphi ,\mu )(\varphi ,\mu )`$ is the rate function describing the large deviations of $`\varphi (\mu )`$ over the ensemble of random graphs ; in these formulae, we reserve greek letters for quantities related to the internal structure and use $`\alpha `$ to index the clusters. The saddle points in Eq. (46) lead to the following Legendre transform relations :
$`(y,\mu )`$ $`=y\varphi (\varphi ,\mu ),`$ $`y=_\varphi (\varphi ,\mu ),`$ (47)
$`\mu \varphi (\mu )`$ $`=\mu ϵ\mathrm{\Sigma }(ϵ),`$ $`\mu =_ϵ\mathrm{\Sigma }(ϵ).`$ (48)
These quantities are computed by applying the standard 1RSB cavity method Mézard and Parisi (2001) to a given set of atypical graphs characterized by their ground-state energy density $`ϵ_0`$. If $`\rho _N(ϵ|ϵ_0)`$ denotes the distribution of their clusters, the corresponding complexity is defined as $`\rho _N(ϵ|ϵ_0)\mathrm{exp}[N\mathrm{\Sigma }(ϵ|ϵ_0)]`$ ; for an energy $`ϵ`$ close to the ground-state reference energy $`ϵ_0`$, it becomes
$$\rho _N(ϵ|ϵ_0)e^{\mu N(ϵϵ_0)}$$
(50)
where $`\mu _ϵ\mathrm{\Sigma }(ϵ=ϵ_0|ϵ_0)`$ defines the 1RSB internal inverse temperature. The shift in energy $`\mathrm{\Delta }E`$ induced by a node addition, which is needed in the recursion at the level of the graph average, is given by the shift in the reference energy, that is,
$$\rho _{N+1}(ϵ|ϵ_0)=\rho _N(ϵ|ϵ_0)e^{\mu \mathrm{\Delta }E}.$$
(51)
The expression for the reweighting term $`\mathrm{\Xi }e^{\mu \mathrm{\Delta }E}`$ is read from the 1RSB cavity recursion which involves $`\mathrm{\Pi }(h)`$, the distribution of cavity fields over the clusters, and $`P[\mathrm{\Pi }]`$, the distribution of the $`\mathrm{\Pi }`$’s over the graphs ; for a given (class of) graph, the connection of a new node to $`k`$ other ones indeed yields
$$\begin{array}{cc}\hfill \mathrm{\Pi }_0=\widehat{\mathrm{\Pi }}^{(k)}[\{\mathrm{\Pi }_i\}],\mathrm{with}& \widehat{\mathrm{\Pi }}^{(k)}[\{\mathrm{\Pi }_i\}](h_0)=\frac{1}{\mathrm{\Xi }}\underset{i=1}{\overset{k}{}}\mathrm{\Pi }_i(h_i)\delta (h_0\widehat{h}^{(k)}(\{h_i\}))e^{\mu \mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_i\})},\hfill \\ & \mathrm{\Xi }=e^{\mu \mathrm{\Delta }E}=\widehat{\mathrm{\Xi }}^{(k)}[\{\mathrm{\Pi }_i\}]=\underset{i=1}{\overset{k}{}}\mathrm{\Pi }_i(h_i)e^{\mu \mathrm{\Delta }\widehat{E}_n^{(k)}(\{h_i\})}.\hfill \end{array}$$
(52)
Therefore, at the level of the graph average, we have for Poissonian graphs
$$P[\mathrm{\Pi }_0]=\frac{1}{Z}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)\underset{i=1}{\overset{k}{}}𝒟\mathrm{\Pi }_iP[\mathrm{\Pi }_i]\delta [\mathrm{\Pi }_0\widehat{\mathrm{\Pi }}^{(k)}[\{\mathrm{\Pi }_i\}]\widehat{\mathrm{\Xi }}^{(k)}(\{\mathrm{\Pi }_i\})]^{y/\mu }e^{z\chi (k,\gamma )},$$
(53)
and for graphs with fixed degree distribution
$$P[\mathrm{\Pi }_0]=\frac{1}{Z^{}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(k+1)p_{k+1}}{k}\underset{i=1}{\overset{k}{}}𝒟\mathrm{\Pi }_iP[\mathrm{\Pi }_i]\delta [\mathrm{\Pi }_0\widehat{\mathrm{\Pi }}^{(k)}(\{\mathrm{\Pi }_i\})]\widehat{\mathrm{\Xi }}^{(k)}[\{\mathrm{\Pi }_i\}]^{y/\mu }e^{kz}.$$
(54)
The 1RSB large deviations equations have thus the structure of a typical factorized 2RSB theory, as the RS large deviations equations resembled a typical factorized 1RSB theory. In particular, for $`y=0`$, the non-factorized 1RSB formalism is exactly recovered.
Replica symmetry breaking (RSB) is relevant to many optimization problems, and the vertex-cover problem already provided us with such an example. For this problem, studying the local stability of the RS Ansatz was enough to locate the continuous transition to a RSB phase. However, other problems may display a different kind of glass transition, known as a discontinuous 1RSB transition, which, due to its discontinuous character, can only be correctly described by implementing a 1RSB formalism. Such a transition is found for instance in the coloring problem Mulet et al. (2002), which we take here as an illustrative example of the broader class of constraint satisfaction problems.
In the statistical physics point of view, a problem is satisfiable (SAT) if it has a zero ground-state energy density, $`ϵ=0`$ foo (b). In presence of a clustered glassy phase however, an alternative characterization is provided by the sign of the complexity $`\mathrm{\Sigma }_0`$, giving, when positive, the number $`\mathrm{exp}[N\mathrm{\Sigma }_0]`$ of clusters with $`ϵ=0`$. This complexity $`\mathrm{\Sigma }_0`$ is obtained in the 1RSB formalism by taking the limit $`\mu \mathrm{}`$ ; for the 3-coloring problem on Erdős-Rényi graphs $`\stackrel{~}{𝒢}_N^{(\gamma )}`$, $`\mathrm{\Sigma }_0`$ is found to be positive only in a restricted interval $`[\gamma _d,\gamma _c]`$, with $`\gamma _d=4.42`$ and $`\gamma _c=4.69`$ Mulet et al. (2002), as schematically represented in Fig. 7. The threshold values $`\gamma _d`$ and $`\gamma _c`$, which also appear in other constraint satisfaction problems such as $`K`$-SAT Mézard and Zecchina (2002), locate two phase transitions, called respectively the clustering and SAT-UNSAT transitions : with probability one when $`N\mathrm{}`$, a graph with $`\gamma <\gamma _d`$ is colorable and RS, a graph with $`\gamma _d<\gamma <\gamma _c`$ is again colorable but RSB, and a graph with $`\gamma >\gamma _c`$ is uncolorable.
As an illustration, we consider here the 3-coloring problem on the Erdős-Rényi ensemble $`\stackrel{~}{𝒢}_N^{(\gamma )}`$. Following Braunstein et al. (2003), Eq. (53), the shift $`\mathrm{\Delta }\varphi _{\mathrm{}}`$ in the 1RSB potential due to a link addition is given by
$$e^{\mu \mathrm{\Delta }\varphi _{\mathrm{}}}=1+(e^\mu 1)q\eta _1\eta _2,$$
(55)
where the $`\eta _j`$’s, with distribution $`\rho (\eta )`$, represent the 1RSB cavity fields for this problem Braunstein et al. (2003). In the LDCM, we need the Laplace transform of the distribution $`P_{\mathrm{}}(\mathrm{\Delta }\varphi )`$, which thus reads
$$𝑑\mathrm{\Delta }\varphi P_{\mathrm{}}(\mathrm{\Delta }\varphi )e^{y\mathrm{\Delta }\varphi }=\underset{i=1,2}{}d\eta _i\rho (\eta _i)\left(1+(e^\mu 1)q\eta _1\eta _2\right)^{y/\mu }.$$
(56)
If only SAT configurations are to be addressed, the general 1RSB-LDCM equations can be simplified by taking the limit $`\mu \mathrm{}`$. This limit enforces $`ϵ0`$ and $`\mu \varphi (\mu )\mathrm{\Sigma }_0`$ and requires to rescale $`y`$ by taking $`y\mathrm{}`$ with $`x=y/\mu `$ fixed, such that $`(y,\mu )(x)`$ with
$$(x)=x\mathrm{\Sigma }_0(\mathrm{\Sigma }_0),x=_{\mathrm{\Sigma }_0}(\mathrm{\Sigma }_0),$$
(57)
where $`(\mathrm{\Sigma }_0)=lim_\mu \mathrm{}(\varphi =\mathrm{\Sigma }_0/\mu ,\mu )`$. In this limit,
$$𝑑\mathrm{\Delta }\varphi P_{\mathrm{}}(\mathrm{\Delta }\varphi )e^{y\mathrm{\Delta }\varphi }\underset{i=1,2}{}d\eta _i\rho (\eta _i)\left(1q\eta _1\eta _2\right)^x.$$
(58)
Similarly for site addition, we have, again in the limit $`\mu \mathrm{}`$,
$$𝑑\mathrm{\Delta }\varphi P_n^{(k)}(\mathrm{\Delta }\varphi )e^{y\mathrm{\Delta }\varphi }\underset{i=1}{\overset{k}{}}d\eta _i\rho (\eta _i)\widehat{\mathrm{\Xi }}^{(k)}(\eta _1,\mathrm{},\eta _k)^x,$$
(59)
with
$$\mathrm{\Xi }^{(k)}\underset{\mu \mathrm{}}{lim}e^{\mu \mathrm{\Delta }\varphi _n^{(k)}}=\underset{\mathrm{}=0}{\overset{q1}{}}(1)^{\mathrm{}}\left(\genfrac{}{}{0pt}{}{q}{\mathrm{}+1}\right)\underset{i=1}{\overset{k}{}}\left(1(\mathrm{}+1)\eta _i\right),$$
(60)
where $`\mathrm{\Delta }\varphi _n^{(k)}`$ refers to the shift in potential due to the connection of a new nodes to $`k`$ old ones (see Eq. (56) in Braunstein et al. (2003)).
The distribution $`\rho (\eta )`$ is determined, in the limit $`\mu \mathrm{}`$, by the self-consistent equation
$$\rho (\eta _0)=\frac{1}{Z}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)\underset{i=1}{\overset{k}{}}d\eta _i\rho (\eta _i)\delta (\eta _0\widehat{\eta }^{(k)}(\{\eta _i\}))\widehat{\mathrm{\Xi }}^{(k)}(\{\eta _i\})^xe^{z\chi (k,\gamma )},$$
(61)
with
$$\widehat{\eta }^{(k)}(\{\eta _i\}))\frac{1}{\widehat{\mathrm{\Xi }}^{(k)}(\{\eta _i\})}_{\mathrm{}=0}^{q1}(1)^{\mathrm{}}\left(\genfrac{}{}{0pt}{}{q1}{\mathrm{}}\right)_{i=1}^k(1(\mathrm{}+1)\eta _i).$$
(62)
This equation can be solved numerically using a population dynamics algorithm Mézard and Parisi (2001) : as for the typical case Braunstein et al. (2003), recovered here by taking $`x=0`$, a peak at $`\eta =0`$ is observed, so that $`\rho (\eta )`$ can be written $`\rho (\eta )=t\delta (\eta )+(1t)\stackrel{~}{\rho }(\eta )`$ where $`\stackrel{~}{\rho }`$ represents a continuous part. Rate functions $`L(\mathrm{\Sigma }_0)`$ obtained with this procedure are presented in Fig. 5.
Interestingly, for any value of $`\gamma `$ clustering and SAT-UNSAT transitions are found to occur within atypical graphs. These phase transitions are found by monitoring the parameter $`x`$, which, roughly speaking, characterizes the degree of frustration, with larger $`x`$ corresponding to less frustrated graphs. For a given $`\gamma `$, we indeed find thresholds $`x_c(\gamma )`$ and $`x_d(\gamma )`$, with $`x_c(\gamma )<x_d(\gamma )`$, such that for $`x<x_c(\gamma )`$ the graphs are UNSAT ($`\mathrm{\Sigma }_0<0`$) while for $`x>x_d(\gamma )`$ no more clustered solution is found ($`\mathrm{\Sigma }_0=0`$). The global phase diagram in the $`(\gamma ,x)`$ plane is presented in Fig. 7. The typical phase diagram of Fig. 7 can be read on the line $`x=0`$, with the thresholds $`\gamma _c`$ and $`\gamma _d`$ determined respectively by $`x_c(\gamma _c)=0`$ and $`x_d(\gamma _d)=0`$. We also expect that, for some values of $`x`$, the 1RSB Ansatz does not hold, in analogy to what is found in the typical case Krzakala et al. (2004) ; we however do not discuss this issue here, which could be handled by extending the stability analysis performed for typical instances Montanari and Ricci-Tersenghi (2003).
In cases such as coloring on regular graphs where the cavity equations are factorized, an interpretation of $`(y,\mu )`$ can be given in terms of the probability $`_N(ϵ=0)`$ for a graph to be SAT while lying in a class of typically UNSAT graphs (or conversely for a graph in a typically SAT ensemble to be UNSAT). For a RS system first, it has been shown in III.3 that the rate function $`L(ϵ)`$ is given by the negative 1RSB complexity $`\mathrm{\Sigma }(ϵ)`$, $`L(ϵ)=\mathrm{\Sigma }(ϵ)`$ ; in particular, the probability to be SAT for a graph in the typically UNSAT phase is $`_N(ϵ=0)e^{N\mathrm{\Sigma }_0}`$, where as before $`\mathrm{\Sigma }_0\mathrm{\Sigma }(ϵ=0)`$. In the typically UNSAT phase of a 1RSB system one still has $`\mathrm{\Sigma }_0<0`$, but the fluctuations of the complexity $`\mathrm{\Sigma }_0`$ from graph to graph described by the rate function $`L(\mathrm{\Sigma }_0)`$ have to be taken into account. The relation with $`_N(ϵ=0)`$ thus becomes
$$_N(ϵ=0)_{\mathrm{}}^0𝑑\mathrm{\Sigma }_0e^{NL(\mathrm{\Sigma }_0)}e^{N\mathrm{\Sigma }_0}$$
(63)
since $`e^{N\mathrm{\Sigma }_0}`$ must now be multiplied by the probability $`e^{NL(\mathrm{\Sigma }_0)}`$ to actually get a complexity $`\mathrm{\Sigma }_0`$. Two cases can then arise : the saddle point can lie at the boundary $`\mathrm{\Sigma }_0=0`$, in which case $`P_N(ϵ=0)e^{NL(\mathrm{\Sigma }_0=0)}`$, or it can be strictly negative, $`\mathrm{\Sigma }_0<0`$, in which case $`P_N(ϵ=0)e^{N[L(\mathrm{\Sigma }_0)\mathrm{\Sigma }_0]}`$ with $`\mathrm{\Sigma }_0`$ given by $`1=_{\mathrm{\Sigma }_0}L(\mathrm{\Sigma }_0)`$. An alternative formulation can be given with $`(x)`$ related to $`L(\mathrm{\Sigma }_0)`$ through
$$e^{N(x)}=𝑑\mathrm{\Sigma }_0e^{N[L(\mathrm{\Sigma }_0)x\mathrm{\Sigma }_0]}.$$
(64)
Since $`\mathrm{\Sigma }_0=_x(x)`$, one has to compute the $`x^{}`$ maximizing $`(x)`$, which is also associated to the saddle point $`x^{}=_{\mathrm{\Sigma }_0}L(\mathrm{\Sigma }_0=0)`$ : $`x^{}<1`$ corresponds to the first case with $`_N(ϵ=0)e^{N(x^{})}`$, while $`x^{}>1`$ corresponds to the second case with $`_N(ϵ=0)e^{N(1)}`$ (the very same mechanism underlies the selection of the 1RSB parameter in the typical cavity method at finite temperature Mézard and Parisi (2001)). Physically, this second situation, $`x^{}>1`$, refers to very rare graphs with extremely small frustration, which are thus in a RS phase \[it can be seen that $`(x=1)`$ indeed gives back the RS rate function $`L(ϵ=0)`$\]. In the phase diagram of a problem like coloring, such a situation is expected only for the largest values of $`\gamma `$.
As such, the interpretation applies only for models where the factorization holds. Otherwise negative complexities $`\mathrm{\Sigma }(ϵ)<0`$ should be interpreted as giving the probability $`e^{N\mathrm{\Sigma }(ϵ)}`$ for a model with typical complexity to have a cluster with energy $`ϵ`$, which, because of the interference between the internal disorder and the local external disorder, is not associated with a rate function relative to the external disorder only. The estimation of $`_N(ϵ=0)`$ along the lines presented above is however still amenable, but one has to perform a two-step large deviations analysis involving a second rate function $`L(L_0)`$ over the first rate function $`L_0`$ estimated under the RS assumption ; technically, this computation is quite similar to what has been done here for $`L(\mathrm{\Sigma }_0)`$, which can be taken as a factorized approximation of $`L(L_0)`$.
### IV.3 Multiple sources of disorder
In problems other that vertex-cover or coloring, the definition of an instance can include some quenched values of the coupling constants, as happens for spin-glass models or $`K`$-SAT. In this case, the energy shifts include a dependence on the couplings $`J`$, with functions $`\mathrm{\Delta }\widehat{E}_n^{(k;\{J_i\})}(\{h_i\})`$ and $`\mathrm{\Delta }\widehat{E}_{\mathrm{}}^{(J_{12})}(h_1,h_2)`$. The situation is described by a two-step large deviations principle,
$$_G[E_{G,J}=Nϵ]e^{N_G(ϵ)},[_G]e^{N𝒦(_G)}.$$
(65)
The rate functions are again computed through their Legendre transforms by considering two temperatures, $`y_J`$ and $`y_G`$,
$`(y_J,y_G)=`$ $`y_Gf𝒦(f,y_J),`$ $`y_G=_f𝒦(f,y_J),`$ (66)
$`y_Jf(y_J)=`$ $`y_Jϵ(ϵ),`$ $`y_J=_ϵ(ϵ).`$ (67)
where the factor $`y_J`$ in front of $`f`$ in the second line is introduced as in Eq. (47) to conform with the traditionnal definition of free energies. Taking Poissonian graphs as an example, we have explicitely the following expressions, to be compared with Eqs. (52) and (53):
$$\begin{array}{cc}& 𝒫[P_0]=\frac{1}{𝒵}\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)\underset{i=1}{\overset{k}{}}𝒟P_i𝒫[P_i]\delta \left[P_0\widehat{P}^{(k)}[\{P_i\}]\right]\widehat{Z}^{(k)}[\{P_i\}]^{y_G/y_J}e^{z\chi (k,\gamma )},\hfill \\ & P_0(h_0)=\widehat{P}^{(k)}[\{P_i\}](h_0)=\frac{1}{Z}𝔼_J\underset{i=1}{\overset{k}{}}dh_iP_i(h_i)\delta (h_0\widehat{h}^{(k,\{J_i\})}(\{h_i\}))e^{y_J\mathrm{\Delta }\widehat{E}_n^{(k;\{J_i\})}(\{h_i\})},\hfill \\ & Z=\widehat{Z}^{(k)}[\{P_i\}]=𝔼_J\underset{i=1}{\overset{k}{}}dh_iP_i(h_i)e^{y_J\mathrm{\Delta }\widehat{E}_n^{(k;\{J_i\})}(\{h_i\})},\hfill \\ & e^z=\underset{k=0}{\overset{\mathrm{}}{}}\sigma (k)\left[\underset{i=1}{\overset{2}{}}𝒟P_i𝒫[P_i]\left(𝔼_J\underset{i=1}{\overset{2}{}}dh_iP_i(h_i)e^{y_J\mathrm{\Delta }\widehat{E}_{\mathrm{}}^{(J_{12})}(h_1,h_2)}\right)^{y_G/y_J}\right]^k,\hfill \\ & (y_J,y_G,\gamma )=\mathrm{ln}𝒵=\mathrm{ln}\left[\underset{k=0}{\overset{\mathrm{}}{}}\pi _\gamma (k)\underset{i=1}{\overset{k}{}}𝒟P_i𝒫[P_i]\delta \left[P_0\widehat{P}^{(k)}[\{P_i\}]\right]\widehat{Z}^{(k)}[\{P_i\}]^{y_G/y_J}e^{z\chi (k,\gamma )}\right].\hfill \end{array}$$
(68)
where $`𝔼_J`$ denotes the average over the couplings $`J`$. Note that with $`y_G=0`$, we obtain large deviations with respect to the couplings on a typical graph, while with $`y_J=0`$, we obtain large deviations with respect to the graphs for typical couplings. With $`y_G=y_J`$, the two sources of disorder are treated at a same level, in analogy with $`y=\beta `$ in Sec. IV.1 and $`\mu =y`$, ($`x=1`$) in Sec. IV.2 (this prescription is sometimes referred to as the ”Nishimori temperature” Nishimori (2001)).
An interesting feature of the cavity equations is the possibility to implement them as a message-passing algorithm Mézard and Zecchina (2002) to study for example here the large deviations with respect to the couplings on a given graph. The message passed along the oriented edge $`(ij)`$ is the distribution $`P^{(ij)}`$, which, in particular cases, can be parameterized by a finite set of real numbers, and the update rule is
$$P^{(i0)}(h_i)=\frac{1}{Z}𝔼_J\underset{ji0}{}dh_jP^{(ji)}(h_j)\delta \left(h_0\widehat{h}^{(k,\{J_{ji}\})}(\{h_j\})\right)e^{y_J\mathrm{\Delta }\widehat{E}_n^{(k;\{J_{ji}\})}(\{h_j\})}$$
(69)
where the notation $`j0i`$ means that, on the particular graph considered, $`j`$ is a neighbor of $`0`$ different from $`i`$. This algorithmic scheme could be used to design graphs with optimal properties, with for instance applications in coding theory, in the context of low-density parity-check codes Richardson et al. (2001).
## V Conclusion
While statistical physics of disordered systems have so far mostly focused on the thermodynamical properties of samples which are typical with respect to the source of quenched disorder, we have shown here that its methods can be extended to describe large deviations, that is, atypical samples. Large deviations are of foremost interest in probability theory and the approach followed here, though admittedly non rigorous, is based on clearly formulated assumptions which should be amenable to mathematical justifications. In its simplest form indeed, the LDCM we exposed assumes no replica symmetry breaking, a situation in which many of its typical predictions have been proved to be exact Talagrand (2003).
From the perspective of algorithmic complexity, the LDCM can be seen as a first step in an attempt to reconcile the worst case and typical case analysis, usually regarded as antagonistic. However the scope of large deviations should not be regarded as restricted to optimization theory, as it notably allows to work out the statistical mechanics of physical systems with adaptive structures. An example of such system is constituted by random networks subject to mechanical constraints where the possibility to adapt leads to the occurrence of an intriguing intermediate phase, preceding a rigidity transition Barré et al. (2004a) ; an other example along the same lines is given by proteins whose structure is shaped by strong constraints Thorpe et al. (2001). Adaptative structures in random graphs are also of interest in the seemingly unrelated field of neural networks Wemmenhove and Skantzos (2004). Finally, in the spirit of the most impressive achievements of its typical counterpart Mézard and Zecchina (2002), it would also be interesting to implement the LDCM on particular ensembles of instances to analyze, and possibly design, graph structures with specific properties.
### Acknowledgments
I am grateful to Andrea Montanari for providing me with his notes Montanari (2002) on the analysis of large deviations by the replica method, and to Marc Mézard for his interest and encouragements.
|
warning/0506/hep-ph0506240.html
|
ar5iv
|
text
|
# 1 Introduction: W-Boson Longitudinal-Transverse Interference
## 1 Introduction: W-Boson Longitudinal-Transverse <br>Interference
In part because of the large top-quark mass and properties of QCD, W-boson polarimetry is a particularly powerful technique for empirical investigation of the $`tW^+b`$ decay mode from top-antitop pair-production data for the “charged-lepton plus jets” channel . For this channel, there is the sequential decay $`tW^+b(l^+\nu )b`$, with $`\overline{t}W^{}\overline{b}`$ in which the $`W^{}`$ decays into hadronic jets. Since the final state is the $`(l^+\nu )`$ decay product of the $`W^+`$, there are observable effects from $`W^+`$ boson longitudinal-transverse interference. For instance, a contribution to the angular-distribution intensity-function is the product of an amplitude in which the $`W^+`$ is longitudinally-polarized with the complex-conjugate of an amplitude in which the $`W^+`$ is transversely polarized, summed with the complex-conjugate of this product. The helicity formalism is a general method for investigating applications of W-boson interference in stage-two spin-correlation functions for describing the charged-lepton plus jets channel, and for the di-lepton plus jets channel.
Most of this paper consists of the derivation of general beam-referenced stage-two
spin-correlation functions (BR-S2SC) \[4-7\] for the analysis of top-antitop pair-production at the Tevatron , at the Large Hadron Collider , and/or at an International Linear Collider . However, as a simple result which illustrates W-boson longitudinal-transverse interference, for the charged-lepton-plus-jets reaction $`q\overline{q}t\overline{t}(W^+b)(W^{}\overline{b})(l^+\nu b)(W^{}\overline{b})`$ we have found that there is a 3-angle spin-correlation function for (i) determination of the relative sign of , or for (ii) measurement of a possible non-trivial phase between the two dominant $`\lambda _b=1/2`$ helicity amplitudes for the $`tW^+b`$ decay mode . For the $`CP`$-conjugate case, there is an analogous function and tests for $`\overline{t}W^{}\overline{b}`$ decay.
Tests for non-trivial phases in top-quark decays are important in searching for possible $`\stackrel{~}{T}_{FS}`$ violation. $`\stackrel{~}{T}_{FS}`$ invariance will be violated if either (i) there is a fundamental violation of canonical time-reversal invariance, and/or (ii) there are absorptive final-state interactions. For instance, unexpected final-state interactions might be associated with additional t-quark decay modes. To keep this assumption of “the absence of final-state interactions” manifest in comparison to a detailed-balance or other direct test for fundamental time-reversal invariance, we refer to this as $`\stackrel{~}{T}_{FS}`$ invariance, see . Measurement of a non-zero primed top-quark decay helicity parameter, such as $`\eta ^{}`$ or $`\omega ^{}`$, would imply $`\stackrel{~}{T}_{FS}`$ violation, see Appendix B. “Explicit $`\stackrel{~}{T}_{FS}`$ violation” will occur if there is an additional complex-coupling $`\frac{g_i}{2\mathrm{\Lambda }_i}`$ associated with a specific single additional Lorentz structure, $`i=S,P,S\pm P,\mathrm{}`$.
For the sequential decay $`tW^+b`$ followed by $`W^+l^+\nu `$, the spherical angles $`\theta _a`$, $`\varphi _a`$ specify the $`l^+`$ momentum in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame (see Fig. 1) when there is first a boost from the $`(t\overline{t})_{c.m.}`$ frame to the $`t_1`$ rest frame, and then a second boost from the $`t_1`$ rest frame to the $`W_{1}^{}{}_{}{}^{+}`$ rest frame. The $`0^o`$ direction for the azimuthal angle $`\varphi _a`$ is defined by the projection of the $`W_{2}^{}{}_{}{}^{}`$ momentum direction. Correspondingly (see Fig. 2) the spherical angles $`\theta _b`$, $`\varphi _b`$ specify the $`l^{}`$ momentum in the $`W_{2}^{}{}_{}{}^{}`$ rest frame when there is first a boost from the $`(t\overline{t})_{c.m.}`$ frame to the $`\overline{t_2}`$ rest frame, and then a second boost from the $`\overline{t_2}`$ rest frame to the $`W_{2}^{}{}_{}{}^{}`$ rest frame. The $`0^o`$ direction for the azimuthal angle $`\varphi _b`$ is defined by the projection of the $`W_{1}^{}{}_{}{}^{+}`$ momentum direction. As shown in Fig. 3, the two angles $`\theta _1^t`$, $`\theta _2^t`$ describe the $`W`$-boson momenta directions in the first stage of the sequential-decays of the $`t\overline{t}`$ system, in which $`t_1W_{1}^{}{}_{}{}^{+}b`$ and $`\overline{t_2}W_{2}^{}{}_{}{}^{}\overline{b}`$. Through out this paper, the subscripts “one” and “two” will be used to distinguish the two sequential-decay chains.
In the $`t_1`$ rest frame, the matrix element for $`t_1W_{1}^{}{}_{}{}^{+}b`$ is
$$\theta _1^t,\varphi _1,\lambda _{W^+},\lambda _b|\frac{1}{2},\lambda _1=D_{\lambda _1,\mu }^{(1/2)}(\varphi _1,\theta _1^t,0)A(\lambda _{W^+},\lambda _b)$$
(1)
where $`\mu =\lambda _{W^+}\lambda _b`$ in terms of the $`W_{1}^{}{}_{}{}^{+}`$ and $`b`$-quark helicities. Through out this paper an asterisk will denote complex conjugation. The final $`W_{1}^{}{}_{}{}^{+}`$ momentum is in the $`\theta _1^t,\varphi _1`$ direction and the $`b`$-quark momentum is in the opposite direction. The variable $`\lambda _1`$ gives the $`t_1`$-quark’s spin component quantized along the $`z_1^t`$ axis in Fig. 3. Upon a boost back to the $`(t\overline{t})_{cm}`$, or on further to the $`\overline{t_2}`$ rest frame, $`\lambda _1`$ also specifies the helicity of the $`t_1`$-quark. For the $`CP`$-conjugate process, $`\overline{t}_2W_{2}^{}{}_{}{}^{}\overline{b}`$, in the $`\overline{t}_2`$ rest frame the matrix element is
$$\theta _2^t,\varphi _2,\lambda _W^{},\lambda _{\overline{b}}|\frac{1}{2},\lambda _2=D_{\lambda _2,\overline{\mu }}^{(1/2)}(\varphi _2,\theta _2^t,0)B(\lambda _W^{},\lambda _{\overline{b}})$$
(2)
with $`\overline{\mu }=\lambda _W^{}\lambda _{\overline{b}}`$. By analogous argument, $`\lambda _2`$ is the $`\overline{t}_2`$ helicity.
In terms of the $`tW^+b`$ helicity amplitudes, the polarized-partial-widths and W-boson-LT-interference-widths are
$`\mathrm{\Gamma }(0,0)`$ $``$ $`\left|A(0,1/2)\right|^2,\mathrm{\Gamma }(1,1)\left|A(1,1/2)\right|^2`$ (3)
$`\mathrm{\Gamma }_R(0,1)`$ $`=`$ $`\mathrm{\Gamma }_R(1,0)Re[A(0,1/2)A(1,1/2)^{}]`$ (4)
$``$ $`|A(0,1/2)||A(1,1/2)|\mathrm{cos}\beta _L`$
$`\mathrm{\Gamma }_I(0,1)`$ $`=`$ $`\mathrm{\Gamma }_I(1,0)Im[A(0,1/2)A(1,1/2)^{}]`$ (5)
$``$ $`|A(0,1/2)||A(1,1/2)|\mathrm{sin}\beta _L`$
where the $`R`$, $`I`$ subscripts denote the real and imaginary parts which define the W-boson-LT-interference. The $`L`$ superscript on the $`\mathrm{\Gamma }^L(\lambda _W,\lambda _W^{^{}})`$’s has been conveniently suppressed in (3-5) for this is the dominant $`\lambda _b`$ helicity channel. By convention, the dominant $`L`$ superscript \[ $`R`$ superscript \] on $`\mathrm{\Gamma }^L(\lambda _W,\lambda _W^{^{}})`$ \[ $`\overline{\mathrm{\Gamma }}^R(\lambda _W,\lambda _W^{^{}})`$ \] will be suppressed in this paper. Note the two important minus-signs in the last two lines of (5). Here, following the conventions in , we define the moduli and phases as
$$A(\lambda _W,\lambda _b)\left|A(\lambda _W,\lambda _b)\right|\mathrm{exp}(ı\phi _{\lambda _W,\lambda _b})$$
(6)
with
$$\beta _L\phi _{1,\frac{1}{2}}\phi _{0,\frac{1}{2}},\beta _R\phi _{1,\frac{1}{2}}\phi _{0,\frac{1}{2}}$$
(7)
In terms of the $`\overline{t}W^{}\overline{b}`$ helicity amplitudes,
$`\overline{\mathrm{\Gamma }}(0,0)`$ $``$ $`\left|B(0,1/2)\right|^2,\overline{\mathrm{\Gamma }}(1,1)\left|B(1,1/2)\right|^2`$ (8)
$`\overline{\mathrm{\Gamma }}_R(0,1)`$ $`=`$ $`\overline{\mathrm{\Gamma }}_R(1,0)Re[B(0,1/2)B(1,1/2)^{}]`$ (9)
$``$ $`|B(0,1/2)||B(1,1/2)|\mathrm{cos}\overline{\beta }_R`$
$`\overline{\mathrm{\Gamma }}_I(0,1)`$ $`=`$ $`\overline{\mathrm{\Gamma }}_I(1,0)Im[B(0,1/2)B(1,1/2)^{}]`$ (10)
$``$ $`|B(0,1/2)||B(1,1/2)|\mathrm{sin}\overline{\beta }_R`$
with the moduli and phases defined by
$$B(\lambda _W,\lambda _{\overline{b}})\left|B(\lambda _W,\lambda \overline{_b})\right|\mathrm{exp}(ı\overline{\phi }_{\lambda _W,\lambda \overline{_b}})$$
(11)
with $`\overline{\beta }_R\overline{\phi }_{1,\frac{1}{2}}\overline{\phi }_{0,\frac{1}{2}}`$ and $`\overline{\beta }_L\overline{\phi }_{1,\frac{1}{2}}\overline{\phi }_{0,\frac{1}{2}}`$.
In this paper, we consider the production-decay sequence
$$q\overline{q},\mathrm{or}e\overline{e}t\overline{t}(W^+b)(W^{}\overline{b})\mathrm{}$$
(12)
At the Tevatron, this is the dominant contribution to $`t\overline{t}`$ production. The contribution from $`ggt\overline{t}(W^+b)(W^{}\overline{b})\mathrm{}`$ can be treated analogously. The latter is the dominant contribution at the LHC. The corresponding BR-S2SC functions for it will be reported separately .
We assume that the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ amplitudes dominate respectively in $`t_1`$ and $`\overline{t_2}`$ decay. In the SM and in the case of an additional large $`t_Rb_L`$ moment , the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ amplitudes are more than $`30`$ times larger than the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ amplitudes. The simple three-angle distribution $`|_0+|_{sig}`$ for $`t_1W_1^+b(l^+\nu )b`$ involves the angles $`\{\theta _2^t`$, $`\theta _a`$, $`\varphi _a\}`$ shown in Figs. 1-3.
$$|_0=\frac{16\pi ^3g^4}{9s^2}(1+\frac{2m_t^2}{s})\left\{\frac{1}{2}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a+\mathrm{\Gamma }(1,1)\mathrm{sin}^4\frac{\theta _a}{2}\right\}[\overline{\mathrm{\Gamma }}(0,0)+\overline{\mathrm{\Gamma }}(1,1)]$$
(13)
$`|_{sig}`$ $`=`$ $`{\displaystyle \frac{4\sqrt{2}\pi ^4g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}[\overline{\mathrm{\Gamma }}(0,0)+\overline{\mathrm{\Gamma }}(1,1)]`$ (14)
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\varphi _a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\varphi _a\right\}K`$
where $`K`$, $``$ are defined below.
The analogous three-angle S2SC function $`\overline{|}_0+\overline{|}_{sig}`$ for the $`CP`$-conjugate channel $`\overline{t}_2W_2^{}\overline{b}(l^{}\overline{\nu })\overline{b}`$ is a distribution versus $`\{\theta _1^t`$, $`\theta _b`$, $`\varphi _b\}`$ :
$`\overline{|}_0={\displaystyle \frac{16\pi ^3g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})\left\{{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}\right\}[\mathrm{\Gamma }(0,0)+\mathrm{\Gamma }(1,1)]`$ (15)
$`\overline{|}_{sig}`$ $`=`$ $`{\displaystyle \frac{4\sqrt{2}\pi ^4g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _1^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}[\mathrm{\Gamma }(0,0)+\mathrm{\Gamma }(1,1)]`$ (16)
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\varphi _b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\varphi _b\right\}K\overline{}`$
Note the important relative plus-sign between $`\overline{\mathrm{\Gamma }}_I(0,1)`$ and $`\overline{\mathrm{\Gamma }}_R(0,1)`$ in (16), in contrast to the relative minus-sign for $`\mathrm{\Gamma }_I(0,1)`$ and $`\mathrm{\Gamma }_R(0,1)`$ in (14).
### 1.1 Structure of three-angle S2SC functions
The “signal” contributions are suppressed by the factor
$$K\frac{(1\frac{2m_t^2}{s})}{(1+\frac{2m_t^2}{s})}$$
(17)
associated with the $`gt\overline{t}`$ production process, and the factor
$$\frac{[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]}{[\overline{\mathrm{\Gamma }}(0,0)+\overline{\mathrm{\Gamma }}(1,1)]},\overline{}\frac{[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]}{[\mathrm{\Gamma }(0,0)+\mathrm{\Gamma }(1,1)]}$$
(18)
associated with the stage-one part of the sequential-decay chains, $`\overline{t}W^{}\overline{b},tW^+b`$. Numerically, $`0.41`$ in both the standard model and in the case of an additional large $`t_Rb_L`$ chiral weak-transition moment . The appearance of the $`=(\mathrm{𝚙𝚛𝚘𝚋}W_L)(\mathrm{𝚙𝚛𝚘𝚋}W_T)`$ factor is not surprising because this is a consequence of the dynamical assumption that the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ amplitudes dominate. In the standard model $`=(1\frac{2m_W^2}{m_{t}^{}{}_{}{}^{2}})/(1+\frac{2m_W^2}{m_{t}^{}{}_{}{}^{2}})`$ whether there is or isn’t a large $`t_Rb_L`$ moment. Fortunately $`m_t\sqrt{2}m_W=+113GeV`$, otherwise many $`W`$-boson polarimetry effects would be absent in top-quark spin-correlation functions. An important exception is the $`\theta _a`$ dependence of $`|_0`$ \[ see (13)\]. Both of the $``$ and $`K`$ suppression factors are absent in purely stage-two $`W`$-boson polarimetry, with or without spin-correlation.
From the $`\theta _{2}^{}{}_{}{}^{t}`$ dependence of the integrated diagonal-elements of the sequential-decay density matrices for $`\overline{t_2}W_{2}^{}{}_{}{}^{}\overline{b}(l^{}\overline{\nu })\overline{b}`$, it follows that $``$’s numerator appears in $`|_{sig}`$ multiplied by $`\mathrm{cos}\theta _2^t`$ and that $``$’s denominator appears in $`|_0`$ multiplied by one \[ see (95-96)\]. Because the t-quark has spin $`\frac{1}{2}`$, there are purely half-angle $`d_{mm^{}}^{\frac{1}{2}}(\theta _{2}^{}{}_{}{}^{t})`$-squared intensity-product-factors in (95-97). The off-diagonal $`\overline{R}_{\lambda _2\lambda _2^{^{}}}`$ elements which describe $`\overline{t_2}`$-helicity interference do not contribute due to the integration over the opening-angle $`\varphi `$ between the $`t_1`$ and $`\overline{t_2}`$ decay planes. The angles $`\theta _{1,2}`$ are respectively equivalent to the $`W_{1,2}^{}{}_{}{}^{\pm }`$-boson energies in the $`(t\overline{t})_{cm}`$ (see Appendix A). In this 3-variable spin-correlation function, the minus sign in the numerator of the $`K`$ suppression factor in $`|_{sig}`$ is a consequence of the minus sign in the sequential-decay density-matrix $`𝐑_{++}^{b_L}`$ of (26) in the helicity-flip contribution (92) for the $`\overline{𝐑}_{++}`$ term, versus the corresponding plus sign in $`𝐑_{}^{b_L}`$ of (27) in the helicity-conserving contribution (72) for the $`\overline{𝐑}_{++}`$ term; and analogously for the $`\overline{𝐑}_{}`$ terms in (92) and (72).
### 1.2 Summary
From the top-quark spin-correlation function (13-14), the two tests for $`t_1W_{1}^{}{}_{}{}^{+}b`$ decay are:
(i) By measurement of $`\mathrm{\Gamma }_R(0,1)`$, the relative sign of the two dominant $`\lambda _b=1/2`$ helicity-amplitudes can be determined if their relative phase is $`0^0`$ or $`180^0`$. Versus the partial-decay-width $`\mathrm{\Gamma }(tW^+b)`$, W-boson longitudinal-transverse interference is a large effect for in the standard model $`\eta _L\frac{\mathrm{\Gamma }_R(0,1)}{\mathrm{\Gamma }}=\pm 0.46`$ without/with a large $`t_Rb_L`$ chiral weak-transition-moment. In both models, the probabilities for longitudinal/transverse W-bosons are large, $`P(W_L)=\frac{\mathrm{\Gamma }(0,0)}{\mathrm{\Gamma }}=0.70`$ and $`P(W_T)=\frac{\mathrm{\Gamma }(1,1)}{\mathrm{\Gamma }}=0.30`$, and so for a trivial relative-phase difference of $`0^0`$ or $`180^0`$, W-boson longitudinal-transverse interference must be a large effect.
(ii) By measurement of both $`\mathrm{\Gamma }_R(0,1)`$ and $`\mathrm{\Gamma }_I(0,1)`$ via the $`\varphi _a`$ dependence, a possible non-trivial phase can be investigated. Tests for non-trivial phases in top-quark decays are important in searching for possible $`\stackrel{~}{T}_{FS}`$ violation.
From (15-16), there are the analogous two tests for $`\overline{t}_2W_{2}^{}{}_{}{}^{}\overline{b}`$ decay. In the standard model $`\overline{\mathrm{\Gamma }}_R(0,1)=\mathrm{\Gamma }_R(0,1)`$, and both $`\overline{\mathrm{\Gamma }}_I(0,1)`$ and $`\mathrm{\Gamma }_I(0,1)`$ vanish whether there is or isn’t a purely-real $`t_Rb_L`$ transition-moment.
Section 2 of this paper contains the derivation of general BR-S2SC functions. For $`t\overline{t}`$ production by $`q\overline{q}`$, or $`e\overline{e}t\overline{t}`$, neither $`CP`$ invariance nor $`\stackrel{~}{T}_{FS}`$ invariance is assumed for the $`T(\lambda _1,\lambda _2)`$ helicity amplitudes in Sec. 2.2. For informative details, see . By $`CP`$ invariance, $`T(++)=T()`$ but $`T(+)`$ and $`T(+)`$ are unrelated. If experiment were to show that one of the primed production-helicity-parameters (76, 82-85, 94) is non-zero, then $`\stackrel{~}{T}_{FS}`$ invariance is violated in the $`gt\overline{t}`$ process.
In Section 3, these results are applied to the lepton-plus-jets channel of the $`t\overline{t}`$ system, assuming that the $`\lambda _b=1/2`$ and $`\lambda _{\overline{b}}=1/2`$ amplitudes dominate. Simple four-angle spin-correlation functions are obtained, which do not involve beam-referencing. These and other additional-angle generalizations might be useful empirically, for instance as checks with respect to the above four tests. Section 4 contains a discussion. The appendices respectively treat (A) kinematic formulas, (B) translation between this paper’s $`\mathrm{\Gamma }(\lambda _W,\lambda _{W}^{}{}_{}{}^{^{}})`$ notation and the helicity parameter’s notation of Refs. , (C) kinematic formulas for beam-referencing versus Figs. 1-2, and (D) formulas for $`e\overline{e}t\overline{t}`$ production.
## 2 Derivation of Beam-Referenced Stage-Two <br>Spin-Correlation Functions
In order to reference stage-two spin-correlation functions (S2SC) to the incident lepton or parton beam , we generalize the derivation of S2SC functions given in . When more data is available for top quark decays, it should be a reasonable further step to consider using the results of to incorporate $`\mathrm{\Lambda }_b`$ polarimetry. $`\mathrm{\Lambda }_b`$ polarimetry could be used to make a complete measurement of the four moduli and the three relative-phases of the helicity amplitudes in $`tW^+b`$ and analogously in $`\overline{t}W^{}\overline{b}`$. In this context, next-to-leading-order QCD, electroweak, and W-boson and t-quark finite-width corrections require further theoretical investigation . If the magnitudes of the two $`\lambda _b=1/2`$ helicity amplitudes are as predicted by the standard model, i.e. at factors of more than $`\frac{1}{30}`$ smaller than the two dominant $`\lambda _b=1/2`$ amplitudes, both detector and background effects will be non-trivial at this level of sensitivity at a hadron collider. Nevertheless, empirical consideration will be warranted if by then, there is compelling evidence for unusual top-quark physics.
In the BR-S2SC functions, we consider the decay sequence $`t_1W_{1}^{}{}_{}{}^{+}b`$ followed by $`W_{1}^{}{}_{}{}^{+}l^+\nu `$, and the $`CP`$-conjugate decay sequence $`\overline{t_2}W_{2}^{}{}_{}{}^{}\overline{b}`$ followed by $`W_{2}^{}{}_{}{}^{}l^{}\overline{\nu }`$. In Figs. 3 and 4, the spherical angles $`\theta _1^t`$ and $`\varphi _1`$ describe the $`W_{1}^{}{}_{}{}^{+}`$ momentum in the “first stage” $`t_1W_{1}^{}{}_{}{}^{+}b`$. Similarly, in Fig. 5 spherical angles $`\theta _a`$ and $`\stackrel{~}{\varphi _a}`$ describe the $`l^+`$ momentum in the “second stage” $`W_{1}^{}{}_{}{}^{+}l^+\nu `$ when there is first a boost from the $`(t\overline{t})_{cm}`$ frame to the $`t_1`$ rest frame, and then a second boost from the $`t_1`$ rest frame to the $`W_{1}^{}{}_{}{}^{+}`$ rest frame. If instead the boost to the $`W_{1}^{}{}_{}{}^{+}`$ rest frame is directly from the $`(t\overline{t})_{cm}`$ frame, one must account for Wigner rotations. Formulas and details about these Wigner rotations are given in Ref. . Analogously, two pairs of spherical angles $`\theta _{2}^{}{}_{}{}^{t},\varphi _2`$ and $`\theta _b`$, $`\stackrel{~}{\varphi _b}`$ specify the two stages in the $`CP`$-conjugate sequential decay $`\overline{t}W^{}\overline{b}`$ followed by $`W^{}l^{}\overline{\nu }`$ when the boost is from the $`\overline{t_2}`$ rest frame.
Note that the charged leptons’ azimuthal angle $`\stackrel{~}{\varphi _a}`$ in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame in Fig. 5, and analogously $`\stackrel{~}{\varphi _b}`$ in the $`W_{2}^{}{}_{}{}^{}`$ rest frame, are referenced respectively by the $`\overline{t_2}`$ and $`t_1`$ momentum directions. Instead of using the anti-top and top quark momenta for this purpose, one can reference these two azimuthal angles in terms of the opposite $`W^{}`$-boson momentum as in the formulas given in the introduction. These azimuthal angles are then denoted without “tilde accents” : $`\varphi _a`$ in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame when the boost is from the $`t_1`$ rest frame, and $`\varphi _b`$ in the $`W_{2}^{}{}_{}{}^{}`$ rest frame when the boost is from the $`\overline{t_2}`$ rest frame.
As discussed in the caption to Fig. 3, the momenta for $`t_1`$, $`W_{1}^{}{}_{}{}^{+}`$, and $`\overline{t_2}`$ lie in the same plane whether the analysis is in the $`t_1`$ rest frame, in the $`\overline{t_2}`$ rest frame, or in the $`t\overline{t}`$ center-of-momentum frame. Therefore, in deriving BR-S2SC functions in the helicity formalism, the angle $`\stackrel{~}{\varphi _a}`$ in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame is theoretically clear and simple. In general in the $`(t\overline{t})_{cm}`$ frame, the momenta for $`t_1`$, $`W_{1}^{}{}_{}{}^{+}`$ and $`W_{2}^{}{}_{}{}^{}`$ do not lie in the same plane. However, from the empirical point of view, the $`W_{2}^{}{}_{}{}^{}`$ momentum direction in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame will often be more precisely known, and so these two azimuthal angles without “tilde accents” will be more useful. From the standpoint of the helicity formalism, in the final S2SC functions either $`\varphi _a`$ or $`\stackrel{~}{\varphi _a}`$ can be used because it is only a matter of referencing the zero direction for the azimuthal angle, i.e. it is an issue concerning the specification of the Euler angles in the $`D`$ function for $`W^+l^+\nu `$ decay.
To simplify the notation, unlike in Refs. , in this paper we do not use “tilde accents” on the polar angles $`\theta _a`$ and $`\theta _b`$. We also do not use “$`t`$” superscripts on $`\varphi _{1,2}`$ for they are Lorentz invariant for each of the three frames considered in Fig. 3. On the other hand, “$`t`$” superscripts on $`\theta _{1,2}^{}{}_{}{}^{t}`$ for the $`t_1`$ and $`\overline{t_2}`$ rest frames, are necessary to distinguish these angles from $`\theta _{1,2}`$ which are defined in the $`(t\overline{t})_{cm}`$.
In the $`W_{1}^{}{}_{}{}^{+}`$ rest frame, the matrix element for $`W_{1}^{}{}_{}{}^{+}l^+\nu `$ \[ or for $`W_{1}^{}{}_{}{}^{+}j_{\overline{d}}j_u`$ \] is
$$\theta _a,\stackrel{~}{\varphi }_a,\lambda _{l^+},\lambda _\nu |1,\lambda _{W^+}=D_{\lambda _{W^+},1}^1(\stackrel{~}{\varphi }_a,\theta _a,0)c$$
(19)
since $`\lambda _\nu =\frac{1}{2},\lambda _{l^+}=\frac{1}{2}`$, neglecting $`(\frac{m_l}{m_W})`$ corrections \[ neglecting $`(\frac{m_{jet}}{m_W})`$ corrections\]. Since the amplitude “$`c`$” in this matrix element is independent of the helicities, we will suppress it in the following formulas since it only affects the overall normalization. We will use below
$$\rho _{\lambda _1\lambda _1^{^{}};\lambda _W\lambda _W^{^{}}}(tW^+b)=\underset{\lambda _b=1/2}{}D_{\lambda _1,\mu }^{(1/2)}(\varphi _1,\theta _1^t,0)D_{\lambda _1^{^{}},\mu ^{^{}}}^{(1/2)}(\varphi _1,\theta _1^t,0)A(\lambda _W,\lambda _b)A^{}(\lambda _W^{^{}},\lambda _b)$$
(20)
where $`\mu =\lambda _{W^+}\lambda _b`$ and $`\mu ^{^{}}=\lambda _{W^+}\lambda _b^{^{}}`$,
$$\rho _{\lambda _W\lambda _W^{^{}}}(W^+l^+\nu )=D_{\lambda _W,1}^1(\stackrel{~}{\varphi _a},\theta _a,0)D_{\lambda _W^{^{}},1}^1(\stackrel{~}{\varphi _a},\theta _a,0)$$
(21)
In the $`W_{2}^{}{}_{}{}^{}`$ rest frame, analogous to (19) the matrix element for $`W_{2}^{}{}_{}{}^{}l^{}\overline{\nu }`$ \[ $`W_{2}^{}{}_{}{}^{}j_{\overline{u}}j_d`$ \] is
$$\theta _b,\stackrel{~}{\varphi }_b,\lambda _l^{},\lambda _{\overline{\nu }}|1,\lambda _W^{}=D_{\lambda _W^{},1}^1(\stackrel{~}{\varphi }_b,\theta _b,0)\overline{c}$$
(22)
and we suppress the “$`\overline{c}`$” factor in the following.
### 2.1 Sequential-decay density matrices
The composite decay-density-matrix for $`t_1W_{1}^{}{}_{}{}^{+}b(l^+\nu )b`$ is
$$R_{\lambda _1\lambda _1^{^{}}}=\underset{\lambda _W,\lambda _W^{^{}}}{}\rho _{\lambda _1\lambda _1^{^{}};\lambda _W\lambda _W^{^{}}}(tW^+b)\rho _{\lambda _W\lambda _W^{^{}}}(W^+l^+\nu )$$
(23)
where $`\lambda _W,\lambda _W^{^{}}=0,\pm 1`$ and the $`\rho `$ density matrices are given in (20-21).
The above composite decay-density-matrix (23) can be expressed
$`𝐑=𝐑^{𝐛_𝐋}+𝐑^{𝐛_𝐑}`$ (24)
The $`\lambda _b=1/2`$ elements are
$`𝐑^{𝐛_𝐋}=\left(\begin{array}{cc}𝐑_{}^{𝐛_𝐋}{}_{++}{}^{}& e^{ı\varphi _1}𝐫_{}^{𝐛_𝐋}{}_{+}{}^{}\\ e^{ı\varphi _1}𝐫_{}^{𝐛_𝐋}{}_{+}{}^{}& 𝐑_{}^{𝐛_𝐋}{}_{}{}^{}\end{array}\right)`$ (27)
where
$`𝐑_{++}^{b_L}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{cos}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^2\theta _a+\mathrm{\Gamma }(1,1)\mathrm{sin}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}`$ (28)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _1^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`𝐑_{}^{b_L}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{sin}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^2\theta _a+\mathrm{\Gamma }(1,1)\mathrm{cos}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}`$ (29)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _1^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\mathrm{𝑅𝑒}(𝐫_+^{b_L})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Gamma }(0,0)\mathrm{sin}\theta _1^t\mathrm{sin}^2\theta _a{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(1,1)\mathrm{sin}\theta _1^t\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}`$ (30)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{cos}\theta _1^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$$Im(𝐫_+^{𝐛_𝐋})=\frac{1}{\sqrt{2}}[\mathrm{\Gamma }_R(0,1)\mathrm{sin}\stackrel{~}{\phi _a}+\mathrm{\Gamma }_I(0,1)\mathrm{cos}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _a\mathrm{sin}^2\frac{\theta _a}{2}$$
(31)
and $`𝐫_+^{b_L}=(𝐫_+^{b_L})^{}`$.
For the subdominant b<sub>R</sub> decay channel,
$$𝐑^{𝐛_𝐑}=\left(\begin{array}{cc}𝐑_{}^{𝐛_𝐑}{}_{++}{}^{}& e^{ı\varphi _1}𝐫_{}^{𝐛_𝐑}{}_{+}{}^{}\\ e^{ı\varphi _1}𝐫_{}^{𝐛_𝐑}{}_{+}{}^{}& 𝐑_{}^{𝐛_𝐑}{}_{}{}^{}\end{array}\right)$$
(32)
$`𝐑_{++}^{b_R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }^R(0,0)\mathrm{sin}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^2\theta _a+\mathrm{\Gamma }^R(1,1)\mathrm{cos}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{cos}^4{\displaystyle \frac{\theta _a}{2}}`$ (33)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R^R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}+\mathrm{\Gamma }_I^R(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _1^t\mathrm{sin}\theta _a\mathrm{cos}^2{\displaystyle \frac{\theta _a}{2}}`$
$`𝐑_{}^{b_R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }^R(0,0)\mathrm{cos}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{sin}^2\theta _a+\mathrm{\Gamma }^R(1,1)\mathrm{sin}^2{\displaystyle \frac{\theta _1^t}{2}}\mathrm{cos}^4{\displaystyle \frac{\theta _a}{2}}`$ (34)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R^R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}+\mathrm{\Gamma }_I^R(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _1^t\mathrm{sin}\theta _a\mathrm{cos}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\mathrm{𝑅𝑒}(𝐫_+^{b_R})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Gamma }^R(0,0)\mathrm{sin}\theta _1^t\mathrm{sin}^2\theta _a+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }^R(1,1)\mathrm{sin}\theta _1^t\mathrm{cos}^4{\displaystyle \frac{\theta _a}{2}}`$ (35)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\mathrm{\Gamma }_R^R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}+\mathrm{\Gamma }_I^R(0,1)\mathrm{sin}\stackrel{~}{\phi _a}]\mathrm{cos}\theta _1^t\mathrm{sin}\theta _a\mathrm{cos}^2{\displaystyle \frac{\theta _a}{2}}`$
$$Im(𝐫_+^{b_R})=\frac{1}{\sqrt{2}}[\mathrm{\Gamma }_R^R(0,1)\mathrm{sin}\stackrel{~}{\phi _a}\mathrm{\Gamma }_I^R(0,1)\mathrm{cos}\stackrel{~}{\phi _a}]\mathrm{sin}\theta _a\mathrm{cos}^2\frac{\theta _a}{2}$$
(36)
and $`𝐫_+^{b_R}=(𝐫_+^{b_R})^{}`$ . The b<sub>R</sub> decay channel’s polarized-partial-widths and
W-boson-LT-interference-widths are
$`\mathrm{\Gamma }^R(0,0)`$ $``$ $`\left|A(0,1/2)\right|^2,\mathrm{\Gamma }^R(1,1)\left|A(1,1/2)\right|^2`$ (37)
$`\mathrm{\Gamma }_R^R(0,1)`$ $`=`$ $`\mathrm{\Gamma }_R^R(1,0)Re[A(0,1/2)A(1,1/2)^{}]`$ (38)
$``$ $`|A(0,1/2)||A(1,1/2)|\mathrm{cos}\beta _R`$
$`\mathrm{\Gamma }_I^R(0,1)`$ $`=`$ $`\mathrm{\Gamma }_I^R(1,0)Im[A(0,1/2)A(1,1/2)^{}]`$ (39)
$``$ $`|A(0,1/2)||A(1,1/2)|\mathrm{sin}\beta _R`$
Note that the superscripts on these $`\mathrm{\Gamma }(\lambda _W,\lambda _{W}^{}{}_{}{}^{})`$’s always denote the $`b`$ or $`\overline{b}`$ helicity, whereas the subscripts denote the real or imaginary part (e.g. alternatively for (36) use $`\mathrm{\Gamma }_{\mathrm{𝑅𝑒}}^R(0,1)`$).
The analogous composite decay-density matrix for the $`CP`$-conjugate process
$`\overline{t}W^{}\overline{b}(l^{}\overline{\nu })\overline{b}`$ is
$`\overline{𝐑}=\overline{𝐑}^{\overline{𝐛}_𝐋}+\overline{𝐑}^{\overline{𝐛}_𝐑}`$ (40)
where the dominant
$$\overline{𝐑}^{\overline{𝐛}_𝐑}=\left(\begin{array}{cc}\overline{𝐑}_{}^{\overline{𝐛}_𝐑}{}_{++}{}^{}& e^{ı\varphi _2}\overline{𝐫}_{}^{\overline{𝐛}_𝐑}{}_{+}{}^{}\\ e^{ı\varphi _2}\overline{𝐫}_{}^{\overline{𝐛}_𝐑}{}_{+}{}^{}& \overline{𝐑}_{}^{\overline{𝐛}_𝐑}{}_{}{}^{}\end{array}\right)$$
(41)
$`\overline{𝐑}_{++}^{\overline{b}_R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}(1,1)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}`$ (42)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\phi _b}+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\overline{𝐑}_{}^{\overline{b}_R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}`$ (43)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\phi _b}+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\mathrm{𝑅𝑒}(\overline{𝐫}_+^{\overline{𝐛}_𝐑})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}\theta _2^t\mathrm{sin}^2\theta _b+{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}\theta _2^t\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}`$ (44)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\phi _b}+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{cos}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$$Im(\overline{𝐫}_+^{\overline{𝐛}_𝐑})=\frac{1}{\sqrt{2}}[\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{sin}\stackrel{~}{\phi _b}\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{cos}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _b\mathrm{sin}^2\frac{\theta _b}{2}$$
(45)
and $`\overline{𝐫}_+^{\overline{𝐛}_𝐑}=(\overline{𝐫}_+^{\overline{𝐛}_𝐑})^{}`$ .
For the subdominant $`\overline{𝐛}_L`$ decay channel,
$$\overline{𝐑}^{𝐛_𝐋}=\left(\begin{array}{cc}\overline{𝐑}_{}^{\overline{𝐛}_𝐋}{}_{++}{}^{}& e^{ı\varphi _2}\overline{𝐫}_{}^{\overline{𝐛}_𝐋}{}_{+}{}^{}\\ e^{ı\varphi _2}\overline{𝐫}_{}^{\overline{𝐛}_𝐋}{}_{+}{}^{}& \overline{𝐑}_{}^{\overline{𝐛}_𝐋}{}_{}{}^{}\end{array}\right)$$
(46)
$`\overline{𝐑}_{++}^{\overline{b}_L}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}^L(0,0)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}^L(1,1)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{cos}^4{\displaystyle \frac{\theta _b}{2}}`$ (47)
$`+{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R^L(0,1)\mathrm{cos}\stackrel{~}{\phi _b}\overline{\mathrm{\Gamma }}_I^L(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{cos}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\overline{𝐑}_{}^{\overline{b}_L}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}^L(0,0)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}^L(1,1)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}\mathrm{cos}^4{\displaystyle \frac{\theta _b}{2}}`$ (48)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R^L(0,1)\mathrm{cos}\stackrel{~}{\phi _b}\overline{\mathrm{\Gamma }}_I^L(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{cos}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\mathrm{𝑅𝑒}(\overline{𝐫}_+^{\overline{𝐛}_𝐋})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\overline{\mathrm{\Gamma }}^L(0,0)\mathrm{sin}\theta _2^t\mathrm{sin}^2\theta _b{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}^L(1,1)\mathrm{sin}\theta _2^t\mathrm{cos}^4{\displaystyle \frac{\theta _b}{2}}`$ (49)
$`{\displaystyle \frac{1}{\sqrt{2}}}[\overline{\mathrm{\Gamma }}_R^L(0,1)\mathrm{cos}\stackrel{~}{\phi _b}\overline{\mathrm{\Gamma }}_I^L(0,1)\mathrm{sin}\stackrel{~}{\phi _b}]\mathrm{cos}\theta _2^t\mathrm{sin}\theta _b\mathrm{cos}^2{\displaystyle \frac{\theta _b}{2}}`$
$$Im(\overline{𝐫}_+^{\overline{𝐛}_𝐋})=\frac{1}{\sqrt{2}}[\overline{\mathrm{\Gamma }}_R^L(0,1)\mathrm{sin}\stackrel{~}{\phi _b}+\overline{\mathrm{\Gamma }}_I^L(0,1)\mathrm{cos}\stackrel{~}{\phi _b}]\mathrm{sin}\theta _b\mathrm{cos}^2\frac{\theta _b}{2}$$
(50)
and $`\overline{𝐫}_+^{\overline{𝐛}_𝐋}=(\overline{𝐫}_+^{\overline{𝐛}_𝐋})^{}`$ .
$`\overline{\mathrm{\Gamma }}^L(0,0)`$ $``$ $`\left|B(0,1/2)\right|^2,\overline{\mathrm{\Gamma }}^L(1,1)\left|B(1,1/2)\right|^2`$ (51)
$`\overline{\mathrm{\Gamma }}_R^L(0,1)`$ $`=`$ $`\overline{\mathrm{\Gamma }}_R^L(1,0)\mathrm{𝑅𝑒}[B(0,1/2)B(1,1/2)^{}]`$ (52)
$``$ $`|B(0,1/2)||B(1,1/2)|\mathrm{cos}\overline{\beta }_L`$ (53)
$`\overline{\mathrm{\Gamma }}_I^L(0,1)`$ $`=`$ $`\overline{\mathrm{\Gamma }}_I^L(1,0)Im[B(0,1/2)B(1,1/2)^{}]`$ (54)
$``$ $`|B(0,1/2)||B(1,1/2)|\mathrm{sin}\overline{\beta }_L`$ (55)
Sometimes in the derivation, we will denote $`𝐫_+=F_a+ıH_a`$ and analogously
$`\overline{𝐫}_+=F_bıH_b`$ . As above, $`b_L`$ and $`b_R`$ superscripts on $`𝐫_+`$, and on $`F_a`$ and $`H_a`$ denote the $`\lambda _b=1/2,1/2`$ contributions, and analogously for $`\overline{𝐫}_+`$, $`F_b`$ and $`H_b`$.
### 2.2 Start of derivation of BR-S2SC functions
The general beam-referenced angular distribution in the $`(t\overline{t})_{cm}`$ is
$$\begin{array}{c}I(\mathrm{\Theta }_B,\mathrm{\Phi }_B;\theta _1^t,\varphi _1;\theta _a,\stackrel{~}{\varphi _a};\theta _2^t,\varphi _2;\theta _b,\stackrel{~}{\varphi _b})=_{\lambda _1\lambda _2\lambda _1^{^{}}\lambda _2^{^{}}}\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}(\mathrm{\Theta }_B,\mathrm{\Phi }_B)\\ \times R_{\lambda _1\lambda _1^{^{}}}(tW^+b\mathrm{})\overline{R}_{\lambda _2\lambda _2^{^{}}}(\overline{t}W^{}\overline{b}\mathrm{})\end{array}$$
(56)
where the summations are over the $`t_1`$ and $`\overline{t}_2`$ helicities. The composite decay-density-matrices $`R_{\lambda _1\lambda _1^{^{}}}`$ for $`tW^+b\mathrm{}`$ and $`\overline{R}_{\lambda _2\lambda _2^{^{}}}`$ for $`\overline{t}W^{}\overline{b}\mathrm{}`$ are given in the preceding subsection. This formula holds for any of the above $`t\overline{t}`$ production channels and for either the lepton-plus-jets, the dilepton-plus-jets, or the all-jets $`t\overline{t}`$ decay channels. The derivation begins in the “home” or starting coordinate system $`(x_h,y_h,z_h)`$ in the $`(t\overline{t})_{c.m.}`$ frame. As shown in Fig. 6-7, the angles $`\mathrm{\Theta }_B,\mathrm{\Phi }_B`$ specify the direction of the incident beam, the $`e`$ momentum, or in the case of $`p\overline{p}t\overline{t}X`$, the $`q`$ momentum arising from the incident $`p`$ in the $`p\overline{p}`$. The $`t_1`$ momentum is chosen to lie along the positive $`z_h`$ axis. The positive $`x_h`$ direction is an arbitrary, fixed perpendicular direction. Because the incident beam is assumed to be unpolarized, there is no dependence on the associated $`\varphi _1`$ angle after the observable azimuthal angles are specified (see below). With respect to the normalization of the various BR-S2SC functions, the $`\varphi _1`$ integration is not explicitly performed in this paper. With (54) there is an associated differential counting rate
$$\begin{array}{c}dN=I(\mathrm{\Theta }_B,\mathrm{\Phi }_B;\mathrm{})d(\mathrm{cos}\mathrm{\Theta }_B)d\mathrm{\Phi }_Bd(\mathrm{cos}\theta _1^t)d\varphi _1\\ d(\mathrm{cos}\theta _a)d\stackrel{~}{\varphi _a}d(\mathrm{cos}\theta _2^t)d\varphi _2d(\mathrm{cos}\theta _b)d\stackrel{~}{\varphi _b}\end{array}$$
(57)
where, for full phase space, the cosine of each polar angle ranges from -1 to 1, and each azimuthal angle ranges over $`2\pi `$.
For $`t\overline{t}`$ production by $`q\overline{q}`$, or $`e\overline{e}t\overline{t}`$ by initial unpolarized particles, the associated production density matrix is derived as in . It is
$`\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}`$ $`=`$ $`({\displaystyle \frac{1}{s^2}})e^{ı(\lambda ^{}\lambda )\mathrm{\Phi }_B}T(\lambda _1,\lambda _2)T^{}(\lambda _1^{^{}},\lambda _2^{^{}})`$ (58)
$`\times {\displaystyle \frac{1}{4}}{\displaystyle \underset{s_1,s_2}{}}|\stackrel{~}{T}(s_1,s_2)|^2d_{\lambda s}^1(\mathrm{\Theta }_B)d_{\lambda ^{^{}}s}^1(\mathrm{\Theta }_B)`$
where $`\lambda =\lambda _1\lambda _2`$, $`\lambda ^{^{}}=\lambda _1^{^{}}\lambda _2^{^{}}`$, and $`s=s_1s_2`$. In the body of this paper we concentrate on results for hadron colliders; formulas for the case of $`e\overline{e}`$ or $`\mu \overline{\mu }`$ production are given in Appendix D. It is convenient to separate the contributions into three parts, depending on the roles of the “helicity-conserving” and “helicity-flip” $`T(\lambda _1,\lambda _2)`$ amplitudes for $`gt_1\overline{t}_2`$ production. Relative to the helicity-conserving amplitudes, the helicity-flip amplitudes are $`(\sqrt{2}m_t/\sqrt{s})`$. We denote by a tilde accent the corresponding helicity-conserving light-quark $`q\overline{q}g`$ annihilation amplitudes. The values $`\lambda _{1,2}=\pm 1/2`$ of the arguments of $`T(\lambda _1,\lambda _2)`$ are denoted by the signs of $`\lambda _1`$, $`\lambda _2`$, and likewise for $`\stackrel{~}{T}(s_1,s_2)`$.
#### 2.2.1 Helicity-conserving contribution
The $`t_1\overline{t}_2`$ helicity-conserving contribution production density matrix is
$`\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}`$ $``$ $`\delta _{\lambda _2,\lambda _1}\delta _{\lambda _2^{^{}},\lambda _1^{^{}}}({\displaystyle \frac{1}{s^2}})e^{ı2(\lambda _1^{^{}}\lambda _1)\mathrm{\Phi }_B}T(\lambda _1,\lambda _1)T^{}(\lambda _1^{^{}},\lambda _1^{^{}})`$ (59)
$`\times {\displaystyle \frac{1}{4}}[|\stackrel{~}{T}(+)|^2d_{\lambda 1}^1(\mathrm{\Theta }_B)d_{\lambda ^{^{}}1}^1(\mathrm{\Theta }_B)+|\stackrel{~}{T}(+)|^2d_{\lambda ,1}^1(\mathrm{\Theta }_B)d_{\lambda ^{^{}},1}^1(\mathrm{\Theta }_B)]`$
where $`\lambda =2\lambda _1`$ and $`\lambda ^{^{}}=2\lambda _1^{^{}}`$. The angular distribution of (57) has four different terms which can be labelled as $`I_{\lambda ,\lambda ^{^{}}}`$ due to the Kronecker $`\delta `$’s. Explicitly, these are
$`I_{++}={\displaystyle \frac{1}{4s^2}}|T(+)|^2𝐑_{++}\overline{𝐑}_{}[|\stackrel{~}{T}(+)|^2\mathrm{cos}^4(\mathrm{\Theta }_B/2)+|\stackrel{~}{T}(+)|^2\mathrm{sin}^4(\mathrm{\Theta }_B/2)]`$ (60)
$`I_{}={\displaystyle \frac{1}{4s^2}}|T(+)|^2𝐑_{}\overline{𝐑}_{++}[|\stackrel{~}{T}(+)|^2\mathrm{sin}^4(\mathrm{\Theta }_B/2)+|\stackrel{~}{T}(+)|^2\mathrm{cos}^4(\mathrm{\Theta }_B/2)]`$ (61)
$`I_+={\displaystyle \frac{1}{4s^2}}T(+)T^{}(+)e^{ı(2\mathrm{\Phi }_R+\varphi )}𝐫_+\overline{𝐫}_+[|\stackrel{~}{T}(+)|^2+|\stackrel{~}{T}(+)|^2]\mathrm{cos}^2(\mathrm{\Theta }_B/2)\mathrm{sin}^2(\mathrm{\Theta }_B/2)`$ (62)
$`I_+={\displaystyle \frac{1}{4s^2}}T(+)T^{}(+)e^{ı(2\mathrm{\Phi }_R+\varphi )}𝐫_+\overline{𝐫}_+[|\stackrel{~}{T}(+)|^2+|\stackrel{~}{T}(+)|^2]\mathrm{cos}^2(\mathrm{\Theta }_B/2)\mathrm{sin}^2(\mathrm{\Theta }_B/2)`$ (63)
where the starting angles $`\varphi _2`$ and $`\mathrm{\Phi }_B`$ have been replaced by the angles $`\varphi =\varphi _1+\varphi _2`$ and $`\mathrm{\Phi }_R=\mathrm{\Phi }_B\varphi _1`$, see Figs. 6-7.
Two rotations are needed to recast the above expressions in terms of the angles of the final $`(t\overline{t})_{c.m.}`$ coordinate system shown in Figs. 1-2:
Step1: We rotate by $`\theta _1`$ so that the new z-axis $`\overline{z}`$ is along the $`W_1^+`$ momentum, as shown in Figs. 8-9.
This replaces the $`\mathrm{\Theta }_B,\mathrm{\Phi }_B`$ referencing of the beam direction by the final polar angle $`\theta _q`$ and an associated azimuthal $`\mathrm{\Phi }_W`$ variable. Since this is simply a coordinate rotation,
$$d(\mathrm{cos}\theta _q)d\mathrm{\Phi }_W=d(\mathrm{cos}\mathrm{\Theta }_B)d\mathrm{\Phi }_R$$
(64)
The Jacobian is 1, and $`\mathrm{cos}\theta _q`$ and $`\mathrm{\Phi }_W`$ have the usual range for spherical coordinates. The formulas for making this change of variables are:
$`\mathrm{cos}\theta _q`$ $`=`$ $`\mathrm{cos}\theta _1\mathrm{cos}\mathrm{\Theta }_B+\mathrm{sin}\theta _1\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Phi }_R`$ (65)
$`\mathrm{sin}\theta _q\mathrm{cos}\mathrm{\Phi }_W`$ $`=`$ $`\mathrm{sin}\theta _1\mathrm{cos}\mathrm{\Theta }_B+\mathrm{cos}\theta _1\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Phi }_R`$ (66)
$`\mathrm{sin}\theta _q\mathrm{sin}\mathrm{\Phi }_W`$ $`=`$ $`\mathrm{sin}\mathrm{\Theta }_B\mathrm{sin}\mathrm{\Phi }_R`$ (67)
and
$$\mathrm{cos}\mathrm{\Theta }_B=\mathrm{cos}\theta _1\mathrm{cos}\theta _q\mathrm{sin}\theta _1\mathrm{sin}\theta _q\mathrm{cos}\mathrm{\Phi }_W$$
(68)
In Fig. 9, the $`W_2^{}`$ momentum is at angles $`\mathrm{\Theta }_2`$ and $`\mathrm{\Phi }_2`$ . Since $`\mathrm{\Theta }_2=\pi \psi `$, $`\mathrm{\Theta }_2`$ can be replaced by the opening angle $`\psi `$ between the $`W_1^+`$ and $`W_2^{}`$ momenta. The opening angle $`\psi `$ is simply related to the important angle $`\varphi =\varphi _1+\varphi _2`$ between the $`t_1`$ and $`\overline{t}_2`$ decay planes:
$`\mathrm{cos}\psi `$ $`=`$ $`\mathrm{cos}\mathrm{\Theta }_2=\mathrm{cos}\theta _1\mathrm{cos}\theta _2+\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\varphi `$ (69)
$`\mathrm{sin}\psi `$ $`=`$ $`\mathrm{sin}\mathrm{\Theta }_2=(1\mathrm{cos}^2\mathrm{\Theta }_2)^{1/2}`$ (70)
On the other hand, $`\mathrm{cos}\mathrm{\Phi }_2`$ and $`\mathrm{sin}\mathrm{\Phi }_2`$ are auxiliary variables that appear in the formulas in Appendix C for transforming the initial beam-referencing spherical angles $`\mathrm{\Theta }_B,\mathrm{\Phi }_R`$ of Figs. 6-7 to the final ones, $`\theta _q,\varphi _q`$ of Figs. 1-2.
$`\mathrm{sin}\psi \mathrm{cos}\mathrm{\Phi }_2`$ $`=`$ $`\mathrm{sin}\theta _1\mathrm{cos}\theta _2+\mathrm{cos}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\varphi `$ (71)
$`\mathrm{sin}\psi \mathrm{sin}\mathrm{\Phi }_2`$ $`=`$ $`\mathrm{sin}\theta _2\mathrm{sin}\varphi `$ (72)
Step 2: We rotate by $`\mathrm{\Phi }_2`$ about $`\overline{z}=\widehat{z}`$ so that the $`W_2^{}`$ momenta is in the positive $`\widehat{x}`$ plane, as shown in Figs. 1-2.
By this rotation,
$$\varphi _q=\mathrm{\Phi }_W+\mathrm{\Phi }_2$$
(73)
so the Jacobian is 1, and $`\varphi _q`$ has the full $`2\pi `$ range.
By these two steps, the above four helicity-conserving contributions are expressed in terms of Figs. 1-2:
$`I_{++}+I_{}`$ $`=`$ $`{\displaystyle \frac{1}{16s^2}}S_q\{|T(+)|^2𝐑_{++}\overline{𝐑}_{}+|T(+)|^2𝐑_{}\overline{𝐑}_{++}\}(1+\mathrm{cos}^2\mathrm{\Theta }_B)`$ (74)
$`+{\displaystyle \frac{1}{8s^2}}T_q\{|T(+)|^2𝐑_{++}\overline{𝐑}_{}|T(+)|^2𝐑_{}\overline{𝐑}_{++}\}\mathrm{cos}\mathrm{\Theta }_B`$
$`I_++I_+`$ $`=`$ $`{\displaystyle \frac{1}{8s^2}}S_q\left\{\overline{\kappa }\left[F_aF_b+H_aH_b\right]+\overline{\kappa }^{^{}}\left[F_aH_bH_aF_b\right]\right\}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )`$ (75)
$`{\displaystyle \frac{1}{8s^2}}S_q\left\{\overline{\kappa }^{^{}}\left[F_aF_b+H_aH_b\right]\overline{\kappa }\left[F_aH_bH_aF_b\right]\right\}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}(2\mathrm{\Phi }_R+\varphi )`$
where
$`S_q`$ $`=`$ $`|\stackrel{~}{T}(+)|^2+|\stackrel{~}{T}(+)|^2`$ (76)
$`T_q`$ $`=`$ $`|\stackrel{~}{T}(+)|^2|\stackrel{~}{T}(+)|^2`$ (77)
$`\overline{\kappa }+ı\overline{\kappa }^{^{}}`$ $`=`$ $`T(+)T^{}(+)`$ (78)
#### 2.2.2 Mixed helicity-properties contribution
The mixed helicity-properties contribution of the $`t_1\overline{t}_2`$ production density matrix is in two parts: The first part is
$`\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}`$ $``$ $`\delta _{\lambda _2,\lambda _1}\delta _{\lambda _2^{^{}},\lambda _1^{^{}}}({\displaystyle \frac{1}{s^2}})e^{ı2\lambda _1^{^{}}\mathrm{\Phi }_B}T(\lambda _1,\lambda _1)T^{}(\lambda _1^{^{}},\lambda _1^{^{}})`$ (79)
$`\times {\displaystyle \frac{1}{4}}[|\stackrel{~}{T}(+)|^2d_{0,1}^1(\mathrm{\Theta }_B)d_{\lambda ^{^{}},1}^1(\mathrm{\Theta }_B)+|\stackrel{~}{T}(+)|^2d_{0,1}^1(\mathrm{\Theta }_B)d_{\lambda ^{^{}},1}^1(\mathrm{\Theta }_B)]`$
where $`\lambda ^{^{}}=2\lambda _1^{^{}}`$.
As in the above subsection for the helicity-conserving contribution, this mixed-helicity properties contribution can be expressed as the sum of
$`I_{++}^{mA}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\eta }^++ı\overline{\eta }^{}_{}{}^{}+)𝐑_{++}(F_b+ıH_b)(S_q\mathrm{cos}\mathrm{\Theta }_B+T_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı(\mathrm{\Phi }_R+\varphi )}`$ (80)
$`I_{}^{mA}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\omega }^{}+ı\overline{\omega }^{{}_{}{}^{}})𝐑_{}(F_bıH_b)(S_q\mathrm{cos}\mathrm{\Theta }_BT_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı(\mathrm{\Phi }_R+\varphi )}`$ (81)
$`I_+^{mA}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\omega }^++ı\overline{\omega }^{}_{}{}^{}+)(F_a+ıH_a)\overline{𝐑}{}_{++}{}^{}(S_q\mathrm{cos}\mathrm{\Theta }_BT_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı\mathrm{\Phi }_R}`$ (82)
$`I_+^{mA}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\eta }^{}+ı\overline{\eta }^{{}_{}{}^{}})(F_aıH_a)\overline{𝐑}{}_{}{}^{}(S_q\mathrm{cos}\mathrm{\Theta }_B+T_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı\mathrm{\Phi }_R}`$ (83)
where
$`\overline{\omega }^++ı\overline{\omega }^{}_{}{}^{}+`$ $`=`$ $`T(++)T^{}(+)`$ (84)
$`\overline{\omega }^{}+ı\overline{\omega }^{{}_{}{}^{}}`$ $`=`$ $`T()T^{}(+)`$ (85)
$`\overline{\eta }^++ı\overline{\eta }^{}_{}{}^{}+`$ $`=`$ $`T(++)T^{}(+)`$ (86)
$`\overline{\eta }^{}+ı\overline{\eta }^{{}_{}{}^{}}`$ $`=`$ $`T()T^{}(+)`$ (87)
The second part of the $`t_1\overline{t}_2`$ mixed helicity-properties part of the production density matrix is
$`\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}`$ $``$ $`\delta _{\lambda _2,\lambda _1}\delta _{\lambda _2^{^{}},\lambda _1^{^{}}}({\displaystyle \frac{1}{s^2}})e^{ı2\lambda _1\mathrm{\Phi }_B}T(\lambda _1,\lambda _1)T^{}(\lambda _1^{^{}},\lambda _1^{^{}})`$ (88)
$`\times {\displaystyle \frac{1}{4}}[|\stackrel{~}{T}(+)|^2d_{\lambda ,1}^1(\mathrm{\Theta }_B)d_{0,1}^1(\mathrm{\Theta }_B)+|\stackrel{~}{T}(+)|^2d_{\lambda ,1}^1(\mathrm{\Theta }_B)d_{0,1}^1(\mathrm{\Theta }_B)]`$
where $`\lambda =2\lambda _1`$ . This mixed-helicity-properties contribution can be expressed as the sum of
$`I_{++}^{mB}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\eta }^+ı\overline{\eta }^{}_{}{}^{}+)𝐑_{++}(F_bıH_b)(S_q\mathrm{cos}\mathrm{\Theta }_B+T_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı(\mathrm{\Phi }_R+\varphi )}`$ (89)
$`I_{}^{mB}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\omega }^{}ı\overline{\omega }^{{}_{}{}^{}})𝐑_{}(F_b+ıH_b)(S_q\mathrm{cos}\mathrm{\Theta }_BT_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı(\mathrm{\Phi }_R+\varphi )}`$ (90)
$`I_+^{mB}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\eta }^{}ı\overline{\eta }^{{}_{}{}^{}})(F_a+ıH_a)\overline{𝐑}{}_{}{}^{}[(S_q\mathrm{cos}\mathrm{\Theta }_B+T_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı\mathrm{\Phi }_R}`$ (91)
$`I_+^{mB}`$ $`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}s^2}}(\overline{\omega }^+ı\overline{\omega }^{}_{}{}^{}+)(F_aıH_a)\overline{𝐑}{}_{++}{}^{}(S_q\mathrm{cos}\mathrm{\Theta }_BT_q)\mathrm{sin}\mathrm{\Theta }_Be^{ı\mathrm{\Phi }_R}`$ (92)
#### 2.2.3 Helicity-flip contribution
The $`t_1\overline{t}_2`$ helicity-flip production density matrix is
$`\rho _{\lambda _1\lambda _2;\lambda _1^{^{}}\lambda _2^{^{}}}^{\mathrm{𝚙𝚛𝚘𝚍}}`$ $``$ $`\delta _{\lambda _2,\lambda _1}\delta _{\lambda _2^{^{}},\lambda _1^{^{}}}({\displaystyle \frac{1}{s^2}})T(\lambda _1,\lambda _1)T^{}(\lambda _1^{^{}},\lambda _1^{^{}})`$ (93)
$`\times {\displaystyle \frac{1}{4}}[|\stackrel{~}{T}(+)|^2d_{01}^1(\mathrm{\Theta }_B)d_{01}^1(\mathrm{\Theta }_B)+|\stackrel{~}{T}(+)|^2d_{0,1}^1(\mathrm{\Theta }_B)d_{0,1}^1(\mathrm{\Theta }_B)]`$
This contribution can be expressed as the sum of
$$I_{++}^{m2}+I_{}^{m2}=\frac{1}{8s^2}S_q\{|T(++)|^2𝐑_{++}\overline{𝐑}_{++}+|T()|^2𝐑_{}\overline{𝐑}_{}\}\mathrm{sin}^2\mathrm{\Theta }_B$$
(94)
and
$`I_+^{m2}+I_+^{m2}`$ $`=`$ $`{\displaystyle \frac{1}{4s^2}}S_q(\{\overline{\zeta }[F_aF_bH_aH_b]+\overline{\zeta }^{^{}}[F_aH_b+H_aF_b]\}\mathrm{cos}\varphi `$ (95)
$`+\{\overline{\zeta }^{^{}}[F_aF_bH_aH_b]+\overline{\zeta }[F_aH_b+H_aF_b]\}\mathrm{sin}\varphi )\mathrm{sin}^2\mathrm{\Theta }_B`$
where
$$\overline{\zeta }+ı\overline{\zeta }^{^{}}=T(++)T^{}()$$
(96)
For $`q\overline{q}t\overline{t}`$, in the Jacob-Wick phase convention, the associated helicity amplitudes are $`\stackrel{~}{T}(+,)=\stackrel{~}{T}(,+)=g,`$ the helicity-conserving $`T(+)=T(+)=g,`$ and the helicity-flip $`T(++)=T()=gm_t\sqrt{2/s}`$.
## 3 Lepton-plus-Jets Channel: $`\lambda _b=1/2`$, $`\lambda _{\overline{b}}=+1/2`$ <br>Dominance
From the perspective of specific helicity amplitude tests, one can use the above results to investigate various BR-S2SC functions for the lepton-plus-jets channel: In this paper, we are interested in tests for the relative sign of, or for measurement of a possible non-trivial phase between the $`\lambda _b=1/2`$ helicity amplitudes for $`tW^+b`$. We assume that the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ contributions dominate.
### 3.1 $`t_1W_1^+b(l^+\nu )b`$
For the case $`t_1W_1^+b(l^+\nu )b`$, with $`W_2^{}`$ decaying into hadronic jets, we separate the intensity contributions into two parts: “signal terms” $`\stackrel{~}{I}|_{sig}`$ which depend on $`\mathrm{\Gamma }_R(0,1)`$ and $`\mathrm{\Gamma }_I(0,1)`$, and “background terms” $`\stackrel{~}{I|}_0`$ which depend on $`\mathrm{\Gamma }(0,0)`$ and $`\mathrm{\Gamma }(1,1)`$. We use a tilde accent on $`\stackrel{~}{I|}_0,\mathrm{}`$ to denote the integration over the $`\theta _b`$, $`\stackrel{~}{\varphi }_b`$ variables. This integration gives
$`{\displaystyle _1^1}d(\mathrm{cos}\theta _b){\displaystyle _0^{2\pi }}𝑑\stackrel{~}{\varphi }_b\overline{𝐑}_{++}^{\overline{b}_R}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}[\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}+\overline{\mathrm{\Gamma }}(1,1)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}]`$ (97)
$`{\displaystyle _1^1}d(\mathrm{cos}\theta _b){\displaystyle _0^{2\pi }}𝑑\stackrel{~}{\varphi }_b\overline{𝐑}_{}^{\overline{b}_R}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}[\overline{\mathrm{\Gamma }}(0,0)\mathrm{cos}^2{\displaystyle \frac{\theta _2^t}{2}}+\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^2{\displaystyle \frac{\theta _2^t}{2}}]`$ (98)
$$_1^1d(\mathrm{cos}\theta _b)_0^{2\pi }𝑑\stackrel{~}{\varphi }_bF_b^{\overline{b}_R}=\frac{2\pi }{3}\mathrm{sin}\theta _2^t[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]$$
(99)
The integration over $`H_b^{\overline{b}_R}`$ vanishes.
We find for the helicity-conserving contribution,
$`(\stackrel{~}{I}_{++}+\stackrel{~}{I}_{})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{12s^2}}(1+\mathrm{cos}^2\mathrm{\Theta }_B)`$
$`\left\{\begin{array}{c}\frac{1}{2}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a[\overline{\mathrm{\Gamma }}(0,0)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\mathrm{\Gamma }(1,1)\mathrm{sin}^4\frac{\theta _a}{2}[\overline{\mathrm{\Gamma }}(0,0)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`(\stackrel{~}{I}_{++}+\stackrel{~}{I}_{})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{6\sqrt{2}s^2}}(1+\mathrm{cos}^2\mathrm{\Theta }_B)\mathrm{sin}\theta _1^t\mathrm{cos}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a+\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`(\stackrel{~}{I}_++\stackrel{~}{I}_+)|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{12s^2}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )\mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a\mathrm{\Gamma }(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`(\stackrel{~}{I}_++\stackrel{~}{I}_+)|_{sig}`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{6\sqrt{2}s^2}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$ (108)
$`\left\{\begin{array}{c}\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )\mathrm{cos}\theta _1^t\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}\\ +\mathrm{sin}(2\mathrm{\Phi }_R+\varphi )\left\{\mathrm{\Gamma }_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a+\mathrm{\Gamma }_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\right\}\end{array}\right\}`$
For the mixed-helicity contribution, the terms with primed coefficients \[see (82-85)\] all vanish. We collect the other mixed-helicity contributions in real sums:
$`\stackrel{~}{I}^{m(\overline{\omega }^++\overline{\eta }^{})}|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Phi }_R\mathrm{sin}\theta _1^t\mathrm{cos}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a\mathrm{\Gamma }(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`\stackrel{~}{I}^{m(\overline{\omega }^++\overline{\eta }^{})}|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$ (113)
$`\left\{\begin{array}{c}\mathrm{cos}\theta _1^t\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}\mathrm{cos}\mathrm{\Phi }_R\\ +\left\{\mathrm{\Gamma }_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a+\mathrm{\Gamma }_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\right\}\mathrm{sin}\mathrm{\Phi }_R\end{array}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`\stackrel{~}{I}^{m(\overline{\omega }^{}+\overline{\eta }^+)}|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\mathrm{cos}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a\mathrm{\Gamma }(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`\stackrel{~}{I}^{m(\overline{\omega }^{}+\overline{\eta }^+)}|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
The helicity-flip contributions are
$`(\stackrel{~}{I}_{++}^{m2}+\stackrel{~}{I}_{}^{m2})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B`$ (119)
$`\left\{\begin{array}{c}\frac{1}{2}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a[\overline{\mathrm{\Gamma }}(0,0)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\mathrm{\Gamma }(1,1)\mathrm{sin}^4\frac{\theta _a}{2}[\overline{\mathrm{\Gamma }}(0,0)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`(\stackrel{~}{I}_{++}^{m2}+\stackrel{~}{I}_{}^{m2})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}\theta _1^t\mathrm{cos}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`(\stackrel{~}{I}_+^{m2}+\stackrel{~}{I}_+^{m2})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}\varphi \mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a+\mathrm{\Gamma }(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
$`(\stackrel{~}{I}_+^{m2}+\stackrel{~}{I}_+^{m2})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}\theta _2^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$ (125)
$`\left\{\begin{array}{c}\mathrm{cos}\varphi \mathrm{cos}\theta _1^t\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a+\mathrm{\Gamma }_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a\right\}\\ +\mathrm{sin}\varphi \left\{\mathrm{\Gamma }_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_a+\mathrm{\Gamma }_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_a\right\}\end{array}\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
### 3.2 $`\overline{t}_2W_2^{}\overline{b}(l^{}\overline{\nu })\overline{b}`$
For the $`CP`$-conjugate process $`\overline{t}_2W_2^{}\overline{b}(l^{}\overline{\nu })\overline{b}`$, with $`W_1^+`$ decaying into hadronic jets, we similarly separate the contributions: “signal terms” $`\stackrel{~}{\overline{I}|}_{sig}`$depending on $`\overline{\mathrm{\Gamma }}_R(0,1)`$ and $`\overline{\mathrm{\Gamma }}_I(0,1)`$, and “background terms” $`\stackrel{~}{\overline{I}}|_0`$ depending on $`\overline{\mathrm{\Gamma }}(0,0)`$ and $`\overline{\mathrm{\Gamma }}(1,1)`$. The integration over $`\theta _a`$, $`\stackrel{~}{\varphi }_a`$ gives
$$_1^1d(\mathrm{cos}\theta _a)_0^{2\pi }𝑑\stackrel{~}{\varphi }_a𝐑_{++}^{b_L}=\frac{4\pi }{3}[\mathrm{\Gamma }(0,0)\mathrm{cos}^2\frac{\theta _1^t}{2}+\mathrm{\Gamma }(1,1)\mathrm{sin}^2\frac{\theta _1^t}{2}]$$
(126)
$$_1^1d(\mathrm{cos}\theta _a)_0^{2\pi }𝑑\stackrel{~}{\varphi }_a𝐑_{}^{b_L}=\frac{4\pi }{3}[\mathrm{\Gamma }(0,0)\mathrm{sin}^2\frac{\theta _1^t}{2}+\mathrm{\Gamma }(1,1)\mathrm{cos}^2\frac{\theta _1^t}{2}]$$
(127)
$$_1^1d(\mathrm{cos}\theta _a)_0^{2\pi }𝑑\stackrel{~}{\varphi }_aF_a^{b_L}=\frac{2\pi }{3}\mathrm{sin}\theta _1^t[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]$$
(128)
The integration over $`H_a^{b_L}`$ vanishes.
We find for the helicity-conserving contribution,
$`(\stackrel{~}{\overline{I}}_{++}+\stackrel{~}{\overline{I}}_{})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{12s^2}}(1+\mathrm{cos}^2\mathrm{\Theta }_B)`$ (132)
$`\left\{\begin{array}{c}\frac{1}{2}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b[\mathrm{\Gamma }(0,0)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4\frac{\theta _b}{2}[\mathrm{\Gamma }(0,0)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`(\stackrel{~}{\overline{I}}_{++}+\stackrel{~}{\overline{I}}_{})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{6\sqrt{2}s^2}}(1+\mathrm{cos}^2\mathrm{\Theta }_B)\mathrm{cos}\theta _1^t\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`(\stackrel{~}{\overline{I}}_++\stackrel{~}{\overline{I}}_+)|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{12s^2}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )\mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`(\stackrel{~}{\overline{I}}_++\stackrel{~}{\overline{I}}_+)|_{sig}`$ $`=`$ $`{\displaystyle \frac{\pi g^4}{6\sqrt{2}s^2}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}\theta _1^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$ (138)
$`\left\{\begin{array}{c}\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )\mathrm{cos}\theta _2^t\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}\\ \mathrm{sin}(2\mathrm{\Phi }_R+\varphi )\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b\right\}\end{array}\right\}`$
The mixed-helicity contributions are
$`\stackrel{~}{\overline{I}}^{m(\overline{\omega }^++\overline{\eta }^{})}|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Phi }_R\mathrm{sin}\theta _1^t\mathrm{cos}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`\stackrel{~}{\overline{I}}^{m(\overline{\omega }^++\overline{\eta }^{})}|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Phi }_R\mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`\stackrel{~}{\overline{I}}^{m(\overline{\omega }^{}+\overline{\eta }^+)}|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\mathrm{cos}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`\stackrel{~}{\overline{I}}^{m(\overline{\omega }^{}+\overline{\eta }^+)}|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t}{3s^2\sqrt{s}}}\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B\mathrm{cos}\theta _1^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$ (145)
$`\left\{\begin{array}{c}\mathrm{cos}\theta _2^t\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\\ +\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b\right\}\mathrm{sin}(\mathrm{\Phi }_R+\varphi )\end{array}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
The helicity-flip contributions are
$`(\stackrel{~}{\overline{I}}_{++}^{m2}+\stackrel{~}{\overline{I}}_{}^{m2})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B`$ (149)
$`\left\{\begin{array}{c}\frac{1}{2}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b[\mathrm{\Gamma }(0,0)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4\frac{\theta _b}{2}[\mathrm{\Gamma }(0,0)(1+\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`(\stackrel{~}{\overline{I}}_{++}^{m2}+\stackrel{~}{\overline{I}}_{}^{m2})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}\theta _1^t\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`(\stackrel{~}{\overline{I}}_+^{m2}+\stackrel{~}{\overline{I}}_+^{m2})|_0`$ $`=`$ $`{\displaystyle \frac{\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}\varphi \mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t`$
$`\left\{{\displaystyle \frac{1}{2}}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4{\displaystyle \frac{\theta _b}{2}}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
$`(\stackrel{~}{\overline{I}}_+^{m2}+\stackrel{~}{\overline{I}}_+^{m2})|_{sig}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi g^4m_t^2}{3s^3}}\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}\theta _1^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$ (155)
$`\left\{\begin{array}{c}\mathrm{cos}\varphi \mathrm{cos}\theta _2^t\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\right\}\\ \mathrm{sin}\varphi \left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{sin}\stackrel{~}{\varphi }_b\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{cos}\stackrel{~}{\varphi }_b\right\}\end{array}\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
### 3.3 $`\mathrm{\Gamma }(\lambda _W,\lambda _{W}^{}{}_{}{}^{^{}})`$ tests versus angular dependence
In summary, with beam-referencing, for the $`t_1W_1^+b(l^+\nu )b`$ case there are six “background terms” depending on $`\mathrm{\Gamma }(0,0)`$ and $`\mathrm{\Gamma }(1,1)`$, and also six “signal terms” depending on $`\mathrm{\Gamma }_{R,I}(0,1)`$. As a consequence of Lorentz invariance, there are associated kinematic factors with simple angular dependence which can be used to isolate and measure these four $`\mathrm{\Gamma }^{}`$s:
(i) $`\theta _a`$ polar-angle dependence:
The coefficients of $`\mathrm{\Gamma }(0,0)/\mathrm{\Gamma }(1,1)/\mathrm{\Gamma }_{R,I}(0,1)`$ vary relatively as the $`W`$-decay
$`d_{mm^{}}^1(\theta _a)`$-squared-intensity-ratios
$`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _a/\left[\mathrm{sin}^4{\displaystyle \frac{\theta _a}{2}}\right]/\left\{{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}\right\}=`$
$`2(1+\mathrm{cos}\theta _a)/\left[1\mathrm{cos}\theta _a\right]/\left\{\sqrt{2(1+\mathrm{cos}\theta _a)(1\mathrm{cos}\theta _a)}=\sqrt{2}\mathrm{sin}\theta _a\right\}`$ (156)
(ii) $`\varphi _a`$ azimuthal-angle dependence in the “signal terms” \[ or $`\stackrel{~}{\varphi }_a`$ dependence if $`\overline{t}_2`$ is used to specify the $`0^o`$ direction\] :
The coefficients of $`\mathrm{\Gamma }_R(0,1)/\mathrm{\Gamma }_I(0,1)`$ vary as
$$\mathrm{cos}\varphi _a/\mathrm{sin}\varphi _a$$
(157)
in each of the signal terms. However, in three terms there are also $`\mathrm{\Gamma }_{R,I}(0,1)`$’s with the opposite association of these $`\mathrm{cos}\varphi _a`$, $`\mathrm{sin}\varphi _a`$ factors. This opposite association occurs in $`(\stackrel{~}{I}_++\stackrel{~}{I}_+)|_{sig}`$, $`\stackrel{~}{I}^{m(\overline{\omega }^++\overline{\eta }^{})}|_{sig}`$, and $`(\stackrel{~}{I}_+^{m2}+\stackrel{~}{I}_+^{m2})|_{sig}`$, along with a different $`\mathrm{\Phi }_R`$ and $`\varphi `$ dependence which might be useful empirically in separation from the terms with the normal $`\varphi _a`$ association.
To reduce the number of angles, we integrate out the two beam-referencing angles, and also $`\varphi `$:
$`\stackrel{~}{}_i{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _1^1}d(\mathrm{cos}\mathrm{\Theta }_B){\displaystyle _0^{2\pi }}𝑑\mathrm{\Phi }_R\stackrel{~}{I}_i`$ (158)
This yields four-angle S2SC functions.
In terms of $`K`$ defined in (17), the four-angle distribution $`\{\theta _1^t`$, $`\theta _2^t`$, $`\theta _a`$, $`\varphi _a\}`$ is
$`\stackrel{~}{}|_0`$ $`=`$ $`{\displaystyle \frac{8\pi ^3g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})`$
$`\left\{\begin{array}{c}\frac{1}{2}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _a[\overline{\mathrm{\Gamma }}(0,0)(1+K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4\frac{\theta _a}{2}[\overline{\mathrm{\Gamma }}(0,0)(1K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\overline{\mathrm{\Gamma }}(1,1)(1+K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`\stackrel{~}{}|_{sig}`$ $`=`$ $`{\displaystyle \frac{8\sqrt{2}\pi ^3g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _2^tK\mathrm{sin}\theta _1^t\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\varphi _a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\varphi _a\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
The terms in these expressions arise from the helicity-conserving $`(\stackrel{~}{I}_{++}+\stackrel{~}{I}_{})`$, and from the helicity-flip $`(\stackrel{~}{I}_{++}^{m2}+\stackrel{~}{I}_{}^{m2})`$. In each case there are contributions to both background and signal parts.
Without the integration over $`\varphi `$, there is a contribution to both the background and signal parts from the helicity-flip $`(\stackrel{~}{I}_+^{m2}+\stackrel{~}{I}_+^{m2})`$ of (108-9). This additional contribution has both the normal and opposite $`\varphi _a`$ dependence as discussed above in (ii). It will be fundamentally significant to empirically demonstrate in both $`\mathrm{cos}\varphi `$ and $`\mathrm{sin}\varphi `$ the presence of this contribution to the spin-correlation because it arises completely from the combination of $`t_1`$-quark L-R interference and $`\overline{t}_2`$-antiquark L-R interference \[see (93)\]. Without the $`\varphi `$ dependence, in the above four-angle function (128-9) there is no contribution from the off-diagonal elements of the $`\lambda _b=1/2`$ and $`\lambda \overline{_b}`$ $`=1/2`$ sequential decay matrices (25) and (39).
For the $`CP`$-conjugate case in terms of $`\{\theta _2^t`$, $`\theta _1^t`$, $`\theta _b`$, $`\varphi _b\}`$, the analogous four-angle distributions are
$`\stackrel{~}{\overline{}}|_0`$ $`=`$ $`{\displaystyle \frac{8\pi ^3g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})`$
$`\left\{\begin{array}{c}\frac{1}{2}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b[\mathrm{\Gamma }(0,0)(1+K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\\ +\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4\frac{\theta _b}{2}[\mathrm{\Gamma }(0,0)(1K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)+\mathrm{\Gamma }(1,1)(1+K\mathrm{cos}\theta _1^t\mathrm{cos}\theta _2^t)]\end{array}\right\}`$
$`\stackrel{~}{\overline{}}|_{sig}`$ $`=`$ $`{\displaystyle \frac{8\sqrt{2}\pi ^3g^4}{9s^2}}(1+{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _1^tK\mathrm{sin}\theta _2^t\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\varphi _b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\varphi _b\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
The still simpler three-angle distributions, which were discussed in the introduction section, then follow if the $`\mathrm{cos}\theta _1^t`$ integration is also performed $`_i_1^1d(\mathrm{cos}\theta _1^t)\stackrel{~}{}_i`$:
$$|_0=\frac{16\pi ^3g^4}{9s^2}(1+\frac{2m_t^2}{s})\left\{\frac{1}{2}\mathrm{\Gamma }(0,0)\mathrm{sin}^2\theta _a+\mathrm{\Gamma }(1,1)\mathrm{sin}^4\frac{\theta _a}{2}\right\}[\overline{\mathrm{\Gamma }}(0,0)+\overline{\mathrm{\Gamma }}(1,1)]$$
(167)
$`|_{sig}`$ $`=`$ $`{\displaystyle \frac{8\pi ^4g^4}{9s^2}}(1{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _2^t{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta _a\mathrm{sin}^2{\displaystyle \frac{\theta _a}{2}}`$
$`\left\{\mathrm{\Gamma }_R(0,1)\mathrm{cos}\varphi _a\mathrm{\Gamma }_I(0,1)\mathrm{sin}\varphi _a\right\}[\overline{\mathrm{\Gamma }}(0,0)\overline{\mathrm{\Gamma }}(1,1)]`$
The analogous three-angle S2SC function for the $`CP`$-conjugate $`\overline{t}_2W_2^{}\overline{b}(l^{}\nu )\overline{b}`$ is
$$\overline{|}_0=\frac{16\pi ^3g^4}{9s^2}(1+\frac{2m_t^2}{s})\left\{\frac{1}{2}\overline{\mathrm{\Gamma }}(0,0)\mathrm{sin}^2\theta _b+\overline{\mathrm{\Gamma }}(1,1)\mathrm{sin}^4\frac{\theta _b}{2}\right\}[\mathrm{\Gamma }(0,0)+\mathrm{\Gamma }(1,1)]$$
(169)
$`\overline{|}_{sig}={\displaystyle \frac{8\pi ^4g^4}{9s^2}}(1{\displaystyle \frac{2m_t^2}{s}})\mathrm{cos}\theta _1^t{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta _b\mathrm{sin}^2{\displaystyle \frac{\theta _b}{2}}`$ (170)
$`\left\{\overline{\mathrm{\Gamma }}_R(0,1)\mathrm{cos}\varphi _b+\overline{\mathrm{\Gamma }}_I(0,1)\mathrm{sin}\varphi _b\right\}[\mathrm{\Gamma }(0,0)\mathrm{\Gamma }(1,1)]`$
## 4 Discussion
In the above derivation of general BR-S2SC functions, in part for greater generality, we include beam-referencing. At hadron colliders, beam-referencing may be useful in some applications. In the case of $`e\overline{e}`$-production, it would probably be useful in investigating possible anomalous initial-state-with-final-state couplings in the $`t_1\overline{t}_2`$ production process. However, the simple three-angle formulas reported in the introduction section do not make use of beam-referencing. Given the conceptual simplicity of the helicity formulation for $`q\overline{q},\mathrm{or}e\overline{e}t\overline{t}(W^+b)(W^{}\overline{b})\mathrm{}`$, such non-beam-referenced functions are ideal for tests of the moduli and phases of the four $`tW^+b`$ helicity amplitudes. While usage of direct boosts from the $`(t\overline{t})_{c.m.}`$ frame to the $`W^+`$ or $`W^{}`$ rest frames might be useful for some purposes, from the perspective of this BR-S2SC helicity formulation, such boosts will be an unnecessary complication. The boosts introduce additional Wigner rotations which obscure the overall simplicity of the helicity formulation which distinctly separates the different physics stages of the $`t\overline{t}`$ production and decay sequences.
In this paper we separate the $`\lambda _b=1/2`$ contributions from the $`\lambda _b=1/2`$ contributions. To display the $`W`$-boson polarization and longitudinal-transverse interference effects, we introduce a transparent $`\mathrm{\Gamma }^{\lambda _b}(\lambda _W,\lambda _W^{^{}})`$ notation. Appendix B relates this notation to the helicity parameters notation used in . At the present time, the $`\lambda _b=1/2`$ amplitudes do indeed appear to dominate in the $`tW^+b`$ decay mode and so the present paper’s $`\mathrm{\Gamma }^{\lambda _b}(\lambda _W,\lambda _W^{^{}})`$ notation is very appropriate. At a later date, in higher precision experiments where effects from all four of the decay amplitudes must be carefully considered, the helicity parameters notation might be useful. It is more analogous to the notation of the Michel-parameters which continue to be used in muon decay data analysis. On the other hand, in the context of a characterization of fundamental “particle properties”, the present $`\mathrm{\Gamma }^{\lambda _b}(\lambda _W,\lambda _W^{^{}})`$ notation is a simple way to precisely specify polarized-partial-width measurements, including $`W`$-boson longitudinal-transverse interference. Since the $`tW^+b`$ decay channel will first be investigated at hadron colliders, such measurements will be of channel polarized-partial-width branching ratios
$`B^{\lambda _b}(\lambda _W,\lambda _W^{^{}})=\mathrm{\Gamma }^{\lambda _b}(\lambda _W,\lambda _W^{^{}})/\mathrm{\Gamma }(tW^+b)`$ (171)
where $`\mathrm{\Gamma }(tW^+b)`$ is the partial width for $`tW^+b`$.
Acknowledgments
One of us (CAN) thanks top-quark experimentalists and theorists for discussions. This work was partially supported by U.S. Dept. of Energy Contract No. DE-FG 02-86ER40291.
## Appendix A Appendix: Kinematic Formulas
In the $`(t\overline{t})_{c.m.}`$ frame, the angles $`\theta _{1,2}`$ of the $`W_1^+`$, $`W_2^{}`$ and their respective energies $`E_{1,2}`$ are related by
$$2\stackrel{~}{P}p_W\mathrm{cos}\theta _{1,2}=2\stackrel{~}{P}_0E_{1,2}m_t^2m_W^2$$
(172)
where $`t`$-energy and magnitude of $`t`$-momentum are $`\stackrel{~}{P}_0=\sqrt{s}/2`$, $`\stackrel{~}{P}=\sqrt{\stackrel{~}{P}_0^2m_t^2}`$, and $`p_W^2=E_{1,2}^2m_W^2`$. In the $`t_1`$ rest frame, $`\overline{t}_2`$ rest frame, respectively
$$\theta _{1,2}^t=\mathrm{arccos}[\frac{\sqrt{s}(m_t^2+m_W^2)+4E_{1,2}m_t^2}{(m_t^2m_W^2)\sqrt{s4m_t^2}}],0\theta _{1,2}^t\pi $$
(173)
which give the kinematic limits
$$E_{1,2}^{\mathrm{max},\mathrm{min}}=\frac{\sqrt{s}(m_t^2+m_W^2)}{4m_t^2}\pm \frac{\sqrt{s}(m_t^2m_W^2)}{4m_t^2}[1\frac{4m_t^2}{s}]^{1/2}$$
(174)
The angles $`\theta _{1,2}`$ are determined uniquely from $`\mathrm{cos}\theta _{1,2}`$ and $`\mathrm{sin}\theta _{1,2}`$ of
$`p_{1,2}\mathrm{cos}\theta _{1,2}`$ $`=`$ $`\gamma (p_{1,2}^t\mathrm{cos}\theta _{1,2}^t+\beta E_{1,2}^t)`$ (175)
$`p_{1,2}\mathrm{sin}\theta _{1,2}`$ $`=`$ $`p_{1,2}^t\mathrm{sin}\theta _{1,2}^t`$ (176)
where $`p_{1,2}^t=(m_t^2m_W^2)/2m_t`$, $`E_{1,2}^t=\sqrt{(p_{1,2}^t)^2+m_W^2}`$, and $`\gamma =\sqrt{s}/(2m_t)`$, $`\beta `$ are for the relativistic boosts between the $`(t\overline{t})_{c.m.}`$ frame and the $`t_1`$, $`\overline{t}_2`$ rest frames. A check is $`E_{1,2}=\gamma (E_{1,2}^t+\beta p_{1,2}^t\mathrm{cos}\theta _{1,2}^t)`$.
From $`\theta _{1,2}`$ there is a unique relation between $`\mathrm{cos}\psi `$ and $`\mathrm{cos}\varphi `$,
$$\mathrm{cos}\psi =\mathrm{cos}\theta _1\mathrm{cos}\theta _2+\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\varphi $$
(177)
or equivalently from $`\theta _{1,2}^t`$
$$\mathrm{sin}\theta _1^t\mathrm{sin}\theta _2^t\mathrm{cos}\varphi =\frac{4m_t^2}{(m_t^2m_W^2)^2}\left\{\begin{array}{c}p_1p_2\mathrm{cos}\psi \\ +\frac{(\sqrt{s}E_1m_t^2m_W^2)(\sqrt{s}E_2m_t^2m_W^2)}{s4m_t^2}\end{array}\right\}$$
(178)
The sign of the quantity $`\mathrm{sin}\varphi `$ is the same as the sign of the auxiliary variable $`\mathrm{sin}\mathrm{\Phi }_2`$.
## Appendix B Appendix: Translation Between $`\mathrm{\Gamma }(\lambda _W,\lambda _{W}^{}{}_{}{}^{^{}})`$’s Notation and Helicity Parameter’s of Refs.
For the $`tW^+b`$ helicity amplitudes, in terms of the helicity-parameters of Refs. , the $`\lambda _b=1/2`$ polarized-partial-widths and W-boson-LT-interference-widths are
$`\mathrm{\Gamma }(0,0)`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }}{4}}\{1+\xi +\zeta +\sigma \}`$ (179)
$`\mathrm{\Gamma }(1,1)`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }}{4}}\{1+\xi \zeta \sigma \}`$ (180)
$`\mathrm{\Gamma }_R(0,1)`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }}{2}}\{\eta +\omega \}=\mathrm{\Gamma }\eta _L`$ (181)
$`\mathrm{\Gamma }_I(0,1)`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }}{2}}\{\eta ^{^{}}+\omega ^{^{}}\}=\mathrm{\Gamma }\eta _{L}^{}{}_{}{}^{^{}}`$ (182)
where the $`L`$ superscript is suppressed, and $`\mathrm{\Gamma }`$ is the partial width for $`tW^+b`$. For $`\overline{t}W^{}\overline{b}`$, the analogous formulas $`\lambda _{\overline{b}}=1/2`$ polarized-partial-widths and W-boson-LT-interference-widths are obtained by replacing $`1`$ $`+1`$ in the $`\mathrm{\Gamma }`$’s on the left-hand-sides, and then barring all of the $`\mathrm{\Gamma }`$’s on both sides and also barring all the helicity parameters.
The important $``$ suppression factor in (18) was denoted as $`S_W`$ in these references.
## Appendix C Appendix: $`\mathrm{\Theta }_B`$ , $`\mathrm{\Phi }_R`$ to $`\theta _q`$ , $`\varphi _q`$ Formulas
The transformation formulas to express the beam spherical angles $`\mathrm{\Theta }_B`$ , $`\mathrm{\Phi }_R`$ in terms of $`\theta _q`$ , $`\varphi _q`$ involve the $`(t\overline{t})_{c.m.}`$ W-boson angles $`\theta _1`$, $`\theta _2`$, and also the auxiliary variables $`\mathrm{sin}\mathrm{\Phi }_2`$ and $`\mathrm{cos}\mathrm{\Phi }_2`$ of (69-70) \[ see Figs. 8-9\]. In the helicity-conserving contributions
$`\mathrm{cos}\mathrm{\Theta }_B`$ $`=`$ $`𝒫_1+𝒬_1`$ (183)
$`𝒫_1`$ $`=`$ $`\mathrm{cos}\theta _1\mathrm{cos}\theta _q\mathrm{cos}\varphi _q\mathrm{sin}\theta _1\mathrm{sin}\theta _q\mathrm{cos}\mathrm{\Phi }_2,𝒬_1=\mathrm{sin}\varphi _q\mathrm{sin}\theta _1\mathrm{sin}\theta _q\mathrm{sin}\mathrm{\Phi }_2`$
$`(1+\mathrm{cos}^2\mathrm{\Theta }_B)`$ $`=`$ $`𝒫_0+𝒬_0`$ (184)
$`𝒫_0`$ $`=`$ $`1+\mathrm{cos}^2\theta _1\mathrm{cos}^2\theta _q+{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _q`$
$`\mathrm{cos}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{sin}\theta _1\mathrm{cos}\mathrm{\Phi }_2+{\displaystyle \frac{1}{2}}\mathrm{cos}2\varphi _q\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _q\mathrm{cos}2\mathrm{\Phi }_2`$
$`𝒬_0`$ $`=`$ $`\mathrm{sin}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{sin}\theta _1\mathrm{sin}\mathrm{\Phi }_2+{\displaystyle \frac{1}{2}}\mathrm{sin}2\varphi _q\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _q\mathrm{sin}2\mathrm{\Phi }_2`$
$`\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{cos}(2\mathrm{\Phi }_R+\varphi )`$ $`=`$ $`𝒫_\kappa +𝒬_\kappa `$ (185)
$`𝒫_\kappa `$ $`=`$ $`𝒞\mathrm{cos}\varphi +𝒮^{^{}}\mathrm{sin}\varphi ,𝒬_\kappa =𝒮\mathrm{cos}\varphi 𝒞^{^{}}\mathrm{sin}\varphi `$
$`\mathrm{sin}^2\mathrm{\Theta }_B\mathrm{sin}(2\mathrm{\Phi }_R+\varphi )`$ $`=`$ $`𝒫_^\kappa +𝒬_^\kappa `$ (186)
$`𝒫_^\kappa `$ $`=`$ $`𝒞^{^{}}\mathrm{cos}\varphi +𝒮\mathrm{sin}\varphi ,𝒬_^\kappa =𝒮^{^{}}\mathrm{cos}\varphi +𝒞\mathrm{sin}\varphi `$
where
$`𝒞`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _1(3\mathrm{cos}^2\theta _q1)`$
$`+\mathrm{cos}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{sin}\theta _1\mathrm{cos}\mathrm{\Phi }_2+{\displaystyle \frac{1}{2}}\mathrm{cos}2\varphi _q\mathrm{sin}^2\theta _q[1+\mathrm{cos}^2\theta _1]\mathrm{cos}2\mathrm{\Phi }_2`$
$`𝒮=\mathrm{sin}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{sin}\theta _1\mathrm{sin}\mathrm{\Phi }_2+{\displaystyle \frac{1}{2}}\mathrm{sin}2\varphi _q\mathrm{sin}^2\theta _q[1+\mathrm{cos}^2\theta _1]\mathrm{sin}2\mathrm{\Phi }_2`$
$`𝒞^{^{}}=\mathrm{sin}\varphi _q\mathrm{sin}2\theta _q\mathrm{sin}\theta _1\mathrm{cos}\mathrm{\Phi }_2+\mathrm{sin}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{cos}\theta _1\mathrm{cos}2\mathrm{\Phi }_2`$
$`𝒮^{^{}}=\mathrm{cos}\varphi _q\mathrm{sin}2\theta _q\mathrm{sin}\theta _1\mathrm{sin}\mathrm{\Phi }_2+\mathrm{cos}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{cos}\theta _1\mathrm{sin}2\mathrm{\Phi }_2`$
For the mixed-helicity contributions, we first define functions of the final angles
$`𝒞_1^m`$ $`=`$ $`\mathrm{sin}\varphi _q\mathrm{sin}\theta _q\mathrm{cos}\mathrm{\Phi }_2,𝒮_1^m=\mathrm{cos}\varphi _q\mathrm{sin}\theta _q\mathrm{sin}\mathrm{\Phi }_2`$
$`𝒞_2^m`$ $`=`$ $`\mathrm{cos}\theta _q\mathrm{sin}\theta _1+\mathrm{cos}\varphi _q\mathrm{sin}\theta _q\mathrm{cos}\theta _1\mathrm{cos}\mathrm{\Phi }_2`$
$`𝒮_2^m`$ $`=`$ $`\mathrm{sin}\varphi _q\mathrm{sin}\theta _q\mathrm{cos}\theta _1\mathrm{sin}\mathrm{\Phi }_2`$
$`𝒞_3^m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{cos}\mathrm{\Phi }_2{\displaystyle \frac{1}{2}}\mathrm{sin}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{sin}\theta _1\mathrm{cos}2\mathrm{\Phi }_2`$
$`𝒮_3^m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{cos}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}\theta _1\mathrm{sin}\mathrm{\Phi }_2{\displaystyle \frac{1}{2}}\mathrm{cos}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{sin}\theta _1\mathrm{sin}2\mathrm{\Phi }_2`$
$`𝒞_4^m`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{sin}2\theta _1(3\mathrm{cos}^2\theta _q1)`$
$`+{\displaystyle \frac{1}{2}}\mathrm{cos}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}2\theta _1\mathrm{cos}\mathrm{\Phi }_2{\displaystyle \frac{1}{4}}\mathrm{cos}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{sin}2\theta _1\mathrm{cos}2\mathrm{\Phi }_2`$
$`𝒮_4^m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}\varphi _q\mathrm{sin}2\theta _q\mathrm{cos}2\theta _1\mathrm{sin}\mathrm{\Phi }_2{\displaystyle \frac{1}{4}}\mathrm{sin}2\varphi _q\mathrm{sin}^2\theta _q\mathrm{sin}2\theta _1\mathrm{sin}2\mathrm{\Phi }_2`$ (187)
Using these definitions,
$`\mathrm{sin}\mathrm{\Phi }_R\mathrm{sin}\mathrm{\Theta }_B`$ $`=`$ $`𝒞_1^m𝒮_1^m,\mathrm{cos}\mathrm{\Phi }_R\mathrm{sin}\mathrm{\Theta }_B=𝒞_2^m+𝒮_2^m`$
$`\mathrm{sin}\mathrm{\Phi }_R\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B`$ $`=`$ $`𝒞_3^m𝒮_3^m`$
$`\mathrm{cos}\mathrm{\Phi }_R\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B`$ $`=`$ $`𝒞_4^m+𝒮_4^m`$ (188)
and
$`\mathrm{sin}(\mathrm{\Phi }_R+\varphi )\mathrm{sin}\mathrm{\Theta }_B`$ $`=`$ $`𝒫_1^m+𝒬_1^m`$
$`𝒫_1^m`$ $`=`$ $`𝒞_1^m\mathrm{cos}\varphi +𝒮_2^m\mathrm{sin}\varphi ,𝒬_1^m=𝒮_1^m\mathrm{cos}\varphi +𝒞_2^m\mathrm{sin}\varphi `$
$`\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\mathrm{sin}\mathrm{\Theta }_B`$ $`=`$ $`𝒫_2^m+𝒬_2^m`$
$`𝒫_2^m`$ $`=`$ $`𝒞_2^m\mathrm{cos}\varphi +𝒮_1^m\mathrm{sin}\varphi ,𝒬_2^m=𝒮_2^m\mathrm{cos}\varphi 𝒞_1^m\mathrm{sin}\varphi `$
$`\mathrm{sin}(\mathrm{\Phi }_R+\varphi )\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B`$ $`=`$ $`𝒫_3^m+𝒬_3^m`$
$`𝒫_3^m`$ $`=`$ $`𝒞_3^m\mathrm{cos}\varphi +𝒮_4^m\mathrm{sin}\varphi ,𝒬_3^m=𝒮_3^m\mathrm{cos}\varphi +𝒞_4^m\mathrm{sin}\varphi `$
$`\mathrm{cos}(\mathrm{\Phi }_R+\varphi )\mathrm{sin}\mathrm{\Theta }_B\mathrm{cos}\mathrm{\Theta }_B`$ $`=`$ $`𝒫_4^m+𝒬_4^m`$
$`𝒫_4^m`$ $`=`$ $`𝒞_4^m\mathrm{cos}\varphi +𝒮_3^m\mathrm{sin}\varphi ,𝒬_4^m=𝒮_4^m\mathrm{cos}\varphi 𝒞_3^m\mathrm{sin}\varphi `$ (189)
For the “helicity-flip” contributions, $`\mathrm{sin}^2\mathrm{\Theta }_B=2𝒫_0𝒬_0`$ .
## Appendix D Appendix: $`e\overline{e}t\overline{t}`$ Production
In $`e\overline{e}t\overline{t}`$ production, as the center-of-mass energy increases, the helicity-flip amplitudes $`T(\lambda _1,\lambda _2)`$ of (56) will be suppressed relative to the helicity-conserving ones by the factor of $`\sqrt{2}m_t/(\sqrt{s})`$. With respect to more accurate and more precise measurements, this could be a useful variable-dependence. We neglect $`m_e/\sqrt{s}`$ corrections. For the case of $`t\overline{t}`$ production via $`\gamma ^{}`$, the formulas in the text apply with the replacement $`g^2\frac{2}{3}e^2`$ with $`e=\sqrt{4\pi \alpha }`$. For $`Z^{}`$ production, $`\stackrel{~}{T}(+)=v_e+a_e`$ and $`\stackrel{~}{T}(+)=v_ea_e`$ with $`v_e=e(1+4\mathrm{sin}^2\theta _W)/(4\mathrm{sin}\theta _W\mathrm{cos}\theta _W)`$ and $`a_e=e/(4\mathrm{sin}\theta _W\mathrm{cos}\theta _W)`$, and $`T(+)=v_t+a_t(2\stackrel{~}{P}/\sqrt{s})`$, $`T(+)=v_ta_t(2\stackrel{~}{P}/\sqrt{s})`$, $`T(++)=T()=\sqrt{2}v_tm_t/\sqrt{s})`$, with $`v_t=e(38\mathrm{sin}^2\theta _W)/(12\mathrm{sin}\theta _W\mathrm{cos}\theta _W)`$ and $`a_t=e/(4\mathrm{sin}\theta _W\mathrm{cos}\theta _W)`$ with $`\stackrel{~}{P}=`$ magnitude of $`t`$-momentum in $`(t\overline{t})_{cm}`$, and $`1/s1/(sM_{Z}^{}{}_{}{}^{2})`$.
Figure Captions
FIG. 1: In the $`(t\overline{t})_{c.m.}`$ frame, the “final coordinate system” $`(\widehat{x},\widehat{y},\widehat{z})`$ for specification of the beam direction by the spherical angles $`\theta _q`$, $`\varphi _q`$. Note that $`\psi `$ is the smaller angle between the $`W_{1}^{}{}_{}{}^{+}`$ and $`W_{2}^{}{}_{}{}^{}`$ momenta. For the sequential decay $`tW^+b`$ followed by $`W^+l^+\nu `$, the spherical angles $`\theta _a`$, $`\varphi _a`$ specify the $`l^+`$ momentum in the $`W_{1}^{}{}_{}{}^{+}`$ rest frame when there is first a boost from the $`(t\overline{t})_{c.m.}`$ frame to the $`t_1`$ rest frame, and then a second boost from the $`t_1`$ rest frame to the $`W_{1}^{}{}_{}{}^{+}`$ rest frame, see Fig. 5 below. The $`0^o`$ direction for the azimuthal angle $`\varphi _a`$ is defined by the projection of the $`W_{2}^{}{}_{}{}^{}`$ momentum direction.
FIG. 2: Supplement to Fig. 1 to specify the $`CP`$-conjugate sequential decay $`\overline{t}W^{}\overline{b}`$ followed by $`W^{}l^{}\overline{\nu }`$. The spherical angles $`\theta _b`$, $`\varphi _b`$ specify the $`l^{}`$ momentum in the $`W_{2}^{}{}_{}{}^{}`$ rest frame when $`W_{1}^{}{}_{}{}^{+}`$ rest frame when there is first a boost from the $`(t\overline{t})_{c.m.}`$ frame to the $`\overline{t_2}`$ rest frame, and then a second boost from the $`\overline{t_2}`$ rest frame to the $`W_{2}^{}{}_{}{}^{}`$ rest frame. The $`0^o`$ direction for the azimuthal angle $`\varphi _b`$ is defined by the projection of the $`W_{1}^{}{}_{}{}^{+}`$ momentum direction. To better display other angles, the values of the angle $`\psi `$ are different in Figs. 1 and 2.
FIG. 3: Summary illustration showing the three angles $`\theta _1^t`$, $`\theta _2^t`$ and $`\varphi `$ describing the first stage in the sequential-decays of the $`t\overline{t}`$ system in which $`t_1W_{1}^{}{}_{}{}^{+}b`$ and $`\overline{t_2}W_{2}^{}{}_{}{}^{}\overline{b}`$. In (a) the $`b`$ momentum, not shown, is back to back with the $`W_{1}^{}{}_{}{}^{+}`$. In (b) the $`\overline{b}`$ momentum, not shown, is back to back with the $`W_{2}^{}{}_{}{}^{}`$. From (a) a boost along the negative $`z_{1}^{}{}_{}{}^{t}`$ axis transforms the kinematics from the $`t_1`$ rest frame to the $`(t\overline{t})_{c.m.}`$ frame and, if boosted further, to the $`\overline{t_2}`$ rest frame shown in (b). In this figure, $`\varphi _1`$ of Fig. 4 is shown equal to zero for simplicity of illustration.
FIG. 4: The usual helicity angles $`\theta _{1}^{}{}_{}{}^{t}`$ and $`\varphi _1`$ specify the $`W_{1}^{}{}_{}{}^{+}`$ momentum, in the $`t_1`$ rest frame, with $`\overline{t_2}`$ moving in the negative $`z`$ direction. The polar angle $`\theta _{2}^{}{}_{}{}^{t}`$ for the $`W_{2}^{}{}_{}{}^{}`$ is defined analogously in the $`\overline{t_2}`$ rest frame, c.f. Fig. 3. The azimuthal angles $`\varphi _1`$ and $`\varphi _2`$ are Lorentz invariant under boosts along the $`z_{1}^{}{}_{}{}^{t}`$ axis. The sum $`\varphi =\varphi _1+\varphi _2`$ is the angle between the $`t_1`$ and $`\overline{t_2}`$ decay planes.
FIG. 5: The two pairs of spherical angles $`\theta _{1}^{}{}_{}{}^{t},\varphi _1`$ and $`\theta _a`$,$`\stackrel{~}{\varphi _a}`$ specify the two stages in the sequential decay $`tW^+b(l^+\nu )b`$ when the boost to the $`W_{1}^{}{}_{}{}^{+}`$ rest frame is from the $`t_1`$ rest frame. In the $`W_{1}^{}{}_{}{}^{+}`$ rest frame, to reference the $`0^o`$ direction for $`\stackrel{~}{\varphi _a}`$ the axis $`x_a`$ lies in the $`\overline{t_2}`$ half-plane. In this figure, $`\varphi _1`$ of Fig. 4 is shown equal to zero for simplicity of illustration. Similarly, two pairs of spherical angles $`\theta _{2}^{}{}_{}{}^{t},\varphi _2`$ and $`\theta _b`$,$`\stackrel{~}{\varphi _b}`$ specify the two stages in the $`CP`$-conjugate sequential decay $`\overline{t}W^{}\overline{b}`$ followed by $`W^{}l^{}\overline{\nu }`$ when the boost is from the $`\overline{t_2}`$ rest frame.
FIG. 6: The derivation of the general “beam referenced stage-two-spin-correlation” function begins in the “home” or starting coordinate system $`(x_h,y_h,z_h)`$ in the $`(t\overline{t})_{c.m.}`$ frame. $`t_1`$ is moving in the positive $`z_h`$ direction, and $`\theta _1,\varphi _1`$ specify the $`W_{1}^{}{}_{}{}^{+}`$ momentum direction. The beam direction is specified by the spherical angles $`\mathrm{\Theta }_B,\mathrm{\Phi }_B`$. Note that $`\mathrm{\Phi }_R=\mathrm{\Phi }_B\varphi _1`$.
FIG. 7: Supplement to previous figure to show $`\theta _2,\varphi _2`$ which specify the $`W_{2}^{}{}_{}{}^{}`$ momentum direction.
FIG. 8: In the derivation, the “barred” coordinate system $`(\overline{x},\overline{y},\overline{z})`$ in the $`(t\overline{t})_{c.m.}`$ frame has $`W_{1}^{}{}_{}{}^{+}`$ along the positive $`\overline{z}`$ axis with the $`t_1`$ in the negative $`\overline{x}`$ half-plane. A rotation by $`\theta _1`$ has transformed the description from the previous “home system” to the one in this “barred” coordinate system.
FIG. 9: Supplement to previous figure, to show specification of the $`W_{2}^{}{}_{}{}^{}`$ by the spherical angles $`\mathrm{\Theta }_2,\mathrm{\Phi }_2`$. Note that $`\psi +\mathrm{\Theta }_2=\pi `$. A further rotation by minus $`\mathrm{\Phi }_2`$ about the $`\overline{z}`$ axis transforms this “barred system” description” into that in the “final coordinate system” shown in Figs. 1 and 2.
|
warning/0506/math0506270.html
|
ar5iv
|
text
|
# The Algebraic Theory of the Fundamental Germ
## Introduction
Let $``$ be a lamination: a space modeled on a “deck of cards” $`^n\times 𝖳`$, where $`𝖳`$ is a topological space and overlap homeomorphisms take cards to cards continuously in the deck direction $`𝖳`$. One thinks of $``$ as a family of manifolds, the leaves, bound by a transversal topology prescribed locally by $`𝖳`$. Using this picture, many constructions familiar to the theory of manifolds can be extended to laminations via the ansatz:
> Replace manifold object $`A`$ by a family of manifold objects $`\{A_L\}`$ existing on the leaves of $``$ and respecting the transverse topology.
For example, one defines a smooth structure to be a family of smooth structures on the leaves in which the card gluing homeomorphisms occurring in a pair of overlapping decks vary transversally in the smooth topology. Continuing in this way, constructions over $``$, such as tensors, de Rham cohomology groups, etc. may be defined.
Identifying those constructions classically defined over $``$ is not as straightforward, especially if one wishes to follow tradition and define them geometrically. To see why this is true, consider the case of an exceptionally well-behaved lamination: an inverse limit $`\widehat{M}=lim_{}M_\alpha `$ of manifolds by covering maps. Such a system induces a direct limit of de Rham cohomology groups, and there is a canonical map from this limit into the tangential cohomology groups $`H^{}(\widehat{M};)`$ with dense image. In fact, here one may use the system to define – by completion of limits – tangential homology groups $`H_{}(\widehat{M};)`$ as well. If one endeavors to use this point of view to define the groups $`\pi _1`$, $`H_{}()`$, $`H^{}()`$, the result is failure since the systems they induce have trivial limits. The purpose of this paper is to introduce for certain classes of laminations $``$ a construction $`[[\pi ]]_1(,x)`$ called the fundamental germ, a generalization of $`\pi _1`$ which represents an attempt to address this omission in the theory of laminations.
The intuition which guides the construction is that of the lamination as irrational manifold. Recall that for a pointed manifold $`(M,x)`$, the deck group of the universal cover $`(\stackrel{~}{M},\stackrel{~}{x})(M,x)`$ – which may be identified with $`\pi _1(M,x)`$ – reveals through its action how to make identifications within $`(\stackrel{~}{M},\stackrel{~}{x})`$ so as to recover $`(M,x)`$ by quotient. Let us imagine that we have disturbed the process of identifying $`\pi _1`$ orbits, so that instead, points in an orbit merely approximate one another through some auxiliary transversal space T. The result is that $`(\stackrel{~}{M},\stackrel{~}{x})`$ does not produce a quotient manifold but rather coils upon itself, perhaps forming a leaf $`(L,x)`$ of a lamination $``$. The germ of the transversal T about $`x`$ may be interpreted as the failed attempt of $`(L,x)`$ to form an identification topology at $`x`$. The fundamental germ $`[[\pi ]]_1(,x)`$ is then a device which records algebraically the dynamics of $`(L,x)`$ as it approaches $`x`$ through the topology of $`𝖳`$. See Figure 1.
One might define an element of $`[[\pi ]]_1(,x)`$ as a tail equivalence class of a sequence of approaches $`\{x_\alpha \}`$, where $`Lx_\alpha x`$ through $`𝖳`$. In this paper, the laminations under consideration (see §2) have the property that there is a group $`G`$ acting on $`L`$ in such a way that every approach is of the form $`\{g_\alpha x\}`$, for $`g_\alpha G`$. We then define $`[[\pi ]]_1(,x)`$ as the set of tail equivalence classes of sequences of the form $`\{g_\alpha h_\alpha ^1\}`$, where $`g_\alpha x,h_\alpha xx`$ in T. A groupoid structure on $`[[\pi ]]_1(,x)`$ is defined by component-wise multiplication of sequences, and $`\pi _1(L,x)`$ is contained in $`[[\pi ]]_1(,x)`$ as a subgroup. In practice, $`[[\pi ]]_1(,x)`$ has no additional structure; but for many reasonably well-behaved laminations such as inverse limit solenoids, Sullivan solenoids and linear foliations of torii, it is a group. See §§3 – 7 for definitions and examples.
When $`=M`$ is a manifold (a lamination with one leaf), $`[[\pi ]]_1(M,x)`$ is equal to $`{}_{}{}^{}\pi _{1}^{}(M,x)`$, the nonstandard version of $`\pi _1(M,x)`$: the group of tail equivalence classes of all sequences in $`\pi _1(M,x)`$. When $``$ is a lamination contained in a manifold $`M`$, under certain circumstances, §7, there is a map $`[[\pi ]]_1(,x){}_{}{}^{}\pi _{1}^{}(M,x)`$ whose image consists of those classes of sequences in $`\pi _1(M,x)`$ that correspond to the holonomy of $``$. Thus, in expanding $`\pi _1`$ to its nonstandard counterpart, it is possible to detect – algebraically – sublaminations invisible to $`\pi _1`$.
One can profitably think of $`[[\pi ]]_1(,x)`$ as made from sequences of “$`G`$-diophantine approximations”. In the case of an irrational foliation $`_r`$ of the torus $`𝕋^2`$ by lines of slope $`r`$, §4.4, this is literally true: the elements of $`[[\pi ]]_1(_r,x)`$ are the equivalence classes of diophantine approximations of $`r`$. More generally, in $`[[\pi ]]_1`$ one finds an algebraic-topological tool which enables systematic translation of the geometry of laminations into the algebra of (non-linear) diophantine approximation.
One can extend the definition of the fundamental germ to include accumulations of $`L`$ on points of other leaves. Thus if $`\widehat{x}`$ is any point of $``$, we define $`[[\pi ]]_1(,x,\widehat{x})`$ as the set of classes of sequences of the form $`\{g_\alpha h_\alpha ^1\}`$ where $`g_\alpha x,h_\alpha x\widehat{x}`$. We suspect that, together with the topological invariants of the leaves, the fundamental germs $`[[\pi ]]_1(,x,\widehat{x})`$ will play a central role in the topological classification of laminations.
By unwrapping the accumulations of $`L`$ implied by the fundamental germ $`[[\pi ]]_1(,x)`$, one obtains the germ universal cover $`[[\stackrel{~}{}]]`$, §9, which is a kind of nonstandard completion of $`\stackrel{~}{L}`$. If $`[[\pi ]]_1(,x)`$ is a group, then under certain circumstances one may associate lamination coverings $`_\text{C}:=\text{C}\backslash [[\stackrel{~}{}]]`$ of $``$ to every conjugacy class of subgroup $`\text{C}<[[\pi ]]_1(,x)`$, and when C is a normal subgroup, the quotient $`[[\pi ]]_1(,x)/\text{C}`$ may be identified with the automorphism group of $`_\text{C}`$. These considerations give rise to the beginnings of a Galois theory of laminations, §10.
This first paper on the fundamental germ is foundational in nature. One should not expect to find in it hard theorems, but rather the description of a complex and mysterious object which reveals the explicit connection between the geometry of laminations and the algebra of diophantine approximation. Due to its somewhat elaborate construction, we shall confine ourselves here to the following themes:
* Basic definitions: §§1 – 3.
* Examples: §§4 – 7.
* Functoriality: §8.
* Covering space theory: §§9,10.
The focus will be on laminations which arise through group actions: suspensions, quasi-suspensions, double coset foliations and locally-free Lie group actions. The exposition will be characterized by a careful exploration of a number of concrete examples which serve not only to illustrate the definitions in action but also to indicate the richness of the algebra they produce. In a second installment , to appear elsewhere, the construction of $`[[\pi ]]_1`$ will be extended to any lamination whose leaves admit a smooth structure.
Acknowledgements: I have benefited from conversations with B. Le Roin, P. Makienko and especially A. Verjovsky. In addition, the referee made many valuable suggestions which helped to considerably improve the presentation. I would also like to thank the Instituto de Matemáticas (Cuernavaca) of the Universidad Nacional Autónoma de México for providing generous financial support and a pleasant work environment.
## 1. Nonstandard Algebra
We review facts concerning nonstandard algebra, proofs of which may found in the literature. References: , .
Let $`=\{0,1,2,\mathrm{}\}`$, $`𝔘2^{}`$ an ultrafilter all of whose elements have infinite cardinality. Given $`𝒮=\{S_i\}`$ a sequence of sets, write $`S_X=_{jX}S_j.`$ The ultraproduct is the direct limit
$$[S_i]:=\underset{}{lim}S_X,$$
where the system maps are the cartesian projections. If $`S_i=S`$ for all $`i`$, the ultraproduct is called the ultrapower of $`S`$, denoted $`{}_{}{}^{}S`$.
If $`𝒮`$ consists of nested sets, denote by $`{}_{}{}^{}𝒮`$ the set of sequences which converge with respect to $`𝒮`$. For each $`X𝔘`$, define a map $`P_X:{}_{}{}^{}𝒮{}_{}{}^{}𝒮`$ by restriction of indices: $`P_X\left(\{x_\alpha \}\right)=\{x_\alpha \}|_{\alpha X}`$. The ultrascope is the direct limit
$$S_i:=\underset{\stackrel{}{P_X}}{lim}{}_{}{}^{}𝒮.$$
There is a canonical inclusion $`[S_i]S_i`$ , and when $`S_i=S`$ for all $`i`$, the ultrascope coincides with the ultrapower. In general, we have $`S_i={}_{}{}^{}S_{i}^{}{}_{}{}^{}(S_i)`$, where the inclusion is an equality if and only if $`S_i`$ is eventually equal to a fixed set.
If $`𝒮`$ is a (nested) sequence of groups or rings, the induced component-wise operations on sequences descend to operations making the ultraproduct (the ultrascope) a group or ring. This is also true if $`𝒮`$ is a (nested) sequence of fields: we remark here that the maximality property of ultrafilters is required to rule out zero divisors.
If one uses a different ultrafilter $`𝔘^{}`$ and if $`𝒮`$ is a (nested) sequence of groups, rings or fields, then assuming the continuum hypothesis, it is classical that the resulting ultraproduct is isomorphic to that formed from $`𝔘`$. The same can shown for the ultrascope, however we shall not pursue this point here.
The ultrapower $`{}_{}{}^{}`$ is called nonstandard $``$. There is a canonical embedding $`{}_{}{}^{}`$ given by the constant sequences, and we will not distinguish between $``$ and its image in $`{}_{}{}^{}`$. For $`{}_{}{}^{}x,{}_{}{}^{}y{}_{}{}^{}`$, we write $`{}_{}{}^{}x<{}_{}{}^{}y`$ if there exists $`X𝔘`$ and representative sequences $`\{x_i\}`$, $`\{y_i\}`$ such that $`x_i<y_i`$ for all $`iX`$. The non-negative nonstandard reals are defined $`{}_{}{}^{}_{+}^{}=\{{}_{}{}^{}x{}_{}{}^{}|{}_{}{}^{}x0\}`$. The Euclidean norm $`||`$ on $``$ extends to a $`{}_{}{}^{}_{+}^{}`$-valued norm on $`{}_{}{}^{}`$. An element $`{}_{}{}^{}x`$ of $`{}_{}{}^{}`$ is called infinite if for all $`r`$, $`|{}_{}{}^{}x|>r`$, otherwise $`{}_{}{}^{}x`$ is called finite. $`{}_{}{}^{}`$ is a totally-ordered, non-archimedian field.
Here are two topologies that we may give $`{}_{}{}^{}`$:
* The enlargement topology $`{}_{}{}^{}\tau `$, generated by sets of the form $`{}_{}{}^{}A`$, where $`A`$ is open. $`{}_{}{}^{}\tau `$ is $`2^{\mathrm{nd}}`$-countable but not Hausdorff.
* The internal topology $`[\tau ]`$, generated by sets of the form $`[A_i]`$ where $`A_i`$ is open for all $`i`$. $`[\tau ]`$ is Hausdorff but not $`2^{\mathrm{nd}}`$-countable.
We have $`{}_{}{}^{}\tau [\tau ]`$, the inclusion being strict. It is not difficult to see that $`[\tau ]`$ is just the order topology.
###### Proposition 1.
$`({}_{}{}^{},[\tau ])`$ is a real, infinite dimensional topological vector space.
We note however that $`{}_{}{}^{}`$ is not a topological group with respect to $`{}_{}{}^{}\tau `$. Let $`{}_{}{}^{}_{\mathrm{fin}}^{}`$ be the set of finite elements of $`{}_{}{}^{}`$.
###### Proposition 2.
$`{}_{}{}^{}_{\mathrm{fin}}^{}`$ is a topological subring of $`{}_{}{}^{}`$ with respect to both the $`{}_{}{}^{}\tau `$ and $`[\tau ]`$ topologies.
The set of infinitesimals is defined $`{}_{}{}^{}_{ϵ}^{}=\{{}_{}{}^{}ϵ||{}_{}{}^{}ϵ|<M\text{ for all }M_+\}`$, a vector subspace of $`{}_{}{}^{}`$. If $`{}_{}{}^{}x{}_{}{}^{}y{}_{}{}^{}_{ϵ}^{}`$, we write $`{}_{}{}^{}x{}_{}{}^{}y`$ and say that $`{}_{}{}^{}x`$ is infinitesimal to $`{}_{}{}^{}y`$.
###### Proposition 3.
$`{}_{}{}^{}_{\mathrm{fin}}^{}`$ is a local ring with maximal ideal $`{}_{}{}^{}_{ϵ}^{}`$ and
$${}_{}{}^{}_{\mathrm{fin}}^{}/{}_{}{}^{}_{ϵ}^{},$$
a homeomorphism with respect to the quotient $`{}_{}{}^{}\tau `$-topology.
We note that $`{}_{}{}^{}_{ϵ}^{}`$ is clopen in the $`[\tau ]`$-topology; the quotient $`[\tau ]`$-topology on $`{}_{}{}^{}_{\mathrm{fin}}^{}/{}_{}{}^{}_{ϵ}^{}`$ is therefore discrete. $`{}_{}{}^{}_{ϵ}^{}`$ is not an ideal in $`{}_{}{}^{}`$. The vector space
$${}_{}{}^{}:={}_{}{}^{}/{}_{}{}^{}_{ϵ}^{},$$
equipped with the quotient $`{}_{}{}^{}\tau `$-topology, is called the extended reals. By Proposition 3, $`{}_{}{}^{}`$ contains a subfield isomorphic to $``$.
The results above show that neither topology $`{}_{}{}^{}\tau `$ or $`[\tau ]`$ can claim to be preferred. The lack of a canonical topology on $`{}_{}{}^{}`$ is a theme we will encounter again in §9, where we will see that $`{}_{}{}^{}`$ may be viewed as the universal cover of a host of 1-dimensional laminations, each one providing a different topology to $`{}_{}{}^{}`$ (and by pull-back to $`{}_{}{}^{}`$).
Now let $`𝔊`$ be any complete topological group. Some of the properties satisfied by $`{}_{}{}^{}`$ also hold for $`{}_{}{}^{}𝔊`$. If $`\tau `$ denotes the topology of $`𝔊`$, then the topologies $`{}_{}{}^{}\tau `$ and $`[\tau ]`$ are defined exactly as above. $`{}_{}{}^{}𝔊`$ is a topological group in the $`[\tau ]`$ topology, but not in the $`{}_{}{}^{}\tau `$ topology. Denote by $`{}_{}{}^{}𝔊_{ϵ}^{}`$ the classes of sequences converging to the unit element 1. $`{}_{}{}^{}𝔊_{ϵ}^{}`$ is a group since a product of sequences converging to 1 in a topological group is again a sequence converging to 1. Let $`{}_{}{}^{}𝔊_{\mathrm{fin}}^{}`$ be the subset of $`{}_{}{}^{}𝔊`$ all of whose elements are represented by sequences which converge to an element of $`𝔊`$. We have the following analogue of Proposition 3:
###### Proposition 4.
$`{}_{}{}^{}𝔊_{ϵ}^{}`$ is a normal subgroup of $`{}_{}{}^{}𝔊_{\mathrm{fin}}^{}`$ and
$${}_{}{}^{}𝔊_{\mathrm{fin}}^{}/{}_{}{}^{}𝔊_{ϵ}^{}𝔊,$$
a homeomorphism with respect to the quotient $`{}_{}{}^{}\tau `$-topology.
The left coset space
$${}_{}{}^{}𝔊:={}_{}{}^{}𝔊/{}_{}{}^{}𝔊_{ϵ}^{},$$
with the quotient $`{}_{}{}^{}\tau `$-topology, is called the extended $`𝔊`$. It contains $`𝔊`$ as a topological subgroup. If $`𝔊`$ is compact or abelian, then $`{}_{}{}^{}𝔊`$ is a group, though in general it need not be. We will avail ourselves of its natural structure as a $`{}_{}{}^{}𝔊`$-set with respect to the left muliplication action.
## 2. Laminations Associated to Group Actions
The laminations for which we shall define the fundamental germ arise from actions of groups: we review them here as a way of fixing notation. References: , , .
Let us begin by reviewing the definitions and terminology surrounding the concept of a lamination. A deck of cards is a product $`^n\times 𝖳`$, where $`𝖳`$ is a topological space. A card is a subset of the form $`C=O\times \{𝗍\}`$, where $`O^n`$ is open and $`𝗍𝖳`$. A lamination of dimension $`n`$ is a space $``$ equipped with a maximal atlas $`𝒜=\{\varphi _\alpha \}`$ consisting of charts with range in a fixed deck of cards $`^n\times 𝖳`$, such that each transition homeomorphism $`\varphi _{\alpha \beta }=\varphi _\beta \varphi _\alpha ^1`$ satisfies the following conditions:
1. For every card $`C\mathrm{𝖣𝗈𝗆}(\varphi _{\alpha \beta })`$, $`\varphi _{\alpha \beta }(C)`$ is a card.
2. The family of homeomorphisms $`\{\varphi _{\alpha \beta }(,𝗍)\}`$ is continuous in $`𝗍`$.
If T is totally disconnected, we say that $``$ is a solenoid.
An open (closed) transversal in $``$ is a subset of the form $`\varphi _\alpha ^1(\{x\}\times 𝖳^{})`$ where $`𝖳^{}`$ is open (closed) in $`𝖳`$. Note that an open (closed) transversal need not be open (closed) in $``$ i.e. if $``$ is a manifold (viewed as a trivial lamination) then every point is an open transversal. An open (closed) flow box is a subset of the form $`\varphi _\alpha ^1(O\times 𝖳^{})`$, where $`O`$ is open and $`𝖳^{}𝖳`$ is open (closed). A plaque in $``$ is a subset of the form $`\varphi _\alpha ^1(C)`$ for $`C`$ a card in the deck $`^n\times 𝖳`$. A leaf $`L`$ is a maximal continuation of overlapping plaques in $``$. Note that $``$ is the disjoint union of its leaves; we denote by $`L_x`$ the leaf containing the point $`x`$. A lamination is weakly minimal if it has a dense leaf; it is minimal if all of its leaves are dense. A transversal which meets every leaf is called complete. Unless we say otherwise, all transversals in this paper will be complete and open. Two laminations $``$ and $`^{}`$ are said to be homeomorphic if there is a homeomorphism $`f:^{}`$ mapping leaves homeomorphically onto leaves and transversals homeomorphically onto transversals.
### 2.1. Suspensions
Let $`B`$ be a manifold in which $`\pi _1B`$ acts without fixed points. Let $`F`$ be a topological space and $`\rho :\pi _1B\mathrm{𝖧𝗈𝗆𝖾𝗈}(F)`$ a representation. The suspension of $`\rho `$ is the space
$$_\rho =\stackrel{~}{B}\times _\rho F$$
defined by quotienting $`\stackrel{~}{B}\times F`$ by the diagonal action of $`\pi _1B`$, $`\alpha (\stackrel{~}{x},t)=(\alpha \stackrel{~}{x},\rho _\alpha (t))`$. The suspension is a fiber bundle over $`B`$ with model fiber $`F`$. If $`F=𝔊`$ is a topological group and $`\phi :\pi _1B𝔊`$ a homomorphism, then the representation $`\rho :\pi _1B\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝔊)`$ defined $`\rho _\gamma (g)=g\phi (\gamma ^1)`$ gives rise to what we call a $`𝔊`$-suspension, denoted $`_\phi `$, a principle $`𝔊`$-bundle over $`B`$.
The action of $`\pi _1B`$ used to define $`_\rho `$ is properly discontinuous and leaf preserving, hence $`_\rho `$ is a lamination modeled on the deck of cards $`\stackrel{~}{B}\times F`$. If $`K=\mathrm{ker}(\rho )`$ and $`(L,x)_\rho `$ is a pointed leaf, we have $`K\pi _1(L,x)`$. $`_\rho `$ is minimal (weakly-minimal) if and only if every (at least one) $`\rho (\pi _1B)`$ orbit is dense.
The restriction $`p|_L`$ of the projection $`p:_\rho B`$ to a leaf $`L`$ is a covering map. Suppose that $`p_L`$ is a Galois covering (we say that $`L`$ is Galois). The deck group $`D_L`$ of $`p|_L`$ has the property that
$$D_Lx=LF_x,$$
where $`F_x`$ is the fiber of $`p`$ through $`x`$. In particular, if we give $`(LF_x)F_x`$ the subspace topology, we have an inclusion
$$D_L\mathrm{𝖧𝗈𝗆𝖾𝗈}(LF_x).$$
A manifold $`B`$ is a suspension with $`F`$ a point and $`\rho :\pi _1BF`$ trivial. The following subsections discuss examples which are more interesting.
#### 2.1.1. Inverse Limit Solenoids
Let $`𝒞=\{\rho _\alpha :M_\alpha M\}`$ be an inverse system of pointed manifolds and Galois covering maps with initial object $`M`$; denote by
$$\widehat{M}=\widehat{M}_𝒞:=\underset{}{lim}M_\alpha $$
the limit. By definition $`\widehat{M}M_\alpha `$, so elements of $`\widehat{M}`$ are denoted $`\widehat{x}=(x_\alpha )`$, where $`x_\alpha M_\alpha `$. The natural projection onto the base surface is denoted $`p:\widehat{M}M`$. We may identify the universal covers $`\stackrel{~}{M}_\alpha `$ with $`\stackrel{~}{M}`$ and choose the universal covering maps $`\stackrel{~}{M}M_\alpha `$ to be compatible with the system $`𝒞`$. By universality, there exists a canonical map $`i:\stackrel{~}{M}\widehat{M}`$.
Let $`H_\alpha =(\rho _\alpha )_{}(\pi _1M_\alpha )<\pi _1M`$. Associated to $`𝒞`$ is the inverse limit of deck groups
$$\widehat{\pi }_1M:=\underset{}{lim}\pi _1M/H_\alpha ,$$
a Cantor group since the $`\pi _1M/H_\alpha `$ are finite. By universality of inverse limits, the projections $`\pi _1M\pi _1M/H_\alpha `$ yield a canonical homomorphism $`\iota :\pi _1M\widehat{\pi }_1M`$ with dense image. The closures of the images $`\iota (H_\alpha )`$ are clopen, and give a neighborhood basis about 1. Let $`_\iota `$ be the associated $`\widehat{\pi }_1M`$-suspension.
###### Proposition 5.
$`\widehat{M}`$ is homeomorphic to $`_\iota `$. In particular, $`\widehat{M}`$ is a solenoid.
###### Proof.
Let $`\mathrm{{\rm Y}}:\stackrel{~}{M}\times \widehat{\pi }_1M\widehat{M}`$ be the map defined $`(\stackrel{~}{x},\widehat{g})\widehat{g}i(\stackrel{~}{x})`$. $`\mathrm{{\rm Y}}`$ is invariant with respect to the diagonal action of $`\pi _1M`$, and descends to a homeomorphism $`\stackrel{~}{M}\times _\rho \widehat{\pi }_1M\widehat{M}`$. ∎
#### 2.1.2. Linear Foliations of Torii
Let $`V`$ be a $`p`$-dimensional subspace of $`^{p+q}`$. Denote by $`\stackrel{~}{}_V`$ the foliation of $`^{p+q}`$ by cosets $`\text{v}+V`$. The image $`_V`$ of $`\stackrel{~}{}_V`$ in the torus $`𝕋^{p+q}=^{p+q}/^{p+q}`$ gives a foliation of the latter by Euclidean manifolds. $`V`$ may be regarded as the graph of a $`q\times p`$ matrix map
$$𝐑:^p^q$$
whose columns are independent. For $`\text{y}^q`$, denote by $`\overline{\text{y}}`$ its image in $`𝕋^q`$. Let $`\phi _\text{R}:^p𝕋^q`$ be the homomorphism defined
$$\phi _\text{R}(\text{n})=\overline{\text{R}\text{n}},$$
and denote by $`_{\phi _\text{R}}`$ the corresponding $`𝕋^q`$-suspension.
###### Proposition 6.
$`_V`$ is homeomorphic to $`_{\phi _\text{R}}`$.
###### Proof.
Let $`P_0:^{p+q}=^p\times ^q^p\times 𝕋^q`$ be the map defined $`(\text{x},\text{y})(\text{x},\overline{\text{y}}\overline{\text{R}\text{x}})`$. Let $`P`$ be the composition of $`P_0`$ with the projection $`\xi :^p\times 𝕋^q_{\phi _\text{R}}`$. Then $`P`$ is a covering homomorphism with kernel $`^{p+q}`$, hence $`_{\phi _\text{R}}𝕋^{p+q}`$. Since $`V=(\text{x},\text{R}\text{x})`$, we have $`P(V)=\xi (^p\times \overline{\text{0}});`$ thus $`P(V)`$ is a leaf of the suspension. It follows that $`P`$ defines a map $`\stackrel{~}{}_V_{\phi _\text{R}}`$ which descends to the desired homeomorphism. ∎
Let $`\text{r}_i`$ be the $`i`$th column vector of $`𝐑`$. If $`\text{r}_i^q`$ for all $`i`$, the leaves of $`_V`$ are homeomorphic to $`𝕋^p`$ and are not dense. If at least one of the $`\text{r}_i`$ has an irrational coordinate, then the leaves of $`_V`$ are non-compact and dense, homeomorphic to the quotient of $`^p`$ by a discrete subgroup with as many generators as rational $`\text{r}_i`$.
#### 2.1.3. Anosov Foliations
Let $`\mathrm{\Sigma }=^2/\mathrm{\Gamma }`$ be a hyperbolic surface and let $`\rho :\mathrm{\Gamma }\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝕊^1)`$ be defined by extending the action of $`\mathrm{\Gamma }`$ on $`^2`$ to $`^2𝕊^1`$. The suspension
$$_\mathrm{\Gamma }=^2\times _\rho 𝕊^1$$
is called an Anosov foliation. Note that $`_\mathrm{\Gamma }`$ is not an $`𝕊^1`$-suspension. It is classical that the underlying space of $`_\mathrm{\Gamma }`$ is homeomorphic to the unit tangent bundle $`\mathrm{T}_{}^1\mathrm{\Sigma }`$.
### 2.2. Quasisuspensions
Let $`_\rho =\stackrel{~}{B}\times _\rho F`$ be a suspension over a manifold $`B`$. We say that $`_\rho `$ is Galois if every leaf of $`_\rho `$ is Galois. Throughout this section, $`_\rho `$ will be a Galois suspension. We define an action of $`\pi _1B`$ on $`_\rho `$ by
$$x\overline{\gamma }x,$$
where, for $`x`$ contained in the leaf $`L`$, $`\overline{\gamma }`$ is the image of $`\gamma \pi _1B`$ in $`\pi _1B/(p_L)_{}(\pi _1L)`$ $``$ $`D_L`$ = the deck group of $`p|_L`$.
Let $`𝒳_\rho `$ be any closed subset which is invariant with respect to the action of $`\pi _1B`$. Let $`_0:=_\rho 𝒳`$, which is a lamination mapping to $`B`$. If $`_\rho `$ is minimal, then $`𝒳`$ is the preimage of a subset $`XB`$, hence $`_0`$ is a fiber bundle over $`B_0=BX`$. In general, we shall define the fibers of $`_0`$ over $`xB`$ to be the preimages of the map $`_0B`$.
A lamination homeomorphism $`f:_0_0`$ is weakly fiber-preserving if for every fiber $`F_x`$ over $`B`$,
(1)
$$f(F_x)=\underset{i=1}{\overset{n}{}}E_{x_i},$$
where $`E_xF_x`$ denotes a subset of the fiber $`F_x`$. The collection $`\mathrm{𝖧𝗈𝗆𝖾𝗈}_{\omega \mathrm{𝖿𝗂𝖻}}(_0)`$ of weakly fiber-preserving homeomorphisms is clearly a group. Since the fibers are disjoint, each $`E_{x_i}`$ occurring in (1) must be open in $`F_{x_i}`$. In particular, if the fibers are connected, a weakly fiber-preserving homeomorphism is fiber-preserving. Thus, the concept of a weakly fiber-preserving homeomorphism differs from that of a fiber-preserving homeomorphism when the fibers are disconnected e.g. when $`_0`$ is a solenoid.
###### Definition 1.
Let $`_0`$ be as above and suppose $`H<\mathrm{𝖧𝗈𝗆𝖾𝗈}_{\omega \mathrm{𝖿𝗂𝖻}}(_0)`$ is a subgroup acting properly discontinuously on $`_0`$. The quotient
$$𝒬=H\backslash _0$$
is a lamination called a quasisuspension (over $`B`$).
We consider now two examples.
#### 2.2.1. The Sullivan Solenoid
The following important example comes from holomorphic dynamics. Let $`U,V`$ be regions conformal to the unit disc, with $`\overline{U}V`$. Recall that a polynomial-like map is a proper conformal map $`f:UV`$. The conjugacy class of $`f`$ is uniquely determined by a pair $`(p,f)`$, where $`p`$ is a complex polynomial of degree $`d`$ and $`f:𝕊^1𝕊^1`$ is a smooth, expanding map of degree $`d`$ . The space
(2)
$$\widehat{𝕊}=\underset{}{lim}\left(𝕊^1\stackrel{f}{}𝕊^1\stackrel{f}{}𝕊^1\stackrel{f}{}\mathrm{}\right)$$
is an inverse limit solenoid which may be identified with the $`\widehat{}_d`$-suspension $`_\iota =\times _\rho \widehat{}_d`$, where $`\widehat{}_d`$ is the group of $`d`$-adic integers and $`ı:\widehat{}_d`$ is the canonical inclusion. Every leaf of $`\widehat{𝕊}`$ is homeomorphic to $``$. $`f`$ defines a self map of the inverse system in (2), inducing a homeomorphism $`\widehat{f}:\widehat{𝕊}\widehat{𝕊}`$.
Consider the suspension
$$\widehat{𝔻}=^2\times _\rho \widehat{}_d$$
obtained by extending to $`^2\times \widehat{}_d`$ the identification used to define $`_\iota `$ e.g. $`(z,\widehat{n})(\gamma ^mz,\rho _m(\widehat{n}))`$ for $`m`$, where
$$\gamma =\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)$$
is the affine extension of the map $`xx+1`$ to $`^2`$. The base of the suspension $`\widehat{𝔻}`$ is the punctured hyperbolic disc $`𝔻^{}=\gamma \backslash ^2`$, and its ideal boundary may be identified with $`\widehat{𝕊}`$.
The map $`\widehat{f}`$ extends to a weakly fiber-preserving homeomorphism $`\widehat{f}:\widehat{𝔻}\widehat{𝔻}`$ which acts properly discontinuously on $`\widehat{𝔻}`$. The quotient
$$\widehat{𝔻}_f:=\widehat{f}\backslash \widehat{𝔻}$$
is a quasisuspension called the Sullivan solenoid .
#### 2.2.2. The Reeb Foliation
Let $`_+=[0,\mathrm{})`$, consider the trivial suspension $`\times _+`$ over $``$, and denote $`(\times _+)^{}=\times _+\{(0,0)\}`$. Fix $`(\mu ,\lambda )(\times _+)^{}`$ with $`|\mu |,\lambda >1`$, $`\mu \lambda `$. Then multiplication by $`(\mu ,\lambda )`$ in $`(\times _+)^{}`$ is a fiber-preserving lamination homeomorphism giving rise to an action by $``$. The resulting quasisuspension
$$_{\mathrm{𝖱𝖾𝖾𝖻}}=\backslash (\times _+)^{}$$
has underlying space a solid torus, and is called the Reeb foliation.
Let $`P:(\times _+)^{}_{\mathrm{𝖱𝖾𝖾𝖻}}`$ denote the projection map. The leaves of $`_{\mathrm{𝖱𝖾𝖾𝖻}}`$ are of the form:
1. $`L_t=P(\times \{t\})`$, for $`t>0`$.
2. $`L_0=P(^{}\times \{0\})^{}/<\mu >`$.
The fiber tranversals of $`_{\mathrm{𝖱𝖾𝖾𝖻}}`$ are of the form:
1. $`T_z=P(\{z\}\times _+)_+`$, $`z>0`$. Every leaf of $`_{\mathrm{𝖱𝖾𝖾𝖻}}`$ intersects $`T_z`$.
2. $`T_0=P(\{0\}\times (0,\mathrm{}))𝕊^1`$. Every leaf except $`L_0`$ intersects $`T_0`$.
There is an action of $``$ on $`_{\mathrm{𝖱𝖾𝖾𝖻}}`$ induced by the map $`(z,t)(\mu ^nz,t)`$. For $`x_{\mathrm{𝖱𝖾𝖾𝖻}}`$, we write this action $`xnx`$. For every $`t`$ we have $`nL_t`$ = $`L_t`$ and for all $`z`$, $`nT_z`$ = $`T_z`$. Note that this action is the identity on $`L_0`$.
### 2.3. Double Coset Foliations
Let $`𝔊`$ be a Lie group, $``$ a closed Lie subgroup, $`\mathrm{\Gamma }<𝔊`$ a discrete subgroup. The foliation of $`𝔊`$ by right cosets $`g`$ descends to a foliation $`_{,\mathrm{\Gamma }}`$ of $`𝔊/\mathrm{\Gamma }`$, called a double coset foliation.
Let $`\mathrm{\Gamma }`$ be a co-finite volume Fuchsian group. Denote by $`\mathrm{\Sigma }=^2/\mathrm{\Gamma }`$ and by $`\mathrm{T}_{}^1\mathrm{\Sigma }`$ the unit tangent bundle of $`\mathrm{\Sigma }`$. Recall that every $`v\mathrm{T}_{}^1^2`$ determines three oriented, parametrized curves: a geodesic $`\gamma `$ and two horocycles $`𝔥_+`$, $`𝔥_{}`$ tangent to, respectively, $`\gamma (\mathrm{})`$ and $`\gamma (\mathrm{})`$. By parallel translating $`v`$ along these curves, we obtain three flows on $`\mathrm{T}_{}^1^2`$. The three flows are $`\mathrm{\Gamma }`$-invariant, and define flows on $`\mathrm{T}_{}^1\mathrm{\Sigma }`$. The corresponding foliations are denoted $`\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma }`$, $`\mathrm{𝖧𝗈𝗋}_\mathrm{\Gamma }^+`$ and $`\mathrm{𝖧𝗈𝗋}_\mathrm{\Gamma }^{}`$.
Now let $`𝔊=SL(2,)`$ and take $``$ to be one of the 1-parameter subgroups $`H^+=\{A_r^+\}`$, $`H^{}=\{A_r^{}\}`$ and $`G=\{B_r\}`$, where
$$A_r^+=\left(\begin{array}{cc}1& r\\ 0& 1\end{array}\right),A_r^{}=\left(\begin{array}{cc}1& 0\\ r& 1\end{array}\right)\text{and}B_r=\left(\begin{array}{cc}e^{r/2}& 0\\ 0& e^{r/2}\end{array}\right)$$
for $`r`$. Then it is classical that the foliations $`_{G,\mathrm{\Gamma }}`$ and $`_{H^\pm ,\mathrm{\Gamma }}`$ are homeomorphic to $`\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma }`$ and $`\mathrm{𝖧𝗈𝗋}_\mathrm{\Gamma }^\pm `$, respectively. Note also that the Anosov foliation $`_\mathrm{\Gamma }`$ is homeomorphic to the sum $`\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma }\mathrm{𝖧𝗈𝗋}_\mathrm{\Gamma }^+`$.
### 2.4. Locally-Free Lie Group Actions
Let $`𝔅`$ be a Lie group of dimension $`k`$, $`M^n`$ an $`n`$-manifold, $`n>k`$, $`X`$ a subspace of $`M^n`$. A continuous representation $`\theta :𝔅\mathrm{𝖧𝗈𝗆𝖾𝗈}(X)`$ is called locally free if for all $`xX`$, the isotropy subgroup $`I_x<𝔅`$ is discrete. If for any pair $`x,yX`$, their $`𝔅`$-orbits are either disjoint or coincide, then $`X`$ has the structure of a lamination $`_𝔅`$ whose leaves are the $`𝔅`$-orbits.
For example, let $`M^n`$ be a Riemannian manifold. Fix a tangent vector $`v\mathrm{T}_xM`$. Let $`lM^n`$ be the complete geodesic determined by $`v`$, $`X`$ its closure (itself a union of geodesics). Then there is a locally free action of $``$ given by geodesic flow along $`X`$, and $`X`$ is a lamination when $`l`$ is simple. When $`M^n=\mathrm{\Sigma }`$ is a hyperbolic surface and $`l`$ is simple, we obtain a geodesic lamination in $`\mathrm{\Sigma }`$ in the sense of , a solenoid since its transversals are totally-disconnected.
## 3. The Fundamental Germ
Let $``$ be any of the laminations considered in the previous section and let $`L`$ be a fixed leaf. If $`=H\backslash _0`$ is a quasisuspension let $`L_0_0`$ be a leaf lying over $`L`$. The diophantine group $`G_L`$ of $``$ with respect to $`L`$ is
* $`\pi _1B`$ if $``$ is a suspension.
* The group generated by $`\pi _1B`$, $`H_L=\{hH|h(L_0)=L_0\}`$ and $`\pi _1L`$ (viewed as groups acting on $`\stackrel{~}{L}`$) if $``$ is a quasisuspension.
* The group $`\stackrel{~}{}`$ if $``$ is a double coset.
* The group $`\stackrel{~}{𝔅}`$ if $``$ is a locally free Lie group action.
Note that in every case, $`\pi _1L<G_L`$.
Let $`\widehat{x}`$ and $`T`$ a transversal containing $`\widehat{x}`$. Denote by $`\stackrel{~}{T}_L\stackrel{~}{L}`$ the set of points lying over $`TL`$. Then $`T`$ is said to be a diophantine transversal if for every leaf $`L`$ and $`\stackrel{~}{x}\stackrel{~}{T}_L`$, any $`\stackrel{~}{y}\stackrel{~}{T}_L`$ may be written in the form $`\stackrel{~}{y}=g\stackrel{~}{x}`$ for some $`gG_L`$. For $`\stackrel{~}{x}\stackrel{~}{T}_L`$ fixed, we call $`\{g_\alpha \}G_L`$ a $`G_L`$-diophantine approximation of $`\widehat{x}`$ along $`T`$ based at $`\stackrel{~}{x}`$ if $`\{g_\alpha \stackrel{~}{x}\}`$ projects in $`L`$ to a sequence converging to $`\widehat{x}`$ in $`T`$. The image of all such $`G_L`$-diophantine approximations in $`{}_{}{}^{}G_{L}^{}`$ is denoted
$${}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T),$$
and when $`\widehat{x}=x`$ we write $`{}_{}{}^{}𝖣(\stackrel{~}{x},T)`$. If there are no $`G_L`$-diophantine approximations of $`\widehat{x}`$ along $`T`$ based at $`\stackrel{~}{x}`$, we define $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)=0`$. Note that if $`\stackrel{~}{x}^{}=\gamma \stackrel{~}{x}`$ for $`\gamma \pi _1L<G_L`$ then
(3)
$${}_{}{}^{}𝖣(\stackrel{~}{x}^{},\widehat{x},T)\gamma ={}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T).$$
Let $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)^1`$ consist of the set of inverses $`{}_{}{}^{}g_{}^{1}`$ of classes belonging to $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$.
###### Definition 2.
Let $``$, $`L`$, $`x`$, $`\widehat{x}`$ and $`T`$ be as above. The fundamental germ of $``$ based at $`\widehat{x}`$ along $`x`$ and $`T`$ is
$$[[\pi ]]_1(,x,\widehat{x},T)={}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T){}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)^1$$
where $`\stackrel{~}{x}`$ is any point in $`\stackrel{~}{L}`$ lying over $`x`$.
By (3), $`[[\pi ]]_1(,x,\widehat{x},T)`$ does not depend on the choice of $`\stackrel{~}{x}`$ over $`x`$. When $`x=\widehat{x}L`$, we write $`[[\pi ]]_1(,x,T)`$. Observe in this case that $`[[\pi ]]_1(,x,T)`$ contains a subgroup isomorphic to $`{}_{}{}^{}\pi _{1}^{}(L,x)`$.
We now describe a groupoid structure on $`[[\pi ]]_1(,x,\widehat{x},T)`$ . To do this, we define a unit space on which it acts: let $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$ be the image of $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$ in $`{}_{}{}^{}G_{L}^{}`$, for any $`\stackrel{~}{x}`$ over $`x`$. We say that $`{}_{}{}^{}u[[\pi ]]_1(,x,\widehat{x},T)`$ is defined on $`{}_{}{}^{}g{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$ if $`{}_{}{}^{}u{}_{}{}^{}g{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$. Here we are using the left action of $`{}_{}{}^{}G_{L}^{}`$ on $`{}_{}{}^{}G_{L}^{}`$. Having defined the domain and range of elements of $`[[\pi ]]_1(,x,\widehat{x},T)`$, it is easy to see that $`[[\pi ]]_1(,x,\widehat{x},T)`$ is a groupoid, as every element has an inverse by construction. This groupoid structure does not depend on the choice of $`\stackrel{~}{x}`$ over $`x`$.
## 4. The Fundamental Germ of a Suspension
In the case of a suspension $`_\rho =\stackrel{~}{B}\times _\rho F`$, any fiber over the base $`B`$ is a diophantine transversal. Conversely, any diophantine transversal is an open subset of a fiber transversal. It follows that any two diophantine transversals $`T,T^{}`$ through a given point $`\widehat{x}`$ define the same set of $`G_L`$-diophantine approximations. Thus
###### Proposition 7.
If $`T`$ and $`T^{}`$ are diophantine transversals through $`\widehat{x}`$ then
$$[[\pi ]]_1(_\rho ,x,\widehat{x},T)=[[\pi ]]_1(_\rho ,x,\widehat{x},T^{}).$$
Accordingly for suspensions we drop mention of the transversal and write $`[[\pi ]]_1(,x,\widehat{x})`$. We note that since the diophantine group $`G_L=\pi _1B`$ is discrete, $`{}_{}{}^{}G_{L}^{}={}_{}{}^{}G_{L}^{}`$ and the unit space for the groupoid structure is just $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$.
### 4.1. Manifolds
A manifold is a lamination with just one leaf, which can be viewed as the suspension of the trivial representation of its fundamental group. Since a fiber transversal is just a point, we have immediately
###### Proposition 8.
If $`M`$ is a manifold then
$$[[\pi ]]_1(M,x)={}_{}{}^{}𝖣(\stackrel{~}{x})={}_{}{}^{}\pi _{1}^{}(M,x).$$
### 4.2. $`𝔊`$-Suspensions
Let $`\phi :\pi _1B𝔊`$ be a homomorphism, $`_\phi `$ the corresponding $`𝔊`$-suspension. Let $`\{U_i\}`$ be a neighborhood basis about 1 in $`𝔊`$ and define a collection of nested sets $`\{G_i\}`$ by $`G_i=\{\gamma \pi _1B|h(\gamma )U_i\}`$. Note that the ultrascope $`G_i`$ is a subgroup of $`{}_{}{}^{}\pi _{1}^{}B`$. In fact, if $`{}_{}{}^{}\phi :{}_{}{}^{}\pi _{1}^{}B{}_{}{}^{}𝔊`$ is the nonstandard version of $`{}_{}{}^{}\phi `$, then
$$G_i={}_{}{}^{}\phi _{}^{1}({}_{}{}^{}𝔊_{ϵ}^{}).$$
###### Theorem 1.
If $`\phi `$ has dense image, then for any pair $`x,\widehat{x}`$ belonging to a diophantine transversal, $`[[\pi ]]_1(_\phi ,x,\widehat{x})`$ is a group isomorphic to $`G_i`$.
###### Proof.
Let $`{}_{}{}^{}g{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$. Then any other element $`{}_{}{}^{}g_{}^{}{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$ may be written in the form $`{}_{}{}^{}g{}_{}{}^{}h`$ where $`{}_{}{}^{}hG_i`$. It follows immediately that
$$[[\pi ]]_1(_\phi ,x,\widehat{x})={}_{}{}^{}g\left(G_i\right){}_{}{}^{}g_{}^{1}G_i.$$
Because the unit space $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$ is invariant under left-multiplication by its elements, it follows that $`[[\pi ]]_1(_\phi ,x,\widehat{x})`$ acts on it as a group, its groupoid law coinciding with multiplication in $`G_i`$. ∎
For $`𝔊`$-suspensions with $`\phi `$ having dense image, we can thus reduce our notation to $`[[\pi ]]_1(_\phi )`$.
Let $`{}_{}{}^{}\phi :{}_{}{}^{}\pi _{1}^{}B{}_{}{}^{}𝔊`$ be the induced map of nonstandard groups, and denote by $`{}_{}{}^{}\pi _{1}^{}B_{\mathrm{fin}}`$ the subgroup $`{}_{}{}^{}\phi _{}^{1}({}_{}{}^{}𝔊_{\mathrm{fin}}^{})`$. The following theorem can be used to display many familiar topological groups as algebraic quotients of nonstandard versions of discrete groups.
###### Theorem 2.
If $`\phi `$ has dense image, then $`[[\pi ]]_1(_\phi )`$ is a normal subgroup of $`{}_{}{}^{}\pi _{1}^{}B_{\mathrm{fin}}`$ with
$${}_{}{}^{}\pi _{1}^{}B_{\mathrm{fin}}/[[\pi ]]_1(_\phi )𝔊.$$
###### Proof.
Since $`\phi `$ has dense image, the composition of homomorphisms $`{}_{}{}^{}\pi _{1}^{}B_{\mathrm{fin}}{}_{}{}^{}𝔊_{\mathrm{fin}}^{}𝔊`$ – where the first arrow is $`{}_{}{}^{}\phi `$ – is surjective with kernel $`{}_{}{}^{}\phi _{}^{1}({}_{}{}^{}𝔊_{ϵ}^{})=[[\pi ]]_1(_\phi )`$. ∎
### 4.3. Inverse Limit Solenoids
Let $`\widehat{M}`$ be an inverse limit solenoid over the base $`M`$, and let $`\{H_i\}`$ be a sequence of subgroups of $`\pi _1M`$ cofinal in the collection of subgroups in the defining inverse system. By the discussion in §2.1.1, the collection of closures $`\{\widehat{H}_i\}\widehat{\pi }_1M`$ defines a neighborhood basis about $`1`$. Since $`\widehat{M}`$ is a $`\widehat{\pi }_1M`$-suspension in which $`\phi `$ is dense, it follows from Theorem 1 that $`[[\pi ]]_1(\widehat{M},x,\widehat{x})`$ is a group isomorphic to $`H_i`$.
For example, consider a solenoid $`\widehat{𝕊}`$ over $`𝕊^1`$. Here, each $`H_i`$ is an ideal in $``$, hence $`[[\pi ]]_1\widehat{𝕊}`$ is an ideal in the ring $`{}_{}{}^{}`$ = nonstandard $``$. When $`H_i=(d^i)`$ for $`d`$ fixed, we denote the resulting germ $`{}_{}{}^{}_{\widehat{ϵ}}^{}(d)`$ and when $`H_i=(i)`$ we write $`{}_{}{}^{}_{\widehat{ϵ}}^{}`$. Being uncountable, these ideals are not principal, so $`{}_{}{}^{}`$, unlike $``$, in not a PID. By Theorem 2, we have $`{}_{}{}^{}/{}_{}{}^{}_{\widehat{ϵ}}^{}\widehat{}`$ and $`{}_{}{}^{}/{}_{}{}^{}_{\widehat{ϵ}}^{}(d)\widehat{}_d`$.
### 4.4. Linear Foliations of Torii and Classical Diophantine Approximation
Let $`_V`$ be the linear foliation of $`𝕋^{p+q}`$ associated to the subspace $`V^{p+q}`$. As in § 2.1.2, we regard $`V`$ as the graph of the $`q\times p`$ matrix $`𝐑`$. Let $`\phi _\text{R}:^p𝕋^q`$ be the homomorphism used to define $`_V`$. Let $`\{U_i\}`$ be a neighborhood basis in $`𝕋^q`$ about $`\overline{\mathrm{𝟎}}`$. We define a nested set $`\{G_i\}^p`$ by $`𝐧G_i`$ if and only if $`\phi _\text{R}(𝐧)U_i`$. Denote
$${}_{}{}^{}_{𝐑}^{p}:=G_i={}_{}{}^{}\phi _{\text{R}}^{1}({}_{}{}^{}𝕋_{ϵ}^{q}),$$
a subgroup of $`{}_{}{}^{}_{}^{p}`$. If $`p=q=1`$ and $`𝐑=r`$, we write instead $`{}_{}{}^{}_{r}^{}`$.
###### Theorem 3.
If $`\text{R}M_{q,p}()`$, then $`[[\pi ]]_1(_V,x,\widehat{x})={}_{}{}^{}_{𝐑}^{p}`$. Otherwise,
$$[[\pi ]]_1(_V,x,\widehat{x})=\{\begin{array}{cc}0& \text{if }x\widehat{x}\hfill \\ {}_{}{}^{}_{𝐑}^{p}& \text{otherwise}\hfill \end{array}$$
###### Proof.
If $`\text{R}M_{q,p}()`$, then $`\phi _\text{R}`$ has dense image and the result follows by Theorem 1. If not, then all of the leaves are torii so $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})=0`$ unless $`x=\widehat{x}`$, in which case, if $`L`$ is the leaf containing $`x`$, $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})={}_{}{}^{}\pi _{1}^{}L={}_{}{}^{}_{𝐑}^{p}`$. ∎
If R is a vector with at least one irrational entry, then Theorems 2 and 3 give:
###### Corollary 1.
Every finite dimensional torus $`𝕋^q`$ is algebraically isomorphic to a quotient of the nonstardard intergers $`{}_{}{}^{}`$.
###### Theorem 4.
$`{}_{}{}^{}_{\text{R}}^{p}`$ is an ideal in $`{}_{}{}^{}_{}^{p}`$ if and only if $`\text{R}M_{q,p}()`$.
###### Proof.
Suppose that $`\text{R}M_{q,p}()`$ and let $`a_k`$ = the l.c.d. of the entries of $`\text{r}_k`$ = the $`k`$th column of R. Write
$$𝔞=(a_1)\mathrm{}(a_p)$$
where $`(a_k)`$ is the ideal generated by $`a_k`$. Note that $`{}_{}{}^{}𝔞{}_{}{}^{}_{\text{R}}^{p}`$. On the other hand, rationality of the entries of the $`\text{r}_k`$ implies that a sequence $`\{\text{n}_\alpha \}^p`$ defines an element of $`{}_{}{}^{}_{\text{R}}^{p}`$ if and only if there exists $`X𝔘`$ such that $`\phi _\text{R}(\text{n}_\alpha )=\overline{\text{0}}`$ for all $`\alpha X`$. This is equivalent to $`\text{n}_\alpha 𝔞`$ for all $`\alpha X`$. Thus $`{}_{}{}^{}_{\text{R}}^{p}={}_{}{}^{}𝔞`$ which is an ideal in $`{}_{}{}^{}_{}^{p}`$.
Suppose now that $`\text{r}=\text{r}_k^q`$ for some $`k`$, $`1kp`$. Let $`\{\text{n}_\alpha \}`$ represent an element $`{}_{}{}^{}\text{n}{}_{}{}^{}_{\text{R}}^{p}`$, and denote by $`\{n_\alpha \}`$ the sequence of $`k`$-th coordinates of the $`\text{n}_\alpha `$ . Note that $`\overline{n_\alpha \text{r}}\overline{\text{0}}`$ for all $`\alpha `$ since r is not rational. In fact, for any $`\delta >0`$ we may find a sequence of integers $`\{m_\alpha \}`$ such that $`\overline{m_\alpha n_\alpha \text{r}}`$ is not within $`\delta `$ of $`\overline{\text{0}}`$. Let $`\text{m}_\alpha ^p`$ be the vector whose $`k`$th coordinate is $`m_\alpha `$ and whose other coordinates are $`0`$. Then the sequence $`\{\text{m}_\alpha \text{n}_\alpha \}`$ does not converge with respect to $`\{G_i\}`$ i.e. $`{}_{}{}^{}\text{m}{}_{}{}^{}\text{n}{}_{}{}^{}_{\text{R}}^{p}`$, so $`{}_{}{}^{}_{\text{R}}^{p}`$ is not an ideal. ∎
Theorem 4 draws another sharp distinction between $``$ and $`{}_{}{}^{}`$: every subgroup of the former is an ideal, while this is false for the latter.
We spend the rest of this section studying $`{}_{}{}^{}_{𝐑}^{p}`$, in and of itself a complicated and inntriguing object. Let us begin with the following alternate description of $`{}_{}{}^{}_{\text{R}}^{p}`$:
(4)
$${}_{}{}^{}_{\text{R}}^{p}=\left\{{}_{}{}^{}\text{n}{}_{}{}^{}_{}^{p}\right|{}_{}{}^{}\text{n}_{}^{}{}_{}{}^{}_{}^{q}\text{ such that }\text{R}({}_{}{}^{}\text{n}){}_{}{}^{}\text{n}_{}^{}{}_{}{}^{}_{ϵ}^{q}\}.$$
Given $`{}_{}{}^{}\text{n}{}_{}{}^{}_{\text{R}}^{p}`$, the corresponding element $`{}_{}{}^{}\text{n}_{}^{}{}_{}{}^{}_{}^{q}`$ is called the dual of $`{}_{}{}^{}\text{n}`$; it is uniquely determined. From (4), it is clear that the set
$$({}_{}{}^{}_{\text{R}}^{p})^{}:=\left\{{}_{}{}^{}\text{n}_{}^{}\right|{}_{}{}^{}\text{n}_{}^{}\text{ is the dual of }{}_{}{}^{}\text{n}{}_{}{}^{}_{\text{R}}^{p}\}$$
is a subgroup of $`{}_{}{}^{}_{}^{q}`$, called the dual of $`{}_{}{}^{}_{\text{R}}^{p}`$. Note that when $`\text{R}M_{q,p}()`$ has a left-inverse S, we have $`({}_{}{}^{}_{\text{R}}^{p})^{}={}_{}{}^{}_{\text{S}}^{q}`$.
Similarly, the set
$${}_{}{}^{}_{\text{R},ϵ}^{q}=\left\{{}_{}{}^{}ϵ{}_{}{}^{}_{ϵ}^{q}\right|{}_{}{}^{}\text{n}{}_{}{}^{}_{\text{R}}^{p}\text{ such that }\text{R}({}_{}{}^{}\text{n}){}_{}{}^{}\text{n}_{}^{}={}_{}{}^{}ϵ\}$$
is a subgroup of $`{}_{}{}^{}_{ϵ}^{q}`$, called the group of rates of R.
The following proposition is an immediate consequence of (4).
###### Proposition 9.
The maps $`{}_{}{}^{}\text{n}{}_{}{}^{}\text{n}_{}^{}`$ and $`{}_{}{}^{}\text{n}{}_{}{}^{}ϵ`$ define isomorphisms
$${}_{}{}^{}_{\text{R}}^{p}({}_{}{}^{}_{\text{R}}^{p})^{}\text{ and }{}_{}{}^{}_{\text{R}}^{p}{}_{}{}^{}_{\text{R},ϵ}^{q}.$$
###### Note 1 (A.Verjovsky).
Using formulation (4) of $`{}_{}{}^{}_{\text{R}}^{p}`$, it follows that every triple
$$({}_{}{}^{}\text{n},{}_{}{}^{}\text{n}_{}^{},{}_{}{}^{}ϵ)$$
represents a diophantine approximation of R. Thus we may regard $`{}_{}{}^{}_{\text{R}}^{p}`$ as the group of diophantine approximations of R.
For example, when $`p=q=1`$ and $`r`$, $`{}_{}{}^{}n`$ and $`{}_{}{}^{}n_{}^{}`$ are equivalence classes of sequences $`\{x_\alpha \}`$ and $`\{y_\alpha \}`$ $``$, and $`{}_{}{}^{}ϵ`$ an equivalence class of sequence $`\{ϵ_\alpha \}`$, $`ϵ_\alpha 0`$, such that
$$\left|r\frac{y_\alpha }{x_\alpha }\right|=\left|\frac{ϵ_\alpha }{x_\alpha }\right|\mathrm{\hspace{0.33em}\hspace{0.33em}0}.$$
Conversely, every diophantine approximation of $`r`$ defines uniquely a triple $`({}_{}{}^{}n,{}_{}{}^{}n_{}^{},{}_{}{}^{}ϵ)`$.
Recall that two irrational numbers $`r,s`$ are equivalent if there exists
$$A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{𝖲𝖫}(2,)$$
such that $`s=A(r)=(ar+b)/(cr+d)`$.
###### Proposition 10.
If $`r`$ and $`s`$ are equivalent irrational numbers, then $`{}_{}{}^{}_{r}^{}{}_{}{}^{}_{s}^{}`$.
###### Proof.
Given $`{}_{}{}^{}n{}_{}{}^{}_{r}^{}`$, observe that $`(cr+d){}_{}{}^{}nc{}_{}{}^{}n_{}^{}+d{}_{}{}^{}n{}_{}{}^{}.`$ Write $`{}_{}{}^{}m=c{}_{}{}^{}n_{}^{}+d{}_{}{}^{}n`$. Then $`{}_{}{}^{}m{}_{}{}^{}_{s}^{}`$, since
$$s{}_{}{}^{}m(ar+b){}_{}{}^{}na{}_{}{}^{}n_{}^{}+b{}_{}{}^{}n{}_{}{}^{}.$$
The association $`{}_{}{}^{}n{}_{}{}^{}m`$ defines an injective homomorphism $`\psi :{}_{}{}^{}_{r}^{}{}_{}{}^{}_{s}^{}`$, with inverse defined $`\psi ^1({}_{}{}^{}m)(cs+a){}_{}{}^{}m`$. ∎
###### Note 2.
Two irrational numbers $`r,s`$ are called virtually equivalent if there exists $`A\mathrm{𝖲𝖫}(2,)`$ such that $`A(r)=s`$. In this case, there exists a pair of monomorphisms
$$\psi _1:{}_{}{}^{}_{r}^{}{}_{}{}^{}_{s}^{}\text{ and }\psi _2:{}_{}{}^{}_{s}^{}{}_{}{}^{}_{r}^{},$$
defined as in Proposition 10. In other words, $`{}_{}{}^{}_{r}^{}`$ and $`{}_{}{}^{}_{s}^{}`$ are virtually isomorphic. These maps are mutually inverse to each other if and only if $`A\mathrm{𝖲𝖫}(2,)`$.
We are led to make the following conjecture.
###### Conjecture 1.
If $`{}_{}{}^{}_{r}^{}{}_{}{}^{}_{s}^{}`$ for irrational numbers $`r`$, $`s`$, then $`r`$ and $`s`$ are equivalent.
A verified Conjecture 1 would augur a group theoretic approach to diophantine approximation.
### 4.5. Anosov Foliations and Hyperbolic Diophantine Approximation
Let $`\mathrm{\Gamma }`$ be a discrete subgroup of $`\mathrm{𝖯𝖲𝖫}(2,)`$ with no elliptics, $`\mathrm{\Sigma }=\mathrm{\Gamma }\backslash ^2`$ the corresponding Riemann surface. Let $`\rho :\mathrm{\Gamma }\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝕊^1)`$ be the representation of $`\mathrm{\Gamma }`$ on $`𝕊^1^2`$ and denote as in § 2.1.3 the associated Anosov foliation by $`_\mathrm{\Gamma }`$. Fix $`t,\xi 𝕊^1`$, consider a neighborhood basis $`\{U_i(\xi )\}`$ about $`\xi `$, and define the nested set $`\{G_i(t;\xi )\}\mathrm{\Gamma }`$ by
$$G_i(t;\xi )=\left\{A\mathrm{\Gamma }\right|\rho _A(t)U_i(\xi )\}.$$
###### Proposition 11.
Let $`\widehat{x}_\mathrm{\Gamma }`$ be contained in a leaf covered by $`\times \{\xi \}`$ and let $`x`$ be contained in a leaf covered by by $`\times \{t\}`$. Then
$$[[\pi ]]_1(,x,\widehat{x})=\left(G_i(t;\xi )G_i(t;\xi )^1\right).$$
###### Proof.
Immediate from the definition of $`[[\pi ]]_1`$. ∎
Classically , given $`\xi 𝕊^1`$ in the limit set of $`\mathrm{\Gamma }`$ and $`t𝕊^1`$, a $`\mathrm{\Gamma }`$-hyperbolic diophantine approximation of $`\xi `$ based at $`t`$ is a sequence $`\{A_\alpha \}\mathrm{\Gamma }`$ such that $`|\xi A_\alpha (t)|0`$, where $`||`$ is the norm induced by the inclusion $`𝕊^1^2`$. It follows from our definitions that $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$ consists precisely of equivalence classes of $`\mathrm{\Gamma }`$-hyperbolic diophantine approximations.
## 5. The Fundamental Germ of a Quasisuspension
Let $`_\rho `$ be a Galois suspension, $`𝒳_\rho `$ a $`\pi _1B`$ invariant closed set, $`_0=_\rho 𝒳`$. Let $`H<\mathrm{𝖧𝗈𝗆𝖾𝗈}_{\omega \mathrm{𝖿𝗂𝖻}}(_0)`$ be a subgroup acting properly discontinuously and let $`𝒬=H\backslash _0`$ be the resulting quasisuspension. See §2.2. We have the following analogue of Proposition 7:
###### Proposition 12.
If $`T`$ and $`T^{}`$ are diophantine transversals containing $`x`$ and $`\widehat{x}`$ then
$$[[\pi ]]_1(𝒬,x,\widehat{x},T)=[[\pi ]]_1(𝒬,x,\widehat{x},T^{}).$$
###### Proof.
First suppose that the leaf $`L`$ containing $`x`$ has the same topology as any leaf $`L_0`$ lying above it in $`_0`$: in other words, $`H_L=1`$. Then the diophantine group $`G_L`$ is generated only by elements of $`\pi _1B`$ and $`\pi _1L_0`$. We may assume that the transversal $`T`$ lifts to an $`H`$ orbit of disjoint $`\pi _1B`$ transversals $`HT_0`$ in $`_0`$, wherein it follows that
(5)
$$[[\pi ]]_1(𝒬,x,\widehat{x},T)=[[\pi ]]_1(_0,x_0,\widehat{x}_0,T_0)$$
where $`(x_0,\widehat{x}_0,T_0)`$ is a triple that covers $`(x,\widehat{x},T)`$. On the other hand, since the $`\pi _1B`$-invariant set $`𝒳`$ which we removed from $`_\rho `$ to get $`_0`$ is closed, we may assume that $`T_0`$ is a diophantine transversal for $`_\rho `$. It follows then that
$$[[\pi ]]_1(_0,x_0,\widehat{x}_0,T_0)={}_{}{}^{}\pi _{1}^{}L_0[[\pi ]]_1(_\rho ,x,\widehat{x},T_0).$$
The same is true for $`T^{}`$ so by Proposition 7 the result follows.
Now suppose that $`H_L1`$. Then there are $`\pi _1B`$ transversals $`T_0,T_0^{}_0`$ covering $`T,T^{}`$ such that every $`G_L`$-diophantine approximation of $`\widehat{x}`$ along $`T`$ resp. $`T^{}`$ is of the form
$${}_{}{}^{}\gamma {}_{}{}^{}h{}_{}{}^{}g,\text{resp.}{}_{}{}^{}\gamma {}_{}{}^{}h{}_{}{}^{}g_{}^{},$$
where $`{}_{}{}^{}g{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x}_0,T_0)`$, $`{}_{}{}^{}g_{}^{}{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x}_0,T_0^{})`$ (here $`\widehat{x}_0T_0T_0^{}`$ covers $`\widehat{x}`$), $`{}_{}{}^{}h{}_{}{}^{}H_{L}^{}`$ and $`{}_{}{}^{}\gamma {}_{}{}^{}\pi _{1}^{}L_0`$. By the previous paragraph, we have $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x}_0,T_0)={}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x}_0,T_0^{})`$ and the result follows. ∎
Accordingly, we drop mention of $`T`$ and write $`[[\pi ]]_1(𝒬,x,\widehat{x})`$.
###### Note 3.
The proof of Proposition 12 shows that $`{}_{}{}^{}\pi _{1}^{}(L)`$ is a subgroup of $`[[\pi ]]_1(𝒬,x,\widehat{x})`$. In addition, there is a monomorphism $`[[\pi ]]_1(_\rho ,x,\widehat{x})[[\pi ]]_1(𝒬,x,\widehat{x})`$, an isomorphism if $`H_L=\{1\}`$.
### 5.1. Sullivan Solenoids and the Baumslag-Solitar Groups
Consider the Baumslag-Solitar group
$$G_{\mathrm{𝖡𝖲}}=G_{\mathrm{𝖡𝖲}}(d)=f,x:fxf^1=x^d.$$
Define a nested set about 1 by
(6)
$$G_i=\left\{f^mx^{rd^i}\right|m,r\},$$
and denote
$$[[G_{\mathrm{𝖡𝖲}}]]:=(G_iG_i^1).$$
###### Theorem 5.
$`[[G_{\mathrm{𝖡𝖲}}]]`$ is a group.
###### Proof.
Observe by induction that in $`G_{\mathrm{𝖡𝖲}}`$,
(7)
$$x^{d^\alpha }f=fx^{d^{\alpha 1}}$$
for all $`\alpha >0`$. To see that $`[[G_{\mathrm{𝖡𝖲}}]]`$ is a group, it suffices to check that $`G_iG_i^1`$ is a group for all $`i`$. Write a generic element $`gG_iG_i^1`$ in the form $`g=f^lx^{rd^i}f^m`$ for $`l,m,r`$. Then an element $`gh^1`$, $`g,hG_iG_i^1`$ may be written (using (7))
$$gh^1=f^lx^{rd^i}f^mx^{sd^i}f^n=\{\begin{array}{cc}f^lx^{(r+sd^m)d^i}f^{m+n}\hfill & \text{if }m>0\hfill \\ & \\ f^{l+m}x^{(rd^m+s)d^i}f^n\hfill & \text{if }m0\hfill \end{array},$$
where $`l,m,n,r,s`$. It follows that $`gh^1G_iG_i^1`$. ∎
###### Note 4.
The ultrascope $`G_i`$ is not even a groupoid as elements do not have inverses. Indeed, consider the sequence $`\{g_\alpha \}=\left\{f^{m_\alpha }x^{d^\alpha }\right\}`$, where $`m_\alpha >\alpha >0`$, $`\alpha =1,2,\mathrm{}`$. Note that $`\{g_\alpha \}`$ defines an element of $`G_i`$. Using (7), we may write the inverse sequence
$$\{g_\alpha ^1\}=\left\{x^{d^\alpha }f^{m_\alpha }\right\}=\left\{f^\alpha x^1f^{m_\alpha \alpha }\right\}.$$
Since $`m_\alpha >\alpha `$, we cannot use the defining relation of $`G_{\mathrm{𝖡𝖲}}`$ to move the remaining $`f^{m_\alpha \alpha }`$ to the left of the $`x`$-term. It follows that $`\{g_\alpha ^1\}`$ does not define an element of $`G_i`$, so the latter does not have the structure of a groupoid.
###### Theorem 6.
For all $`x,\widehat{x}𝔻_f`$ with $`xL`$,
$$[[\pi ]]_1(\widehat{𝔻}_f,x,\widehat{x})\{\begin{array}{cc}[[G_{\mathrm{𝖡𝖲}}]]\hfill & \text{if }L\text{ is an annulus}\hfill \\ & \\ {}_{}{}^{}_{\widehat{d}}^{}\hfill & \text{if }L\text{ is a disk}\hfill \end{array}$$
In either event, $`[[\pi ]]_1(\widehat{𝔻}_f,x,\widehat{x})`$ is a group.
###### Proof.
First suppose $`L`$ is an annulus. The action of $`\pi _1𝔻^{}`$ on $`\widehat{𝔻}`$ is generated by $`(z,\widehat{n})(z,\widehat{n}+1)`$, where $`(z,\widehat{n})^2\times \widehat{}_d`$. Then if $`\gamma `$ is the generator of $`\pi _1𝔻^{}`$, we have $`\widehat{f}\gamma \widehat{f}^1=\gamma ^d`$. It follows that the diophantine group is isomorphic to $`G_{\mathrm{𝖡𝖲}}`$. The set of diophantine approximations $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})`$ is equal to $`G_i`$, where $`G_i`$ is the nested set (6). The result now follows by definition of $`[[\pi ]]_1`$. If $`L`$ is a disk, then $`[[\pi ]]_1(\widehat{𝔻}_f,x,\widehat{x})=[[\pi ]]_1(\widehat{𝔻},x_0,\widehat{x}_0)`$ where $`(x_0,\widehat{x}_0)`$ covers $`(x,\widehat{x})`$. By the results of § 4.3 we have $`[[\pi ]]_1(\widehat{𝔻},x_0,\widehat{x}_0)={}_{}{}^{}_{\widehat{ϵ}}^{}(d)`$. ∎
The example of $`\widehat{𝔻}_f`$ illustrates the advantage of the “nonabelian Grothendieck group” type construction used in Definition 2: by Note 4, the naive choice “$`[[\pi ]]_1={}_{}{}^{}𝖣`$” would not even have produced a groupoid.
### 5.2. Reeb Foliations
Let $`_{\mathrm{𝖱𝖾𝖾𝖻}}`$ be a Reeb foliation. The diophantine group here is $``$. Recall that $`L_0`$ is the torus leaf.
###### Theorem 7.
For any pair $`x,\widehat{x}_{\mathrm{𝖱𝖾𝖾𝖻}}`$ contained in a diophantine transversal with $`xL`$,
$$[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})\{\begin{array}{cc}{}_{}{}^{}_{}^{2}\hfill & \text{if }x=\widehat{x}L_0=L\hfill \\ & \\ {}_{}{}^{}\hfill & \text{if }\widehat{x}L_0L\hfill \\ & \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
In every case, $`[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})`$ is a group.
###### Proof.
Suppose first that $`x=\widehat{x}L_0`$. Then $`[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})={}_{}{}^{}\pi _{1}^{}L_0={}_{}{}^{}_{}^{2}`$. If $`x\widehat{x}`$ and $`L=L_0`$, there is no diophantine transversal containing the two points hence the fundamental germ is undefined. Now if $`xL`$, $`\widehat{x}L_0L`$ are contained in a diophantine transversal, then a sequence $`\{n_\alpha \}`$ is a diophantine approximation if and only if it is infinite. Thus $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x})={}_{}{}^{}_{\mathrm{}}^{}:={}_{}{}^{}{}_{}{}^{}_{\mathrm{fin}}^{}`$, the infinite nonstandard integers. Then
$$[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})={}_{}{}^{}_{\mathrm{}}^{}{}_{}{}^{}_{\mathrm{}}^{}={}_{}{}^{}.$$
If $`\widehat{x}L^{}L_0`$, there are no accumulations of $`L`$ on $`L^{}`$ so the fundamental germ is 0. ∎
Intuitively, when $`\widehat{x}L_0L`$, $`[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})`$ records the approximation by the dense leaf of the circumferential cycle $`cL_0`$ through $`\widehat{x}`$. On the other hand, $`[[\pi ]]_1(_{\mathrm{𝖱𝖾𝖾𝖻}},x,\widehat{x})`$ does not predict the meridian cycle $`c^{}L_0`$. Instead, $`c^{}`$ is approximated by a sequence of inessential loops in $`L`$ that move off to infinity, and such sequences are not the stuff of $`[[\pi ]]_1`$.
## 6. The Fundamental Germ of a Double Coset Foliation
Let $`𝔊`$ be a Lie group, $`<𝔊`$ a closed subgroup, $`\mathrm{\Gamma }<𝔊`$ a discrete subgroup and $`_{,\mathrm{\Gamma }}`$ the associated double coset foliation. The situation is considerably more subtle due to the fact that the diophantine group is no longer discrete. Thus two choices of diophantine transversal $`T_1`$, $`T_2`$ through $`x,\widehat{x}`$ yield distinct sets of diophantine approximations, in contrast with the case of a (quasi)suspension. Note on the other hand that every transversal is diophantine, since the universal covers of the leaves are homogeneous with respect to the left action of the diophantine group $`\stackrel{~}{}`$. In fact, if $`x_1`$ and $`x_2`$ are contained in the same leaf, then $`\stackrel{~}{a}\stackrel{~}{x}_1=\stackrel{~}{x}_2`$ for some $`\stackrel{~}{a}\stackrel{~}{}`$. This yields a bijection of diophantine sets
$${}_{}{}^{}𝖣(\stackrel{~}{x}_1,\widehat{x},T_1){}_{}{}^{}𝖣(\stackrel{~}{x}_2,\widehat{x},T_2)$$
defined $`{}_{}{}^{}g_{1}^{}{}_{}{}^{}g_{2}^{}`$ if $`{}_{}{}^{}g_{1}^{}={}_{}{}^{}g_{2}^{}\stackrel{~}{a}`$ in $`{}_{}{}^{}\stackrel{~}{}`$. That is, the bijection is given by the equality $`{}_{}{}^{}𝖣(\stackrel{~}{x}_1,\widehat{x},T_1)={}_{}{}^{}𝖣(\stackrel{~}{x}_2,\widehat{x},T_2)\stackrel{~}{a}`$. However, it is not clear that the following prescription for a map of fundamental germs:
(8)
$${}_{}{}^{}u_{1}^{}{}_{}{}^{}u_{2}^{}\text{ iff }{}_{}{}^{}u_{1}^{}={}_{}{}^{}g_{1}^{}{}_{}{}^{}h_{1}^{1},{}_{}{}^{}u_{2}^{}={}_{}{}^{}g_{2}^{}{}_{}{}^{}h_{2}^{1}\text{ and }{}_{}{}^{}g_{1}^{}={}_{}{}^{}g_{2}^{}\stackrel{~}{a},{}_{}{}^{}h_{1}^{}={}_{}{}^{}h_{2}^{}\stackrel{~}{a}$$
is well-defined since there might be, say, another representation $`{}_{}{}^{}u_{1}^{}={}_{}{}^{}g_{1}^{}({}_{}{}^{}h_{1}^{})^1`$ which leads to a different assignment. Even if (8) were well-defined, there is no reason to expect that it should respect the groupoid structure. When $`{}_{}{}^{}\stackrel{~}{}`$ is a group, one can say more:
###### Lemma 1.
If $`{}_{}{}^{}\stackrel{~}{}`$ is a group then $`{}_{}{}^{}u{}_{}{}^{}v={}_{}{}^{}w`$ in $`[[\pi ]]_1(_{,\mathrm{\Gamma }},x,\widehat{x},T)`$ implies $`{}_{}{}^{}u{}_{}{}^{}v={}_{}{}^{}w`$ in $`{}_{}{}^{}\stackrel{~}{}`$.
###### Proof.
This follows immediately since the groupoid structure of the fundamental germ is defined in terms of left multiplication on the unit space $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$. ∎
###### Proposition 13.
If $`{}_{}{}^{}\stackrel{~}{}`$ is a group and $`T_1`$ and $`T_2`$ are diophantine transversals through $`x_1,\widehat{x}`$ and $`x_2,\widehat{x}`$, respectively, where $`x_1,x_2`$ belong to the same leaf $`L`$, then
$$[[\pi ]]_1(_{,\mathrm{\Gamma }},x_1,\widehat{x},T_1)[[\pi ]]_1(_{,\mathrm{\Gamma }},x_2,\widehat{x},T_2).$$
###### Proof.
It is clear now that the bijection (8) is well-defined: in fact, since $`{}_{}{}^{}\stackrel{~}{}`$ is a group, we have $`{}_{}{}^{}u_{1}^{}={}_{}{}^{}u_{2}^{}`$. From this it follows that $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u_{1}^{})=\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u_{2}^{})\stackrel{~}{a}`$, and that the bijection (8) defines a groupoid isomorphism. ∎
We shall assume from this moment on that $`{}_{}{}^{}\stackrel{~}{}`$ is a group. We will then not mention the base point $`x`$ and the transversal $`T`$ and write $`[[\pi ]]_1(_{,\mathrm{\Gamma }},L,\widehat{x})`$ where $`L`$ is the leaf along which diophantine approximations are taking place. If $`\widehat{x}L`$ we write simply $`[[\pi ]]_1(_{,\mathrm{\Gamma }},L)`$.
We now give a “diophantine” description of $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$, similar in spirit to that of $`{}_{}{}^{}_{𝐑}^{p}`$ appearing in (4). Denote by $`p:\stackrel{~}{}`$ the universal cover of $``$. Suppose that $`L`$ is covered by a coset $`g`$ and $`\widehat{g}𝔊`$ is an element covering $`\widehat{x}`$. A subset $`𝒯^{\widehat{g}}𝔊`$ is called a local section at $`\widehat{g}`$ for the quotient map $`𝔊\backslash 𝔊`$ if $`𝒯^{\widehat{g}}`$ maps homeomorphically onto an open subset containing $`\widehat{g}`$. We may assume without loss of generality that the transversal $`T`$ through $`\widehat{x}`$ lifts to a local section $`𝒯^{\widehat{g}}`$ through $`\widehat{g}`$. As our interest is in sequences which converge to $`\widehat{g}`$ in $`𝒯^{\widehat{g}}`$, we may assume also that $`𝒯^{\widehat{g}}=\widehat{g}𝒯`$ for some local section $`𝒯`$ about $`1`$. Let $`{}_{}{}^{}𝒯_{ϵ}^{}{}_{}{}^{}𝔊_{ϵ}^{}`$ denote the set of infinitesimals which are represented by sequences in $`𝒯`$.
Now let $`{}_{}{}^{}\stackrel{~}{h}`$ be a diophantine approximation of $`\widehat{x}`$ based at $`\stackrel{~}{x}`$ along $`T`$, which is characterized by the property that $`\{p({}_{}{}^{}\stackrel{~}{h})g\}`$ lies in $`\widehat{g}{}_{}{}^{}𝒯_{ϵ}^{}{}_{}{}^{}\mathrm{\Gamma }`$. This gives the following diophantine description of $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$:
(9)
$${}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)=\left\{{}_{}{}^{}\stackrel{~}{h}{}_{}{}^{}\stackrel{~}{}\right|{}_{}{}^{}\gamma {}_{}{}^{}\mathrm{\Gamma },{}_{}{}^{}ϵ{}_{}{}^{}𝒯_{ϵ}^{}\text{ such that }\widehat{g}^1p({}_{}{}^{}\stackrel{~}{h})g{}_{}{}^{}\gamma ={}_{}{}^{}ϵ\}.$$
The element $`{}_{}{}^{}\stackrel{~}{h}_{}^{}:={}_{}{}^{}\gamma `$ associated to $`{}_{}{}^{}\stackrel{~}{h}`$ in (9) is called the dual of $`{}_{}{}^{}\stackrel{~}{h}`$. When $`\widehat{g}=g`$, we let $`{}_{}{}^{}\stackrel{~}{}_{g}^{}:={}_{}{}^{}𝖣(\stackrel{~}{x},T)`$ denote the set of diophantine approximations and let $`{}_{}{}^{}\stackrel{~}{}_{g}^{}`$ denote the set of duals. Thus if $`\sigma _g`$ denotes the conjugation map $`ag^1ag`$,
(10)
$${}_{}{}^{}\stackrel{~}{}_{g}^{}=\left\{{}_{}{}^{}\stackrel{~}{h}{}_{}{}^{}\stackrel{~}{}\right|{}_{}{}^{}\gamma {}_{}{}^{}\mathrm{\Gamma },{}_{}{}^{}ϵ{}_{}{}^{}𝒯_{ϵ}^{}\text{ such that }\sigma _g(p({}_{}{}^{}\stackrel{~}{h})){}_{}{}^{}\gamma ={}_{}{}^{}ϵ\}.$$
In general, whether $`g=\widehat{g}`$ or not, it follows that
$$[[\pi ]]_1(_{,\mathrm{\Gamma }},L,\widehat{x})\left\{{}_{}{}^{}\stackrel{~}{u}{}_{}{}^{}\stackrel{~}{}\right|{}_{}{}^{}\gamma ,{}_{}{}^{}\eta {}_{}{}^{}\mathrm{\Gamma },{}_{}{}^{}\omega [[𝒯]]\text{ s.t. }{}_{}{}^{}\gamma \sigma _g(p({}_{}{}^{}\stackrel{~}{u})){}_{}{}^{}\eta ={}_{}{}^{}\omega \},$$
where $`[[𝒯]]={}_{}{}^{}𝒯_{ϵ}^{}{}_{}{}^{}𝒯_{ϵ}^{1}`$. The inclusion is in general strict as the following example shows:
###### Example 1.
Consider the double coset foliation $`\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma }`$, which possesses a noncompact leaf $`L`$ and a pair of cycles $`c_{}`$, $`c_+`$ such that $`L`$ coils about $`c_{}`$ (about $`c_+`$) as one goes to negative (positive) infinity in $`L`$, and has no other accumulations. If $`\widehat{x}`$ belongs to either $`c_{}`$ or $`c_+`$, we have
$$[[\pi ]]_1(\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma },L,\widehat{x}){}_{}{}^{},$$
but if $`\widehat{x}L`$, we have $`[[\pi ]]_1(\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma },L,\widehat{x})=0`$.
###### Note 5.
Since $`p^1(e)\pi _1`$, we have $`{}_{}{}^{}\pi _{1}^{}<[[\pi ]]_1(_{,\mathrm{\Gamma }},L,\widehat{x})`$.
One can understand the description of $`{}_{}{}^{}\stackrel{~}{}_{g}^{}`$ appearing in (10) as a nonlinear version of (4). In fact, if $`𝔊`$ is a linear group of $`p\times p`$ matrices and $`g𝔊`$, then one can think of $`{}_{}{}^{}_{g}^{p}`$ as defined in (4) as the set of linear diophantine approximations of $`g`$ (approximations of $`g`$ by pairs of vectors with respect to linear algebra), whereas $`{}_{}{}^{}\stackrel{~}{}_{g}^{}`$ can be thought of as a set of nonlinear diophantine approximations of $`g`$ (approximations of $`g`$ by pairs of matrices with respect to matrix algebra).
We now consider the horocyclic and geodesic flows on the unit tanget bundle of a riemannian surface, which are, as is widely appreciated, deep mathematical objects. It should come as no suprise that this deepness is reflected in their fundamental germs, which present the most complex and intractable diophantine algebra we have encountered thus far. In the remainder of this section, we will attempt to give the reader a feel for the complexity of these fundamental germs by walking through a sample calculation.
We restrict to the case $`𝔊=SL(2,)`$ and $`\mathrm{\Gamma }=SL(2,)`$. See § 2.3 for the relevant notation. Consider first the case of the (positive) horocyclic flow $`\mathrm{𝖧𝗈𝗋}=\mathrm{𝖧𝗈𝗋}_{\mathrm{𝖲𝖫}(2,)}^+`$, that is, $`=H=H^+`$. If $`D`$ is the subgroup of matrices of the form
$$\left(\begin{array}{cc}e^{s/2}& 0\\ t& e^{s/2}\end{array}\right)$$
$`s,t`$, then $`D`$ defines a local section about $`1`$ so we take $`𝒯=D`$. Finally, since $`H(,+)`$, we shall simplify notation by identifying $`r`$ with the matrix $`A_r`$ and write $`{}_{}{}^{}_{g}^{}={}_{}{}^{}H_{g}^{}`$ for the set of diophantine approximations.
Let us consider the relatively simple choice
$$g=\left(\begin{array}{cc}\sqrt{2}& 1\\ 1& \sqrt{2}\end{array}\right).$$
The right coset of $`g`$ is
$$Hg=\left\{\left(\begin{array}{cc}r+\sqrt{2}& \sqrt{2}r+1\\ 1& \sqrt{2}\end{array}\right)\right|r\}.$$
Since $`Hg`$ does not define a cycle in $`SL(2,)/SL(2,)`$ it must be dense by a theorem of Hedlund , so we can expect from $`g`$ a nontrivial set of diophantine approximations. The conjugate of $`H`$ by $`g`$ is
$$\sigma _g(H)=\left\{\left(\begin{array}{cc}1+\sqrt{2}r& 2r\\ r& 1\sqrt{2}r\end{array}\right)\right|r\}.$$
In order to characterize the elements of $`{}_{}{}^{}_{g}^{}`$, we shall need the following generalization of $`{}_{}{}^{}_{r}^{}`$. Let $`𝕆`$ be the ring of integers of a number field. For $`{}_{}{}^{}r{}_{}{}^{}`$, define
$${}_{}{}^{}𝕆_{{}_{}{}^{}r}^{}=\{{}_{}{}^{}n{}_{}{}^{}𝕆|{}_{}{}^{}n_{}^{}{}_{}{}^{}𝕆\text{ such that }{}_{}{}^{}r{}_{}{}^{}n{}_{}{}^{}n_{}^{}{}_{}{}^{}_{ϵ}^{}\}.$$
Clearly $`{}_{}{}^{}𝕆_{{}_{}{}^{}r}^{}`$ is a subgroup of $`{}_{}{}^{}𝕆`$.
###### Theorem 8.
Let $`𝕆`$ be the ring of integers in $`(\sqrt{2})`$. Then $`{}_{}{}^{}r{}_{}{}^{}_{g}^{}`$ if and only if there exists $`{}_{}{}^{}\gamma =\left(\begin{array}{cc}{}_{}{}^{}a& {}_{}{}^{}b\\ {}_{}{}^{}c& {}_{}{}^{}d\end{array}\right)\mathrm{𝖲𝖫}(2,{}_{}{}^{})`$ for which
* $`\sqrt{2}{}_{}{}^{}a+2{}_{}{}^{}c,\sqrt{2}{}_{}{}^{}b+2{}_{}{}^{}d{}_{}{}^{}𝕆_{{}_{}{}^{}r}^{}`$ and $`(\sqrt{2}{}_{}{}^{}a+2{}_{}{}^{}c)^{}=1{}_{}{}^{}a`$,
$`(\sqrt{2}{}_{}{}^{}b+2{}_{}{}^{}d)^{}={}_{}{}^{}b`$.
* $`{}_{}{}^{}c,{}_{}{}^{}d{}_{}{}^{}_{\sqrt{2}}^{}`$ and $`{}_{}{}^{}c_{}^{}=1{}_{}{}^{}a`$, $`{}_{}{}^{}d_{}^{}=1{}_{}{}^{}b`$.
* $`{}_{}{}^{}b=(\sqrt{2}{}_{}{}^{}b+2{}_{}{}^{}d){}_{}{}^{}r`$ and $`({}_{}{}^{}a+(\sqrt{2}{}_{}{}^{}a+2{}_{}{}^{}c){}_{}{}^{}r)({}_{}{}^{}d({}_{}{}^{}b+\sqrt{2}{}_{}{}^{}d){}_{}{}^{}r)=1`$.
###### Proof.
From (10), $`{}_{}{}^{}r_g`$ if and only if there exists $`{}_{}{}^{}\gamma {}_{}{}^{}\mathrm{\Gamma }`$ and $`{}_{}{}^{}ϵ,{}_{}{}^{}\delta {}_{}{}^{}_{ϵ}^{}`$ with
$$\left(\begin{array}{cc}{}_{}{}^{}a(1+\sqrt{2}{}_{}{}^{}r)+2{}_{}{}^{}c{}_{}{}^{}r& {}_{}{}^{}b(1+\sqrt{2}{}_{}{}^{}r)+2{}_{}{}^{}d{}_{}{}^{}r\\ & \\ {}_{}{}^{}a{}_{}{}^{}r+{}_{}{}^{}c(1\sqrt{2}{}_{}{}^{}r)& {}_{}{}^{}b{}_{}{}^{}r+{}_{}{}^{}d(1\sqrt{2}{}_{}{}^{}r)\end{array}\right)=\left(\begin{array}{cc}1+{}_{}{}^{}ϵ& 0\\ & \\ {}_{}{}^{}\delta & (1+{}_{}{}^{}ϵ)^1\end{array}\right).$$
The first and third items follow immediately. The second item follows upon noting that we may eliminate $`{}_{}{}^{}r`$ by multiplying the second row equations by $`\sqrt{2}`$ and adding them to the first row equations. ∎
Theorem 8 illustrates why it is so difficult to say anything about the algebraic structure of $`{}_{}{}^{}_{g}^{}`$ or $`[[\pi ]]_1(\mathrm{𝖧𝗈𝗋},L)`$. In order to determine whether the sum $`{}_{}{}^{}r+{}_{}{}^{}s`$ defines an element of $`{}_{}{}^{}_{g}^{}`$, we must find a way to “compose” the corresponding duals $`{}_{}{}^{}r_{}^{},{}_{}{}^{}s_{}^{}_g^{}`$ to obtain one for their sum, and it is not even clear what this operation on matrices should be. One could reverse the logic and ask if the product $`{}_{}{}^{}r_{}^{}{}_{}{}^{}s_{}^{}`$ defines an element of $`_g^{}`$, however this seems just as hopeless since the diophantine conditions spelled out in the statement of Theorem 8 are not stable with respect to matrix multiplication.
As for the geodesic flow, we leave it to the reader to formulate the appropriate analogue of Theorem 8 e.g. using the local section $`𝒯`$ for which
$${}_{}{}^{}𝒯_{ϵ}^{}=\left\{\left(\begin{array}{cc}1& {}_{}{}^{}\delta \\ {}_{}{}^{}\delta _{}^{}& 1+{}_{}{}^{}\delta {}_{}{}^{}\delta _{}^{}\end{array}\right)\right|{}_{}{}^{}\delta ,{}_{}{}^{}\delta _{}^{}{}_{}{}^{}_{ϵ}^{}\}.$$
The result would be a set of diophantine conditions at least as daunting as that obtained for the horocyclic flow.
## 7. The Fundamental Germ of a Locally Free Lie Group Action
The discussion here is very similar to that for a double coset, so we will be brief. Let $`𝔅`$ be a Lie group of dimension $`k`$, $`M^n`$ an $`n`$-manifold, $`n>k`$, $`XM^n`$. Let $`\theta :𝔅\mathrm{𝖧𝗈𝗆𝖾𝗈}(X)`$ be a locally-free representation whose orbits either coincide or are disjoint and let $`_𝔅`$ be the associated lamination on $`X`$. Any diophantine transversal through $`x,\widehat{x}`$ may be obtained as the intersection of $`_𝔅`$ with a submanifold $`T`$ of $`M^n`$ of dimension $`nk`$ such that $`x,\widehat{x}T`$ and $`T(\theta (𝔅)x)`$ is discrete in $`\theta (𝔅)x`$. As in the case of a double coset foliation, when $`{}_{}{}^{}𝔅`$ is group,
1. Groupoid multiplication in the fundamental germ corresponds to multiplication in $`{}_{}{}^{}𝔅`$.
2. If $`T_1`$, $`T_2`$ are transversals through $`x_1,\widehat{x}`$ and $`x_2,\widehat{x}`$ where $`x_1,x_2`$ belong to the same leaf $`L`$ then
$$[[\pi ]]_1(_𝔅,x_1,\widehat{x},T_1)[[\pi ]]_1(_𝔅,x_2,\widehat{x},T_2).$$
Accordingly we shorten to $`[[\pi ]]_1(_𝔅,L,\widehat{x})`$.
###### Theorem 9.
Let $`\mathrm{\Sigma }=\mathrm{\Gamma }\backslash ^2`$ be a compact hyperbolic surface, $`𝔩\mathrm{\Sigma }`$ a geodesic lamination, $`\widehat{x}𝔩`$ and $`l𝔩`$ a leaf. Then
$$[[\pi ]]_1(𝔩,l,\widehat{x})=[[\pi ]]_1(\mathrm{𝖦𝖾𝗈𝖽}_\mathrm{\Gamma },L,\widehat{v})$$
where $`L`$ is a leaf covering $`l`$ and $`\widehat{v}`$ is a tangent vector to $`l`$ at $`\widehat{x}`$.
###### Proof.
This follows immediately from the fact that any diophantine approximation of $`\widehat{v}`$ along $`L`$ canonically defines a diophantine approximation of $`\widehat{x}`$ along $`l`$ and vice verca. ∎
## 8. Functoriality
We begin by recalling the notion of morphism in the category of laminations. A lamination map $`F:^{}`$ is a map satisfying the following conditions:
1. For every leaf $`L`$, there exists a leaf $`L^{}^{}`$ with $`F(L)L^{}`$.
2. For all $`x`$, there exist open transversals $`Tx`$, $`T^{}F(x)`$, such that $`F(T)T^{}`$.
The projection $`P:B`$ of a suspension onto its base is a lamination map. On the other hand, let $``$ be a foliation, $`M`$ the underlying manifold. Then the canonical inclusion $`ı:M`$ is a map which maps leaves into the unique leaf $`M`$, yet is not a lamination map since no open transversal of $``$ is mapped into a point, an open transversal of $`M`$.
Let
$$F:(,x,\widehat{x})(^{},x^{},\widehat{x}^{})$$
be a lamination map. We say that $`F`$ is diophantine if there exist diophantine transversals $`Tx,\widehat{x}`$ and $`T^{}x^{},\widehat{x}^{}`$ such that $`F(T)T^{}`$. Note that this condition is always satisfied if either $``$ or $`^{}`$ are laminations defined by double cosets or locally free Lie group actions. Denote by $`L`$ and $`L^{}`$ the leaves containing $`x,x^{}`$ and let $`\stackrel{~}{F}:\stackrel{~}{L}\stackrel{~}{L}^{}`$ be the lift of the restriction $`F|_L`$. Let $`\stackrel{~}{T}\stackrel{~}{L}`$, $`\stackrel{~}{T}^{}\stackrel{~}{L}^{}`$ be the pre-images of $`TL`$, $`T^{}L^{}`$. Then for $`F`$ diophantine there is a well-defined map
$${}_{}{}^{}𝖣F:{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T){}_{}{}^{}𝖣(\stackrel{~}{x}^{},\widehat{x}^{},T^{})$$
of diophantine approximations. If the assigment
$${}_{}{}^{}u={}_{}{}^{}g{}_{}{}^{}h_{}^{1}{}_{}{}^{}𝖣F({}_{}{}^{}g)({}_{}{}^{}𝖣F({}_{}{}^{}h))^1$$
leads to a well-defined map
$$[[F]]:[[\pi ]]_1(,x,\widehat{x})[[\pi ]]_1(^{},x^{},\widehat{x}^{}),$$
we say that $`F`$ is germ.
###### Proposition 14.
Let $`=\stackrel{~}{B}\times _\rho F`$ be a suspension with $`x,\widehat{x}`$ lying over $`x_0B`$. Then the projection $`\xi :(,x,\widehat{x})(B,x_0)`$ is germ, and the induced map $`[[\xi ]]`$ is a groupoid monomorphism.
###### Proof.
It is clear from the definitions that $`{}_{}{}^{}𝖣\xi `$ is the inclusion
$${}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x}){}_{}{}^{}\pi _{1}^{}(B,x).$$
In particular, it follows that $`[[\xi ]]`$ is well-defined. Since the product in $`[[\pi ]]_1(,\widehat{x},L)`$ is induced by multiplication in $`{}_{}{}^{}\pi _{1}^{}(B,x)`$, $`[[\xi ]]`$ is a groupoid homomorphism as well. ∎
Unfortunately, we cannot assert in general that the map $`[[F]]`$ induced by a germ lamination map $`F`$ defines a groupoid homomorphism. We now introduce a class of lamination maps which is sufficiently well-behaved so as to allow us to say more.
Let $``$ be a foliation, $`M`$ the underlying space of $``$, and $`ı:M`$ the inclusion. Although $`ı`$ is not a lamination map, we may nevertheless define a map of diophantine approximations as follows. An element $`{}_{}{}^{}g{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$, represented say by $`\{g_\alpha \}`$, may be regarded as made up from an equivalence class of sequence $`\{\gamma _{g_\alpha }\}`$ where the $`\gamma _{g_\alpha }`$ are homotopy classes of curves lying within $`L`$ whose endpoints converge to $`\widehat{x}`$. One may assume that there is an open disc $`OM`$ about $`\widehat{x}`$ such that the endpoints of these sequences lie entirely in $`O`$. By connecting their endpoints to $`\widehat{x}`$ by a paths contained in $`O`$, we obtain a sequence of homotopy classes of curves $`\{\eta _{g_\alpha }\}\mathrm{\Pi }_1(M,x,\widehat{x})`$ = the set of homotopy classes of paths from $`x`$ and $`\widehat{x}`$, hence a map
$${}_{}{}^{}𝖣ı:{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T){}_{}{}^{}\mathrm{\Pi }_{1}^{}(M,x,\widehat{x}),{}_{}{}^{}g\eta _{}_{}{}^{}g$$
which depends neither on $`O`$ nor on the choice of connecting paths. More generally, given $``$ a lamination and $`ı:X`$ a map into a path-connected space, we may define a map $`{}_{}{}^{}𝖣ı:{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T){}_{}{}^{}\mathrm{\Pi }_{1}^{}(X,ı(x),ı(\widehat{x}))`$. We say that the map $`ı`$ is germ if $`{}_{}{}^{}𝖣ı`$ induces a well-defined map
$$[[ı]]:[[\pi ]]_1(,x,\widehat{x},T){}_{}{}^{}\pi _{1}^{}(X,x),{}_{}{}^{}u={}_{}{}^{}g{}_{}{}^{}h_{}^{1}{}_{}{}^{}𝖣ı({}_{}{}^{}g)({}_{}{}^{}𝖣ı({}_{}{}^{}h))^1.$$
###### Definition 3.
Let $``$ be a lamination arising from a group action, $`X`$ a path connected space. A map $`ı:(,x,\widehat{x})(X,ı(x),ı(\widehat{x}))`$ is called a fidelity if it is germ and $`[[ı]]`$ is a groupoid monomorphism. We say that $``$ is faithful if it has a fidelity.
For example, by Proposition 14 any suspension is faithful, however if the underlying space of a suspension $``$ is a manifold $`M`$, we shall see that it is much more useful to be able to assert that the inclusion $`M`$ is a fidelity.
For the remainder of the section, the base points $`x`$ and $`\widehat{x}`$ will be supressed in order to simplify notation.
###### Proposition 15.
Let $`_V`$ be the foliation of $`𝕋^{p+q}`$ induced by the $`p`$-plane $`V^{p+q}`$. Then the inclusion $`ı:_V𝕋^{p+q}`$ is a fidelity.
###### Proof.
Recall that for some $`q\times p`$ matrix $`𝐑`$, $`[[\pi ]]_1(_V)={}_{}{}^{}_{𝐑}^{p}`$. Then for $`{}_{}{}^{}𝐧{}_{}{}^{}_{𝐑}^{p}`$, the map $`[[ı]]`$ is
$$[[ı]]\left({}_{}{}^{}𝐧\right)=({}_{}{}^{}𝐧,{}_{}{}^{}𝐧_{}^{}){}_{}{}^{}_{}^{p+q}={}_{}{}^{}\pi _{1}^{}𝕋^{p+q},$$
where $`{}_{}{}^{}𝐧_{}^{}`$ is the dual to $`{}_{}{}^{}𝐧`$. $`[[ı]]`$ is then clearly an injective homomorphism. ∎
The problem of the existence of fidelities for laminations arising from group actions is interesting but seems difficult.
###### Conjecture 2.
Every lamination arising form a group action is faithful.
###### Definition 4.
A germ lamination map $`F:^{}`$ is trained if $``$ and $`^{}`$ are faithful, and there exist fidelities $`ı:X`$, $`ı^{}:^{}X^{}`$ and a map $`f:XX^{}`$ such that
(11)
$${}_{}{}^{}f[[ı]]=[[ı^{}]][[F]].$$
The triple $`(ı,ı^{},f)`$ is called a training for $`F`$.
###### Theorem 10.
Let $`F:^{}`$ be a trained lamination map. Then the induced map $`[[F]]`$ is a groupoid homomorphism.
###### Proof.
Let $`(ı,ı^{},f)`$ be a training for $`F`$. Then for all $`{}_{}{}^{}u,{}_{}{}^{}v[[\pi ]]_1()`$ such that $`{}_{}{}^{}u{}_{}{}^{}v`$ is defined we have
$$[[ı^{}]][[F]]\left({}_{}{}^{}u{}_{}{}^{}v\right)=[[ı^{}]]\left([[F]]{}_{}{}^{}u[[F]]{}_{}{}^{}v\right).$$
Since $`[[ı^{}]]`$ is injective, $`[[F]]\left({}_{}{}^{}u{}_{}{}^{}v\right)=[[F]]{}_{}{}^{}u[[F]]{}_{}{}^{}v`$. ∎
###### Corollary 2.
Let $`F:(,x)(^{},x^{})`$ be a map of foliations. Suppose that the inclusions into the underlying manifolds $`ı:M`$, $`ı^{}:^{}M^{}`$ are fidelities. Then $`[[F]]`$ is a groupoid homomorphism.
###### Proof.
Take $`f:MM^{}`$ to be $`F`$, viewed as a map on underlying manifolds. Then $`(ı,ı^{},f)`$ is a training. ∎
###### Corollary 3.
Any map $`F:_V_V^{}`$ of linear foliations of torii induces a homomorphism $`[[F]]`$ of fundamental germs.
## 9. The Germ Universal Cover
We assume throughout this section that
1. $``$ is a weakly-minimal lamination arising from a group action.
2. $`x=\widehat{x}L`$ a fixed dense leaf.
We abreviate the associated fundamental germ to $`[[\pi ]]_1()`$. An ultrafilter $`𝔘`$ is fixed throughout.
Let $`p:\stackrel{~}{L}L`$ be the universal cover. A sequence $`\{\stackrel{~}{x}_\alpha \}\stackrel{~}{L}`$ is called $``$-convergent if it projects to a sequence in $`L`$ converging to some $`\widehat{x}`$. Two $``$-convergent sequences $`\{\stackrel{~}{x}_\alpha \}`$ and $`\{\stackrel{~}{x}_\alpha ^{}\}\stackrel{~}{L}`$ are called $``$-asymptotic if their projections converge to the same point $`\widehat{x}`$ and if for every flowbox $`O`$ in $``$ about $`\widehat{x}`$, there exists $`X𝔘`$ such that $`\stackrel{~}{x}_\alpha `$ and $`\stackrel{~}{x}_\alpha ^{}`$ lie in a common lift of a plaque of $`O`$, for all $`\alpha X`$. The asymptotic class corresponding to $`\{\stackrel{~}{x}_\alpha \}`$ is denoted $`{}_{}{}^{}\stackrel{~}{x}`$; we refer to $`\widehat{x}`$ as the limit of $`{}_{}{}^{}\stackrel{~}{x}`$ and write $`lim{}_{}{}^{}\stackrel{~}{x}=\widehat{x}`$. The set of $`{}_{}{}^{}\stackrel{~}{x}`$ with limit $`\widehat{x}`$ is denoted $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$.
###### Definition 5.
The germ universal cover of $``$ with respect to $`L`$ is
$$[[\stackrel{~}{}]]=\left\{\text{classes }{}_{}{}^{}\stackrel{~}{x}\text{ of }\text{-convergent sequences in }\stackrel{~}{L}\right\}.$$
Note that for any $`\widehat{x}`$, every $`G_L`$-diophantine approximation $`{}_{}{}^{}g`$ of $`\widehat{x}`$ determines an element of $`[[\stackrel{~}{}]]`$, and the sets $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$ and $`{}_{}{}^{}𝖣(\stackrel{~}{x},\widehat{x},T)`$ are in bijective correspondence, for any diophantine transversal $`T`$ through $`x,\widehat{x}`$.
###### Proposition 16.
Let $``$ be compact and suppose that $`L=𝔊`$ is a topological group for which $`{}_{}{}^{}\stackrel{~}{b},{}_{}{}^{}\stackrel{~}{c}{}_{}{}^{}\stackrel{~}{𝔊}`$ are $``$-asymptotic if and only if $`{}_{}{}^{}\stackrel{~}{b}{}_{}{}^{}\stackrel{~}{c}_{}^{1}{}_{}{}^{}\stackrel{~}{𝔊}_{ϵ}^{}`$. Then $`[[\stackrel{~}{}]]={}_{}{}^{}\stackrel{~}{𝔊}`$.
###### Proof.
Suppose that there is some $`{}_{}{}^{}\stackrel{~}{b}{}_{}{}^{}\stackrel{~}{𝔊}`$ represented by a sequence $`\{\stackrel{~}{b}_\alpha \}`$ which is not $``$-convergent. Thus if $`\{b_\alpha \}`$ is the projection of this sequence to $`𝔊`$, then for all $`\widehat{x}`$, $`\widehat{x}`$ has a neighborhood $`U_{\widehat{x}}`$ for which there is no $`X𝔘`$ with $`\{b_\alpha \}|_XU_{\widehat{x}}`$. The $`U_{\widehat{x}}`$ cover $``$ so that there is a subcover $`U_{\widehat{x}_1},\mathrm{},U_{\widehat{x}_n}`$; this implies that there exists a partition $`X_1\mathrm{}X_n`$ of $``$ with $`\{b_\alpha \}|_{X_i}U_{\widehat{x}_i}`$. Since $`𝔘`$ is an ultrafilter, one of the $`X_i`$ belongs to $`𝔘`$, contradiction. Thus every element $`{}_{}{}^{}\stackrel{~}{b}{}_{}{}^{}\stackrel{~}{𝔊}`$ defines an element of $`[[\stackrel{~}{}]]`$. Since the relation of being $``$-asymptotic coincides with differing by an infinitesimal, we are done. ∎
For example, if $`_V`$ is a linear $`n`$-foliation of a torus, $`[[\stackrel{~}{_V}]]={}_{}{}^{}_{}^{n}`$.
Denote by
$${}_{}{}^{}p:[[\stackrel{~}{}]]$$
the natural projection defined $`{}_{}{}^{}\stackrel{~}{x}lim{}_{}{}^{}\stackrel{~}{x}`$. The leaf $`L_{{}_{}{}^{}\stackrel{~}{x}}`$ through $`{}_{}{}^{}\stackrel{~}{x}`$ is defined to be the set of $`{}_{}{}^{}\stackrel{~}{y}`$ such that
1. If $`\widehat{x}=lim{}_{}{}^{}\stackrel{~}{x}`$ and $`\widehat{y}=lim{}_{}{}^{}\stackrel{~}{y}`$ then $`L_{\widehat{x}}=L_{\widehat{y}}`$.
2. There are representative sequences $`\{\stackrel{~}{x}_\alpha \}`$, $`\{\stackrel{~}{y}_\alpha \}`$, and paths $`\stackrel{~}{\eta }_\alpha `$ connecting $`\stackrel{~}{x}_\alpha `$ to $`\stackrel{~}{y}_\alpha `$ so that $`p(\stackrel{~}{\eta }_\alpha )`$ converges to a path connecting $`\widehat{x}`$ to $`\widehat{y}`$.
###### Theorem 11.
$`[[\stackrel{~}{}]]`$ may be given the structure of a lamination whose leaves are nowhere dense and for which $`{}_{}{}^{}p`$ is an open lamination map.
###### Proof.
Denote by $`[[T]][[\stackrel{~}{}]]`$ the pre-image of a transversal $`T`$ and well-order each $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$ for $`\widehat{x}T`$. Note that the cardinalities of the $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$ are the same: that of the continuum, since $`L`$ is dense and $`LT`$ is countable. We define a decomposition
(12)
$$[[T]]=T_\alpha $$
where $`T_\alpha `$ is the section over $`T`$ defined by $`\widehat{x}`$ the $`\alpha `$th element of $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$. By definition of the leaves of $`[[\stackrel{~}{}]]`$, given $`\widehat{x},\widehat{y}T`$,
(13)
$$\left(\underset{{}_{}{}^{}\stackrel{~}{x}\mathrm{𝖫𝗂𝗆}_{\widehat{x}}}{}L_{{}_{}{}^{}\stackrel{~}{x}}\right)\left(\underset{{}_{}{}^{}\stackrel{~}{y}\mathrm{𝖫𝗂𝗆}_{\widehat{y}}}{}L_{{}_{}{}^{}\stackrel{~}{y}}\right)\mathrm{}$$
if and only if $`L_{\widehat{x}}=L_{\widehat{y}}`$. In the latter event the two unions of leaves appearing in (13) are equal, so in particular, given $`{}_{}{}^{}\stackrel{~}{x}\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$, there is a unique $`{}_{}{}^{}\stackrel{~}{y}\mathrm{𝖫𝗂𝗆}_{\widehat{y}}`$ for which $`L_{{}_{}{}^{}\stackrel{~}{x}}=L_{{}_{}{}^{}\stackrel{~}{y}}`$. Since $`TL_{\widehat{x}}`$ is countable, we may thus choose the ordering of each $`\mathrm{𝖫𝗂𝗆}_{\widehat{y}}`$, $`\widehat{y}TL_{\widehat{x}}`$, so that all of the $`\alpha `$th elements lie on distinct leaves. In this way we may asume that the associated section $`T_\alpha `$ intersects any leaf of $`[[\stackrel{~}{}]]`$ no more than once. We topologize each section $`T_\alpha `$ through its identification with $`T`$, and give $`[[\stackrel{~}{}]]`$ the associated product lamination structure. By construction of this topology, $`{}_{}{}^{}p`$ becomes an open lamination map. ∎
The topology constructed in Theorem 11 is called a germ universal cover topology: it is not unique and depends on the choice of decomposition (12). From now on, we assume that $`[[\stackrel{~}{}]]`$ has been equipped with such a topology.
There is a canonical simply connected leaf corresponding to the inclusion $`\stackrel{~}{L}[[\stackrel{~}{}]]`$, however the other leaves need not be simply connected. For example, if $`L`$ is one of the simply connected leaves of the Sullivan solenoid $`\widehat{𝔻}_f`$, then leaves of the associated germ universal cover that correspond to accumulations of $`L`$ on an annular leaf will not be simply connected. Thus $`[[\stackrel{~}{}]]`$ can be thought of as the ordinary universal cover $`\stackrel{~}{L}`$ surrounded by a nonstandard cloud of leaves corresponding to the laminar accumulations of $`L`$; since these leaves are nowhere dense, one might say that on passing to $`[[\stackrel{~}{}]]`$ all of the diophantine approximations within $``$ have been “unwrapped”.
We now posit $`[[\stackrel{~}{}]]`$ as the unit space of an enhanced groupoid structure for $`[[\pi ]]_1()`$. Let $`{}_{}{}^{}u[[\pi ]]_1()`$ and $`{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]`$. We say that $`{}_{}{}^{}u`$ acts on $`{}_{}{}^{}\stackrel{~}{x}`$ if there exist representative sequences such that $`\{u_\alpha \stackrel{~}{x}_\alpha \}`$ defines an $``$-convergent sequence $`{}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x}`$ with
$$lim({}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x})=lim{}_{}{}^{}\stackrel{~}{x}.$$
Defining the domain $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$ and range $`\mathrm{𝖱𝖺𝗇}({}_{}{}^{}u)`$ of $`{}_{}{}^{}u`$ through this notion of action, we see that $`[[\stackrel{~}{}]]`$ yields a new groupoid structure on $`[[\pi ]]_1()`$, called the geometric groupoid structure. It is clear that both $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$ and $`\mathrm{𝖱𝖺𝗇}({}_{}{}^{}u)`$ are sublaminations of $`[[\stackrel{~}{}]]`$, since $`{}_{}{}^{}\stackrel{~}{x}\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$ implies that $`L_{{}_{}{}^{}\stackrel{~}{x}}\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$. Thus we may view $`[[\pi ]]_1()`$ as a groupoid of partially defined bijections of $`[[\stackrel{~}{}]]`$. Note that the unit space for the old groupoid structure, $`{}_{}{}^{}𝖣(\stackrel{~}{x},T)`$, maps into the new unit space $`[[\stackrel{~}{}]]`$ via its bijection with $`\mathrm{𝖫𝗂𝗆}_x`$. There is a canonical inclusion of the old groupoid structure into the geometric groupoid structure, given by extension of domain and range, however in general this map need not be a groupoid homomorphism.
###### Assumption.
For the remainder of the paper, we will assume that $`[[\pi ]]_1()`$ is endowed with the geometric groupoid structure.
###### Definition 6.
We say that $`[[\pi ]]_1()`$ is tame if whenever $`lim{}_{}{}^{}\stackrel{~}{x}=lim{}_{}{}^{}\stackrel{~}{y}`$, there exists $`{}_{}{}^{}u[[\pi ]]_1()`$ such that $`{}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x}={}_{}{}^{}\stackrel{~}{y}`$.
###### Proposition 17.
If $`[[\pi ]]_1()`$ is tame, then the quotient
$$[[\pi ]]_1()\backslash [[\stackrel{~}{}]]$$
is homeomorphic to $``$.
###### Proof.
The equivalence relation enacted by the action of $`[[\pi ]]_1()`$ identifies precisely those points of $`[[\stackrel{~}{}]]`$ which map to the same point $`\widehat{x}`$ by $`{}_{}{}^{}p`$. Since $`{}_{}{}^{}p`$ is open, it follows that quotient topology is that the of $``$. ∎
###### Theorem 12.
If $`[[\pi ]]_1()`$ is tame and a group, then there is a germ universal cover topology on $`[[\stackrel{~}{}]]`$ for which $`[[\pi ]]_1()`$ acts as a group of homeomorphisms.
###### Proof.
Let $`[[T]]`$ be the preimage of a transversal $`T`$ of $``$. As $`[[\pi ]]_1()`$ is a group, $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)=[[\stackrel{~}{}]]`$ for every element $`{}_{}{}^{}u[[\pi ]]_1()`$, and moreover $`{}_{}{}^{}u([[T]])=[[T]]`$. Let $`i:T[[T]]`$ be a section so that for all $`{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]`$, $`i(T)L_{{}_{}{}^{}\stackrel{~}{x}}`$ contains at most one point. Since $`[[\pi ]]_1()`$ acts without fixed points and is tame, we have a decomposition as disjoint union
$$[[T]]=\underset{{}_{}{}^{}u[[\pi ]]_1()}{}{}_{}{}^{}u(i(T)).$$
Now construct as in Theorem 11 a lamination structure on $`[[\stackrel{~}{}]]`$ based on this decomposition. It follows then that each $`{}_{}{}^{}u[[\pi ]]_1()`$ acts homeomorphically on $`[[\stackrel{~}{}]]`$. ∎
###### Proposition 18.
Let $`F:(,L)(^{},L^{})`$ be a lamination map, where $`L`$ and $`L^{}`$ are dense leaves. Then $`F`$ induces a map
$$[[\stackrel{~}{F}]]:[[\stackrel{~}{}]][[\stackrel{~}{^{}}]],$$
continuous with respect to appropriate choices of germ universal cover topologies.
###### Proof.
Denote by $`p^{}:\stackrel{~}{L}^{}L^{}`$ the universal cover. The map $`[[\stackrel{~}{F}]]`$ is defined by representing $`{}_{}{}^{}\stackrel{~}{x}`$ by a sequence $`\{\stackrel{~}{x}_\alpha \}`$ and taking $`[[\stackrel{~}{F}]]({}_{}{}^{}\stackrel{~}{x})`$ to be the asymptotic class of $`\{\stackrel{~}{F}(\stackrel{~}{x}_\alpha )\}`$. Now let $`[[\tau ^{}]]`$ be any germ universal cover topology on $`[[\stackrel{~}{^{}}]]`$, say constructed from a transversal $`T^{}`$. Since $`F`$ is a lamination map, there exists a transversal $`T`$ with $`F(T)T^{}`$. We may thus find a decomposition $`[[T]]=T_\alpha `$ compatible with that of $`[[T^{}]]`$ i.e. so that $`[[\stackrel{~}{F}]](T_\alpha )T_\alpha ^{}`$ for all $`\alpha `$. Let $`[[\tau ]]`$ to be the associated germ universal cover topology. Then $`[[\stackrel{~}{F}]]`$ is continuous with respect to $`[[\tau ]]`$ and $`[[\tau ^{}]]`$. ∎
We now return to the question of functoriality, which we must address in view of our adoption of a new groupoid structure. If we reconsider the notions of fidelities and trainings with regard to the geometric groupoid structure, then the analogue of Theorem 10 – as well as its corollaries – remain true with identical proofs. For the remainder of the paper, the concepts of fidelity and training will be understood in the context of the geometric groupoid structure.
The classical universal cover enjoys the property that the lift $`\stackrel{~}{f}:\stackrel{~}{X}\stackrel{~}{Y}`$ of a map $`f:XY`$ is $`\pi _1X`$-equivariant. We now describe conditions under which the same can be said for a lamination map. A germ lamination map $`F:^{}`$ is said to be geometric if for all $`{}_{}{}^{}u[[\pi ]]_1()`$, $`[[\stackrel{~}{F}]](\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u))\mathrm{𝖣𝗈𝗆}\left([[F]]({}_{}{}^{}u)\right)`$, $`[[F]]:[[\pi ]]_1()[[\pi ]]_1(^{})`$ is a homomorphism and
$$[[\stackrel{~}{F}]]({}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x})=[[F]]({}_{}{}^{}u)[[\stackrel{~}{F}]]({}_{}{}^{}\stackrel{~}{x}).$$
Examples of geometric maps are the projection $`_\rho B`$ of a suspension onto its base and any map of manifolds $`f:MM^{}`$.
We say that a lamination $``$ is geometrically faithful if it has a geometric fidelity: a fidelity $`ı:X`$ which is geometric and for which $`[[\stackrel{~}{ı}]]:[[\stackrel{~}{}]][[\stackrel{~}{X}]]`$ is injective. In addition $`F:^{}`$ is said to be geometrically trained if it possesses a training $`(ı,ı^{},f)`$ where $`ı,ı^{}`$ are geometric fidelities. For example, the fidelity $`ı:_V𝕋^{p+q}`$ of a linear foliation of a torus is geometric, as well as the projection of a suspension onto a compact base.
###### Theorem 13.
Let $`F:^{}`$ be geometrically trained. Then $`F`$ is geometric.
###### Proof.
Let $`(ı,ı^{},f)`$ be a geometric training. Then we have
$$[[\stackrel{~}{ı}]][[\stackrel{~}{F}]]({}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x})=[[\stackrel{~}{ı}]]\left([[F]]({}_{}{}^{}u)[[\stackrel{~}{F}]]({}_{}{}^{}\stackrel{~}{x})\right)$$
which implies the result as $`[[\stackrel{~}{ı}]]`$ is injective. ∎
###### Corollary 4.
Suppose $`F:^{}`$ is a lamination map of foliations such that the inclusions into the underlying manifolds are geometric fidelities. Then $`F`$ is geometric. In particular any lamination map of linear foliations of torii is geometric.
## 10. Covering Space Theory
A surjective lamination map $`P:^{}`$ is called a lamination covering if $`P|_L`$ is a covering map for every leaf $`L`$. A lamination map which is a covering map in the classical sense is a lamination covering but not all lamination coverings occur this way e.g. the projection $`\xi :B`$ of a suspension onto its base. We say that $`P`$ is cover trained if it has a training $`(\iota ,\iota ^{},p)`$ in which $`p:XX^{}`$ is a covering map.
###### Theorem 14.
Let $`P:^{}`$ be a germ lamination covering that is cover trained. Then
1. The induced map of fundamental germs
$$[[P]]:[[\pi ]]_1()[[\pi ]]_1(^{})$$
is a groupoid monomorphism.
2. The induced map of germ universal covers
$$[[\stackrel{~}{P}]]:[[\stackrel{~}{}]][[\stackrel{~}{^{}}]]$$
is an open, injective map with respect to appropriate choices of germ universal cover topologies.
###### Proof.
The first statement follows from the definition of training and the fact that $`{}_{}{}^{}p`$ is injective on $`{}_{}{}^{}\pi _{1}^{}`$. Let $`L`$, $`L^{}`$ be dense leaves in $``$, $`^{}`$ containing $`x`$, $`x^{}`$. Then the lift of the restriction $`P|_L`$, $`\stackrel{~}{P}|_L:\stackrel{~}{L}\stackrel{~}{L}^{}`$, is a homeomorphism. It follows that the induced map $`[[\stackrel{~}{P}]]`$ is injective. $`[[\stackrel{~}{P}]]`$ is automatically open with respect to the germ universal cover topologies constructed as in Proposition 18. ∎
###### Note 6.
Here is an example when the map $`[[\stackrel{~}{P}]]`$ is not surjective. Take $`=`$, $`^{}=𝕊^1`$ and $`P:𝕊^1`$ the universal cover. Then $`[[\stackrel{~}{}]]`$ but $`[[\stackrel{~}{^{}}]]{}_{}{}^{}`$.
Thus when $`P`$ is a cover trained, the image
$$\text{C}=[[P]]\left([[\pi ]]_1()\right)$$
is a subgroupoid of $`[[\pi ]]_1(^{})`$. We shall now construct lamination coverings from subgroups, restricting attention to the case where $`[[\pi ]]_1()`$ is tame and a group. Assume that $`[[\stackrel{~}{}]]`$ has been given a germ universal cover topology $`[[\tau ]]`$ of the type guaranteed by Theorem 12. Consider a subgroup $`\text{C}<[[\pi ]]_1()`$ and denote by $`_\text{C}`$ the quotient $`\text{C}\backslash [[\stackrel{~}{}]]`$. Note that $`_\text{C}`$ decomposes into a disjoint union of leaves. Let $`L_\text{C}`$ be any leaf over the dense leaf $`L`$. Consider the set X of topologies $`[[\tau _𝖢]]`$ on $`[[\stackrel{~}{}]]`$ that satisfy the following conditions.
1. The induced topology $`\tau _𝖢`$ on $`_\text{C}`$ defines a (possibly non Hausdorff) lamination structure for which $`_\text{C}`$ is dense and $`_\text{C}`$ is a lamination map.
2. Let $`T_\text{C}`$ be any transversal of $`_\text{C}`$, and denote by $`[[T_\text{C}]]`$ its preimage with the induced topology. If $`\text{C}{}_{}{}^{}\stackrel{~}{x}`$ is contained in $`[[T_\text{C}]]`$, then $`\text{C}{}_{}{}^{}\stackrel{~}{x}`$ is not open in the topology of $`[[T_\text{C}]]`$.
3. The identity map $`([[\stackrel{~}{}]],[[\tau _𝖢]])([[\stackrel{~}{}]],[[\tau ]])`$ is open and $`[[\pi ]]_1()`$ acts by homeomorphisms on $`([[\stackrel{~}{}]],[[\tau _𝖢]])`$.
X is not empty, as it contains $`[[\tau ]]`$. If we order the elements of X with respect to inclusion, then X is closed under chains and so contains a maximal element which we also denote $`[[\tau _𝖢]]`$, called a covering topology. Denote by $`\tau _𝖢`$ the quotient topology induced by $`[[\tau _𝖢]]`$ on $`_\text{C}`$.
###### Theorem 15.
$`_\text{C}`$ is Hausdorff with respect to $`\tau _𝖢`$ and the map $`_\text{C}`$ is a lamination covering.
###### Proof.
Since $`[[\tau _𝖢]]`$ is maximal, for any $`{}_{}{}^{}\stackrel{~}{x}`$,
$$\text{C}{}_{}{}^{}\stackrel{~}{x}=\underset{\text{C}{}_{}{}^{}\stackrel{~}{x}𝒰[[\tau _𝖢]]}{}𝒰.$$
It follows that $`_\text{C}`$ is Hausdorff. By construction, $`_\text{C}`$ is surjective and a covering when restricted to any leaf. ∎
Two lamination coverings $`P_i:_i`$, $`i=1,2`$, are isomorphic if there exists a geometric homeomorphism $`F:_1_2`$ such that $`P_1=P_2F`$. The group of automorphisms of a lamination cover $`P`$ is denoted $`\mathrm{𝖠𝗎𝗍}(P)`$.
###### Proposition 19.
Let $`{}_{}{}^{}u[[\pi ]]_1(,x)`$ and $`\text{C}^{}={}_{}{}^{}u\text{C}{}_{}{}^{}u_{}^{1}`$. Then there exist covering topologies $`[[\tau _𝖢]]`$ and $`[[\tau _𝖢^{}]]`$ so that $`_\text{C}`$ and $`_\text{C}^{}`$ are isomorphic.
###### Proof.
Choose $`[[\tau _𝖢]]`$ a covering topology for $`_\text{C}`$ and let $`[[\tau _𝖢^{}]]`$ be the image of $`[[\tau _𝖢]]`$ by $`{}_{}{}^{}u`$. Then $`[[\tau _𝖢^{}]]`$ is a covering topology for $`_\text{C}^{}`$. With respect to these choices, the bijection $`{}_{}{}^{}\stackrel{~}{x}{}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x}`$ defines a homeomorphism
$$([[\stackrel{~}{}]],[[\tau _𝖢]])([[\stackrel{~}{}]],[[\tau _𝖢^{}]])$$
which descends to a geometric homeomorphism of covers. ∎
Now suppose $`\text{C}\mathrm{}[[\pi ]]_1(,x)`$ is a normal subgroup, $`[[\tau _𝖢]]`$ the covering topology and $`P_\text{C}:_\text{C}`$ the associated covering.
###### Theorem 16.
$`\mathrm{𝖠𝗎𝗍}(P_\text{C})`$ is isomorphic to the quotient $`[[\pi ]]_1(,x)/\text{C}`$. The quotient of $`_\text{C}`$ by $`[[\pi ]]_1(,x)/\text{C}`$ is $``$.
###### Proof.
Every element of $`_\text{C}`$ is a class $`\text{C}{}_{}{}^{}\stackrel{~}{x}`$, for $`{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]`$. The action of $`[[\pi ]]_1(,x)/\text{C}`$ on such classes is well-defined and yields a subgroup of $`\mathrm{𝖠𝗎𝗍}(P_\text{C})`$. On the other hand, the set $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$ is a $`[[\pi ]]_1()`$-set on which any geometric automorphism acts automorphically. However the automorphism group of $`\mathrm{𝖫𝗂𝗆}_{\widehat{x}}`$ is $`[[\pi ]]_1(,x)/\text{C}`$, so it follows that $`\mathrm{𝖠𝗎𝗍}(P_\text{C})[[\pi ]]_1(,x)/\text{C}`$. It is clear that the quotient of $`_\text{C}`$ by $`[[\pi ]]_1(,x)/\text{C}`$ is $``$. ∎
|
warning/0506/cond-mat0506456.html
|
ar5iv
|
text
|
# Nonaffine Correlations in Random Elastic Media
## I Introduction
In the classical theory of elasticity Landau and Lifshitz (1986); Born and Huang (1954); Love (1944); Chaikin and Lubensky (1995), an elastic material is viewed as a spatially homogeneous medium characterized by a spatially constant elastic-modulus tensor $`K_{ijkl}`$. When such a medium is subjected to uniform stresses at its boundaries, it will undergo a homogeneous deformation with a constant strain. Such homogeneous deformations are called affine. This picture of affine strain is generally valid at length scales large compared to any characteristic inhomogeneities: displacements averaged over a sufficiently large volume are affine (at least in dimensions greater than two). It applies not only to regular periodic crystals, but also to polycrystalline materials like a typical bar of steel. At more microscopic scales, however, individual particles in an elastic medium do not necessarily follow trajectories defined by uniform strain in response to external stress: they undergo nonaffine rather than affine displacements. The only systems that are guaranteed to exhibit affine distortions at the microscopic scale are periodic solids with a single atom per unit cell. Atoms within a multi-atom unit cell of a periodic solid will in general undergo nonaffine distortions Jaric and Mohanty (1988), and atoms in random and amorphous solids will certainly undergo nonaffine distortions. Such distortions can lead to substantial corrections to the Born-Huang Born and Huang (1954) expression for macroscopic elastic moduli.
Research on fragile Durian (1995, 1997); Langer and Liu (1997); Tewari et al. (1999); Evans and Cates (2000), granular Jaeger et al. (1996); Halsey and Mehta (2002), crosslinked polymeric Rubinstein and Panyukov (1997, 2002); Glatting et al. (1997); Everaers (1998); Svaneborg et al. (2004); Sommer and Lay (2002), and biological materials Mackintosh et al. (1995); Head et al. (2003a, b, c); Wilhelm and Frey (2003), particularly in small samples, has sparked a renewed interest in the nature of nonaffine response and its ramifications. Liu and Langer Langer and Liu (1997) introduced various measures of nonaffinity, in particular the mean-square deviation from affinity of individual particles in model foams subjected to shear. Tanguy et al. Tanguy et al. (2002) in their simulation of amorphous systems of Lennard-Jones beads found substantial nonaffine response and a resultant size-dependence to the macroscopic elastic moduli. Lemaitre and Maloney Lemaitre and Maloney (2005) relate nonaffinity to a random force field induced by an initial affine response. Head et al. Head et al. (2003a, b, c) studied models of crosslinked semi-flexible rods in two-dimensions and found two types of behavior depending on the density of rods. In dense systems, the response is close to affine and is dominated by rod compression, whereas in more dilute systems, the response is strongly nonaffine and dominated by rod bending.
The recent work discussed above provides valuable insight into the nature of nonaffine response. It does not, however, provide a general framework in which to describe it. In this paper, we provide a such a framework for describing the long-wavelength properties of nonaffinity, and we verify its validity with numerical calculations of these properties on a number of zero-temperature central-force lattice models specifically designed to demonstrate our ideas. Our hope is that this framework will prove a useful tool for studying more realistic models of amorphous glasses, granular material, and jammed systems, particularly at zero temperature just above the jamming transition Liu and Nagel (1998); O’Hern et al. (2001, 2003). We are currently applying them to jammed systems Vernon et al. (2005) and to networks of semi-flexible polymers Didonna and Levine (2005).
Though nonaffinity concerns the displacement of individual particles at the microscopic scale, we show that general aspects of nonaffine response in random and amorphous systems can be described in terms of a continuum elastic model characterized by a local elastic-modulus tensor $`K_{ijkl}(𝐱)`$ at point $`𝐱`$, consisting of a spatially uniform average part $`K_{ijkl}`$ and a locally fluctuating part $`\delta K_{ijkl}(𝐱)`$, and possibly a local stress tensor $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ with vanishing mean. We show that under stress leading to a macroscopic strain $`\gamma _{ij}`$, the random part of the elastic-modulus tensor, in conjunction with the strain $`\gamma _{ij}`$, acts as a source of nonaffine displacement $`u_i^{}(𝐱)`$ proportional to $`_j\delta K_{ijkl}\gamma _{kl}`$. For small $`\delta K_{ijkl}`$ and $`\gamma _{ij}`$, the Fourier transform of the correlation function $`G_{ij}(𝐱,0)`$ of the displacement $`𝐮^{}(𝐱)`$ can be expressed schematically as $`\gamma ^2\mathrm{\Delta }^K(𝐪)/(q^2K^2)`$ where $`\mathrm{\Delta }^K(𝐪)`$ represents the Fourier transform of relevant components of the correlation function of the random part of the elastic-modulus tensor and $`K`$ represents the average elastic-modulus tensor. At length scales large compared to the correlation length $`\xi `$ of the random elastic modulus, $`\mathrm{\Delta }^K(𝐪)`$ is a constant $`\mathrm{\Delta }^K`$, and the nonaffinity correlation function in $`d`$ dimensions scales as $`(\mathrm{\Delta }^K/K^2)\gamma ^2|𝐱|^{(d2)}`$, which exhibits, in particular, a logarithmic divergence in two dimensions; at length scales smaller than $`\xi `$, $`\mathrm{\Delta }^K(𝐪)q^\varphi `$, where $`\varphi `$ can be viewed as a critical exponent, and the nonaffinity correlation function scales as $`|𝐱|^{\varphi +2d}`$ for $`\varphi <d`$. For simplicity, we focus on zero-temperature systems. Our analytic approach is, however, easily generalized to nonzero temperature in systems with unbreakable bonds. At nonzero temperature, the dominant, long-distance behavior of nonaffinity correlation functions is the same as at zero temperature.
Our numerical studies were carried out on systems composed of sites either on regular periodic lattices or on random lattices constructed by sampling a Lennard-Jones liquid and connecting nearest-neighbor sites with unbreakable central force springs. We allowed the spring constants of the springs, their preferred lengths, or both to vary randomly. The local elastic modulus at a particular site in these models depends on the strength and length of springs connected to that site as well as on the number of springs connected to it. Thus, a periodic lattice with random spring constants and an amorphous lattice with random site coordination numbers both have a random local elastic constant. Their nonaffinity correlation function should, therefore, exhibit similar behavior, as our calculations and simulations verify. It is important to note that macroscopically isotropic systems are always amorphous and, therefore, always have a random elastic-modulus tensor and exhibit nonaffine response. For simplicity, we do not consider systems in which any spring is infinitely rigid (i.e., has an infinite spring constant). With appropriate coarse graining of $`\delta K_{ijkl}(𝐱)`$, however, our primary analytical results are expected to apply to this more general case.
The outline of this paper is as follows. In Sec. II, we derive familiar formulae for the elastic energy of central-force lattices and introduce our continuum model, giving special attention to the nature of random stresses. In Sec. III, we use the continuum model to calculate nonaffine response functions in different dimensions for systems with random elastic moduli with both short- and long-range correlations and with random stress tensors relative to a uniform state, and we calculate the correlation function of local rotations induced by nonaffine distortions. In Sec. IV, we present numerical results for the four model systems we consider: periodic lattices with random elastic constants without (Model A) and with (Model B) random stress, and amorphous lattices with random elastic constants without (Model C) and with (Model D) random stress. Four appendices present calculational details: Appendix A derives the independent components of the $`8`$th rank modulus correlator in an isotropic medium, App. B calculates the general form of the nonaffinity correlation function as a function of wavevector, App. C calculates the asymptotic forms as a function of separation $`𝐱`$ of the nonaffinity correlation function, and App. D calculates the correlation function of local vorticity.
## II Models and Definitions
### II.1 Notation and Model Energy
We consider model elastic networks in which particles occupy sites on periodic or random lattices in their force-free equilibrium state. Thus, particle $`\mathrm{}`$ is at lattice position $`𝐑_\mathrm{}0`$ in equilibrium. When the lattices are distorted, particle $`\mathrm{}`$ undergoes a displacement $`𝐮_{\mathrm{}}`$ to a new position
$$𝐑_{\mathrm{}}=𝐑_\mathrm{}0+𝐮_{\mathrm{}}.$$
(1)
We will refer to the equilibrium lattice, with lattice positions $`𝐑_\mathrm{}0`$, as the reference lattice or reference space, and the space into which the lattice is distorted via the displacements $`𝐮_{\mathrm{}}`$ as the target space. Pairs of particles $`\mathrm{}`$ and $`\mathrm{}^{}`$ are connected by unbreakable central-force springs on the bond $`b<\mathrm{}^{},\mathrm{}>`$. The coordination number of each particle (or site) is equal to the number of particles (or sites) to which it is connected by bonds. The potential energy, $`V_b(R_b)`$, of the spring on bond $`b`$ depends only on the magnitude,
$$R_b=|𝐑_{\mathrm{}^{}}𝐑_{\mathrm{}}|,$$
(2)
of the vector connecting particles $`\mathrm{}`$ and $`\mathrm{}^{}`$. The total potential energy is thus
$$U_T=\underset{b}{}V_b(R_b)\frac{1}{2}\underset{\mathrm{},\mathrm{}^{}}{}V_{<\mathrm{}^{},\mathrm{}>}(|𝐑_{\mathrm{}^{}}𝐑_{\mathrm{}}|).$$
(3)
We will consider anharmonic potentials
$$V_b=\frac{1}{2}k_b(R_bR_{bR})^2+\frac{1}{4}g_b(R_bR_{bR})^4,$$
(4)
with both harmonic and quartic components, where $`R_{bR}`$ is the rest length of bond $`b`$. We assume that both $`k_b`$ and $`g_b`$ are finite. The harmonic limit is obtained when the quartic coefficient $`g_b`$ vanishes, in which case, $`k_b`$ is the harmonic spring constant.
We will only study systems in which there is an equilibrium reference state with particle positions $`\{𝐑_\mathrm{}0\}`$ in which the force on each site is zero. The length $`R_{b0}|𝐑_\mathrm{}^{}0𝐑_\mathrm{}0|`$ of each bond $`b`$ in this configuration does not have to coincide with its rest length $`R_{bR}`$. As we shall see in more detail shortly, it is possible to have the total force on every site be zero but still have nonzero forces on each bond.
The potential energy of the lattice can be expanded in terms of the discrete lattice nonlinear strain Born and Huang (1954),
$$v_b=\frac{1}{2}(R_b^2R_{b0}^2)=𝐑_{b0}\mathrm{\Delta }𝐮_b+\frac{1}{2}(\mathrm{\Delta }𝐮_b\mathrm{\Delta }𝐮_b)$$
(5)
relative to the reference state, where $`\mathrm{\Delta }𝐮_b=𝐮_{\mathrm{}^{}}𝐮_{\mathrm{}}`$. The discrete strain variable, $`v_b`$, is by construction invariant with respect to rigid rotations of the sample, i.e., it is invariant under $`R_\mathrm{}iU_{ij}R_\mathrm{}j`$, where $`U_{ij}`$ is any $`\mathrm{}`$-independent rotation matrix. To second order in $`v_b`$ in an expansion about a reference lattice with lattice sites $`𝐑_\mathrm{}0`$, the potential energy is Born and Huang (1954)
$$\mathrm{\Delta }U_T=\underset{b}{}R_{b0}^1\stackrel{~}{F}(b)v_b+\frac{1}{2}\underset{b}{}R_{b0}^2k(b)v_b^2,$$
(6)
where $`\stackrel{~}{F}(b)=|\stackrel{~}{𝐅}(b)|`$ is the magnitude of the force,
$$\stackrel{~}{𝐅}(b)=V_b^{}(R_{b0})𝐑_{b0}/R_{b0},$$
(7)
acting on bond $`b`$ and
$$k(b)=V_{b0}^{\prime \prime }(R_{b0})R_{b0}^1V_{b0}^{}(R_{b0})$$
(8)
is the effective spring constant of bond $`b`$, which reduces to $`k_b`$ when $`R_{b0}=R_{bR}`$. $`k(b)`$ is never infinite because we we assume $`k_b`$ and $`g_b`$ are finite. The equilibrium bond-length $`R_{b0}`$ for each bond is determined by the condition that the total force at each site $`\mathrm{}`$ vanish at $`𝐮_{\mathrm{}}=0`$:
$$F_i(\mathrm{})=\frac{\mathrm{\Delta }U_T}{u_\mathrm{}i}|_{𝐮_{\mathrm{}}=0}=\underset{\mathrm{}^{}}{}\stackrel{~}{F}_i(<\mathrm{}^{},\mathrm{}>).$$
(9)
This equilibrium condition only requires that the total force on each site, arising from all of the springs attached to it, be zero Alexander (1998). It does not require that the force $`\stackrel{~}{𝐅}(b)`$ be equal to zero on every bond $`b`$.
In equilibrium, when Eq. (9) is satisfied, the part of $`v_b`$ linear in $`\mathrm{\Delta }𝐮_b`$ disappears from $`\mathrm{\Delta }U_T`$. In this case, it is customary to express $`\mathrm{\Delta }U_T`$ to harmonic order in $`\mathrm{\Delta }𝐮_b`$:
$$\mathrm{\Delta }U_T^{\mathrm{har}}=\frac{1}{2}\underset{b}{}[V_b^{\prime \prime }e_{b0i}e_{b0j}+R_{b0}^1V_b^{}(\delta _{ij}e_{b0i}e_{b0j})]\mathrm{\Delta }u_{bi}\mathrm{\Delta }u_{bj},$$
(10)
where $`e_{b0i}=R_{b0i}/R_{b0}`$ is the unit vector directed along bond $`b`$. Thus the harmonic potential on each bond decomposes into a parallel part, proportional to $`V_b^{\prime \prime }`$, directed along the bond and a transverse part, proportional to $`R_{b0}^1V_b^{}`$, directed perpendicular to the bond. The transverse part vanishes when the force on the bond vanishes.
The harmonic energy $`\mathrm{\Delta }U_T^{\mathrm{har}}`$ does not preserve the invariance with respect to arbitrary rotations of the full nonlinear strain energy $`\mathrm{\Delta }U_T`$ of Eq. (6), under which
$$\mathrm{\Delta }u_{bi}\mathrm{\Delta }u_{bi}^{}(U_{ij}\delta _{ij})R_{b0j}+U_{ij}\mathrm{\Delta }u_{bj},$$
(11)
where $`U_{ij}`$ is a rotation matrix. It does, however preserve this invariance up to order $`\theta ^2`$ but not order $`\theta ^2\mathrm{\Delta }u_b`$ and $`\theta (\mathrm{\Delta }u_b)^2`$, where $`\theta `$ is a rotation angle. For small $`𝜽`$,
$$\mathrm{\Delta }𝐮_b^{}=\mathrm{\Delta }𝐮_b+𝜽\times 𝐑_{b0}+O(\theta ^2,\theta \mathrm{\Delta }u_b),$$
(12)
and $`𝐞_{b0}\mathrm{\Delta }𝐮_b^{}=𝐞_{b0}\mathrm{\Delta }𝐮_b+O(\theta ^2,\theta \mathrm{\Delta }u_b)`$. Thus, the part of the harmonic energy arising from the $`k(b)`$ term in Eq. (6) is invariant to the order stated above. The invariance of the force term of Eq. (6) is more subtle. Under the above transformation of Eq. (12), $`(\mathrm{\Delta }u_b^{})^2=(\mathrm{\Delta }u_b)^2+2𝜽\times 𝐑_b\mathrm{\Delta }𝐮_b+(𝜽\times 𝐑_b)^2+O(\theta ^2\mathrm{\Delta }u_b,\theta (\mathrm{\Delta }u_b)^2)`$, and it would seem that there are terms of order $`\theta `$, and $`\theta ^2`$ in $`\mathrm{\Delta }U_T^{\mathrm{har}}`$. These terms vanish, however, upon summation over $`\mathrm{}`$ and $`\mathrm{}^{}`$ because of the equilibrium force condition of Eq. (9). Thus, the full $`\mathrm{\Delta }U_T^{\mathrm{har}}`$ is invariant under rotations up to order $`\theta ^2`$.
### II.2 Definition of Models
We will consider the following simple models of random lattices.
Model A: Random, zero-force bonds on a periodic lattice. In this model, all sites lie on a periodic Bravais lattice with all bond lengths constant and equal to $`R_{b0}`$, and the rest length $`R_{bR}`$ of each bond is equal to $`R_{b0}`$. The force $`\stackrel{~}{𝐅}(b)`$ on each bond is zero, but the spring constant $`k_b`$ and other properties of the potential $`V_b`$ can vary from site to site. Each lattice site has the same coordination number.
Model B: Random, finite-force bonds on an originally periodic lattice. In this model, sites are originally on a regular periodic lattice, but rest bond lengths $`R_{bR}`$ are not equal to the initial constant bond length on the lattice. Sites in this model will relax to positions $`R_\mathrm{}0`$ with bond lengths $`R_{b0}=|𝐑_\mathrm{}^{}0𝐑_\mathrm{}0|`$ such that the force $`𝐅(\mathrm{})`$ at each site $`\mathrm{}`$ is zero but the force $`\stackrel{~}{𝐅}(b)`$ exerted by each bond $`b`$ is in general not. This model has random stresses and, as we shall see, random elastic moduli as well. The bond vectors $`𝐑_{b0}`$ and spring constant $`k_b`$ are random variables, but the coordination number of each site is not. Random stresses in an originally periodic lattice necessarily induce randomness in the elastic moduli relative to the relaxed lattices with zero force at each site.
Model C: Random, zero-force bonds on a random lattice. In this model, lattice sites are at random positions and have random coordination numbers. The equilibrium length $`R_{b0}`$ varies from bond to bond. The rest length $`R_{bR}`$ of each bond is equal to its equilibrium length so that the force $`\stackrel{~}{𝐅}(b)`$ of each bond is zero. This model, which is meant to describe an amorphous material, is macroscopically but not microscopically homogeneous and isotropic.
Model D: Random finite-force bonds on a random lattice. This is the most general model, and it is the one that provides the best description of glassy and random granular materials. In it, the rest lengths $`R_{bR}`$, the spring constants $`k_b`$, and the coordination number are all random variables. Like Model C, this model describes macroscopically isotropic and homogeneous amorphous material.
Though Models A, B, and C can be viewed as subsets of the most general model D, we find it useful to treat them as distinct models because they each isolate separate causes of randomness in the local elastic modulus or stress. One of our goals, for example, is to show analytically and numerically that the non-affinity correlations arising from structural randomness in models C and D have exactly the same form as those arising from the more controlled periodic models A and B. Another is to study the different effects of random elastic moduli and random stress.
In all of these models the random elastic-modulus tensor can in principle exhibit either short- or long-range correlations in space. To investigate the effects of such long-range correlations, we explicitly construct spring constant distributions with long-range correlations in model A. We will also find evidence of long-range correlations in model C when the reference lattice has correlated crystalline domains.
### II.3 Continuum Models
In the continuum limit, when spatial variations are slow on a scale set by the lattice spacing, the equilibrium lattice positions become continuous positions $`𝐱`$ in the reference space: $`𝐑_\mathrm{}0𝐱`$; and the target-space position and displacement vectors become functions of $`𝐱`$: $`𝐑_{\mathrm{}}𝐑(𝐱)`$ and $`𝐮_{\mathrm{}}𝐮(𝐱)`$. In this limit, the lattice strain $`v_b`$ becomes
$$v_bR_{bi}^0R_{bj}^0u_{ij}(𝐱),$$
(13)
where
$$u_{ij}(𝐱)=\frac{1}{2}(_iu_j+_ju_i+_i𝐮_j𝐮)$$
(14)
is the full Green-Saint Venant Lagrangian nonlinear strain Love (1944); Landau and Lifshitz (1986); Chaikin and Lubensky (1995), which is invariant with respect to rigid rotations in the target space \[i.e., with respect to rigid rotations of $`𝐑(𝐱)`$\]. Sums over lattice sites of the form $`_{\mathrm{}}S(\mathrm{})`$, for any function $`S(\mathrm{})`$, can be replaced by integrals $`d^dxS(𝐱)/v(𝐱)`$ where $`v(𝐱)`$ is the volume of the Voronoi cell centered at position $`𝐱=𝐑_\mathrm{}0`$. The continuum energy is then
$$=d^dx\left[\frac{1}{2}K_{ijkl}(𝐱)u_{ij}(𝐱)u_{kl}(𝐱)+\stackrel{~}{\sigma }_{ij}(𝐱)u_{ij}(𝐱)\right],$$
(15)
where
$$\stackrel{~}{\sigma }_{ij}(𝐱)=\frac{1}{2v(𝐱)}\underset{\mathrm{}^{}}{}\stackrel{~}{F}_i(b)R_{b0j}|_{b=<\mathrm{}^{},\mathrm{}>}$$
(16)
is a local symmetric stress tensor at $`𝐱`$ where the sum over $`\mathrm{}^{}`$ is over all bonds with one end at $`\mathrm{}`$ and
$$K_{ijkl}(𝐱)=\frac{1}{2v(𝐱)}\underset{\mathrm{}^{}}{}k(b)R_{b0}^2R_{b0i}R_{b0j}R_{b0k}R_{b0l}|_{b=<\mathrm{}^{},\mathrm{}>}$$
(17)
is the local elastic-modulus tensor pressure . Because it depends only on the full nonlinear strain $`u_{ij}(𝐱)`$, the continuum energy $``$ of Eq. (15) is invariant with respect to rigid rotations in the target space. This is a direct result of the fact that we consider only internal forces between particles. The stress tensor $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ is generated by these internal forces, and as a result, it multiplies $`u_{ij}`$ in $``$. It is necessarily symmetric, and it transforms like a tensor in the reference space. (It is not, however, the second Piola-Kirchoff tensor Marsden and Hughes (1968), $`\sigma _{ij}^{II}=\delta /\delta u_{ij}(𝐱)=K_{ijkl}u_{kl}+\stackrel{~}{\sigma }_{ij}`$, which also transforms in this way.) External stresses, on the other hand, specify a force direction in the target space and couple to the linear part of the strain.
Since $`K_{ijkl}(𝐱)`$ in Eq. (17) arises from central forces on bonds, it and its average over randomness obey the Cauchy relations Love (1944); Born and Huang (1954), $`K_{ijkl}(𝐱)=K_{ikjl}(𝐱)=K_{iljk}(𝐱)`$, in addition to the more general symmetry relations, $`K_{ijkl}(𝐱)=K_{jikl}(𝐱)=K_{ijlk}(𝐱)=K_{klij}(𝐱)`$. The Cauchy relations reduce the number of independent elastic moduli in the average modulus $`K_{ijkl}=K_{ijkl}(𝐱)`$ below the maximum number permitted for a given point-group symmetry (for the lowest symmetry, from 21 to 15). In particular, they reduce the number of independent moduli in isotropic and hexagonal systems from two to one, setting the Lamé coefficients $`\lambda `$ and $`\mu `$ equal to each other. In our analytical calculations, we will, however, treat $`\lambda `$ and $`\mu `$ as independent. The Cauchy limit is easily obtained by setting $`\lambda =\mu `$.
The stress tensor $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ is generated entirely by internal forces on bonds. The elastic-modulus tensor $`K_{ijkl}(𝐱)`$ depends on the local effective spring constant $`k(b)`$, the length and direction of the bond vectors $`𝐑_{b0}`$, and the site coordination number; and it will be a random function of position if any of these variables are random functions of position. Thus $`K_{ijkl}(𝐱)`$ is a random function of position in Models A to D. The stress tensor $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ is nonzero only if the bond forces are nonzero. It is thus a random function of position only in Models B and D.
We require that the continuum limit of our lattice models be in mechanical equilibrium when $`𝐮(𝐱)=0`$. This means that the linear variation of $``$ with respect to $`𝐮(𝐱)`$ must be zero, i.e., that
$$\delta =d^dx\stackrel{~}{\sigma }_{ij}(𝐱)_j\delta u_i(𝐱)=0$$
(18)
for any $`\delta u_i(𝐱)`$. $`\delta u_i(𝐱)`$ can be decomposed into a constant strain part and a part whose average strain vanishes: $`\delta u_i(𝐱)=\delta \gamma _{ij}x_j+\delta 𝐮^{}(𝐱)`$ where $`d^d_j\delta u_i^{}(𝐱)=0`$. Equilibrium with respect to variations in $`\gamma _{ij}`$ implies that the spatial average of $`\stackrel{~}{\sigma }_{ij}`$ is zero. Equilibrium with respect to $`\delta 𝐮^{}(𝐱)`$ implies that when $`𝐱`$ is in the interior of the sample,
$$f_i(𝐱)=_j\stackrel{~}{\sigma }_{ij}(𝐱)=_j\stackrel{~}{\sigma }_{ji}=0,$$
(19)
where $`𝐟`$ is the force density that is a vector in the target space. In addition, $`𝑑S_j\stackrel{~}{\sigma }_{ij}(𝐱)\delta u_j^{}(𝐱)=0`$ for any $`\delta u_j(𝐱)`$, where the integral is over the surface of the sample, implying that $`\stackrel{~}{\sigma }_{ij}(𝐱_B)=0`$ for points $`𝐱_B`$ on the surface.
Thus, we see that equilibrium conditions in the reference space impose stringent constraints on the random stress tensor $`\stackrel{~}{\sigma }_{ij}(𝐱)`$: its spatial average must be zero, its values on sample surfaces must be zero, and it must be purely transverse, i.e., it must have no longitudinal components parallel to the gradient operator. Though the linear part of $`u_{ij}`$ does not contribute to the stress term in $``$, the nonlinear part still does, and $``$ can be written as
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^dx[K_{ijkl}(𝐱)u_{ij}(𝐱)u_{kl}(𝐱)`$ (20)
$`+\stackrel{~}{\sigma }_{ij}(𝐱)_iu_k(𝐱)_ju_k(𝐱)].`$
Because of the constraints on $`\stackrel{~}{\sigma }_{ij}`$, this free energy is identical to that of Eq. (15). It is invariant with respect to rotations in the target space even though it is written so that the explicit dependence on the rotationally invariant strain is not so evident rot\_inv .
As we have seen, the spatial average of $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ is zero; it only has a random fluctuating part in models we consider. The elastic-modulus tensor $`K_{ijkl}(𝐱)`$, on the other hand, has an average part and a random part with zero mean:
$$K_{ijkl}(𝐱)=K_{ijkl}+\delta K_{ijkl}(𝐱).$$
(21)
We will view both $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ and $`\delta K_{ijkl}(𝐱)`$ as quenched random variables with zero mean.
## III Strains and nonaffinity
Consider a reference elastic body in the shape of a regular parallelepiped. When such a body is subjected to stresses that are uniform across each of its faces, it will undergo a strain deformation in which its boundary sites at positions $`𝐱_B`$ distort to new positions
$$R_i(𝐱_B)=\mathrm{\Lambda }_{ij}x_{Bj},$$
(22)
where $`\mathrm{\Lambda }_{ij}`$ is the deformation gradient tensor Marsden and Hughes (1968). If the medium is spatially homogeneous, then $`\mathrm{\Lambda }_{ij}\delta _{ij}+\gamma _{ij}`$ determines the displacements of all points in the medium: $`R_i(𝐱)=\mathrm{\Lambda }_{ij}x_j`$ or $`u_i(𝐱)=\gamma _{ij}x_j`$. Such a distortion is called affine. In inhomogeneous elastic media, there will be local deviations from affinity \[Fig. 1\] described by a displacement variable $`𝐮^{}(𝐱)`$ defined via
$$R_i(𝐱)=\mathrm{\Lambda }_{ij}x_j+u_i^{}(𝐱)$$
(23)
or, equivalently,
$`u_i(𝐱)`$ $`=`$ $`\gamma _{ij}x_j+u_i^{}(𝐱)`$ (24)
$`u_{ij}(𝐱)`$ $``$ $`\gamma _{ij}^Sx_j+(_iu_j^{}+_ju_i^{}`$ (25)
$`+\gamma _{ip}_ju_p^{}+\gamma _{jp}_iu_p^{})/2,`$
where the final equation contains only terms up to linear order in $`𝐮^{}`$ and where $`\gamma _{ij}^S=(\gamma _{ij}+\gamma _{ji}+\gamma _{ik}\gamma _{jk})/2`$. Since distortions at the boundary are constrained to satisfy Eq. (22), $`u_i^{}(𝐱_B)`$ is zero for all points $`𝐱_B`$ on the boundary. It is often useful to consider periodic boundary conditions in which $`𝐮^{}(𝐱)`$ has the same value (possibly not zero) on opposite sides of the parallelepiped. This condition implies
$$d^dx_ju_i^{}(𝐱)=𝑑S_ju_i^{}=0.$$
(26)
### III.1 Nonaffinity in $`1d`$
To develop quantitative measures of nonaffinity, it is useful to consider a simple one-dimensional model, which can be solved exactly. We study a one-dimensional periodic lattice, depicted in Fig. 2 with sites labelled by $`\mathrm{}=0,\mathrm{},N`$, whose equilibrium positions are $`R_\mathrm{}0=a\mathrm{}`$, where $`a`$ is the rest bond length. Harmonic springs with spring constant $`k_{\mathrm{}}k+\delta k_{\mathrm{}}`$ connect sites $`\mathrm{}`$ and $`\mathrm{}1`$, where $`k=(_{\mathrm{}}k_{\mathrm{}})/N`$ is the average spring constant and $`_{\mathrm{}}\delta k_{\mathrm{}}=0`$. The lattice is stretched from its equilibrium length $`Na`$ to a new length $`L=\gamma Na`$. If all $`k_{\mathrm{}}`$’s were equal, the lattice would undergo an affine distortion with $`R_{\mathrm{}}=\gamma a\mathrm{}`$. When the $`k_{\mathrm{}}`$’s are random, sites undergo an additional nonaffine displacement $`u_{\mathrm{}}^{}`$ so that $`R_{\mathrm{}}=\gamma a\mathrm{}+u_{\mathrm{}}^{}`$. The energy is thus
$$=\frac{1}{2}\underset{\mathrm{}=1}{\overset{N}{}}k_{\mathrm{}}(\gamma a+u_{\mathrm{}}^{}u_\mathrm{}1^{})^2.$$
(27)
In equilibrium, the force $`F_{\mathrm{}}=/u_{\mathrm{}}^{}`$ on each bond is zero. The resulting equation for $`u_{\mathrm{}}^{}`$ is
$$F_{\mathrm{}}=k_{\mathrm{}+1}(\gamma a+u_{\mathrm{}+1}^{}u_{\mathrm{}}^{})k_{\mathrm{}}(\gamma a+u_{\mathrm{}}^{}u_\mathrm{}1^{})=0,$$
(28)
which can be rewritten as
$$\mathrm{\Delta }_+k_{\mathrm{}}\mathrm{\Delta }_{}u_{\mathrm{}}^{}=\gamma a\mathrm{\Delta }_+\delta k_l,$$
(29)
where $`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_{}`$ are difference operators defined via $`\mathrm{\Delta }_+A_{\mathrm{}}=A_{\mathrm{}+1}A_{\mathrm{}}`$ and $`\mathrm{\Delta }_{}=A_{\mathrm{}}A_\mathrm{}1`$ for any function $`A_{\mathrm{}}`$. The Fourier transforms of $`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_{}`$ are, respectively, $`\mathrm{\Delta }_+(q)=e^{iq}1`$ and $`\mathrm{\Delta }_{}(q)=1e^{iq}`$. Equations (28) and (29) must be supplemented with boundary conditions. We use periodic boundary conditions for which $`u_N^{}=u_0^{}`$ or equivalently
$$\underset{\mathrm{}=1}{\overset{N}{}}\mathrm{\Delta }_{}u_{\mathrm{}}^{}=0.$$
(30)
The solution to Eq. (29) can be written as the sum of a solution,
$$u_{\mathrm{}}^I=(\mathrm{\Delta }_+k_{\mathrm{}}\mathrm{\Delta }_{})^1\gamma a\mathrm{\Delta }_+\delta k_{\mathrm{}}=\gamma a\mathrm{\Delta }_{}^1\frac{\delta k_{\mathrm{}}}{k_{\mathrm{}}},$$
(31)
to the inhomogeneous equation and a solution,
$$\mathrm{\Delta }_+k_{\mathrm{}}\mathrm{\Delta }_{}u_{\mathrm{}}^H=0,$$
(32)
to the homogeneous one. The latter solution is $`u_{\mathrm{}}^H=\mathrm{\Delta }_{}^1C/k_{\mathrm{}}`$ where $`C`$ is an as yet undetermined constant. Adding the two solutions we obtain
$$u_{\mathrm{}}^{}=\mathrm{\Delta }_{}^1\left(\gamma a\frac{\delta k_{\mathrm{}}}{k_{\mathrm{}}}+\frac{C}{k_{\mathrm{}}}\right),$$
(33)
which implies $`\mathrm{\Delta }_{}u_{\mathrm{}}^{}=(\gamma a\delta k_{\mathrm{}}+C)/k_{\mathrm{}}`$. The boundary condition of Eq. (30) determines $`C`$, and the final solution for $`u_{\mathrm{}}^{}`$ is
$`u_{\mathrm{}}^{}`$ $`=\gamma a\mathrm{\Delta }_{}^1{\displaystyle \frac{1}{k_{\mathrm{}}}}\left[\delta k_{\mathrm{}}({\displaystyle k_{\mathrm{}}^1})^1{\displaystyle k_{\mathrm{}}^1\delta k_{\mathrm{}}}\right]`$
$`=\gamma a\mathrm{\Delta }_{}^1\left({\displaystyle \frac{k_{\mathrm{}}^1}{N^1k_{\mathrm{}}^1}}1\right)\gamma a\mathrm{\Delta }_{}^1S_{\mathrm{}}.`$ (34)
The quantity
$$S_{\mathrm{}}=1\frac{(1+p_{\mathrm{}})^1}{N^1_1^{\mathrm{}}(1+p_{\mathrm{}})^1}$$
(35)
depends only on the ratio $`p_{\mathrm{}}=\delta k_{\mathrm{}}/k`$.
Equation (34) is the complete solution for $`u_{\mathrm{}}^{}`$ for an arbitrary set of spring constants $`k_{\mathrm{}}`$. In our model, these spring constants are taken to be random variables, and information about the nonaffinity is best represented by correlation functions of the nonaffine displacement, averaged over the ensemble of random $`k_{\mathrm{}}`$’s. The simplest of these is the two-point function $`G(\mathrm{}\mathrm{}^{})=u_{\mathrm{}^{}}^{}u_{\mathrm{}}^{}`$, where $``$ represents an average over $`k_{\mathrm{}}`$. $`G(\mathrm{}\mathrm{}^{})`$ is easily calculated from Eq. (34); its Fourier transform is
$$G(q)=(\gamma a)^2\frac{\mathrm{\Delta }^S(q)}{2(1\mathrm{cos}q)},$$
(36)
where $`\mathrm{\Delta }^S(q)`$ is the Fourier transform of $`\mathrm{\Delta }^S(\mathrm{}^{}\mathrm{})=S_{\mathrm{}^{}}S_{\mathrm{}}`$.
There are several important observations that follow from the expression Eq. (36) and that generalize to higher dimensions.
$``$ $`\mathrm{\Delta }^S(\mathrm{}^{},\mathrm{})`$ depends only on the ratios $`\delta k_{\mathrm{}}/k`$ and $`\delta k_{\mathrm{}^{}}/k`$, and it increases with increasing width of the distribution of $`\delta k_{\mathrm{}}`$. To lowest order in averages in $`\delta k_{\mathrm{}}`$,
$`\mathrm{\Delta }^S(\mathrm{}^{},\mathrm{})`$ $`=`$ $`k^2[\mathrm{\Delta }^k(\mathrm{}^{},\mathrm{})N^1{\displaystyle \underset{\mathrm{}_1}{}}\mathrm{\Delta }^k(\mathrm{},\mathrm{}_1)]`$ (37)
$``$ $`k^2\mathrm{\Delta }^k(\delta _\mathrm{},\mathrm{}^{}N^1),`$
where $`\mathrm{\Delta }^k(\mathrm{}^{},\mathrm{})=\delta k_{\mathrm{}^{}}\delta k_{\mathrm{}}`$. The final form applies to uncorrelated distributions in which spring constants on different bonds are independent and $`\delta k_{\mathrm{}}\delta k_{\mathrm{}^{}}=\mathrm{\Delta }^k\delta _\mathrm{},\mathrm{}^{}`$. As the width of the distribution increases, higher moments in $`\delta k_{\mathrm{}}`$ become important in $`\mathrm{\Delta }^S(\mathrm{}^{},\mathrm{})`$. If we assume that the only nonvanishing fourth order moments are of the form $`\mathrm{\Delta }^{(k,4)}(\mathrm{}^{},\mathrm{})=(\delta k_{\mathrm{}^{}})^2(\delta k_{\mathrm{}})^2`$, then the fourth-order contributions to $`\mathrm{\Delta }^S`$ are
$`k^4\mathrm{\Delta }^{(S,4)}(\mathrm{}^{},\mathrm{})`$
$`=\left(1{\displaystyle \frac{4}{N}}+{\displaystyle \frac{6}{N^3}}\right)\left[\mathrm{\Delta }^{(k,4)}(\mathrm{}^{},\mathrm{}){\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{}_1}{}}\mathrm{\Delta }^{(k,4)}(\mathrm{}^{},\mathrm{}_1)\right]`$
$`+2\left(1{\displaystyle \frac{3}{N^2}}\right)\mathrm{\Delta }^{(k,4)}(0)\left(\delta _\mathrm{}^{},\mathrm{}{\displaystyle \frac{1}{N}}\right)`$
$`{\displaystyle \frac{1}{N}}\left(2{\displaystyle \frac{3}{N}}\right){\displaystyle \underset{\mathrm{}_1}{}}\mathrm{\Delta }^{(k,4)}(\mathrm{},\mathrm{}_1)\left(\delta _\mathrm{},\mathrm{}^{}{\displaystyle \frac{1}{N}}\right).`$ (38)
For uncorrelated distributions, $`\mathrm{\Delta }^{(k,4)}(\mathrm{},\mathrm{}^{})=\mathrm{\Delta }^{(k,4)}\delta _\mathrm{},\mathrm{}^{}+(\mathrm{\Delta }^k)^2(1\delta _\mathrm{},\mathrm{}^{})`$. Thus, for uncorrelated distributions in the limit $`N\mathrm{}`$,
$`\mathrm{\Delta }^S(q)=\left({\displaystyle \frac{\mathrm{\Delta }^k}{k^2}}+{\displaystyle \frac{3\mathrm{\Delta }^{(k,4)}(\mathrm{\Delta }^k)^2}{k^4}}\right)(1\delta _{q,0})`$ (39)
to fourth order in $`\delta k_{\mathrm{}}`$. Note that the constraint $`_{\mathrm{}}\delta k_{\mathrm{}}=0`$ requires $`_{\mathrm{}}\mathrm{\Delta }^S(\mathrm{},\mathrm{}^{})=_{\mathrm{}^{}}\mathrm{\Delta }^S(\mathrm{},\mathrm{}^{})=0`$ and thus that $`\mathrm{\Delta }^S(q=0)=0`$. This condition is imposed by the factor $`1\delta _{q,0}`$ in Eq. (39), which implies that $`lim_{q0}\mathrm{\Delta }^S(q)\mathrm{\Delta }^S(q=0)=0`$, i.e., $`\mathrm{\Delta }^S(q)`$ does not approach zero as some power of $`q`$ as $`q0`$.
It is easy to verify that Eq. (39) is exactly the same result that would have been obtained using only the solution \[Eq. (31)\] to the inhomogeneous equation for $`u_{\mathrm{}}^{}`$ with $`\delta k_{\mathrm{}}/k_{\mathrm{}}`$ replaced by $`\delta k_{\mathrm{}}/k_{\mathrm{}}N^1_{\mathrm{}}\delta k_{\mathrm{}}/k_{\mathrm{}}`$. Thus, to obtain the solution for $`G(q)`$ to leading order $`1/N`$, we can ignore the boundary condition, Eq. (30), and use the solution to the inhomogeneous equation with the constraint that $`q^2G(q)`$ be zero at $`q=0`$. This observation will considerably simplify our analysis of the more complicated higher-dimensional problem.
$``$ If correlations in $`\delta k_{\mathrm{}}`$ are of finite range, then $`\mathrm{\Delta }^S(q)`$ has a well defined $`q0`$ limit. In this limit,
$$G(q)=(\gamma a)^2\frac{\mathrm{\Delta }^S(0)}{q^2}\frac{(\gamma a)^2}{q^2k^2}\mathrm{\Delta }^k(q=0),$$
(40)
where $`\mathrm{\Delta }^A(0)lim_{q0}\mathrm{\Delta }^A(q)`$ for $`A=S,k`$. Thus, there is a $`q^2`$ divergence in $`G(q)`$, and the spatial correlation function $`𝒢(\mathrm{}^{},\mathrm{})=(u_{\mathrm{}}^{}u_{\mathrm{}})^2`$ diverges linearly in separation
$$𝒢(\mathrm{},\mathrm{}^{})(\gamma a)^2\mathrm{\Delta }^S(0)|\mathrm{}^{}\mathrm{}|(\gamma a)^2\frac{\mathrm{\Delta }^k(0)}{k^2}|\mathrm{}^{}\mathrm{}|.$$
(41)
$``$ If correlations in $`S_{\mathrm{}}`$ extend out to a distance $`\xi `$, then $`\mathrm{\Delta }^S(q\xi )`$ becomes a function of $`q\xi `$. Long-range correlations in $`S_{\mathrm{}}`$ will lead to long range correlations in $`G(q)q^2\mathrm{\Delta }^S(q\xi )`$, and $`𝒢(\mathrm{},0)`$ will grow more rapidly than $`\mathrm{}`$ for $`1\mathrm{}\xi `$. It is possible that this is the correlation length that diverges at the jamming point $`J`$ in granular media Wyart et al. (2005); Silbert et al. (2005). We will discuss this point further in Sec. III.5.
### III.2 Nonaffinity for $`d>1`$
The nonaffinity correlation function,
$$G_{ij}(𝐱,𝐱^{})=u_i^{}(𝐱)u_j^{}(𝐱^{}),$$
(42)
for $`d>1`$ has a form very similar to that for $`d=1`$, except that it has more complex tensor indices. We will be primarily interested in the scalar part of this function, obtained by tracing over the indices $`i`$ and $`j`$. The Fourier transform of this function scales as
$$G(𝐪)G_{ii}(𝐪)\frac{\gamma ^2}{q^2}\mathrm{\Delta }^S(𝐪)\frac{\gamma ^2}{q^2}\frac{\mathrm{\Delta }^K(𝐪)}{K^2},$$
(43)
where $`\gamma `$ represents the appropriate components of the applied strain and $`\mathrm{\Delta }^S(𝐱,𝐱^{})`$ is in general a nonlinear function of the ratio of the fluctuating components $`\delta K_{ijkl}(𝐱)`$ of the elastic-modulus tensor to its uniform components $`K_{ijkl}`$. To lowest order in the variance, $`\mathrm{\Delta }^S\mathrm{\Delta }^K/K^2`$ where $`\mathrm{\Delta }^K`$ represents components of the variance of the elastic-modulus tensor and $`K`$ components of its average. Thus, the nonaffinity correlation function in coordinate space is proportional to $`|𝐱|^{(d2)}`$ in dimension $`d`$, or
$`𝒢(𝐱)`$ $`=`$ $`(𝐮^{}(𝐱)𝐮^{}(0))^2`$ (44a)
$``$ $`A\mathrm{ln}(|𝐱|/B)d=2`$ (44b)
$``$ $`CD|𝐱|^1d=3,`$ (44c)
where
$`A`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\gamma ^2\mathrm{\Delta }^S(0){\displaystyle \frac{1}{\pi }}\gamma ^2{\displaystyle \frac{\mathrm{\Delta }^K(0)}{K^2}}`$ (45a)
$`B`$ $`=`$ $`(\alpha \mathrm{\Lambda })^1`$ (45b)
$`C`$ $`=`$ $`\gamma ^2\mathrm{\Delta }^S(0){\displaystyle \frac{\mathrm{\Lambda }}{\pi ^2}}`$ (45c)
$`D`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\gamma ^2\mathrm{\Delta }^S(0),`$ (45d)
where $`\mathrm{\Lambda }=2\pi /a`$ is the upper momentum cutoff for a spherical Brillouin zone with $`a`$ the short distance cutoff and $`\alpha =0.8905`$ is evaluated in App. C. The length $`B`$ depends on the spatial form and range $`\xi `$ of local elastic-modulus correlations. We will derive explicit forms for it shortly. In our numerical simulations, we allow the bond spring constant $`k_b`$ to be a random variable with variance $`\mathrm{\Delta }^k=(\delta k_b)^2`$. Variations in $`k_b`$ in general induce changes in all of the components of $`\delta K_{ijkl}`$, and $`\mathrm{\Delta }^S`$ is an average of a function of $`\delta k_b/k`$ where $`k`$ is the average of $`k_b`$.
In general $`𝒢(𝐱)`$ also has anisotropic contributions whose angular average is zero. We will not consider these contributions in detail, but we do evaluate them analytically in App. C.
When a sample is subjected to a distortion via stresses at its boundaries, the strains can be expressed in terms of an affine strain and deviations from it. Using the expressions in Eq. (25) for these strains, we obtain the energy
$`\delta `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }\{K_{ijkl}_ju_i^{}_lu_k^{}`$ (46)
$`+[\delta K_{ijkl}(𝐱)+\delta _{ik}\stackrel{~}{\sigma }_{jl}(𝐱)]_ju_i^{}_lu_k^{}`$
$`+2\delta K_{ijkl}(𝐱)\gamma _{kl}_ju_i^{}\}`$
to lowest order in $`\gamma _{ij}`$. Minimizing $`\delta `$ with respect to $`𝐮^{}`$, we obtain
$$_j[K_{ijkl}+\delta K_{ijkl}(𝐱)+\delta _{ik}\stackrel{~}{\sigma }_{jl}(𝐱)]_lu_k^{}=_j\delta K_{ijkl}(𝐱)\gamma _{kl}.$$
(47)
This equation shows that the random part of the elastic-modulus tensor times the affine strain acts as a source that drives nonaffine distortions. The random stress, which is transverse, does not drive nonaffinity; it is the continuum limit of the random force. The operator $`_jK_{ijkl}^T(𝐱)_l\delta (𝐱𝐱^{})\chi _{ik}^1(𝐱,𝐱^{})`$, where $`K_{ijkl}^T(𝐱)=K_{ijkl}+\delta K_{ijkl}(𝐱)+\delta _{ik}\stackrel{~}{\sigma }_{jl}(𝐱)`$ is the continuum limit of the dynamical matrix or Hessian discussed in Refs. Lemaitre and Maloney (2005) and Tanguy et al. (2002). The matrix $`\chi _{ij}(𝐱,𝐱^{})`$ is the response $`\delta u_i(𝐱)/\delta f_i(𝐱^{})`$ of the displacement to an external force. The formal solution to Eq. (47) for $`u_i^{}(𝐱)`$ in terms of $`\delta K_{ijkl}(𝐱)`$ and $`\stackrel{~}{\sigma }_{ij}(𝐱)`$ is trivially obtained by operating on both sides with $`\chi _{pi}(𝐱,𝐱^{})`$:
$$u_i^{}(𝐱)=d^dx^{}\chi _{ip}(𝐱𝐱^{})_j^{}\delta K_{pjkl}(𝐱^{})\gamma _{kl},$$
(48)
The random component of the elastic modulus appears both explicitly and in a hidden form in $`\chi _{ip}`$ in this equation.
Equation (48) is the solution to the inhomogeneous equation, Eq. (47). Solutions to the homogeneous equation should be added to Eq. (48) to ensure that the boundary condition $`𝐮^{}(𝐱_B)=0`$ for points $`𝐱_B`$ on the sample boundary is met. As in the $`1D`$ case, however, the contribution from the homogeneous solution vanishes in the infinite volume limit and can be ignored.
To lowest order in the randomness, we replace $`\chi _{ip}`$ in Eq. (48) with its nonrandom counterpart, $`\chi _{ip}^0(𝐱𝐱^{})`$, the harmonic elastic response function $`\delta u_i(𝐱)/\delta f_j(𝐱^{})`$ of a spatially uniform system with elastic-modulus tensor $`K_{ijkl}`$ to an external force $`f_j(𝐱^{})`$. Thus, to lowest order in $`\gamma _{ij}`$,
$$G_{ii^{}}(𝐪)=\chi _{ip}^0(𝐪)\chi _{i^{}p^{}}^0(𝐪)q_jq_j^{}\mathrm{\Delta }_{pjkl;p^{}j^{}k^{}l^{}}^K(𝐪)\gamma _{kl}\gamma _{k^{}l^{}},$$
(49)
where
$$\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K(𝐱,𝐱^{})=\delta K_{ijkl}(𝐱)\delta K_{i^{}j^{}k^{}l^{}}(𝐱^{})$$
(50)
is the variance of the elastic-modulus tensor, which we simply call the modulus correlator. Equation (49) contains all relevant information about nonaffine correlations to lowest order in the imposed strain. It applies to any system with random elastic moduli and stresses regardless of the symmetry of its average macroscopic state.
Our primary interest is in systems whose elastic-modulus tensor is macroscopically isotropic. In these systems, which include two-dimensional hexagonal lattices, $`K_{ijkl}=\lambda \delta _{ij}\delta _{kl}+\mu (\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})`$ is characterized by only two elastic moduli, the shear modulus $`\mu `$ and the bulk modulus $`B=\lambda +(2\mu /d)`$, where $`d`$ is the dimension of the reference space. The Fourier transform of $`\chi _{ij}^0(𝐱,𝐱^{})`$ in an isotropic system is
$$\chi _{ij}^0(𝐪)=\frac{1}{(\lambda +2\mu )q^2}\widehat{q}_i\widehat{q}_j+\frac{1}{\mu q^2}(\delta _{ij}\widehat{q}_i\widehat{q}_j).$$
(51)
The modulus correlator is an eighth-rank tensor. At $`𝐪=0`$, it has eight independent components in an isotropic medium (See App. A) and more in media with lower symmetry, including hexagonal symmetry. As discussed above, however, all components of $`\delta K_{ijkl}`$ are proportional to $`\delta k_b`$.
We show in App. B that $`G(𝐪)`$ has the general form
$$G(𝐪)=\frac{\gamma _{xy}^2}{\mu ^2q^2}(\mathrm{\Delta }_A+\mathrm{\Delta }_B\widehat{q}_{}^2\mathrm{\Delta }_C\widehat{q}_x^2\widehat{q}_y^2),$$
(52)
where $`\widehat{q}_i=q_i/q`$, $`\widehat{q}_{}^2=\widehat{q}_x^2+\widehat{q}_y^2`$ and $`\mathrm{\Delta }_A`$, $`\mathrm{\Delta }_B`$ and $`\mathrm{\Delta }_C`$ are linear combinations of the independent components of $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K`$ times a function of $`\lambda /\mu `$. Thus, in general $`𝒢(𝐱)`$ will have anisotropic parts that depend on the direction of $`𝐱`$ in addition to an isotropic part that depends only on the magnitude of $`𝐱`$. In App. C, we derive expressions for the full form of $`𝒢(𝐱)`$. Here we discuss only the isotropic part, which has the from of Eq. (44) with
$`A`$ $`={\displaystyle \frac{\gamma _{xy}^2}{\pi \mu ^2}}[\mathrm{\Delta }_A+\mathrm{\Delta }_B{\displaystyle \frac{1}{8}}\mathrm{\Delta }_C]`$ (53)
$`B`$ $`=(\alpha \mathrm{\Lambda })^1`$ (54)
$`C`$ $`={\displaystyle \frac{\gamma _{xy}^2}{\pi ^2\mu ^2}}[\mathrm{\Delta }_A+{\displaystyle \frac{2}{3}}\mathrm{\Delta }_B{\displaystyle \frac{1}{15}}\mathrm{\Delta }_C]\mathrm{\Lambda }`$ (55)
$`D`$ $`={\displaystyle \frac{\gamma _{xy}^2}{2\pi \mu ^2}}[\mathrm{\Delta }_A+{\displaystyle \frac{2}{3}}\mathrm{\Delta }_B{\displaystyle \frac{1}{15}}\mathrm{\Delta }_C].`$ (56)
In two dimensions, the anisotropic term is proportional to $`\mathrm{cos}4\psi `$ where $`\psi `$ is the angle that $`𝐱`$ makes with the $`x`$-axis. In the limit of large $`|𝐱|`$, the coefficient of $`\mathrm{cos}4\psi `$ is a constant. In three dimensions, the anisotropic terms are more complicated. In both two and three dimensions, however, the average of the anisotropic terms over angles are zero.
### III.3 Other Measures of Nonaffinity
The nonaffinity correlation function $`G_{ij}`$ (and its cousin $`𝒢`$) is not the only measures of nonafinity, though other measures can usually be represented in terms of it. Perhaps the simplest measure of nonaffinity is simply the mean-square fluctuation in the local value of of $`𝐮^{}(𝐱)`$, which is the equal-argument limit of the trace of $`G_{ij}(𝐱^{},𝐱)`$:
$$[u^{}(𝐱)]^2=G_{ii}(𝐱,𝐱).$$
(57)
This measure was used in Ref. Langer and Liu (1997) to measure nonaffinity in models for foams. In three dimensions, it is a number that depends on the cutoff, $`a^1`$: $`[u^{}(𝐱)]^2\gamma ^2(\mathrm{\Delta }^K/K^2)a^1`$; in two dimensions, it diverges logarithmically with the size of the sample $`L`$: $`[u^{}(𝐱)]^2\gamma ^2(\mathrm{\Delta }^K/K^2)\mathrm{ln}(L/a)`$.
References Head et al. (2003a, b, c), which investigate a two-dimensional model of crosslinked semi-flexible rods designed to describe crosslinked networks of actin and other semi-flexible biopolymers, introduce \[Fig. 3\] a measure based on comparing the angle $`\theta \theta (𝐱^{},𝐱)`$ that the vector connecting two sites originally at $`𝐱`$ and $`𝐱^{}`$ makes with some fixed axis after nonaffine distortion under shear to the angle $`\theta _0\theta _0(𝐱^{})\theta _0(𝐱)`$ that that vector would make if the points were affinely distorted:
$$𝒢_\theta (𝐱^{}𝐱)=[\theta (𝐱^{},𝐱)\theta _0(𝐱^{},𝐱)]^2.$$
(58)
Under affine distortion, the vector connecting points $`𝐱^{}`$ and $`𝐱`$ is $`r_i=x_i^{}x_i+\gamma _{ij}(x_j^{}x_j)`$; under nonaffine distortion, the separation is $`𝐫^{}=𝐫+𝐮^{}(𝐱^{})𝐮^{}(𝐱)`$. In two dimensions,
$$𝐫\times 𝐫^{}=rr^{}\mathrm{sin}(\theta \theta _0)𝐞_z=𝐫\times [𝐮^{}(𝐱^{})𝐮^{}(𝐱)],$$
(59)
where $`𝐞_z`$ is the unit vector along the $`z`$ direction perpendicular to the two-dimensional plane and $`r=|𝐫|`$. If both $`\gamma `$ and $`|𝐮^{}(𝐱^{})𝐮^{}(𝐱)|/|𝐱^{}𝐱|`$ are small,
$$\theta (𝐱^{},𝐱)\theta _0(𝐱^{},𝐱)\frac{𝐞_z[(𝐱^{}𝐱)\times (𝐮^{}(𝐱^{})𝐮^{}(𝐱))]}{|𝐱^{}𝐱|^2},$$
(60)
and
$$𝒢_\theta (𝐱)=ϵ_{ij}ϵ_{kl}\frac{x_ix_j}{|𝐱|^4}𝒢_{kl}(𝐱)\frac{1}{|𝐱|^2}𝒢(𝐱)$$
(61)
where $`ϵ_{ij}=ϵ_{zij}`$ is the two-dimensional antisymmetric symbol, and $`𝒢_{ij}(𝐱)=[u_i^{}(𝐱)u_i^{}(0)][u_j^{}(𝐱)u_j^{}(0)]`$.
### III.4 Generation of Random Stresses
As we have discussed, a system of particles in mechanical equilibrium can be characterized by random elastic moduli and a random local stress tensor with only transverse components. To better understand random stresses, it is useful to consider a model in which random stress is introduced in a material that is initially stress free. We begin with a system with a local elastic-modulus tensor $`K_{ijkl}(𝐱)`$ that can in general be random but with $`\stackrel{~}{\sigma }_{ij}(𝐱)=0`$, and to this we add a local random stress $`\overline{\sigma }_{ij}(𝐱)`$ with zero mean that couples to the rotationally invariant nonlinear strain and that has longitudinal components so that its variance in an isotropic system is
$$\mathrm{\Delta }_{ijkl}^{\overline{\sigma }}=\mathrm{\Delta }_1^{\overline{\sigma }}\delta _{ij}\delta _{kl}+\mathrm{\Delta }_2^{\overline{\sigma }}(\delta _{ik}\delta _{jl}+\delta _{il}+\delta _{jk}).$$
(62)
For simplicity, we assume that the spatial average of $`\overline{\sigma }_{ij}(𝐱)`$ is zero. A random stress of this sort can be generated in a lattice model by making the rest bond length $`R_{bR}`$ a random variable in a system in which initially the rest and equilibrium bond lengths are equal. In the continuum limit, our elastic energy is thus
$$\delta =[\frac{1}{2}K_{ijkl}^\sigma (𝐱)u_{ij}(𝐱)u_{kl}(𝐱)+\overline{\sigma }_{ij}(𝐱)u_{ij}(𝐱)],$$
(63)
where the superscipt $`\sigma `$ on $`K_{ijkl}^\sigma `$ indicates that this is an elastic modulus prior to relaxation in the presence of $`\sigma _{ij}`$.
Sites that were in equilibrium at positions $`𝐱`$ in the original reference space in the absence of $`\overline{\sigma }_{ij}`$ are no longer so in its presence. These sites will undergo displacements to new equilibrium sites $`𝐱^{}𝐑_0(𝐱)=𝐱+𝐮_0(𝐱)`$, which define a new reference space. Positions $`𝐑(𝐱)`$ in the target space can be expressed as displacements relative to the new reference space: $`𝐑(𝐱^{})=𝐱^{}+𝐮^{}(𝐱^{})`$. Then, strains relative to the original reference space can be expressed as the sum of a strain relative to the new reference space and one describing the distortion of the original references space to the new one:
$`u_{ij}(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{R_k(𝐱)}{x_i}}{\displaystyle \frac{R_k(𝐱)}{x_j}}\delta _{ij}\right)`$ (64)
$`=`$ $`u_{ij}^0(𝐱)+\mathrm{\Lambda }_{0ik}^T(𝐱^{})u_{kl}^{}(𝐱^{})\mathrm{\Lambda }_{0lj}(𝐱^{}),`$
where
$`\mathrm{\Lambda }_{0ij}(𝐱)`$ $`=`$ $`{\displaystyle \frac{R_{0i}(𝐱)}{x_j}}=\delta _{ij}+_ju_{0i},`$ (65)
$`u_{ij}^0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Lambda }_{0ki}\mathrm{\Lambda }_{0kj}\delta _{ij}\right),`$ (66)
and
$$u_{ij}^{}(𝐱^{})=(_i^{}u_j^{}+_j^{}u_i^{}+_i^{}u_k^{}_j^{}u_k^{})/2,$$
(67)
where $`_i^{}/x_i^{}`$. Using Eq. (64) in Eq. (63), we obtain $`\delta [𝐮]=\delta [𝐮_0]+\delta ^{}[𝐮^{}]`$, where
$$\delta ^{}=\frac{1}{2}d^dx^{}[K_{ijkl}u_{ij}^{}u_{kl}^{}+\stackrel{~}{\sigma }_{jl}u_{ij}^{}]$$
(68)
with
$$K_{ijkl}(𝐱^{})=(det\mathrm{\Lambda }_0)^1\mathrm{\Lambda }_{0ia}\mathrm{\Lambda }_{0jb}K_{abcd}^\sigma \mathrm{\Lambda }_{0ck}^T\mathrm{\Lambda }_{0dl}^T$$
(69)
and
$$\stackrel{~}{\sigma }_{ij}(𝐱^{})=(det\mathrm{\Lambda }_0)^1\mathrm{\Lambda }_{0ia}(K_{abcd}^\sigma u_{0cd}+\overline{\sigma }_{ab})\mathrm{\Lambda }_{0bj}^T,$$
(70)
where we have not displayed explicitly the dependence of $`\mathrm{\Lambda }_{0ij}`$ on $`𝐱^{}`$. The displacement field $`𝐮_0(𝐱)`$ is determined by the condition that the force density at each point in the new reference state be zero, i.e., so that $`_j^{}\stackrel{~}{\sigma }_{ij}(𝐱^{})=0`$. To linear order in displacement and $`\overline{\sigma }_{ij}`$, this condition is
$$_j(K_{ijkl}u_{0kl}+\overline{\sigma }_{ij})=0,$$
(71)
where to this linearized order, we can ignore the difference between $`𝐱`$ and $`𝐱^{}`$. For an initially isotropic medium, this equation can be solved for $`𝐮_0`$ to yield
$$u_{0i}(𝐪)=\frac{1}{\mu q^2}\left(\delta _{ik}\frac{\lambda +\mu }{\lambda +2\mu }\frac{q_iq_k}{q^2}\right)iq_l\overline{\sigma }_{kl}.$$
(72)
To lowest order in $`𝐮_0`$, the elastic moduli and stress tensors in the new reference state are
$`\stackrel{~}{\sigma }_{ij}(𝐪)`$ $`=`$ $`\delta _{ik}^T\delta _{jl}^T\overline{\sigma }_{kl}{\displaystyle \frac{\lambda }{\lambda +2\mu }}\delta _{ij}^T\widehat{q}_k\widehat{q}_l\overline{\sigma }_{kl}`$ (73)
$`\delta K_{ijkl}`$ $`=`$ $`2\lambda (\delta _{ij}v_{kl}+\delta _{kl}v_{ij})`$
$`+2\mu (\delta _{ik}v_{jl}+\delta _{jl}v_{ik}+\delta _{il}v_{jk}+\delta _{jk}v_{il}),`$
where
$`v_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Lambda }_{0ik}\mathrm{\Lambda }_{0kj}^T\delta _{ij}\right)`$ (75)
$`=`$ $`(_iu_{0j}+_ju_{0i}+_ku_{0i}_ku_{0j})/2`$
$``$ $`(_iu_{0j}+_ju_{0i})/2`$
is the left Cauchy strain tensor relative to the original reference state.
Note that $`\stackrel{~}{\sigma }_{ij}(𝐱^{})`$ is transverse and random as it should be. The elastic-modulus tensor is a random variable via its dependence on $`\mathrm{\Lambda }_{0ij}(𝐱)`$. Thus, a random stress added to an initially homogeneous elastic medium (with $`K_{ijkl}^0`$ nonrandom and independent of $`𝐱`$) produces both a random transverse stress and a random elastic-modulus tensor in the new relaxed reference frame. The statistical properties of $`K_{ijkl}(𝐱^{})`$ are determined in this model entirely by those of $`\overline{\sigma }_{ij}(𝐱)`$, and $`\mathrm{\Delta }^K\mathrm{\Delta }^{\overline{\sigma }}`$. In general, of course, the randomness in $`K_{ijkl}(𝐱^{})`$ arises both from randomness in the original $`K_{ijkl}^\sigma (𝐱)`$ and $`\overline{\sigma }_{ij}(𝐱)`$.
The nonaffinity correlation function can be calculated exactly to lowest order in $`\mathrm{\Delta }_1^{\overline{\sigma }}`$ and $`\mathrm{\Delta }_2^{\overline{\sigma }}`$ when the initial reference state is homogeneous and nonrandom. It has exactly the same form as Eq. (43) when expressed in terms of $`\mathrm{\Delta }^K`$. When expressed in terms of $`\mathrm{\Delta }_1^{\overline{\sigma }}`$ and $`\mathrm{\Delta }_2^{\overline{\sigma }}`$, it has a similar form, which in an isotopic elastic medium can be expressed as
$$G(𝐪)\frac{\mathrm{\Delta }_2^{\overline{\sigma }}\gamma ^2}{\mu ^2q^2}f(\widehat{𝐪},\lambda /\mu ,\mathrm{\Delta }_1^{\overline{\sigma }}/\mathrm{\Delta }_2^{\overline{\sigma }}).$$
(76)
Thus, $`𝒢(𝐱)`$ has the same form in this model as Eq. (44).
### III.5 Long-range Correlations in Elastic Moduli
Long-range correlations in random elastic moduli can significantly modify the behavior of $`𝒢(𝐱)`$. To illustrate this, we consider a simple scaling form for $`\mathrm{\Delta }^K(𝐪)`$ inspired by critical phenomena:
$`\mathrm{\Delta }^K(𝐪)=\xi ^\varphi g(q\xi )`$ (77)
$`\{\begin{array}{cc}\xi ^\varphi g_0[1+(q\xi )^s+\mathrm{}],\hfill & \text{for }q\xi 0,\hfill \\ g_{\mathrm{}}q^\varphi [1+b(q\xi )^t+\mathrm{}],\hfill & \text{as }q\xi \mathrm{},\hfill \end{array}`$ (78)
where $`\xi `$ is a correlation length, $`\varphi `$ is the dominant critical exponent, and $`s`$ and $`t`$ are corrections to scaling exponents. It is possible in principle for each of the components of $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K`$ to be described by difference scaling lengths $`\xi `$ and functions $`g(u)`$. We will assume, however, that $`\xi `$ and the functional form of $`g`$ is the same for all components, but we will allow for the zero-argument value $`g_0`$ to vary. $`G(𝐪)`$ is thus given by Eq. (52) with $`\mathrm{\Delta }_A`$, $`\mathrm{\Delta }_B`$, and $`\mathrm{\Delta }_C`$ replaced by $`\mathrm{\Delta }_A(𝐪)`$, $`\mathrm{\Delta }_B(𝐪)`$, and $`\mathrm{\Delta }_C(𝐪)`$ with scaling forms given by Eq. (77). In this case, $`𝒢(𝐱)`$ can be written as $`(\gamma _{xy}^2/\mu ^2)(𝐱)`$ with $`(𝐱)=_A(𝐱)+_B(𝐱)_C(𝐱)`$, where
$$_\alpha (𝐱)=2\xi ^\varphi \frac{d^dq}{(2\pi )^d}f_\alpha (𝐪)g_\alpha (q\xi )\frac{1}{q^2}(1e^{i𝐪𝐱}),$$
(79)
with $`f_A=1`$, $`f_B(𝐪)=\widehat{q}_{}^2`$, and $`f_C(𝐪)=\widehat{q}_x^2\widehat{q}_y^2`$. There are two important observations to make about the functions $`_\alpha `$. First, for $`q\xi 1`$, $`g(q\xi )`$ can be replaced by its zero $`q`$ limit, $`g_0`$. Thus, as long as $`\xi `$ is not infinite, the asymptotic behavior of $`𝒢(𝐱)`$ for $`|𝐱|\xi `$ is identical to those of Eq. (44) but with amplitudes that increase as $`\xi ^\varphi `$. Second, when $`\xi \mathrm{}`$, the $`q^{(d3\varphi )}`$ behavior of the integrand leads to modified power-law behavior in $`|𝐱|`$ for $`ax\xi `$, where $`a=2\pi /\mathrm{\Lambda }`$ is the short distance cutoff, depending on dimension.
In two dimensions, which is the focus of most of our simulations, the isotropic part of $``$ is
$$_I(𝐱)=\frac{\xi ^\varphi }{\pi }_0^\mathrm{\Lambda }\frac{dq}{q}g(q\xi )[1J_0(q|𝐱|)],$$
(80)
where $`g(y)=g_A(y)+g_B(y)\frac{1}{8}g_C(y)`$ and $`J_0(y)`$ is the zeroth order Bessel function. In the limit $`|𝐱|\xi `$,
$$_I(𝐱)\frac{1}{\pi }g_0\xi ^\varphi \mathrm{ln}\frac{\beta (\mathrm{\Lambda }\xi ,\varphi )}{\xi }|𝐱|,$$
(81)
where $`\beta (\mathrm{\Lambda }\xi ,\varphi )`$ is evaluated in Appendix C. The behavior of $`_I(𝐱)`$ when $`\mathrm{\Lambda }^1|𝐱|\xi `$ depends on the value of $`\varphi `$
$$_I(𝐱)\{\begin{array}{cc}\frac{1}{\pi }g_{\mathrm{}}|𝐱|^\varphi 𝒜_2(\varphi ),\hfill & \text{if }\varphi <2,\hfill \\ \frac{1}{4\pi }g_{\mathrm{}}|𝐱|^2\mathrm{ln}(\nu \xi /|𝐱|),\hfill & \text{if }\varphi =2,\hfill \\ \frac{1}{4\pi }g_{\mathrm{}}|𝐱|^2\xi ^{\varphi 2}𝒞_2(\varphi ),\hfill & \text{if }\varphi >2.\hfill \end{array}$$
(82)
The quantities $`𝒜_2(\varphi )`$, $`𝒞_2(\varphi )`$, and $`\nu `$ are evaluated in Appendix C.
The function $`g(u)`$ can have any form provided its large- and small-$`u`$ limits are given by Eq. (77). A useful model form to consider, of course, is the simple Lorentzian for which $`\varphi =2`$ and
$$g(u)=\frac{g_0}{1+u^2}.$$
(83)
For the purposes of illustration, in Fig. 4 we plot $`_i(|𝐱|)`$ for a family of functions parameterized by the exponent $`\varphi `$:
$$g(u)=\frac{g_0(\varphi )}{(1+u^2)^{\varphi /2}}.$$
(84)
these curves clearly show the crossover from $`|𝐱|^\varphi `$ behavior for $`\mathrm{\Lambda }^1|𝐱|\xi `$ to the characteristic log behavior for $`|𝐱|\xi `$. The correlation length $`\xi `$ and the amplitude $`g_0(\varphi )`$ were set so that the large $`𝐱`$ log behavior is the same for every $`\varphi `$. For this family of crossover functions, the value of $`|𝐱|`$ at which $`(|𝐱|)`$ crosses over from $`|𝐱|^\varphi `$ to logarithmic behavior increases with decreasing $`\varphi `$, and curves with smaller $`\varphi `$ systematically lie above those with large $`\varphi `$.
The limiting forms for $`_I(|𝐱|)`$ in one and three dimensions are given in App. C.
### III.6 Rotational Correlations
The nonaffine displacements generated in random elastic media by external strains contain rotational as well as irrotational components as is evident from Fig. (5). The local nonaffine rotation angle is $`\omega _k(𝐱)=\frac{1}{2}ϵ_{ijk}_ju_k^{}`$, where $`ϵ_{ijk}`$ is the anti-symmetric Levi-Civita tensor, and rotational correlations are measured by the correlation function $`G_{\omega _i\omega _j}(𝐱)=\omega _i(𝐱)\omega _j(0)`$. In two dimensions, there is only one angle $`\omega (𝐱)=\frac{1}{2}ϵ_{ri}_ru_i`$, where $`ϵ_{ri}ϵ_{zri}`$. The Fourier transform of the correlation function $`G_\omega =\omega (𝐱)\omega (0)`$ will then scale as $`\gamma ^2\mathrm{\Delta }^K/\mu ^2`$, approaching a constant rather than diverging as $`𝐪0`$. We show in App. D that
$$G_\omega (𝐪)=\frac{\gamma _{xy}^2}{\mu ^2}[\mathrm{\Delta }_A^\omega (q)\mathrm{\Delta }_C^\omega (q)\widehat{q}_x^2\widehat{q}_y^2]$$
(85)
in two dimensions, where $`\mathrm{\Delta }_A^\omega (q)`$ and $`\mathrm{\Delta }_C^\omega `$ are linear combinations of the independent components of $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K`$. Thus, the rotation correlation function contains direct information about elastic-modulus correlations. If these correlations are short range, and there is no $`q`$ dependence in either $`\mathrm{\Delta }_A^\omega (q)`$ or $`\mathrm{\Delta }_C^\omega `$, the spatial correlation function has an isotropic short-range part and an anisotropic power-law part:
$$G_\omega (𝐱)=\frac{\gamma _{xy}^2}{\mu ^2}\left[\mathrm{\Delta }_A^\omega \delta (𝐱)\mathrm{\Delta }_C^\omega \left(16\frac{x^2y^2}{|𝐱|^6}\frac{2}{|𝐱|^2}\right)\right].$$
(86)
If there are long-range correlations in the elastic moduli with the Lorentzian form of Eq. (83), then
$`G_\omega (𝐱)`$ $`=`$ $`{\displaystyle \frac{\gamma _{xy}^2}{2\pi \mu ^2}}[(\stackrel{~}{\mathrm{\Delta }}_A^\omega {\displaystyle \frac{1}{8}}\mathrm{\Delta }_C^\omega )K_0(|𝐱|/\xi )`$
$`+{\displaystyle \frac{1}{8}}\mathrm{cos}4\psi \stackrel{~}{\mathrm{\Delta }}_C^\omega ({\displaystyle \frac{48\xi ^4}{|𝐱|^4}}+{\displaystyle \frac{4\xi ^2}{|𝐱|^2}}+K_4(|𝐱|/\xi ))],`$
where $`K_n(y)`$ is the Bessel function of imaginary argument. The $`\mathrm{cos}4\psi `$ behavior is for isotropic systems. There will be $`\mathrm{cos}6\psi `$ and higher order terms present in a hexagonal lattice. In Sec. IV, we verify in numerical simulations the exponential decay of the isotropic part of $`G_\omega (|𝐱|)`$ in Model A with long-range correlations in spring constants and the $`|𝐱|^2`$ behavior of the $`\mathrm{cos}4\psi `$ part of $`G_\omega (|𝐱|)`$ in Model C, which is isotropic.
## IV Numerical Minimizations
To further our understanding of nonaffinity in random lattices and to verify our analytic predictions about them, we carried out a series of numerical studies on models A-D described in Sec. II.1. To carry out these studies, we began with an initial lattice – a periodic hexagonal or FCC lattice for models A and B and a randomly tesallated lattice for models C and D. We assigned spring potentials $`V_b(R_b)`$ and rest bond lengths $`R_{bR}`$ to each bond. To study nonaffinity, we subjected lattices to shear and then numerically determined the minimum-energy positions of all sites subject to periodic boundary condition. The elastic energy of the lattice was linearized about the affine shear state. Interestingly, in this linearization the value of the imposed shear, $`\gamma `$, factored out of our calculation, so that $`𝐮^{}(𝐱)`$ was linear in $`\gamma `$ and thus $`𝒢(𝐱)`$ was automatically quadratic in $`\gamma `$. We present below the procedures and results for each model.
### IV.1 Model A
In this model, the initial reference lattice is periodic, and the rest bond length $`R_{bR}`$ is equal to the equilibrium lattice parameter $`R_{b0}`$ for every bond, which we set equal to one. Each bond is assigned an anharmonic potential
$$V_b(R_b)=\frac{1}{2}k_b(\delta R_b^2+\delta R_b^4),$$
(88)
where $`\delta R_R=R_bR_{b0}R_b1`$ and the spring constant $`k_b`$ is a random variable. We chose $`k_b=1+\delta k_b`$ where $`\delta k_b`$ is a random variable with zero mean lying between $`\delta \overline{k}`$ and $`+\delta \overline{k}`$ with $`\delta \overline{k}<1`$.
#### IV.1.1 Independent bonds on hexagonal and FCC lattice
In the simplest versions of model A, the spring constant $`k_b`$ is an independent random variable on each bond of a two-dimensional hexagonal or a three-dimensional FCC lattice. We assign each bond a random value of $`\delta k_b`$ chosen from a flat distribution lying between $`\delta \overline{k}`$ and $`+\delta \overline{k}`$. Randomly distributed spring constants give rise to random local elastic moduli as defined by Eq. (17). We verified that the distribution of the values of the local shear modulus $`K_{xyxy}`$ on a hexagonal lattice for different $`\delta \overline{k}`$ was well fit by a Gaussian function with width linearly proportional to $`\delta \overline{k}`$.
The nonaffinity correlation function $`𝒢(𝐱)=\left|𝐮^{}(𝐱)𝐮^{}(0)\right|^2`$ \[Eq. (44)\] measured on the numerically relaxed lattices is shown in Fig. 6(a). The averages were calculated by summing the differences in deviation for every pair of nodes on the lattice and binning according to the nodes’ separation in the undeformed (reference) state. Note that this process automatically averages over angle, so it produces only the isotropic part of $`𝒢(𝐱)`$. The separation between nodes was taken as the least distance between the nodes across any periodic boundaries. The curves were well fit by the $`A\mathrm{ln}(\left|𝐱\right|/B)`$ dependence on $`\left|𝐱\right|`$ predicted by Eq. (44b). The excellent data collapse achieved by plotting the rescaled function $`𝒢(𝐱)/(\gamma \delta \overline{k})^2`$ demonstrates the quadratic dependence of the amplitude $`A`$ on $`\delta \overline{k}`$. Figure 7 shows the quadratic plus quartic dependence of the amplitude $`A`$ on and $`\delta \overline{k}`$ at larger values. It is worth noting that while all correlation functions were independently fit with a two-parameter function $`A\mathrm{ln}(\left|𝐱\right|/B)`$, the optimal values of $`B`$ in all cases fell within $`10\%`$ of one another.
Figure 6(b) displays $`𝒢(𝐱)/(\delta \overline{k})^2`$ on an FCC lattice as a function of $`\left|𝐱\right|`$ for different $`\delta \overline{k}`$ for $`\gamma =0.1\%`$ fit to the function $`CD/\left|𝐱\right|`$ predicted by Eq. (44). The data collapse verifies the expected dependence of $`C`$ on $`(\delta \overline{k})^2`$.
#### IV.1.2 Correlated random bonds on an hexagonal lattice
As discussed in Sec. III.5, random lattices can exhibit long-range correlations, characterized by a correlation length $`\xi `$, in local elastic moduli that can significantly modify the behavior of nonaffinity correlation functions at distances less than $`\xi `$. To verify the prediction of Sec. III.4, we numerically constructed hexagonal lattices with long-range correlations in bond spring constants. To do this, we set $`k_b=1+\delta k_b`$ where $`\delta k_b`$ was set equal to a small, randomly generated scalar field with proper spatial correlations. This scalar field was created by taking the reverse Fourier transform of the function $`\mathrm{exp}(i\varphi _r)/\sqrt{q^2+\xi ^2}`$, where $`\xi `$ is a variable decay length and $`\varphi _r`$ is a random complex phase. The scalar field in these simulations was normalized to have constant mean squared value and peak values of $`\pm 0.1`$, so that the variation to the local spring constants was at most $`10\%`$. This method of generation yields a clean exponential decay in the two-point correlation function $`\mathrm{\Delta }^K(𝐱,0)\delta K_{xyxy}(𝐱)\delta K_{xyxy}(0)`$ which persists for separations up to three times the correlation length. Figure 8 shows the two-point correlation function $`\mathrm{\Delta }^K(𝐱,0)`$ as a function of separation. The region of exponential correlation was followed by a small region of anti-correlation, which is not pictured. By construction, the distributions of the local shear elastic modulus $`K_{xyxy}`$ were essentially constant, independent of $`\xi `$; thus $`\mathrm{\Delta }^K(0,0)`$ is equal for all curves in Fig. 8.
According to Eq. (111), the growth of the correlation function $`𝒢(𝐱)`$ for large $`\left|𝐱\right|`$ is logarithmic with prefactor proportional to $`g_0\xi ^\varphi `$, where $`\varphi =2`$ for the Lorentzian case we are now considering. The quantity $`g_0\xi ^\varphi `$ is equivalent to $`\mathrm{\Delta }^K(𝐪=0)`$, but this quantity is difficult to measure numerically. However, $`g_0`$ can also be expressed in terms of the coordinate space correlation $`\mathrm{\Delta }^K(𝐱=0)`$. The latter quantity is easily measured by averaging $`\left(\delta K_{xyxy}(𝐱)\right)^2`$ over all nodes. For the form of the correlation function $`g(u)`$ given in Eq. (84),
$$\mathrm{\Delta }^K(𝐱=0)g_0\xi ^{\varphi 2}\{\begin{array}{cc}\frac{1}{\varphi 2}\left(1\left(1+\left(\mathrm{\Lambda }\xi \right)^2\right)^{1\frac{\varphi }{2}}\right)\hfill & \varphi 2,\hfill \\ \frac{1}{2}\mathrm{ln}\left(1+\left(\mathrm{\Lambda }\xi \right)^2\right)\hfill & \varphi =2.\hfill \end{array}$$
(89)
In the limit $`\xi \mathrm{\Lambda }1`$, $`g_0\mathrm{\Delta }^K(𝐱=0)`$ and the large separation form of the correlation function $`𝒢(𝐱)`$ is logarithmic with prefactor $`\gamma ^2\mathrm{\Delta }^K(𝐱=0)\xi ^\varphi /\mu ^2`$.
We have already established that for this set of simulations, $`\mathrm{\Delta }^K(𝐱=0)/\mu ^2`$ is a constant, independent of $`\xi `$ \[see Fig.8\]. In Fig. 9, we plot $`𝒢(𝐱)/\xi ^2`$ versus $`\left|𝐱\right|/\xi `$ for different values of $`\xi `$. We also plot the function $`(\left|𝐱\right|)`$ calculated from Eq. (80) with a Lorentzian $`g(y)`$ \[Eq. (83)\]. The agreement between the numerical and analytical results is excellent with both showing $`\left|𝐱\right|^2`$ behavior for $`\left|𝐱\right|<\xi `$ and $`\mathrm{ln}\left|𝐱\right|`$ behavior for $`\left|𝐱\right|>\xi `$. In Fig. 10 we plot the vorticity correlation function $`G_\omega `$ versus separation rescaled by the correlation length, $`\left|𝐱\right|/\xi `$. The vorticity correlation function decreases exponentially away from zero separation with a decay length $`1.1\times \xi `$; our framework predicted decay with an exponent of $`\xi `$ exactly. The slight discrepency between theory and simulation is not understood.
### IV.2 Model B: Internal stresses
In this model, random stresses are introduced in a periodic lattice via a random distribution of rest bond lengths. We study hexagonal lattices in which the rest lengths of the bonds are multiplied by a factor $`(1+\beta _b)`$ where $`\beta _b`$ is chosen randomly from a flat distribution lying between $`\beta `$ and $`\beta `$ with $`\beta <0.1`$. Once again, the spring constants are set to $`k_b=1+\delta k_b`$, with $`\delta k_b`$ chosen randomly from a flat distribution lying between $`\delta \overline{k}`$ and $`+\delta \overline{k}`$. After specifying the rest length of each bond, we numerically determined the equilibrium state of this random lattice with zero applied stress by minimizing the rest energy over lattice positions and the size of the simulation box (for a system of 40,000 particles, the minimization over box size was only a fraction of a percent). The resulting equilibrium configuration has zero net force on each node. This relaxed state constitutes the reference state of our random system with lattice positions $`𝐑_\mathrm{}0=𝐱`$.
The original lattice before relaxation is characterized by random stresses $`\overline{\sigma }_{ij}`$, which can be calculated from Eq. (16),
$$\overline{\sigma }_{ij}(\mathrm{})=\frac{1}{2v}\underset{\mathrm{}^{}}{}R_{bIi}R_{bIj}k_b\delta R_b/a,$$
(90)
where $`𝐑_{bI}`$ is the bond vector of length $`a`$ (independent of $`b`$) for bond $`b`$ in the initial undistorted hexagonal lattice and $`\delta R_b=aR_{bR}=\beta _ba`$. The average over of $`\overline{\sigma }_{ij}`$ over $`\beta _b`$ is zero: $`\overline{\sigma }_{ij}_b=0`$, and its variance is
$$\overline{\sigma }_{ij}(\mathrm{})\overline{\sigma }_{kl}(\mathrm{})_b=\frac{\beta ^2}{27}(1+\frac{(\delta \overline{k})^2}{3})(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}).$$
(91)
As discussed in Sec. III.4, randomness in $`\overline{\sigma }_{ij}`$ generates a random elastic moduli in the relaxed reference lattice. Figure 11 shows how the random stress broadens the distribution of local elastic moduli. For lattices with $`\delta \overline{k}=0`$, $`\sqrt{\mathrm{\Delta }^K(𝐱=0)}`$ is linearly proportional to $`\beta `$ as predicted by Eqs. (72) to (75).
After constructing the relaxed state, we sheared it in the $`xy`$ plane as before and measured the nonaffinity correlation function. The measurements were well fit by the functional form $`𝒢(𝐱)A\mathrm{ln}(\left|𝐱\right|/B)`$. Figure 12 shows that for $`\delta \overline{k}=0`$ the measured ratio $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ is nearly independent of $`\beta `$, as predicted in Section III.4. For $`\delta \overline{k}>0`$, the ratio $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ is $`20\%`$ lower at small $`\beta `$, but asymptotes to the $`\delta \overline{k}=0`$ value as $`\beta `$ increases, approaching the asymptote more quickly for smaller $`\delta \overline{k}`$. The difference in $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ between stressed and stress-free lattices is most likely a higher order effect due to the breaking of hexagonal symmetry as $`\beta `$ is increased.
### IV.3 Model C: Random lattice
In this model, the initial reference lattice is geometrically random. The rest bond length $`R_{bR}`$ is equal to the equilibrium lattice parameter $`R_{b0}`$ for every bond, so that the reference state is stress free. Our method of generating reference lattices of varying randomness is detailed below. Each bond is assigned the anharmonic potential of Eq. (88), where the spring constant $`k_b(R_b)=k_0/R_b`$ is a constant per unit length of the rest bond length.
We use the approach followed in Chung et al. (2002) to generate networks with a tunable degree of randomness. We begin by simulating a $`2`$-dimensional gas of $`40000`$ point particles interacting through a Lennard-Jones potential. The procedure outlined in Berendsen et al. (1984) is used to equilibrate the gas at a prescribed temperature and pressure, with periodic boundary conditions. The gas is equilibrated for $`10000`$ time steps, after which the particle configurations are sampled every $`1000`$ time steps. In this manner we obtain $`40`$ uncorrelated configurations of the gas at thirteen different temperature-pressure combinations, with $`T=8.0`$ and $`P=0.025`$, $`0.05`$, $`0.1`$, $`0.2`$, $`0.25`$, $`0.3`$, $`0.35`$, $`0.4`$, $`0.5`$ $`0.6`$, $`0.7`$, $`0.8`$, and $`1.0`$, all in units of the Lennard-Jones potential.
We use the particle positions from the snapshots of the equilibrated gas as the positions of nodes in our random lattice. Each sampled configuration is rescaled to have a box length of $`1`$ on each side. The point configurations are then tesselated using the Delaunay triangulation, which places a bond between each node and its nearest neighbors. The Delaunay triangulation produces networks with an average of $`6`$ bonds per node. A resulting lattice is pictured in Fig. 13.
The randomness in local elastic moduli as calculated from Eq. (17) is proportional to the distribution of bond lengths and bonds per node. In principle, as we take the equilibrium gas pressure to zero, the distribution of bond lengths will become completely random. Conversely, as we increase the pressure past a critical point the simulated gas begins to crystalize, forming spatial domains of hexagonal order separated by grain boundaries. This transition should be marked by a growth in the two-point correlation of local shear moduli. We fit the non-affinity correlation data for a broad range of pressures which cross this transition and compare it to the framework developed in previous sections. We used Eq. (17) to calculate $`\mathrm{\Delta }^K(𝐱=0)/\mu ^2`$ for each ensemble of random lattices, while the crystalline correlation length is fit as an unknown.
The lattice is sheared by $`0.1\%`$ and the energy is minimized as a function of node position as before. Figure 14 shows the displacement correlation function $`𝒢(𝐱)`$ as a function of separation $`\left|𝐱\right|`$ for lattices with different degrees of randomness. This correlation function shows the same logarithmic growth at large $`\left|𝐱\right|`$ as it does in the random spring constant lattices from the last section.
We fit the measurements of $`𝒢(𝐱)`$ to the functional form $`A\mathrm{ln}\left(\left|𝐱\right|/B\right)`$ at large $`\left|𝐱\right|`$. This data is shown in Figure 15. For the very random lattices generated at low Lennard-Jones pressure ($`T=8.0`$, $`P<0.3`$) the values of $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ and $`B`$ are nearly constant, as our framework predicts for the simple case of delta-function spatial correlations. However, lattices created at higher pressure values ($`T=8.0`$, $`P0.3`$) showed significant growth of both $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ and $`B`$ with increasing pressure, reaching a saturation point at around $`P=8.0`$. Visual inspection of the lattices in question revealed subdomains of hexagonal crystalline ordering. Long range correlations in the connectivity implies long-range correlations in the elastic moduli, so we must apply the framework developed in Sec. III.5 and App. C.1 in order to fit the data for partially crystalline lattices. Once again, we try the functional form in Eq. (84) for the spatial correlations in the elastic modulus. The numerical value of the factor $`g_0(\varphi )`$ can be calculated from the measured modulus autocorrelation $`\mathrm{\Delta }^K(𝐱=0)`$ using Eq. (89).
The fitting line in Fig. 15 represents a best fit of both the correlation exponent $`\varphi `$ and the cutoff length $`\mathrm{\Lambda }^1`$ to the form
$`A`$ $`{\displaystyle \frac{\mathrm{\Delta }^K(𝐱=0)\gamma ^2}{\mu ^2}}\xi ^2\left(1\left(1+\left(\mathrm{\Lambda }\xi \right)^2\right)^{1\frac{\varphi }{2}}\right)^1`$
$`B`$ $`{\displaystyle \frac{\xi }{\beta (\xi \mathrm{\Lambda },\varphi )}},`$ (92)
as suggested by Eq. (111). The best fit was achieved for $`\varphi =0.4`$ and a cutoff length of $`\mathrm{\Lambda }^11.25`$ lattice spacings. The corresponding analytic form of $`𝒢(𝐱)`$ calculated from Eq. (80) using $`g(u)`$ from Eq. (84) is shown by the solid line in Fig. 14.
To test the predicted \[Eq. (III.6)\] $`\mathrm{cos}4\psi `$ anisotropy in vorticity correlations, we measured $`G_\omega (𝐱)`$ as a function of the angle $`𝐱`$ makes with the $`x`$-axis. Figure 16 shows a polar plot of $`G_\omega (𝐱)`$, which clearly shows $`\mathrm{cos}4\psi `$ behavior, and the dependence of the $`\mathrm{cos}4\psi `$ term on $`|𝐱|`$, which shows the expected $`|𝐱|^2`$ behavior.
### IV.4 Model D: Random lattice with internal stresses
Finally, we simulate the most general model for random lattices, in which the rest bond length $`R_{bR}`$ is not equal to the equilibrium lattice parameter $`R_{b0}`$, and the lattice parameters $`R_{b0}`$ along with the number of bonds per node are random to within some finite distribution. Each bond is assigned the anharmonic potential of Eq. (88), where the spring constant $`k_b(R_b)=k_0/R_b`$ is a constant per unit length of the rest bond length. We used the same geometrically random lattices from Section IV.3 as staring points, then we add bond length frustration using the technique from Section IV.2: We multiply the rest lengths of all bonds by a factor $`(1+\beta _b)`$ where $`\beta _b`$ is chosen randomly from a flat distribution lying between $`\beta `$ and $`\beta `$ with $`\beta <0.1`$. We find the equilibrium configuration of the lattice by minimizing the elastic energy over node positions and box size. We then shear the lattice by $`0.1\%`$, minimize the energy over node positions, and measure the non-affinity correlation function $`𝒢(𝐱)`$.
In all these simulations, the correlation function $`𝒢(𝐱)`$ was well fit by the functional form $`A\mathrm{ln}\left(\left|𝐱\right|/B\right)`$. Figure 17 shows a plot of $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ for all data sets as a function of $`\beta `$. The data points for $`\beta =0`$ correspond to the data from Section IV.3; their deviation from the expected constancy of $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ was explained in that section by the growth of a correlation length scale as the system acquires partial hexagonal crystalline ordering. Here we see that as $`\beta `$ is increased, the long length scale ordering is disrupted by the additional randomness, and the ratio $`A/(\gamma ^2\mathrm{\Delta }^K(𝐱=0)/\mu ^2)`$ decreases toward the value for completely disordered lattices.
## V Summary and Conclusions
Nonaffine distortions are always present in random elastic networks subjected to external stress. In this paper, using both analytical and numerical techniques, we study properties of nonaffinity in these systems manifested in correlation functions of the deviation, $`𝐮^{}(𝐱)`$, of local displacements from their affine form. We introduce four models of random elastic networks with random local elastic moduli and possibly local random stress arising either from randomness in the form of the central force potentials between nearest neighbor sites or from random connectivity of the the network. In all cases, we show analytically and verify with numerical simulations that random elastic modulus times imposed strain and not random stress act as sources for nonaffine distortions. We calculate the nonaffinity displacement correlation function, $`𝒢(𝐱)=[𝐮^{}(𝐱)𝐮^{}(0)]^2`$, and the vorticity correlation function, $`G_\omega (𝐱)=\omega (𝐱)\omega (0)`$ analytically and verify their form in numerical simulations for systems with both short- and long-range correlations in local elastic moduli. We show in particular that $`𝒢(𝐱)\gamma ^2((\delta K)^2/K^2)\mathrm{ln}|𝐱|`$ at large $`𝐱`$ in two dimensions, where $`\gamma `$ is the imposed strain, $`K`$ is the average of elastic modulus, and $`(\delta K)^2`$ is it variance.
The formalism we develop is general and should be applicable to any elastic system that has a well defined average shear modulus. It should provide a basis for studying nonaffinity in granular media, foams, networks of semi-flexible polymers, and related systems. It should, in particular, provide a method of calculating correlation lengths near percolation-like thresholds such as the $`J`$-point in jammed systems or the rigidity percolation point. We have begun Vernon et al. (2005) to use these techniques to calculate correlation lengths in the former systems which we will eventually compared with those calculated from the density of states Wyart et al. (2005); Silbert et al. (2005) and to study nonaffinity in networks of semi-flexible polymers Didonna and Levine (2005).
###### Acknowledgements.
We are grateful to Peter Sollich and Dan Vernon for careful readings of the manuscript and their resultant useful suggestions and identification of misprints. BD gratefully acknowledges helpful discussions with Eric van der Giessen, Mitchell Luskin, Fred Mackintosh and Michael Rubinstein. This work was supported in part by the National Science Foundation under DMR 04-04670 (TCL), the National Institutes of Health under grant R01 GM056707 (BD and TCL), and the Institute for Mathematics and its Applications with funds provided by the National Science Foundation.
## Appendix A Properties of the modulus correlator
The modulus correlator $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K(𝐪)`$ is an 8th rank tensor. The number of its independent components depends on the symmetry of the reference space. In this appendix, we will determine the number and form of its independent components at $`𝐪=0`$ (strictly speaking $`𝐪0`$), or, equivalently, at all $`𝐪`$ when correlations are short range and it is independent of $`𝐪`$, when the reference space is isotropic. In this case, the general form of $`\mathrm{\Delta }_{ijkl;i^{}j^{}kj^{}l^{}}^K\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K(𝐪=0)`$ must be constructed from products of Kroneker $`\delta `$’s that distinctly pair all indices while respecting all symmetries.
It is useful to recall how this process is carried out for the simpler case of the 4th-rank elastic-modulus tensor $`K_{ijkl}`$, which is symmetric under interchange of $`i`$ and $`j`$, of $`k`$ and $`l`$, and of the pairs $`ij`$ and $`kl`$. Since any index can be paired with any of the remaining three and there is only one way to pair the remaining two, there are three distinct Kroneker-$`\delta `$ pairings, which we will call contractions, of the four indices: $`\delta _{ij}\delta _{kl}`$, $`\delta _{ik}\delta _{jl}`$, and $`\delta _{il}\delta _{jk}`$. The first of these satisfies all of the symmetry constraints, but the second two do not; their sum, however, does. The elastic-modulus tensor, therefore, has two independent components in an isotropic medium: $`K_{ijkl}=\lambda \delta _{ij}\delta _{kl}+\mu (\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})`$.
$`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K`$ is symmetric under interchange of $`i`$ and $`j`$, $`k`$ and $`l`$, $`i^{}`$ and $`j^{}`$, and $`k^{}`$ and $`l^{}`$; under the interchange of the pairs $`ij`$ and $`kl`$ and of the pairs $`i^{}j^{}`$ and $`k^{}l^{}`$; and under the interchange of the four-plets $`ijkl`$ and $`i^{}j^{}k^{}l^{}`$. The total number of possible contractions of these 8 indices is $`N_T=7\times 5\times 3\times 1=105`$ because any index can be contracted with any of the seven remaining indices, any one of the six remaining indices can then be contracted with any of the other five remaining, etc. Most of the individual realizations of these 105 possible contractions will not satisfy symmetry constraints; it is necessary to find the linear combinations of them that do. Figure 18 provides a graphical representation of the eight distinct contraction groups the sum over whose elements satisfy all constraints. The elastic-modulus correlation function in an isotropic medium can thus be written as
$`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K={\displaystyle \underset{\alpha }{}}\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^{K\alpha }`$
$`=\mathrm{\Delta }_1\delta _{ij}\delta _{kl}\delta _{i^{}j^{}}\delta _{k^{}l^{}}`$
$`+\mathrm{\Delta }_2\delta _{ij}\delta _{kl}(\delta _{i^{}k^{}}\delta _{j^{}l^{}}+\delta _{i^{}l^{}}\delta _{j^{}k^{}})+\text{prime}\text{unprime}`$
$`+\mathrm{\Delta }_3(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})(\delta _{i^{}k^{}}\delta _{j^{}l^{}}+\delta _{i^{}l^{}}\delta _{j^{}l^{}})`$
$`+\mathrm{\Delta }_4\delta _{ij}\delta _{i^{}j^{}}\delta _{kk^{}}\delta _{ll^{}}+\text{7 perm.}`$
$`+\mathrm{\Delta }_5\delta _{ij}\delta _{i^{}k^{}}\delta _{kj^{}}\delta _{ll^{}}+\text{31 perm.}`$
$`+\mathrm{\Delta }_6\delta _{ik}\delta _{i^{}k^{}}\delta _{jj^{}}\delta _{ll^{}}+\text{31 perm.}`$
$`+\mathrm{\Delta }_7\delta _{ii^{}}\delta _{jj^{}}\delta _{kk^{}}\delta _{ll^{}}+\text{7 perm.}`$
$`+\mathrm{\Delta }_8\delta _{ii^{}}\delta _{jk^{}}\delta _{kj^{}}\delta _{ll^{}}+\text{15 perm.}`$ (93)
where $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^{K1}=\mathrm{\Delta }_1\delta _{ij}\delta _{kl}\delta _{i^{}j^{}}\delta _{k^{}l^{}}`$, etc. The first three terms in $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^K`$ describe correlations in the isotropic Lamé coefficients: $`\mathrm{\Delta }_1=(\delta \lambda )^2`$, $`\mathrm{\Delta }_2=\delta \lambda \delta \mu `$, and $`\mathrm{\Delta }_3=(\delta \mu )^2`$. The other terms represent fluctuations away from local isotropy.
## Appendix B Evaluation of $`G(𝐪)`$
We outline here the calculation of $`G(𝐪)`$ to lowest order in $`\mathrm{\Delta }^K`$ in isotropic systems. We use Eq. (49) for $`G_{ij}(𝐪)`$ and sum over $`i=j`$. Using Eq. (51), we find
$$\chi _{ip}^0(𝐪)\chi _{ip^{}}^0(𝐪)=\frac{1}{\mu ^2q^4}\delta _{pp^{}}^T+\frac{1}{(\lambda +2\mu )^2q^4}\widehat{q}_p\widehat{q}_p^{},$$
(94)
where $`\widehat{q}_p=q_p/q`$ and $`\delta _{pp^{}}^T=\delta _{pp^{}}\widehat{q}_p\widehat{q}_p^{}`$. Then
$$G(𝐪)=\gamma _{xy}^2\underset{\alpha =1}{\overset{8}{}}\left(\frac{1}{\mu ^2q^2}S_\alpha ^T+\frac{1}{(\lambda +2\mu )^2q^2}S_\alpha ^L\right),$$
(95)
where
$$S_\alpha ^T=\delta _{pp^{}}^TS_{\alpha pp^{}}S_\alpha ^L=\widehat{q}_p\widehat{q}_p^{}S_{\alpha pp^{}},$$
(96)
with
$$S_{\alpha pp^{}}=\mathrm{\Delta }_{pjxy;p^{}j^{}xy}^{K\alpha }\widehat{q}_j\widehat{q}_j^{}$$
(97)
where $`\mathrm{\Delta }_{ijkl;i^{}j^{}k^{}l^{}}^{K\alpha }`$ is defined in Eq. (93). It is straightforward but tedious to calculate $`S_\alpha ^T`$ and $`S_\alpha ^L`$ from Eq. (93). The results are
$$\begin{array}{cc}S_1^T=S_2^T=0,\hfill & S_1^L=S_2^L=0;\hfill \\ S_3^T=\mathrm{\Delta }_3^T(\widehat{q}_{}^24\widehat{q}_x^2\widehat{q}_y^2),\hfill & S_3^L=\mathrm{\Delta }_3^L\widehat{q}_x^2\widehat{q}_y^2;\hfill \\ S_5^T=0,\hfill & S_5^L=\mathrm{\Delta }_5^L4\widehat{q}_{}^2\hfill \\ S_6^T=\mathrm{\Delta }_6^T(2+d\widehat{q}_{}^216\widehat{q}_x^2),\hfill & S_6^L=\mathrm{\Delta }_6^L(4\widehat{q}_{}^2+16\widehat{q}_x^2\widehat{q}_y^2);\hfill \\ S_7^T=\mathrm{\Delta }_7^T[(d1)+\widehat{q}_{}^24\widehat{q}_x^2\widehat{q}_y^2];\hfill & S_7^L=\mathrm{\Delta }_7^L(2+4\widehat{q}_x^2\widehat{q}_y^2)\hfill \\ S_8^T=\mathrm{\Delta }_8^T[2+(d2)\widehat{q}_{}^2],\hfill & S_8^L=\mathrm{\Delta }_8^L4\widehat{q}_{}^2.\hfill \end{array}$$
(98)
## Appendix C Evaluation of $`𝒢(𝐱)`$
In this appendix, we will evaluate the integrals \[Eq. (79)\]
$$_\alpha =2\xi ^\varphi \frac{d^dq}{(2\pi )^d}\frac{1}{q^2}f_\alpha (𝐪)g_\alpha (q\xi )\left(1e^{i𝐪𝐱}\right)$$
(99)
in $`2D`$ and $`3D`$ that make up the function $`𝒢(𝐱)`$, where $`f_A(𝐪)=1`$, $`f_B(𝐪)=\widehat{q}_{}^2`$, and $`f_C(𝐪)=\widehat{q}_x^2\widehat{q}_y^2`$.
### C.1 Two dimensions
In two dimensions, $`f_B(𝐪)=1=f_A(𝐪)`$, and
$$f_C(𝐪)=\mathrm{sin}^2\varphi _q\mathrm{cos}^2\varphi _q=\frac{1}{8}(1\mathrm{cos}4\varphi _q),$$
(100)
where $`\varphi _q`$ is the angle between $`𝐪`$ and the $`x`$-axis. Using the plane-wave decomposition relation
$$e^{i𝐪𝐱}=J_0(q|𝐱|)+2\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{cos}n\mathrm{\Theta }J_n(q|𝐱|),$$
(101)
where $`J_n(x)`$ is the $`n`$th order Bessel function, $`\mathrm{\Theta }=\varphi _q\psi `$, and $`\psi `$ is the angle between $`𝐱`$ and the $`x`$-axis, and the orthogonality relation
$$\frac{1}{2\pi }_0^{2\pi }𝑑\varphi \mathrm{cos}n\psi \mathrm{cos}m\mathrm{\Theta }=\{\begin{array}{cc}\frac{1}{2}\delta _{nm}\mathrm{cos}n\psi \hfill & n0\hfill \\ \delta _{nm}\hfill & n=0\hfill \end{array},$$
(102)
we find
$$_\alpha F_I[g_\alpha ]=\frac{\xi ^\varphi }{\pi }_0^\mathrm{\Lambda }\frac{dq}{q}g_\alpha (q\xi )[1J_0(q|𝐱|)]$$
(103)
for $`\alpha =A,B`$ and
$$_C=_I[g_C]\mathrm{cos}4\psi F_A[g_C],$$
(104)
where
$$F_A[g_\alpha ]=\frac{\xi }{\pi }_0^\mathrm{\Lambda }\frac{dq}{q}g_\alpha (q\xi )J_4(q|𝐱|).$$
(105)
Thus,
$$(𝐱)=F_I[g]+\frac{1}{8}\mathrm{cos}4\psi F_A[g_C],$$
(106)
where
$$g(q\xi )=g_A(q\xi )+g_B(q\xi )\frac{1}{8}g_C(q\xi ).$$
(107)
We now evaluate the integrals $`F_I`$ and $`F_A`$ in the limits $`|𝐱|\xi \mathrm{\Lambda }^1`$ and $`\mathrm{\Lambda }^1|𝐱|\xi `$.
#### C.1.1 $`|𝐱|\xi \mathrm{\Lambda }^1`$ in Two Dimensions
To evaluate the first limit of $`F_I`$, we set $`y=q|𝐱|`$ in Eq. (103):
$`F_I[g]={\displaystyle \frac{\xi ^\varphi }{\pi }}\{{\displaystyle _1^{\mathrm{\Lambda }|𝐱|}}{\displaystyle \frac{dy}{y}}g(y\xi /|𝐱|)`$
$`+{\displaystyle _0^1}{\displaystyle \frac{dy}{y}}g(y\xi /|𝐱|)[1J_0(y)]{\displaystyle _1^{\mathrm{\Lambda }|𝐱|}}{\displaystyle \frac{dy}{y}}g(y\xi /|𝐱|)J_0(y)\}.`$ (108)
In the limit $`|𝐱|/\xi \mathrm{}`$, we can safely replace $`g(y\xi /|𝐱|)`$ by $`g_0`$ in the second and third integrals in this expression, and we can let $`\mathrm{\Lambda }|𝐱|\mathrm{}`$ in the third integral. The first integral diverges as $`\mathrm{ln}|𝐱|`$ if we replace $`g(y\xi /|𝐱|)`$ by $`g_0`$ in it, and we have to be more careful to extract the constant term beyond the log:
$`{\displaystyle _1^{\mathrm{\Lambda }|𝐱|}}`$ $`{\displaystyle \frac{dy}{y}}g(y\xi /|𝐱|)={\displaystyle _{\xi /|𝐱|}^1}{\displaystyle \frac{du}{u}}g(u)+{\displaystyle _1^{\mathrm{\Lambda }\xi }}{\displaystyle \frac{du}{u}}g(u)`$ (109)
$`g_0\mathrm{ln}|𝐱|/\xi +{\displaystyle _0^1}{\displaystyle \frac{du}{u}}[g(u)g_0]+{\displaystyle _1^{\mathrm{\Lambda }\xi }}{\displaystyle \frac{du}{u}}g(u).`$ (110)
Thus in the limit $`|𝐱|\xi \mathrm{\Lambda }^1`$,
$$F_I[g]=\frac{\xi ^\varphi }{\pi }g_0\mathrm{ln}\beta (\mathrm{\Lambda }\xi ,\varphi )\frac{|𝐱|}{\xi },$$
(111)
where
$$\mathrm{ln}\beta (\mathrm{\Lambda }\xi ,\varphi )=\mathrm{ln}\alpha +_0^1\frac{dy}{y}\left[\frac{g(y)}{g_0}1\right]+_1^{\mathrm{\Lambda }\xi }\frac{dy}{y}\frac{g(y)}{g_0},$$
(112)
where
$$\mathrm{ln}\alpha =_0^1\frac{dy}{y}[1J_0(y)]_1^{\mathrm{}}\frac{dy}{y}J_0(y)=0.1159$$
(113)
and $`\alpha =0.8905`$. When $`g(y)=g_0`$, independent of $`y`$,
$$\mathrm{ln}\beta (\mathrm{\Lambda }\xi )\mathrm{ln}\alpha +\mathrm{ln}\mathrm{\Lambda }\xi ,$$
(114)
and Eq. (108) reduces to Eq. (44b) when $`g_0\xi ^\varphi `$ is identified with $`\mathrm{\Delta }^S`$.
The $`|𝐱|\xi `$ limit of $`F_A`$ is obtained by setting $`y=q|𝐱|`$ and noting that letting the upper limit of the integral go to infinity and replacing $`g(y|𝐱|/\xi )`$ by $`g_0`$ introduces no singularities. The result is
$$F_A[g]=\frac{\xi }{\pi }g_0_0^{\mathrm{}}\frac{dy}{y}J_4(y)=\frac{g_0\xi }{4}.$$
(115)
#### C.1.2 $`\mathrm{\Lambda }^1|𝐱|\xi `$ in Two Dimensions
To evaluate integrals when $`\mathrm{\Lambda }^1|𝐱|\xi `$, we introduce a new function
$$h(u)=u^\varphi \frac{g(u)}{g_{\mathrm{}}}\{\begin{array}{cc}1\hfill & \text{as }u\mathrm{},\hfill \\ (g_0/g_{\mathrm{}})u^\varphi \hfill & \text{as }u0,\hfill \end{array}$$
(116)
where $`g_{\mathrm{}}`$ is defined in Eq. (78) Then $`_I[g]=g_{\mathrm{}}(\xi ^\varphi /\pi )(|𝐱|/\xi )^\varphi (𝐱)`$, where
$$=_0^{\mathrm{\Lambda }|𝐱|}\frac{dy}{y}y^\varphi h(y\xi /|𝐱|)[1J_0(y)].$$
(117)
This integral has a potential infrared divergence as $`\xi /|𝐱|\mathrm{}`$ when $`\varphi 2`$. To isolate it, we break up the integral from $`0`$ to $`\mathrm{\Lambda }|𝐱|`$ into one from $`0`$ to $`1`$ and another from $`1`$ to $`\mathrm{\Lambda }|𝐱|`$. There are no troubles with ultraviolet divergences in the second integral, and in it, we can let $`\mathrm{\Lambda }|𝐱|\mathrm{}`$ and replace $`h(y\xi /|𝐱|)`$ by its infinite argument limit of one. In the integral from $`0`$ to $`1`$, we extract the small $`y`$ behavior of $`1J_0(y)`$ via $`1J_0(y)=(y^2/4)+[1J_0(y)(y^2/4)]`$. The second part of this expression vanishes as $`y^4`$ at small $`y`$, and there is no infrared divergence in the integral involving it so long as $`\varphi <4`$. Thus, we have
$`=_1`$ $`+{\displaystyle _0^1}{\displaystyle \frac{dy}{y}}y^\varphi [1J_0(y)(y^2/4)]`$ (118)
$`+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dy}{y}}y^\varphi [1J_0(y)],`$ (119)
where
$`_1`$ $`={\displaystyle \frac{1}{4}}{\displaystyle _0^1}𝑑yy^{1\varphi }h(y\xi /|𝐱|),`$ (120)
$`={\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\xi }{|𝐱|}}\right)^{\varphi 2}{\displaystyle _0^{\xi /|𝐱|}}𝑑uu^{1\varphi }h(u)`$ (121)
$`={\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\xi }{|𝐱|}}\right)^{\varphi 2}\{{\displaystyle _0^1}duu^{1\varphi }h(u)`$ (122)
$`+{\displaystyle _1^{\xi /|𝐱|}}dyu^{1\varphi }+{\displaystyle _1^{\xi /|𝐱|}}u^{1\varphi }[h(u)1]\}.`$ (123)
Using $`_1^\eta u^{1\varphi }=[\eta ^{2\varphi }1]/(2\varphi )`$, we arrive at Eq. (82) with
$`𝒜_2(\varphi )`$ $`={\displaystyle _0^{\mathrm{}}}𝑑yy^{(1+\varphi )}[1J_0(y)]`$ (124)
$`\mathrm{ln}\nu `$ $`={\displaystyle _0^1}{\displaystyle \frac{du}{u}}h(u)+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{du}{u}}[h(u)1]`$ (125)
$`𝒞_2(\varphi )`$ $`={\displaystyle _0^{\mathrm{}}}𝑑uu^{1\varphi }h(u).`$ (126)
The evaluation of the $`\xi |𝐱|`$ limit of $`F_A[g]`$ is straightforward. The result is
$$F_A=g_{\mathrm{}}\frac{|𝐱|^\varphi }{\pi }_0^{\mathrm{}}\frac{dy}{y}y^\varphi J_4(y).$$
(127)
### C.2 Three Dimensions
To evaluate the integrals $`_\alpha `$ in $`3D`$, we make use of the $`3D`$ plane-wave decomposition:
$$e^{i𝐪𝐱}=4\pi \underset{l=0}{\overset{\mathrm{}}{}}i^lj_l(q|𝐱|)\underset{m=l}{\overset{l}{}}Y_{lm}(\mathrm{\Omega }_q)Y_{lm}^{}(\mathrm{\Omega }_x),$$
(128)
where $`\mathrm{\Omega }_q=(\theta _q,\varphi _q)`$ and $`\mathrm{\Omega }_x=(\theta _x,\varphi _x)`$ are, respectively, the polar angles of $`𝐪`$ and $`𝐱`$, $`Y_{lm}(\mathrm{\Omega })`$ are spherical harmonics, and $`j_n(u)`$ is the $`n`$th order spherical Bessel function. Then, noting that
$`\widehat{q}_{}^2`$ $`=\mathrm{sin}^2\theta _q={\displaystyle \frac{2}{3}}\left[1P_2(\mathrm{cos}\theta _q)\right],`$ (129)
$`\widehat{q}_x^2\widehat{q}_y^2`$ $`={\displaystyle \frac{1}{8}}\mathrm{sin}^2\theta _q(1\mathrm{cos}4\varphi _q)`$ (130)
$`={\displaystyle \frac{1}{105}}[710P_2(\mathrm{cos}\theta _q)+3P_4(\mathrm{cos}\theta _q)]`$ (131)
$`{\displaystyle \frac{1}{16}}\sqrt{{\displaystyle \frac{4\pi }{9!}}}2^44![Y_{44}(\mathrm{\Omega }_q)+Y_{4,4}(\mathrm{\Omega }_q)],`$ (132)
where $`P_n(x)`$ is the $`n`$th-order Legendre Polynomial, we find
$`_A`$ $`={\displaystyle \frac{\xi ^\varphi }{\pi }}I_0[g_A]`$ (133)
$`_B`$ $`={\displaystyle \frac{\xi ^\varphi }{\pi }}\left\{{\displaystyle \frac{2}{3}}I_2[g_B]+{\displaystyle \frac{2}{3}}P_2(\mathrm{cos}\theta _x)I_2[g_B]\right\}`$ (134)
$`_C`$ $`={\displaystyle \frac{\xi ^\varphi }{\pi }}\{{\displaystyle \frac{1}{15}}I_0[g_C]{\displaystyle \frac{2}{21}}P_2(\mathrm{cos}\theta _x)I_2[g_C]`$ (135)
$`\left[{\displaystyle \frac{1}{35}}P_2(\mathrm{cos}\theta _x)+{\displaystyle \frac{1}{8}}\mathrm{sin}^2\theta _x\mathrm{cos}4\varphi _x\right]I_4[g_C],`$ (136)
where
$`I_0[g]`$ $`={\displaystyle _0^\mathrm{\Lambda }}𝑑qg(q\xi )[1j_0(q|𝐱|)]`$ (137)
$`I_2[g]`$ $`={\displaystyle _0^\mathrm{\Lambda }}𝑑qg(q\xi )j_2(q|𝐱|)`$ (138)
$`I_4[g]`$ $`={\displaystyle _0^\mathrm{\Lambda }}𝑑qg(q\xi )j_4(q|𝐱|).`$ (139)
Thus,
$``$ $`=_A+_B_C`$ (140)
$`={\displaystyle \frac{\xi ^\varphi }{\pi }}\{I_1[g_1]+P_2(\mathrm{cos}\theta _x)I_2[g_2]`$ (141)
$`\left[{\displaystyle \frac{1}{35}}P_2(\mathrm{cos}\theta _x)+{\displaystyle \frac{1}{8}}\mathrm{sin}^2\theta _x\mathrm{cos}4\varphi _x\right]I_4[g_C],`$ (142)
where $`g_1=g_A+\frac{2}{3}g_B\frac{1}{15}g_C`$ and $`g_2=\frac{2}{3}g_B+\frac{2}{21}g_C`$. Thus we need only evaluate the three integrals $`I_1`$, $`I_2`$, and $`I_3`$.
#### C.2.1 $`|𝐱|\xi >\mathrm{\Lambda }^1`$ in Three Dimensions
In this limit, in integrals with integrands proportional to $`j_n(q|𝐱|)`$, we set $`y=q|𝐱|`$, replace $`g(y\xi /|𝐱|)`$ by $`g_0`$ and replace the upper limit, $`\mathrm{\Lambda }|𝐱|`$, of integration by $`\mathrm{}`$. In the part of the integral $`I_1`$ not proportional to $`j_0(q|𝐱|)`$, we set $`y=q\xi `$. The result is
$`I_1`$ $`g_0{\displaystyle \frac{\pi }{2}}\left[{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{\Lambda }\xi }}{\displaystyle \frac{g(y)}{g_0}}𝑑y{\displaystyle \frac{1}{|𝐱|}}\right]`$ (143)
$`I_2`$ $`{\displaystyle \frac{g_0}{|𝐱|}}{\displaystyle _0^{\mathrm{}}}𝑑yj_2(y)={\displaystyle \frac{g_0\pi }{4|𝐱|}}`$ (144)
$`I_3`$ $`{\displaystyle \frac{g_0}{|𝐱|}}{\displaystyle _0^{\mathrm{}}}j_4(y)={\displaystyle \frac{3g_0\pi }{16|𝐱|}}.`$ (145)
#### C.2.2 $`\mathrm{\Lambda }^1|𝐱|\xi `$ in Three Dimensions
To treat this limit, as in $`2D`$, we use the function $`h(u)`$ \[Eq. ( 116)\]. To evaluate $`I_1`$, we break up the limits of integration in much the same way we did in $`2D`$. The result is
$`I_1`$ $`=g_{\mathrm{}}\xi ^\varphi |𝐱|^{\varphi 1}`$ (146)
$`\times \{{\displaystyle _0^1}dyy^\varphi [1j_0(y)]{\displaystyle _1^{\mathrm{}}}dyy^\varphi j_0(y)`$ (147)
$`+{\displaystyle \frac{1}{1\varphi }}[(\mathrm{\Lambda }|𝐱|)^{1\varphi }1]`$ (148)
$`+\left({\displaystyle \frac{\xi }{|𝐱|}}\right)^{1\varphi }{\displaystyle _0^{\mathrm{\Lambda }\xi }}duu^\varphi [h(u)1]\}`$ (149)
for $`0<\varphi <3`$. The dominant behavior for $`1<\varphi <3`$ and $`0<\varphi <1`$ is then
$$I_1\{\begin{array}{cc}g_{\mathrm{}}\xi ^\varphi |𝐱|^{\varphi 1}𝒜_3(\varphi )\hfill & 1<\varphi <3\hfill \\ \frac{1}{\xi }_0^{\mathrm{\Lambda }\xi }𝑑ug(u)g_{\mathrm{}}\xi ^\varphi |𝐱|^{(1\varphi )}𝒞_3(\varphi )\hfill & 0<\varphi <1.\hfill \end{array},$$
(150)
where
$`𝒜_3(\varphi )`$ $`={\displaystyle _0^{\mathrm{}}}y^\varphi [1j_0(y)]`$ (151a)
$`𝒞_3(\varphi )`$ $`={\displaystyle _0^{\mathrm{}}}y^\varphi j_0(y).`$ (151b)
### C.3 One Dimension
In $`1D`$, there is only one function to evaluate
$`(x)`$ $`=2\xi ^\varphi {\displaystyle _\mathrm{\Lambda }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{2\pi }}{\displaystyle \frac{1}{q^2}}g(q\xi )[1\mathrm{cos}(q|x|)]`$
$`={\displaystyle \frac{2\xi ^\varphi }{\pi }}{\displaystyle _0^{\mathrm{\Lambda }|x|}}𝑑y{\displaystyle \frac{1}{y^2}}g(y\xi /|x|)(1\mathrm{cos}y).`$ (152)
The limit $`|x|\xi `$ is obtained as before by replacing $`g(y\xi /|x|)`$ with $`g_0`$ and letting $`\mathrm{\Lambda }|x|\mathrm{}`$:
$$(x)=\frac{2}{\pi }|x|g_0\xi ^\varphi _0^{\mathrm{}}𝑑y\frac{1\mathrm{cos}y}{y^2}=g_0\xi ^\varphi |x|.$$
(153)
In the limit $`\mathrm{\Lambda }^1|x|\xi `$, we introduce $`h(y)`$ as in $`2D`$ and $`3D`$: $`(x)=(2g_{\mathrm{}}|x|/\pi )𝒦`$, where
$`𝒦`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }|x|}}𝑑yh(y\xi /|x|)y^{(2+\varphi )}(1\mathrm{cos}y)`$
$`=`$ $`𝒦_1+{\displaystyle _0^1}𝑑yy^{(2+\varphi )}[1\mathrm{cos}y(y^2/2)]`$
$`+{\displaystyle _1^{\mathrm{}}}y^{(2+\varphi )}[1\mathrm{cos}y],`$ (154)
where
$`𝒦_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑yy^\varphi h(y\xi /|x|)`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\xi }{|x|}}\right)^{\varphi 1}\{{\displaystyle \frac{1}{1\varphi }}(\left({\displaystyle \frac{\xi }{|x|}}\right)^{1\varphi }1]`$
$`+{\displaystyle _0^2}duh(u)u^\varphi +{\displaystyle _1^{\mathrm{}}}duu^\varphi [h(u)1]\}.`$
Combining Eqs. (154) with (C.3), we find
$$(x)\{\begin{array}{cc}\frac{2}{\pi }g_{\mathrm{}}|x|^{1+\varphi }𝒜_1(\varphi )\hfill & \text{if }\varphi <1,\hfill \\ \frac{1}{\pi }\mathrm{ln}\nu \xi /|x|\hfill & \text{if }\varphi =1,\hfill \\ \frac{1}{\pi }\xi ^{\varphi 1}|x|^2𝒞_2(\varphi )\hfill & \text{if }\varphi >1,\hfill \end{array}$$
(156)
where $`\nu `$ is given by Eq. ( 125) and $`𝒞_2(\varphi )`$ is given by Eq. (126) and where
$$𝒜_1=_0^{\mathrm{}}\frac{1\mathrm{cos}y}{y^{2+\varphi }}.$$
(157)
The $`|𝐱|\xi `$ limits of both $`I_2`$ and $`I_3`$ can be obtained by simply by replacing $`g(q\xi )`$ by $`(q\xi )^\varphi g_{\mathrm{}}`$:
$`I_2`$ $`g_{\mathrm{}}\xi ^\varphi |𝐱|^{\varphi 1}{\displaystyle _0^{\mathrm{}}}𝑑yy^\varphi j_2(y)`$ $`0<\varphi <3`$ (158)
$`I_3`$ $`g_{\mathrm{}}\xi ^\varphi |𝐱|^{\varphi 1}{\displaystyle _0^{\mathrm{}}}𝑑yy^\varphi j_4(y)`$ $`0<\varphi <5.`$ (159)
## Appendix D Evaluation of $`𝒢_\omega (𝐱)`$
In this appendix we will evaluate the rotational correlation function $`G_\omega (𝐱)`$ in two dimensions. To lowest order in $`\mathrm{\Delta }^K`$,
$$G_\omega (𝐪)=\frac{1}{4}ϵ_{ri}ϵ_{r^{}i^{}}q_rq_r^{}\chi _{ip}^0(𝐪)\chi _{i^{}p^{}}^0(𝐪)\underset{\alpha }{}S_{\alpha pp^{}},$$
(160)
where $`S_{\alpha pp^{}}`$ is defined in Eq. (97). The product $`ϵ_{ri}ϵ_{r^{}i^{}}`$ is simply $`\delta _{rr^{}}\delta _{ii^{}}\delta _{ri^{}}\delta _{ir^{}}`$, and $`ϵ_{ri}ϵ_{r^{}i^{}}q_rq_r^{}=q^2\delta _{ii^{}}^T`$. When this operates on $`\chi _{ip}^0\chi _{i^{}p^{}}^0`$, it projects out the transverse part leaving $`\delta _{pp^{}}^T/(\mu ^2q^2)`$. Thus
$$G_\omega (𝐪)=\frac{\gamma _{xy}^2}{\mu ^2}\underset{\alpha }{}S_\alpha ^T=\frac{\gamma _{xy}^2}{\mu ^2}(\mathrm{\Delta }_A^\omega \mathrm{\Delta }_C^\omega \widehat{q}_x^2\widehat{q}_y^2).$$
(161)
where $`\mathrm{\Delta }_A^\omega =\mathrm{\Delta }_3+4\mathrm{\Delta }_6+2\mathrm{\Delta }_7+2\mathrm{\Delta }_8`$ and $`\mathrm{\Delta }_C^\omega =4\mathrm{\Delta }_3+15\mathrm{\Delta }_6+4\mathrm{\Delta }_7`$.
When $`\mathrm{\Delta }_{A,C}^\omega (𝐪)`$ have a Lorentizan form, we need to evaluate two integrals to determine $`G_\omega (𝐱)`$:
$$F_1(𝐱)=\frac{d^2q}{(2\pi )^2}\frac{1}{q^2+\xi ^2}e^{i𝐪𝐱}=\frac{1}{2\pi }K_0(|𝐱|/\xi )$$
(162)
and
$`F_2(𝐱)`$ $`=`$ $`{\displaystyle \frac{d^2q}{(2\pi )^2}\frac{\widehat{q}_x^2\widehat{q}_y^2}{q^2+\xi ^2}e^{i𝐪𝐱}}`$ (163)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{qdq}{q^2+\xi ^2}}{\displaystyle \frac{d\varphi }{2\pi }\mathrm{cos}^2\varphi \mathrm{sin}^2\varphi e^{iq|𝐱|\mathrm{cos}(\varphi \psi )}}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi }}[K_0(|𝐱|/\xi )`$
$`\mathrm{cos}4\psi ({\displaystyle \frac{48\xi ^4}{|𝐱|^4}}+{\displaystyle \frac{4\xi ^2}{|𝐱|^2}}+K_4(|𝐱|/\xi ))],`$
where $`𝐪=q(\mathrm{cos}\varphi ,\mathrm{sin}\varphi )`$, $`𝐱=|𝐱|(\mathrm{cos}\psi ,\mathrm{sin}\psi )`$ and $`K_n(y)`$ is the Bessel function of imaginary argument.
|
warning/0506/hep-lat0506029.html
|
ar5iv
|
text
|
# Group-theoretical construction of extended baryon operators in lattice QCD
## I Introduction
Spectroscopy is a powerful tool for uncovering the important degrees of freedom of a physical system and the interaction forces between them. The spectrum of quantum chromodynamics (QCD) is indeed very rich: conventional baryons (nucleons, $`\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Xi },\mathrm{\Omega },\mathrm{}`$) and mesons $`(\pi ,K,\rho ,\varphi ,\mathrm{})`$ have been known for nearly half a century, but other higher-lying exotic states, such as glueballs, four-quark states, and so-called hybrid mesons and hybrid baryons bound by an excited gluon field, have proved more elusive, mainly because our theoretical understanding of such states is insufficient, making their identification problematical.
Recent discoveries of several new hadronic resonances have generated much excitement in the field of hadron spectroscopy. The E852 collaborationAdams et al. (1998) reported a signal for a $`1^+`$ hybrid meson decaying into $`\rho \pi `$ with a mass near 1.6 GeV, though the significance of this result has been questionedDzierba et al. . Another exotic $`1^+`$ candidate at 1.4 GeV has been tentatively identified in the $`\eta \pi `$ channel by E852Chung et al. (1999), VESBeladidze et al. (1993), KEKAoyagi et al. (1993), and Crystal BarrelAmsler et al. (1994). The first observation of a doubly-charmed baryon has been reportedMattson et al. (2002), and evidence for the possible existence of a strangeness $`S=1`$ pentaquark state has emerged at SPring-8, JLab, and elsewhereStepanyan et al. (2003); Nakano et al. (2003); Barth et al. (2003); Barmin et al. (2003). Interest in excited baryon resonances has also been raised by experiments dedicated to mapping out the $`N^{}`$ spectrum in Hall B at the Thomas Jefferson National Accelerator Facility (JLab), and the search for hybrid mesons and glueballs is intensifying due to the Hall D initiative at JLab and experiments at CLEO-c.
Unfortunately, our understanding of many conventional and excited hadron resonances is vestigial and comes only from QCD-inspired phenomenological models, such as the bag model, the nonrelativistic quark model, and quark-diquark models, or from approaches, such as QCD sum rules and methods based on Schwinger-Dyson equations, which use approximations whose justifications are unclear. There are a growing number of resonances which cannot be easily accommodated within quark models. States bound by an excited gluon field, such as hybrid mesons and baryons, are still poorly understood. The natures of the Roper resonance and the $`\mathrm{\Lambda }(1405)`$ remain controversial. Experiment shows that the first excited positive-parity spin-$`1/2`$ baryon lies below the lowest-lying negative-parity spin-$`1/2`$ resonance, a fact which is difficult to reconcile in quark models. Given the current surge in experimental activity, the need for an understanding of such states from QCD itself has never been greater.
This need has motivated us to undertake a comprehensive ab initio study of the hadron spectrum in QCD. Presently, Monte Carlo estimates of QCD path integrals defined on a space-time lattice offer the best way to make progress in this regard, so this is the calculational approach we have adopted. The Monte Carlo approach has been used to investigate hadrons throughout the past two decades, but the number of states studied to date has been somewhat limited. Also, prior works have not strived to identify the continuum spin $`J`$ of the states studied, simply assuming that the lowest allowed $`J`$ would be the lowest-lying state. The novel features of our approach are its comprehensiveness and its use of techniques to identify spin. Given the vast amount of experimental data being generated at accelerator facilities, such as Jefferson Lab, there is an urgency to investigate the spectrum (masses, widths, transition rates, form factors, and so on) as completely as possible. The number of hadron eigenstates which can be reliably extracted in Monte Carlo computations is not currently known, so our undertaking will be partly an exploration of the limits of possibility. Another aim is to discover whether a very large number of interpolating operators are needed to extract the low-lying spectrum, or whether a handful of carefully chosen ones is sufficient, and this work outlines a systematic means of finding such operators.
Our first goal is to calculate the masses of as many low-lying hadron resonances as possible in QCD using Monte Carlo techniques. The masses and widths of resonances (unstable hadrons) cannot be calculated directly in finite-volume Monte Carlo computations, but must be deduced from the discrete spectrum of finite-volume stationary states for a range of box sizesDeWitt (1956); Wiese (1989); Lüscher (1991); Rummukainen and Gottlieb (1995). The rigorous application of such techniques to obtain the resonance parameters to high accuracy would require vast computational resources, but our goal here is merely to obtain a first exploratory scan of the spectrum of QCD, not to pin down each mass to very high precision. Hence, simply obtaining the finite-volume spectrum for a few judiciously-chosen volumes should suffice for ferreting out the hadron resonances from the less interesting scattering states and may even give qualitative information about preferred decay modes.
To compute the finite-volume stationary state energies, the temporal correlations $`C_{ij}(t)=0|TO_i(t)\overline{O}_j(0)|0`$, where $`T`$ denotes time-ordering, of a set of operators $`\overline{O}_j(0)`$ which create the states of interest at an initial time $`t=0`$ with a corresponding set of operators $`O_i(t)`$ which annihilate the states of interest at a later time $`t`$ must be determined. The correlation functions $`C_{ij}(t)`$ can be expressed in terms of path integrals over the quark and gluon fields, and when formulated on a Euclidean space-time lattice, such path integrals can be estimated using the Monte Carlo method with Markov-chain importance sampling. Incorporating the quark-field effects into the Monte Carlo updating for realistically light quark masses remains a challenge, but there is steady progress with improving algorithms and increasing computational power.
The procedure for extracting the lowest stationary-state energies $`E_0,E_1,E_2,\mathrm{}`$ from the hermitian matrix of correlation functions $`C_{ij}(t)`$ is well knownMichael (1985); Lüscher and Wolff (1990). Let $`\lambda _n(t,t_0)`$ denote the eigenvalues of the hermitian matrix $`C(t_0)^{1/2}C(t)C(t_0)^{1/2}`$, where $`t_0`$ is some fixed reference time (typically small) and the eigenvalues, also known as the principal correlation functions, are ordered such that $`\lambda _0\lambda _1\mathrm{}`$ as $`t`$ becomes large. Then one can show that
$`\underset{t\mathrm{}}{lim}\lambda _n(t,t_0)`$ $`=`$ $`e^{E_n(tt_0)}\left(1+O(e^{\mathrm{\Delta }_n(tt_0)})\right),`$ (1)
$`\mathrm{\Delta }_n`$ $`=`$ $`\underset{kn}{\mathrm{min}}|E_kE_n|.`$ (2)
Determinations of the principal correlators $`\lambda _n(t,t_0)`$ for sufficiently large temporal separations $`t`$ yield the desired energies $`E_n`$.
The above equations illustrate the difficulties which must be faced in order to extract the stationary state energies $`E_0,E_1,\mathrm{}`$ from the temporal correlations of the hadronic operators. In each principal correlator, there are contaminating contributions from all other states which can be created and annihilated by the operators used. In order to reliably extract the single decaying exponential of interest, the obscuring contributions from all of these other states must somehow be suppressed.
There are two crucial ingredients in reducing the unwanted contributions in $`\lambda _n(t,t_0)`$. The first is to use sufficiently large values of $`t`$. However, there are often practical considerations which limit how large $`t`$ can be, and the statistical uncertainties in the Monte Carlo estimates of the correlation functions generally increase with $`t`$. The second, and more important, consideration in suppressing the contamination in $`\lambda _n(t,t_0)`$ is to use cleverly-devised operators which couple minimally to the unwanted states.
Successfully extracting the spectrum of QCD in our Monte Carlo computations will hinge crucially on using carefully designed hadronic operators. Excited meson and baryon resonances are expected to be large objects, so the use of spatially-extended operators is important. Since our calculations will be carried out on hypercubic space-time lattices, the energies of states in all irreducible representations (irreps) of the $`O_h`$ cubic point group must be determined in order to identify the continuum-limit spin $`J`$ of each physical state. Because determining the mass of a particular resonance requires determining the energies of all lower-lying stationary states, including scattering states, the set of operators we use must include not only meson and baryon operators, but also multi-hadron operators. Another very important fact to keep in mind is the computational cost of evaluating quark propagators, especially for light quark masses. Hence, our operators must be devised with an eye towards minimizing the number of sources needed to calculate the required quark propagators.
Designing the hadronic operators is an important first step in our comprehensive study of the mass spectrum in lattice QCD. From the above considerations, the guiding principles in devising our operators are maximizing overlaps with the states of interest while minimizing the number of quark-propagator sources. Although operators for baryons, mesons, and their scattering states will be needed, we restrict our attention to the construction of the three-quark baryon operators in this first paper. Due to the complexity of these calculations and the importance of providing checks on our final results, we have been pursuing two different approaches to constructing the baryon interpolating field operators. Only one approach is described here; an alternative approach based on Clebsch-Gordan techniques will appear elsewhereBasak et al. . Meson and multi-hadron operators will be detailed in subsequent works. Furthermore, this paper deals only with issues related to the construction and utilization of these operators. Results from Monte Carlo calculations using these operators will be presented in later publications.
This paper is organized as follows. An overview of our approach to constructing the hadronic operators is first outlined in Sec. II. Our operators are assemblages of basic building blocks, which are described in Sec. III, along with the conventions we use for the Dirac $`\gamma `$-matrices and in Wick rotating into Euclidean space-time. We use quark fields which are Dirac spinors, so our hadron operators apply to Monte Carlo calculations involving Wilson, domain-wall, and overlap fermion actions, but not to computations involving staggered fermions. The basic building blocks are then combined into gauge-invariant three-quark operators (referred to as elemental operators) having appropriate flavor structure in Sec. IV. Projections onto the rotation-reflection symmetry sectors produce the final operators in Sec. V. Issues related to the use of these projections in constructing the baryon propagators are discussed in Sec. VI. Concluding remarks and plans for future work are outlined in Sec. VII.
## II Overview of operator construction
Devising relativistic hadronic operators in continuous space-time usually involves combining Dirac spinors to form Lorentz scalars, pseudoscalars, vectors, axial-vectors, and so on, using the Dirac $`\gamma `$-matrices and the charge conjugation matrix $`C`$ satisfying $`C\gamma _\mu C^1=\gamma _\mu ^T`$. It has been common practice in lattice QCD simulations to use hadron operators built in a similar fashion through simply discretizing their continuum analogs. However, such an approach becomes very cumbersome when constructing higher spin states or complicated extended operators. Also, the above operators generally couple to states belonging to different $`J^P`$ (spin-parity) symmetry sectors. Since the hypercubic lattice breaks Lorentz covariance, there is really no reason to construct operators according to Lorentz symmetries. Finally, we wish to extract a large portion of the low-lying spectrum, which means that large sets of operators will be needed to compute complete correlation matrices. Hence, the usual approach which mimics that used in continuous space-time is not feasible for our purposes.
Instead, we advocate an approach which more directly combines the physical characteristics of baryons with the symmetries of the lattice regularization of QCD. Recall that baryon states are characterized by their total momentum $`𝒑`$, their total (half-integral) spin $`J`$, a projection $`\lambda `$ of this spin onto some axis (the $`z`$-axis or the momentum, say), their parity $`P=\pm 1`$, and their quark flavor content. The masses of the light $`u`$ and $`d`$ quarks are very nearly equal, and therefore, we work in the approximation that $`m_u=m_d`$. In this approximation, the theory has an exact isotopic spin symmetry, and states carry two more labels, total isospin $`I`$ and its projection $`I_3`$ onto a given axis. The other flavor quantum numbers which label the states are strangeness $`S`$, charm $`C`$, and bottomness $`B`$ (we do not consider the top quark). In our calculations, isospin remains an exact symmetry since we neglect electromagnetic interactions.
Since our simulations will be performed using a hypercubic space-time lattice, our operators should be classified according to the symmetries of the lattice, rather than the full rotational symmetries of continuous space-time. Of course, we expect to recover the symmetries of the space-time continuum as the lattice spacing is made small. Since we are interested only in determining the masses of the baryon states, we restrict our attention to representations corresponding to zero total three-momentum $`𝒑=\mathrm{𝟎}`$. Hence, our operators must be invariant under all allowed spatial translations and we require that they transform under spatial rotation-reflection symmetry operations according to the irreducible representations of the octahedral point group $`O_h`$. These irreducible representations are the lattice analogs of the continuum $`J^P`$ labels, and the row of the representation is the analog of the spin projection $`\lambda `$. Thus, our (annihilation) operators can be written
$$B_i^{\mathrm{\Lambda }\lambda F}(t)=\underset{𝒙}{}B_i^{\mathrm{\Lambda }\lambda F}(𝒙,t),$$
(3)
where $`\mathrm{\Lambda }`$ indicates the irreducible representation of $`O_h`$, $`\lambda `$ is the row of the $`\mathrm{\Lambda }`$ representation, $`F`$ denotes all of the quantum numbers associated with the flavor content of the operator, and $`i`$ labels the different operators in the $`\mathrm{\Lambda }\lambda F`$ symmetry sector. Under a symmetry operation $`R`$, these operators transform according to
$`U_RB_i^{\mathrm{\Lambda }\lambda F}(t)U_R^{}`$ $`=`$ $`{\displaystyle \underset{\mu }{}}B_i^{\mathrm{\Lambda }\mu F}(t)\mathrm{\Gamma }_{\mu \lambda }^{(\mathrm{\Lambda })}(R)^{},`$ (4)
$`U_R\overline{B}_i^{\mathrm{\Lambda }\lambda F}(t)U_R^{}`$ $`=`$ $`{\displaystyle \underset{\mu }{}}\overline{B}_i^{\mathrm{\Lambda }\mu F}(t)\mathrm{\Gamma }_{\mu \lambda }^{(\mathrm{\Lambda })}(R),`$ (5)
where $`U_R`$ denotes the quantum operator which effects the symmetry operation corresponding to group element $`R`$ (not to be confused with the gauge link variables), and $`\mathrm{\Gamma }_{\mu \lambda }^{(\mathrm{\Lambda })}(R)`$ are the elements of the $`\mathrm{\Lambda }`$ representation matrix corresponding to group element $`R`$. For baryons, we are only interested in states corresponding to half-integral spin $`J`$ in the continuum limit, so we can restrict our attention to the six spinorial representations of $`O_h`$. There are four two-dimensional irreducible representations $`G_{1g},G_{1u},G_{2g}`$, and $`G_{2u}`$ (adopting a Mulliken-like naming convention), and two four-dimensional representations $`H_g`$ and $`H_u`$. These representations will be discussed in greater detail later.
Our general approach to constructing the $`B_i^{\mathrm{\Lambda }\lambda F}(t)`$ operators is to (1) first identify appropriate basic “building blocks” to use in constructing all baryon operators, (2) devise simple elemental operators containing the appropriate flavor and color structure, then (3) apply appropriate group-theoretical projection operators to find linear combinations of the elemental operators with the desired transformation properties under the symmetry group of a spatial cubic lattice. Let $`B_i^F(t)=_𝒙B_i^F(𝒙,t)`$ denote a gauge-invariant elemental operator with the appropriate quark flavor content and which is invariant under allowed spatial translations, then an operator which transforms according to the row $`\lambda `$ of the $`\mathrm{\Lambda }`$ irreducible representation is obtained using
$$_i^{\mathrm{\Lambda }\lambda F}(t)=\frac{d_\mathrm{\Lambda }}{g_{O_h^D}}\underset{RO_h^D}{}\mathrm{\Gamma }_{\lambda \lambda }^{(\mathrm{\Lambda })}(R)U_RB_i^F(t)U_R^{},$$
(6)
where $`O_h^D`$ is the double group of $`O_h`$, $`R`$ denotes an element of $`O_h^D`$, $`g_{O_h^D}`$ is the number of elements in $`O_h^D`$, and $`d_\mathrm{\Lambda }`$ is the dimension of the $`\mathrm{\Lambda }`$ irreducible representation. Projections onto the double-valued irreps of $`O_h`$ require using the double group $`O_h^D`$ in Eq. (6). Given $`M_B`$ elemental $`B_i^F`$ operators, many of the projections in Eq. (6) vanish or lead to linearly-dependent operators, so one must then choose suitable linear combinations of the projected operators to obtain a final set of independent baryon operators. Thus, in each symmetry channel, one ends up with a set of $`r`$ operators given in terms of a linear superposition of the $`M_B`$ elemental operators:
$$B_i^{\mathrm{\Lambda }\lambda F}(t)=\underset{j=1}{\overset{M_B}{}}c_{ij}^{\mathrm{\Lambda }\lambda F}B_j^F(t),i=1\mathrm{}r.$$
(7)
Note that the expansion coefficients in the $`\overline{B}_i^{\mathrm{\Lambda }\lambda F}`$ operators are the complex conjugates of those in $`B_i^{\mathrm{\Lambda }\lambda F}`$:
$$\overline{B}_i^{\mathrm{\Lambda }\lambda F}(t)=\underset{j=1}{\overset{M_B}{}}c_{ij}^{\mathrm{\Lambda }\lambda F}\overline{B}_j^F(t),i=1\mathrm{}r.$$
(8)
## III The basic building blocks
The oscillating path integral weight $`e^{iS_M}`$ in quantum field theory, where $`S_M`$ is the action defined in Minkowski space-time and using natural units $`\mathrm{}=c=1`$, is unsuitable for applying the Monte Carlo method to evaluate the correlation functions of the theory via Feynman path integrals. A rotation to imaginary time $`ti\tau `$, where $`\tau `$ is real, leads to path integrals with weight $`e^S`$, where $`S`$ is the action defined in Euclidean space-time. If $`S`$ is real, the path integral weight is real and positive and, hence, can be interpreted as a probability, allowing the application of Monte Carlo methods with suitable importance sampling. The Euclidean action $`S`$ is defined such that $`S`$ is invariant under all symmetries of Euclidean space-time and all Green’s functions of the theory are identical to the Green’s functions of the Minkowski theory, analytically continued to imaginary time $`ti\tau `$. Although our simulations employ a path integral quantization of the field theory, a canonical quantization viewpoint can be adopted when discussing the quantum operators.
Our conventions for the continuation from Minkowski space-time with metric $`g_{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ into Euclidean space-time (imaginary time) are as follows. We define the following Euclidean space-time coordinates and derivatives (a subscript or superscript $`M`$ indicates a Minkowski space-time quantity):
$`x^4=x_4=ix_M^0,x^j=x_j=x_M^j=x_j^M,`$ (9)
$`^4=_4=i_0^M,^j=_j=_M^j=_j^M,`$ (10)
for spatial directions $`j=1,2,3`$. The metric in Euclidean space-time is $`\delta ^{\mu \nu }`$ so there is no distinction between covariant and contravariant indices. Our Monte Carlo calculations will be carried out using a hypercubic space-time lattice, and we require that the lattice spacings in the three spatial directions are the same, denoted by $`a_s`$; the temporal lattice spacing $`a_t`$ may differ from $`a_s`$, allowing us to exploit the known benefits of anisotropic latticesMorningstar and Peardon (1999). Throughout this paper, we set $`a_s=1`$ to simplify the notation. Four-vectors of unit length pointing along the spatial axes of the lattice with a vanishing temporal component will be denoted by $`\widehat{j}`$, $`\widehat{k}`$, and so on, for $`j,k=\pm 1,\pm 2,\pm 3`$.
As usual in lattice gauge theory, the gluon field is introduced using the parallel transporter $`U_\mu (x)`$ given by the path-ordered exponential of the gauge field along each link connecting neighboring sites of the hypercubic lattice. We also introduce the Dirac spinor field $`\psi _{a\alpha }^A(x)`$ which annihilates a quark and creates an antiquark, where $`A`$ refers to the quark flavor, $`a`$ refers to color, and $`\alpha `$ is the Dirac spin index, and the field $`\overline{\psi }_{a\alpha }^A(x)`$ which annihilates an antiquark and creates a quark. Unlike in Minkowski space-time, $`\psi `$ and $`\overline{\psi }`$ must be treated as independent fields, so we emphasize that $`\overline{\psi }\psi ^{}\gamma _4`$. This is required in order to simultaneously satisfy Euclidean covariance of the fields, the canonical anticommutation relations, and the equality of the Euclidean two-point function with the relativistic Feynman propagator continued to imaginary timesOsterwalder and Schrader (1973); Williams (1974). Our Euclidean space-time Dirac-$`\gamma `$ matrices are related to their Minkowski counterparts by
$`\gamma ^4=\gamma _4=\gamma _M^0,\gamma _k=\gamma ^k=i\gamma _M^k,`$ (11)
$`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu },\gamma _\mu ^{}=\gamma _\mu ,`$ (12)
$`\gamma ^5=\gamma _5=\gamma _4\gamma _1\gamma _2\gamma _3=\gamma _M^5.`$ (13)
Throughout this paper, we use the standard Dirac-Pauli representation for the $`\gamma `$-matrices:
$$\gamma _k=(\begin{array}{cc}0& i\sigma _k\\ i\sigma _k& 0\end{array}),\gamma _4=(\begin{array}{cc}I& 0\\ 0& I\end{array}),$$
(14)
where the Pauli spin matrices are given by
$$\sigma _1=(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}),\sigma _2=(\begin{array}{cc}\hfill 0& \hfill i\\ \hfill i& \hfill 0\end{array}),\sigma _3=(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}).$$
(15)
It has long been known that operators constructed out of smeared fields have dramatically reduced mixings with the high frequency modes of the theory. Thus, our operators are constructed using spatially-smoothed link variables $`\stackrel{~}{U}_j(x)`$ and spatially-smeared quark fields $`\stackrel{~}{\psi }(x)`$. The spatial links are smeared using either the stout-link procedure described in Ref. Morningstar and Peardon (2004) or the method introduced in Ref. Albanese et al. (1987). Note that only spatial staples are used in the link smoothening; no temporal staples are used, and the temporal link variables are not smeared. The smeared quark fields are defined byAlford et al. (1996)
$`\stackrel{~}{\psi }(x)`$ $`=`$ $`\left(1+{\displaystyle \frac{\sigma _s^2}{4n_\sigma }}\stackrel{~}{\mathrm{\Delta }}\right)^{n_\sigma }\psi (x),`$ (16)
$`\stackrel{~}{\overline{\psi }}(x)`$ $`=`$ $`\overline{\psi }(x)\left(1+{\displaystyle \frac{\sigma _s^2}{4n_\sigma }}\stackrel{}{\stackrel{~}{\mathrm{\Delta }}}\right)^{n_\sigma },`$ (17)
where $`\sigma _s`$ and $`n_\sigma `$ are tunable parameters ($`n_\sigma `$ is a positive integer) and the three-dimensional covariant Laplacian operators are defined in terms of the smeared link variables $`\stackrel{~}{U}_j(x)`$ as follows:
$`\stackrel{~}{\mathrm{\Delta }}O(x)`$ $`=`$ $`{\displaystyle \underset{k=\pm 1,\pm 2,\pm 3}{}}\left(\stackrel{~}{U}_k(x)O(x+\widehat{k})O(x)\right),`$ (18)
$`\overline{O}(x)\stackrel{}{\stackrel{~}{\mathrm{\Delta }}}`$ $`=`$ $`{\displaystyle \underset{k=\pm 1,\pm 2,\pm 3}{}}\left(\overline{O}(x+\widehat{k})\stackrel{~}{U}_k^{}(x)\overline{O}(x)\right),`$ (19)
where $`O(x)`$ and $`\overline{O}(x)`$ are operators defined at lattice site $`x`$ with appropriate color structure, and noting that $`\stackrel{~}{U}_k(x)=\stackrel{~}{U}_k^{}(x\widehat{k})`$. The smeared fields $`\stackrel{~}{\psi }`$ and $`\stackrel{~}{\overline{\psi }}`$ are Grassmann-valued; in particular, these fields anticommute in the same way that the original fields do, and the square of each smeared field vanishes.
Our operators are designed with one eye on capturing the states of interest and the other eye on facilitating the efficient computation of the hadron correlation functions. Since the baryon resonances are expected to be large objects, the use of extended operators is crucial. Our choice of basic building blocks is motivated by the need to incorporate spatial extensions so as to facilitate the efficient and gauge-invariant assembly of many large hadron operators. To capture orbital structure, the quarks must be displaced in different directions, and to capture radial structure, the quarks must be displaced by different distances. To maintain gauge invariance, covariant displacements in terms of the gauge-field parallel transporters must be used. To simplify matters, we consider only straight-path displacements in the six directions along the axes of the spatial cubic lattice: $`j=\pm 1,\pm 2,\pm 3`$. Hence, we shall assemble our baryon (and later, meson) operators using the following basic building blocks:
$$\stackrel{~}{\psi }_{a\alpha }^A,\stackrel{~}{\overline{\psi }}_{a\alpha }^A,\left(\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi }\right)_{a\alpha }^A,\left(\stackrel{~}{\overline{\psi }}\stackrel{~}{D}_j^{(p)}\right)_{a\alpha }^A,j=\pm 1,\pm 2,\pm 3,$$
where $`A`$ is a flavor index, $`a`$ is a color index, $`\alpha `$ is a Dirac spin index, and the $`p`$-link gauge-covariant displacement operator in the $`j`$-th direction is defined by
$$\stackrel{~}{D}_j^{(p)}(x,x^{})=\stackrel{~}{U}_j(x)\stackrel{~}{U}_j(x+\widehat{j})\mathrm{}\stackrel{~}{U}_j(x+(p1)\widehat{j})\delta _{x^{},x+p\widehat{j}},$$
(20)
for $`j=\pm 1,\pm 2,\pm 3`$ and $`p1`$. In what follows, we can achieve a significant economy in the notation by defining a zero-displacement operator
$$\stackrel{~}{D}_0^{(p)}(x,x^{})=\delta _{xx^{}},$$
(21)
to indicate no displacement. In this way, our basic building blocks can be listed more succinctly:
$$\left(\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi }\right)_{a\alpha }^A,\left(\stackrel{~}{\overline{\psi }}\stackrel{~}{D}_j^{(p)}\right)_{a\alpha }^A,3j3.$$
(22)
Sometimes it will be convenient to make the flavor quantum number more apparent by writing
$$\stackrel{~}{u}_{a\alpha }(x)\stackrel{~}{\psi }_{a\alpha }^u(x),\stackrel{~}{d}_{a\alpha }(x)\stackrel{~}{\psi }_{a\alpha }^d(x),$$
(23)
and similarly for the $`s,c,b`$ quarks. It is important to remember that the $`\stackrel{~}{u},\stackrel{~}{d},\stackrel{~}{s},`$ etc. operators so defined refer to smeared quark fields.
## IV Three-quark elemental operators
Having chosen the basic building blocks for our hadron operators listed in Eq. (22), the next step is to devise elemental baryon operators $`B_i^F(t)`$ having the appropriate color and flavor structure. These operators are chosen such that they are gauge-invariant and transform irreducibly under the isotopic spin symmetry. Explicitly dealing with $`SU(3)`$ flavor symmetry is not necessary, as will be discussed below.
Gauge invariance is easily handled. Given three quarks with color indices $`a,b,c`$ associated with the same lattice site $`x`$, there is only one way of combining the color indices to arrive at a locally gauge-invariant object: the use of the antisymmetric Levi-Civita symbol $`\epsilon _{abc}`$. Similarly, covariantly-displaced quark fields must also be connected by an $`\epsilon _{abc}`$ coupling at a single lattice site. To simplify the operator construction, we consider only combinations of displaced quarks having the same displacement length $`p`$. Thus, all of our three-quark baryon operators are linear superpositions of gauge-invariant terms of the form
$`\mathrm{\Phi }_{\alpha \beta \gamma ;ijk}^{ABC}(t)={\displaystyle \underset{𝒙}{}}\epsilon _{abc}((\stackrel{~}{D}_i^{(p)}\stackrel{~}{\psi })_{a\alpha }^A(𝒙,t)`$
$`\times (\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi })_{b\beta }^B(𝒙,t)(\stackrel{~}{D}_k^{(p)}\stackrel{~}{\psi })_{c\gamma }^C(𝒙,t)).`$ (24)
The “barred” operators have the form
$`\overline{\mathrm{\Phi }}_{\alpha \beta \gamma ;ijk}^{ABC}(t)={\displaystyle \underset{𝒙}{}}\epsilon _{abc}((\stackrel{~}{\overline{\psi }}\stackrel{~}{D}_k^{(p)})_{c\gamma ^{}}^C(𝒙,t)\gamma _{\gamma ^{}\gamma }^4`$
$`\times (\stackrel{~}{\overline{\psi }}\stackrel{~}{D}_j^{(p)})_{b\beta ^{}}^B(𝒙,t)\gamma _{\beta ^{}\beta }^4(\stackrel{~}{\overline{\psi }}\stackrel{~}{D}_i^{(p)})_{a\alpha ^{}}^A(𝒙,t)\gamma _{\alpha ^{}\alpha }^4).`$ (25)
Note that the barred composite operator $`\overline{\mathrm{\Phi }}`$ is defined differently than the barred fermion field operator $`\overline{\psi }`$ due to the presence of the $`\gamma _4`$ spin matrices in Eq. (25). Their presence is needed to ensure that the resulting correlation matrices satisfy the desirable property of hermiticity. These $`\gamma _4`$ matrices do not affect the transformation properties of these operators under proper spatial rotations and parity, although they do affect the Lorentz boost properties.
The simplest way to combine the basic building blocks previously described is to combine three quark fields at a single lattice site, corresponding to $`i=j=k=0`$ in Eq. (24). We refer to these operators as single-site operators. Next, one of the three quarks can be displaced; such singly-displaced operators correspond to $`i=j=0,k0`$ in Eq. (24). In these operators, the two quarks which are not displaced may be viewed as forming a localized diquark, so such operators may be important if baryon formation is dominated by a quark-diquark mechanism.
Two of the quarks can be displaced; they can be displaced either in opposite directions (doubly-displaced-I), or in orthogonal directions (doubly-displaced-L). In such operators, one can in general choose different lengths for the two displacements, but to simplify matters, we restrict our attention to the case in which both displacements have the same length $`p`$. Since the displacement of two quark fields essentially results in an object in which all three quarks are at different lattice sites, one may be tempted to exclude these operators in favor of triply-displaced operators. However, the relative importance of so-called $`Y`$-flux and $`\mathrm{\Delta }`$-flux formation of the gluon field in three-quark systems (the $`\mathrm{\Delta }`$ is actually a quantum superposition of $`V`$-flux forms) is a much-discussed issue (see, for example, Refs. Sommer and Wosiek (1986); Bali (2001); Alexandrou et al. (2002, 2003); Jahn and de Forcrand (2004); Takahashi et al. (2001, 2002)). Hence, it is important to include some operators with significant overlap with a $`\mathrm{\Delta }`$-flux configuration, as well as operators with strong mixing with a $`Y`$-flux configuration. The doubly-displaced operators above are suitable for $`\mathrm{\Delta }`$-flux formation.
Lastly, all three quarks can be displaced (again, restricting attention to the case of equal distances). If two of the quarks are displaced in opposite directions, this produces a coplanar T-shape (triply-displaced-T); alternatively, the three quarks can be displaced in mutually orthogonal directions (triply-displaced-O). These operators are suitable for producing $`Y`$-flux configurations.
These operators, summarized and illustrated in Table 1, allow a large number of baryon operators to be constructed using a relatively small number of quark propagator sources. For a given reference source site $`x`$, quark propagators must be evaluated using only a handful of different sources: for each quark mass value, we need a local source and displaced sources in at most each of the six directions. However, rotational invariance of the baryon correlation functions can be exploited to reduce the number of source displacement directions. For singly-displaced operators at the source, a simultaneous rotation of the source and sink can always be used to align the source displacement along the $`+z`$ direction, say. For doubly-displaced-I sources, a rotation can always be done to align one displacement along the $`+z`$ direction and the other along the $`z`$ direction. Doubly-displaced-L sources can be rotated so the source displacements are along the $`+y`$ and $`+z`$ directions, triply-displaced-T sources can be rotated to align the displacements along the $`+y,+z,`$ and $`z`$ directions, and triply-displaced-O sources can always be rotated so the displacements are along the $`+x,+y,+z`$ directions. In total, only source displacements along the $`+x,+y,+z,z`$ directions (four directions) are required. Hence, the number of conjugate gradient inversions needed is
$$N_{\mathrm{inversions}}=N_cN_{sp}N_\kappa (1+4N_p),$$
(26)
where $`N_c=3`$ is the number of colors, $`N_{sp}=4`$ is the number of Dirac spin components, $`N_p`$ is the number of displacement lengths $`p`$, and $`N_\kappa `$ is the number of quark masses to be used.
Incorporating the isospin symmetry is also straightforwardly done. Let $`\tau _1,\tau _2,\tau _3`$ denote the three hermitian generators of the isospin symmetry satisfying the commutation relations $`[\tau _i,\tau _j]=i\epsilon _{ijk}\tau _k`$. A set of quantum operators $`\overline{O}_{I_3}^{(I)}`$ transforms under isospin according to the irreducible representation $`I`$ if
$`[\tau _3,\overline{O}_{I_3}^{(I)}]`$ $`=`$ $`I_3\overline{O}_{I_3}^{(I)},`$ (27)
$`[\tau _+,\overline{O}_{I_3}^{(I)}]`$ $`=`$ $`\sqrt{(II_3)(I+I_3+1)}\overline{O}_{I_3+1}^{(I)},`$ (28)
$`[\tau _{},\overline{O}_{I_3}^{(I)}]`$ $`=`$ $`\sqrt{(I+I_3)(II_3+1)}\overline{O}_{I_31}^{(I)},`$ (29)
where $`\tau _\pm =\tau _1\pm i\tau _2`$. It follows from these relations that
$`[\tau _3,[\tau _3,\overline{O}_{I_3}^{(I)}]]+\frac{1}{2}[\tau _+,[\tau _{},\overline{O}_{I_3}^{(I)}]]`$
$`+\frac{1}{2}[\tau _{},[\tau _+,\overline{O}_{I_3}^{(I)}]]=I(I+1)\overline{O}_{I_3}^{(I)}.`$ (30)
Baryonic operators of definite isospin $`I`$ and $`I_3`$ are easily constructed using the above relations and the following commutation rules involving the isospin generators and the barred $`\overline{u},\overline{d},\overline{s}`$ quark field operators (suppressing all indices except flavor):
$$\begin{array}{ccccccccc}\hfill [\tau _3,\overline{u}]& =& \frac{1}{2}\overline{u},\hfill & \hfill [\tau _3,\overline{d}]& =& \frac{1}{2}\overline{d},\hfill & \hfill [\tau _3,\overline{s}]& =& 0,\hfill \\ \hfill [\tau _+,\overline{u}]& =& 0,\hfill & \hfill [\tau _+,\overline{d}]& =& \overline{u},\hfill & \hfill [\tau _+,\overline{s}]& =& 0,\hfill \\ \hfill [\tau _{},\overline{u}]& =& \overline{d},\hfill & \hfill [\tau _{},\overline{d}]& =& 0,\hfill & \hfill [\tau _{},\overline{s}]& =& 0,\hfill \end{array}$$
(31)
and for the unbarred field operators,
$$\begin{array}{ccccccccc}\hfill [\tau _3,u]& =& \frac{1}{2}u,\hfill & \hfill [\tau _3,d]& =& \frac{1}{2}d,\hfill & \hfill [\tau _3,s]& =& 0,\hfill \\ \hfill [\tau _{},u]& =& 0,\hfill & \hfill [\tau _{},d]& =& u,\hfill & \hfill [\tau _{},s]& =& 0,\hfill \\ \hfill [\tau _+,u]& =& d,\hfill & \hfill [\tau _+,d]& =& 0,\hfill & \hfill [\tau _+,s]& =& 0.\hfill \end{array}$$
(32)
Due to the isospin symmetry of QCD in the $`m_u=m_d`$ approximation, the particle masses do not depend on $`I_3`$, so we limit our attention to only one value of $`I_3`$ in each isospin $`I`$, strangeness $`S`$ sector, choosing the maximal $`I_3=I`$ value. Using the above isospin relations, we first write down all possible flavor combinations appropriate for each isospin channel for the three-quark elemental operators from Table 1. These are listed in Table 2. For each flavor sector and quark-displacement type, we then determine a maximal set of linearly independent operators, giving us the final set of elemental operators we use. A symbolic computer program capable of manipulating Grassmann fields was written using Maple 9.5 and utilized to identify linearly independent operators. The independent elemental operators we chose in the different flavor sections are described in Tables 3, 4, 5, and 6.
Due to an approximate $`SU(3)`$ $`uds`$-flavor symmetry, quark flavor combinations in baryon operators are usually chosen according to the irreducible representations of $`SU(3)`$ flavor. Such combinations are simply linear superpositions of the operators presented above. Since we plan to obtain Monte Carlo estimates of the complete correlation matrices of operators including all allowed flavor combinations, the use of linear superpositions which transform irreducibly under $`SU(3)`$ flavor is unnecessary. Our choice of operators described above is dictated by computational simplicity and efficiency.
Lastly, note that the baryons $`\mathrm{\Lambda }_c,\mathrm{\Sigma }_c,\mathrm{\Xi }_{cc}`$, and $`\mathrm{\Omega }_{ccc}`$ can be studied using the operators presented above if the $`s`$ quark is replaced with a $`c`$ quark. Similarly, replacing the $`s`$ quark by a $`b`$ quark in the above operators allows us to investigate the $`\mathrm{\Lambda }_b,\mathrm{\Sigma }_b,\mathrm{\Xi }_{bb}`$, and $`\mathrm{\Omega }_{bbb}`$ baryons. Some other baryons, such as $`\mathrm{\Omega }_{cc}`$, can be studied using other suitable flavor replacements in the above operators, whereas the investigations of baryons such as $`\mathrm{\Xi }_c`$ containing $`usc`$ quarks cannot directly exploit the above tables, requiring slight modifications.
## V Projections onto symmetry sectors
The final step in our operator construction is to apply group-theoretical projections to obtain operators which transform irreducibly under all lattice rotation and reflection symmetries. The basic building blocks used to assemble our baryon operators transform under the allowed spatial rotations and reflections of the point group $`O_h`$ according to
$`U_R\left(\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi }(x)\right)_{a\alpha }^AU_R^{}=S(R)_{\alpha \beta }^1\left(\stackrel{~}{D}_{R\widehat{j}}^{(p)}\stackrel{~}{\psi }(Rx)\right)_{a\beta }^A,`$ (33)
$`U_R\left(\stackrel{~}{\overline{\psi }}(x)\stackrel{~}{D}_j^{(p)}\right)_{a\alpha }^AU_R^{}=\left(\stackrel{~}{\overline{\psi }}(Rx)\stackrel{~}{D}_{R\widehat{j}}^{(p)}\right)_{a\beta }^AS(R)_{\beta \alpha },`$ (34)
where the transformation matrices for spatial inversion $`I_s`$ and proper rotations $`C_{nj}`$ through angle $`2\pi /n`$ about axis $`Oj`$ are given by
$`S(C_{nj})`$ $`=`$ $`\mathrm{exp}\left(\frac{1}{8}\omega _{\mu \nu }[\gamma _\mu ,\gamma _\nu ]\right),`$ (35)
$`S(I_s)`$ $`=`$ $`\gamma _4,`$ (36)
with $`\omega _{kl}=2\pi \epsilon _{jkl}/n`$ and $`\omega _{4k}=\omega _{k4}=0`$ ($`\omega _{\mu \nu }`$ is an antisymmetric tensor which parametrizes rotations and boosts). A rotation by $`\pi /2`$ about the $`y`$-axis is conventionally denoted by $`C_{4y}`$, and $`C_{4z}`$ denotes a rotation by $`\pi /2`$ about the $`z`$-axis. These particular group elements are given by
$$S(C_{4y})=\frac{1}{\sqrt{2}}(1+\gamma _1\gamma _3),S(C_{4z})=\frac{1}{\sqrt{2}}(1+\gamma _2\gamma _1).$$
(37)
The allowed rotations on a three-dimensional spatially-isotropic cubic lattice form the octahedral group $`O`$ which has 24 elements. Inclusion of spatial inversion yields the point group $`O_h`$ which has 48 elements occurring in ten conjugacy classes. All elements of $`O_h`$ can be generated from appropriate products of only $`C_{4y}`$, $`C_{4z}`$, and $`I_s`$.
Charge conjugation is another symmetry of our theory. Under charge conjugation $`𝒞`$, the link variables $`UU^{}`$ and our basic building blocks transform according to
$`𝒞\left(\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi }(x)\right)_{a\alpha }^A𝒞^{}`$ $`=`$ $`\left(\stackrel{~}{\overline{\psi }}(x)\stackrel{~}{D}_j^{(p)}\right)_{a\beta }^AC_{\beta \alpha }^{},`$ (38)
$`𝒞\left(\stackrel{~}{\overline{\psi }}(x)\stackrel{~}{D}_j^{(p)}\right)_{a\alpha }^A𝒞^{}`$ $`=`$ $`C_{\alpha \beta }\left(\stackrel{~}{D}_j^{(p)}\stackrel{~}{\psi }(x)\right)_{a\beta }^A,`$ (39)
where the charge conjugation matrix $`C`$ must be unitary, antisymmetric, and satisfy $`C\gamma _\mu C^{}=\gamma _\mu ^T`$. Our choice for $`C`$ in the Dirac-Pauli representation is $`C=\gamma _4\gamma _2`$.
Operators which transform according to the irreducible representations of $`O_h`$ can then be constructed using the well-known group-theoretical projections given in Eq. (6). Orthogonality relations and hence, projection techniques, in group theory apply only to single-valued irreducible representations. However, the fermionic representations are double-valued representations of $`O_h`$. The commonly-used trick to circumvent this difficulty is to exploit the equivalence of the double-valued irreps of $`O_h`$ with the extra single-valued irreps of the so-called double point group $`O_h^D`$. This group is formed by introducing a new element $`\overline{E}`$ which represents a rotation by an angle $`2\pi `$ about any axis, such that $`\overline{E}^2=E`$ (the identity). By including such an element, the total number of elements in $`O_h^D`$ is double the number of elements in $`O_h`$. The 96 elements of $`O_h^D`$ occur in sixteen conjugacy classes.
Since baryons are fermions, we need only be concerned with the six double-valued irreps of $`O_h`$. There are four two-dimensional irreps $`G_{1g},G_{1u},G_{2g}`$, and $`G_{2u}`$, and two four-dimensional irreps $`H_g`$ and $`H_u`$. The subscript $`g`$ refers to even parity states, whereas the subscript $`u`$ refers to odd parity states. The irreps $`G_{1g}`$ and $`G_{1u}`$ contain the spin-1/2 states, spin-3/2 states reside in the $`H_g`$ and $`H_u`$, and two of the spin projections of the spin-5/2 states occur in the $`G_{2g}`$ and $`G_{2u}`$ irreps, while the remaining four projections reside in the $`H_g`$ and $`H_u`$ irreps. The spin content of each $`O_h`$ irrep in the continuum limit is summarized in Table 7. This table lists the number of times that each of the $`O_h`$ irreps occurs in the $`J=\frac{1}{2},\frac{3}{2},\frac{5}{2},\mathrm{}`$ representations of $`SU(2)`$ subduced to $`O_h`$.
To carry out the projections in Eq. (6), explicit representation matrices are needed. Our choice of representation matrices is summarized in Table 8. Matrices for only the group elements $`C_{4y},C_{4z}`$, and $`I_s`$ are given in Table 8 since the representation matrices for all other group elements can be obtained by suitable multiplications of the matrices for the three generating elements. For baryons, the representation matrix for $`\overline{E}`$ in each of the $`O_h^D`$ extra irreps is $`1`$ times the identity matrix.
Note that Table 7 is the key to identifying the continuum-limit spin $`J`$ corresponding to the masses extracted in our Monte Carlo calculations. For example, to identify an even parity baryon as having $`J=1/2`$, a level must be observed in the $`G_{1g}`$ channel, and there must be no degenerate partners in either of the $`G_{2g}`$ or $`H_g`$. A level observed in the $`H_g`$ channel with no degenerate partners in the $`G_{1g}`$ and $`G_{2g}`$ channels (in the continuum limit) is a $`J=3/2`$ state. Degenerate partners observed in the $`G_{2g}`$ and $`H_g`$ channels with no partner in the $`G_{1g}`$ channel indicates a $`J=5/2`$ baryon. In other words, Table 7 details the patterns of continuum-limit degeneracies corresponding to each half-integral $`J`$ value.
Our operators are constructed using fermion fields $`\psi (x)`$ which annihilate a quark and create an antiquark. Hence, each of our baryon operators annihilates a three-quark system of a given parity $`P`$ and creates a three-antiquark system of the same parity $`P`$. This means that in the baryon propagator, a baryon of parity $`P`$ propagates forward in time while an antibaryon of parity $`P`$ propagates backwards in time. Unlike boson fields, a fermion and its antifermion have opposite intrinsic parities, so that the antibaryon propagating backwards in time is the antiparticle of the parity partner of the baryon propagating forwards in time. Since chiral symmetry is spontaneously broken, the masses of parity partners may differ. The forward propagating baryon will have a mass different from that of the antibaryon propagating backwards in time. If the even and odd parity baryon operators are carefully designed with respect to one another, it is possible to arrange a definite relationship between the correlation matrix elements of one parity for $`t>0`$ and the opposite-parity matrix elements for $`t<0`$, allowing an increase in statistics. Our operators are designed to take advantage of this symmetry (see below).
Our construction of the irreducible $`B_i^{\mathrm{\Lambda }\lambda F}(t)`$ baryon operators is done in the following sequence of steps.
(a) A set of $`M_B`$ linearly-independent elemental operators $`B_j^F(t)`$ that transform among one another under $`O_h^D`$ is identified with the help of the computer software mentioned earlier. This is done by starting with all possible operators of a given type (single-site, singly-displaced, and so on), then using the Maple program to detect dependencies between the operators. To find such dependencies, the computer program expresses each $`B_j^F(t)`$ operator as a sum of products of Grassmann fields (or gauge-covariantly displaced Grassmann fields) with explicit color, flavor, and spin indices: $`B_j^F=_kg_{jk}^FO_k`$. The coefficients $`g_{jk}^F`$ form a $`M_B\times M_O`$ matrix, where $`M_O`$ is the total number of $`O_k`$ operators encountered, which is much larger than $`M_B`$. Since each $`O_k`$ operator is a single product of three displaced Grassmann fields, these operators are linearly independent, so the linear independence of the $`B_j^F`$ operators boils down to the linear independence of the rows of the $`g_{jk}^F`$ matrix. In this way, the set of all possible operators of a given type is easily reduced to a set containing only linearly-independent operators.
(b) We obtain the $`M_B\times M_B`$ representation matrices $`W_{ij}(R)`$ which satisfy
$$U_RB_i^F(t)U_R^{}=\underset{j=1}{\overset{M_B}{}}B_j^F(t)W_{ji}(R).$$
(40)
Our Maple program determines the $`i`$-th column of the $`W_{ji}(R)`$ matrix for a given symmetry transformation $`R`$ as follows. First, the $`B_i^F(t)`$ operator is expressed as a sum of products of displaced Grassmann fields with explicit color, flavor, and spin indices as in the previous step: $`B_i^F=_kg_{ik}^FO_k`$. Next, the $`R`$ symmetry transformation is applied to the displaced Grassmann fields in each $`O_k`$ term in this sum of products using Eq. (33). The resulting sum of terms $`U_RB_i^F(t)U_R^{}=_kh_{ik}^FO_k`$ is then expressed as a linear superposition of the original $`M_B`$ operators using the Moore-Penrose pseudoinversePenrose (1955) of the $`g_{ik}^F`$ matrix. If the transformed operator contains Grassmann products which are not in the original set of $`O_k`$ operators, then this signals that the starting basis of $`B_j^F`$ operators is incomplete, but with our method in the previous step of choosing linearly-independent operators, this did not occur. In this way, the $`W_{ji}(R)`$ matrices for the generating group elements $`C_{4y}`$, $`C_{4z}`$, and $`I_s`$ are obtained. The matrices for all other group elements are then determined by appropriate products of these three matrices. At this point, we have a set of $`M_B`$ operators $`B_i^F(t)`$ which form the basis of a reducible representation given by the $`W_{ij}(R)`$ matrices. Our remaining task is to find a change of basis such that the resulting representation matrices are block-diagonal with the blocks given by the irreducible representations of $`O_h`$.
(c) Since the resulting $`W_{ji}(R)`$ matrices may not be unitary, we compute the hermitian metric matrix $`M`$
$$M_{ij}=\frac{1}{g_{O_h^D}}\underset{RO_h^D}{}\underset{k=1}{\overset{M_B}{}}W_{ki}(R)^{}W_{kj}(R).$$
(41)
This matrix will be needed in a later step where it will facilitate the full block-diagonalization of the $`W_{ij}(R)`$ matrices.
(d) For each even-parity irrep $`\mathrm{\Lambda }`$, we compute the large $`M_B\times M_B`$ projection matrix for row $`\lambda =1`$:
$$P_{ij}^{\mathrm{\Lambda }\lambda F}=\frac{d_\mathrm{\Lambda }}{g_{O_h^D}}\underset{RO_h^D}{}\left[\mathrm{\Gamma }_{\lambda \lambda }^{(\mathrm{\Lambda })}(R)W_{ji}(R)\right]_{\lambda =1}.$$
(42)
This is one of the most important steps in our operator construction. Applying the group-theoretical projection of Eq. (6) to the operator $`B_i^F`$, then utilizing Eq. (40), produces a new operator $`_i^{\mathrm{\Lambda }\lambda F}(t)=_jP_{ij}^{\mathrm{\Lambda }\lambda F}B_j^F(t)`$ which resides in the subspace of operators which transform according to the row $`\lambda =1`$ of the given irrep $`\mathrm{\Lambda }`$. In other words, the $`i`$-th row of the projection matrix $`P`$ contains the superposition coefficients of the projected $`_i^{\mathrm{\Lambda }\lambda F}`$ operator.
(e) Although group theory guarantees that the resulting projected $`_i^{\mathrm{\Lambda }\lambda F}`$ operators reside in the subspace associated with row $`\lambda `$ of the $`\mathrm{\Lambda }`$ irrep, it does not guarantee that all of the resulting operators are linearly independent. The maximum number of independent operators in the projected subspace is given by the rank $`r`$ of the projection matrix. Hence, the next step is to form $`r`$ superpositions of the $`_i^{\mathrm{\Lambda }\lambda F}`$ operators such that the resulting $`r`$ operators are linearly independent. The choice of these operators is not unique. In practice, these linear combinations are obtained using the well-known Gram-Schmidt procedure, but with a modified inner product to incorporate the metric matrix $`M`$. The use of the metric matrix $`M`$ ensures full block-diagonalization of the original $`W_{ij}(R)`$ matrices. Using such a procedure, the final operators, expressed in terms of the original set of $`B_i^F`$ operators by
$$B_i^{\mathrm{\Lambda }\lambda F}(t)=\underset{j=1}{\overset{M_B}{}}c_{ij}^{\mathrm{\Lambda }\lambda F}B_j^F(t),(\lambda =1)$$
(43)
have superposition coefficients $`c_{ij}^{\mathrm{\Lambda }\lambda F}`$ that satisfy
$$\underset{k,l=1}{\overset{M_B}{}}c_{ik}^{\mathrm{\Lambda }\lambda F}M_{kl}c_{jl}^{\mathrm{\Lambda }\lambda F}=\delta _{ij},(i=1\mathrm{}r).$$
(44)
(f) For each of the $`r`$ operators $`B_i^{\mathrm{\Lambda }\lambda F}(t)`$ in the first row $`\lambda =1`$, we obtain partner operators in all other rows $`\mu >1`$ using
$$c_{ik}^{\mathrm{\Lambda }\mu F}=\underset{j=1}{\overset{M_B}{}}c_{ij}^{\mathrm{\Lambda }\lambda F}\frac{d_\mathrm{\Lambda }}{g_{O_h^D}}\underset{RO_h^D}{}\mathrm{\Gamma }_{\mu \lambda }^{(\mathrm{\Lambda })}(R)W_{kj}(R).$$
(45)
The use of operators belonging to other rows will be important for increasing the statistics of our Monte Carlo calculations, as will be discussed below.
(g) Although the odd-parity operators can be obtained using the same procedure described above, we instead utilize charge conjugation to construct the odd-parity operators. Consider the correlation matrix element of even-parity operators for $`t0`$. Suppressing flavor and displacement indices, one sees that invariance under charge conjugation implies that
$`C_{ij}(t)`$ $`=`$ $`c_{\alpha \beta \gamma }^{(i)}c_{\overline{\alpha }\overline{\beta }\overline{\gamma }}^{(j)}0|\mathrm{\Phi }_{\alpha \beta \gamma }(t)\overline{\mathrm{\Phi }}_{\overline{\alpha }\overline{\beta }\overline{\gamma }}(0)|0,`$
$`=`$ $`c_{\alpha \beta \gamma }^{(i)}c_{\overline{\alpha }\overline{\beta }\overline{\gamma }}^{(j)}0|𝒞^{}𝒞\mathrm{\Phi }_{\alpha \beta \gamma }(t)𝒞^{}𝒞\overline{\mathrm{\Phi }}_{\overline{\alpha }\overline{\beta }\overline{\gamma }}(0)𝒞^{}𝒞|0,`$
$`=`$ $`c_{\alpha \beta \gamma }^{(i)}c_{\overline{\alpha }\overline{\beta }\overline{\gamma }}^{(j)}0|\overline{\mathrm{\Phi }}_{\alpha ^{}\beta ^{}\gamma ^{}}(t)\mathrm{\Phi }_{\overline{\alpha }^{}\overline{\beta }^{}\overline{\gamma }^{}}(0)|0`$
$`\times \gamma _{\alpha ^{}\alpha }^2\gamma _{\beta ^{}\beta }^2\gamma _{\gamma ^{}\gamma }^2\gamma _{\overline{\alpha }^{}\overline{\alpha }}^2\gamma _{\overline{\beta }^{}\overline{\beta }}^2\gamma _{\overline{\gamma }^{}\overline{\gamma }}^2,`$
$`=`$ $`c_{\alpha \beta \gamma }^{(i)}c_{\overline{\alpha }\overline{\beta }\overline{\gamma }}^{(j)}0|\overline{\mathrm{\Phi }}_{\alpha ^{}\beta ^{}\gamma ^{}}(0)\mathrm{\Phi }_{\overline{\alpha }^{}\overline{\beta }^{}\overline{\gamma }^{}}(t)|0`$
$`\times \gamma _{\alpha ^{}\alpha }^2\gamma _{\beta ^{}\beta }^2\gamma _{\gamma ^{}\gamma }^2\gamma _{\overline{\alpha }^{}\overline{\alpha }}^2\gamma _{\overline{\beta }^{}\overline{\beta }}^2\gamma _{\overline{\gamma }^{}\overline{\gamma }}^2,`$
using invariance under time translations of the above expectation value and invariance of the vacuum under charge conjugation. The last line above represents the correlation of odd-parity operators propagating temporally backwards. Hence, for a given even-parity operator $`B_i^g(t)`$, we can define an odd-parity operator $`B_i^u(t)`$ by rotating the three Dirac indices using the $`\gamma _2`$ matrix and replacing the expansion coefficients by their complex conjugates such that the correlation matrices of the even and odd parity operators are related by
$`C_{ij}^{G_{1g}}(t)=C_{ij}^{G_{1u}}(t)^{},C_{ij}^{G_{2g}}(t)=C_{ij}^{G_{2u}}(t)^{},`$
$`C_{ij}^{H_g}(t)=C_{ij}^{H_u}(t)^{}.`$ (46)
For
a lattice of $`N_t`$ sites in the time direction with periodic $`(\eta _t=1)`$ or antiperiodic $`(\eta _t=1)`$ boundary conditions, this means that
$$C_{ij}^{G_{1g}}(t)=\eta _tC_{ij}^{G_{1u}}(N_tt)^{},$$
(47)
and similarly for the other irreps. This allows us to appropriately average over forward and backward temporal propagations for increased statistics.
In the absence of any external applied fields, the energies of the baryons do not depend on the row $`\lambda `$ of a given irrep $`\mathrm{\Lambda }`$, so we can increase statistics by averaging over all rows. The construction of the operators in the different rows of the irreps as described above leads to correlation matrices which satisfy $`C_{ij}^{\mathrm{\Lambda }\lambda F}(t)=C_{ij}^{\mathrm{\Lambda }\mu F}(t),`$ for all rows $`\lambda ,\mu `$. Hence, the correlation matrix elements themselves can be averaged over rows.
The single-site operators produced by the above procedure are presented in Tables 9-12. It is not possible to list all of the singly-displaced, doubly-displaced, and triply-displaced operators in this paper. Instead, we simply list the numbers of operators of each type which project into each row of the irreps for the $`\mathrm{\Delta }^{++}`$, $`\mathrm{\Sigma }^+,N^+,`$ and $`\mathrm{\Lambda }^0`$ baryons in Table 13. Further details about these operators are available upon request.
Note that we have not yet attempted to remove redundant operators related to others by a total lattice derivative, and that there may exist relationships between the correlation matrix elements of these operators from the equations of motion. Such relationships can be easily identified using a singular value decomposition of the correlation matrices in small-lattice low-statistics Monte Carlo calculations. Since the Schwinger-Dyson equations may relate operators of different types, such as singly-displaced operators with single-site operators, identifying these relationships must be done at a later stage in the calculations.
## VI Baryon propagators
To extract the baryon masses, we need to compute the correlations $`C_{ij}^{\mathrm{\Lambda }\lambda F}(t)=0|TB_i^{\mathrm{\Lambda }\lambda F}(t)\overline{B}_j^{\mathrm{\Lambda }\lambda F}(0)|0,`$ using the operators constructed as described above. In this section, several issues related to computing these correlation matrix elements are discussed. In particular, we discuss the use of symmetry to minimize the number of quark-propagator sources, the use of gauge-invariant three-quark propagators as an intermediate step in the baryon propagator determinations, and detail the application of Wick’s theorem.
The baryon propagators may be expressed in terms of the correlations of the elemental operators by
$$C_{ij}^{\mathrm{\Lambda }\lambda F}(t)=\underset{k,l=1}{\overset{M_B}{}}c_{ik}^{\mathrm{\Lambda }\lambda F}c_{jl}^{\mathrm{\Lambda }\lambda F}0|TB_k^F(t)\overline{B}_l^F(0)|0.$$
(48)
Note that correlations between operators in different rows of the same irrep vanish. Since the number of elemental operators is large and the quark propagators are rather expensive to compute, it is very important to use symmetry to reduce the number of quark-propagator sources. Given the invariance of the vacuum and the unitarity of the symmetry transformation operators, we know that
$$\begin{array}{c}0|TB_k^F(t)\overline{B}_l^F(0)|0\hfill \\ =0|TU_RB_k^F(t)U_R^{}U_R\overline{B}_l^F(0)U_R^{}|0,\hfill \\ =\underset{k^{},l^{}=1}{\overset{M_B}{}}W_{k^{}k}(R)W_{l^{}l}(R)^{}0|TB_k^{}^F(t)\overline{B}_l^{}^F(0)|0,\hfill \end{array}$$
for any group element $`R`$ of $`O_h`$. Hence, for each source $`\overline{B}_l^F(0)`$, we can choose a group element $`R_l`$ such that we minimize the total number of source elemental operators which must be considered. For example, consider the singly-displaced operators. We can choose an $`R_l`$ such that the displaced quark in the source is always displaced in the $`+z`$ direction. Similarly, a group element $`R_l`$ can always be chosen to rotate each of the other types of operators into a specific orientation.
The coefficients $`c_{ij}^{\mathrm{\Lambda }\lambda F}`$ in the baryon operators involve only the Dirac spin components and the quark displacement directions and are independent of the color indices and spatial sites. Thus, in calculating the baryon correlators, it is convenient to first calculate gauge-invariant three-quark propagators in which all summations over color indices and spatial sites have been done. A three-quark propagator is defined by
$`\stackrel{~}{G}_{(\alpha i|\overline{\alpha }\overline{i})(\beta j|\overline{\beta }\overline{j})(\gamma k|\overline{\gamma }\overline{k})}^{(ABC)(p\overline{p})}(t)`$ (49)
$`=`$ $`{\displaystyle \underset{𝒙}{}}\epsilon _{abc}\epsilon _{\overline{a}\overline{b}\overline{c}}\stackrel{~}{Q}_{a\alpha ip|\overline{a}\overline{\alpha }\overline{i}\overline{p}}^{(A)}(𝒙,t|𝒙_0,0)`$
$`\times `$ $`\stackrel{~}{Q}_{b\beta jp|\overline{b}\overline{\beta }\overline{j}\overline{p}}^{(B)}(𝒙,t|𝒙_0,0)\stackrel{~}{Q}_{c\gamma kp|\overline{c}\overline{\gamma }\overline{k}\overline{p}}^{(C)}(𝒙,t|𝒙_0,0),`$
where $`\stackrel{~}{Q}_{a\alpha ip|\overline{a}\overline{\alpha }\overline{i}\overline{p}}^{(A)}(𝒙,t|𝒙_0,0)`$ denotes the propagator for a single smeared quark field of flavor $`A`$ from source site $`𝒙_0`$ at time $`t=0`$ to sink site $`𝒙`$ at time $`t`$. At the sink, $`a`$ denotes color, $`\alpha `$ is the Dirac spin index, $`i`$ is the displacement direction, and $`p`$ is the displacement length, and similarly at the source for $`\overline{a},\overline{\alpha },\overline{i},`$ and $`\overline{p}`$, respectively. Notice that the three-quark propagator is symmetric under interchange of all indices associated with the same flavor. As usual, translation invariance is invoked at the source so that summation over spatial sites is done only at the sink. These three-quark propagators are computed for all possible values of the six Dirac spin indices.
Each baryon correlator is simply a linear superposition of elements of the three-quark propagators. These superposition coefficients are calculated as follows: first, the baryon operators at the source and sink are expressed in terms of the elemental operators; next, Wick’s theorem is applied to express the correlator as a large sum of three-quark propagator components; finally, symmetry operations are applied to minimize the number of source orientations, and the results are averaged over the rows of the representations. C++ code was written to perform these computations, and the resulting superposition coefficients are stored in computer files which are subsequently used as input to the Monte Carlo runs.
Wick’s theorem is an important part of expressing the baryon correlators in terms of the three-quark propagators. To simplify the notation in the following, let the indices $`\mu ,\nu ,\tau `$ each represent a Dirac spin index and a displacement direction, and suppress the displacement lengths. Define $`\overline{c}_{\mu \nu \tau }^{(i)}=c_{\mu ^{}\nu ^{}\tau ^{}}^{(i)}\gamma _{\mu \mu ^{}}^4\gamma _{\nu \nu ^{}}^4\gamma _{\tau \tau ^{}}^4,`$ then the elements of the baryon correlation matrix in the $`\mathrm{\Delta }^{++}`$ channel are given in terms of three-quark propagator components (before source-minimizing rotations) by
$`C_{ij}^{(\mathrm{\Delta })}(t)`$ $`=`$ $`c_{\mu \nu \tau }^{(i)}\overline{c}_{\overline{\mu }\overline{\nu }\overline{\tau }}^{(j)}\{\stackrel{~}{G}_{(\tau |\overline{\mu })(\nu |\overline{\nu })(\mu |\overline{\tau })}^{(uuu)}(t)`$ (50)
$`+`$ $`\stackrel{~}{G}_{(\tau |\overline{\mu })(\nu |\overline{\tau })(\mu |\overline{\nu })}^{(uuu)}(t)+\stackrel{~}{G}_{(\tau |\overline{\nu })(\nu |\overline{\mu })(\mu |\overline{\tau })}^{(uuu)}(t)`$
$`+`$ $`\stackrel{~}{G}_{(\tau |\overline{\nu })(\nu |\overline{\tau })(\mu |\overline{\mu })}^{(uuu)}(t)+\stackrel{~}{G}_{(\tau |\overline{\tau })(\nu |\overline{\nu })(\mu |\overline{\mu })}^{(uuu)}(t)`$
$`+`$ $`\stackrel{~}{G}_{(\tau |\overline{\tau })(\nu |\overline{\mu })(\mu |\overline{\nu })}^{(uuu)}(t)\}.`$
The $`N^+`$ correlators are expressed in terms of components of three-quark propagators by
$`C_{ij}^{(N)}(t)`$ $`=`$ $`c_{\mu \nu \tau }^{(i)}\overline{c}_{\overline{\mu }\overline{\nu }\overline{\tau }}^{(j)}\{\stackrel{~}{G}_{(\mu |\overline{\mu })(\nu |\overline{\nu })(\tau |\overline{\tau })}^{(uud)}`$ (51)
$`+`$ $`\stackrel{~}{G}_{(\mu |\overline{\nu })(\nu |\overline{\mu })(\tau |\overline{\tau })}^{(uud)}\stackrel{~}{G}_{(\mu |\overline{\tau })(\nu |\overline{\nu })(\tau |\overline{\mu })}^{(uud)}`$
$``$ $`\stackrel{~}{G}_{(\mu |\overline{\nu })(\nu |\overline{\tau })(\tau |\overline{\mu })}^{(uud)}\stackrel{~}{G}_{(\nu |\overline{\nu })(\tau |\overline{\mu })(\mu |\overline{\tau })}^{(uud)}`$
$``$ $`\stackrel{~}{G}_{(\nu |\overline{\mu })(\tau |\overline{\nu })(\mu |\overline{\tau })}^{(uud)}+\stackrel{~}{G}_{(\tau |\overline{\tau })(\nu |\overline{\nu })(\mu |\overline{\mu })}^{(uud)}`$
$`+`$ $`\stackrel{~}{G}_{(\tau |\overline{\nu })(\nu |\overline{\tau })(\mu |\overline{\mu })}^{(uud)}\},`$
and for the $`\mathrm{\Sigma }^+`$ and $`\mathrm{\Lambda }^0`$ channels, one finds
$`C_{ij}^{(\mathrm{\Sigma })}(t)`$ $`=`$ $`c_{\mu \nu \tau }^{(i)}\overline{c}_{\overline{\mu }\overline{\nu }\overline{\tau }}^{(j)}\{\stackrel{~}{G}_{(\mu |\overline{\mu })(\nu |\overline{\nu })(\tau |\overline{\tau })}^{(uus)}(t)`$ (52)
$`+`$ $`\stackrel{~}{G}_{(\mu |\overline{\nu })(\nu |\overline{\mu })(\tau |\overline{\tau })}^{(uus)}(t)\},`$
$`C_{ij}^{(\mathrm{\Lambda })}(t)`$ $`=`$ $`c_{\mu \nu \tau }^{(i)}\overline{c}_{\overline{\mu }\overline{\nu }\overline{\tau }}^{(j)}\{\stackrel{~}{G}_{(\mu |\overline{\mu })(\nu |\overline{\nu })(\tau |\overline{\tau })}^{(uds)}`$ (53)
$``$ $`\stackrel{~}{G}_{(\mu |\overline{\nu })(\nu |\overline{\mu })(\tau |\overline{\tau })}^{(uds)}\stackrel{~}{G}_{(\nu |\overline{\mu })(\mu |\overline{\nu })(\tau |\overline{\tau })}^{(uds)}`$
$`+`$ $`\stackrel{~}{G}_{(\nu |\overline{\nu })(\mu |\overline{\mu })(\tau |\overline{\tau })}^{(uds)}\}.`$
## VII Conclusion
We plan to undertake a comprehensive study of the spectrum of QCD using Monte Carlo computations. Our first goal in this study is to calculate the masses of as many low-lying hadron resonances as possible. Successfully extracting these masses will depend crucially on using carefully designed spatially-extended hadronic operators. In this first paper, the construction of three-quark $`\mathrm{\Delta },N,\mathrm{\Sigma },\mathrm{\Lambda },\mathrm{\Omega },\mathrm{\Xi }`$ baryon operators using group-theoretical projections was detailed. The operators were assembled out of gauge-covariantly displaced quark fields and transform according to the irreducible representations of the symmetry group of a spatial simple cubic lattice. Single-site, singly-displaced, doubly-displaced, and triply-displaced three-quark operators were considered. The guiding principles in devising our operators were maximizing overlaps with the states of interest while minimizing the number of quark-propagator sources. Identifying the continuum-limit spins $`J`$ of the states was addressed, and various issues related to computing the correlation matrix elements of the baryon operators were discussed.
Due to the complexity of these calculations and the importance of providing checks on our final results, we have been pursuing two different approaches to constructing the baryon operators. An alternative method of building the baryon operators based on Clebsch-Gordan techniques will be presented elsewhereBasak et al. . The construction of meson and multi-hadron operators will be described in subsequent papers.
The Monte Carlo software to evaluate the correlation matrix elements of these baryons operators has been written and thoroughly tested using a large variety of checks, including comparison with known results in the case of a uniform constant background gauge field. This software uses the Chroma Software System for Lattice QCDEdwards and Joo (2005) with QDP++ and QMP, developed under the Scientific Discovery through Advanced Computing initiative of the U.S. Department of Energy. Preliminary results have already been presented in Refs. Basak et al. (2004, 2005). Although a very large number of baryon operators have been devised, it is not our intent to evaluate correlation matrices using all of these operators. Such calculations would not be feasible. The next important step in our study is to remove dynamically-redundant and ineffective operators using low-statistics Monte Carlo calculations, with the goal of finding some reasonably small subsets of operators adequate for extracting the low-lying masses of interest. Such calculations are currently in progress.
###### Acknowledgements.
This work was supported by the U.S. National Science Foundation through grants PHY-0354982 and PHY-0300065, and by the U.S. Dept. of Energy under contracts DE-AC05-84ER40150 and DE-FG02-93ER-40762.
|
warning/0506/hep-th0506030.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The Kaluza-Klein hypothesis has been discussed in theoretical physics for three quarters of a century. In accordance with this hypothesis, space-time may have extra dimensions, which are unobservable for certain reasons. The explanation of this unobservability, which was put forward in the original papers by Kaluza and Klein, implies that the extra dimensions are compactified and have a very small size of the order of the Planck’s length $`l_{Pl}=1/M_{Pl}`$.
In 1983 Rubakov and Shaposhnikov put forward a new scenario for Kaluza-Klein theories, which was based on the idea of localization of fields on a domain wall . They have also proposed an ansatz for multidimensional metric, which is compatible with this hypothesis .
In the last years there appeared indications that scenarios of this type can arise in the theory of strings (see Ref. for a review and references). In this approach our three spatial dimensions are supposed to be realized as a three-dimensional hypersurface embedded into a mutidimensional space-time. Such hypersurfaces are called 3-branes, or just branes. The main goal of such scenarios was to find a solution to the hierarchy problem. It was solved either due to the sufficiently large characteristic size of extra dimensions , or due to exponential warp factor appearing in the metric . In both approaches gravity in multidimensional space-time becomes ”strong” not at the energies of the order of $`10^{19}`$GeV, but at much lower energies, maybe of the order of $`1÷10`$TeV. An attractive feature of these models is that they predict new effects which can be observed at the coming collider experiments.
In paper an exact solution for a system of two branes interacting with gravity in a five-dimensional space-time $`E`$ was found. This model is called the Randall-Sundrum model (usually abbreviated as RS1 model), and it is widely discussed in the literature (see Refs. for reviews and references). Let us denote the coordinates by $`\{x^M\}\{x^\mu ,y\}`$, $`M=0,1,2,3,4,\mu =0,1,2,3`$, the coordinate $`x^4y`$ parameterizing the fifth dimension. It forms the orbifold $`S^1/Z_2`$, which is realised as the circle of the circumference $`2R`$ with points $`y`$ and $`y`$ identified. Correspondingly, we have the usual periodicity condition in space-time $`E`$, which identifies points $`(x,y)`$ and $`(x,y+2nR)`$, and the metric $`g_{MN}`$ satisfies the orbifold symmetry conditions
$`g_{\mu \nu }(x,y)=g_{\mu \nu }(x,y),`$ (1)
$`g_{\mu 4}(x,y)=g_{\mu 4}(x,y),`$
$`g_{44}(x,y)=g_{44}(x,y).`$
The branes are located at the fixed points of the orbifold, $`y=0`$ and $`y=R`$.
The action of the model is
$$S=S_g+S_1+S_2,$$
(2)
where $`S_g`$, $`S_1`$ and $`S_2`$ are given by
$`S_g`$ $`=`$ $`{\displaystyle \frac{1}{16\pi \widehat{G}}}{\displaystyle _E}\left(R\mathrm{\Lambda }\right)\sqrt{g}d^4x𝑑y,`$ (3)
$`S_1`$ $`=`$ $`V_1{\displaystyle _E}\sqrt{\stackrel{~}{g}}\delta (y)d^4x𝑑y,`$
$`S_2`$ $`=`$ $`V_2{\displaystyle _E}\sqrt{\stackrel{~}{g}}\delta (yR)d^4x𝑑y.`$
Here $`\stackrel{~}{g}_{\mu \nu }`$ is the induced metric on the branes and the subscripts 1 and 2 label the branes. We also note that the signature of the metric $`g_{MN}`$ is chosen to be $`(,+,+,+,+)`$.
The Randall-Sundrum solution for the metric is given by
$$ds^2=g_{MN}dx^Mdx^N=\gamma _{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(4)
where $`\gamma _{\mu \nu }=e^{2\sigma (y)}\eta _{\mu \nu }`$, $`\eta _{\mu \nu }`$ is the Minkowski metric and the function $`\sigma (y)=k|y|`$ in the interval $`RyR`$. The parameter $`k`$ is positive and has the dimension of mass, the parameters $`\mathrm{\Lambda }`$ and $`V_{1,2}`$ are related to it as follows:
$$\mathrm{\Lambda }=12k^2,V_1=V_2=\frac{3k}{4\pi \widehat{G}}.$$
We see that brane 1 has a positive energy density, whereas brane 2 has a negative one. The function $`\sigma `$ has the properties
$$_4\sigma =ksign(y),\frac{^2\sigma }{y^2}=2k(\delta (y)\delta (yR))2k\stackrel{~}{\delta }.$$
(5)
Here and in the sequel $`_4\frac{}{y}`$.
We denote $`\widehat{\kappa }=\sqrt{16\pi \widehat{G}}`$, where $`\widehat{G}`$ is the five-dimensional gravitational constant, and parameterize the metric $`g_{MN}`$ as
$$g_{MN}=\gamma _{MN}+\widehat{\kappa }h_{MN},$$
(6)
$`h_{MN}`$ being the metric fluctuations. Substituting this parameterization into (2) and retaining the terms of the zero order in $`\widehat{\kappa }`$, we get the second variation action of this model . It is invariant under the gauge transformations
$`h_{MN}^{}(x,y)=h_{MN}(x,y)(_M\xi _N(x,y)+_N\xi _M(x,y)),`$ (7)
where $`_M`$ is the covariant derivative with respect to the background metric $`\gamma _{MN}`$, and the functions $`\xi _N(x,y)`$ satisfy the orbifold symmetry conditions
$`\xi ^\mu (x,y)`$ $`=`$ $`\xi ^\mu (x,y),`$ (8)
$`\xi ^4(x,y)`$ $`=`$ $`\xi ^4(x,y).`$
With the help of these gauge transformations we can impose the gauge
$$h_{\mu 4}=0,h_{44}=h_{44}(x)\varphi (x),$$
(9)
which will be called the unitary gauge (see ). We would like to emphasize once again that the branes remain straight in this gauge, i.e. we do not use the bent-brane formulation, which allegedly destroys the structure of the model, i.e. the orbifold symmetry under the reflection $`yy`$ (this problem was discussed in ).
In the linear approximation interaction with matter looks like
$$\frac{\widehat{\kappa }}{2}h^{MN}(x,y)T_{MN}\sqrt{\gamma }d^4x𝑑y,$$
(10)
where $`T_{MN}`$ is the energy-momentum tensor of the matter:
$$T_{MN}=2\frac{\delta L}{\delta \gamma ^{MN}}\gamma _{MN}L.$$
Thus, we are considering perturbations about the background, for which $`T_{MN}=0`$.
We will examine the $`h_{\mu \nu }`$-components of the metric fluctuations $`h_{MN}`$, since all four-dimensional physical effects can be described in terms of this field. Obviously, the unitary gauge conditions (9) do not fix the gauge of this field. In fact, after imposing these gauge conditions there remain gauge transformations of the form
$$\xi _\mu =e^{2\sigma }ϵ_\mu (x),$$
(11)
which change the longitudinal components of the field $`h_{\mu \nu }`$. Nevertheless, it turns out that it is convenient to solve the equations of motion for linearized gravity in the unitary gauge and then to choose an appropriate gauge in our four-dimensional world on the brane. We will use the de Donder gauge for the field $`h_{\mu \nu }`$ on the brane, which corresponds to the choice of harmonic coordinates.
## 2 Equations of motion
The equations of motion for different components of the metric fluctuations in the unitary gauge take the form (see ):
1) $`\mu \nu `$-component
$`{\displaystyle \frac{1}{2}}\left(_\rho ^\rho h_{\mu \nu }_\mu ^\rho h_{\rho \nu }_\nu ^\rho h_{\rho \mu }+{\displaystyle \frac{^2h_{\mu \nu }}{x_{}^{4}{}_{}{}^{2}}}\right)2k^2h_{\mu \nu }+{\displaystyle \frac{1}{2}}_\mu _\nu \stackrel{~}{h}+`$
$`+`$ $`{\displaystyle \frac{1}{2}}_\mu _\nu \varphi +{\displaystyle \frac{1}{2}}\gamma _{\mu \nu }\left(^\rho ^\sigma h_{\rho \sigma }_\rho ^\rho \stackrel{~}{h}{\displaystyle \frac{^2\stackrel{~}{h}}{x_{}^{4}{}_{}{}^{2}}}4_4\sigma _4\stackrel{~}{h}_\rho ^\rho \varphi +12k^2\varphi \right)+`$
$`+`$ $`\left[2kh_{\mu \nu }3k\gamma _{\mu \nu }\varphi \right]\stackrel{~}{\delta }={\displaystyle \frac{\widehat{\kappa }}{2}}T_{\mu \nu },`$
2) $`\mu 4`$-component,
$$_4(_\mu \stackrel{~}{h}^\nu h_{\mu \nu })3_4\sigma _\mu \varphi =\widehat{\kappa }T_{\mu 4},$$
(13)
which plays the role of a constraint,
3) $`44`$-component
$$\frac{1}{2}(^\mu ^\nu h_{\mu \nu }_\mu ^\mu \stackrel{~}{h})\frac{3}{2}_4\sigma _4\stackrel{~}{h}+6k^2\varphi =\frac{\widehat{\kappa }}{2}T_{44},$$
(14)
with $`T_{MN}`$ being the energy-momentum tensor of the matter and $`\stackrel{~}{h}=\gamma ^{\mu \nu }h_{\mu \nu }`$. In what follows, we will also use an auxiliary equation, which is obtained by multiplying the equation for $`44`$-component by 2 and subtracting it from the contracted equation for $`\mu \nu `$-component. This equation contains $`\stackrel{~}{h}`$ and $`\varphi `$ only and has the form:
$$\frac{^2\stackrel{~}{h}}{x_{}^{4}{}_{}{}^{2}}+2_4\sigma _4\stackrel{~}{h}8k^2\varphi +8k\varphi \stackrel{~}{\delta }+_\mu ^\mu \varphi =\frac{\widehat{\kappa }}{3}\left(T_\mu ^\mu 2T_{44}\right).$$
(15)
For example, if $`T_{MN}=0,`$ the physical degrees of freedom of the model can be extracted by the substitution
$$h_{\mu \nu }=b_{\mu \nu }+\gamma _{\mu \nu }(\sigma c)\varphi +\frac{1}{2k^2}\left(\sigma c+\frac{1}{2}+\frac{c}{2}e^{2\sigma }\right)_\mu _\nu \varphi .$$
(16)
with $`c=\frac{kR}{e^{2kR}1}`$. It turns out that the field $`b_{\mu \nu }(x^\mu ,y)`$ describes the massless graviton and massive Kaluza-Klein spin-2 fields, whereas $`\varphi (x)`$ describes a scalar field called the radion.
However, the situation is rather different, when there is matter on the branes. The cases, in which matter is located on the branes, i.e. the energy-momentum tensor is of the form $`T_{\mu \nu }=t_{\mu \nu }(x)\delta (y)`$ or $`T_{\mu \nu }=t_{\mu \nu }(x)\delta (yR)`$ (and $`T_{\mu 4}=0`$, $`T_{44}=0`$) were discussed in detail in . It is well known that if we live on the negative tension brane (at $`y=R`$), the contribution of the radion is $`e^{2kR}`$ times stronger than the contribution of the massless graviton . It means that in the case of the massless radion scalar gravity is realized on brane 2, and it is necessary to have a mechanism for generating the radion mass, for example, the Goldberger-Wise mechanism to make gravity in the zero mode approximation tensor, i.e. to suppress the contribution of the scalar radion component. But we will show below, that this is not the only solution for this problem.
## 3 A simple example
In the original formulation of the Randall-Sundrum model the mechanism of localization of fields is not taken into account. The branes are treated as infinitely thin objects, and it is assumed that energy-momentum tensor has $`\delta (y)`$-like profile in the extra coordinate (matter is located on the brane only). Nevertheless in some papers (see, for example, ) considering models with extra dimensions, energy-momentum tensors with dependence on the extra coordinates were used for constructing non-trivial background metric. One can also recall models with universal extra dimensions, in which physical fields can propagate in the extra dimension , and fat brane scenarios . Thus, the existence of energy-momentum tensors of this type seems to be quite reasonable. Here we consider some simple examples, which are sufficient to show attractive features of the non-delta-like localization of matter.
Let us choose as an example the following form of the energy-momentum tensor
$`T_{\mu \nu }`$ $`=`$ $`t_{\mu \nu }(x){\displaystyle \frac{e^{2\sigma }c}{R}},`$ (17)
$`T_{\mu 4}`$ $``$ $`0,`$ (18)
$`T_{44}`$ $`=`$ $`t_{\mu \nu }(x)\gamma ^{\mu \nu }e^{2\sigma }{\displaystyle \frac{c}{R}}\left(\sigma +e^{2kR}cϵc\right),`$ (19)
where $`ϵ`$ is some arbitrary constant, which will be defined later. Here $`t_{\mu \nu }(x)`$ is the energy-momentum tensor of matter on the brane, it depends on the four-dimensional coordinates $`x`$ only. One can see that the matter is localized near the negative tension brane at $`y=R`$. We would like to note that the energy-momentum tensor of the form (17), (18) and (19) satisfies the covariant energy conservation law, which has the form
$$^NT_{MN}=0,$$
(20)
and the function of localization in (17) is normalized to unity. Constructions similar to that used in this paper, but in the case of six dimensions, were utilized in . The explanation of such constructions depends on the method of localization, and this issue will not be discussed in this paper.
The substitution, which allows one to decouple the equations (2), (13), (14), (15) with $`T_{MN}`$ of the form (17), (18), (19) looks like
$`b_{\mu \nu }=u_{\mu \nu }+{\displaystyle \frac{1}{2k^2}}\left({\displaystyle \frac{ce^{2kR}}{ϵ}}e^{2\sigma }\left[{\displaystyle \frac{ce^{2kR}}{2ϵ}}+{\displaystyle \frac{1}{8ϵ}}\right]\sigma e^{2\sigma }{\displaystyle \frac{1}{2ϵ}}\right)_\mu _\nu \varphi ,`$ (21)
where $`b_{\mu \nu }`$ was defined in (16) and $`u_{\mu \nu }(x^\mu ,y)`$ describes the massless graviton and massive Kaluza-Klein spin-2 fields. Substituting (16) and (21) into equations (13), (14) and (15), we get:
$`_4(e^{2\sigma }(^\nu u_{\mu \nu }_\mu u))=0,`$ (22)
$`e^{4\sigma }(^\mu ^\nu u_{\mu \nu }\mathrm{}u)3_4\sigma _4(e^{2\sigma }u)+e^{2\sigma }{\displaystyle \frac{3ce^{2kR}}{ϵ}}\mathrm{}\varphi `$ (23)
$`{\displaystyle \frac{3}{ϵ}}e^{4\sigma }\left(\sigma +\stackrel{~}{c}e^{2kR}\right)\mathrm{}\varphi =\widehat{\kappa }t{\displaystyle \frac{c}{R}}e^{4\sigma }\left(\sigma +\stackrel{~}{c}e^{2kR}\right),`$
$`_4(e^{2\sigma }_4(e^{2\sigma }u))+{\displaystyle \frac{1}{ϵ}}\mathrm{}\varphi e^{2\sigma }\left(12\left(\sigma +\stackrel{~}{c}e^{2kR}\right)\right)=`$ (24)
$`={\displaystyle \frac{\widehat{\kappa }c}{3R}}te^{2\sigma }\left(12\left(\sigma +\stackrel{~}{c}e^{2kR}\right)\right),`$
where $`u=\eta ^{\mu \nu }u_{\mu \nu }`$, $`t=\eta ^{\mu \nu }t_{\mu \nu }`$, $`^\mu =\eta ^{\mu \nu }_\nu `$, $`\mathrm{}=\eta ^{\mu \nu }_\mu _\nu `$ and
$$\stackrel{~}{c}=c\left(1ϵe^{2kR}\right)$$
Let us consider the Fourier expansion of all terms of equation (24) with respect to coordinate $`y`$. Since the term with the derivative $`_4`$ has no zero mode, this equation implies that
$$\mathrm{}\varphi =\frac{\widehat{\kappa }ϵc}{3R}t,$$
(25)
$$_4(e^{2\sigma }u)=0.$$
(26)
Now let us consider $`\mu \nu `$-equation. It is well known that the field $`u_{\mu \nu }`$ in the presence of matter is a combination of zero and massive modes, whose eigenfunctions are orthogonal . In particular, the zero mode can be represented as $`u_{\mu \nu }^0=e^{2\sigma }\alpha _{\mu \nu }`$, where $`\alpha _{\mu \nu }`$ depends on $`x`$ only. It also means that with the help of the residual gauge transformations (11) it is possible to impose the gauge condition
$`^\nu u_{\mu \nu }^m=0,`$ (27)
$`u_\mu ^{m\mu }=0,`$
on the massive modes $`u_{\mu \nu }^m`$ and the de Donder gauge condition on the zero mode
$$^\nu \left(\alpha _{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\alpha \right)=0.$$
(28)
Having imposed this gauge, we are still left with residual gauge transformation
$$\xi _\mu =e^{2\sigma }ϵ_\mu (x),\mathrm{}ϵ_\mu =0.$$
(29)
The gauge transformations with $`\xi _\mu `$ satisfying these conditions are important for determining the number of degrees of freedom of the massless graviton. It follows from equation (23) that
$$\mathrm{}\alpha =2\frac{\widehat{\kappa }e^{2kR}c^2}{R}t.$$
(30)
Substituting (16) and (21) into equation (2) with condition (25) and passing to gauge (27), (28) we get
$`{\displaystyle \frac{1}{2}}\mathrm{}(\alpha _{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }\alpha )+{\displaystyle \frac{1}{2}}e^{2\sigma }\mathrm{}u_{\mu \nu }^m+{\displaystyle \frac{1}{2}}_4_4u_{\mu \nu }^m2k^2u_{\mu \nu }^m_4_4\sigma u_{\mu \nu }^m=`$ (31)
$`={\displaystyle \frac{\widehat{\kappa }c}{2R}}t_{\mu \nu }e^{2\sigma }+{\displaystyle \frac{\widehat{\kappa }c}{6R}}\left(2ce^{2kR}e^{2\sigma }\right)\left({\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\eta _{\mu \nu }\right)t`$
We are going to calculate the equations of motion in the zero mode approximation. Thus, we have to find an equation for the field $`\alpha _{\mu \nu }`$. If we multiply equation (31) by $`e^{2\sigma }`$, integrate it over $`y`$ and take into account the orthonormality condition for the wave functions of the modes , we get
$$\mathrm{}(\alpha _{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\alpha )=2\frac{\widehat{\kappa }c^2e^{2kR}}{R}t_{\mu \nu }.$$
(32)
The equation for massive modes takes the form
$`{\displaystyle \frac{1}{2}}e^{2\sigma }\mathrm{}u_{\mu \nu }^m+{\displaystyle \frac{1}{2}}_4_4u_{\mu \nu }^m2k^2u_{\mu \nu }^m_4_4\sigma u_{\mu \nu }^m=`$ (33)
$`={\displaystyle \frac{\widehat{\kappa }c}{2R}}t_{\mu \nu }\left(e^{2\sigma }2ce^{2kR}\right)+{\displaystyle \frac{\widehat{\kappa }c}{6R}}\left(2ce^{2kR}e^{2\sigma }\right)\left({\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\eta _{\mu \nu }\right)t`$
Now we are ready to find equations of motion for the zero mode part of $`h_{\mu \nu }`$. Since the matter is localized near the brane, we will calculate the equation for $`h_{\mu \nu }|_{y=R}`$. It is easy to see that $`h_{\mu \nu }|_{y=R}`$ (16), (21) does not satisfy the de Donder gauge condition. The residual gauge transformations (29) are not sufficient to pass to this gauge. But since we consider only the effective theory on brane 2 (at $`y=R`$), we can drop the condition $`\mathrm{}ϵ_\mu =0`$, which fixes the gauge for the field $`h_{\mu \nu }`$. Then we can pass to the de Donder gauge condition for the field $`h_{\mu \nu }`$ on brane 2 with the help of these gauge functions. Making these transformations (analogously to what was made in ), we get in the zero mode approximation
$`\left(h_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }h\right)|_{y=R}=`$ (34)
$`=e^{2kR}\left(\alpha _{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }\alpha \right)2e^{2kR}(kRc)\left(\eta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\right)\varphi .`$
An important point is that these equations are written in the coordinates $`\{x^\mu \}`$, which are Galilean on brane 1 (not on brane 2) and are inappropriate for studying physical effects on brane 2 (we recall that coordinates are called Galilean, if $`g_{\mu \nu }=diag(1,1,1,1)`$ ). It is necessary to pass to Galilean coordinates on brane 2 to get a correct result. This problem was discussed in detail in papers . In Galilean coordinates on brane 2 equation (34) looks like
$`\left(h_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }h\right)|_{y=R}=`$ (35)
$`=e^{2kR}\left(\alpha _{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }\alpha \right)2(kRc)\left(\eta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\right)\varphi ,`$
and equations (25), (32) take the form
$$\mathrm{}\varphi =\frac{\widehat{\kappa }ϵc}{3R}t,$$
(36)
$$\mathrm{}(\alpha _{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\alpha )=2\frac{\widehat{\kappa }c^2e^{4kR}}{R}t_{\mu \nu }.$$
(37)
Thus, we can get
$`\mathrm{}\left(h_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }h\right)|_{y=R}={\displaystyle \frac{2\widehat{\kappa }c^2e^{2kR}}{R}}\left(t_{\mu \nu }{\displaystyle \frac{ϵ}{3}}\left(\eta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\right)t\right),`$ (38)
where $`t_{\mu \nu }`$ is the energy-momentum tensor of matter in the coordinates, which are Galilean at $`y=R`$. One can see, that this equation coincides with equation for the fluctuations of metric in the linearized Brans-Dicke theory, which looks like
$`\mathrm{}\left(\delta g_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }\delta g\right)=16\pi G\left(t_{\mu \nu }{\displaystyle \frac{1}{2\omega +3}}\left(\eta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\right)t\right),`$ (39)
where $`\omega `$ is the BD-parameter and $`G`$ is the gravitational constant. Thus, comparing equations (38) and (39), we get
$$G=\widehat{G}\frac{2c^2e^{2kR}}{R},$$
(40)
$$\omega =\frac{3(1ϵ)}{2ϵ}.$$
(41)
It is easy to see, that in the the Randall-Sundrum model with the $`\delta `$-function-like localization of matter on the brane there is a factor $`e^{2kR}`$ instead of $`ϵ`$ in (38) (compare with the analogous formula in ).
From the recent experimental data (see, for example, ) we know, that $`\omega >3500`$. It means, that
$$ϵ<4,310^4$$
(42)
With this value of $`ϵ`$ there are already no problems with the radion field.
## 4 Localization of matter
In (17) the factor $`e^{2\sigma }`$, describing distribution of matter in the extra dimension, was utilized. However one can say that matter is spread in the whole bulk rather than localized on the brane (in analogy to the massless graviton whose wave function is $`e^{2\sigma }`$), and it is not evident, why we make all calculations for $`h_{\mu \nu }`$ at $`y=R`$. It seems that it is more reasonable to consider matter, which is confined to brane 2 much stronger, since the description of gravity on brane 2 by $`h_{\mu \nu }|_{y=R}`$ is more justified in this case. Let us consider a ”toy model” with energy-momentum tensor of the form
$$T_{\mu \nu }=t_{\mu \nu }(x)\frac{k(1+N)}{e^{2kR+2NkR}1}e^{2\sigma (1+N)},$$
(43)
where the function of localization is normalized to unity as well. Other components of the five-dimensional energy-momentum tensor are chosen to be
$$T_{\mu 4}=0,$$
(44)
$`T_{44}=t(x)(k+kN)e^{4\sigma }({\displaystyle \frac{c}{2NkR}}`$ (45)
$`{\displaystyle \frac{cϵe^{2NkR}}{e^{2kR+2NkR}1}}{\displaystyle \frac{k}{2Nk\left(e^{2kR+2NkR}1\right)}}e^{2N\sigma }).`$
Such a complicated form of (45) is caused by the requirements to satisfy the energy conservation law (20) and to get a weak coupling constant for the radion.
The substitution, which allows one to decouple and solve equations (13), (14) and $`\mu \nu `$-equation with (43), (44) and (45), has the following form:
$`b_{\mu \nu }=u_{\mu \nu }+({\displaystyle \frac{e^{2kR}\left(1e^{2NkR}\right)}{4k^2Nϵ\left(e^{2kR}1\right)}}{\displaystyle \frac{\left(e^{2kR}e^{2NkR}\right)}{8k^2Nϵ\left(e^{2kR}1\right)}}e^{2\sigma }+`$ (46)
$`+{\displaystyle \frac{e^{2NkR}}{4Nk\left(Nk+2k\right)ϵ}}e^{2\sigma 2N\sigma })_\mu _\nu \varphi ,`$
Analogously to what was made in Section 3, we can get
$$\mathrm{}\varphi =\frac{\widehat{\kappa }ϵ\left(k+kN\right)e^{2NkR}}{3\left(e^{2kR+2NkR}1\right)}t,$$
(47)
$$\mathrm{}(\alpha _{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\alpha )=\frac{\widehat{\kappa }k\left(k+Nk\right)\left(1e^{2NkR}\right)}{Nk\left(e^{2kR}1\right)\left(1e^{2kR2NkR}\right)}t_{\mu \nu }.$$
(48)
Using (46) and passing to Galilean coordinates at $`y=R`$, in the zero mode approximation we get
$`\mathrm{}\left(h_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }h\right)|_{y=R}=`$ (49)
$`={\displaystyle \frac{\widehat{\kappa }c\left(1+N\right)\left(1e^{2NkR}\right)}{NR\left(1e^{2kR2NkR}\right)}}\left(t_{\mu \nu }{\displaystyle \frac{2NkRϵ}{3\left(1e^{2NkR}\right)}}\left(\eta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{\mathrm{}}}\right)t\right).`$
One can see that in the limit $`N0`$ equation (49) passes into (38). The four-dimensional gravitational constant has the form
$$G=\widehat{G}\frac{c\left(1+N\right)\left(1e^{2NkR}\right)}{NR\left(1e^{2kR2NkR}\right)}$$
(50)
(compare with (39)). For a large $`N`$ we get
$$G\widehat{G}\frac{c}{R}.$$
(51)
A choice of relatively small $`ϵ`$ makes the contribution of the radion to be not in contradiction with the present-day experimental data.
## 5 Conclusion and final remarks
There may arise a question about generation of the new hierarchy instead of the one solved in the Randall-Sundrum model. Actually, the factor $`ϵ`$ in (19) is much smaller than the factor $`e^{2kR}`$ in (19). But for the energy-momentum tensor on brane 2 (at $`y=R`$) we get
$`T_{\mu \nu }|_{y=R}=t_{\mu \nu }{\displaystyle \frac{e^{2kR}c}{R}}kt_{\mu \nu },`$
$`T_{44}|_{y=R}=t{\displaystyle \frac{e^{4kR}c^2}{R}}\left(14.310^4\right)tk^2R\left(14.310^4\right),`$
and the correction $`ϵ`$ is not so small in comparison with the main value $`1`$ (analogous calculations can be made in the case (43) and (45) too). It seems to be quite reasonable. One should take into account that the energy-momentum tensors of the form (17), (18), (19) and (43), (44), (45) do not correspond to any real field action, and must be interpreted as ”toy models”. There also arises the question about the origin of such correction $`ϵ`$. Explicitly the term $`e^{4\sigma }`$ in (19) and (45) corresponds to the homogeneous solution of equation (20) (i.e. to the case of $`T_{\mu \nu }=0`$). Nevertheless such term can be taken into consideration, since its existence does not evidently contradict our assumptions. Surely it would be better to derive similar form of the energy-momentum tensor in a more natural way. For example, simple examples with scalar and gauge fields in the bulk in the Randall-Sundrum model were discussed in papers . Wave functions of these fields have their maximum values on brane 2, and these solutions can be treated as a basis for our choice of the energy-momentum tensor. But even in the simplest case of the scalar field only, there arises a nonzero $`T_{\mu 4}`$-component of the energy-momentum tensor, which does not satisfy the condition (18). Moreover, the wave functions are different not only for different fields, but even for different modes. Thus, formulas (17), (18), (19) and (43), (44), (45) must be used as ”toy models” for classical objects, but not for the fields (for example, one can choose $`t_{00}m\delta (\stackrel{}{x})`$, $`t_{i0}=0`$, $`t_{ij}=0`$ for a static point-like source).
In this paper we have discussed only one aspect of the radion field problem. Indeed, though the coupling constant of the radion to matter can be small enough to fit the experimental data, the size of extra dimension in the RS1 model should be stabilized in any case. It is well known that in this model gravity alone cannot stabilize the size of extra dimension due to the Casimir effect . But if matter is not localized exactly on the brane, i.e. it has wave functions in the extra dimension, it is possible that the contribution of matter fields to the Casimir force between two branes could stabilize the size of the extra dimension. Moreover, even in the case of classical stabilization with the help of scalar field (such as in ) the coupling constant of the massive radion to matter can vary essentially depending on the profile of matter in the extra dimension.
There is another question, which may arise in connection with our choice for matter distribution: why do we choose non-delta-like profile for matter, whereas the branes have delta-like profile? Of course, such profile of the branes is an idealization. But one can recall kink-like solution found in , which describes a delta-like brane in the thin brane limit. Thus in the case under consideration the RS1 setup can be regarded as a phenomenological model for describing the branes in the case when the thickness of branes is much smaller than the effective width of the wave functions of matter. Another advantage of the RS1 background is its relative simplicity, which allowed us to decouple and solve exactly equations of motion for the radion field and tensor modes at least in the zero mode approximation and to show possible effects which can be produced by a non-standard matter distribution. In a more realistic model the non-zero width of the branes should be definitely taken into account, but it is evident that effect analogous to that described above should exist in this case too.
It should be also noted that the longitudinal parts of substitutions (16), (21) and (46) contain some terms, which can lead (but it is not necessarily so) to the strong coupling effect, analogously to what happens in the DGP model . These terms are pure gauge from the four-dimensional point of view from the brane and are not dangerous in the linear order. Nevertheless it is necessary to consider the non-linear corrections to the equations of motion to discover the contribution of these longitudinal terms. But this problem deserves an additional detailed investigation.
Acknowledgments
The authors are grateful to Yu.A. Kubyshin for valuable discussions. The work was supported by the RFBR grants 04-02-16476 and 04-02-17448, by the grant UR.02.02.503 of the scientific program ”Universities of Russia”, and by the grant NS.1685.2003.2 of the Russian Federal Agency for Science.
|
warning/0506/hep-ph0506159.html
|
ar5iv
|
text
|
# FROM STATISTICAL PHYSICS TO HIGH ENERGY QCD
I discuss recent progress in understanding the high–energy evolution in QCD, which points towards a remarkable correspondence with the reaction–diffusion problem of statistical physics.
Recently, there has been significant progress in our understanding of QCD at high energy, based on the observations that (i) the gluon number fluctuations play an important role in the evolution towards saturation and the unitarity limit $`^{\mathrm{?},\mathrm{?}}`$ and (ii) the QCD evolution in the presence of fluctuations and saturation is in the same universality class as a series of problems in statistical physics, the prototype of which being the ‘reaction–diffusion’ problem $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. These observations have developed into a rich correspondence between high–energy QCD and modern problems in statistical physics, which relates topics of current research in both fields, and which has already allowed us to deduce some insightful results in QCD by properly translating the corresponding results from statistical physics $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
To put such theoretical developments into a specific physical context, let us consider $`\gamma ^{}`$–proton deep inelastic scattering (DIS) at high energy, or small Bjorken–$`x`$. We shall view this process in a special frame in which most of the total energy is carried by the proton, whose wavefunction is therefore highly evolved, while the virtual photon has just enough energy to dissociate long before splitting into a quark–antiquark pair in a colorless state (a ‘color dipole’), which then scatters off the gluon distribution in the proton. The transverse size $`r`$ of the dipole is controlled by the virtuality $`Q^2`$ of $`\gamma ^{}`$ (roughly, $`r^21/Q^2`$), so for $`Q^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ one can treat the dipole scattering in perturbation theory. But for sufficiently small $`x`$, even such a small dipole can see a high–density gluonic system, and thus undergo strong scattering.
Specifically, the small–$`x`$ gluons to which couple the projectile form a color glass condensate $`^\mathrm{?}`$ (CGC), i.e., a multigluonic state which is characterized by high quantum occupancy, of order $`1/\alpha _s`$, for transverse momenta $`k_{}`$ below the saturation momentum $`Q_s(x)`$, but which becomes rapidly dilute when increasing $`k_{}`$ above $`Q_s`$. The saturation scale rises very fast with the energy, $`Q_s^2(x)x^\lambda `$, and is the fundamental scale in QCD at high energy. In particular, the external dipole is strongly absorbed provided its size $`r`$ is large on the scale set by $`1/Q_s`$, whereas for $`r1/Q_s`$ one rather has weak scattering, or ‘color transparency’.
In turn, the small–$`x`$ gluons are produced through quantum evolution, i.e., through radiation from color sources (typically, other gluons) with larger values of $`x`$, whose internal dynamics is ‘frozen’ by Lorentz time dilation. Let $`\tau =\mathrm{ln}1/x`$ denote the rapidity ; it takes, roughly, a rapidity interval $`\mathrm{\Delta }\tau 1/\alpha _s`$ to emit one small–$`x`$ gluon; thus, in the high energy regime where $`\alpha _s\tau 1`$, the dipole meets with well developed gluon cascades, as illustrated in Fig. 1. Three types of processes can be distinguished in Fig. 1, which for more clarity are singled out in Fig. 2.
The first process, Fig. 2.a, represents one step in the standard BFKL evolution $`^\mathrm{?}`$; by iterating this step, one generates gluon ladders which are resummed in the solution to the BFKL equation $`^\mathrm{?}`$. However, by itself, the latter is well known to suffer from conceptual difficulties in the high energy limit : (i) The BFKL estimate for the dipole scattering amplitude $`T_\tau (r)`$ grows exponentially with $`\tau `$ (i.e., like a power of the energy), and thus eventually violates the unitarity bound $`T_\tau (r)1`$. (The upper limit $`T_\tau =1`$ corresponds to the black disk limit, in which the dipole is totally absorbed by the target.) (ii) The BFKL ladder is not protected from deviations towards the non–perturbative domain at low transverse momenta $`k_{}^2<\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ (‘infrared diffusion’). With increasing energy, the BFKL solution receives larger and larger contributions from such soft intermediate gluons, and thus becomes unreliable.
These ‘small–$`x`$ problems’ of the BFKL equation are both cured by the gluon recombination processes $`(n2)`$ illustrated in Fig. 2.b which are important at high energy, when the gluon density in the target is large, and lead to gluon saturation and the formation of the CGC. Such processes are included in the Balitsky–JIMWLK equation $`^\mathrm{?}`$, a non–linear, functional, generalization of the BFKL evolution which describes the approach towards gluon saturation in the target and preserves the unitarity bound in the evolution of the scattering amplitudes.
However, the Balitsky–JIMWLK equation misses $`^\mathrm{?}`$ the process in Fig. 2.c — the $`2n`$ gluon splitting — which describes the bremsstrahlung of additional small–$`x`$ gluons in one step of the evolution. By itself, this process is important in the dilute regime, where it leads to the construction of higher–point gluon correlation functions from the dominant 2–point function. But once generated, the $`n`$–point functions with $`n>2`$ are rapidly amplified by the subsequent BFKL evolution (the faster the larger is $`n`$) and then they play an important role in the non–linear dynamics leading to saturation. Thus, such splitting processes are in fact important for the evolution towards high gluon density, as originally observed in numerical simulations $`^\mathrm{?}`$ of Mueller’s ‘color dipole picture’ $`^\mathrm{?}`$, and more recently reiterated in the analysis in Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
Equations including both merging and splitting in the limit where the number of colors $`N_c`$ is large have recently became available $`^\mathrm{?}`$ (see also Refs. $`^{\mathrm{?},\mathrm{?}}`$), but their general solutions have not yet been investigated (except under some additional approximations $`^{\mathrm{?},\mathrm{?}}`$). Still, as we shall argue now, the asymptotic behaviour of the corresponding solutions — where by ‘asymptotic’ we mean both the high–energy limit $`\tau \mathrm{}`$ and the weak coupling limit $`\alpha _s0`$ — can be a priori deduced from universality considerations, by exploiting the correspondence between high–energy QCD and the reaction–diffusion problem of statistical physics $`^\mathrm{?}`$.
To that aim, it is convenient to rely on the event–by–event description $`^\mathrm{?}`$ of the scattering between the external dipole and the hadronic target (cf. Fig. 1) and to use the large–$`N_c`$ approximation to replace the gluons in the target wavefunction by color dipoles $`^\mathrm{?}`$. Then, the dipole–target scattering amplitude corresponding to a given event can be estimated as
$`T_\tau (r,b)\alpha _s^2f_\tau (r,b),`$ (1)
where $`\alpha _s^2`$ is the scattering amplitude between two dipoles with comparable sizes and nearby impact parameters, and $`f_\tau (r,b)`$ is the occupation number for target dipoles with size $`r`$ at impact parameter $`b`$, and is a discrete quantity: $`f=0,1,2,\mathrm{}`$. Thus, in a given event, the scattering amplitude is a multiple integer of $`\alpha _s^2`$.
In this dipole language, the $`24`$ gluon splitting depicted in Fig. 2.c is tantamount to $`12`$ dipole splitting, and generates fluctuations in the dipole occupation number and hence in the scattering amplitude. Thus, the evolution of the amplitude $`T_\tau (r,b)`$ with increasing $`\tau `$ represents a stochastic process characterized by an expectation value $`T(r,b)_\tau \alpha _s^2f(r,b)_\tau `$, and also by fluctuations $`\delta T\alpha _s^2\delta f\sqrt{\alpha _s^2T}`$ (we have used the fact that $`\delta f\sqrt{f}`$ for fluctuations in the particle number). Clearly, such fluctuations are relatively important (in the sense that $`\delta T>T`$) only in the very dilute regime where $`f<1`$, or $`T<\alpha _s^2`$.
Eq. (1) applies so long as the scattering is weak, $`T1`$, but by extrapolation it shows that the unitarity corrections are expected to be important when the dipole occupation factor becomes of order $`1/\alpha _s^2`$. Consider first the formal limit $`\alpha _s^20`$, in which the maximal occupation number $`N1/\alpha _s^2`$ becomes arbitrarily large. Then one can neglect the particle number fluctuations and follow the evolution of the scattering amplitude in the mean field approximation (MFA). This is described by the Balitsky–Kovchegov equation $`^\mathrm{?}`$, a non–linear version of the BFKL equation which, as shown in Ref. $`^\mathrm{?}`$, lies in the same universality class as the Fisher–Kolmogorov–Petrovsky–Piscounov (FKPP) equation (the MFA for the reaction–diffusion process and related phenomena in biology, chemistry, astrophysics, etc; see $`^\mathrm{?}`$ for recent reviews and more references). The FKPP equation reads, schematically,
$`_\tau T(\rho ,\tau )=_\rho ^2T(\rho ,\tau )+T(\rho ,\tau )T^2(\rho ,\tau ),`$ (2)
in notations appropriate for the dipole scattering problem: $`T(\rho ,\tau )T(r)_\tau `$ and $`\rho \mathrm{ln}(r_0^2/r^2)`$, with $`r_0`$ a scale introduced by the initial conditions at low energy. Note that weak scattering ($`T1`$) corresponds to small dipole sizes ($`r1/Q_s`$), and thus to large values of $`\rho `$. The three terms on the r.h.s. of Eq. (2) describe, respectively, diffusion, growth and recombination. The first two among them represent (an approximate version of) the BFKL dynamics, while the latter is the non–linear term which describes multiple scattering and thus ensures that the evolution is consistent with the unitarity bound $`T1`$.
Specifically, the solution $`T_\tau (\rho )`$ to Eq. (2) is a front which interpolates between two fixed points : the saturation fixed point $`T=1`$ at $`\rho \mathrm{}`$ and the unstable fixed point $`T=0`$ at $`\rho \mathrm{}`$ (see Fig. 3). The position of the front, which marks the transition between strong scattering ($`T1`$) and, respectively, weak scattering ($`T1`$), defines the saturation scale : $`\rho _s(\tau )\mathrm{ln}(r_0^2Q_s^2(\tau ))`$. With increasing $`\tau `$, the front moves towards larger values of $`\rho `$, as illustrated in Fig. 3. Note that the dominant mechanism for propagation is the BFKL growth in the tail of the distribution at large $`\rho `$ : the front is pulled by the rapid growth of a small perturbation around the unstable state. In view of that, the velocity of the front $`\lambda d\rho _s/d\tau `$ is fully determined by the linearized version of Eq. (2), which describes the dynamics in the tail. Specifically, by solving the BFKL equation one finds $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ that, for $`\rho >\rho _s(\tau )`$ and sufficiently large $`\tau `$,
$`T_\tau (\rho )\mathrm{e}^{\omega \overline{\alpha }_s\tau }\mathrm{e}^{\gamma \rho }=\mathrm{e}^{\gamma (\rho \rho _s(\tau ))},\rho _s(\tau )c\overline{\alpha }_s\tau ,`$ (3)
where $`\overline{\alpha }_s=\alpha _sN_c/\pi `$, $`\gamma =0.63..`$, and $`c\omega /\gamma =4.88..`$. From Eq. (3) one can immediately identify the velocity of the front in the MFA as $`\lambda _0=c\overline{\alpha }_s`$. Since $`Q_s^2(\tau )Q_0^2\mathrm{e}^{\lambda _0\tau }`$, it is furthermore clear that $`\lambda _0`$ plays also the role of the saturation exponent (here, in the MFA).
What is the validity of the mean field approximation ? We have earlier argued that the gluon splitting processes (cf. Fig. 2.c) responsible for dipole number fluctuations should play an important role in the dilute regime. This is further supported by the above considerations on the pulled nature of the front: Since the propagation of the front is driven by the dynamics in its tail where the fluctuations are a priori important, the front properties should be strongly sensitive to fluctuations. This is indeed known to be the case for the corresponding problem in statistical physics $`^{\mathrm{?},\mathrm{?}}`$, as it can be understood from the following, qualitative, argument:
Consider a particular realization of the stochastic evolution of the target, and the corresponding scattering amplitude which is discrete (in steps of $`\alpha _s^2`$). Because of discreteness, the microscopic front looks like a histogram and thus is necessarily compact : for any $`\tau `$, there is only a finite number of bins in $`\rho `$ ahead of $`\rho _s`$ where $`T_\tau `$ is non–zero (see Fig. 4). This property has important consequences for the propagation of the front. In the empty bins on the right of the tip of the front, the local, BFKL–like, growth is not possible anymore (as this would require a seed). Thus, the only way for the front to progress there is via diffusion, i.e., via radiation from the occupied bins at $`\rho <\rho _{\mathrm{tip}}`$ (compare in that respect Figs. 3 and 4). But since diffusion is less effective than the local growth, we expect the velocity of the microscopic front (i.e., the saturation exponent) to be reduced as compared to the respective prediction of the MFA.
To obtain an estimate for this effect, we shall rely again on the universality of the asymptotic ($`\tau \mathrm{}`$ and $`N1/\alpha _s^21`$) behaviour$`^\mathrm{?}`$. Namely, from the experience with the reaction–diffusion process and related problems in statistical physics $`^{\mathrm{?},\mathrm{?}}`$, one knows indeed that the dominant behaviour for large evolution ‘time’ and large (but finite) occupancy $`N1`$ is independent of the details of the microscopic dynamics, and thus is the same for all the processes whose mean field limit ($`N\mathrm{}`$) is governed by the FKPP equation (2). In particular, the dominant contribution to the correction $`\lambda _N\lambda _0`$ to the front velocity is known to be universal, and can be obtained through the following, intuitive, argument, due to Brunet and Derrida $`^\mathrm{?}`$:
For a given microscopic front and $`N1`$, the MFA should work reasonably well everywhere except in the vicinity of the tip of the front, where the occupation number $`f`$ becomes of order one (corresponding to $`T\alpha _s^2`$ in the QCD problem) and thus the linear growth term becomes ineffective. Accordingly, Brunet and Derrida suggested a modified version of the FKPP equation (2) in which the ‘BFKL–like’ growth term is switched off when $`T<\alpha _s^2`$ :
$`_\tau T(\rho ,\tau )=_\rho ^2T+\mathrm{\Theta }\left(T\alpha _s^2\right)T(1T).`$ (4)
By solving this equation in the linear regime, they have obtained the first correction to the front velocity as compared to the MFA (in notations adapted to QCD; see Ref. $`^\mathrm{?}`$ for details):
$`\lambda \overline{\alpha }_s\left[c{\displaystyle \frac{\kappa }{\mathrm{ln}^2(1/\alpha _s^2)}}+𝒪\left(1/\mathrm{ln}^3\alpha _s^2\right)\right],`$ (5)
where the numbers $`c4.88`$ and $`\kappa 150`$ are fully determined by the linear (BFKL) equation. In QCD, the same result has been first obtained through a different but related argument by Mueller and Shoshi $`^\mathrm{?}`$. Note the extremely slow convergence of this result towards its mean field limit: the corrective term vanishes only logarithmically with decreasing $`1/N=\alpha _s^2`$, rather than the power–like suppression usually found for the effects of fluctuations. This is related to the high sensitivity of the pulled fronts to fluctuations, as alluded to above. Such a slow convergence, together with the relatively large value of the numerical factor $`\kappa `$, make that the above estimate for $`\lambda `$, although exact for asymptotically small $`\alpha _s^2`$, is in fact useless for practical applications. To understand the subasymptotic corrections and, more generally, the behaviour of the saturation momentum and of the scattering amplitudes for realistic values of $`\tau `$ and $`\alpha _s`$, one needs to solve the actual evolution equations of QCD $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, a program which is currently under way $`^\mathrm{?}`$.
|
warning/0506/hep-th0506180.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
One of the outstanding questions facing string theory is how to describe a cosmological singularity like the big bang. For recent cosmological scenarios where the big bang singularity plays a crucial role, see for instance . What prior work has taught us is that perturbative string theory breaks down on many toy model space-times that include space-like and light-like curvature singularities . See for some related work. To capture the physics of the singularity, a complete non-perturbative description of string theory appears to be necessary.<sup>4</sup><sup>4</sup>4See, however, the very recent paper , which claims that certain space-like singularities are replaced by a tachyon condensate phase within perturbative string theory. In this work, we will present a particularly clean example of a cosmological singularity that admits a holographic dual description via Matrix Theory (for reviews, see for instance ). Prior examples of holographic descriptions, in the sense of AdS/CFT, appear in . See for some other ideas relating Matrix theory and cosmology.
The backgrounds that we will consider are linear dilaton backgrounds which are key ingredients in some of the oldest exact solutions of string theory . The dilaton, $`\varphi `$, is identified with a direction in space-time. If the direction is time-like, the solution is cosmological and non-supersymmetric. By definition, we lose perturbative control over these backgrounds when the string coupling
$$g_s=e^\varphi $$
(1)
becomes large.
In this work, we want to consider a simple variant of these cosmologies where we choose light-cone coordinates in space-time and identify,
$$\varphi =QX^+,$$
(2)
where $`Q`$ is a constant. This kind of dilaton profile appears as an ingredient in many supergravity solutions like some plane-wave backgrounds. To construct a solution of string theory, we also need to specify a $`10`$-dimensional space-time metric. This metric could describe some non-trivial compactification. For simplicity, we will take flat Minkowski space with coordinates $`X^\mu =(X^+,X^{},X^i)`$ and metric,
$$ds_{10}^2=2dX^+dX^{}+\underset{i}{}(dX^i)^2$$
(3)
as our $`10`$-dimensional string metric.
This background is a remarkably simple, time-dependent string solution. Note that the parameter $`Q`$ appearing in $`(\text{2})`$ can be scaled to any non-zero value using the boost symmetry
$$X^+\alpha X^+,X^{}\alpha ^1X^{},$$
where $`\alpha `$ is non-zero. To see that flat space is still a string solution in the presence of this linear dilaton, we only need to note that a light-like linear dilaton (unlike a space-like or time-like linear dilaton) makes no contribution to the conformal anomaly.
From the perspective of string frame, the only time-dependence appears in the coupling constant. However, the corresponding Einstein frame metric is given by
$$ds_E^2=e^{QX^+/2}ds_{10}^2.$$
(4)
Viewed in Einstein frame, this space-time originates at a big bang as $`X^+\mathrm{}`$ since the scale factor goes to zero. In the following section, we will study this solution as a model for a big bang, and describe perturbative string quantization in this background. In section 3, we derive the Matrix description of this background which involves Matrix strings propagating on a time-dependent world-sheet. The world-sheet is described by the two-dimensional metric
$$ds^2=e^{2\eta }\left(d\eta ^2+dx^2\right)$$
(5)
with $`xx+2\pi `$. This metric describes the future quadrant of the Milne orbifold, which can be thought of as flat space with a boost identification. The curvature singularity of the metric $`(\text{4})`$ corresponds via Matrix theory to the Milne singularity at $`\eta =\mathrm{}`$, where the $`x`$-circle shrinks to zero size. This is the “big bang”. Time evolution from the big bang to the asymptotic regime corresponds in the Matrix description to renormalization group flow from the UV to the IR. The physics near the big bang is described by weakly-coupled Yang-Mills theory. In section 3.2, we argue that our Matrix description remains decoupled from gravity even at the singularity. At late times, the matrix degrees of freedom re-organize themselves into weakly-coupled strings, and a conventional space-time picture emerges.
It is worth noting that the Milne orbifold has been studied as a space-time background for closed string propagation. What can be concluded from this work is that perturbative string theory breaks down because of large gravitational backreaction from the singularity. In our case, the Milne orbifold appears as the Matrix string world-sheet. Matrix string theory should capture the non-perturbative physics of the space-time singularity.
In the final section, we mention a few of the possible generalizations of the light-like linear dilaton background.
## 2 The Light-like Linear Dilaton
### 2.1 The light-like linear dilaton as a big bang cosmology
We begin by considering the light-like linear dilaton as a background of type IIA string theory. This defines an exact CFT that describes string propagation in flat space-time with a varying string coupling given by
$$g_s=e^{QX^+}.$$
(6)
The space-time theory is free at late times ($`X^+\mathrm{}`$), and strongly coupled at early times ($`X^+\mathrm{}`$).
This background preserves one-half of the $`32`$ flat space supersymmetries. To see this, we need to check the supersymmetry variations of the gravitino and dilatino. Only the dilatino, $`\lambda `$, feels the presence of this linear dilaton background via the term in its supersymmetry variation,
$$\delta \lambda \mathrm{\Gamma }^+_+\varphi ϵ,$$
(7)
where $`ϵ`$ is the supersymmetry parameter. However, there are $`16`$ solutions to the condition
$$\mathrm{\Gamma }^+ϵ=0,$$
(8)
so half the supersymmetry is preserved. This is rather crucial since as $`X^+\mathrm{}`$, we have enough control from supersymmetry to determine a good strong coupling description. The spectrum in the weak coupling regime where $`X^+\mathrm{}`$ is determined from the perturbative string quantization to be described in section 2.2.
As $`g_s`$ becomes large, we expect this background to lift to a solution of M-theory. The $`11`$-dimensional metric
$$ds^2=e^{2QX^+/3}ds_{10}^2+e^{4QX^+/3}(dY)^2,$$
(9)
with $`Y`$ the eleventh direction, governs the strong coupling limit of this background. We should check that the background defined by $`(\text{9})`$ is not trivial. We define an orthonormal basis of 1-forms
$$e^i=e^{QX^+/3}dX^i,e^+=e^{QX^+/3}dX^+,e^{}=e^{QX^+/3}dX^{},e^y=e^{2QX^+/3}dY$$
(10)
with respect to which the metric takes the canonical form
$$ds_{11}^2=2e^+e^{}+(e^i)^2+(e^y)^2.$$
(11)
Up to symmetry, the corresponding spin connection has non-vanishing components
$$\omega _{i+}=\frac{Q}{3}e^{QX^+/3}e^i,\omega _{y+}=\frac{2Q}{3}e^{QX^+/3}e^y,\omega _+=\frac{Q}{3}e^{QX^+/3}e^+.$$
(12)
The non-vanishing curvature $`2`$-forms are
$`R_{+i}`$ $`=`$ $`{\displaystyle \frac{Q^2}{9}}e^{2QX^+/3}e^+e^i,`$
$`R_{y+}`$ $`=`$ $`{\displaystyle \frac{8Q^2}{9}}e^{2QX^+/3}e^+e^y,`$ (13)
with respect to the orthonormal basis $`(\text{10})`$. The non-vanishing components of the Riemann tensor in a coordinate basis are (up to symmetry)
$`R_{+i+i}`$ $`=`$ $`{\displaystyle \frac{Q^2}{9}}e^{2QX^+/3},`$ (14)
$`R_{+y+y}`$ $`=`$ $`{\displaystyle \frac{8Q^2}{9}}e^{4QX^+/3}.`$ (15)
It is easy to see from these equations that the Ricci tensor vanishes, as it should for a purely gravitational M-theory solution.
It might appear that the $`11`$-dimensional metric $`(\text{9})`$ has a ‘singularity’ at both $`X^++\mathrm{}`$ and $`X^+\mathrm{}`$ since in both limits some metric components go to zero. The difference between these two limits is, however, that the $`X^+\mathrm{}`$ singularity occurs at finite geodesic distance, while the $`X^+\mathrm{}`$ singularity is at infinite distance.
The presence of the finite distance $`X^+\mathrm{}`$ singularity implies that the space-time is geodesically incomplete. Namely, some geodesics terminate at finite affine parameter. This is most easily seen for the lines $`X^{}=\mathrm{const}.`$, $`X^i=\mathrm{const}.`$, which are geodesics. The geodesic equation in this case is
$$\frac{d^2X^+}{d\lambda ^2}+\mathrm{\Gamma }_{++}^+\left(\frac{dX^+}{d\lambda }\right)^2=0,$$
(16)
where $`\mathrm{\Gamma }_{++}^+=2Q/3`$. This can be integrated to give
$$e^{\frac{2}{3}QX^+}\left(\frac{dX^+}{d\lambda }\right)=\mathrm{const}.$$
(17)
and hence the affine parameter is (up to an affine transformation)
$$\lambda =e^{\frac{2}{3}QX^+}.$$
(18)
We thus find that the point $`X^+\mathrm{}`$ corresponds to $`\lambda =0`$, and hence it has finite affine distance to all points in the interior. Note that the other ‘singularity’ at $`X^+\mathrm{}`$ is indeed at infinite affine parameter, $`\lambda =\mathrm{}`$, so it represents an asymptotic region in which the eleventh dimension happens to curl up to zero size.
One can write the metric in terms of the affine parameter $`\lambda `$ for $`\lambda >0`$ as
$$ds^2=\frac{3}{Q}d\lambda dX^{}+\lambda ds_8^2+\lambda ^2dy^2.$$
(19)
In terms of these coordinates, the non-vanishing components of the Riemann tensor are
$`R_{\lambda i\lambda i}`$ $`=`$ $`{\displaystyle \frac{1}{4\lambda }},`$
$`R_{\lambda y\lambda y}`$ $`=`$ $`{\displaystyle \frac{2}{\lambda ^4}},`$ (20)
which clearly shows that there is a curvature singularity at $`\lambda =0`$, where an inertial observer experiences divergent tidal forces.
It does not make sense to consider the metric (19) for $`\lambda <0`$ because the signature of the eight transverse dimensions changes sign. To extend to $`\lambda <0`$ in a sensible way, one might try replacing the $`\lambda `$ in front of $`ds_8^2`$ by its absolute value $`|\lambda |`$. However, this extension is ad hoc without some additional input beyond general relativity about how to treat the curvature singularity. What we can conclude is that there is truly a singularity in the classical gravity description of the light-like linear dilaton background.
In fact, this same conclusion also applies to the 10-d desciption in Einstein frame: namely, there exists a singularity at finite geodesic distance. To see this, rewrite the Einstein metric (4) in terms of its affine parameter
$$u=e^{\frac{1}{2}QX^+}$$
(21)
and a new coordinate $`v=X^{}`$:
$$ds_E^2=\frac{4}{Q}dudv+u\underset{i}{}(dX^i)^2.$$
(22)
Defining an orthonormal basis
$$e^i=u^{1/2}dX^i,e^u=\frac{2}{Q}du,e^v=dv,$$
(23)
we find the following non-vanishing components of the spin connection:
$$\omega ^i{}_{u}{}^{}=\frac{Q}{4u}e^i,$$
(24)
and of the curvature two form:
$$R^i{}_{u}{}^{}=\frac{Q^2}{16u^2}e^ie^u.$$
(25)
In a coordinate basis, we find
$$R_{iuiu}=\frac{1}{4u},$$
(26)
which is indeed singular at $`u=0`$.
In Einstein frame, the Ricci tensor is non-zero:
$$R_{uu}=\frac{2}{u^2}.$$
(27)
This non-vanishing Ricci tensor is supported by the dilaton, which, unlike in the string frame, has a non-zero stress-energy tensor and thus contributes to Einstein’s equations. We have
$$\varphi =QX^+=2\mathrm{log}u,$$
(28)
and thus
$$T_{uu}=\frac{1}{2}(_u\varphi )^2=\frac{2}{u^2}.$$
(29)
This can be interpreted as follows: the singular nature of the $`R_{\lambda y\lambda y}`$ component of the 11-d Riemann tensor (20) is transferred to the stress-energy tensor of the dilaton, and hence by Einstein’s equations to the Ricci tensor (27).
### 2.2 Perturbative string theory
We now describe some of the properties of the light-like linear dilaton solution in perturbative string theory. We start with the bosonic string in $`D`$ dimensions. The world-sheet fields $`X^\mu `$ are free, but the time-dependence of the dilaton is reflected in a modified world-sheet stress tensor,
$$T(z)=X_iX^i+2X^+X^{}Q^2X^+.$$
(30)
The central charge is unmodified,
$$c=D,$$
(31)
and $`Q`$ is a free parameter.
Since the world-sheet theory is free and the string coupling is small at late times, it is not difficult to construct the physical states. As in flat space-time with constant dilaton, states are labeled by momentum $`p_\mu `$ and an oscillator contribution. The corresponding vertex operators have the form
$$V=e^{ip_\mu X^\mu }P_N(X^\mu ,\overline{}X^\mu ,\mathrm{}),$$
(32)
where $`P_N`$ is a polynomial in derivatives of the world-sheet fields $`X^\mu `$ of total (left and right) scaling dimension $`N`$. We take the zero mode part of the vertex operator to be a plane wave. Physical states correspond to Virasoro primaries of the form $`(\text{32})`$ with scaling dimension one. From $`(\text{30})`$, it follows that the scaling dimension of $`V`$ is
$$L_0=\frac{1}{4}p_i^2\frac{1}{2}p^+(p^{}+iQ)+N.$$
(33)
The non-standard contribution of $`p^+`$ to the scaling dimension is easy to understand. In string theory, the zero mode part of vertex operators for the emission of string modes have the general form
$$V=g_s\mathrm{\Psi },$$
(34)
where $`\mathrm{\Psi }`$ is the wavefunction of the state. Usually, the factor of $`g_s`$ in (34) can be neglected since it is constant, but here according to $`(\text{6})`$ it is time-dependent and therefore needs to be retained. The vertex operator (32) corresponds to the wavefunction
$$\mathrm{\Psi }(\stackrel{}{X},X^+,X^{})=e^{i\stackrel{}{p}\stackrel{}{X}ip^+X^{}i(p^{}+iQ)X^+}.$$
(35)
Thus, the light-cone energy is
$$E^{}=p^{}+iQ$$
(36)
and $`(\text{33})`$ reads
$$L_0=\frac{1}{4}\left(\stackrel{}{p}^22p^+E^{}\right)+N.$$
(37)
The mass shell condition then becomes
$$m_{\mathrm{eff}}^22p^+E^{}\stackrel{}{p}^2=4(N1).$$
(38)
The classical evolution of fields in this background is easy to describe. For example, consider a scalar field $`T`$ with mass $`m`$ in the light-like linear dilaton background. The Lagrangian is proportional to
$$=\frac{1}{2}e^{2QX^+}(2_+T_{}T_iT_iTm^2T^2).$$
(39)
The equation of motion of $`T`$ is
$$(2_+_{}_i_i+2Q_{}+m^2)T=0.$$
(40)
A basis of solutions is given by
$$T(X^+,X^{},\stackrel{}{X})=e^{QX^+}e^{ip^+X^{}iE^{}X^++i\stackrel{}{k}\stackrel{}{X}},$$
(41)
with
$$2p^+E^{}+\stackrel{}{p}^2+m^2=0.$$
(42)
### 2.3 Light-cone string field theory
The calculation of the perturbative string amplitudes of the light-like linear dilaton background becomes particularly simple in the light-cone gauge. Of course, given that the dilaton itself picks a preferred light-cone direction, one does not even break Lorentz invariance by making the usual gauge choice $`X^+=p^+\tau `$ on the world-sheet. The world-sheet theory for the transverse coordinates is completely identical to that of flat space superstring theory.
As is well known from the old literature, one can represent a perturbative string amplitude in terms of a sum over light-cone diagrams. The contribution of each diagram is expressed as an integral over the positions $`\tau _i`$ of the joining and splitting operators on the world-sheet. For a genus $`g`$ contribution to a $`n`$-string scattering amplitude the number of these vertex operators is $`2g2+n`$. The effect of the linear dilaton is that the coupling constants now becomes a function of the light-cone coordinate $`\tau `$ on the world-sheet. Specifically, every joining/splitting operator gets multiplied by $`e^{Qp^+\tau _i}`$. Hence the overall amplitude, before integrating over the $`\tau _i`$, gets multiplied by
$$\underset{i=1}{\overset{2g2+n}{}}e^{Qp^+\tau _i}e^{(2g2+n)Qp^+\tau _{}}.$$
(43)
Here $`\tau _{}`$ is the average of the insertion points $`\tau _i`$. Because the world-sheet theory is translation invariant in $`\tau `$, the rest of the integrand only depends on the relative differences of the positions $`\tau _i`$ of the joining/splitting vertices. This fact can be exploited by separating the integration over the $`\tau _i`$ into the integral over the relative positions multiplied by the integral over $`\tau _{}`$. The integral over the relative positions precisely gives the usual amplitude in flat space. We thus obtain the following simple relation between the string amplitudes in the light-like linear dilaton background and the corresponding flat space amplitudes:
$$A^{g,n}=A_{\mathrm{flat}}^{g,n}_{\mathrm{}}^+\mathrm{}𝑑\tau _{}e^{(2g2+n)Qp^+\tau _{}}.$$
(44)
Clearly, the integral diverges even before summing over the genus. Hence, one clearly has to introduce a cut-off in the $`\tau ^{}`$ integral, keeping it away from $`\tau =\mathrm{}`$. But even when we take $`\tau _{}>\tau _c`$ one finds that the effective coupling $`g_s^{eff}e^{Qp^+\tau _c}`$ can become large when $`\tau _c`$ is negative. Therefore another description is required in this region. We will provide such a description in the next section.
## 3 Matrix String Description
We will begin our discussion of Matrix theory by taking the flat space Matrix string action and inserting the time-dependent string coupling, $`g_s=\mathrm{exp}(QX^+)`$. This leads directly to supersymmetric Yang-Mills on the Milne orbifold as the Matrix description of the light-like linear dilaton space-time. In section 3.1, we will provide an independent derivation leading to this same conclusion. This derivation will allow us to describe the regime of validity of the Matrix description.
Matrix string theory is described by a $`(1+1)`$-d super-Yang-Mills (SYM) theory with 16 supercharges. The action follows from dimensional reduction of $`(9+1)`$-d SYM theory. It contains eight matrix-valued fields $`X^i`$ representing the transverse bosonic coordinates, as well as eight matrix-valued spinor coordinates $`\mathrm{\Theta }^a`$. The action is
$$S=\frac{1}{2\pi \mathrm{}_s^2}\mathrm{tr}\left(\frac{1}{2}(D_\mu X^i)^2+\theta ^TD\text{/}\theta +g_s^2\mathrm{}_s^4\pi ^2F_{\mu \nu }^2\frac{1}{4\pi ^2g_s^2\mathrm{}_s^4}[X^i,X^j]^2+\frac{1}{2\pi g_s\mathrm{}_s^2}\theta ^T\gamma _i[X^i,\theta ]\right).$$
(45)
The metric on the world-sheet is flat, i.e. $`\eta _{\mu \nu }=\text{diag}(1,1)`$, and the spatial coordinate $`\sigma `$ on the world-sheet has a fixed periodicity equal to $`2\pi \mathrm{}_s`$. Notice that the Yang-Mills coupling constant, which is dimensionful in $`(1+1)`$ dimensions, is here identified with the inverse product of the string length and the string coupling,
$$g_{YM}\frac{1}{g_s\mathrm{}_s}.$$
(46)
In the IR limit the SYM theory become strongly coupled and, as shown in , reduces to the perturbative description of the type IIA superstring. In the UV, however, the SYM theory is weakly coupled.
In the light-cone gauge the world-sheet time coordinate is proportional to the space-time null coordinate $`X^+`$. We can thus describe the light-like linear dilaton background in a simple way by allowing the string coupling to depend on the world-sheet time $`\tau `$ via a relation like $`g_s=e^{Q\tau }.`$ In section 3.1, we will determine the precise proportionality constant between $`X^+`$ and $`\tau `$ leading to $`(\text{83})`$.
It thus appears that we are dealing with a SYM theory with a time-dependent coupling constant. However, there is another way to view this result; namely, by changing the geometry on the world-sheet. Unlike the usual string action, the matrix string action is not conformally invariant: rescaling the metric by a function $`f(\tau )^2`$ changes the terms involving the coupling $`g_s`$ in such a way that $`g_s`$ gets multiplied by $`f(\tau )^1`$. Thus, we conclude that the matrix string description of the light-like linear dilaton background is given by $`(1+1)`$-d SYM theory with fixed coupling, but on a world sheet with geometry
$$ds^2=e^{2Q\tau }(d\tau ^2+d\sigma ^2).$$
(47)
In fact, this metric is flat since it reduces to the usual Minkowski metric $`ds^2=2d\xi ^+d\xi ^{}`$ through the substitution
$$\xi ^\pm =\frac{1}{\sqrt{2}Q}e^{Q(\tau \pm \sigma )}.$$
(48)
However, since the coordinate $`\sigma `$ is periodic modulo $`2\pi \mathrm{}_s`$, the $`(1+1)`$-d Minkowski space described by the coordinates $`(\xi ^+,\xi ^{})`$ turns into the Milne orbifold because of the identifications
$$\xi ^\pm e^{\pm 2\pi Q\mathrm{}_s}\xi ^\pm .$$
(49)
### 3.1 A more detailed derivation
We would like to extend the derivation of Matrix theory given in (see also ) to our time-dependent example. We start with the ten-dimensional string metric
$$ds^2=2dX^+dX^{}+\underset{i=1}{\overset{8}{}}(dX^i)^2$$
(50)
and the light-like linear dilaton
$$\varphi =QX^+.$$
(51)
In discrete light-cone quantization (DLCQ), we make the identification
$$X^{}X^{}+R$$
(52)
and focus on a sector with $`N`$ units of light-cone momentum,
$$p^+=\frac{2\pi N}{R}.$$
(53)
In , the theory with the identification $`(\text{52})`$ was defined as a limit of a space-like compactification, where the shift $`(\text{52})`$ of $`X^{}`$ is accompanied by a small shift of $`X^+`$. However, shifting $`X^+`$ is not a symmetry of our background $`(\text{51})`$, so we have to define the DLCQ in a different way. To that effect, we single out one direction $`X^1`$ from among the $`X^i`$ and make the identification
$$(X^+,X^{},X^1)(X^+,X^{},X^1)+(0,R,ϵR),$$
(54)
where in the end we will take $`ϵ0`$. The Lorentz transformation
$`X^+`$ $`=`$ $`ϵx^+,`$
$`X^{}`$ $`=`$ $`{\displaystyle \frac{x^+}{2ϵ}}+{\displaystyle \frac{x^{}}{ϵ}}+{\displaystyle \frac{x^1}{ϵ}},`$
$`X^1`$ $`=`$ $`x^++x^1`$ (55)
puts the background in the form
$`ds^2`$ $`=`$ $`2dx^+dx^{}+{\displaystyle \underset{i=1}{\overset{8}{}}}(dx^i)^2,`$ (56)
$`\varphi `$ $`=`$ $`Qϵx^+,`$ (57)
with the identification
$$x^1x^1+ϵR.$$
(58)
In this background, we focus on a sector with $`N`$ units of momentum in the $`x^1`$ direction. After a T and an S duality, and introducing
$$r\frac{ϵR}{2\pi \mathrm{}_s},$$
(59)
we are studying a sector with $`N`$ D1-branes wrapped around $`x^1`$ in the type IIB background
$`ds^2`$ $`=`$ $`re^{ϵQx^+}\left\{2dx^+dx^{}+{\displaystyle \underset{i=1}{\overset{8}{}}}(dx^i)^2\right\},`$ (60)
$`\varphi `$ $`=`$ $`ϵQx^++\mathrm{log}r,`$ (61)
with the identification
$$x^1x^1+\frac{2\pi \mathrm{}_s}{r}.$$
(62)
This is now a theory of D1-branes in a background where the string coupling becomes weak near the big-bang, and strong at late times. It is worth stressing that this is opposite to the behavior of the string coupling in our original background (50, 51).
We now need to find a ground state of the D1-brane theory, and study fluctuations about this ground state. First consider a single D1-brane. If in the Dirac-Born-Infeld action,
$$S_{D1}=\frac{1}{2\pi \mathrm{}_s^2}𝑑\tau 𝑑\sigma e^\varphi \sqrt{det\left(_\alpha X^\mu _\beta X^\nu G_{\mu \nu }+2\pi \mathrm{}_s^2F_{\alpha \beta }\right)},$$
(63)
we first put $`F_{\alpha \beta }`$ to zero, then the $`x^+`$ dependence cancels between the inverse string coupling and the determinant of the metric. So a simple classical solution is
$`x^1`$ $`=`$ $`{\displaystyle \frac{1}{r}}\sigma ,`$ (64)
$`x^+`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{\tau }{\sqrt{2}}},`$ (65)
$`x^{}`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{\tau }{\sqrt{2}}},`$ (66)
$`x^i`$ $`=`$ $`0,i=2,\mathrm{},8.`$ (67)
If we now choose the gauge
$`x^1`$ $`=`$ $`{\displaystyle \frac{1}{r}}\sigma ,`$
$`x^+`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{\tau }{\sqrt{2}}},`$ (68)
and define a new coordinate $`y`$ by
$$x^{}=\frac{1}{r}\frac{\tau }{\sqrt{2}}+\sqrt{2}y,$$
(69)
then, ignoring a total derivative, $`(\text{63})`$ can be expanded to give
$`S_{D1}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{}_s^2}}{\displaystyle }d\tau d\sigma ({\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{2}}[(_\tau y)^2+(_\tau x^i)^2(_\sigma y)^2(_\sigma x^i)^2]`$ (70)
$`+2\pi ^2\mathrm{}_s^4\mathrm{exp}({\displaystyle \frac{\sqrt{2}ϵQ\tau }{r}})F_{\tau \sigma }^2+\mathrm{}),`$
with
$$\sigma \sigma +2\pi \mathrm{}_s.$$
(71)
This agrees with the $`N=1`$ case of (45) after rescaling the fields. Note that the coordinate $`y`$ appears on the same footing as the $`x^i`$ ($`i=2,\mathrm{},8`$), so it plays the role that $`x^i`$ used to play before we made the compactification space-like.
For $`N`$ D1-branes, the world-volume theory is given by (45).
### 3.2 Regime of validity of the Matrix string description
The modes (41) of a scalar field in the lightlike linear dilaton background (50, 51) take the following form in terms of the new coordinates (55):
$$T(x^+,x^{},\stackrel{}{x})=e^{ϵQx^+}\mathrm{exp}\left\{i(ϵE^{}+\frac{p^+}{2ϵ}k_1)x^+i\frac{p^+}{ϵ}x^{}+i(k_1\frac{p^+}{ϵ})x^1+i\underset{j=2}{\overset{8}{}}k_jx^j\right\}.$$
(72)
The identification (54), or equivalently (58), implies the momentum quantization condition
$$p^+=ϵk_1\frac{2\pi n}{R}.$$
(73)
In DLCQ, we focus on a sector with given $`N`$ and study fluctuations that stay within that sector. Such fluctuations have $`n=0`$, so that
$$p^+=ϵk_1.$$
(74)
Using (74), the mass shell condition (42) for non-negative $`m^2`$ implies
$$|k_1|2ϵ|E^{}|.$$
(75)
As a consequence, (72) shows that the energy and momentum in the new coordinate system $`(x^+,x^{},x^i)`$ are at most of order
$$ϵE^{}.$$
(76)
Taking into account the identifications (68), we conclude that the world-sheet energy and momentum appearing in the action (70) are at most of order
$$E_{typical}\frac{ϵE^{}}{r}\frac{E^{}\mathrm{}_s}{R}.$$
(77)
The effective time-dependent string length $`\mathrm{}_s^{\mathrm{eff}}`$ can be read from the metric $`(\text{60})`$,
$$\mathrm{}_s^{\mathrm{eff}}=\frac{\mathrm{}_se^{ϵQx^+/2}}{\sqrt{r}}.$$
(78)
The condition for open string oscillators to decouple is given by
$$ϵE^{}\mathrm{}_s^{\mathrm{eff}}=\sqrt{\frac{2\pi ϵ\mathrm{}_s^3}{R}}E^{}e^{ϵQx^+/2}1,$$
(79)
which is satisfied in the $`ϵ0`$ limit.
We must also check that gravity decouples from the Matrix description. The effective ten-dimensional Newton “constant” can be determined from the metric $`(\text{60})`$ and dilaton $`(\text{61})`$,
$$G_N^{\mathrm{eff}}g_s^2(\mathrm{}_s^{\mathrm{eff}})^8=\frac{\mathrm{}_s^8e^{2ϵQx^+}}{r^2}=\frac{4\pi ^2\mathrm{}_s^{10}e^{2ϵQx^+}}{ϵ^2R^2}.$$
(80)
Taking into account the fact that the energies of fluctuations are given by (76), we see that the fluctuations interact gravitationally with strength
$$G_N(ϵE^{})^8\frac{4\pi ^2\mathrm{}_s^{10}ϵ^6(E^{})^8e^{2ϵQx^+}}{R^2},$$
(81)
so closed strings also decouple for $`ϵ0`$.
So complete decoupling of closed and massive open strings can be achieved by strictly setting
$$ϵ=0.$$
(82)
This strict limit makes perfect sense from the SYM point of view, since its coupling
$$g_{YM}=\frac{1}{g_s\mathrm{}_s}=\frac{1}{\mathrm{}_s}\mathrm{exp}\left(\frac{\sqrt{2}\pi \mathrm{}_sQ\tau }{R}\right)$$
(83)
and the typical energies (77) are $`ϵ`$-independent. In other words, by using the limit (82), we reach the remarkable conclusion that our Matrix description is valid all the way to the singularity, which opens up the perspective of using it to study the fate of the singularity.
Now that we have argued that the Matrix description is a complete description of the physics of the singularity, we should ask whether the description is weakly coupled. The relevant dimensionless parameter is the ratio of the SYM coupling to the typical energy of processes we are interested in. Using (83), (77), (84) and
$$p^+=2\pi N/R,$$
(84)
we see that this parameter equals
$$\frac{g_{YM}}{E_{typical}}\frac{N\mathrm{exp}\left(\frac{\sqrt{2}\pi \mathrm{}_sQ\tau }{R}\right)}{p^+E^{}\mathrm{}_s^2}.$$
(85)
Thus for any fixed finite $`N`$ the Matrix description is strongly coupled for late times and weakly coupled for early times. Note, however, that in DLCQ, we eventually want to take the decompactification limit $`N\mathrm{}`$, $`R\mathrm{}`$ with fixed $`p^+=2\pi N/R`$. This is corresponds to the limit where we take $`N\mathrm{}`$ holding $`g_{YM}^2`$ fixed while considering energies of order $`1/N`$. In a strict $`N\mathrm{}`$ limit, which undoes the DLCQ by decompactifying the lightlike circle, one finds a strong coupling behavior for all times. One could also imagine having $`N`$ depend on time and keeping an appropriate combination of the parameter (85) and $`N`$ almost fixed close to the singularity, thus effectively decompactifying the DLCQ circle near the singularity while keeping the SYM theory weakly coupled near the singularity. For some range of times, it might also be useful to perform an analysis along the lines of .
### 3.3 Cosmological evolution and the emergence of space-time
Time evolution from the big bang to late times corresponds beautifully to renormalization group flow in the Yang-Mills theory. At the big bang, the Yang-Mills theory is weakly coupled. Since the coupling has positive mass dimension, weak coupling corresponds to the UV sector of the theory. As time evolves, the coupling increases and the Yang-Mills theory flows to the IR.
At early times, we have weakly coupled Yang-Mills theory. Since the coupling $`(\text{46})`$ is small, the potential terms in $`(\text{45})`$ turn off as we approach the big bang. We are left with a theory of non-commuting matrices. These appear to be the correct degrees of freedom near the singularity, replacing our conventional notion of space-time. At very late times, we recover light-cone quantized perturbative string field theory in the light-like linear dilaton background, along the lines described in .
In Matrix theory descriptions of flat space-time, supersymmetry plays a critical role. Roughly speaking, diagonal matrix elements are interpreted as positions in space-time (they are the positions of the D1-branes in their transverse space), while off-diagonal matrix elements correspond to strings connecting the various D1-branes. When two well-separated clusters of D1-branes are considered, for instance corresponding to two well-separated supergravitons, the off-diagonal modes between the two clusters are heavy and can be integrated out, a priori giving rise to an effective potential for the separation modulus of the two clusters. Supersymmetric cancellations ensure that the effective potential vanishes and that the moduli space metric is flat . This is crucial for the space-time interpretation of Matrix theory: if supergravitons interacted with a static rather than velocity-dependent potential, the model would not describe gravity in flat space-time.
Indeed, a Matrix theory description of a non-supersymmetric flat string background with a closed string tachyon was given in . This Matrix theory develops a potential that lifts the flat directions, a pathology that was given the interpretation that the original non-supersymmetric space-time is not a solution of non-perturbative string theory.
Our model, the light-like linear dilaton background of type IIA string theory, preserves 16 supersymmetries, so one might have hoped to find a supersymmetric Matrix theory description. However, it turns out that supersymmetry is spontaneously broken in any sector with non-zero light-cone momentum $`p^+=2\pi N/R`$. Since in our discrete light-cone quantization we focus on a sector with a fixed non-zero value for $`N`$, we are bound to find a Matrix string description in which supersymmetry is broken. This would have been disastrous if the breaking were explicit.
However, what we actually found is maximally supersymmetric two-dimensional Yang-Mills theory in flat space with a boost identification. The boost identification breaks all the supersymmetry but only via boundary conditions on the Milne circle: namely, the action of the boost transformation acts differently on fields of different spin leading to a different quantization condition for bosons and fermions . This is an effect that becomes less relevant as the Milne circle grows in time. At late times, the potential becomes small as supersymmetry is restored. It would be interesting to understand this potential in a quantitative way.
### 3.4 Does time begin?
The deepest question that we would hope to address with this formalism is whether the big bang should be thought of as the beginning of time, or whether space-time exists prior to the singularity. From a space-time point of view, it is not clear that it makes sense to continue the spacetime metric (19) beyond the singularity at $`\lambda =0`$.
The same question has been addressed for the Milne orbifold as (two dimensions of) a space-time background in perturbative string theory . The conclusion was that $`22`$ scattering amplitudes across the singularity diverge at tree-level because of large tree-level gravitational backreaction from the region close to the singularity .
Our Matrix description a priori contains only the future quadrant of the Milne orbifold as the Matrix string world-sheet. One could try to include the other three quadrants of the Milne orbifold and see whether states can be propagated across the singularity. Although there is now no gravitational backreaction, one can show along the lines of that large gauge backreaction gives rise to UV enhanced IR divergences similar to the ones found in . These divergences are in some sense milder than those encountered in the gravitational case, but they might still be problematic in our low dimension field theory. However, it is not entirely clear that these divergences are associated with the singularity.
It is also possible that the correct prescription involves selecting an initial condition at the big bang and considering only the future quadrant of the Milne orbifold. There are two natural states to choose as an initial vacuum: the conformal vacuum annihilated by positive frequency modes with respect to the conformal time $`\eta `$ given in $`(\text{5})`$, and the adiabatic vacuum inherited from the underlying Minkowski space. For a more detailed discussion, see . At the big bang, we have a conformal field theory since the Yang-Mills theory is free. The natural vacuum from the perspective of conformal field theory would be the conformal vacuum which is $`SL(2)`$ invariant. However, the adiabatic vacuum has the advantage of better high-energy behavior at the singularity, and is the vacuum state that is usually used in string theory computations . At late times, where the perturbative string description is good, it is natural to study states with reference to conformal time which is identified with $`X^+`$ in space-time via $`(\text{68})`$.
It is an interesting question to determine the precise observables in this Matrix model. For example, one possibility could be to take a natural initial Yang-Mills state and ask how it evolves into a collection of excited strings in the space-time that emerges at late times.
## 4 Some Generalizations
### 4.1 A dual type IIA background
There are many interesting ways to generalize the light-like linear dilaton solution. For example, given the M-theory metric (9), we can compactify the $`X^9`$ direction and interpret it as the M-theory circle. This gives rise to an alternate type IIA description, with metric
$$ds_{10}^2=e^{QX^+}[(dX^0)^2+(dX^1)^2+\mathrm{}+(dX^8)^2]+e^{QX^+}(dY)^2$$
(86)
and light-like linear dilaton
$$\varphi =\frac{QX^+}{2}.$$
(87)
In a coordinate basis, the metric (86) has non-vanishing curvature components
$`R_{+i+i}`$ $`=`$ $`{\displaystyle \frac{Q^2}{4}}e^{QX^+},`$ (88)
$`R_{+y+y}`$ $`=`$ $`{\displaystyle \frac{3Q^2}{4}}e^{QX^+}.`$ (89)
Thus we find that the Ricci tensor has a non-vanishing component
$$R_{++}=Q^2.$$
(90)
The non-vanishing Ricci tensor is supported by the dilaton, which in the presence of the metric (86) makes a non-vanishing contribution to Einstein’s equations:
$$_+_+\varphi =\mathrm{\Gamma }_{++}^+_+\varphi =\frac{Q^2}{2}.$$
(91)
Like (9), the space-time (86) is singular as $`X^+\mathrm{}`$ because the metric components $`\stackrel{~}{g}_{ii}`$ go to zero in that limit.
If we now compactify the $`X^8`$-direction, $`X^8X^8+2\pi R^8`$, and T-dualize, we find a type IIB solution with metric
$$ds_{10}^2=e^{QX^+}[(dX^0)^2+(dX^1)^2+\mathrm{}+(dX^7)^2]+e^{QX^+}[(dX^8)^2+(dY)^2]$$
(92)
and constant dilaton,
$$\varphi =\mathrm{log}\left(\frac{\mathrm{}_s}{R^8}\right)$$
(93)
with $`X^8X^8+2\pi \mathrm{}_s^2/R^8`$.
### 4.2 The light-like linear dilaton in type IIB
We could also directly consider the light-like linear dilaton in type IIB string theory. The new ingredient in type IIB that we should consider as $`g_s\mathrm{}`$ is S-duality. The Einstein frame metric $`(\text{4})`$ is invariant under S-duality, but we obtain a new string frame metric valid in the strong coupling regime. We can view the resulting S-dual description as either string theory with a time-independent string length $`\mathrm{}_s`$ but with a coupling and metric
$$\stackrel{~}{g_s}=e^{QX^+},ds^2=e^{QX^+}\left\{2dX^+dX^{}+\underset{i}{}(dX^i)^2\right\},$$
(94)
or as string theory with a coupling and metric
$$\stackrel{~}{g_s}=e^{QX^+},ds^2=\left\{2dX^+dX^{}+\underset{i}{}(dX^i)^2\right\},$$
(95)
and a time-dependent string length $`\stackrel{~}{\mathrm{}_s}^2=e^{QX^+}\mathrm{}_s^2.`$ As we approach the singularity at $`X^+\mathrm{}`$, the effective string tension goes to zero. Once again, we see that the resolution of the singularity requires physics beyond the $`\stackrel{~}{\mathrm{}_s}`$ expansion.
It might appear that because the string coupling is small near the singularity, we should have a good perturbative string description. If this were the case, we might hope to resolve the cosmological singularity in string perturbation theory. However, this is not the case: graviton perturbation theory, which is controlled by the (duality invariant) effective Newton constant, still breaks down near the singularity. We therefore expect a new description involving new degrees of freedom in the strong coupling regime. Such a description should follow from type IIB Matrix string theory .
### 4.3 The light-like linear dilaton from little string theory
Finally, it is worth mentioning that the light-like linear dilaton can be obtained as an unusual Penrose limit of the near horizon geometry of type II NS5-branes studied in .<sup>5</sup><sup>5</sup>5We would like to thank Daniel Robbins for collaboration on the material in this subsection. This geometry is described by
$$ds^2=N\mathrm{}_s^2\left[d\stackrel{~}{t}^2+\frac{dr^2}{r^2}+d\mathrm{\Omega }_3^2\right]+dy_5^2,e^{2\mathrm{\Phi }}=\frac{N\mathrm{}_s^2g_s^2}{r^2}.$$
(96)
The conventional time coordinate $`t`$ has been rescaled, $`t=\sqrt{N}\mathrm{}_s\stackrel{~}{t}`$, to uniformize the factors of $`N`$ appearing in the metric. There is also an NS three-form field strength $`H_3`$ which is $`N`$ times the volume form of the three-sphere, but this will vanish when we take the limit.
We are interested in boosting along a radial rather than angular null geodesic in this space. To do this we switch coordinates $`(\stackrel{~}{t},r)(u,v)`$ by $`\stackrel{~}{t}=uv`$, $`r=\sqrt{N}\mathrm{}_se^u`$. This gives
$$ds^2=N\mathrm{}_s^2\left[2dudvdv^2+d\theta ^2+\mathrm{cos}^2\theta d\psi ^2+\mathrm{sin}^2\theta d\varphi ^2\right]+dy_5^2,e^{2\mathrm{\Phi }}=g_s^2e^{2u}.$$
(97)
Finally, to take the limit , we rescale $`vv/N`$, $`\theta \theta /\sqrt{N}`$, $`\psi \psi /\sqrt{N}`$, and take $`N\mathrm{}`$. This sends $`H_3`$ to zero since
$$H_3N\mathrm{sin}\theta d\theta d\psi d\varphi ,$$
(98)
which goes like $`N^{1/2}`$. We are left with the metric,
$$ds^2=\mathrm{}_s^2\left[2dudv+dx_3^2\right]+dy_5^2,\mathrm{\Phi }=\mathrm{\Phi }_0u,$$
(99)
which describes the light-like linear dilaton in flat space-time.
If we were now to study perturbative states in light-cone string theory on this background, we would like to understand to which states these correspond in the original coordinates. After the limit, the light cone energy and momentum are given by
$`2p^{}`$ $`=`$ $`i{\displaystyle \frac{}{u}}=i\left({\displaystyle \frac{}{\stackrel{~}{t}}}+r{\displaystyle \frac{}{r}}\right)=i\left(\mathrm{}_s\sqrt{N}{\displaystyle \frac{}{t}}+r{\displaystyle \frac{}{r}}\right)`$
$`2p^+`$ $`=`$ $`{\displaystyle \frac{i}{N}}{\displaystyle \frac{}{v}}={\displaystyle \frac{i}{N}}{\displaystyle \frac{}{\stackrel{~}{t}}}={\displaystyle \frac{i}{\sqrt{N}}}\mathrm{}_s{\displaystyle \frac{}{t}}.`$ (100)
The states of interest to us are those with finite $`p^{}`$ and $`p^+`$. Mapping these states back to the original variables, we see that these states comprise a certain sector of high-energy states in little string theory with energies of order $`\sqrt{N}`$.
The main problem with this approach for obtaining a holographic description of the null dilaton is that little is still known about little string theory beyond its bulk definition. However, many other geometries give rise to the light-like linear dilaton via similar limits, and perhaps one of those geometries will provide a more tractable holographic dual in the spirit of the AdS/CFT correspondence.
## Acknowledgements
The work of B. C. is supported by Stichting FOM. The work of S. S. is supported in part by NSF CAREER Grant No. PHY-0094328, and by the Alfred P. Sloan Foundation.
B. C. and S. S. would like to thank the Aspen Center for Physics for hospitality during the early stages of this work. B.C. thanks the organizers of the IPAM “Conformal Field Theory 2nd Reunion Conference” at Lake Arrowhead, where part of this work was carried out. S. S. would also like to thank the organizers of the 2005 Amsterdam String Theory Workshop, where this work was completed.
|
warning/0506/math0506418.html
|
ar5iv
|
text
|
# Mixable Shuffles, Quasi-shuffles and Hopf Algebras
## 1. Introduction
This paper studies the relationship between the mixable shuffle product and the quasi-shuffle product, both generalizations of the shuffle product.
Mixable shuffles arise from the study of Rota-Baxter algebras. Let $`𝐤`$ be a commutative ring and let $`\lambda 𝐤`$ be fixed. A Rota-Baxter $`𝐤`$-algebra of weight $`\lambda `$ (previously called a Baxter algebra) is a pair $`(R,P)`$ in which $`R`$ is a $`𝐤`$-algebra and $`P:RR`$ is a $`𝐤`$-linear map, such that
(1)
$$P(x)P(y)=P(xP(y))+P(P(x)y)+\lambda P(xy),x,yR.$$
Rota-Baxter algebra was introduced by the mathematician Glen Baxter in 1960 to study the theory of fluctuations in probability. Rota greatly contributed to the study of the Rota-Baxter algebra by his pioneer work in the late 1960s and early 1970s and by his survey articles in late 1990s . Unaware of these works, in the early 1980s the school around Faddeev, especially Semenov-Tian-Shansky , developed a whole theory for the Lie algebraic version of equation (1), which is nowadays well-know in the realm of the theory of integrable systems under the name of (modified) classical Yang-Baxter equation.<sup>1</sup><sup>1</sup>1The latter Baxter is the Australian physicist Rodney Baxter. In recent years, Rota-Baxter algebras have found applications in quantum field theory , dendriform algebras , number theory , Hopf algebras and combinatorics .
Key to much of these applications is the realization of the free objects in which the product is defined by mixable shuffles as a generalization of the shuffle product. The shuffle product is a natural generalization of the integration by parts formula and its construction can be traced back to Chen’s path integrals in 1950s. It has been defined and studied in many areas of mathematics, such as Lie and Hopf algebras, algebraic $`K`$-theory, algebraic topology and combinatorics. Its applications can also be found in chemistry and biology. It naturally carries the notion of a Rota-Baxter operator of weight zero.
Another paper on a generalization of the shuffle product was published in the same year as the papers on mixable shuffle products. It was on the quasi-shuffle product by Hoffman<sup>2</sup><sup>2</sup>2Hoffman mentioned in that there was also a generalization in the thesis of F. Fares .. Hoffman’s quasi-shuffle product plays a prominent rôle in the recent studies of harmonic functions, quasi-symmetric functions, multiple zeta values (where in special cases it is also called stuffle product or harmonic product) and $`q`$-multiple zeta values .
Despite the extensive works on the two generalizations of shuffle products, it appears that they were carried out without being aware of each other. In particular, the relation of quasi-shuffles with Rota-Baxter algebras seems unnoticed. For example, in the numerous applications of quasi-shuffles in multiple zeta values in the current literature, no connections with Rota-Baxter algebras and mixable shuffles have been mentioned. In fact, concepts and results on Rota-Baxter algebras were rediscovered in the study of multiple zeta values. For instance, the construction of the stuffle product in follows easily from the construction of free Rota-Baxter algebras in , while the generalized shuffle product in is the same as the mixable shuffle product in .
The situation is similar in the theory of dendriform algebras. Even though both quasi-shuffles and Rota-Baxter algebras have been used to give examples of dendriform algebras , no connection of the two have been made. Also, in the work of Kreimer, and Connes and Kreimer on renormalization theory in perturbative quantum field theory, both the shuffle and its generalization in terms of the quasi-shuffle, and Rota-Baxter algebras appeared, under different contexts.
It was noted in that the two constructions should be related. Our first goal of this paper is to make this connection precise. We show that the recursive formula for the quasi-shuffle product has its explicit form in the mixable shuffle product. Both can be derived from the Baxter relation (1) that defines a Rota-Baxter algebra of weight 1. We further show that the quasi-shuffle algebra on a locally finite set is a subalgebra of a mixable shuffle algebra on the corresponding locally finite algebra. With this connection, the concept of quasi-shuffle algebras can be defined for a larger class of algebras.
This connection allows us to use the Hopf algebra structure on quasi-shuffle algebras to obtain Hopf algebra structures on free Rota-Baxter algebras, generalizing a previous work on this topic. In the other direction, considering the critical rôle played by the quasi-shuffle (stuffle) product in the recent studies of multiple zeta values and quasi-symmetric functions, this connection should allow us to use the theory of Rota-Baxter algebras in the studies of these exciting areas .
The paper is organized as follows. In the next section, we recall the concepts of shuffles, quasi-shuffles and mixable shuffles, and describe their relations (Theorem 2.5). In Section 3, we use these connections to obtain Hopf algebra structures on free Rota-Baxter algebras (Theorem 3.3).
## 2. Shuffles, quasi-shuffles, and mixable shuffles
For the convenience of the reader and for the ease of later references, we recall the definition of each product before giving the relation among them.
### 2.1. Shuffle product
The shuffle product can be defined in two ways, one recursively, one explicitly. We will see that Hoffman’s quasi-shuffle product is a generalization of the recursive definition and the mixable shuffle product is a generalization of the explicit definition.
Let $`𝐤`$ be a commutative ring with identity $`\mathrm{𝟏}_𝐤`$. Let $`V`$ be a $`𝐤`$-module. Consider the $`𝐤`$-module
$$T(V)=\underset{n0}{}V^n.$$
Here the tensor products are taken over $`𝐤`$ and we take $`V^0=𝐤`$.
Usually the shuffle product on $`T(V)`$ starts with the shuffles of permutations . For $`m,n_+`$, define the set of $`(m,n)`$-shuffles by
$$S(m,n)=\{\sigma S_{m+n}\begin{array}{cc}& \end{array}|\begin{array}{c}\sigma ^1(1)<\sigma ^1(2)<\mathrm{}<\sigma ^1(m),\hfill \\ \sigma ^1(m+1)<\sigma ^1(m+2)<\mathrm{}<\sigma ^1(m+n)\hfill \end{array}\}.$$
Here $`S_{m+n}`$ is the symmetric group on $`m+n`$ letters.
For $`a=a_1\mathrm{}a_mV^m`$, $`b=b_1\mathrm{}b_nV^n`$ and $`\sigma S(m,n)`$, the element
$$\sigma (ab)=u_{\sigma (1)}u_{\sigma (2)}\mathrm{}u_{\sigma (m+n)}V^{(m+n)},$$
where
$$u_k=\{\begin{array}{cc}a_k,\hfill & 1km,\hfill \\ b_{km},\hfill & m+1km+n,\hfill \end{array}$$
is called a shuffle of $`a`$ and $`b`$. The sum
(2)
$$a\text{X}b:=\underset{\sigma S(m,n)}{}\sigma (ab)$$
is called the shuffle product of $`a`$ and $`b`$. Also, by convention, $`a\text{X}b`$ is the scalar product if either $`m=0`$ or $`n=0`$. The operation X extends to a commutative and associative binary operation on $`T(V)`$, making $`T(V)`$ into a commutative algebra with identity, called the shuffle product algebra generated by $`V`$.
The shuffle product on $`T(V)`$ can also be recursively defined as follows. As above we choose two elements $`a_1\mathrm{}a_mV^m`$ and $`b_1\mathrm{}b_nV^n`$, and define
$`a_0\text{X}(b_1b_2\mathrm{}b_n)`$ $`=`$ $`a_0b_1b_2\mathrm{}b_n,`$
$`(a_1a_2\mathrm{}a_m)\text{X}b_0`$ $`=`$ $`b_0a_1a_2\mathrm{}a_m,a_0,b_0V^0=𝐤,`$
and
$`(a_1\mathrm{}a_m)\text{X}(b_1\mathrm{}b_n)`$
$`=`$ $`a_1\left((a_2\mathrm{}a_m)\text{X}(b_1\mathrm{}b_n)\right)`$
$`+b_1\left((a_1\mathrm{}a_m)\text{X}(b_2\mathrm{}b_n)\right),a_i,b_jV.`$
###### Lemma 2.1.
For every element $`vV`$, the $`𝐤`$-linear map $`P_{(v)}:(T(V),\text{X})(T(V),\text{X})`$, $`P_{(v)}(a):=va`$ is a Rota-Baxter operator of weight zero.
###### Proof.
This is evident from the recursive definition of the shuffle product. ∎
### 2.2. Quasi-shuffle product
We recall the construction of quasi-shuffle algebras . Let $`X`$ be a locally finite set, that is, $`X`$ is the disjoint union of finite sets $`X_n,n1`$. The elements of $`X_n`$ are defined to have degree $`n`$. Elements in $`X`$ are called letters and monomials in the letters are called words. Even though the original paper only considered $`𝐤`$ to be a subfield of $``$, much of the construction goes through for any commutative ring $`𝐤`$. So we will work in this generality whenever possible. Consider the $`𝐤`$-module underlying the noncommutative polynomial algebra $`𝔄=𝐤X`$, that is, the free $`𝐤`$-algebra generated by $`X`$. The identity 1 of $`𝔄`$ is called the empty word. Define $`\overline{X}=X\{0\}`$. Suppose that there is a pairing
(4)
$$[,]:\overline{X}\times \overline{X}\overline{X}$$
with the properties
* $`[a,0]=0`$ for all $`a\overline{X}`$;
* $`[a,b]=[b,a]`$ for all $`a,b\overline{X}`$;
* $`[[a,b],c]=[a,[b,c]]`$ for all $`a,b,c\overline{X}`$;
* either $`[a,b]=0`$ for all $`a,b\overline{X}`$, or $`\mathrm{deg}([a,b])=\mathrm{deg}(a)+\mathrm{deg}(b)`$ for all $`a,b\overline{X}`$.
We define a Hoffman set to be a locally finite set $`X`$ with a pairing (4) that satisfies conditions S0-S3.
###### Definition 2.2.
Let $`𝐤`$ be a commutative ring and let $`X`$ be a Hoffman set. The quasi-shuffle product $``$ on $`𝔄`$ is defined recursively by
* $`1w=w1=w`$ for any word $`w`$;
* $`(aw_1)(bw_2)=a(w_1(bw_2))+b((aw_1)w_2)+[a,b](w_1w_2)`$, for any words $`w_1,w_2`$ and letters $`a,b`$.
When $`[,]`$ is identically zero, $``$ is the usual shuffle product X defined in Eq. (2.1).
###### Theorem 2.3 ((Hoffman)).
1. $`(𝔄,)`$ is a commutative graded $`𝐤`$-algebra.
2. When $`[,]0`$, $`(𝔄,)`$ is the shuffle product algebra $`(T(V),\text{X})`$, where $`V`$ is the vector space generated by $`X`$.
3. Suppose further $`𝐤`$ is subfield of $``$. Together with the coconcatenation comultiplication
$$\mathrm{\Delta }:𝔄𝔄𝔄,w\underset{uv=w}{}uv$$
where $`uv`$ is the concatenation of words, and counit
$$ϵ:𝔄𝐤,w\delta _{w,1},$$
$`(𝔄,)`$ becomes a graded, connected bialgebra, in fact a Hopf algebra.
### 2.3. Mixable shuffle product
We next turn to the construction of mixable shuffle algebras and their properties . The adjective mixable suggests that certain elements in the shuffles can be mixed or merged. We first give an explicit formula of the product before giving a recursive definition which, under proper restrictions, will be seen to be equivalent to Hoffman’s quasi-shuffle product.
Intuitively, to form the shuffle product, one starts with two decks of cards and puts together all possible shuffles of the two decks. Suppose a shuffle of the two decks is taken and suppose a card from the first deck is followed immediately by a card from the second deck, we allow the option to merge the two cards and call the result a mixable shuffle. To get the mixable shuffle product of the two decks of cards, one puts together all possible mixable shuffles.
Given an $`(m,n)`$-shuffle $`\sigma S(m,n)`$, a pair of indices $`(k,k+1)`$, $`1k<m+n`$, is called an admissible pair for $`\sigma `$ if $`\sigma (k)m<\sigma (k+1)`$. Denote $`𝒯^\sigma `$ for the set of admissible pairs for $`\sigma `$. For a subset $`T`$ of $`𝒯^\sigma `$, call the pair $`(\sigma ,T)`$ a mixable $`(m,n)`$-shuffle. Let $`T`$ be the cardinality of $`T`$. By convention, $`(\sigma ,T)=\sigma `$ if $`T=\mathrm{}`$. Denote
(5)
$$\overline{S}(m,n)=\{(\sigma ,T)\sigma S(m,n),T𝒯^\sigma \}$$
for the set of mixable $`(m,n)`$-shuffles.
Let $`A`$ be a commutative $`𝐤`$-algebra not necessarily having an identity. For $`a=a_1\mathrm{}a_mA^m`$, $`b=b_1\mathrm{}b_nA^n`$ and $`(\sigma ,T)\overline{S}(m,n)`$, the element
$$\sigma (ab;T)=u_{\sigma (1)}\widehat{}u_{\sigma (2)}\widehat{}\mathrm{}\widehat{}u_{\sigma (m+n)}A^{(m+nT)},$$
where for each pair $`(k,k+1)`$, $`1k<m+n`$,
$$u_{\sigma (k)}\widehat{}u_{\sigma (k+1)}=\{\begin{array}{cc}u_{\sigma (k)}u_{\sigma (k+1)},\hfill & (k,k+1)T\hfill \\ u_{\sigma (k)}u_{\sigma (k+1)},\hfill & (k,k+1)T,\hfill \end{array}$$
is called a mixable shuffle of the words $`a`$ and $`b`$.
Now fix $`\lambda 𝐤`$. Define, for $`a`$ and $`b`$ as above, the mixable shuffle product
(6)
$$a^+b:=a_\lambda ^+b=\underset{(\sigma ,T)\overline{S}(m,n)}{}\lambda ^T\sigma (ab;T)\underset{km+n}{}A^k.$$
As in the case of the shuffle product, the operation $`^+`$ extends to a commutative and associative binary operation on
$$\underset{k1}{}A^k=AA^2\mathrm{}$$
Making it a commutative algebra without identity. Note that this is so even when $`A`$ has an identity $`\mathrm{𝟏}_A`$. In any case, we extend the product to
(7)
$$\text{X}^+(A):=\text{X}_{𝐤,\lambda }^+(A):=\underset{k}{}A^k=𝐤AA^2\mathrm{},$$
making $`\text{X}^+(A)`$ a commutative algebra with identity $`\mathrm{𝟏}𝐤`$ .
Suppose $`A`$ has an identity $`\mathrm{𝟏}_A`$. Define
(8)
$$\text{X}(A):=\text{X}_{𝐤,\lambda }(A):=A\text{X}_{𝐤,\lambda }^+(A)$$
to be the tensor product algebra, i.e., the augmented mixable shuffle product $``$ on $`\text{X}(A)`$ is defined by:
(9)
$$(a_0a)(b_0b):=(a_0b_0)(a^+b),a_0,b_0A,a,b\text{X}^+(A).$$
Thus we have the algebra isomorphism (embedding of the second tensor factor)
(10)
$$\alpha :(\text{X}^+(A),^+)(\mathrm{𝟏}_A\text{X}^+(A),).$$
Define the $`𝐤`$-linear endomorphism $`P_A`$ on $`\text{X}_𝐤(A)`$ by assigning
$`P_A(a_0a)`$ $`=`$ $`\mathrm{𝟏}_Aa_0a,aA^n,n1,`$
$`P_A(a_0c)`$ $`=`$ $`\mathrm{𝟏}_Aca_0,cA^0=𝐤`$
and extending by additivity. Let $`j_A:A\text{X}_𝐤(A)`$ be the canonical inclusion map. Call $`(\text{X}_𝐤(A),P_A)`$ the (mixable) shuffle Rota-Baxter $`𝐤`$-algebra on $`A`$ of weight $`\lambda `$. The following theorem was proved in .
###### Theorem 2.4.
The shuffle Rota-Baxter algebra $`(\text{X}_𝐤(A),P_A)`$, together with the natural embedding $`j_A`$, is a free Rota-Baxter $`𝐤`$-algebra on $`A`$ of weight $`\lambda `$. More precisely, for any Rota-Baxter $`𝐤`$-algebra $`(R,P)`$ of weight $`\lambda `$ and algebra homomorphism $`f:AR`$, there is a Rota-Baxter $`𝐤`$-algebra homomorphism $`\stackrel{~}{f}:(\text{X}_𝐤(A),P_A)(R,P)`$ such that $`f=\stackrel{~}{f}j_A`$.
We will suppress $`𝐤`$ and $`\lambda `$ from $`\text{X}_{𝐤,\lambda }(A)`$ when there is no danger of confusion.
### 2.4. The connection
We now establish the connection between quasi-shuffle product and mixable shuffle product.
Let $`𝐤`$ be a commutative ring with identity. Let $`X=_{n1}X_n`$ be a Hoffman set. Then the pairing $`[,]`$ in (4) extends by $`𝐤`$-linearity to a binary operation on the free $`𝐤`$-module $`A=𝐤\{X\}`$ on $`X`$, making $`A`$ into a commutative $`𝐤`$-algebra without identity, with grading $`A_n=𝐤\{X_n\}`$, the free $`𝐤`$-module generated by $`X_n`$. Let $`\stackrel{~}{A}=𝐤A`$ be the unitary $`𝐤`$-algebra spanned by $`A`$. Then $`\stackrel{~}{A}=𝐤\{\stackrel{~}{X}\}`$ where $`\stackrel{~}{X}=\{\mathrm{𝟏}\}X`$ with $`\mathrm{𝟏}_{\stackrel{~}{A}}:=(\mathrm{𝟏}_𝐤,0)`$ the identity of $`\stackrel{~}{A}`$. Here and in the rest of the paper, we will use $`\mathrm{𝟏}_A`$ (instead of $`\mathrm{𝟏}_{\stackrel{~}{A}}`$) to denote this identity of $`\stackrel{~}{A}`$. We will call $`A`$ (resp. $`\stackrel{~}{A}`$) the algebra (resp. unitary algebra) generated by $`X`$.
With the notations in Eq. (7) and (8), we have embeddings
(11)
$$\begin{array}{cccccc}\beta :& \text{X}^+(A)& & \text{X}^+(\stackrel{~}{A})& & \text{X}(\stackrel{~}{A}),\\ & a& & a& & \mathrm{𝟏}_Aa.\end{array}$$
of $`𝐤`$-algebras. Here the first embedding is induced by the embedding $`A\stackrel{~}{A}`$ and the second embedding is the natural one, $`\text{X}^+(\stackrel{~}{A})\text{X}(\stackrel{~}{A}):=\stackrel{~}{A}\text{X}^+(\stackrel{~}{A}).`$
###### Theorem 2.5.
For a Hoffman set $`X`$, the quasi-shuffle algebra $`𝔄=𝐤X`$ is isomorphic to the algebra $`\text{X}^+(A)`$ and thus to the subalgebra $`\mathrm{𝟏}_A\text{X}^+(A)`$ of $`\text{X}(\stackrel{~}{A})`$ where the weight $`\lambda `$ is 1.
###### Proof.
We define
$$f:XXA=A^1\text{X}^+(A)$$
to be the canonical embedding. We note that both $`𝔄`$, with the concatenation product, and $`\text{X}^+(A)`$, with the tensor concatenation, are the free unitary non-commutative $`𝐤`$-algebra on $`X`$. Thus $`f`$ extends uniquely to a bijective map $`\overline{f}:𝔄\text{X}^+(A)`$ such that for any letters $`a_1,\mathrm{},a_nX`$, we have
$$\overline{f}(a_1\mathrm{}a_n)=a_1\mathrm{}a_nA^n.$$
To distinguish elements, we will use $`a=a_1\mathrm{}a_n`$ for an element in $`𝔄`$ and use $`a^{}:=a_1\mathrm{}a_n`$ for an element in $`\text{X}^+(A)`$ in the rest of the proof.
To prove that $`\overline{f}`$ is also an isomorphism between $`𝔄`$, with the quasi-shuffle product $``$, and $`\text{X}^+(A)`$, with the mixable shuffle product $`^+`$, we just need to show that both products satisfy the same recursive relations. We first note that the recursive relation of $``$ in Definition 2.2 can be rewritten as follows. For any $`m,n1`$ and $`a:=a_1\mathrm{}a_m`$, $`b:=b_1\mathrm{}b_n`$ with $`a_i,b_jX,1im,1jn,`$ then
(12)
$$\begin{array}{c}(1).ab=a_1b_1+b_1a_1+[a_1,b_1],\mathrm{when}m,n=1,\hfill \\ (2).ab=a_1b_1\mathrm{}b_n+b_1\left(a_1(b_2\mathrm{}b_n)\right)\hfill \\ +[a_1,b_1]b_2\mathrm{}b_n,\mathrm{when}m=1,n2,\hfill \\ (3).ab=a_1\left((a_2\mathrm{}a_m)b_1\right)+b_1a_1\mathrm{}a_m\hfill \\ +[a_1,b_1]a_2\mathrm{}a_m,\mathrm{when}m2,n=1,\hfill \\ (4).ab=a_1\left((a_2\mathrm{}a_m)(b_1\mathrm{}b_n)\right)+b_1\left((a_1\mathrm{}a_m)(b_2\mathrm{}b_n)\right)\hfill \\ +[a_1,b_1]\left((a_2\mathrm{}a_m)(b_2\mathrm{}b_n)\right),\mathrm{when}m,n2.\hfill \end{array}$$
On the other hand, consider the set $`\overline{S}(m,n)`$ of mixable $`(m,n)`$-shuffles defined in Eq. (5). By \[23, Eq. (4)\], we have
$$\overline{S}(m,n)=\overline{S}_{1,0}(m,n)\stackrel{}{}\overline{S}_{0,1}(m,n)\stackrel{}{}\overline{S}_{1,1}(m,n)$$
where
$`\overline{S}_{1,0}(m,n)`$ $`=`$ $`\{(\sigma ,T)\overline{S}(m,n)(1,2)T,\sigma ^1(1)=1\},`$
$`\overline{S}_{0,1}(m,n)`$ $`=`$ $`\{(\sigma ,T)\overline{S}(m,n)(1,2)T,\sigma ^1(m+1)=1\},`$
$`\overline{S}_{1,1}(m,n)`$ $`=`$ $`\{(\sigma ,T)\overline{S}(m,n)(1,2)T\}.`$
Further, by \[23, p. 137\], for any $`\lambda 𝐤`$,
$`{\displaystyle \underset{(\sigma ,T)\overline{S}_{1,0}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{})`$ $`=`$ $`a_1{\displaystyle \underset{(\sigma ,T)\overline{S}(m1,n)}{}}\lambda ^{|T|}\sigma ((a_2\mathrm{}a_m)b^{};T)`$
$`=`$ $`\{\begin{array}{cc}a_1\left((a_2\mathrm{}a_m)^+b^{}\right),\hfill & m2,\hfill \\ a_1b^{},\hfill & m=1.\hfill \end{array}`$
Similarly,
$`{\displaystyle \underset{(\sigma ,T)\overline{S}_{0,1}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{})`$ $`=`$ $`b_1{\displaystyle \underset{(\sigma ,T)\overline{S}(m,n1)}{}}\lambda ^{|T|}\sigma (a^{}(b_2\mathrm{}b_n);T)`$
$`=`$ $`\{\begin{array}{cc}b_1\left(a^{}^+(b_2\mathrm{}b_n)\right),\hfill & n2,\hfill \\ b_1a^{},\hfill & n=1;\hfill \end{array}`$
and
$`{\displaystyle \underset{(\sigma ,T)\overline{S}_{1,1}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{};T)`$
$`=`$ $`{\displaystyle \underset{(\sigma ,T)\overline{S}(m1,n1)}{}}\lambda [a_1,b_1]\lambda ^{|T|}\sigma ((a_2\mathrm{}a_m)(b_2\mathrm{}b_n);T)`$
$`=`$ $`\{\begin{array}{cc}\lambda [a_1,b_1]\left((a_2\mathrm{}a_m)^+(b_2\mathrm{}b_n)\right),\hfill & m2,n2,\hfill \\ \lambda [a_1,b_1](b_2\mathrm{}b_n),\hfill & m=1,n2,\hfill \\ \lambda [a_1,b_1](a_2\mathrm{}a_m),\hfill & m2,n=1,\hfill \\ \lambda [a_1,b_1],\hfill & m=n=1.\hfill \end{array}`$
Since
$`a^{}^+b^{}:={\displaystyle \underset{(\sigma ,T)\overline{S}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{})`$
$`=`$ $`{\displaystyle \underset{(\sigma ,T)\overline{S}_{1,0}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{})+{\displaystyle \underset{(\sigma ,T)\overline{S}_{0,1}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{})`$
$`+{\displaystyle \underset{(\sigma ,T)\overline{S}_{1,1}(m,n)}{}}\lambda ^{|T|}\sigma (a^{}b^{}),`$
we get
(16)
$$\begin{array}{c}(1).a^{}^+b^{}=a_1b_1+b_1a_1+\lambda [a_1,b_1],\mathrm{when}m,n=1,\hfill \\ (2).a^{}^+b^{}=a_1b_1\mathrm{}b_n+b_1\left(a_1^+(b_2\mathrm{}b_n)\right)\hfill \\ +\lambda [a_1,b_1]b_2\mathrm{}b_n,\mathrm{when}m=1,n2,\hfill \\ (3).a^{}^+b^{}=a_1\left((a_2\mathrm{}a_m)^+b_1\right)+a_1\mathrm{}a_mb_1\hfill \\ +\lambda [a_1,b_1]a_2\mathrm{}a_m,\mathrm{when}m2,n=1,\hfill \\ (4).a^{}^+b^{}=a_1\left((a_2\mathrm{}a_m)^+(b_1\mathrm{}b_n)\right)\hfill \\ +b_1\left((a_1\mathrm{}a_m)^+(b_2\mathrm{}b_n)\right)\hfill \\ +\lambda [a_1,b_1]\left((a_2\mathrm{}a_m)^+(b_2\mathrm{}b_n)\right),\mathrm{when}m,n2.\hfill \end{array}$$
Thus when $`\lambda =1`$, we have $`\overline{f}(ab)=a^{}^+b^{}`$ for all words $`a`$ and $`b`$ with $`m,n1`$, and hence for all $`a`$ and $`b`$ with $`m,n0`$ since when $`m=0`$ or $`n=0`$, we have $`a=1`$ or $`b=1`$ and the multiplications through $``$ and $`^+`$ are both given by the identity. This proves the first isomorphism. The second one then follows from Eq. (11). ∎
###### Corollary 2.6.
Under the same assumptions of Theorem 2.5 and the additional assumption that $`𝐤`$ is a subfield of $``$, for any $`\lambda 𝐤`$, the subalgebra $`\text{X}^+(A)`$ of $`\text{X}^+(\stackrel{~}{A})`$ and the subalgebra $`\mathrm{𝟏}_A\text{X}^+(A)`$ of $`\text{X}(\stackrel{~}{A})`$ are Hopf algebras.
In the next section, we will address the question on whether this Hopf algebra can be extended to a larger Hopf algebra in $`\text{X}(\stackrel{~}{A})`$.
###### Proof.
Because of the isomorphism (10), we only need to prove the first part of the statement. When $`\lambda =1`$, this follows from Theorem 2.5 and Theorem 2.3. When $`\lambda =0`$, this is well-known (see , for example).
Now assume $`\lambda 1,0`$. We define a map
$$g:\text{X}_\lambda ^+(A)\text{X}_1^+(A)$$
by $`g(a_1\mathrm{}a_n)=\lambda ^na_1\mathrm{}a_n,n1`$ and $`g(\mathrm{𝟏}_𝐤)=\mathrm{𝟏}_𝐤`$. Then by the first equation in (16), we have
$$g(a_1_\lambda ^+b_1)=g(a_1b_1+b_1a_1+\lambda [a_1,b_1])=\lambda ^2(a_1b_1+b_1a_1+[a_1,b_1])=\lambda ^2(a_1_1^+b_1)$$
which is just $`g(a_1)_1^+g(b_1).`$ Using the other three equations in (16) and the induction on $`m+n`$, we verify that
$$g\left((a_1\mathrm{}a_m)_\lambda ^+(b_1\mathrm{}b_n)\right)=g(a_1\mathrm{}a_m)_1^+g(b_1\mathrm{}b_n)$$
for all $`m,n1`$. Thus we have $`\text{X}_\lambda ^+(A)`$ is isomorphic to $`\text{X}_1^+(A)`$ and thus carries a Hopf algebra structure. ∎
## 3. Hopf algebras in Rota-Baxter algebras
We first recall the following theorem from .
###### Theorem 3.1 (Andrews-Guo-Keigher-Ono).
For any commutative ring $`𝐤`$ with identity and for any $`\lambda 𝐤`$, the free Rota-Baxter algebra $`\text{X}_\lambda (𝐤)`$ is a Hopf $`𝐤`$-algebra.
As shown in , when $`\lambda =0`$, we have the divided power Hopf algebra.
We now extend this result to $`\text{X}(\stackrel{~}{A})`$ for a $`𝐤`$-algebra $`\stackrel{~}{A}`$ coming from a Hoffman set $`X`$. To avoid confusion, we will use $`\mathrm{𝟏}_𝐤`$ for the identity of $`𝐤`$ and $`\mathrm{𝟏}_A`$ for the identity of $`\stackrel{~}{A}`$ even though they are often identified under the structure map $`𝐤\stackrel{~}{A}`$ of the unitary $`𝐤`$-algebra $`\stackrel{~}{A}`$.
Fix a $`\lambda 𝐤`$. First note that, as a $`𝐤`$-module,
$$\text{X}^+(𝐤)=\underset{n0}{}𝐤^n=𝐤𝐤𝐤^2+\mathrm{}.$$
There are two copies of $`𝐤`$ in the sum since $`𝐤^0=𝐤=𝐤^1`$. The identity of $`\text{X}^+(𝐤)`$ is the identity in the first copy, which we denote by $`\mathrm{𝟏}_𝐤^0=\mathrm{𝟏}`$ as we did in (7). Thus we have
$$\text{X}^+(𝐤)=\mathrm{𝐤𝟏}\mathrm{𝐤𝟏}_𝐤\mathrm{𝐤𝟏}_𝐤^2\mathrm{}=\underset{n0}{}\mathrm{𝐤𝟏}_𝐤^n.$$
Then
$$\text{X}(𝐤)=𝐤\text{X}^+(𝐤)=\underset{n0}{}𝐤(\mathrm{𝟏}_𝐤\mathrm{𝟏}_𝐤^n)$$
with the identity $`\mathrm{𝟏}_𝐤\mathrm{𝟏}`$. Since $`𝐤`$ is the base ring, the algebra homomorphism (11) gives
$$\beta :(\text{X}^+(𝐤),^+)(\text{X}(𝐤),).$$
Thus, by Theorem 3.1 we get
###### Lemma 3.2.
For any $`\lambda 𝐤`$, $`(\text{X}^+(𝐤),^+)`$ is a Hopf algebra.
For now let $`\stackrel{~}{A}`$ be any unitary $`𝐤`$-algebra with unit $`\mathrm{𝟏}_A`$. Then
$$\text{X}^+(\stackrel{~}{A})=\underset{n0}{}\stackrel{~}{A}^n=\mathrm{𝐤𝟏}\stackrel{~}{A}\stackrel{~}{A}^2\mathrm{}$$
and
$$\text{X}(\stackrel{~}{A})=\stackrel{~}{A}\text{X}^+(\stackrel{~}{A})=(\stackrel{~}{A}\mathrm{𝐤𝟏})\stackrel{~}{A}^2\stackrel{~}{A}^3\mathrm{}.$$
Since $`\text{X}(\stackrel{~}{A})`$ is an $`\stackrel{~}{A}`$-algebra, and hence a $`𝐤`$-algebra, we have the structure map $`\gamma :𝐤\text{X}(\stackrel{~}{A})`$ given by $`\gamma (c)=c\mathrm{𝟏}_A\mathrm{𝟏}`$. By the universal property of the free $`𝐤`$-Rota-Baxter algebra $`\text{X}(𝐤)`$, we have an induced homomorphism $`\gamma :\text{X}(𝐤)\text{X}(A)`$ of Rota-Baxter algebras. It is given by
$$\gamma (\mathrm{𝟏}_𝐤\mathrm{𝟏}_𝐤^n)=\mathrm{𝟏}_A\mathrm{𝟏}_A^n,n0.$$
Let
$$\gamma ^+:\text{X}^+(𝐤)\text{X}^+(\stackrel{~}{A}),\mathbf{1}_𝐤^n\mathrm{𝟏}_A^n,n0.$$
We have the following commutative diagram
$$\begin{array}{ccc}\text{X}^+(𝐤)& \stackrel{\gamma ^+}{}& \text{X}^+(\stackrel{~}{A})\\ & & \\ \text{X}(𝐤)=𝐤\text{X}^+(𝐤)& \stackrel{\gamma }{}& \text{X}(\stackrel{~}{A})=\stackrel{~}{A}\text{X}^+(\stackrel{~}{A})\end{array}$$
where the vertical arrow are the injective maps to the second tensor factors.
###### Theorem 3.3.
Let $`𝐤`$ be a field. Let $`X`$ be a Hoffman set and let the algebras $`A`$ and $`\stackrel{~}{A}`$ be the algebra and unitary algebra generated by $`X`$ (as defined before Theorem 2.5). Let $`\lambda 𝐤`$.
1. The algebra product of $`\gamma ^+(\text{X}^+(𝐤))`$ and $`\text{X}^+(A)`$ in $`\text{X}^+(\stackrel{~}{A})`$ has a Hopf algebra structure that expands the Hopf algebra structures on $`\gamma ^+(\text{X}^+(𝐤))`$ (see Lemma 3.2) and $`\text{X}^+(A)`$ (see Corollary 2.6).
2. The algebra product of $`\gamma (\text{X}(𝐤))`$ and $`\mathrm{𝟏}_A\text{X}^+(A)`$ in $`\text{X}(\stackrel{~}{A})`$ has a Hopf algebra structure that expands the Hopf algebra structures on $`\gamma (\text{X}(𝐤))`$ (see Theorem 3.1) and $`\mathrm{𝟏}_A\text{X}^+(A)`$ (see Corollary 2.6).
See Theorem 3.6 for a characterization of the elements in these Hopf algebras.
###### Proof.
Since the restriction of the isomorphism $`\alpha :\text{X}^+(\stackrel{~}{A})\mathrm{𝟏}_A\text{X}^+(\stackrel{~}{A})`$ in (10) restricts to isomorphisms $`\gamma ^+(\text{X}^+(𝐤))\gamma (\text{X}(𝐤))`$ and $`\text{X}^+(A)\mathrm{𝟏}_A\text{X}^+(A)`$, we only need to prove the first statement. Since the tensor product of two commutative, cocommutative Hopf algebras is a Hopf algebra , assuming Proposition 3.4 which is stated and proved below, we see that $`\gamma ^+(\text{X}^+(𝐤))^+(\mathrm{𝟏}\text{X}^+(A))`$ is a Hopf algebra for any $`\lambda 𝐤`$ by Theorem 3.1 and Corollary 2.6. ∎
###### Proposition 3.4.
For any weight $`\lambda 𝐤`$, let $`^+`$ be the mixable shuffle product of weight $`\lambda `$. The two subalgebras $`\gamma ^+(\text{X}^+(𝐤))`$ and $`\text{X}^+(A)`$ of $`\text{X}^+(\stackrel{~}{A})`$ are linearly disjoint. Therefore, $`\gamma ^+(\text{X}^+(𝐤))^+\text{X}^+(A)`$ is isomorphic to the tensor product $`\gamma (\text{X}(𝐤))(\mathrm{𝟏}_A\text{X}^+(A))`$.
###### Proof.
Let $`𝐤,X,\stackrel{~}{X},A`$ and $`\stackrel{~}{A}`$ be as in Theorem 3.3. Since $`X`$ is locally finite, it is countable. So we can write $`X=\{y_n|n1\}`$. Also denote $`y_0=\mathrm{𝟏}_A`$, the unit of $`\stackrel{~}{A}`$. Thus $`\stackrel{~}{X}=\{y_n|n\}`$ and $`\stackrel{~}{A}=_{n0}𝐤y_n.`$ For $`r1`$ and $`I=(i_1,\mathrm{},i_r)^r`$, denote $`y_I=y_{i_1}\mathrm{}y_{i_r}`$. Then
$$A^r=\underset{I_{>0}^r}{}𝐤y_I,\stackrel{~}{A}^r=\underset{I^r}{}𝐤y_I.$$
By convention, we define $`^0=_{>0}^0=\{\mathrm{}\}`$, and $`y_{\mathrm{}}=\mathrm{𝟏}`$. Let $`=_{r0}_{>0}^r`$ and $`\stackrel{~}{}=_{r0}^r`$. We then have
$$\text{X}^+(A)=_I𝐤y_I,\text{X}^+(\stackrel{~}{A})=_{I\stackrel{~}{}}𝐤y_I.$$
Recall that
$$\gamma ^+(\text{X}^+(𝐤))=_{n0}\mathrm{𝐤𝟏}_A^n.$$
So to prove that $`_{n0}𝐤\mathbf{\hspace{0.17em}1}_A^n`$ and $`\text{X}^+(A)`$ are linearly disjoint under the product $`^+`$, we only need to prove
###### Claim 3.1.
The set $`\{\mathrm{𝟏}_A^n^+y_I|n0,I\}`$ is linearly independent.
Before proceeding further, we give a formula for the product $`\mathrm{𝟏}_A^n^+y_I`$ which express a mixable shuffle product as a sum of shuffle products in Eq. (2).
###### Lemma 3.5.
For any $`m0`$ and $`I`$, we have
$$\mathrm{𝟏}_A^m^+y_I=\underset{i=0}{\overset{m}{}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y_I.$$
###### Proof.
Define the length of $`I^r`$ to be $`\mathrm{}(I)=r`$. We will prove by induction on $`w=m+\mathrm{}(I)`$. When $`w=0`$, we have $`m=\mathrm{}(I)=0`$. Then $`\mathrm{𝟏}_A^m`$ and $`y_I`$ are both $`\mathrm{𝟏}`$, so the lemma is clear, as is if either $`m=0`$ or $`\mathrm{}(I)=0`$. Suppose it holds for all $`\mathrm{𝟏}_A^m^+y_I`$ with $`m+\mathrm{}<w`$. For given $`\mathrm{𝟏}_A^m`$ and $`y_I=y_{i_1}\mathrm{}y_{i_r}`$ with $`m1,r1`$ and $`m+r=w`$, let $`y^{}=y_{i_2}\mathrm{}y_{i_r}`$ if $`r>1`$ and $`y^{}=\mathrm{𝟏}`$ if $`r=1`$. Applying the recursive relation of $`^+`$ in Eq. (16), the induction hypothesis, the Pascal equality and the recursive relation of X in Eq. (2.1), we have
$`\mathrm{𝟏}_A^m^+y_I=\mathrm{𝟏}_A(\mathrm{𝟏}_A^{(m1)}^+y_I)+y_{i_1}(\mathrm{𝟏}_A^m^+y^{})+\lambda y_{i_1}(\mathrm{𝟏}_A^{(m1)}^+y^{})`$
$`=`$ $`\mathrm{𝟏}_A\left({\displaystyle \underset{i=0}{\overset{m1}{}}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(m1i)}\text{X}y_I\right)+y_{i_1}\left({\displaystyle \underset{i=0}{\overset{m}{}}}\lambda ^i(\begin{array}{c}n1\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y^{}\right)`$
$`+\lambda y_{i_1}\left({\displaystyle \underset{i=0}{\overset{m1}{}}}\lambda ^i(\begin{array}{c}n1\\ i\end{array})\mathrm{𝟏}_A^{(m1i)}\text{X}y^{}\right)`$
$`=`$ $`\mathrm{𝟏}_A\left({\displaystyle \underset{i=0}{\overset{m1}{}}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(m1i)}\text{X}y_I\right)+y_{i_1}\left({\displaystyle \underset{i=0}{\overset{m}{}}}\lambda ^i(\begin{array}{c}n1\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y^{}\right)`$
$`+y_{i_1}\left({\displaystyle \underset{i=1}{\overset{m}{}}}\lambda ^i(\begin{array}{c}n1\\ i1\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y^{}\right)`$
$`=`$ $`\mathrm{𝟏}_A\left({\displaystyle \underset{i=0}{\overset{m1}{}}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(m1i)}\text{X}y_I\right)+y_{i_1}\left({\displaystyle \underset{i=0}{\overset{m}{}}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y^{}\right)`$
$`\stackrel{(\text{2.1})}{=}`$ $`{\displaystyle \underset{i=0}{\overset{m1}{}}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y_I+\lambda ^my(\begin{array}{c}n\\ m\end{array})\mathrm{𝟏}\text{X}y^{}.`$
Since $`y_{i_1}(\mathrm{𝟏}\text{X}y^{})=y_{i_1}y^{}=y=\mathrm{𝟏}_A^0\text{X}y`$, we get exactly what we want. ∎
We continue with the proof of Proposition 3.4. For $`r1`$, let $`[r]=(1,\mathrm{},r)`$. For a sequence $`I=(i_1,\mathrm{},i_r)^r`$, denote $`\mathrm{S}Supp(I)`$ (called sequential support) for the subsequence (with ordering) of $`I`$ of non-zero entries. For an all zero sequence $`I=(0,\mathrm{},0)`$ and the empty sequence $`\mathrm{}`$, we define $`\mathrm{S}Supp(I)=\mathrm{}`$. We then get a map
$$\mathrm{S}Supp:\stackrel{~}{}\stackrel{~}{}.$$
Clearly, $`=\{I\stackrel{~}{}|\mathrm{S}Supp(I)=I\}`$. So
$$\stackrel{~}{}=\stackrel{}{}_I\mathrm{S}Supp^1(I).$$
For each $`I`$, consider the subset $`𝒪_I=\{y_J|J\mathrm{S}Supp^1(I)\}`$. Then we have
$$\{y_I|I\stackrel{~}{}\}=\stackrel{}{}_I𝒪_I.$$
So $`𝒪_I`$ span linearly independent subspaces of $`\text{X}(\stackrel{~}{A})`$.
By Lemma 3.5, $`\mathrm{𝟏}_A^n^+y_I,n0,`$ is in the linear span of $`𝒪_I`$. Thus to prove Claim 3.1 and hence Proposition 3.4, we only need to prove that, for a fix $`I`$, the subset $`\{\mathrm{𝟏}_A^n^+y_I|n0\}`$ is linearly independent.
Suppose the contrary. Then there are integers $`n_1>n_2>\mathrm{}>n_r0`$ and $`c_10`$ in $`𝐤`$ such that $`_{i=1}^rc_i\mathrm{𝟏}_A^{n_i}^+y_I=0`$. Express this sum as a linear combination in terms of the basis $`𝒪_I`$. By Lemma 3.5, the coefficient of $`\mathrm{𝟏}_A^{n_1}y_I`$ is $`c_1`$, so we must have $`c_1=0`$, a contradiction. ∎
It is desirable to characterize the elements in the Hopf algebra $`\gamma ^+(\text{X}^+(𝐤))^+\text{X}^+(A)`$. This is our last goal in this article. Recall that the length of $`y_I`$ with $`I^r,r0,`$ is defined to be $`\mathrm{}(y_r)=\mathrm{}(I)=r`$. For a given $`I^n`$, the sum $`y_J`$ over $`J^n`$ with $`\mathrm{S}Supp(J)=\mathrm{S}Supp(I)`$, is called the one-shuffled element of $`y_I`$, denoted by $`O(y_I)`$. So $`O(y_I)`$ is the subset of $`𝒪_I`$ consisting of elements of length $`\mathrm{}(I)`$. For example, if $`I=(2,0,1)`$, then the corresponding one-shuffle element of $`y_I=y_2\mathrm{𝟏}_Ay_1`$ is $`O(y_I)=y_2\mathrm{𝟏}_Ay_1+\mathrm{𝟏}_Ay_2y_1+y_2y_1\mathrm{𝟏}_A.`$ On the other hand, $`O(y_I)`$ is itself if $`I`$ is either an all zero sequence or an all non-zero sequence. It is so named because the sum can be obtained from shuffling the subsequence of $`y_I`$ of the $`\mathrm{𝟏}_A`$-entries with the subsequence of $`I`$ of the non-$`\mathrm{𝟏}_A`$ entries (from $`\mathrm{S}Supp(I)`$). To put it in another way, define a relation $``$ on $`\stackrel{~}{}`$ by $`I_1I_2`$ if $`\mathrm{}(I_1)=\mathrm{}(I_2)`$ and $`\mathrm{S}Supp(I_1)=\mathrm{S}Supp(I_2)`$. Then it is easy to check that $``$ is an equivalence relation and a one-shuffled element is of the form $`y_J`$ where the sum is taken over all $`J`$ in an equivalence class.
We now give another version of Theorem 3.3.
###### Theorem 3.6.
Under the hypotheses of Theorem 3.3, the subspace of $`\text{X}^+(\stackrel{~}{A})`$ spanned by one-shuffled elements form a Hopf algebra that contains the Hopf algebras $`\gamma ^+(\text{X}^+(𝐤))`$ and $`\text{X}^+(A)`$.
By Theorem 3.3, we only need to prove the following lemma.
###### Lemma 3.7.
The product of $`\gamma ^+(\text{X}^+(𝐤))`$ and $`\text{X}^+(A)`$ in $`\text{X}^+(\stackrel{~}{A})`$ is given by the subspace generated by one-shuffled elements.
###### Proof.
To prove the proposition, let $`U`$ be the product of $`\gamma ^+(\text{X}^+(𝐤))`$ and $`\text{X}^+(A)`$ in $`\text{X}^+(\stackrel{~}{A})`$, and let $`V`$ be the subspace of one-shuffled elements of $`\text{X}^+(\stackrel{~}{A})`$. Then by Lemma 3.5 and the comments before the theorem, we have $`UV`$. To prove $`VU`$, we only need to show that, for each $`k0`$ and $`I(^+)^n,n0`$, the one-shuffled element $`\mathrm{𝟏}_A^k\text{X}x_I`$ is in $`U`$. When $`n=0`$, $`x_I=\mathrm{𝟏}`$. So $`\mathrm{𝟏}_A^k\text{X}x_I=\mathrm{𝟏}_A^k`$ which is in $`\gamma ^+(\text{X}^+(𝐤))`$ and hence in $`U`$. When $`n1`$, we use induction on $`k`$. When $`k=0`$, then $`\mathrm{𝟏}_A^k\text{X}x_I=x_I`$ which is in $`\text{X}^+(A)`$, hence is in $`U`$. Assume that it is true for $`\mathrm{𝟏}_A^k,k<m`$ and consider $`\mathrm{𝟏}_A^m\text{X}x_I`$. By Lemma 3.5, we have
$$\mathrm{𝟏}_A^m^+y_I=\underset{i=0}{\overset{m}{}}\lambda ^i(\begin{array}{c}n\\ i\end{array})\mathrm{𝟏}_A^{(mi)}\text{X}y_I.$$
The left hand side of the equation is in $`U`$ and, by induction, every term on the right hand side except the first one (with $`i=0`$) is also in $`U`$. Thus the first term, which is $`\mathrm{𝟏}_A^m\text{X}y_{}`$, is also in $`U`$. This completes the induction. ∎
Acknowledgements. We thank Zongzhu Lin for helpful discussions. The first author thanks the I.H.É.S. for the its warm hospitality, and the Ev. Studienwerk for financial support.
|
warning/0506/quant-ph0506082.html
|
ar5iv
|
text
|
# Between classical and quantum11footnote 1To appear in Elsevier’s forthcoming Handbook of the Philosophy of Science, Vol. 2: Philosophy of Physics (eds. John Earman & Jeremy Butterfield). The author is indebted to Stephan de Bièvre, Jeremy Butterfield, Dennis Dieks, Jim Hartle, Gijs Tuynman, Steven Zelditch, and Wojciech Zurek for detailed comments on various drafts of this paper. The final version has greatly benefited from the 7 Pines Meeting on ‘The Classical-Quantum Borderland’ (May, 2005); the author wishes to express his gratitude to Lee Gohlike and the Board of the 7 Pines Meetings for the invitation, and to the other speakers (M. Devoret, J. Hartle, E. Heller, G. ‘t Hooft, D. Howard, M. Gutzwiller, M. Janssen, A. Leggett, R. Penrose, P. Stamp, and W. Zurek) for sharing their insights with him.
## 1 Introduction
Most modern physicists and philosophers would agree that a decent interpretation of quantum mechanics should fullfil at least two criteria. Firstly, it has to elucidate the physical meaning of its mathematical formalism and thereby secure the empirical content of the theory. This point (which we address only in a derivative way) was clearly recognized by all the founders of quantum theory.<sup>3</sup><sup>3</sup>3The history of quantum theory has been described in a large number of books. The most detailed presentation is in Mehra & Rechenberg (1982–2001), but this multi-volume series has by no means superseded smaller works such as Jammer (1966), vander Waerden (1967), Hendry (1984), Darrigol (1992), and Beller (1999). Much information may also be found in biographies such as Heisenberg (1969), Pais (1982), Moore (1989), Pais (1991), Cassidy (1992), Heilbron (2000), Enz (2002), etc. See also Pauli (1979). A new project on the history of matrix mechanics led by Jürgen Renn is on its way. Secondly (and this is the subject of this paper), it has to explain at least the appearance of the classical world.<sup>4</sup><sup>4</sup>4That these point are quite distinct is shown by the Copenhagen Interpretation, which exclusively addresses the first at utter neglect of the second. Nonetheless, in most other approaches to quantum mechanics there is substantial overlap between the various mechanisms that are proposed to fullfil the two criteria in question. As shown by our second quotation above, Planck saw the difficulty this poses, and as a first contribution he noted that the high-temperature limit of his formula for black-body radiation converged to the classical expression. Although Bohr believed that quantum mechanics should be interpreted through classical physics, among the founders of the theory he seems to have been unique in his lack of appreciation of the problem of deriving classical physics from quantum theory. Nonetheless, through his correspondence principle (which he proposed in order to address the first problem above rather than the second) Bohr made one of the most profound contributions to the issue. Heisenberg initially recognized the problem, but quite erroneously came to believe he had solved it in his renowned paper on the uncertainty relations.<sup>5</sup><sup>5</sup>5‘One can see that the transition from micro- to macro-mechanics is now very easy to understand: classical mechanics is altogether part of quantum mechanics.’ (Heisenberg to Bohr, 19 March 1927, just before the submission on 23 March of Heisenberg (1927). See Bohr’s Scientific Correspondence in the Archives for the History of Quantum Physics). Einstein famously did not believe in the fundamental nature of quantum theory, whereas Schrödinger was well aware of the problem from the beginning, later highlighted the issue with his legendary cat, and at various stages in his career made important technical contributions towards its resolution. Ehrenfest stated the well-known theorem named after him. Von Neumann saw the difficulty, too, and addressed it by means of his well-known analysis of the measurement procedure in quantum mechanics.
The problem is actually even more acute than the founders of quantum theory foresaw. The experimental realization of Schrödinger’s cat is nearer than most physicists would feel comfortable with (Leggett, 2002; Brezger et al., 2002; Chiorescu et al., 2003; Marshall et al., 2003; Devoret et al., 2004). Moreover, awkward superpositions are by no means confined to physics laboratories: due to its chaotic motion, Saturn’s moon Hyperion (which is about the size of New York) has been estimated to spread out all over its orbit within 20 years if treated as an isolated quantum-mechanical wave packet (Zurek & Paz, 1995). Furthermore, decoherence theorists have made the point that “measurement” is not only a procedure carried out by experimental physicists in their labs, but takes place in Nature all the time without any human intervention. On the conceptual side, parties as diverse as Bohm & Bell and their followers on the one hand and the quantum cosmologists on the other have argued that a “Heisenberg cut” between object and observer cannot possibly lie at the basis of a fundamental theory of physics.<sup>6</sup><sup>6</sup>6Not to speak of the problem, also raised by quantum cosmologists, of deriving classical space-time from some theory of quantum gravity. This is certainly part of the general program of deriving classical physics from quantum theory, but unfortunately it cannot be discussed in this paper. These and other remarkable insights of the past few decades have drawn wide attention to the importance of the problem of interpreting quantum mechanics, and in particular of explaining classical physics from it.
We will discuss these ideas in more detail below, and indeed our discussion of the relationship between classical and quantum mechanics will be partly historical. However, other than that it will be technical and mathematically rigorous. For the problem at hand is so delicate that in this area sloppy mathematics is almost guaranteed to lead to unreliable physics and conceptual confusion (notwithstanding the undeniable success of poor man’s math elsewhere in theoretical physics). Except for von Neumann, this was not the attitude of the pioneers of quantum mechanics; but while it has to be acknowledged that many of their ideas are still central to the current discussion, these ideas per se have not solved the problem. Thus we assume the reader to be familiar with the Hilbert space formalism of quantum mechanics,<sup>7</sup><sup>7</sup>7 Apart from seasoned classics such as Mackey (1963), Jauch (1968), Prugovecki (1971), Reed & Simon (1972), or Thirring (1981), the reader might consult more recent books such as Gustafson & Sigal (2003) or Williams (2003). See also Dickson (2005). and for some parts of this paper (notably Section 6 and parts of Section 4) also with the basic theory of $`C^{}`$-algebras and its applications to quantum theory.<sup>8</sup><sup>8</sup>8For physics-oriented introductions to $`C^{}`$-algebras see Davies (1976), Roberts & Roepstorff (1969), Primas (1983), Thirring (1983), Emch (1984), Strocchi (1985), Sewell (1986), Roberts (1990), Haag (1992), Landsman (1998), Araki (1999), and Sewell (2002). Authoratitive mathematical texts include Kadison & Ringrose (1983, 1986) and Takesaki (2003). In addition, some previous encounter with the conceptual problems of quantum theory would be helpful.<sup>9</sup><sup>9</sup>9Trustworthy books include, for example, Scheibe (1973), Jammer (1974), van Fraassen (1991), d Espagnat (1995), Peres (1995), Omnès (1994, 1999), Bub (1997), and Mittelstaedt (2004).
Which ideas have solved the problem of explaining the appearance of the classical world from quantum theory? In our opinion, none have, although since the founding days of quantum mechanics a number of new ideas have been proposed that almost certainly will play a role in the eventual resolution, should it ever be found. These ideas surely include:
* The limit $`\mathrm{}0`$ of small Planck’s constant (coming of age with the mathematical field of microlocal analysis);
* The limit $`N\mathrm{}`$ of a large system with $`N`$ degrees of freedom (studied in a serious only way after the emergence of $`C^{}`$-algebraic methods);
* Decoherence and consistent histories.
Mathematically, the second limit may be seen as a special case of the first, though the underlying physical situation is of course quite different. In any case, after a detailed analysis our conclusion will be that none of these ideas in isolation is capable of explaining the classical world, but that there is some hope that by combining all three of them, one might do so in the future.
Because of the fact that the subject matter of this review is unfinished business, to date one may adopt a number of internally consistent but mutually incompatible philosophical stances on the relationship between classical and quantum theory. Two extreme ones, which are always useful to keep in mind whether one holds one of them or not, are:
1. Quantum theory is fundamental and universally valid, and the classical world has only “relative” or “perspectival” existence.
2. Quantum theory is an approximate and derived theory, possibly false, and the classical world exists absolutely.
An example of a position that our modern understanding of the measurement problem<sup>10</sup><sup>10</sup>10See the books cited in footnote 9, especially Mittelstaedt (2004). has rendered internally inconsistent is:
> 3. Quantum theory is fundamental and universally valid, and (yet) the classical world exists absolutely.
In some sense stance 1 originates with Heisenberg (1927), but the modern era started with Everett (1957).<sup>11</sup><sup>11</sup>11 Note, though, that stance 1 by no means implies the so-called Many-Worlds Interpretation, which also in our opinion is ‘simply a meaningless collage of words’ (Leggett, 2002). These days, most decoherence theorists, consistent historians, and modal interpreters seem to support it. Stance 2 has a long and respectable pedigree unequivocally, including among others Einstein, Schrödinger, and Bell. More recent backing has come from Leggett as well as from “spontaneous collapse” theorists such as Pearle, Ghirardi, Rimini, Weber, and others. As we shall see in Section 3, Bohr’s position eludes classification according to these terms; our three stances being of an ontological nature, he probably would have found each of them unattractive.<sup>12</sup><sup>12</sup>12To the extent that it was inconclusive, Bohr’s debate with Einstein certainly suffered from the fact that the latter attacked strawman 3 (Landsman, 2006). The fruitlessness of discussions such as those between Bohm and Copenhagen (Cushing, 1994) or between Bell (1987, 2001) and Hepp (1972) has the same origin.
Of course, one has to specify what the terminology involved means. By quantum theory we mean standard quantum mechanics including the eigenvector-eigenvalue link.<sup>13</sup><sup>13</sup>13Let $`A`$ be a selfadjoint operator on a Hilbert space $``$, with associated projection-valued measure $`P(\mathrm{\Delta })`$, $`\mathrm{\Delta }`$, so that $`A=𝑑P(\lambda )\lambda `$ (see also footnote 100 below). The eigenvector-eigenvalue link states that a state $`\mathrm{\Psi }`$ of the system lies in $`P(\mathrm{\Delta })`$ if and only if $`A`$ takes some value in $`\mathrm{\Delta }`$ for sure. In particular, if $`\mathrm{\Psi }`$ is an eigenvector of $`A`$ with eigenvalue $`\lambda `$ (so that $`P(\{\lambda \})0`$ and $`\mathrm{\Psi }P(\{\lambda \})`$), then $`A`$ takes the value $`\lambda `$ in the state $`\mathrm{\Psi }`$ with probability one. In general, the probability $`p_\mathrm{\Psi }(\mathrm{\Delta })`$ that in a state $`\mathrm{\Psi }`$ the observable $`a`$ takes some value in $`\mathrm{\Delta }`$ (“upon measurement”) is given by the Born–von Neumann rule $`p_\mathrm{\Psi }(\mathrm{\Delta })=(\mathrm{\Psi },P(\mathrm{\Delta })\mathrm{\Psi })`$. Modal interpretations of quantum mechanics (Dieks (1989a,b; van Fraassen, 1991; Bub, 1999; Vermaas, 2000; Bene & Dieks, 2002; Dickson, 2005) deny this link, and lead to positions close to or identical to stance 1. The projection postulate is neither endorsed nor denied when we generically speak of quantum theory.
It is a bit harder to say what “the classical world” means. In the present discussion we evidently can not define the classical world as the world that exists independently of observation - as Bohr did, see Subsection 3.1 \- but neither can it be taken to mean the part of the world that is described by the laws of classical physics full stop; for if stance 1 is correct, then these laws are only approximately valid, if at all. Thus we simply put it like this:
> The classical world is what observation shows us to behave - with appropriate accuracy - according to the laws of classical physics.
There should be little room for doubt as to what ‘with appropriate accuracy’ means: the existence of the colour grey does not imply the nonexistence of black and white!
We can define the absolute existence of the classical world à la Bohr as its existence independently of observers or measuring devices. Compare with Moore’s (1939) proof of the existence of the external world:
> How? By holding up my two hands, and saying, as I make a certain gesture with the right hand, ‘Here is one hand’, and adding, as I make a certain gesture with the left, ‘and here is another’.
Those holding position 1, then, maintain that the classical world exists only as an appearance relative to a certain specification, where the specification in question could be an observer (Heisenberg), a certain class of observers and states (as in decoherence theory), or some coarse-graining of the Universe defined by a particular consistent set of histories, etc. If the notion of an observer is construed in a sufficiently abstract and general sense, one might also formulate stance 1 as claiming that the classical world merely exists from the perspective of the observer (or the corresponding class of observables).<sup>14</sup><sup>14</sup>14The terminology “perspectival” was suggested to the author by Richard Healey. For example, Schrödinger’s cat “paradox” dissolves at once when the appropriate perspective is introduced; cf. Subsection 6.6.
Those holding stance 2, on the other hand, believe that the classical world exists in an absolute sense (as Moore did). Thus stance 2 is akin to common-sense realism, though the distinction between 1 and 2 is largely independent of the issue of scientific realism.<sup>15</sup><sup>15</sup>15See Landsman (1995) for a more elaborate discussion of realism in this context. Words like “objective” or “subjective” are not likely to be helpful in drawing the distinction either: the claim that ‘my children are the loveliest creatures in the world’ is at first glance subjective, but it can trivially be turned into an objective one through the reformulation that ‘Klaas Landsman finds his children the loveliest creatures in the world’. Similarly, the proposition that (perhaps due to decoherence) ‘local observers find that the world is classical’ is perfectly objective, although it describes a subjective experience. See also Davidson (2001). For defendants of stance 1 usually still believe in the existence of some observer-independent reality (namely somewhere in the quantum realm), but deny that this reality incorporates the world observed around us. This justifies a pretty vague specification of such an important notion as the classical world: one of the interesting outcomes of the otherwise futile discussions surrounding the Many Worlds Interpretation has been the insight that if quantum mechanics is fundamental, then the notion of a classical world is intrinsically vague and approximate. Hence it would be self-defeating to be too precise at this point.<sup>16</sup><sup>16</sup>16See Wallace (2002, 2003); also cf. Butterfield (2002). This point was not lost on Bohr and Heisenberg either; see Scheibe (1973).
Although stance 1 is considered defensive if not cowardly by adherents of stance 2, it is a highly nontrivial mathematical fact that so far it seems supported by the formalism of quantum mechanics. In his derision of what he called ‘FAPP’ (= For All Practical Purposes) solutions to the measurement problem (and more general attempts to explain the appearance of the classical world from quantum theory), Bell (1987, 2001) and others in his wake mistook a profound epistemological stance for a poor defensive move.<sup>17</sup><sup>17</sup>17The insistence on “precision” in such literature is reminiscent of Planck’s long-held belief in the absolute nature of irreversibility (Darrigol, 1992; Heilbron, 2002). It should be mentioned that although Planck’s stubbornness by historical accident led him to take the first steps towards quantum theory, he eventually gave it up to side with Boltzmann. It is, in fact, stance 2 that we would recommend to the cowardly: for proving or disproving stance 1 seems the real challenge of the entire debate, and we regard the technical content of this paper as a survey of progress towards actually proving it. Indeed, to sum up our conclusions, we claim that there is good evidence that:
1. Classical physics emerges from quantum theory in the limit $`\mathrm{}0`$ or $`N\mathrm{}`$ provided that the system is in certain “classical” states and is monitored with “classical” observables only;
2. Decoherence and consistent histories will probably explain why the system happens to be in such states and has to be observed in such a way.
However, even if one fine day this scheme will be made to work, the explanation of the appearance of the classical world from quantum theory will be predicated on an external solution of the notorious ‘from “and” to “or” problem’: If quantum mechanics predicts various possible outcomes with certain probabilities, why does only one of these appear to us?<sup>18</sup><sup>18</sup>18It has to be acknowledged that we owe the insistence on this question to the defendants of stance 2. See also footnote 11.
For a more detailed outline of this paper we refer to the table of contents above. Most philosophical discussion will be found in Section 3 on the Copenhagen interpretation, since whatever its merits, it undeniably set the stage for the entire discussion on the relationship between classical and quantum.<sup>19</sup><sup>19</sup>19We do not discuss the classical limit of quantum mechanics in the philosophical setting of theory reduction and intertheoretic relations; see, e.g., Scheibe (1999) and Batterman (2002). The remainder of the paper will be of an almost purely technical nature. Beyond this point we will try to avoid controversy, but when unavoidable it will be confined to the Epilogues appended to Sections 3-6. The final Epilogue (Section 8) expresses our deepest thoughts on the subject.
## 2 Early history
This section is a recapitulation of the opinions and contributions of the founders of quantum mechanics regarding the relationship between classical and quantum. More detail may be found in the books cited in footnote 3 and in specific literature to be cited; for an impressive bibliography see also Gutzwiller (1998). The early history of quantum theory is of interest in its own right, concerned as it is with one of the most significant scientific revolutions in history. Although this history is not a main focus of this paper, it is of special significance for our theme. For the usual and mistaken interpretation of Planck’s work (i.e. the idea that he introduced something like a “quantum postulate”, see Subsection 3.2 below) appears to have triggered the belief that quantum theory and Planck’s constant are related to a universal discontinuity in Nature. Indeed, this discontinuity is sometimes even felt to mark the basic difference between classical and quantum physics. This belief is particularly evident in the writings of Bohr, but still resonates even today.
### 2.1 Planck and Einstein
The relationship between classical physics and quantum theory is so subtle and confusing that historians and physicists cannot even agree about the precise way the classical gave way to the quantum! As Darrigol (2001) puts it: ‘During the past twenty years, historians \[and physicists\] have disagreed over the meaning of the quanta which Max Planck introduced in his black-body theory of 1900. The source of this confusion is the publication (…) of Thomas Kuhn’s \[(1978)\] iconoclastic thesis that Planck did not mean his energy quanta to express a quantum discontinuity.’
As is well known (cf. Mehra & Rechenberg, 1982a, etc.), Planck initially derived Wien’s law for blackbody radiation in the context of his (i.e. Planck’s) program of establishing the absolute nature of irreversibility (competing with Boltzmann’s probabilistic approach, which eventually triumphed). When new high-precision measurements in October 1900 turned out to refute Wien’s law, Planck first guessed his expression
$$E_\nu /N_\nu =h\nu /(e^{h\nu /kT}1)$$
(2.1)
for the correct law, en passant introducing two new constants of nature $`h`$ and $`k`$,<sup>20</sup><sup>20</sup>20Hence Boltzmann’s constant $`k`$ was introduced by Planck, who was the first to write down the formula $`S=k\mathrm{log}W`$. and subsequently, on December 14, 1900, presented a theoretical derivation of his law in which he allegedly introduced the idea that the energy of the resonators making up his black body was quantized in units of $`\epsilon _\nu =h\nu `$ (where $`\nu `$ is the frequency of a given resonator). This derivation is generally seen as the birth of quantum theory, with the associated date of birth just mentioned.
However, it is clear by now (Kuhn, 1978; Darrigol, 1992, 2001; Carson, 2000; Brush, 2002) that Planck was at best agnostic about the energy of his resonators, and at worst assigned them a continuous energy spectrum. Technically, in the particular derivation of his empirical law that eventually turned out to lead to the desired result (which relied on Boltzmann’s concept of entropy),<sup>21</sup><sup>21</sup>21Despite the fact that Planck only converted to Boltzmann’s approach to irreversibility around 1914. Planck had to count the number of ways a given amount of energy $`E_\nu `$ could be distributed over a given number of resonators $`N_\nu `$ at frequency $`\nu `$. This number is, of course, infinite, hence in order to find a finite answer Planck followed Boltzmann in breaking up $`E_\nu `$ into a large number $`A_\nu `$ of portions of identical size $`\epsilon _\nu `$, so that $`A_\nu \epsilon _\nu =E_\nu `$.<sup>22</sup><sup>22</sup>22The number in question is then given by $`(N+A1)!/(N1)!A!`$, dropping the dependence on $`\nu `$ in the notation. Now, as we all know, whereas Boltzmann let $`\epsilon _\nu 0`$ at the end of his corresponding calculation for a gas, Planck discovered that his empirical blackbody law emerged if he assumed the relation $`\epsilon _\nu =h\nu `$.
However, this postulate did not imply that Planck quantized the energy of his resonators. In fact, in his definition of a given distribution he counted the number of resonators with energy between say $`(k1)\epsilon _\nu `$ and $`k\epsilon _\nu `$ (for some $`k`$), as Boltzmann did in an analogous way for a gas, rather than the number of resonators with energy $`k\epsilon _\nu `$, as most physicists came to interpret his procedure. More generally, there is overwhelming textual evidence that Planck himself by no means believed or implied that he had quantized energy; for one thing, in his Nobel Prize Lecture in 1920 he attributed the correct interpretation of the energy-quanta $`\epsilon _\nu `$ to Einstein. Indeed, the modern understanding of the earliest phase of quantum theory is that it was Einstein rather than Planck who, during the period 1900–1905, clearly realized that Planck’s radiation law marked a break with classical physics (Büttner, Renn, & Schemmel, 2003). This insight, then, led Einstein to the quantization of energy. This he did in a twofold way, both in connection with Planck’s resonators - interpreted by Einstein as harmonic oscillators in the modern way - and, in a closely related move, through his concept of a photon. Although Planck of course introduced the constant named after him, and as such is the founding father of the theory characterized by $`\mathrm{}`$, it is the introduction of the photon that made Einstein at least the mother of quantum theory. Einstein himself may well have regarded the photon as his most revolutionary discovery, for what he wrote about his pertinent paper is not matched in self-confidence by anything he said about relativity: ‘Sie handelt über die Strahlung und die energetischen Eigenschaften des Lichtes und ist sehr revolutionär.’<sup>23</sup><sup>23</sup>23‘\[This paper\] is about radiation and the energetic properties of light, and is very revolutionary.’ See also the Preface to Pais (1982).
Finally, in the light of the present paper, it deserves to be mentioned that Einstein (1905) and Planck (1906) were the first to comment on the classical limit of quantum theory; see the preamble to Section 5 below.
### 2.2 Bohr
Bohr’s brilliant model of the atom reinforced his idea that quantum theory was a theory of quanta.<sup>24</sup><sup>24</sup>24Although at the time Bohr followed practically all physicists in their rejection of Einstein’s photon, since he believed that during a quantum jump the atom emits electromagnetic radiation in the form of a spherical wave. His model probably would have gained in consistency by adopting the photon picture of radiation, but in fact Bohr was to be the last prominent opponent of the photon, resisting the idea until 1925. See also Blair Bolles (2004) and footnote 35 below. Since this model simultaneously highlighted the clash between classical and quantum physics and carried the germ of a resolution of this conflict through Bohr’s equally brilliant correspondence principle, it is worth saying a few words about it here.<sup>25</sup><sup>25</sup>25Cf. Darrigol (1992) for a detailed treatment; also see Liboff (1984) and Steiner (1998). Bohr’s atomic model addressed the radiative instability of Rutherford’s solar-system-style atom:<sup>26</sup><sup>26</sup>26The solar system provides the popular visualization of Rutherford’s atom, but his own picture was more akin to Saturn’ rings than to a planet orbiting the Sun. according to the electrodynamics of Lorentz, an accelerating electron should radiate, and since the envisaged circular or elliptical motion of an electron around the nucleus is a special case of an accelerated motion, the electron should continuously lose energy and spiral towards the nucleus.<sup>27</sup><sup>27</sup>27In addition, any Rutherford style atom with more than one electron is mechanically unstable, since the electrons repel each other, as opposed to planets, which attract each other. Bohr countered this instability by three simultaneous moves, each of striking originality:
1. He introduced a quantization condition that singled out only a discrete number of allowed electronic orbits (which subsequently were to be described using classical mechanics, for example, in Bohr’s calculation of the Rydberg constant $`R`$).
2. He replaced the emission of continuous radiation called for by Lorentz by quantum jumps with unpredictable destinations taking place at unpredictable moments, during which the atom emits light with energy equal to the energy difference of the orbits between which the electron jumps.
3. He prevented the collapse of the atom through such quantum jumps by introducing the notion of ground state, below which no electron could fall.
With these postulates, for which at the time there existed no foundation whatsoever,<sup>28</sup><sup>28</sup>28What has hitherto been mathematically proved of Bohr’s atomic model is the existence of a ground state (see Griesemer, Lieb, & Loss, 2001, and references therein for the greatest generality available to date) and the metastability of the excited states of the atom after coupling to the electromagnetic field (cf. Bach, Fröhlich, & Sigal, 1998, 1999 and Gustafson & Sigal, 2003). The energy spectrum is discrete only if the radiation field is decoupled, leading to the usual computation of the spectrum of the hydrogen atom first performed by Schrödinger and Weyl. See also the end of Subsection 5.4. Bohr explained the spectrum of the hydrogen atom, including an amazingly accurate calculation of $`R`$. Moreover, he proposed what was destined to be the key guiding principle in the search for quantum mechanics in the coming decade, viz. the correspondence principle (cf. Darrigol, 1992, passim, and Mehra & Rechenberg, 1982a, pp. 249–257).
In general, there is no relation between the energy that an electron loses during a particular quantum jump and the energy it would have radiated classically (i.e. according to Lorentz) in the orbit it revolves around preceding this jump. Indeed, in the ground state it cannot radiate through quantum jumps at all, whereas according to classical electrodynamics it should radiate all the time. However, Bohr saw that in the opposite case of very wide orbits (i.e. those having very large principal quantum numbers $`n`$), the frequency $`\nu =(E_nE_{n1})/h`$ (with $`E_n=R/n^2`$) of the emitted radiation approximately corresponds to the frequency of the lowest harmonic of the classical theory, applied to electron motion in the initial orbit.<sup>29</sup><sup>29</sup>29Similarly, higher harmonics correspond to quantum jumps $`nnk`$ for $`k>1`$. Moreover, the measured intensity of the associated spectral line (which theoretically should be related to the probability of the quantum jump, a quantity out of the reach of early quantum theory), similarly turned out to be given by classical electrodynamics. This property, which in simple cases could be verified either by explicit computation or by experiment, became a guiding principle in situations where it could not be verified, and was sometimes even extended to low quantum numbers, especially when the classical theory predicted selection rules.
It should be emphasized that Bohr’s correspondence principle was concerned with the properties of radiation, rather than with the mechanical orbits themselves.<sup>30</sup><sup>30</sup>30As such, it remains to be verified in a rigorous way. This is not quite the same as what is usually called the correspondence principle in the modern literature.<sup>31</sup><sup>31</sup>31A typical example of the modern version is: ‘Non-relativistic quantum mechanics was founded on the correspondence principle of Bohr: “When the Planck constant $`\mathrm{}`$ can be considered small with respect to the other parameters such as masses and distances, quantum theory approaches classical Newton theory.”’ (Robert, 1998, p. 44). The reference to Bohr is historically inaccurate! In fact, although also this modern correspondence principle has a certain range of validity (as we shall see in detail in Section 5), Bohr never endorsed anything like that, and is even on record as opposing such a principle:<sup>32</sup><sup>32</sup>32Quoted from Miller (1984), p. 313.
> ‘The place was Purcell’s office where Purcell and others had taken Bohr for a few minutes of rest \[during a visit to the Physics Department at Harvard University in 1961\]. They were in the midst of a general discussion when Bohr commented: “People say that classical mechanics is the limit of quantum mechanics when $`h`$ goes to zero.” Then, Purcell recalled, Bohr shook his finger and walked to the blackboard on which he wrote $`e^2/hc`$. As he made three strokes under $`h`$, Bohr turned around and said, “you see $`h`$ is in the denominator.”’
### 2.3 Heisenberg
Heisenberg’s (1925) paper Über die quantentheoretische Umdeutung kinematischer und mechanischer Beziehungen<sup>33</sup><sup>33</sup>33On the quantum theoretical reinterpretation of kinematical and mechanical relations. English translation in vander Waerden, 1967. is generally seen as a turning point in the development of quantum mechanics. Even A. Pais, no friend of Heisenberg’s,<sup>34</sup><sup>34</sup>34For example, in Pais (2000), claiming to portray the ‘genius of science’, Heisenberg is conspicously absent. conceded that Heisenberg’s paper marked ’one of the great jumps - perhaps the greatest - in the development of twentieth century physics.’ What did Heisenberg actually accomplish? This question is particularly interesting from the perspective of our theme.
At the time, atomic physics was in a state of crisis, to which various camps responded in different ways. Bohr’s approach might best be described as damage control: his quantum theory was a hybrid of classical mechanics adjusted by means of ad hoc quantization rules, whilst keeping electrodynamics classical at all cost.<sup>35</sup><sup>35</sup>35 Continuing footnote 24, we quote from Mehra & Rechenberg, 1982a, pp 256–257: ‘Thus, in the early 1920s, Niels Bohr arrived at a definite point of view how to proceed forward in atomic theory. He wanted to make maximum use of what he called the “more dualistic prescription” (…) In it the atom was regarded as a mechanical system having discrete states and emitting radiation of discrete frequencies, determined (in a nonclassical way) by the energy differences between stationary states; radiation, on the other hand, had to be described by the classical electrodynamic theory.’ Einstein, who had been the first physicist to recognize the need to quantize classical electrodynamics, in the light of his triumph with General Relativity nonetheless dreamt of a classical field theory with singular solutions as the ultimate explanation of quantum phenomena. Born led the radical camp, which included Pauli: he saw the need for an entirely new mechanics replacing classical mechanics,<sup>36</sup><sup>36</sup>36It was Born who coined the name quantum mechanics even before Heisenberg’s paper. which was to be based on discrete quantities satisfying difference equations.<sup>37</sup><sup>37</sup>37This idea had earlier occurred to Kramers. This was a leap in the dark, especially because of Pauli’s frowning upon the correspondence principle (Hendry, 1984; Beller, 1999).
It was Heisenberg’s genius to interpolate between Bohr and Born.<sup>38</sup><sup>38</sup>38Also literally! Heisenberg’s traveled between Copenhagen and Göttingen most of the time. The meaning of his Umdeutung was to keep the classical equations of motion,<sup>39</sup><sup>39</sup>39This crucial aspect of Umdeutung was appreciated at once by Dirac (1926): ‘In a recent paper Heisenberg puts forward a new theory which suggests that it is not the equations of classical mechanics that are in any way at fault, but that the mathematical operations by which physical results are deduced from them require modification. (…) The correspondence between the quantum and classical theories lies not so much in the limiting agreement when $`\mathrm{}0`$ as in the fact that the mathematical operations on the two theories obey in many cases the same laws.’ whilst reinterpreting the mathematical symbols occurring therein as (what were later recognized to be) matrices. Thus his Umdeutung $`xa(n,m)`$ was a precursor of what now would be called a quantization map $`fQ_{\mathrm{}}(f)`$, where $`f`$ is a classical observable, i.e. a function on phase space, and $`Q_{\mathrm{}}(f)`$ is a quantum mechanical observable, in the sense of an operator on a Hilbert space or, more abstractly, an element of some $`C^{}`$-algebra. See Section 4 below. As Heisenberg recognized, this move implies the noncommutativity of the quantum mechanical observables; it is this, rather than something like a “quantum postulate” (see Subsection 3.2 below), that is the defining characteristic of quantum mechanics. Indeed, most later work on quantum physics and practically all considerations on the connection between classical and quantum rely on Heisenberg’s idea of Umdeutung. This even applies to the mathematical formalism as a whole; see Subsection 2.5.
We here use the term “observable” in a loose way. It is now well recognized (Mehra & Rechenberg, 1982b; Beller, 1999; Camilleri, 2005) that Heisenberg’s claim that his formalism could be physically interpreted as the replacement of atomic orbits by observable quantities was a red herring, inspired by his discussions with Pauli. In fact, in quantum mechanics any mechanical quantity has to be “reinterpreted”, whether or not it is observable. As Heisenberg (1969) recalls, Einstein reprimanded him for the illusion that physics admits an a priori notion of an observable, and explained that a theory determines what can be observed. Rethinking the issue of observability then led Heisenberg to his second major contribution to quantum mechanics, namely his uncertainty relations.
These relations were Heisenberg’s own answer to the quote opening this paper. Indeed, matrix mechanics was initially an extremely abstract and formal scheme, which lacked not only any visualization but also the concept of a state (see below). Although these features were initially quite to the liking of Born, Heisenberg, Pauli, and Jordan, the success of Schrödinger’s work forced them to renege on their radical stance, and look for a semiclassical picture supporting their mathematics; this was a considerable U-turn (Beller, 1999; Camilleri, 2005). Heisenberg (1927) found such a picture, claiming that his uncertainty relations provided the ‘intuitive content of the quantum theoretical kinematics and mechanics’ (as his paper was called). His idea was that the classical world emerged from quantum mechanics through observation: ‘The trajectory only comes into existence because we observe it.’ <sup>40</sup><sup>40</sup>40‘Die Bahn entsteht erst dadurch, daß wir sie beobachten.’ This idea was to become extremely influential, and could be regarded as the origin of stance 1 in the Introduction.
### 2.4 Schrödinger
The history of quantum mechanics is considerably clarified by the insight that Heisenberg and Schrödinger did not, as is generally believed, discover two equivalent formulations of the theory, but rather that Heisenberg (1925) identified the mathematical nature of the observables, whereas Schrödinger (1926a) found the description of states.<sup>41</sup><sup>41</sup>41See also Muller (1997). Matrix mechanics lacked the notion of a state, but by the same token wave mechanics initially had no observables; it was only in his attempts to relate wave mechanics to matrix mechanics that Schrödinger (1926c) introduced the position and momentum operators<sup>42</sup><sup>42</sup>42Here $`j=1,2,3`$. In modern terms, the expressions on the right-hand side are unbounded operators on the Hilbert space $`=L^2(^n)`$. See Section 4 for more details. The expression $`x^i`$ is a multiplication operator, i.e. $`(x^j\mathrm{\Psi })(x)=x^j\mathrm{\Psi }(x)`$, whereas, obviously, $`(/x^j\mathrm{\Psi })(x)=(\mathrm{\Psi }/x^j)(x)`$.
$`𝒬_{\mathrm{}}(q^j)`$ $`=`$ $`x^j;`$
$`𝒬_{\mathrm{}}(p_j)`$ $`=`$ $`i\mathrm{}{\displaystyle \frac{}{x^j}}.`$ (2.2)
This provided a new basis for Schrödinger’s equation<sup>43</sup><sup>43</sup>43Or the corresponding time-independent one, with $`E\mathrm{\Psi }`$ on the right-hand side.
$$\left(\frac{\mathrm{}^2}{2m}\underset{j=1}{\overset{n}{}}\frac{^2}{x_j^2}+V(x)\right)\mathrm{\Psi }=i\mathrm{}\frac{\mathrm{\Psi }}{t},$$
(2.3)
by interpreting the left-hand side as $`H\mathrm{\Psi }`$, with $`H=𝒬_{\mathrm{}}(h)`$ in terms of the classical Hamiltonian $`h(p,q)=_jp_j^2/2m+V(q)`$. Thus Schrödinger founded the theory of the operators now named after him,<sup>44</sup><sup>44</sup>44 See Reed & Simon (1972, 1975, 1987, 1979), Cycon et al. (1987), Hislop & Sigal (1996), Hunziker & Sigal (2000), Simon (2000), Gustafson & Sigal (2003). For the mathematical origin of the Schrödinger equation also cf. Simon (1976). and in doing so gave what is still the most important example of Heisenberg’s idea of Umdeutung of classical observables.
Subsequently, correcting and expanding on certain ideas of Dirac, Pauli, and Schrödinger, von Neumann (1932) brilliantly glued these two parts together through the concept of a Hilbert space. He also gave an abstract form of the formulae of Born, Pauli, Dirac, and Jordan for the transition probabilities, thus completing the mathematical formulation of quantum mechanics.
However, this is not how Schrödinger saw his contribution. He intended wave mechanics as a full-fledged classical field theory of reality, rather than merely as one half (namely in modern parlance the state space half) of a probabilistic description of the world that still incorporated the quantum jumps he so detested (Mehra & and Rechenberg, 1987; Götsch, 1992; Bitbol & Darrigol, 1992; Bitbol, 1996; Beller, 1999). Particles were supposed to emerge in the form of wave packets, but it was immediately pointed out by Heisenberg, Lorentz, and others that in realistic situations such wave packets tend to spread in the course of time. This had initially been overlooked by Schrödinger (1926b), who had based his intuition on the special case of the harmonic oscillator. On the positive side, in the course of his unsuccessful attempts to derive classical particle mechanics from wave mechanics through the use of wave packets, Schrödinger (1926b) gave the first example of what is now called a coherent state. Here a quantum wave function $`\mathrm{\Psi }_z`$ is labeled by a ‘classical’ parameter $`z`$, in such a way that the quantum-mechanical time-evolution $`\mathrm{\Psi }_z(t)`$ is approximately given by $`\mathrm{\Psi }_{z(t)}`$, where $`z(t)`$ stands for some associated classical time-evolution; see Subsections 4.2 and 5.2 below. This has turned out to be a very important idea in understanding the transition from quantum to classical mechanics.
Furthermore, in the same paper Schrödinger (1926b) proposed the following wave-mechanical version of Bohr’s correspondence principle: classical atomic states should come from superpositions of a very large number (say at least 10,000) of highly excited states (i.e. energy eigenfunctions with very large quantum numbers). After decades of limited theoretical interest in this idea, interest in wave packets in atomic physics was revived in the late 1980s due to the development of modern experimental techniques based on lasers (such as pump-probing and phase-modulation). See Robinett (2004) for a recent technical review, or Nauenberg, Stroud, & Yeazell (1994) for an earlier popular account. Roughly speaking, the picture that has emerged is this: a localized wave packet of the said type initially follows a time-evolution with almost classical periodicity, as Schrödinger hoped, but subsequently spreads out after a number of orbits. Consequently, during this second phase the probability distribution approximately fills the classical orbit (though not uniformly). Even more surprisingly, on a much longer time scale there is a phenomenon of wave packet revival, in which the wave packet recovers its initial localization. Then the whole cycle starts once again, so that one does see periodic behaviour, but not of the expected classical type. Hence even in what naively would be thought of as the thoroughly classical regime, wave phenomena continue to play a role, leading to quite unusual and unexpected behaviour. Although a rigorous mathematical description of wave packet revival has not yet been forthcoming, the overall picture (based on both “theoretical physics” style mathematics and experiments) is clear enough.
It is debatable (and irrelevant) whether the story of wave packets has evolved according to Schrödinger’s intentions (cf. Littlejohn, 1986); what is certain is that his other main idea on the relationship between classical and quantum has been extremely influential. This was, of course, Schrödinger’s (1926a) “derivation” of his wave equation from the Hamilton–Jacobi formalism of classical mechanics. This gave rise to the WKB approximation and related methods; see Subsection 5.5.
In any case, where Schrödinger hoped for a classical interpretation of his wave function, and Heisenberg wanted to have nothing to do with it whatsoever (Beller, 1999), Born and Pauli were quick to realize its correct, probabilistic significance. Thus they deprived the wave function of its naive physical nature, and effectively degraded it to the purely mathematical status of a probability amplitude. And in doing so, Born and Pauli rendered the connection between quantum mechanics and classical mechanics almost incomprehensible once again! It was this incomprehensibility that Heisenberg addressed with his uncertainty relations.
### 2.5 von Neumann
Through its creation of the Hilbert space formalism of quantum mechanics, von Neumann’s book (1932) can be seen as a mathematical implementation of Heisenberg’s idea of Umdeutung. Von Neumann in effect proposed the following quantum-theoretical reinterpretations:
* Phase space $`M`$ $``$ Hilbert space $``$;
* Classical observable (i.e. real-valued measurable function on $`M`$) $``$ self-adjoint operator on $``$;
* Pure state (seen as point in $`M`$) $``$ unit vector (actually ray) in $``$;
* Mixed state (i.e. probability measure on $`M`$) $``$ density matrix on $``$;
* Measurable subset of $`M`$ $``$ closed linear subspace of $``$;
* Set complement $``$ orthogonal complement;
* Union of subsets $``$ closed linear span of subspaces;
* Intersection of subsets $``$ intersection of subspaces;
* Yes-no question (i.e. characteristic function on $`M`$) $``$ projection operator.<sup>45</sup><sup>45</sup>45Later on, he of course added the Umdeutung of a Boolean lattice by a modular lattice, and the ensuing Umdeutung of classical logic by quantum logic (Birkhoff & von Neumann, 1936).
Here we assume for simplicity that quantum observables $`R`$ on a Hilbert space $``$ are bounded operators, i.e. $`R()`$. Von Neumann actually derived his Umdeutung of classical mixed states as density matrices from his axiomatic characterization of quantum-mechanical states as linear maps $`\mathrm{Exp}:()`$ that satisfy $`\mathrm{Exp}(R)0`$ when $`R0`$,<sup>46</sup><sup>46</sup>46I.e., when $`R`$ is self-adjoint with positive spectrum, or, equivalently, when $`R=S^{}S`$ for some $`S()`$. $`\mathrm{Exp}(1)=1`$,<sup>47</sup><sup>47</sup>47Where the $`1`$ in $`\mathrm{Exp}(1)`$ is the unit operator on $``$., and countable additivity on a commuting set of operators. For he proved that such a map $`\mathrm{Exp}`$ is necessarily given by a density matrix $`\rho `$ according to $`\mathrm{Exp}(R)=\text{Tr}(\rho R)`$.<sup>48</sup><sup>48</sup>48This result has been widely misinterpreted (apparently also by von Neumann himself) as a theorem excluding hidden variables in quantum mechanics. See Scheibe (1991). However, Bell’s characterization of von Neumann’s linearity assumption in the definition of a state as “silly” is far off the mark, since it holds both in classical mechanics and in quantum mechanics. Indeed, von Neumann’s theorem does exclude all hidden variable extensions of quantum mechanics that are classical in nature, and it is precisely such extensions that many physicists were originally looking for. See Rédei & Stöltzner (2001) and Scheibe (2001) for recent discussions of this issue. A unit vector $`\mathrm{\Psi }`$ defines a pure state in the sense of von Neumann, which we call $`\psi `$, by $`\psi (R)=(\mathrm{\Psi },R\mathrm{\Psi })`$ for $`R()`$. Similarly, a density matrix $`\rho `$ on $``$ defines a (generally mixed) state, called $`\rho `$ as well, by $`\rho (R)=\text{Tr}(\rho R)`$. In modern terminology, a state on $`()`$ as defined by von Neumann would be called a normal state. In the $`C^{}`$-algebraic formulation of quantum physics (cf. footnote 8), this axiomatization has been maintained until the present day; here $`()`$ is replaced by more general algebras of observables in order to accommodate possible superselection rules (Haag, 1992).
Beyond his mathematical axiomatization of quantum mechanics, which (along with its subsequent extension by the $`C^{}`$-algebraic formulation) lies at the basis of all serious efforts to relate classical and quantum mechanics, von Neumann contributed to this relationship through his analysis of the measurement problem.<sup>49</sup><sup>49</sup>49Von Neumann (1932) refrained from discussing either the classical limit of quantum mechanics or (probably) the notion of quantization. In the latter direction, he declares that ‘If the quantity $``$ has the operator $`R`$, then the quantity $`f()`$ has the operator $`f(R)`$’, and that ‘If the quantities $``$, $`𝔖`$, $`\mathrm{}`$ have the operators $`R`$, $`S`$, $`\mathrm{}`$, then the quantity $`+𝔖+\mathrm{}`$ has the operator $`R+S+\mathrm{}`$’. However, despite his legendary clarity and precision, von Neumann is rather vague about the meaning of the transition $`R`$. It is tempting to construe $``$ as a classical observable whose quantum-mechanical counterpart is $`R`$, so that the above quotations might be taken as axioms for quantization. However, such an interpretation is neither supported by the surrounding text, nor by our current understanding of quantization (cf. Section 4). For example, a quantization map $`𝒬_{\mathrm{}}()`$ cannot satisfy $`f()f(𝒬_{\mathrm{}}())`$ even for very reasonable functions such as $`f(x)=x^2`$. Since here the apparent clash between classical and quantum physics comes to a head, it is worth summarizing von Neumann’s analysis of this problem here. See also Wheeler & Zurek (1983), Busch, Lahti & Mittelstaedt (1991), Auletta (2001) and Mittelstaedt (2004) for general discussions of the measurement problem.
The essence of the measurement problem is that certain states are never seen in nature, although they are not merely allowed by quantum mechanics (on the assumption of its universal validity), but are even predicted to arise in typical measurement situations. Consider a system $`S`$, whose pure states are mathematically described by normalized vectors (more precisely, rays) in a Hilbert space $`_S`$. One wants to measure an observable $`𝒪`$, which is mathematically represented by a self-adjoint operator $`O`$ on $`_S`$. Von Neumann assumes that $`O`$ has discrete spectrum, a simplification which does not hide the basic issues in the measurement problem. Hence $`O`$ has unit eigenvectors $`\mathrm{\Psi }_n`$ with real eigenvalues $`o_n`$. To measure $`𝒪`$, one couples the system to an apparatus $`A`$ with Hilbert space $`_A`$ and “pointer” observable $`𝒫`$, represented by a self-adjoint operator $`P`$ on $`_A`$, with discrete eigenvalues $`p_n`$ and unit eigenvectors $`\mathrm{\Phi }_n`$. The pure states of the total system $`S+A`$ then correspond to unit vectors in the tensor product $`_S_A`$. A good (“first kind”) measurement is then such that after the measurement, $`\mathrm{\Psi }_n`$ is correlated to $`\mathrm{\Phi }_n`$, that is, for a suitably chosen initial state $`I_A`$, a state $`\mathrm{\Psi }_nI`$ (at $`t=0`$) almost immediately evolves into $`\mathrm{\Psi }_n\mathrm{\Phi }_n`$. This can indeed be achieved by a suitable Hamiltonian.
The problem, highlighted by Schrödinger’s cat, now arises if one selects the initial state of $`S`$ to be $`_nc_n\mathrm{\Psi }_n`$ (with $`|c_n|^2=1`$), for then the superposition principle leads to the conclusion that the final state of the coupled system is $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n`$. Now, basically all von Neumann said was that if one restricts the final state to the system $`S`$, then the resulting density matrix is the mixture $`_n|c_n|^2[\mathrm{\Psi }_n]`$ (where $`[\mathrm{\Psi }]`$ is the orthogonal projection onto a unit vector $`\mathrm{\Psi }`$),<sup>50</sup><sup>50</sup>50I.e., $`[\mathrm{\Psi }]f=(\mathrm{\Psi },f)\mathrm{\Psi }`$; in Dirac notation one would have $`[\mathrm{\Psi }]=|\mathrm{\Psi }\mathrm{\Psi }|`$. so that, from the perspective of the system alone, the measurement appears to have caused a transition from the pure state $`_{n,m}c_n\overline{c_m}\mathrm{\Psi }_n\mathrm{\Psi }_m^{}`$ to the mixed state $`_n|c_n|^2[\mathrm{\Psi }_n]`$, in which interference terms $`\mathrm{\Psi }_n\mathrm{\Psi }_m^{}`$ for $`nm`$ are absent. Here the operator $`\mathrm{\Psi }_n\mathrm{\Psi }_m^{}`$ is defined by $`\mathrm{\Psi }_n\mathrm{\Psi }_m^{}f=(\mathrm{\Psi }_m,f)\mathrm{\Psi }_n`$; in particular, $`\mathrm{\Psi }\mathrm{\Psi }^{}=[\mathrm{\Psi }]`$.<sup>51</sup><sup>51</sup>51In Dirac notation one would have $`\mathrm{\Psi }_n\mathrm{\Psi }_m^{}=|\mathrm{\Psi }_n\mathrm{\Psi }_m|`$. Similarly, the apparatus, taken by itself, has evolved from the pure state $`_{n,m}c_n\overline{c_m}\mathrm{\Phi }_n\mathrm{\Phi }_m^{}`$ to the mixed state $`_n|c_n|^2[\mathrm{\Phi }_n]`$. This is simply a mathematical theorem (granted the possibility of coupling the system to the apparatus in the desired way), rather than a proposal that there exist two different time-evolutions in Nature, viz. the unitary propagation according to the Schrödinger equation side by side with the above “collapse” process.
In any case, by itself this move by no means solves the measurement problem.<sup>52</sup><sup>52</sup>52Not even in an ensemble-interpretation of quantum mechanics, which was the interpretation von Neumann unfortunately adhered to when he wrote his book. Firstly, in the given circumstances one is not allowed to adopt the ignorance interpretation of mixed states (i.e. assume that the system really is in one of the states $`\mathrm{\Psi }_n`$); cf., e.g., Mittelstaedt (2004). Secondly, even if one were allowed to do so, one could restore the problem (i.e. the original superposition $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n`$) by once again taking the other component of the system into account.
Von Neumann was well aware of at least this second point, to which he responded by his construction of a chain: one redefines $`S+A`$ as the system, and couples it to a new apparatus $`B`$, etc. This eventually leads to a post-measurement state $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n\chi _n`$ (in hopefully self-explanatory notation, assuming the vectors $`\chi _n`$ form an orthonormal set), whose restriction to $`S+A`$ is the mixed state $`_n|c_n|^2[\mathrm{\Psi }_n][\mathrm{\Phi }_n]`$. The restriction of the latter state to $`S`$ is, once again, $`_n|c_n|^2[\mathrm{\Psi }_n]`$. This procedure may evidently be iterated; the point of the construction is evidently to pass on superpositions in some given system to arbitrary systems higher up in the chain. It follows that for the final state of the original system it does not matter where one “cuts the chain” (that is, which part of the chain one leaves out of consideration), as long as it is done somewhere. Von Neumann (1932, in beautiful prose) and others suggested identifying the cutting with the act of observation, but it is preferable and much more general to simply say that some end of the chain is omitted in the description.
The burden of the measurement problem, then, is to
1. Construct a suitable chain along with an appropriate cut thereof; it doesn’t matter where the cut is made, as long as it is done.
2. Construct a suitable time-evolution accomplishing the measurement.
3. Justify the ignorance interpretation of mixed states.
As we shall see, these problems are addressed, in a conceptually different but mathematically analogous way, in the Copenhagen interpretation as well as in the decoherence approach. (The main conceptual difference will be that the latter aims to solve also the more ambitious problem of explaining the appearance of the classical world, which in the former seems to be taken for granted).
We conclude this section by saying that despite some brilliant ideas, the founders of quantum mechanics left wide open the problem of deriving classical mechanics as a certain regime of their theory.
## 3 Copenhagen: a reappraisal
The so-called “Copenhagen interpretation” of quantum mechanics goes back to ideas first discussed and formulated by Bohr, Heisenberg, and Pauli around 1927. Against the idea that there has been a “party line” from the very beginning, it has frequently been pointed out that in the late 1920s there were actually sharp differences of opinion between Bohr and Heisenberg on the interpretation of quantum mechanics and that they never really arrived at a joint doctrine (Hooker, 1972; Stapp, 1972; Hendry, 1984; Beller, 1999; Howard, 2004; Camilleri, 2005). For example, they never came to agree about the notion of complementarity (see Subsection 3.3). More generally, Heisenberg usually based his ideas on the mathematical formalism of quantum theory, whereas Bohr’s position was primarily philosophically oriented. Nonetheless, there is a clearly identifiable core of ideas on which they did agree, and since this core has everything to do with the relationship between classical and quantum, we are going to discuss it in some detail.
The principal primary sources are Bohr’s Como Lecture, his reply to epr, and his essay dedicated to Einstein (Bohr, 1927, 1935, 1949).<sup>53</sup><sup>53</sup>53 These papers were actually written in collaboration with Pauli (after first attempts with Klein), Rosenfeld, and Pais, respectively. Historical discussions of the emergence and reception of these papers are given in Bohr (1985, 1996) and in Mehra & Rechenberg (2001). As a selection of the enormous literature these papers have given rise to, we mention among relatively recent works Hooker (1972), Scheibe (1973), Folse (1985), Murdoch (1987), Lahti & Mittelstaedt (1987), Honner (1987), Chevalley (1991, 1999), Faye (1991), Faye & Folse (1994), Held (1994), Howard (1994), Beller (1999), Faye (2002), and Saunders (2004). For Bohr’s sparring partners see Heisenberg (1930, 1942, 1958, 1984a,b, 1985) with associated secondary literature (Heelan, 1965; Hörz, 1968; Geyer et al., 1993; Camilleri, 2005), and Pauli (1933, 1949, 1979, 1985, 1994), along with Laurikainen (1988) and Enz (2002).
As with Wittgenstein (and many other thinkers), it helps to understand Bohr if one makes a distinction between an “early” Bohr and a “later” Bohr.<sup>54</sup><sup>54</sup>54Here we side with Held (1994) and Beller (1999) against Howard (1994) and Suanders (2004). See also Pais (2000), p. 22: ‘Bohr’s Como Lecture did not bring the house down, however. He himself would later frown on expressions he used there, such as “disturbing the phenomena by observation”. Such language may have contributed to the considerable confusion that for so long has reigned around this subject.’ Despite a good deal of continuity in his thought (see below), the demarcation point is his response to epr (Bohr, 1935),<sup>55</sup><sup>55</sup>55This response is problematic, as is epr itself. Consequently, there exists a considerable exegetical literature on both, marked by the fact that equally competent and well-informed pairs of commentators manage to flatly contradict each other while at the same time both claiming to explain or reconstruct what Bohr “really” meant. and the main shift he made afterwards lies in his sharp insistence on the indivisible unity of object and observer after 1935, focusing on the concept of a phenomenon. Before epr, Bohr equally well believed that object and observer were both necessary ingredients of a complete description of quantum theory, but he then thought that although their interaction could never be neglected, they might at least logically be considered separately. After 1935, Bohr gradually began to claim that object and observer no longer even had separate identities, together forming a “phenomenon”. Accordingly, also his notion of complementarity changed, increasingly focusing on the idea that the specification of the experimental conditions is crucial for the unambiguous use of (necessarily) classical concepts in quantum theory (Scheibe, 1973; Held, 1994). See also Subsection 3.3 below. This development culminated in Bohr’s eventual denial of the existence of the quantum world:
> ‘There is no quantum world. There is only an abstract quantum-physical description. It is wrong to think that the task of physics is to find out how nature is. Physics concerns what we can say about nature. (…) What is it that we humans depend on? We depend on our words. Our task is to communicate experience and ideas to others. We are suspended in language.’ (quoted by Petersen (1963), p. 8.)<sup>56</sup><sup>56</sup>56See Mermin (2004) for a witty discussion of this controversial quotation.
### 3.1 The doctrine of classical concepts
Despite this shift, it seems that Bohr stuck to one key thought throughout his career:
> ‘However far the phenomena transcend the scope of classical physical explanation, the account of all evidence must be expressed in classical terms. (…) The argument is simply that by the word experiment we refer to a situation where we can tell others what we have done and what we have learned and that, therefore, the account of the experimental arrangements and of the results of the observations must be expressed in unambiguous language with suitable application of the terminology of classical physics.’ (Bohr, 1949, p. 209).
This is, in a nutshell, Bohr’s doctrine of classical concepts. Although his many drawings and stories may suggest otherwise, Bohr does not quite express the idea here that the goal of physics lies in the description of experiments.<sup>57</sup><sup>57</sup>57Which often but misleadingly has been contrasted with Einstein’s belief that the goal of physics is rather to describe reality. See Landsman (2006) for a recent discussion. In fact, he merely points out the need for “unambiguous” communication, which he evidently felt threatened by quantum mechanics.<sup>58</sup><sup>58</sup>58Here “unambiguous” means “objective” (Scheibe, 1973; Chevalley, 1991). The controversial part of the quote lies in his identification of the means of unambiguous communication with the language of classical physics, involving particles and waves and the like. We will study Bohr’s specific argument in favour of this identification shortly, but it has to be said that, like practically all his foundational remarks on quantum mechanics, Bohr presents his reasoning as self-evident, necessary, and not in need of any further analysis (Scheibe, 1973; Beller, 1999). Nonetheless, young Heisenberg clashed with Bohr on precisely this point, for Heisenberg felt that the abstract mathematical formalism of quantum theory (rather than Bohr’s world of words and pictures) provided those means of unambiguous communication.<sup>59</sup><sup>59</sup>59It is hard to disagree with Beller’s (1999) conclusion that Bohr was simply not capable of understanding the formalism of post-1925 quantum mechanics, turning his own need of understanding this theory in terms of words and pictures into a deep philosophical necessity.
By classical physics Bohr undoubtedly meant the theories of Newton, Maxwell, and Lorentz, but that is not the main point.<sup>60</sup><sup>60</sup>60Otherwise, one should wonder why one shouldn’t use the physics of Aristotle and the scholastics for this purpose, which is a much more effective way of communicating our naive impressions of the world. In contrast, the essence of physics since Newton has been to unmask a reality behind the phenomena. Indeed, Newton himself empasized that his physics was intended for those capable of natural philosophy, in contrast to ye vulgar who believed naive appearances. The fact that Aristotle’s physics is now known to be wrong should not suffice to disqualify its use for Bohr’s purposes, since the very same comment may be made about the physics of Newton etc. For Bohr, the defining property of classical physics was the property that it was objective, i.e. that it could be studied in an observer-independent way:
> ‘All description of experiences so far has been based on the assumption, already inherent in ordinary conventions of language, that it is possible to distinguish sharply between the behaviour of objects and the means of observation. This assumption is not only fully justified by everyday experience, but even constitutes the whole basis of classical physics’ (Bohr, 1958, p. 25; italics added).<sup>61</sup><sup>61</sup>61Despite the typical imperative tone of this quotation, Bohr often regarded certain other properties as essential to classical physics, such as determinism, the combined use of space-time concepts and dynamical conservation laws, and the possibility of pictorial descriptions. However, these properties were in some sense secondary, as Bohr considered them to be consequences of the possibility of isolating an object in classical physics. For example: ‘The assumption underlying the ideal of causality \[is\] that the behaviour of the object is uniquely determined, quite independently of whether it is observed or not’ (Bohr, 1937), and then again, now negatively: ‘the renunciation of the ideal of causality \[in quantum mechanics\] is founded logically only on our not being any longer in a position to speak of the autonomous behaviour of a physical object’ (Bohr, 1937). See Scheibe (1973).
See also Hooker (1972), Scheibe (1973) and Howard (1994). Heisenberg (1958, p. 55) shared this view:<sup>62</sup><sup>62</sup>62As Camilleri (2005, p. 161) states: ‘For Heisenberg, classical physics is the fullest expression of the ideal of objectivity.’
> ‘In classical physics science started from the belief - or should one say from the illusion? - that we could describe the world or at least part of the world without any reference to ourselves. This is actually possible to a large extent. We know that the city of London exists whether we see it or not. It may be said that classical physics is just that idealization in which we can speak about parts of the world without any reference to ourselves. Its success has led to the general idea of an objective description of the world.’
On the basis of his “quantum postulate” (see Subsection 3.2), Bohr came to believe that, similarly, the defining property of quantum physics was precisely the opposite, i.e. the necessity of the role of the observer (or apparatus - Bohr did not distinguish between the two and never assigned a special role to the mind of the observer or endorsed a subjective view of physics). Identifying unambiguous communication with an objective description, in turn claimed to be the essence of classical physics, Bohr concluded that despite itself quantum physics had to be described entirely in terms of classical physics. Thus his doctrine of classical concepts has an epistemological origin, arising from an analysis of the conditions for human knowledge.<sup>63</sup><sup>63</sup>63See, for example, the very title of Bohr (1958)! In that sense it may be said to be Kantian in spirit (Hooker, 1972; Murdoch, 1987; Chevalley, 1991, 1999).
Now, Bohr himself is on record as saying: ‘They do it smartly, but what counts is to do it right’ (Rosenfeld, p. 129).<sup>64</sup><sup>64</sup>64‘They’ refers to epr. The doctrine of classical concepts is certainly smart, but is it right? As we have seen, Bohr’s argument starts from the claim that classical physics is objective (or ‘unambiguous’) in being independent of the observer. In fact, nowadays it is widely believed that quantum mechanics leads to the opposite conclusion that “quantum reality” (whatever that may be) is objective (though “veiled” in the terminology of d Espagnat (1995)), while “classical reality” only comes into existence relative to a certain specification: this is stance 1 discussed in the Introduction.<sup>65</sup><sup>65</sup>65 Indeed, interesting recent attempts to make Bohr’s philosophy of quantum mechanics precise accommodate the a priori status of classical observables into some version of the modal interpretation; see Dieks (1989b), Bub (1999), Halvorson & Clifton (1999, 2002), and Dickson (2005). It should give one some confidence in the possibility of world peace that the two most hostile interpretations of quantum mechanics, viz. Copenhagen and Bohm (Cushing, 1994) have now found a common home in the modal interpretation in the sense of the authors just cited! Whether or not one agrees with Bub’s (2004) criticism of the modal interpretation, Bohr’s insistence on the necessity of classical concepts is not vindicated by any current version of it. Those who disagree with stance 1 cannot use stance 2 (of denying the fundamental nature of quantum theory) at this point either, as that is certainly not what Bohr had in mind. Unfortunately, in his most outspoken defence of Bohr, even Heisenberg (1958, p. 55) was unable to find a better argument for Bohr’s doctrine than the lame remark that ‘the use of classical concepts is finally a consequence of the general human way of thinking.’<sup>66</sup><sup>66</sup>66And similarly: ‘We are forced to use the language of classical physics, simply because we have no other language in which to express the results.’ (Heisenberg, 1971, p. 130). This in spite of the fact that the later Heisenberg thought about this matter very deeply; see, e.g., his (1942), as well as Camilleri (2005). Murdoch (1987, pp. 207–210) desperately tries to boost the doctrine of classical concepts into a profound philosophical argument by appealing to Strawson (1959).
In our opinion, Bohr’s motivation for his doctrine has to be revised in the light of our current understanding of quantum theory; we will do so in Subsection 3.4. In any case, whatever its motivation, the doctrine itself seems worth keeping: apart from the fact that it evidently describes experimental practice, it provides a convincing explanation for the probabilistic nature of quantum mechanics (cf. the next subsection).
### 3.2 Object and apparatus: the Heisenberg cut
Describing quantum physics in terms of classical concepts sounds like an impossible and even self-contradictory task (cf. Heisenberg, 1958). For one, it precludes a completely quantum-mechanical description of the world: ‘However far the phenomena transcend the scope of classical physical explanation, the account of all evidence must be expressed in classical terms.’ But at the same time it precludes a purely classical description of the world, for underneath classical physics one has quantum theory.<sup>67</sup><sup>67</sup>67This peculiar situation makes it very hard to give a realist account of the Copenhagen interpretation, since quantum reality is denied whereas classical reality is neither fundamental nor real. The fascination of Bohr’s philosophy of quantum mechanics lies precisely in his brilliant resolution of this apparently paradoxical situation.
The first step of this resolution that he and Heisenberg proposed is to divide the system whose description is sought into two parts: one, the object, is to be described quantum-mechanically, whereas the other, the apparatus, is treated as if it were classical. Despite innumerable claims to the contrary in the literature (i.e. to the effect that Bohr held that a separate realm of Nature was intrinsically classical), there is no doubt that both Bohr and Heisenberg believed in the fundamental and universal nature of quantum mechanics, and saw the classical description of the apparatus as a purely epistemological move without any counterpart in ontology, expressing the fact that a given quantum system is being used as a measuring device.<sup>68</sup><sup>68</sup>68See especially Scheibe (1973) on Bohr, and Heisenberg (1958). The point in question has also been made by R. Haag (who knew both Bohr and Heisenberg) in most of his talks on quantum mechanics in the 1990s. In this respect we disagree with Howard (1994), who claims that according to Bohr a classical description of an apparatus amounts to picking a particular (maximally) abelian subalgebra of its quantum-mechanical algebra of ‘beables’, which choice is dictated by the measurement context. But having a commutative algebra falls far short of a classical description, since in typical examples one obtains only half of the canonical classical degrees of freedom in this way. Finding a classical description of a quantum-mechanical system is a much deeper problem, to which we shall return throughout this paper. For example: ‘The construction and the functioning of all apparatus like diaphragms and shutters, serving to define geometry and timing of the experimental arrangements, or photographic plates used for recording the localization of atomic objects, will depend on properties of materials which are themselves essentially determined by the quantum of action’ (Bohr, 1948), as well as: ‘We are free to make the cut only within a region where the quantum mechanical description of the process concerned is effectively equivalent with the classical description’ (Bohr, 1935).<sup>69</sup><sup>69</sup>69This last point suggests that the cut has something to do with the division between a microscopic and a macroscopic realm in Nature, but although this division often facilitates making the cut when it is well defined, this is by no means a matter of principle. Cf. Howard (1994). In particular, all objections to the Copenhagen interpretation to the effect that the interpretation is ill-defined because the micro-macro distinction is blurred are unfounded.
The separation between object and apparatus called for here is usually called the Heisenberg cut, and it plays an absolutely central role in the Copenhagen interpretation of quantum mechanics.<sup>70</sup><sup>70</sup>70Pauli (1949) went as far as saying that the Heisenberg cut provides the appropriate generalization modern physics offers of the old Kantian opposition between a knowable object and a knowing subject: ’Auf diese Weise verallgemeinert die moderne Physik die alte Gegenüberstellung von erkennenden Subjekt auf der einen Seite und des erkannten Objektes auf der anderen Seite zu der Idee des Schnittes zwischen Beobachter oder Beobachtungsmittel und dem beobachten System.’ (‘In this way, modern physics generalizes the old opposition between the knowing subject on the one hand and the known object on the other to the idea of the cut between observer or means of observation and the observed system.’) He then continued calling the cut a necessary condition for human knowledge: see footnote 74. The idea, then, is that a quantum-mechanical object is studied exclusively through its influence on an apparatus that is described classically. Although described classically, the apparatus is a quantum system, and is supposed to be influenced by its quantum-mechanical coupling to the underlying (quantum) object.
The alleged necessity of including both object and apparatus in the description was initially claimed to be a consequence of the so-called “quantum postulate”. This notion played a key role in Bohr’s thinking: his Como Lecture (Bohr, 1927) was even entitled ‘The quantum postulate and the recent development of atomic theory’. There he stated its contents as follows: ‘The essence of quantum theory is the quantum postulate: every atomic process has an essential discreteness - completely foreign to classical theories - characterized by Planck s quantum of action.’<sup>71</sup><sup>71</sup>71Instead of ‘discreteness’, Bohr alternatively used the words ‘discontinuity’ or ‘individuality’ as well. He rarely omitted amplifications like ‘essential’. Even more emphatically, in his reply to epr (Bohr, 1935): ‘Indeed the finite interaction between object and measuring agencies conditioned by the very existence of the quantum of action entails - because of the impossibility of controlling the reaction of the object on the measurement instruments if these are to serve their purpose - the necessity of a final renunication of the classical ideal of causality and a radical revision of our attitude towards the problem of physical reality.’ Also, Heisenberg’s uncertainty relations were originally motivated by the quantum postulate in the above form. According to Bohr and Heisenberg around 1927, this ‘essential discreteness’ causes an ‘uncontrollable disturbance’ of the object by the apparatus during their interaction. Although the “quantum postulate” is not supported by the mature mathematical formalism of quantum mechanics and is basically obsolete, the intuition of Bohr and Heisenberg that a measurement of a quantum-mechanical object causes an ‘uncontrollable disturbance’ of the latter is actually quite right.<sup>72</sup><sup>72</sup>72Despite the fact that Bohr later distanced himself from it; cf. Beller (1999) and footnote 54 above. In a correct analysis, what is disturbed upon coupling to a classical apparatus is the quantum-mechanical state of the object (rather than certain sharp values of classical observables such as position and momentum, as the early writings of Bohr and Heisenberg suggest).
In actual fact, the reason for this disturbance does not lie in the “quantum postulate”, but in the phenomenon of entanglement, as further discussed in Subsection 3.4. Namely, from the point of view of von Neumann’s measurement theory (see Subsection 2.5) the Heisenberg cut is just a two-step example of a von Neumann chain, with the special feature that after the quantum-mechanical interaction has taken place, the second link (i.e. the apparatus) is described classically. The latter feature not only supports Bohr’s philosophical agenda, but, more importantly, also suffices to guarantee the applicability of the ignorance interpretation of the mixed state that arises after completion of the measurement.<sup>73</sup><sup>73</sup>73In a purely quantum-mechanical von Neumann chain the final state of system plus apparatus is pure, but if the apparatus is classical, then the post-measurement state is mixed. All of von Neumann’s analysis of the arbitrariness of the location of the cut applies here, for one may always extend the definition of the quantum-mechanical object by coupling the original choice to any other purely quantum-mechanical system one likes, and analogously for the classical part. Thus the two-step nature of the Heisenberg cut includes the possibility that the first link or object is in fact a lengthy chain in itself (as long as it is quantum-mechanical), and similarly for the second link (as long as it is classical).<sup>74</sup><sup>74</sup>74 Pauli (1949) once more: ’Während die Existenz eines solchen Schnittes eine notwendige Bedingung menschlicher Erkenntnis ist, faßt sie die Lage des Schnittes als bis zu einem gewissen Grade willkürlich und als Resultat einer durch Zweckmäßigkeitserwägungen mitbestimmten, also teilweise freien Wahl auf.’ (‘While the existence of such a \[Heisenberg\] cut is a necessary condition for human knowledge, its location is to some extent arbitrary as a result of a pragmatic and thereby partly free choice.’) This arbitrariness, subject to the limitation expressed by the second (1935) Bohr quote in this subsection, was well recognized by Bohr and Heisenberg, and was found at least by Bohr to be of great philosophical importance.
It is the interaction between object and apparatus that causes the measurement to ‘disturb’ the former, but it is only and precisely the classical description of the latter that (through the ignorance interpretation of the final state) makes the disturbance ‘uncontrollable’.<sup>75</sup><sup>75</sup>75These points were not clearly separated by Heisenberg (1927) in his paper on the uncertainty relations, but were later clarified by Bohr. See Scheibe (1973). In the Copenhagen interpretation, probabilities arise solely because we look at the quantum world through classical glasses.
> ‘Just the necessity of accounting for the function of the measuring agencies on classical lines excludes in principle in proper quantum phenomena an accurate control of the reaction of the measuring instruments on the atomic objects.’ (Bohr, 1956, p. 87)
> ‘One may call these uncertainties objective, in that they are simply a consequence of the fact that we describe the experiment in terms of classical physics; they do not depend in detail on the observer. One may call them subjective, in that they reflect our incomplete knowledge of the world.’ (Heisenberg, 1958, pp. 53–54)
Thus the picture that arises is this: Although the quantum-mechanical side of the Heisenberg cut is described by the Schrödinger equation (which is deterministic), while the classical side is subject to Newton’s laws (which are equally well deterministic),<sup>76</sup><sup>76</sup>76But see Earman (1986, 2005). unpredictability arises because the quantum system serving as an apparatus is approximated by a classical system. The ensuing probabilities reflect the ignorance arising from the decision (or need) to ignore the quantum-mechanical degrees of freedom of the apparatus. Hence the probabilistic nature of quantum theory is not intrinsic but extrinsic, and as such is entirely a consequence of the doctrine of classical concepts, which by the same token explains this nature.
Mathematically, the simplest illustration of this idea is as follows. Take a finite-dimensional Hilbert space $`=^n`$ with the ensuing algebra of observables $`𝒜=M_n()`$ (i.e. the $`n\times n`$ matrices). A unit vector $`\mathrm{\Psi }^n`$ determines a quantum-mechanical state in the usual way. Now describe this quantum system as if it were classical by ignoring all observables except the diagonal matrices. The state then immediately collapses to a probability measure on the set of $`n`$ points, with probabilities given by the Born rule $`p(i)=|(e_i,\mathrm{\Psi })|^2`$, where $`(e_i)_{i=1,\mathrm{},n}`$ is the standard basis of $`^n`$.
Despite the appeal of this entire picture, it is not at all clear that it actually applies! There is no a priori guarantee whatsoever that one may indeed describe a quantum system “as if it were classical”. Bohr and Heisenberg apparently took this possibility for granted, probably on empirical grounds, blind to the extremely delicate theoretical nature of their assumption. It is equally astounding that they never reflected in print on the question if and how the classical worlds of mountains and creeks they loved so much emerges from a quantum-mechanical world. In our opinion, the main difficulty in making sense of the Copenhagen interpretation is therefore not of a philosophical nature, but is a mathematical one. This difficulty is the main topic of this paper, of which Section 6 is of particular relevance in the present context.
### 3.3 Complementarity
The notion of a Heisenberg cut is subject to a certain arbitrariness even apart from the precise location of the cut within a given chain, for one might in principle construct the chain in various different and incompatible ways. This arbitrariness was analyzed by Bohr in terms of what he called complementarity.<sup>77</sup><sup>77</sup>77Unfortunately and typically, Bohr once again presented complementarity as a necessity of thought rather than as the truly amazing possible mode of description it really is.
Bohr never gave a precise definition of complementarity,<sup>78</sup><sup>78</sup>78Perhaps he preferred this approach because he felt a definition could only reveal part of what was supposed to be defined: one of his favourite examples of complementarity was that between definition and observation. but restricted himself to the analysis of a number of examples.<sup>79</sup><sup>79</sup>79We refrain from discussing the complementarity between truth and clarity, science and religion, thoughts and feelings, and objectivity and introspection here, despite the fact that on this basis Bohr’s biographer Pais (1997) came to regard his subject as the greatest philosopher since Kant. A prominent such example is the complementarity between a “causal” <sup>80</sup><sup>80</sup>80 Bohr’s use the word “causal” is quite confusing in view of the fact that in the British empiricist tradition causality is often interpreted in the sense of a space-time description. But Bohr’s “causal” is meant to be complementary to a space-time description! description of a quantum system in which conservation laws hold, and a space-time description that is necessarily statistical in character. Here Bohr’s idea seems to have been that a stationary state (i.e. an energy eigenstate) of an atom is incompatible with an electron moving in its orbit in space and time - see Subsection 5.4 for a discussion of this issue. Heisenberg (1958), however, took this example of complementarity to mean that a system on which no measurement is performed evolves deterministically according to the Schrödinger equation, whereas a rapid succession of measurements produces a space-time path whose precise form quantum theory is only able to predict statistically (Camilleri, 2005). In other words, this example reproduces precisely the picture through which Heisenberg (1927) believed he had established the connection between classical and quantum mechanics; cf. Subsection 2.3.
Bohr’s other key example was the complementarity between particles and waves. Here his principal aim was to make sense of Young’s double-slit experiment. The well-known difficulty with a classical visualization of this experiment is that a particle description appears impossible because a particle has to go through a single slit, ruining the interference pattern gradually built up on the detection screen, whereas a wave description seems incompatible with the point-like localization on the screen once the wave hits it. Thus Bohr suggested that whilst each of these classical descriptions is incomplete, the union of them is necessary for a complete description of the experiment.
The deeper epistemological point appears to be that although the completeness of the quantum-mechanical description of the microworld systems seems to be endangered by the doctrine of classical concepts, it is actually restored by the inclusion of two “complementary” descriptions (i.e. of a given quantum system plus a measuring device that is necessarily described classicaly, ‘if it is to serve its purpose’). Unfortunately, despite this attractive general idea it is unclear to what precise definition of complementarity Bohr’s examples should lead. In the first, the complementary notions of determinism and a space-time description are in mutual harmony as far as classical physics is concerned, but are apparently in conflict with each other in quantum mechanics. In the second, however, the wave description of some entity contradicts a particle description of the same entity precisely in classical physics, whereas in quantum mechanics these descriptions somehow coexist.<sup>81</sup><sup>81</sup>81On top of this, Bohr mixed these examples in conflicting ways. In discussing bound states of electrons in an atom he jointly made determinism and particles one half of a complementary pair, waves and space-time being the other. In his description of electron-photon scattering he did it the other way round: this time determinism and waves formed one side, particles and space-time the other (cf. Beller, 1999).
Scheibe (1973, p. 32) notes a ‘clear convergence \[in the writings of Bohr\] towards a preferred expression of a complementarity between phenomena’, where a Bohrian phenomenon is an indivisible union (or “whole”) of a quantum system and a classically described experimental arrangement used to study it; see item 2 below. Some of Bohr’s early examples of complementarity can be brought under this heading, others cannot (Held, 1994). For many students of Bohr (including the present author), the fog has yet to clear up.<sup>82</sup><sup>82</sup>82Even Einstein (1949, p. 674) conceded that throughout his debate with Bohr he had never understood the notion of complementarity, ‘the sharp formulation of which, moreover, I have been unable to achieve despite much effort which I have expended on it.’ See Landsman (2006) for the author’s view on the Bohr–Einstein debate. Nonetheless, the following mathematical interpretations might assign some meaning to the idea of complementarity in the framework of von Neumann’s formalism of quantum mechanics.<sup>83</sup><sup>83</sup>83This exercise is quite against the spirit of Bohr, who is on record as saying that ‘von Neumann’s approach (…) did not solve problems but created imaginary difficulties (Scheibe, 1973, p. 11, quoting Feyerabend; italics in original).
1. Heisenberg (1958) identified complementary pictures of a quantum-mechanical system with equivalent mathematical representations thereof. For example, he thought of the complementarity of $`x`$ and $`p`$ as the existence of what we now call the Schrödinger representations of the canonical commutation relations (CCR) on $`L^2(^n)`$ and its Fourier transform to momentum space. Furthermore, he felt that in quantum field theory particles and waves gave two equivalent modes of description of quantum theory because of second quantization. Thus for Heisenberg complementary pictures are classical because there is an underlying classical variable, with no apparatus in sight, and such pictures are not mutually contradictory but (unitarily) equivalent. See also Camilleri (2005, p. 88), according to whom ‘Heisenberg never accepted Bohr’s complementarity arguments’.
2. Pauli (1933) simply stated that two observables are complementary when the corresponding operators fail to commute.<sup>84</sup><sup>84</sup>84 More precisely, one should probably require that the two operators in question generate the ambient algebra of observables, so that complementarity in Pauli’s sense is really defined between two commutative subalgebras of a given algebra of observables (again, provided they jointly generate the latter). Consequently, it then follows from Heisenberg’s uncertainty relations that complementary observables cannot be measured simultaneously with arbitrary precision. This suggests (but by no means proves) that they should be measured independently, using mutually exclusive experimental arrangements. The latter feature of complementarity was emphasized by Bohr in his later writings.<sup>85</sup><sup>85</sup>85We follow Held (1994) and others. Bohr’s earlier writings do not quite conform to Pauli’s approach. In Bohr’s discussions of the double-slit experiment particle and wave form a complementary pair, whereas Pauli’s complementary observables are position and momentum, which refer to a single side of Bohr’s pair. This approach makes the notion of complementarity unambiguous and mathematically precise, and perhaps for this reason the few physicists who actually use the idea of complementarity in their work tend to follow Pauli and the later Bohr.<sup>86</sup><sup>86</sup>86Adopting this point of view, it is tempting to capture the complementarity between position and momentum by means of the following conjecture: Any normal pure state $`\omega `$ on $`(L^2(^n))`$ (that is, any wave function seen as a state in the sense of $`C^{}`$-algebras) is determined by the pair $`\{\omega |L^{\mathrm{}}(^n),\omega |FL^{\mathrm{}}(^n)F^1\}`$ (in other words, by its restrictions to position and momentum). Here $`L^{\mathrm{}}(^n)`$ is the von Neumann algebra of multiplication operators on $`L^2(^n)`$, i.e. the von Neumann algebra generated by the position operator, whereas $`FL^{\mathrm{}}(^n)F^1`$ is its Fourier transform, i.e. the von Neumann algebra generated by the momentum operator. The idea is that each of its restrictions $`\omega |L^{\mathrm{}}(^n)`$ and $`\omega |FL^{\mathrm{}}(^n)F^1`$ gives a classical picture of $`\omega `$. These restrictions are a measure on $`^n`$ interpreted as position space, and another measure on $`^n`$ interpreted as momentum space. Unfortunately, this conjecture is false. The following counterexample was provided by D. Buchholz (private communication): take $`\omega `$ as the state defined by the wave function $`\mathrm{\Psi }(x)\mathrm{exp}(ax^2/2)`$ with $`\mathrm{Re}(a)>0`$, $`\mathrm{Im}(a)0`$, and $`|a|^2=1`$. Then $`\omega `$ depends on $`\mathrm{Im}(a)`$, whereas neither $`\omega |L^{\mathrm{}}(^n)`$ nor $`\omega |FL^{\mathrm{}}(^n)F^1`$ does. There is even a counterexample to the analogous conjecture for the $`C^{}`$-algebra of $`2\times 2`$ matrices, found by H. Halvorson: if $`A`$ is the commutative $`C^{}`$-algebra generated by $`\sigma _x`$, and $`B`$ the one generated by $`\sigma _y`$, then the two different eigenstates of $`\sigma _z`$ coincide on $`A`$ and on $`B`$. One way to improve our conjecture might be to hope that if, in the Schrödinger picture, two states coincide on the two given commutative von Neumann algebras for all times, then they must be equal. But this can only be true for certain “realistic” time-evolutions, for the trivial Hamiltonian $`H=0`$ yields the above counterexample. We leave this as a problem for future research. At the time of writing, Halvorson (2004) contains the only sound mathematical interpretation of the complementarity between position and momentum, by relating it to the representation theory of the CCR. He shows that in any representation where the position operator has eigenstates, there is no momentum operator, and vice versa.
3. The present author proposes that observables and pure states are complementary. For in the Schrödinger representation of elementary quantum mechanics, the former are, roughly speaking, generated by the position and momentum operators, whereas the latter are given by wave functions. Some of Bohr’s other examples of complementarity also square with this interpretation (at least if one overlooks the collapse of the wavefunction upon a measurement). Here one captures the idea that both ingredients of a complementary pair are necessary for a complete description, though the alleged mutual contradiction between observables and states is vague. Also, this reading of complementarity relies on a specific representation of the canonical commutation relations. It is not quite clear what one gains with this ideology, but perhaps it deserves to be developed in some more detail. For example, in quantum field theory it is once more the observables that carry the space-time description, especially in the algebraic description of Haag (1992).
### 3.4 Epilogue: entanglement to the rescue?
Bohr’s “quantum postulate” being obscure and obsolete, it is interesting to consider Howard’s (1994) ‘reconstruction’ of Bohr’s philosophy of physics on the basis of entanglement.<sup>87</sup><sup>87</sup>87We find little evidence that Bohr himself ever thought along those lines. With approval we quote Zeh, who, following a statement of the quantum postulate by Bohr similar to the one in Subsection 3.2 above, writes: ‘The later revision of these early interpretations of quantum theory (required by the important role of entangled quantum states for much larger systems) seems to have gone unnoticed by many physicists.’ (Joos et al., 2003, p. 23.) See also Howard (1990) for an interesting historical perspective on entanglement, and cf. Raimond, Brune, & Haroche (2001) for the experimental situation. His case can perhaps be strengthened by an appeal to the analysis Primas (1983) has given of the need for classical concepts in quantum physics.<sup>88</sup><sup>88</sup>88See also Amann & Primas (1997) and Primas (1997). Primas proposes to define a “quantum object” as a physical system $`𝒮`$ that is free from what he calls “epr-correlations” with its environment. Here the “environment” is meant to include apparatus, observer, the rest of the universe if necessary, and what not. In elementary quantum mechanics, quantum objects in this sense exist only in very special states: if $`_S`$ is the Hilbert space of the system $`S`$, and $`_E`$ that of the environment $`E`$, any pure state of the form $`_ic_i\mathrm{\Psi }_i\mathrm{\Phi }_i`$ (with more than one term) by definition correlates $`S`$ with $`E`$; the only uncorrelated pure states are those of the form $`\mathrm{\Psi }\mathrm{\Phi }`$ for unit vectors $`\mathrm{\Psi }_S`$, $`\mathrm{\Phi }_E`$. The restriction of an epr-correlated state on $`S+E`$ to $`S`$ is mixed, so that the (would-be) quantum object ‘does not have its own pure state’; equivalently, the restriction of an epr-correlated state $`\omega `$ to $`S`$ together with its restriction to $`E`$ do not jointly determine $`\omega `$. Again in other words, if the state of the total $`S+E`$ is epr-correlated, a complete characterization of the state of $`S`$ requires $`E`$ (and vice versa). But (against Bohr!) mathematics defeats words: the sharpest characterization of the notion of epr-correlations can be given in terms of operator algebras, as follows. In the spirit of the remainder of the paper we proceed in a rather general and abstract way.<sup>89</sup><sup>89</sup>89Though Summers & Werner (1987) give even more general results, where the tensor product $`𝒜\widehat{}`$ below is replaced by an arbitrary $`C^{}`$-algebra $`𝒞`$ containing $`𝒜`$ and $``$ as $`C^{}`$-subalgebras.
Let $`𝒜`$ and $``$ be $`C^{}`$-algebras,<sup>90</sup><sup>90</sup>90 Recall that a $`C^{}`$-algebra is a complex algebra $`𝒜`$ that is complete in a norm $``$ that satisfies $`ABAB`$ for all $`A,B𝒜`$, and has an involution $`AA^{}`$ such that $`A^{}A=A^2`$. A basic examples is $`𝒜=()`$, the algebra of all bounded operators on a Hilbert space $``$, equipped with the usual operator norm and adjoint. By the Gelfand–Naimark theorem, any $`C^{}`$-algebra is isomorphic to a norm-closed self-adjoint subalgebra of $`()`$, for some Hilbert space $``$. Another key example is $`𝒜=C_0(X)`$, the space of all continuous complex-valued functions on a (locally compact Hausdorff) space $`X`$ that vanish at infinity (in the sense that for every $`\epsilon >0`$ there is a compact subset $`KX`$ such that $`|f(x)|<\epsilon `$ for all $`xK`$), equipped with the supremum norm $`f_{\mathrm{}}:=sup_{xX}|f(x)|`$, and involution given by (pointwise) complex conjugation. By the Gelfand–Naimark lemma, any commutative $`C^{}`$-algebra is isomorphic to $`C_0(X)`$ for some locally compact Hausdorff space $`X`$. with tensor product $`𝒜\widehat{}`$.<sup>91</sup><sup>91</sup>91 The tensor product of two (or more) $`C^{}`$-algebras is not unique, and we here need the so-called projective tensor product $`𝒜\widehat{}`$, defined as the completion of the algebraic tensor product $`𝒜`$ in the maximal $`C^{}`$-cross-norm. The choice of the projective tensor product guarantees that each state on $`𝒜`$ extends to a state on $`𝒜\widehat{}`$ by continuity; conversely, since $`𝒜`$ is dense in $`𝒜\widehat{}`$, each state on the latter is uniquely determined by its values on the former. See Wegge-Olsen (1993), Appendix T, or Takesaki (2003), Vol. i, Ch. iv. In particular, product states $`\rho \sigma `$ and mixtures $`\omega =_ip_i\rho _i\sigma _i`$ thereof as considered below are well defined on $`𝒜\widehat{}`$. If $`𝒜(_S)`$ and $`(_E)`$ are von Neumann algebras, as in the analysis of Raggio (1981, 1988), it is easier (and sufficient) to work with the spatial tensor product $`𝒜\overline{}`$, defined as the double commutant (or weak completion) of $`𝒜`$ in $`(_S_E)`$. For any normal state on $`𝒜`$ extends to a normal state on $`𝒜\overline{}`$ by continuity. Less abstractly, just think of two Hilbert spaces $`_S`$ and $`_E`$ as above, with tensor product $`_S_E`$, and assume that $`𝒜=(_S)`$ while $``$ is either $`(_E)`$ itself or some (norm-closed and involutive) commutative subalgebra thereof. The tensor product $`𝒜\widehat{}`$ is then a (norm-closed and involutive) subalgebra of $`(_S_E)`$, the algebra of all bounded operators on $`_S_E`$.
A product state on $`𝒜\widehat{}`$ is a state of the form $`\omega =\rho \sigma `$, where the states $`\rho `$ on $`𝒜`$ and $`\sigma `$ on $``$ may be either pure or mixed.<sup>92</sup><sup>92</sup>92 We use the notion of a state that is usual in the algebraic framework. Hence a state on a $`C^{}`$-algebra $`𝒜`$ is a linear functional $`\rho :𝒜`$ that is positive in that $`\rho (A^{}A)0`$ for all $`A𝒜`$ and normalized in that $`\rho (1)=1`$, where $`1`$ is the unit element of $`𝒜`$. If $`𝒜`$ is a von Neumann algebra, one has the notion of a normal state, which satisfies an additional continuity condition. If $`𝒜=()`$, then a fundamental theorem of von Neumann states that each normal state $`\rho `$ on $`𝒜`$ is given by a density matrix $`\widehat{\rho }`$ on $``$, so that $`\rho (A)=\text{Tr}(\widehat{\rho }A)`$ for each $`A𝒜`$. In particular, a normal pure state on $`()`$ (seen as a von Neumann algebra) is necessarily of the form $`\psi (A)=(\mathrm{\Psi },A\mathrm{\Psi })`$ for some unit vector $`\mathrm{\Psi }`$. We say that a state $`\omega `$ on $`𝒜\widehat{}`$ is decomposable when it is a mixture of product states, i.e. when $`\omega =_ip_i\rho _i\sigma _i`$, where the coefficients $`p_i>0`$ satisfy $`_ip_i=1`$.<sup>93</sup><sup>93</sup>93Infinite sums are allowed here. More precisely, $`\omega `$ is decomposable if it is in the $`w^{}`$-closure of the convex hull of the product states on $`𝒜\widehat{}`$. A decomposable state $`\omega `$ is pure precisely when it is a product of pure states. This has the important consequence that both its restrictions $`\omega _{|𝒜}`$ and $`\omega _|`$ to $`𝒜`$ and $``$, respectively, are pure as well.<sup>94</sup><sup>94</sup>94The restriction $`\omega _{|𝒜}`$ of a state $`\omega `$ on $`𝒜\widehat{}`$ to, say, $`𝒜`$ is given by $`\omega _{|𝒜}(A)=\omega (A1)`$, where $`1`$ is the unit element of $``$, etc. On the other hand, a state on $`𝒜\widehat{}`$ may be said to be epr-correlated (Primas, 1983) when it is not decomposable. An epr-correlated pure state has the property that its restriction to $`𝒜`$ or $``$ is mixed.
Raggio (1981) proved that each normal state on $`𝒜\widehat{}`$ is decomposable if and only if $`𝒜`$ or $``$ is commutative. In other words, epr-correlated states exist if and only if $`𝒜`$ and $``$ are both noncommutative.<sup>95</sup><sup>95</sup>95Raggio (1981) proved this for von Neumann algebras and normal states. His proof was adapted to $`C^{}`$-algebras by Bacciagaluppi (1993). As one might expect, this result is closely related to the Bell inequalities. Namely, the Bell-type (or Clauser–Horne) inequality
$$sup\{\omega (A_1(B_1+B_2)+A_2(B_1B_2))\}2,$$
(3.1)
where for a fixed state $`\omega `$ the supremum is taken over all self-adjoint operators $`A_1,A_2𝒜`$, $`B_1,B_2`$, each of norm $`1`$, holds if and only if $`\omega `$ is decomposable (Baez, 1987; Raggio, 1988). Consequently, the inequality (3.1) can only be violated in some state $`\omega `$ when the algebras $`𝒜`$ and $``$ are both noncommutative. If, on the other hand, (3.1) is satisfied, then one knows that there exists a classical probability space and probability measure (and hence a “hidden variables” theory) reproducing the given correlations (Pitowsky, 1989). As stressed by Bacciagaluppi (1993), such a description does not require the entire setting to be classical; as we have seen, only one of the algebras $`𝒜`$ and $``$ has to be commutative for the Bell inequalities to hold.
Where does this leave us with respect to Bohr? If we follow Primas (1983) in describing a (quantum) object as a system free from epr-correlations with its environment, then the mathematical results just reviewed leave us with two possibilities. Firstly, we may pay lip-service to Bohr in taking the algebra $``$ (interpreted as the algebra of observables of the environment in the widest possible sense, as above) to be commutative as a matter of description. In that case, our object is really an “object” in any of its states. But clarly this is not the only possibility. For even in the case of elementary quantum mechanics - where $`𝒜=(_S)`$ and $`=(_E)`$ \- the system is still an “object” in the sense of Primas as long as the total state $`\omega `$ of $`S+E`$ is decomposable. In general, for pure states this just means that $`\omega =\psi \varphi `$, i.e. that the total state is a product of pure states. To accomplish this, one has to define the Heisenberg cut in an appropriate way, and subsequently hope that the given product state remains so under time-evolution (see Amann & Primas (1997) and Atmanspacher, Amann & Müller-Herold, 1999, and references therein). This selects certain states on $`𝒜`$ as “robust” or “stable”, in much the same way as in the decoherence approach. We therefore continue this discussion in Section 7 (see especially point 6 in Subsection 7.1).
## 4 Quantization
Heisenberg’s (1925) idea of Umdeutung (reinterpretation) suggests that it is possible to construct a quantum-mechanical description of a physical system whose classical description is known. As we have seen, this possibility was realized by Schrödinger (1925c), who found the simplest example (2.2) and (2.3) of Umdeutung in the context of atomic physics. This early example was phenomenally successful, as almost all of atomic and molecular physics is still based on it.
Quantization theory is an attempt to understand this example, make it mathematically precise, and generalize it to more complicated systems. It has to be stated from the outset that, like the entire classical-quantum interface, the nature of quantization is not yet well understood. This fact is reflected by the existence of a fair number of competing quantization procedures, the most transparent of which we will review below.<sup>96</sup><sup>96</sup>96The path integral approach to quantization is still under development and so far has had no impact on foundational debates, so we will not discuss it here. See Albeverio & Høegh-Krohn (1976) and Glimm & Jaffe (1987). Among the first mathematically serious discussions of quantization are Mackey (1968) and Souriau (1969); more recent and comprehensive treatments are, for example, Woodhouse (1992), Landsman (1998), and Ali & Englis (2004).
### 4.1 Canonical quantization and systems of imprimitivity
The approach based on (2.2) is often called canonical quantization. Even apart from the issue of mathematical rigour, one can only side with Mackey (1992, p. 283), who wrote: ‘Simple and elegant as this model is, it appears at first sight to be quite arbitrary and ad hoc. It is difficult to understand how anyone could have guessed it and by no means obvious how to modify it to fit a model for space different from $`^r`$.’
One veil of the mystery of quantization was lifted by von Neumann (1931), who (following earlier heuristic proposals by Heisenberg, Schrödinger, Dirac, and Pauli) recognized that (2.2) does not merely provide a representation of the canonical commutation relations
$$[𝒬_{\mathrm{}}(p_j),𝒬_{\mathrm{}}(q^k)]=i\mathrm{}\delta _j^k,$$
(4.1)
but (subject to a regularity condition)<sup>97</sup><sup>97</sup>97It is required that the unbounded operators $`𝒬_{\mathrm{}}(p_j)`$ and $`𝒬_{\mathrm{}}(q^k)`$ integrate to a unitary representation of the $`2n+1`$-dimensional Heisenberg group $`H_n`$, i.e. the unique connected and simply connected Lie group with $`2n+1`$-dimensional Lie algebra with generators $`X_i,Y_i,Z`$ ($`i=1,\mathrm{},n`$) subject to the Lie brackets $`[X_i,X_j]=[Y_i,Y_j]=0`$, $`[X_i,Y_j]=\delta _{ij}Z`$, $`[X_i,Z]=[Y_i,Z]=0`$. Thus von Neumann’s uniqueness theorem for representations of the canonical commutation relations is (as he indeed recognized himself) really a uniqueness theorem for unitary representations of $`H_n`$ for which the central element $`Z`$ is mapped to $`i\mathrm{}^11`$, where $`\mathrm{}0`$ is a fixed constant. See, for example, Corwin & Greenleaf (1989) or Landsman (1998). is the only such representation that is irreducible (up to unitary equivalence). In particular, the seemingly different formulations of quantum theory by Heisenberg and Schrödinger (amended by the inclusion of states and of observables, respectively - cf. Section 2) simply involved superficially different but unitarily equivalent representations of (4.1): the difference between matrices and waves was just one between coordinate systems in Hilbert space, so to speak. Moreover, any other conceivable formulation of quantum mechanics - now simply defined as a (regular) Hilbert space representation of (4.1) - has to be equivalent to the one of Heisenberg and Schrödinger.<sup>98</sup><sup>98</sup>98This is unrelated to the issue of the Heisenberg picture versus the Schrödinger picture, which is about the time-evolution of observables versus that of states.
This, then, transfers the quantization problem of a particle moving on $`^n`$ to the canonical commutation relations (4.1). Although a mathematically rigorous theory of these commutation relations (as they stand) exists (Jørgensen,& Moore, 1984; Schmüdgen, 1990), they are problematic nonetheless. Firstly, technically speaking the operators involved are unbounded, and in order to represent physical observables they have to be self-adjoint; yet on their respective domains of self-adjointness the commutator on the left-hand side is undefined. Secondly, and more importantly, (4.1) relies on the possibility of choosing global coordinates on $`^n`$, which precludes a naive generalization to arbitrary configuration spaces. And thirdly, even if one has managed to quantize $`p`$ and $`q`$ by finding a representation of (4.1), the problem of quantizing other observables remains - think of the Hamiltonian and the Schrödinger equation.
About 50 years ago, Mackey set himself the task of making good sense of canonical quantization; see Mackey (1968, 1978, 1992) and the brief exposition below for the result. Although the author now regards Mackey’s reformulation of quantization in terms of induced representations and systems of imprimitivity merely as a stepping stone towards our current understanding based on deformation theory and groupoids (cf. Subsection 4.3 below), Mackey’s approach is (quite rightly) often used in the foundations of physics, and one is well advised to be familiar with it. In any case, Mackey (1992, p. 283 - continuing the previous quotation) claims with some justification that his approach to quantization ‘removes much of the mystery.’
Like most approaches to quantization, Mackey assigns momentum and position a quite different role in quantum mechanics, despite the fact that in classical mechanics $`p`$ and $`q`$ can be interchanged by a canonical transformation:<sup>99</sup><sup>99</sup>99Up to a minus sign, that is. This is true globally on $`^n`$ and locally on any symplectic manifold, where local Darboux coordinates do not distinguish between position and momentum.
1. The position operators $`𝒬_{\mathrm{}}(q^j)`$ are collectively replaced by a single projection-valued measure $`P`$ on $`^n`$,<sup>100</sup><sup>100</sup>100 A projection-valued measure $`P`$ on a space $`\mathrm{\Omega }`$ with Borel structure (i.e. equipped with a $`\sigma `$-algebra of measurable sets defined by the topology) with values in a Hilbert space $``$ is a map $`EP(E)`$ from the Borel subsets $`E\mathrm{\Omega }`$ to the projections on $``$ that satisfies $`P(\mathrm{})=0`$, $`P(\mathrm{\Omega })=1`$, $`P(E)P(F)=P(F)P(E)=P(EF)`$ for all measurable $`E,F\mathrm{\Omega }`$, and $`P(_{i=1}^{\mathrm{}}E_i)=_{i=1}^{\mathrm{}}P(E_i)`$ for all countable collections of mutually disjoint $`E_i\mathrm{\Omega }`$. which on $`L^2(^n)`$ is given by $`P(E)=\chi _E`$ as a multiplication operator. Given this $`P`$, any multiplication operator defined by a (measurable) function $`f:^n`$ can be represented as $`_^n𝑑P(x)f(x)`$, which is defined and self-adjoint on a suitable domain.<sup>101</sup><sup>101</sup>101 This domain consists of all $`\mathrm{\Psi }`$ for which $`_^nd(\mathrm{\Psi },P(x)\mathrm{\Psi })|f(x)|^2<\mathrm{}`$. In particular, the position operators $`𝒬_{\mathrm{}}(q^j)`$ can be reconstructed from $`P`$ by choosing $`f(x)=x^j`$, i.e.
$$𝒬_{\mathrm{}}(q^j)=_^n𝑑P(x)x^j.$$
(4.2)
2. The momentum operators $`𝒬_{\mathrm{}}(p_j)`$ are collectively replaced by a single unitary group representation $`U(^n)`$, defined on $`L^2(^n)`$ by
$$U(y)\mathrm{\Psi }(x):=\mathrm{\Psi }(xy).$$
Each $`𝒬_{\mathrm{}}(p_j)`$ can be reconstructed from $`U`$ by means of
$$𝒬_{\mathrm{}}(p_j)\mathrm{\Psi }:=i\mathrm{}\underset{t_j0}{lim}t_j^1(U(t_j)1)\mathrm{\Psi },$$
(4.3)
where $`U(t_j)`$ is $`U`$ at $`x^j=t_j`$ and $`x^k=0`$ for $`kj`$.<sup>102</sup><sup>102</sup>102By Stone’s theorem (cf. Reed & Simon, 1972), this operator is defined and self-adjoint on the set of all $`\mathrm{\Psi }H`$ for which the limit exists.
Consequently, it entails no loss of generality to work with the pair $`(P,U)`$ instead of the pair $`(𝒬_{\mathrm{}}(q^k),𝒬_{\mathrm{}}(p_j))`$. The commutation relations (4.1) are now replaced by
$$U(x)P(E)U(x)^1=P(xE),$$
(4.4)
where $`E`$ is a (Borel) subset of $`^n`$ and $`xE=\{x\omega \omega E\}`$. On the basis of this reformulation, Mackey proposed the following sweeping generalization of the the canonical commutation relations:<sup>103</sup><sup>103</sup>103All groups and spaces are supposed to be locally compact, and actions and representations are assumed continuous.
> A system of imprimitivity $`(,U,P)`$ for a given action of a group $`G`$ on a space $`Q`$ consists of a Hilbert space $``$, a unitary representation $`U`$ of $`G`$ on $``$, and a projection-valued measure $`EP(E)`$ on $`Q`$ with values in $``$, such that (4.4) holds for all $`xG`$ and all Borel sets $`EQ`$.
In physics such a system describes the quantum mechanics of a particle moving on a configuration space $`Q`$ on which $`G`$ acts by symmetry transformations; see Subsection 4.3 for a more detailed discussion. When everything is smooth,<sup>104</sup><sup>104</sup>104I.e. $`G`$ is a Lie group, $`Q`$ is a manifold, and the $`G`$-action is smooth. each element $`X`$ of the Lie algebra $`𝔤`$ of $`G`$ defines a generalized momentum operator
$$𝒬_{\mathrm{}}(X)=i\mathrm{}dU(X)$$
(4.5)
on $``$.<sup>105</sup><sup>105</sup>105This operator is defined and self-adjoint on the domain of vectors $`\mathrm{\Psi }`$ for which $`dU(X)\mathrm{\Psi }:=lim_{t0}t^1(U(\mathrm{exp}(tX))1)\mathrm{\Psi }`$ exists. These operators satisfy the generalized canonical commutation relations<sup>106</sup><sup>106</sup>106As noted before in the context of (4.1), the commutation relations (4.6), (4.8) and (4.9) do not hold on the domain of self-adjointness of the operators involved, but on a smaller common core.
$$[𝒬_{\mathrm{}}(X),𝒬_{\mathrm{}}(Y)]=i\mathrm{}𝒬_{\mathrm{}}([X,Y]).$$
(4.6)
Furthermore, in terms of the operators<sup>107</sup><sup>107</sup>107For the domain of $`𝒬_{\mathrm{}}(f)`$ see footnote 101.
$$𝒬_{\mathrm{}}(f)=_Q𝑑P(x)f(x),$$
(4.7)
where $`f`$ is a smooth function on $`Q`$ and $`X𝔤`$, one in addition has
$$[𝒬_{\mathrm{}}(X),𝒬_{\mathrm{}}(f)]=i\mathrm{}𝒬_{\mathrm{}}(\xi _X^Qf),$$
(4.8)
where $`\xi _X^Q`$ is the canonical vector field on $`Q`$ defined by the $`G`$-action,<sup>108</sup><sup>108</sup>108I.e. $`\xi _X^Qf(y)=d/dt|_{t=0}[f(\mathrm{exp}(tX)y)]`$. and
$$[𝒬_{\mathrm{}}(f_1),𝒬_{\mathrm{}}(f_2)]=0.$$
(4.9)
Elementary quantum mechanics on $`^n`$ corresponds to the special case $`Q=^n`$ and $`G=^n`$ with the usual additive group structure. To see this, we denote the standard basis of $`^3`$ (in its guise as the Lie algebra of $`^3`$) by the name $`(p_j)`$, and furthermore take $`f_1(q)=q^j`$, $`f_2(q)=f(q)=q^k`$. Eq. (4.6) for $`X=p_j`$ and $`Y=p_k`$ then reads $`[𝒬_{\mathrm{}}(p_j),𝒬_{\mathrm{}}(p_k)]=0`$, eq. (4.8) yields the canonical commutation relations (4.1), and (4.9) states the commutativity of the position operators, i.e. $`[𝒬_{\mathrm{}}(q^j),𝒬_{\mathrm{}}(q^k)]=0`$.
In order to incorporate spin, one picks $`G=E(3)=SO(3)^3`$ (i.e. the Euclidean motion group), acting on $`Q=^3`$ in the obvious (defining) way. The Lie algebra of $`E(3)`$ is $`^6=^3\times ^3`$ as a vector space; we extend the basis $`(p_j)`$ of the second copy of $`^3`$ (i.e. the Lie algebra of $`^3`$) by a basis $`(J_i)`$ of the first copy of $`^3`$ (in its guise as the Lie algebra of $`SO(3)`$) , and find that the $`𝒬_{\mathrm{}}(J_i)`$ are just the usual angular momentum operators.<sup>109</sup><sup>109</sup>109The commutation relations in the previous paragraph are now extended by the familiar relations $`[𝒬_{\mathrm{}}(J_i),𝒬_{\mathrm{}}(J_j)]=i\mathrm{}ϵ_{ijk}𝒬_{\mathrm{}}(J_k)`$, $`[𝒬_{\mathrm{}}(J_i),𝒬_{\mathrm{}}(p_j)]=i\mathrm{}ϵ_{ijk}𝒬_{\mathrm{}}(p_k)`$, and $`[𝒬_{\mathrm{}}(J_i),𝒬_{\mathrm{}}(q^j)]=i\mathrm{}ϵ_{ijk}𝒬_{\mathrm{}}(q^k)`$.
Mackey’s generalization of von Neumann’s (1931) uniqueness theorem for the irreducible representations of the canonical commutation relations (4.1) is his imprimitivity theorem. This theorem applies to the special case where $`Q=G/H`$ for some (closed) subgroup $`HG`$, and states that (up to unitary equivalence) there is a bijective correspondence between:
1. Systems of imprimitivity $`(,U,P)`$ for the left-translation of $`G`$ on $`G/H`$;
2. Unitary representations $`U_\chi `$ of $`H`$.
This correspondence preserves irreducibility.<sup>110</sup><sup>110</sup>110Specifically, given $`U_\chi `$ the triple $`(^\chi ,U^\chi ,P^\chi )`$ is a system of imprimitivity, where $`^\chi =L^2(G/H)_\chi `$ carries the representation $`U^\chi (G)`$ induced by $`U_\chi (H)`$, and the $`P^\chi `$ act like multiplication operators. Conversely, if $`(,U,P)`$ is a system of imprimitivity, then there exists a unitary representation $`U_\chi (H)`$ such that the triple $`(,U,P)`$ is unitarily equivalent to the triple $`(^\chi ,U^\chi ,P^\chi )`$ just described. For example, for $`G=E(3)`$ and $`H=SO(3)`$ one has $`\chi =j=0,1,2,\mathrm{}`$ and $`^j=L^2(^3)_j`$ (where $`_j=^{2j+1}`$ carries the given representation $`U_j(SO(3))`$).
For example, von Neumann’s theorem is recovered as a special case of Mackey’s by making the choice $`G=^3`$ and $`H=\{e\}`$ (so that $`Q=^3`$, as above): the uniqueness of the (regular) irreducible representation of the canonical commutation relations here follows from the uniqueness of the irreducible representation of the trivial group. A more illustrative example is $`G=E(3)`$ and $`H=SO(3)`$ (so that $`Q=^3`$), in which case the irreducible representations of the associated system of imprimitivity are classified by spin $`j=0,1,\mathrm{}`$.<sup>111</sup><sup>111</sup>111By the usual arguments (Wigner’s theorem), one may replace $`SO(3)`$ by $`SU(2)`$, so as to obtain $`j=0,1/2,\mathrm{}`$. Mackey saw this as an explanation for the emergence of spin as a purely quantum-mechanical degree of freedom. Although the opinion that spin has no classical analogue was widely shared also among the pioneers of quantum theory,<sup>112</sup><sup>112</sup>112This opinion goes back to Pauli (1925), who talked about a ‘klassisch nicht beschreibbare Zweideutigkeit in den quantentheoretischen Eigenschaften des Elektrons,’ (i.e. an ‘ambivalence in the quantum theoretical properties of the electron that has no classical description’) which was later identified as spin by Goudsmit and Uhlenbeck. Probably the first person to draw attention to the classical counterpart of spin was Souriau (1969). Another misunderstanding about spin is that its ultimate explanation must be found in relativistic quantum mechanics. it is now obsolete (see Subsection 4.3 below). Despite this unfortunate misinterpretation, Mackey’s approach to canonical quantization is hard to surpass in power and clarity, and has many interesting applications.<sup>113</sup><sup>113</sup>113This begs the question about the ‘best’ possible proof of Mackey’s imprimitivity theorem. Mackey’s own proof was rather measure-theoretic in flavour, and did not shed much light on the origin of his result. Probably the shortest proof has been given by Ørsted (1979), but the insight brevity gives is still rather limited. Quite to the contrary, truly transparent proofs reduce a mathematical claim to a tautology. Such proofs, however, tend to require a formidable machinery to make this reduction work; see Echterhoff et al. (2002) and Landsman (2005b) for two different approaches to the imprimitivity theorem in this style.
We mention one of specific interest to the philosophy of physics, namely the Newton–Wigner position operator (as analyzed by Wightman, 1962).<sup>114</sup><sup>114</sup>114Fleming & Butterfield (2000) give an up-to-date introduction to particle localization in relativistic quantum theory. See also De Bièvre (2003). Here the general question is whether a given unitary representation $`U`$ of $`G=E(3)`$ on some Hilbert space $``$ may be extended to a system of imprimitivity with respect to $`H=SO(3)`$ (and hence $`Q=^3`$, as above); in that case, $`U`$ (or rather the associated quantum system) is said to be localizable in $`^3`$. Following Wigner’s (1939) suggestion that a relativistic particle is described by an irreducible representation $`U`$ of the Poincaré group $`P`$, one obtains a representation $`U(E(3))`$ by restricting $`U(P)`$ to the subgroup $`E(3)P`$.<sup>115</sup><sup>115</sup>115Strictly speaking, this hinges on the choice of an inertial frame in Minkowski space, with associated adapted co-ordinates such that the configuration space $`^3`$ in question is given by $`x^0=0`$. It then follows from the previous analysis that the particle described by $`U(P)`$ is localizable if and only if $`U(E(3))`$ is induced by some representation of $`SO(3)`$. This can, of course, be settled, with the result that massive particles of arbitrary spin can be localized in $`^3`$ (the corresponding position operator being precisely the one of Newton and Wigner), whereas massless particles may be localized in $`^3`$ if and only if their helicity is less than one. In particular, the photon (and the graviton) cannot be localized in $`^3`$ in the stated sense.<sup>116</sup><sup>116</sup>116 Seeing photons as quantized light waves with two possible polarizations transverse to the direction of propagation, this last result is physically perfectly reasonable.
To appreciate our later material on both phase space quantization and deformation quantization, it is helpful to give a $`C^{}`$-algebraic reformulation of Mackey’s approach. Firstly, by the spectral theorem (Reed & Simon, 1972; Pedersen, 1989), a projection-valued measure $`EP(E)`$ on a space $`Q`$ taking values in a Hilbert space $``$ is equivalent to a nondegenerate representation $`\pi `$ of the commutative $`C^{}`$-algebra $`C_0(Q)`$ on $``$ through the correspondence (4.7).<sup>117</sup><sup>117</sup>117A representation of a $`C^{}`$-algebra $`𝒜`$ on a Hilbert space $``$ is a linear map $`\pi :𝒜()`$ such that $`\pi (AB)=\pi (A)\pi (B)`$ and $`\pi (A^{})=\pi (A)^{}`$ for all $`A,B𝒜`$. Such a representation is called nondegenerate when $`\pi (A)\mathrm{\Psi }=0`$ for all $`A𝒜`$ implies $`\mathrm{\Psi }=0`$. Secondly, if $``$ in addition carries a unitary representation $`U`$ of $`G`$, the defining condition (4.4) of a system of imprimitivity (given a $`G`$-action on $`Q`$) is equivalent to the covariance condition
$$U(x)𝒬_{\mathrm{}}(f)U(x)^1=𝒬_{\mathrm{}}(L_xf)$$
(4.10)
for all $`xG`$ and $`fC_0(Q)`$, where $`L_xf(m)=f(x^1m)`$. Thus a system of imprimitivity for a given $`G`$-action on $`Q`$ is “the same” as a covariant nondegenerate representation of $`C_0(Q)`$. Thirdly, from a $`G`$-action on $`Q`$ one can construct a certain $`C^{}`$-algebra $`C^{}(G,Q)`$, the so-called transformation group $`C^{}`$-algebra defined by the action, which has the property that its nondegenerate representations correspond bijectively (and “naturally”) to covariant nondegenerate representations of $`C_0(Q)`$, and therefore to systems of imprimitivity for the given $`G`$-action (Effros & Hahn, 1967; Pedersen, 1979; Landsman, 1998). In the $`C^{}`$-algebraic approach to quantum physics, $`C^{}(G,Q)`$ is the algebra of observables of a particle moving on $`Q`$ subject to the symmetries defined by the $`G`$-action; its inequivalent irreducible representations correspond to the possible superselection sectors of the system (Doebner & Tolar, 1975; Majid, 1988, 1990; Landsman, 1990a, 1990b, 1992).<sup>118</sup><sup>118</sup>118Another reformulation of Mackey’s approach, or rather an extension of it, has been given by Isham (1984). In an attempt to reduce the whole theory to a problem in group representations, he proposed that the possible quantizations of a particle with configuration space $`G/H`$ are given by the inequivalent irreducible representations of a “canonical group” $`G_c=GV`$, where $`V`$ is the lowest-dimensional vector space that carries a representation of $`G`$ under which $`G/H`$ is an orbit in the dual vector space $`V^{}`$. All pertinent systems of imprimitivity then indeed correspond to unitary representations of $`G_c`$, but this group has many other representations whose physical interpretation is obscure. See also footnote 158.
### 4.2 Phase space quantization and coherent states
In Mackey’s approach to quantization, $`Q`$ is the configuration space of the system; the associated position coordinates commute (cf. (4.9)). This is reflected by the correspondence just discussed between projection-valued measures on $`Q`$ and representations of the commutative $`C^{}`$-algebra $`C_0(Q)`$. The noncommutativity of observables (and the associated uncertainty relations) typical of quantum mechanics is incorporated by adding the symmetry group $`G`$ to the picture and imposing the relations (4.4) (or, equivalently, (4.8) or (4.10)). As we have pointed out, this procedure upsets the symmetry between the phase space variables position and momentum in classical mechanics.
This somewhat unsatisfactory feature of Mackey’s approach may be avoided by replacing $`Q`$ by the phase space of the system, henceforth called $`M`$.<sup>119</sup><sup>119</sup>119Here the reader may think of the simplest case $`M=^6`$, the space of $`p`$’s and $`q`$’s of a particle moving on $`^3`$. More generally, if $`Q`$ is the configuration space, the associated phase space is the cotangent bundle $`M=T^{}Q`$. Even more general phase spaces, namely arbitrary symplectic manifolds, may be included in the theory as well. References for what follows include Busch, Grabowski, & Lahti, 1998, Schroeck, 1996, and Landsman, 1998, 1999a. In this approach, noncommutativity is incorporated by a treacherously tiny modification to Mackey’s setup. Namely, the projection-valued measure $`EP(M)`$ on $`M`$ with which he starts is now replaced by a positive-operator-valued measure or POVM on $`M`$, still taking values in some Hilbert space $`𝒦`$. This is a map $`EA(E)`$ from the (Borel) subsets $`E`$ of $`M`$ to the collection of positive bounded operators on $`𝒦`$,<sup>120</sup><sup>120</sup>120A bounded operator $`A`$ on $`𝒦`$ is called positive when $`(\mathrm{\Psi },A\mathrm{\Psi })0`$ for all $`\mathrm{\Psi }𝒦`$. Consequently, it is self-adjoint with spectrum contained in $`^+`$. satisfying $`A(\mathrm{})=0`$, $`A(M)=1`$, and $`A(_iE_i)=_iA(E_i)`$ for any countable collection of disjoint Borel sets $`E_i`$.<sup>121</sup><sup>121</sup>121Here the infinite sum is taken in the weak operator topology. Note that the above conditions force $`0A(E)1`$, in the sense that $`0(\mathrm{\Psi },A(E)\mathrm{\Psi })(\mathrm{\Psi },\mathrm{\Psi })`$ for all $`\mathrm{\Psi }𝒦`$. A POVM that satisfies $`A(EF)=A(E)A(F)`$ for all (Borel) $`E,FM`$ is precisely a projection-valued measure, so that a POVM is a generalization of the latter.<sup>122</sup><sup>122</sup>122This has given rise to the so-called operational approach to quantum theory, in which observables are not represented by self-adjoint operators (or, equivalently, by their associated projection-valued measures), but by POVM’s. The space $`M`$ on which the POVM is defined is the space of outcomes of the measuring instrument; the POVM is determined by both $`A`$ and a calibration procedure for this instrument. The probability that in a state $`\rho `$ the outcome of the experiment lies in $`EM`$ is taken to be $`\text{Tr}(\rho A(E))`$. See Davies (1976), Holevo (1982), Ludwig (1985), Schroeck (1996), Busch, Grabowski, & Lahti (1998), and De Muynck (2002). The point, then, is that a given POVM defines a quantization procedure by the stipulation that a classical observable $`f`$ (i.e. a measurable function on the phase space $`M`$, for simplicity assumed bounded) is quantized by the operator<sup>123</sup><sup>123</sup>123The easiest way to define the right-hand side of (4.11) is to fix $`\mathrm{\Psi }𝒦`$ and define a probability measure $`p_\mathrm{\Psi }`$ on $`M`$ by means of $`p_\mathrm{\Psi }(E)=(\mathrm{\Psi },A(E)\mathrm{\Psi })`$. One then defines $`𝒬(f)`$ as an operator through its expectation values $`(\mathrm{\Psi },𝒬(f)\mathrm{\Psi })=_M𝑑p_\mathrm{\Psi }(x)f(x)`$. The expression (4.11) generalizes (4.7), and also generalizes the spectral resolution of the operator $`f(A)=_{}𝑑P(\lambda )f(\lambda )`$, where $`P`$ is the projection-valued measure defined by a self-adjoint operator $`A`$.
$$𝒬(f)=_M𝑑A(x)f(x).$$
(4.11)
Thus the seemingly slight move from projection-valued measures on configuration space to positive-operator valued measures on phase space gives a wholly new perspective on quantization, actually reducing this task to the problem of finding such POVM’s.<sup>124</sup><sup>124</sup>124An important feature of $`𝒬`$ is that it is positive in the sense that if $`f(x)0`$ for all $`xM`$, then $`(\mathrm{\Psi },𝒬(f)\mathrm{\Psi })0`$ for all $`\mathrm{\Psi }𝒦`$. In other words, $`𝒬`$ is positive as a map from the $`C^{}`$-algebra $`C_0(M)`$ to the $`C^{}`$-algebra $`()`$.
The solution to this problem is greatly facilitated by Naimark’s dilation theorem.<sup>125</sup><sup>125</sup>125See, for example, Riesz and Sz.-Nagy (1990). It is better, however, to see Naimark’s theorem as a special case of Stinesprings’s, as explained e.g. in Landsman, 1998, and below. This states that, given a POVM $`EA(E)`$ on $`M`$ in a Hilbert space $`𝒦`$, there exists a Hilbert space $``$ carrying a projection-valued measure $`P`$ on $`M`$ and an isometric injection $`𝒦`$, such that
$$A(E)=[𝒦]P(E)[𝒦]$$
(4.12)
for all $`EM`$ (where $`[𝒦]`$ is the orthogonal projection from $``$ onto $`𝒦`$). Combining this with Mackey’s imprimitivity theorem yields a powerful generalization of the latter (Poulsen, 1970; Neumann, 1972; Scutaru, 1977; Cattaneo, 1979; Castrigiano & Henrichs, 1980).
First, define a generalized system of imprimitivity $`(𝒦,U,A)`$ for a given action of a group $`G`$ on a space $`M`$ as a POVM $`A`$ on $`M`$ taking values in a Hilbert space $`𝒦`$, along with a unitary representation $`V`$ of $`G`$ on $`𝒦`$ such that
$$V(x)A(E)V(x)^1=A(xE)$$
(4.13)
for all $`xG`$ and $`EM`$; cf. (4.4). Now assume $`M=G/H`$ (and the associated canonical left-action on $`M`$). The generalized imprimitivity theorem states that a generalized system of imprimitivity $`(𝒦,V,A)`$ for this action is necessarily (unitarily equivalent to) a reduction of a system of imprimitivity $`(,U,P)`$ for the same action. In other words, the Hilbert space $``$ in Naimark’s theorem carries a unitary representation $`U(G)`$ that commutes with the projection $`[𝒦]`$, and the representation $`V(G)`$ is simply the restriction of $`U`$ to $`𝒦`$. Furthermore, the POVM $`A`$ has the form (4.12). The structure of $`(,U,P)`$ is fully described by Mackey’s imprimitivity theorem, so that one has a complete classification of generalized systems of imprimitivity.<sup>126</sup><sup>126</sup>126Continuing footnote 110: $`V(G)`$ is necessarily a subrepresentation of some representation $`U^\chi (G)`$ induced by $`U_\chi (H)`$. One has
$$𝒦=p;=L^2(M)_\chi ,$$
(4.14)
where $`L^2`$ is defined with respect to a suitable measure on $`M=G/H`$,<sup>127</sup><sup>127</sup>127In the physically relevant case that $`G/H`$ is symplectic (so that it typically is a coadjoint orbit for $`G`$) one should take a multiple of the Liouville measure. the Hilbert space $`_\chi `$ carries a unitary representation of $`H`$, and $`p`$ is a projection in the commutant of the representation $`U^\chi (G)`$ induced by $`U_\chi (G)`$.<sup>128</sup><sup>128</sup>128The explicit form of $`U^\chi (g)`$, $`gG`$, depends on the choice of a cross-section $`\sigma :G/HG`$ of the projection $`\pi :GG/H`$ (i.e. $`\pi \sigma =\text{id}`$). If the measure on $`G/H`$ defining $`L^2(G/H)`$ is $`G`$-invariant, the explicit formula is $`U^\chi (g)\mathrm{\Psi }(x)=U_\chi (s(x)^1gs(g^1x))\mathrm{\Psi }(g^1x)`$. The quantization (4.11) is given by
$$𝒬(f)=pfp,$$
(4.15)
where $`f`$ acts on $`L^2(M)_\chi `$ as a multiplication operator, i.e. $`(f\mathrm{\Psi })(x)=f(x)\mathrm{\Psi }(x)`$. In particular, one has $`P(E)=\chi _E`$ (as a multiplication operator) for a region $`EM`$ of phase space, so that $`𝒬(\chi _E)=A(E)`$. Consequently, the probability that in a state $`\rho `$ (i.e. a density matrix on $`𝒦`$) the system is localized in $`E`$ is given by $`\text{Tr}(\rho A(E))`$.
In a more natural way than in Mackey’s approach, the covariant POVM quantization method allows one to incorporate space-time symmetries ab initio by taking $`G`$ to be the Galilei group or the Poincaré group, and choosing $`H`$ such that $`G/H`$ is a physical phase space (on which $`G`$, then, canonically acts). See Ali et al. (1995) and Schroeck (1996).
Another powerful method of constructing POVM’s on phase space (which in the presence of symmetries overlaps with the preceding one)<sup>129</sup><sup>129</sup>129 Suppose there is a vector $`\mathrm{\Omega }𝒦`$ such that $`_{G/H}𝑑\mu (x)|(\mathrm{\Omega },V(\sigma (x))\mathrm{\Omega })|^2<\mathrm{}`$ with respect to some cross-section $`\sigma :G/HG`$ and a $`G`$-invariant measure $`\mu `$, as well as $`V(h)\mathrm{\Omega }=U_\chi (h)\mathrm{\Omega }`$ for all $`hH`$, where $`U_\chi :H`$ is one-dimensional. Then (taking $`\mathrm{}=1`$) the vectors $`V(\sigma (x))\mathrm{\Omega }`$ (suitably normalized) form a family of coherent states on $`G/H`$ (Ali et al., 1995; Schroeck, 1996; Ali, Antoine, & Gazeau, 2000). For example, the coherent states (4.20) are of this form for the Heisenberg group. is based on coherent states.<sup>130</sup><sup>130</sup>130See Klauder & Skagerstam, 1985, Perelomov, 1986, Odzijewicz, 1992, Paul & Uribe, 1995, 1996, Ali et al., 1995, and Ali, Antoine, & Gazeau, 2000, for general discussions of coherent states. The minimal definition of coherent states in a Hilbert space $``$ for a phase space $`M`$ is that (for some fixed value of Planck’s constant $`\mathrm{}`$, for the moment) one has an injection<sup>131</sup><sup>131</sup>131This injection must be continuous as a map from $`M`$ to $``$, the projective Hilbert space of $``$. $`M`$, $`z\mathrm{\Psi }_z^{\mathrm{}}`$, such that
$$\mathrm{\Psi }_z^{\mathrm{}}=1$$
(4.16)
for all $`zM`$, and
$$c_{\mathrm{}}_M𝑑\mu _L(z)|(\mathrm{\Psi }_z^{\mathrm{}},\mathrm{\Phi })|^2=1,$$
(4.17)
for each $`\mathrm{\Phi }`$ of unit norm (here $`\mu _L`$ is the Liouville measure on $`M`$ and $`c_{\mathrm{}}>0`$ is a suitable constant).<sup>132</sup><sup>132</sup>132Other measures might occur here; see, for example, Bonechi & De Bièvre (2000). Condition (4.17) guarantees that we may define a POVM on $`M`$ in $`𝒦`$ by<sup>133</sup><sup>133</sup>133Recall that $`[\mathrm{\Psi }]`$ is the orthogonal projection onto a unit vector $`\mathrm{\Psi }`$.
$$A(E)=c_{\mathrm{}}_E𝑑\mu _L(z)[\mathrm{\Psi }_z^{\mathrm{}}].$$
(4.18)
Eq. (4.11) then simply reads (inserting the $`\mathrm{}`$-dependence of $`𝒬`$ and a suffix $`B`$ for later use)
$$𝒬_{\mathrm{}}^B(f)=c_{\mathrm{}}_M𝑑\mu _L(z)f(z)[\mathrm{\Psi }_z^{\mathrm{}}].$$
(4.19)
The time-honoured example, due to Schrödinger (1926b), is $`M=^{2n}`$, $`=L^2(^n)`$, and
$$\mathrm{\Psi }_{(p,q)}^{\mathrm{}}(x)=(\pi \mathrm{})^{n/4}e^{ipq/2\mathrm{}}e^{ipx/\mathrm{}}e^{(xq)^2/2\mathrm{}}.$$
(4.20)
Eq. (4.17) then holds with $`d\mu _L(p,q)=(2\pi )^nd^npd^nq`$ and $`c_{\mathrm{}}=\mathrm{}^n`$. One may verify that $`𝒬_{\mathrm{}}^B(p_j)`$ and $`𝒬_{\mathrm{}}^B(q^j)`$ coincide with Schrödinger’s operators (2.2). This example illustrates that coherent states need not be mutually orthogonal; in fact, in terms of $`z=p+iq`$ one has for the states in (4.20)
$$|(\mathrm{\Psi }_z^{\mathrm{}},\mathrm{\Psi }_w^{\mathrm{}})|^2=e^{|zw|^2/2\mathrm{}};$$
(4.21)
the significance of this result will emerge later on.
In the general case, it is an easy matter to verify Naimark’s dilation theorem for the POVM (4.18): changing notation so that the vectors $`\mathrm{\Psi }_z^{\mathrm{}}`$ now lie in $`𝒦`$, one finds
$$=L^2(M,c_{\mathrm{}}\mu _L),$$
(4.22)
the embedding $`W:𝒦`$ being given by $`(W\mathrm{\Phi })(z)=(\mathrm{\Psi }_z^{\mathrm{}},\mathrm{\Phi })`$. The projection-valued measure $`P`$ on $``$ is just $`P(E)=\chi _E`$ (as a multiplication operator), and the projection $`p`$ onto $`W𝒦`$ is given by
$$p\mathrm{\Psi }(z)=c_{\mathrm{}}_M𝑑\mu _L(w)(\mathrm{\Psi }_z^{\mathrm{}},\mathrm{\Psi }_w^{\mathrm{}})\mathrm{\Psi }(w).$$
(4.23)
Consequently, (4.19) is unitarily equivalent to (4.15), where $`f`$ acts on $`L^2(M)`$ as a multiplication operator.<sup>134</sup><sup>134</sup>134This leads to a close relationship between coherent states and Hilbert spaces with a reproducing kernel; see Landsman (1998) or Ali, Antoine, & Gazeau (2000).
> Thus (4.15) and (4.22) (or its possible extension (4.14)) form the essence of phase space quantization.<sup>135</sup><sup>135</sup>135See also footnote 173 below.
We close this subsection in the same fashion as the previous one, namely by pointing out the $`C^{}`$-algebraic significance of POVM’s. This is extremely easy: whereas a projection-valued measure on $`M`$ in $``$ is the same as a nondegenerate representation of $`C_0(M)`$ in $``$, a POVM on $`M`$ in a Hilbert space $`𝒦`$ is nothing but a nondegenerate completely positive map $`\phi :C_0(M)(𝒦)`$.<sup>136</sup><sup>136</sup>136A map $`\phi :𝒜`$ between $`C^{}`$-algebras is called positive when $`\phi (A)0`$ whenever $`A0`$; such a map is called completely positive if for all $`n`$ the map $`\phi _n:𝒜M_n()M_n()`$, defined by linear extension of $`\phi \text{id}`$ on elementary tensors, is positive (here $`M_n()`$ is the $`C^{}`$-algebra of $`n\times n`$ complex matrices). When $`𝒜`$ is commutative a nondegenerate positive map $`𝒜`$ is automatically completely positive for any $``$. Consequently, Naimark’s dilation theorem becomes a special case of Stinespring’s (1955) theorem: if $`𝒬:𝒜(𝒦)`$ is a completely positive map, there exists a Hilbert space $``$ carrying a representation $`\pi `$ of $`C_0(M)`$ and an isometric injection $`𝒦`$, such that $`𝒬(f)=[𝒦]\pi (f)[𝒦]`$ for all $`fC_0(M)`$. In terms of $`𝒬(C_0(M))`$, the covariance condition (4.13) becomes $`U(x)𝒬(f)U(x)^1=𝒬(L_xf)`$, just like (4.10).
### 4.3 Deformation quantization
So far, we have used the word ‘quantization’ in a heuristic way, basing our account on historical continuity rather than on axiomatic foundations. In this subsection and the next we set the record straight by introducing two alternative ways of looking at quantization in an axiomatic way. We start with the approach that historically came last, but which conceptually is closer to the material just discussed. This is deformation quantization, originating in the work of Berezin (1974, 1975a, 1975b), Vey (1975), and Bayen et al. (1977). We here follow the $`C^{}`$-algebraic approach to deformation quantization proposed by Rieffel (1989a, 1994), since it is not only mathematically transparent and rigorous, but also reasonable close to physical practice.<sup>137</sup><sup>137</sup>137See also Landsman (1998) for an extensive discussion of the $`C^{}`$-algebraic approach to deformation quantization. In other approaches to deformation quantization, such as the theory of star products, $`\mathrm{}`$ is a formal parameter rather than a real number. In particular, the meaning of the limit $`\mathrm{}0`$ is obscure. Due to the mathematical language used, this method of course naturally fits into the general $`C^{}`$-algebraic approach to quantum physics.
The key idea of deformation quantization is that quantization should be defined through the property of having the correct classical limit. Consequently, Planck’s “constant” $`\mathrm{}`$ is treated as a variable, so that for each of its values one should have a quantum theory. The key requirement is that this family of quantum theories converges to the underlying classical theory as $`\mathrm{}0`$.<sup>138</sup><sup>138</sup>138Cf. the preamble to Section 5 for further comments on this limit. The mathematical implementation of this idea is quite beautiful, in that the classical algebra of observables is “glued” to the family of quantum algebras of observables in such a way that the classical theory literally forms the boundary of the space containing the pertinent quantum theories (one for each value of $`\mathrm{}>0`$). Technically, this is done through the concept of a continuous field of $`C^{}`$-algebras.<sup>139</sup><sup>139</sup>139See Dixmier (1977), Fell & Doran (1988), and Kirchberg & Wassermann (1995) for three different approaches to the same concept. Our definition follows the latter; replacing $`I`$ by an arbitrary locally compact Hausdorff space one finds the general definition. What follows may sound unnecessarily technical, but the last 15 years have indicated that this yields exactly the right definition of quantization.
Let $`I`$ be the set in which $`\mathrm{}`$ takes values; one usually has $`I=[0,1]`$, but when the phase space is compact, $`\mathrm{}`$ often takes values in a countable subset of $`(0,1]`$.<sup>140</sup><sup>140</sup>140Cf. Landsman (1998) and footnote 205, but also see Rieffel (1989a) for the example of the noncommutative torus, where one quantizes a compact phase space for each $`\mathrm{}(0,1]`$. The same situation occurs in the theory of infinite systems; see Section 6. In any case, $`I`$ should contain zero as an accumulation point. A continuous field of $`C^{}`$-algebras over $`I`$, then, consists of a $`C^{}`$-algebra $`𝒜`$, a collection of $`C^{}`$-algebras $`\{𝒜_{\mathrm{}}\}_\mathrm{}I`$, and a surjective morphism $`\phi _{\mathrm{}}:𝒜𝒜_{\mathrm{}}`$ for each $`\mathrm{}I`$ , such that:
1. The function $`\mathrm{}\phi _{\mathrm{}}(A)_{\mathrm{}}`$ is in $`C_0(I)`$ for all $`A𝒜`$;<sup>141</sup><sup>141</sup>141Here $`_{\mathrm{}}`$ is the norm in the $`C^{}`$-algebra $`𝒜_{\mathrm{}}`$.
2. The norm of any $`A𝒜`$ is $`A=sup_\mathrm{}I\phi _{\mathrm{}}(A)`$;
3. For any $`fC_0(I)`$ and $`A𝒜`$ there is an element $`fA𝒜`$ for which $`\phi _{\mathrm{}}(fA)=f(\mathrm{})\phi _{\mathrm{}}(A)`$ for all $`\mathrm{}I`$.
The idea is that the family $`(𝒜_{\mathrm{}})_\mathrm{}I`$ of $`C^{}`$-algebras is glued together by specifying a topology on the bundle $`_{\mathrm{}[0,1]}𝒜_{\mathrm{}}`$ (disjoint union). However, this topology is in fact defined rather indirectly, via the specification of the space of continuous sections of the bundle.<sup>142</sup><sup>142</sup>142This is reminiscent of the Gelfand–Naimark theorem for commutative $`C^{}`$-algebras, which specifies the topology on a locally compact Hausdorff space $`X`$ via the $`C^{}`$-algebra $`C_0(X)`$. Similarly, in the theory of (locally trivial) vector bundles the Serre–Swan theorem allows one to reconstruct the topology on a vector bundle $`E\stackrel{\pi }{}X`$ from the space $`\mathrm{\Gamma }_0(E)`$ of continuous sections of $`E`$, seen as a (finitely generated projective) $`C_0(X)`$-module. See, for example, Gracia-Bondía, Várilly, & Figueroa (2001). The third condition in our definition of a continuous field of $`C^{}`$-algebras makes $`𝒜`$ a $`C_0(I)`$-module in the precise sense that there exits a nondegenerate morphism from $`C_0(I)`$ to the center of the multiplier of $`𝒜`$. This property may also replace our condition 3. Namely, a continuous section of the field is by definition an element $`\{A_{\mathrm{}}\}_\mathrm{}I`$ of $`_\mathrm{}I𝒜_{\mathrm{}}`$ (equivalently, a map $`\mathrm{}A_{\mathrm{}}`$ where $`A_{\mathrm{}}𝒜_{\mathrm{}}`$) for which there is an $`A𝒜`$ such that $`A_{\mathrm{}}=\phi _{\mathrm{}}(A)`$ for all $`\mathrm{}I`$. It follows that the $`C^{}`$-algebra $`𝒜`$ may actually be identified with the space of continuous sections of the field: if we do so, the morphism $`\phi _{\mathrm{}}`$ is just the evaluation map at $`\mathrm{}`$.<sup>143</sup><sup>143</sup>143The structure of $`𝒜`$ as a $`C^{}`$-algebra corresponds to the operations of pointwise scalar multiplication, addition, adjointing, and operator multiplication on sections.
Physically, $`𝒜_0`$ is the commutative algebra of observables of the underlying classical system, and for each $`\mathrm{}>0`$ the noncommutative $`C^{}`$-algebra $`𝒜_{\mathrm{}}`$ is supposed to be the algebra of observables of the corresponding quantum system at value $`\mathrm{}`$ of Planck’s constant. The algebra $`𝒜_0`$, then, is of the form $`C_0(M)`$, where $`M`$ is the phase space defining the classical theory. A phase space has more structure than an arbitrary topological space; it is a manifold on which a Poisson bracket $`\{,\}`$ can be defined. For example, on $`M=^{2n}`$ one has the familiar expression
$$\{f,g\}=\underset{j}{}\frac{f}{p_j}\frac{g}{q^j}\frac{f}{q^j}\frac{g}{p_j}.$$
(4.24)
Technically, $`M`$ is taken to be a Poisson manifold. This is a manifold equipped with a Lie bracket $`\{,\}`$ on $`C^{\mathrm{}}(M)`$ with the property that for each $`fC^{\mathrm{}}(M)`$ the map $`g\{f,g\}`$ defines a derivation of the commutative algebra structure of $`C^{\mathrm{}}(M)`$ given by pointwise multiplication. Hence this map is given by a vector field $`\xi _f`$, called the Hamiltonian vector field of $`f`$ (i.e. one has $`\xi _fg=\{f,g\}`$). Symplectic manifolds are special instances of Poisson manifolds, characterized by the property that the Hamiltonian vector fields exhaust the tangent bundle. A Poisson manifold is foliated by its symplectic leaves: a given symplectic leaf $`L`$ is characterized by the property that at each $`xL`$ the tangent space $`T_xLT_xM`$ is spanned by the collection of all Hamiltonian vector fields at $`x`$. Consequently, the flow of any Hamiltonian vector field on $`M`$ through a given point lies in its entirety within the symplectic leaf containing that point. The simplest example of a Poisson manifold is $`M=^{2n}`$ with Poisson bracket (4.24); this manifold is even symplectic.<sup>144</sup><sup>144</sup>144See Marsden & Ratiu (1994) for a mechanics-oriented introduction to Poisson manifolds; also cf. Landsman (1998) or Butterfield (2005) for the basic facts. A classical mathematical paper on Poisson manifolds is Weinstein (1983).
After this preparation, our basic definition is this:<sup>145</sup><sup>145</sup>145Here $`C_c^{\mathrm{}}(M)`$ stands for the space of smooth functions on $`M`$ with compact support; this is a norm-dense subalgebra of $`𝒜_0=C_0(M)`$. The question whether the maps $`𝒬_{\mathrm{}}`$ can be extended from $`C_c^{\mathrm{}}(M)`$ to $`C_0(M)`$ has to be answered on a case by case basis. Upon such an extension, if it exists, condition (4.25) will lose its meaning, since the Poisson bracket $`\{f,g\}`$ is not defined for all $`f,gC_0(M)`$.
> A deformation quantization of a phase space $`M`$ consists of a continuous field of $`C^{}`$-algebras $`(𝒜_{\mathrm{}})_{\mathrm{}[0,1]}`$ (with $`𝒜_0=C_0(M)`$), along with a family of self-adjoint<sup>146</sup><sup>146</sup>146I.e. $`𝒬_{\mathrm{}}(\overline{f})=𝒬_{\mathrm{}}(f)^{}`$. linear maps $`𝒬_{\mathrm{}}:C_c^{\mathrm{}}(M)𝒜_{\mathrm{}}`$, $`\mathrm{}(0,1]`$, such that:
1. For each $`fC_c^{\mathrm{}}(M)`$ the map defined by $`0f`$ and $`\mathrm{}𝒬_{\mathrm{}}(f)`$ ($`\mathrm{}0`$) is a continuous section of the given continuous field;<sup>147</sup><sup>147</sup>147Equivalently, one could extend the family $`(𝒬_{\mathrm{}})_{\mathrm{}(0,1]}`$ to $`\mathrm{}=0`$ by $`𝒬_0=\text{id}`$, and state that $`\mathrm{}𝒬_{\mathrm{}}(f)`$ is a continuous section. Also, one could replace this family of maps by a single map $`𝒬:C_c^{\mathrm{}}(M)𝒜`$ and define $`𝒬_{\mathrm{}}=\phi _{\mathrm{}}𝒬:C_c^{\mathrm{}}(M)𝒜_{\mathrm{}}`$.
2. For all $`f,gC_c^{\mathrm{}}(M)`$ one has
$$\underset{\mathrm{}0}{lim}\frac{i}{\mathrm{}}[𝒬_{\mathrm{}}(f),𝒬_{\mathrm{}}(g)]𝒬_{\mathrm{}}(\{f,g\})_{\mathrm{}}=0.$$
(4.25)
Obvious continuity properties one might like to impose, such as
$$\underset{\mathrm{}0}{lim}𝒬_{\mathrm{}}(f)𝒬_{\mathrm{}}(g)𝒬_{\mathrm{}}(fg)=0,$$
(4.26)
or
$$\underset{\mathrm{}0}{lim}𝒬_{\mathrm{}}(f)=f_{\mathrm{}},$$
(4.27)
turn out to be an automatic consequence of the definitions.<sup>148</sup><sup>148</sup>148That they are automatic should not distract from the fact that especially (4.27) is a beautiful connection between classical and quantum mechanics. See footnote 90 for the meaning of $`f_{\mathrm{}}`$. Condition (4.25), however, transcends the $`C^{}`$-algebraic setting, and is the key ingredient in proving (among other things) that the quantum dynamics converges to the classical dynamics;<sup>149</sup><sup>149</sup>149This insight is often attributed to Dirac (1930), who was the first to recognize the analogy between the commutator in quantum mechanics and the Poisson bracket in classical mechanics. see Section 5. The map $`𝒬_{\mathrm{}}`$ is the quantization map at value $`\mathrm{}`$ of Planck’s constant; we feel it is the most precise formulation of Heisenberg’s original Umdeutung of classical observables known to date. It has the same interpretation as the heuristic symbol $`𝒬_{\mathrm{}}`$ used so far: the operator $`𝒬_{\mathrm{}}(f)`$ is the quantum-mechanical observable whose classical counterpart is $`f`$.
This has turned out to be an fruitful definition of quantization, firstly because most well-understood examples of quantization fit into it (Rieffel, 1994; Landsman, 1998), and secondly because it has suggested various fascinating new ones (Rieffel, 1989a; Natsume& Nest, 1999; Natsume, Nest, & Ingo, 2003). Restricting ourselves to the former, we note, for example, that (4.19) with (4.20) defines a deformation quantization of the phase space $`^{2n}`$ (with standard Poisson bracket) if one takes $`𝒜_{\mathrm{}}`$ to be the $`C^{}`$-algebra of compact operators on the Hilbert space $`L^2(^n)`$. This is called the Berezin quantization of $`^{2n}`$ (as a phase space);<sup>150</sup><sup>150</sup>150 In the literature, Berezin quantization on $`^{2n}`$ is often called anti-Wick quantization (or ordering), whereas on compact complex manifolds it is sometimes called Toeplitz or Berezin–Toeplitz quantization. Coherent states based on other phase spaces often define deformation quantizations as well; see Landsman, 1998. explicitly, for $`\mathrm{\Phi }L^2(^n)`$ one has
$$𝒬_{\mathrm{}}^B(f)\mathrm{\Phi }(x)=_{^{2n}}\frac{d^npd^nqd^ny}{(2\pi \mathrm{})^n}f(p,q)\overline{\mathrm{\Psi }_{(p,q)}^{\mathrm{}}(y)}\mathrm{\Phi }(y)\mathrm{\Psi }_{(p,q)}^{\mathrm{}}(x).$$
(4.28)
This quantization has the distinguishing feature of positivity,<sup>151</sup><sup>151</sup>151Cf. footnote 124. As a consequence, (4.28) is valid not only for $`fC_c^{\mathrm{}}(^{2n})`$, but even for all $`fL^{\mathrm{}}(^{2n})`$, and the extension of $`𝒬_{\mathrm{}}^B`$ from $`C_c^{\mathrm{}}(^{2n})`$ to $`L^{\mathrm{}}(^{2n})`$ is continuous. a property not shared by its more famous sister called Weyl quantization.<sup>152</sup><sup>152</sup>152The original reference is Weyl (1931). See, for example, Dubin, Hennings, & Smith (2000) and Esposito, Marmo, & Sudarshan (2004) for a modern physics-oriented yet mathematically rigorous treatment. See also Rieffel (1994) and Landsman (1998) for a discussion from the perspective of deformation quantization. The latter is a deformation quantization of $`^{2n}`$ as well, having the same continuous field of $`C^{}`$-algebras, but differing from Berezin quantization in its quantization map
$$𝒬_{\mathrm{}}^W(f)\mathrm{\Phi }(x)=_{^{2n}}\frac{d^npd^nq}{(2\pi \mathrm{})^n}e^{ip(xq)/\mathrm{}}f(p,\frac{1}{2}(x+q))\mathrm{\Phi }(q).$$
(4.29)
Although it lacks good positivity and hence continuity properties,<sup>153</sup><sup>153</sup>153Nonetheless, Weyl quantization may be extended from $`C_c^{\mathrm{}}(^{2n})`$ to much larger function spaces using techniques from the theory of distributions (leaving the Hilbert space setting typical of quantum mechanics). The classical treatment is in Hörmander (1979, 1985a). Weyl quantization enjoys better symmetry properties than Berezin quantization.<sup>154</sup><sup>154</sup>154 Weyl quantization is covariant under the affine symplectic group $`\mathrm{Sp}(n,)^{2n}`$, whereas Berezin quantization is merely covariant under its subgroup $`\mathrm{O}(2n)^{2n}`$. Despite these differences, which illustrate the lack of uniqueness of concrete quantization procedures, Weyl and Berezin quantization both reproduce Schrödinger’s position and momentum operators (2.2).<sup>155</sup><sup>155</sup>155This requires a formal extension of the maps $`𝒬_{\mathrm{}}^W`$ and $`𝒬_{\mathrm{}}^B`$ to unbounded functions on $`M`$ like $`p_j`$ and $`q^j`$. Furthermore, if $`fL^1(^{2n})`$, then $`𝒬_{\mathrm{}}^B(f)`$ and $`𝒬_{\mathrm{}}^W(f)`$ are trace class, with
$$\text{Tr}𝒬_{\mathrm{}}^B(f)=\text{Tr}𝒬_{\mathrm{}}^W(f)=_{^{2n}}\frac{d^npd^nq}{(2\pi \mathrm{})^n}f(p,q).$$
(4.30)
Weyl and Berezin quantization are related by
$$𝒬_{\mathrm{}}^B(f)=𝒬_{\mathrm{}}^W(e^{\frac{\mathrm{}}{4}\mathrm{\Delta }_{2n}}f),$$
(4.31)
where $`\mathrm{\Delta }_{2n}=_{j=1}^n(^2/p_j^2+^2/(q^j)^2)`$, from which it may be shown that Weyl and Berezin quantization are asymptotically equal in the sense that for any $`fC_c^{\mathrm{}}(^{2n})`$ one has
$$\underset{\mathrm{}0}{lim}𝒬_{\mathrm{}}^B(f)𝒬_{\mathrm{}}^W(f)=0.$$
(4.32)
Mackey’s approach to quantization also finds its natural home in the setting of deformation quantization. Let a Lie group $`G`$ act on a manifold $`Q`$, interpreted as a configuration space, as in Subsection 4.1. It turns out that the corresponding classical phase space is the manifold $`𝔤^{}\times Q`$, equipped with the so-called semidirect product Poisson structure (Marsden, Raţiu & Weinstein, 1984; Krishnaprasad & Marsden, 1987). Relative to a basis $`(T_a)`$ of the Lie algebra $`𝔤`$ of $`G`$ with structure constants $`C_{ab}^c`$ (i.e. $`[T_a,T_b]=_cC_{ab}^cT_c`$), the Poisson bracket in question is given by
$$\{f,g\}=\underset{a}{}\left(\xi _a^Mf\frac{g}{\theta _a}\frac{f}{\theta _a}\xi _a^Mg\right)\underset{a,b,c}{}C_{ab}^c\theta _c\frac{f}{\theta _a}\frac{g}{\theta _b},$$
(4.33)
where $`\xi _a^M=\xi _{T_a}^M`$. To illustrate the meaning of this lengthy expression, we consider a few special cases. First, take $`f=X𝔤`$ and $`g=Y𝔤`$ (seen as linear functions on the dual $`𝔤^{}`$). This yields
$$\{X,Y\}=[X,Y].$$
(4.34)
Subsequently, assume that $`g`$ depends on position $`q`$ alone. This leads to
$$\{X,g\}=\xi _X^Mg.$$
(4.35)
Finally, assume that $`f=f_1`$ and $`g=f_2`$ depend on $`q`$ only; this clearly gives
$$\{f_1,f_2\}=0.$$
(4.36)
The two simplest physically relevant examples, already considered at the quantum level in Subsection 4.1, are as follows. First, take $`G=^n`$ (as a Lie group) and $`Q=^n`$ (as a manifold), with $`G`$ acting on $`Q`$ by translation. Eqs. (4.34) - (4.36) then yield the Poisson brackets $`\{p_j,p_k\}=0`$, $`\{p_j,q^k\}=\delta _j^k`$, and $`\{q^j,q^k\}=0`$, showing that in this case $`M=𝔤^{}\times Q=^{2n}`$ is the standard phase space of a particle moving in $`^n`$; cf. (4.24). Second, the case $`G=E(3)`$ and $`Q=^3`$ yields a phase space $`M=^3\times ^6`$, where $`^6`$ is the phase space of a spinless particle just considered, and $`^3`$ is an additional internal space containing spin as a classical degree of freedom. Indeed, beyond the Poisson brackets on $`^6`$ just described, (4.34) - (4.36) give rise to the additional Poisson brackets $`\{J_i,J_j\}=ϵ_{ijk}J_k`$, $`\{J_i,p_j\}=ϵ_{ijk}p_k`$, and $`\{J_i,q^j\}=ϵ_{ijk}q^k`$.<sup>156</sup><sup>156</sup>156These are the classical counterparts of the commutation relations for angular momentum written down in footnote 109.
The analogy between (4.34), (4.35), (4.36) on the one hand, and (4.6), (4.8), (4.9), respectively, on the other, is no accident: the Poisson brackets in question are the classical counterpart of the commutation relations just referred to. This observation is made precise by the fundamental theorem relating Mackey’s systems of imprimitivity to deformation quantization (Landsman, 1993, 1998): one can equip the family of $`C^{}`$-algebras
$`𝒜_0`$ $`=`$ $`C_0(𝔤^{}\times Q);`$
$`𝒜_{\mathrm{}}`$ $`=`$ $`C^{}(G,Q),`$ (4.37)
where $`C^{}(G,Q)`$ is the transformation grouo $`C^{}`$-algebra defined by the given $`G`$-action on $`Q`$ (cf. the end of Subsection 4.1), with the structure of a continuous field, and one can define quantization maps $`𝒬_{\mathrm{}}:C_c^{\mathrm{}}(𝔤^{}\times Q)C^{}(G,Q)`$ so as to obtain a deformation quantization of the phase space $`𝔤^{}\times Q`$. It turns out that for special functions of the type $`X,Y𝔤`$, and $`f=f(q)`$ just considered, the equality
$$\frac{i}{\mathrm{}}[𝒬_{\mathrm{}}(f),𝒬_{\mathrm{}}(g)]𝒬_{\mathrm{}}(\{f,g\})=0$$
(4.38)
holds exactly (and not merely asymptotically for $`\mathrm{}0`$, as required in the fundamental axiom (4.25) for deformation quantization).
This result clarifies the status of Mackey’s quantization by systems of imprimitivity. The classical theory underlying the relations (4.4) is not the usual phase space $`T^{}Q`$ of a structureless particle moving on $`Q`$, but $`M=𝔤^{}\times Q`$. For simplicity we restrict ourselves to the transitive case $`Q=G/H`$ (with canonical left $`G`$-action). Then $`M`$ coincides with $`T^{}Q`$ only when $`H=\{e\}`$ and hence $`Q=G`$;<sup>157</sup><sup>157</sup>157For a Lie group $`G`$ one has $`T^{}G𝔤^{}\times G`$. in general, the phase space $`𝔤^{}\times (G/H)`$ is locally of the form $`T^{}(G/H)\times 𝔥^{}`$ (where $`𝔥^{}`$ is the dual of the Lie algebra of $`H`$). The internal degree of freedom described by $`𝔥^{}`$ is a generalization of classical spin, which, as we have seen, emerges in the case $`G=E(3)`$ and $`H=SO(3)`$. All this is merely a special case of a vast class of deformation quantizations described by Lie groupoids; see Bellisard & Vittot (1990), Landsman (1998, 1999b, 2005b) and Landsman & Ramazan (2001).<sup>158</sup><sup>158</sup>158A similar analysis can be applied to Isham’s (1984) quantization scheme mentioned in footnote 118. The unitary irreducible representations of the canonical group $`G_c`$ stand in bijective correspondence with the nondegenerate representations of the group $`C^{}`$-algebra $`C^{}(G_c)`$ (Pedersen, 1979), which is a deformation quantization of the Poisson manifold $`𝔤_c^{}`$ (i.e. the dual of the Lie algebra of $`G_c`$). This Poisson manifold contains the coadjoint orbits of $`G_c`$ as “irreducible” classical phase spaces, of which only one is the cotangent bundle $`T^{}(G/H)`$ one initially thought one was quantizing (see Landsman (1998) for the classification of the coadjoint orbits of semidirect products). All other orbits are mere lumber that one should avoid. See also Robson (1996). If one is ready for groupoids, there is no need for the canonical group approach.
### 4.4 Geometric quantization
Because of its use of abstract $`C^{}`$-algebras, deformation quantization is a fairly sophisticated and recent technique. Historically, it was preceded by a more concrete and traditional approach called geometric quantization.<sup>159</sup><sup>159</sup>159 Geometric quantization was independently introduced by Kostant (1970) and Souriau (1969). Major later treatments on the basis of the original formalism are Guillemin & Sternberg (1977), Śniatycki (1980), Kirillov (1990), Woodhouse (1992), Puta (1993), Chernoff (1995), Kirillov (2004), and Ali & Englis (2004). The modern era (based on the use of Dirac operators and $`K`$-theory) was initiated by unpublished remarks by Bott in the early 1990s; see Vergne (1994) and Guillemin, Ginzburg & Karshon (2002). The postmodern (i.e. functorial) epoch was launched in Landsman (2005a). Here the goal is to firstly “quantize” a phase space $`M`$ by a concretely given Hilbert space $`(M)`$, and secondly to map the classical observables (i.e. the real-valued smooth functions on $`M`$) into self-adjoint operators on $``$ (which after all play the role of observables in von Neumann’s formalism of quantum mechanics).<sup>160</sup><sup>160</sup>160In geometric quantization phase spaces are always seen as symplectic manifolds (with the sole exception of Vaisman, 1991); the reason why it is unnatural to start with the more general class of Poisson manifolds will become clear in the next subsection. In principle, this program should align geometric quantization much better with the fundamental role unbounded self-adjoint operators play in quantum mechanics than deformation quantization, but in practice geometric quantization continues to be plagued by problems.<sup>161</sup><sup>161</sup>161 Apart from rather technical issues concerning the domains and self-adjointness properties of the operators defined by geometric quantization, the main point is that the various mathematical choices one has to make in the geometric quantization procedure cannot all be justified by physical arguments, although the physical properties of the theory depend on these choices. (The notion of a polarization is the principal case in point; see also footnote 174 below.) Furthermore, as we shall see, one cannot quantize sufficiently many functions in standard geometric quantization. Our functorial approach to geometric quantization in Subsection 4.5 was partly invented to alleviate these problems. However, it would be wrong to see deformation quantization and geometric quantization as competitors; as we shall see in the next subsection, they are natural allies, forming “complementary” parts of a conjectural quantization functor.
In fact, in our opinion geometric quantization is best compared and contrasted with phase space quantization in its concrete formulation of Subsection 4.2 (i.e. before its $`C^{}`$-algebraic abstraction and subsequent absorption into deformation quantization as indicated in Subsection 4.3).<sup>162</sup><sup>162</sup>162See also Tuynman (1987). For geometric quantization equally well starts with the Hilbert space $`L^2(M)`$,<sup>163</sup><sup>163</sup>163Defined with respect to the Liouville measure times a suitable factor $`c_{\mathrm{}}`$, as in (4.17) etc.; in geometric quantization this factor is not very important, as it is unusual to study the limit $`\mathrm{}0`$. For $`M=^{2n}`$ the measure on $`M`$ with respect to which $`L^2(M)`$ is defined is $`d^npd^nq/(2\pi \mathrm{})^n`$. and subsequently attempts to construct $`(M)`$ from it, though typically in a different way from (4.14).
Before doing so, however, the geometric quantization procedure first tries to define a linear map $`𝒬_{\mathrm{}}^{pre}`$ from $`C^{\mathrm{}}(M)`$ to the class of (generally unbounded) operators on $`L^2(M)`$ that formally satisfies
$$\frac{i}{\mathrm{}}[𝒬_{\mathrm{}}^{pre}(f),𝒬_{\mathrm{}}^{pre}(g)]𝒬_{\mathrm{}}^{pre}(\{f,g\})=0,$$
(4.39)
i.e. (4.38) with $`𝒬=𝒬_{\mathrm{}}^{pre}`$, as well as the nondegeneracy property
$$𝒬_{\mathrm{}}^{pre}(\chi _M)=1,$$
(4.40)
where $`\chi _M`$ is the function on $`M`$ that is identically equal to 1, and the 1 on the right-hand side is the unit operator on $`L^2(M)`$. Such a map is called prequantization.<sup>164</sup><sup>164</sup>164The idea of prequantization predates geometric quantization; see van Hove (1951) and Segal (1960). For $`M=^{2n}`$ (equipped with its standard Poisson bracket (4.24)), a prequantization map is given (on $`\mathrm{\Phi }L^2(M)`$) by
$$𝒬_{\mathrm{}}^{pre}(f)\mathrm{\Phi }=i\mathrm{}\{f,\mathrm{\Phi }\}+\left(f\underset{j}{}p_j\frac{f}{p_j}\right)\mathrm{\Phi }.$$
(4.41)
This expression is initially defined for $`\mathrm{\Phi }C_c^{\mathrm{}}(M)L^2(M)`$, on which domain $`𝒬_{\mathrm{}}^{pre}(f)`$ is symmetric when $`f`$ is real-valued;<sup>165</sup><sup>165</sup>165An operator $`A`$ defined on a dense subspace $`𝒟`$ of a Hilbert space $``$ is called symmetric when $`(A\mathrm{\Psi },\mathrm{\Phi })=(\mathrm{\Psi },A\mathrm{\Phi })`$ for all $`\mathrm{\Psi },\mathrm{\Phi }𝒟`$. note that the operator in question is unbounded even when $`f`$ is bounded.<sup>166</sup><sup>166</sup>166As mentioned, self-adjointness is a problem in geometric quantization; we will not address this issue here. Berezin quantization is much better behaved than geometric quantization in this respect, since it maps bounded functions into bounded operators. This looks complicated; the simpler expression $`𝒬_{\mathrm{}}(f)\mathrm{\Phi }=i\mathrm{}\{f,\mathrm{\Phi }\}`$, however, would satisfy (4.38) but not (4.40), and the goal of the second term in (4.41) is to satisfy the latter condition while preserving the former.<sup>167</sup><sup>167</sup>167One may criticize the geometric quantization procedure for emphasizing (4.39) against its equally natural counterpart $`𝒬(fg)=𝒬(f)𝒬(g)`$, which fails to be satisified by $`𝒬_{\mathrm{}}^{pre}`$ (and indeed by any known quantization procedure, except the silly $`𝒬(f)=f`$ (as a multiplication operator on $`L^2(M)`$). For example, one has
$`𝒬_{\mathrm{}}^{pre}(q^k)`$ $`=`$ $`q^k+i\mathrm{}{\displaystyle \frac{}{p_k}}`$
$`𝒬_{\mathrm{}}^{pre}(p_j)`$ $`=`$ $`i\mathrm{}{\displaystyle \frac{}{q^j}}.`$ (4.42)
For general phase spaces $`M`$ one may construct a map $`𝒬_{\mathrm{}}^{pre}`$ that satisfies (4.39) and (4.40) when $`M`$ is “prequantizable”; a full explanation of this notion requires some differential geometry.<sup>168</sup><sup>168</sup>168A symplectic manifold $`(M,\omega )`$ is called prequantizable at some fixed value of $`\mathrm{}`$ when it admits a complex line bundle $`LM`$ (called the prequantization line bundle) with connection $``$ such that $`F=i\omega /\mathrm{}`$ (where $`F`$ is the curvature of the connection, defined by $`F(X,Y)=[_X,_Y]_{[X,Y]}`$). This is the case iff $`[\omega ]/2\pi \mathrm{}H^2(M,)`$, where $`[\omega ]`$ is the de Rham cohomology class of the symplectic form. If so, prequantization is defined by the formula $`𝒬_{\mathrm{}}^{pre}(f)=i\mathrm{}_{\xi _f}+f`$, where $`\xi _f`$ is the Hamiltonian vector field of $`f`$ (see Subsection 4.3). This expression is defined and symmetric on the space $`C_c^{\mathrm{}}(M,L)L^2(M)`$ of compactly supported smooth sections of $`L`$, and is easily checked to satisfy (4.39) and (4.40). To obtain (4.41) as a special case, note that for $`M=^{2n}`$ with the canonical symplectic form $`\omega =_kdp_kdq^k`$ one has $`[\omega ]=0`$, so that $`L`$ is the trivial bundle $`L=^{2n}\times `$. The connection $`=d+A`$ with $`A=\frac{i}{\mathrm{}}_kp_kdq^k`$ satisfies $`F=i\omega /\mathrm{}`$, and this eventually yields (4.41). Assuming this to be the case, then for one thing prequantization is a very effective tool in constructing unitary group representations of the kind that are interesting for physics. Namely, suppose a Lie group $`G`$ acts on the phase space $`M`$ in “canonical” fashion. This means that there exists a map $`\mu :M𝔤^{}`$, called the momentum map, such that $`\xi _{\mu _X}=\xi _X^M`$ for each $`X𝔤`$,<sup>169</sup><sup>169</sup>169 Here $`\mu _XC^{\mathrm{}}(M)`$ is defined by $`\mu _X(x)=\mu (x),X`$, and $`\xi _X^M`$ is the vector field on $`M`$ defined by the $`G`$-action (cf. footnote 108). Hence this condition means that $`\{\mu _X,f\}(y)=d/dt_{|t=0}[f(\mathrm{exp}(tX)y)]`$ for all $`fC^{\mathrm{}}(M)`$ and all $`yM`$. and in addition $`\{\mu _X,\mu _Y\}=\mu _{[X,Y]}`$. See Abraham & Marsden (1985), Marsden & Ratiu (1994), Landsman (1998), Butterfield (2005), etc. On then obtains a representation $`\pi `$ of the Lie algebra $`𝔤`$ of $`G`$ by skew-symmetric unbounded operators on $`L^2(M)`$ through
$$\pi (X)=i\mathrm{}𝒬_{\mathrm{}}^{pre}(\mu _X),$$
(4.43)
which often exponentiates to a unitary representation of $`G`$.<sup>170</sup><sup>170</sup>170An operator $`A`$ defined on a dense subspace $`𝒟`$ of a Hilbert space $``$ is called skew-symmetric when $`(A\mathrm{\Psi },\mathrm{\Phi })=(\mathrm{\Psi },A\mathrm{\Phi })`$ for all $`\mathrm{\Psi },\mathrm{\Phi }𝒟`$. If one has a unitary representation $`U`$ of a Lie group $`G`$ on $``$, then the derived representation $`dU`$ of the Lie algebra $`𝔤`$ (see footnote 105) consists of skew-symmetric operators, making one hopeful that a given representation of $`𝔤`$ by skew-symmetric operators can be integrated (or exponentiated) to a unitary representation of $`G`$. See Barut & Raçka (1977) or Jørgensen & Moore (1984) and references therein.
As the name suggests, prequantization is not yet quantization. For example, the prequantization of $`M=^{2n}`$ does not reproduce Schrödinger’s wave mechanics: the operators (4.42) are not unitarily equivalent to (2.2). In fact, as a carrier of the representation (4.42) of the canonical commutation relations (4.1), the Hilbert space $`L^2(^{2n})`$ contains $`L^2(^n)`$ (carrying the representation (2.2)) with infinite multiplicity (Ali & Emch, 1986). This situation is often expressed by the statement that “prequantization is reducible” or that the prequantization Hilbert space $`L^2(M)`$ is ‘too large’, but both claims are misleading: $`L^2(M)`$ is actually irreducible under the action of $`𝒬_{\mathrm{}}^{pre}(C^{\mathrm{}}(M))`$ (Tuynman, 1998), and saying that for example $`L^2(^n)`$ is “larger” than $`L^2(^n)`$ is unmathematical in view of the unitary isomorphism of these Hilbert spaces. What is true is that in typical examples $`L^2(M)`$ is generically reducible under the action of some Lie algebra where one would like it to be irreducible. This applies, for example, to (2.2), which defines a representation of the Lie algebra of the Heisenberg group. More generally, in the case where a phase space $`M`$ carries a transitive action of a Lie group $`G`$, so that one would expect the quantization of this $`G`$-action by unitary operators on a Hilbert space to be irreducible, $`L^2(M)`$ is typically highly reducible under the representation (4.43) of $`𝔤`$.<sup>171</sup><sup>171</sup>171This can be made precise in the context of the so-called orbit method, cf. the books cited in footnote 159.
Phase space quantization encounters this problem as well. Instead of the complicated expression (4.41), through (4.11) it simply “phase space prequantizes” $`fC^{\mathrm{}}(M)`$ on $`L^2(M)`$ by $`f`$ as a multiplication operator.<sup>172</sup><sup>172</sup>172For unbounded $`f`$ this operator is defined on the set of all $`\mathrm{\Phi }L^2(M)`$ for which $`f\mathrm{\Phi }L^2(M)`$. Under this action of $`C^{\mathrm{}}(M)`$ the Hilbert space $`L^2(M)`$ is of course highly reducible.<sup>173</sup><sup>173</sup>173 Namely, each (measurable) subset $`EM`$ defines a projection $`\chi _E`$, and $`\chi _EL^2(M)`$ is stable under all multiplication operators $`f`$. One could actually decide not to be bothered by this problem and stop here, but then one is simply doing classical mechanics in a Hilbert space setting (Koopman, 1931). This formalism even turns out to be quite useful for ergodic theory (Reed & Simon, 1972). The identification of an appropriate subspace
$$(M)=pL^2(M)$$
(4.44)
of $`L^2(M)`$ (where $`p`$ is a projection) as the Hilbert space carrying the “quantization” of $`M`$ (or rather of $`C^{\mathrm{}}(M)`$) may be seen as a solution to this reducibility problem, for if the procedure is successful, the projection $`p`$ is chosen such that $`pL^2(M)`$ is irreducible under $`pC^{\mathrm{}}(M)p`$. Moreover, in this way practically any function on $`M`$ can be quantized, albeit at the expense of (4.38) (which, as we have seen, gets replaced by its asymptotic version (4.25)). See Subsection 6.3 for a discussion of reducibility versus irreducibility of representations of algebras of observables in classical and quantum theory.
We restrict our treatment of geometric quantization to situations where it adopts the same strategy as above, in assuming that the final Hilbert space has the form (4.44) as well.<sup>174</sup><sup>174</sup>174 Geometric quantization has traditionally been based on the notion of a polarization (cf. the references in footnote 159). This device produces a final Hilbert space $`(M)`$ which may not be a subspace of $`L^2(M)`$, except in the so-called (anti-) holomorphic case. But it crucially differs from phase space quantization in that its first step is (4.41) (or its generalization to more general phase spaces) rather than just having $`f\mathrm{\Phi }`$ on the right-hand side.<sup>175</sup><sup>175</sup>175It also differs from phase space quantization in the ideology that the projection $`p`$ ought to be constructed solely from the geometry of $`M`$: hence the name ‘geometric quantization’. Moreover, in geometric quantization one merely quantizes a subspace of the set $`C^{\mathrm{}}(M)`$ of classical observables, consisting of those functions that satisfy
$$[𝒬_{\mathrm{}}^{pre}(f),p]=0.$$
(4.45)
If a function $`fC^{\mathrm{}}(M)`$ satisfies this condition, then one defines the “geometric quantization” of $`f`$ as
$$𝒬_{\mathrm{}}^G(f)=𝒬_{\mathrm{}}^{pre}(f)(M).$$
(4.46)
This is well defined, since because of (4.45) the operator $`𝒬_{\mathrm{}}^{pre}(f)`$ now maps $`pL^2(M)`$ onto itself. Hence (4.38) holds for $`𝒬_{\mathrm{}}=𝒬_{\mathrm{}}^G`$ because of (4.39); in geometric quantization one simply refuses to quantize functions for which (4.38) is not valid.
Despite some impressive initial triumphs,<sup>176</sup><sup>176</sup>176Such as the orbit method for nilpotent groups and the newly understood Borel–Weil method for compact groups, cf. Kirillov (2004) and most other books cited in footnote 159. there is no general method that accomplishes the goals of geometric quantization with guaranteed success. Therefore, geometric quantization has remained something like a hacker’s tool, whose applicability largely depends on the creativity of the user.
In any case, our familiar example $`M=^{2n}`$ is well understood, and we illustrate the general spirit of the method in its setting, simplifying further by taking $`n=1`$. It is convenient to replace the canonical coordinates $`(p,q)`$ on $`M`$ by $`z=p+iq`$ and $`\overline{z}=piq`$, and the mathematical toolkit of geometric quantization makes it very natural to look at the space of solutions within $`L^2(^2)`$ of the equations<sup>177</sup><sup>177</sup>177Using the formalism explained in footnote 168, we replace the 1-form $`A=\frac{i}{\mathrm{}}_kp_kdq^k`$ defining the connection $`=d+A`$ by the gauge-equivalent form $`A=\frac{i}{2\mathrm{}}(_kq^kdp_k_kp_kdq^k)=\frac{i}{\mathrm{}}_kp_kdq^k+\frac{i}{2\mathrm{}}d(_kp_kq^k)`$, which has the same curvature. In terms of this new $`A`$, which in complex coordinates reads $`A=_k(z_kd\overline{z}_k\overline{z}_kdz_k)/4\mathrm{}`$, eq. (4.47) is just $`_{/\overline{z}}\mathrm{\Phi }=0`$. This is an example of the so-called holomorphic polarization in the formalism of geometric quantization.
$$\left(\frac{}{\overline{z}}+\frac{z}{4\mathrm{}}\right)\mathrm{\Phi }(z,\overline{z})=0.$$
(4.47)
The general solution of these equations that lies in $`L^2(^2)=L^2()`$ is
$$\mathrm{\Phi }(z,\overline{z})=e^{|z|^2/4\mathrm{}}f(z),$$
(4.48)
where $`f`$ is a holomorphic function such that
$$_{}\frac{dzd\overline{z}}{2\pi \mathrm{}i}e^{|z|^2/2\mathrm{}}|f(z)|^2<\mathrm{}.$$
(4.49)
The projection $`p`$, then, is the projection onto the closed subspace of $`L^2()`$ consisting of these solutions.<sup>178</sup><sup>178</sup>178 The collection of all holomorphic functions on $``$ satisfying (4.49) is a Hilbert space with respect to the inner product $`(f,g)=(2\pi \mathrm{}i)^1_{}𝑑z𝑑\overline{z}\mathrm{exp}(|z|^2/2\mathrm{})\overline{f(z)}g(z)`$, called the Bargmann–Fock space $`_{BF}`$. This space may be embedded in $`L^2()`$ by $`f(z)\mathrm{exp}(|z|^2/2\mathrm{})f(z)`$, and the image of this embedding is of course just $`pL^2()`$. The Hilbert space $`pL^2()`$ is unitarily equivalent to $`L^2()`$ in a natural way (i.e. without the choice of a basis). The condition (4.45) boils down to $`^2f(z,\overline{z})/\overline{z}_i\overline{z}_j=0`$; in particular, the coordinate functions $`q`$ and $`p`$ are quantizable. Transforming to $`L^2()`$, one finds that the operators $`𝒬_{\mathrm{}}^G(q)`$ and $`𝒬_{\mathrm{}}^G(p)`$ coincide with Schrödinger’s expressions (2.2). In particular, the Heisenberg group $`H_1`$, which acts with infinite multiplicity on $`L^2()`$, acts irreducibly on $`pL^2()`$.
### 4.5 Epilogue: functoriality of quantization
A very important aspect of quantization is its interplay with symmetries and constraints. Indeed, the fundamental theories describing Nature (viz. electrodynamics, Yang–Mills theory, general relativity, and possibly also string theory) are a priori formulated as constrained systems. The classical side of constraints and reduction is well understood,<sup>179</sup><sup>179</sup>179See Gotay, Nester, & Hinds (1978), Binz, Śniatycki and Fischer (1988), Marsden (1992), Marsden & Ratiu (1994), Landsman (1998), Butterfield (2005), and Belot (2005). a large class of important examples being codified by the procedure of symplectic reduction. A special case of this is Marsden–Weinstein reduction: if a Lie group $`G`$ acts on a phase space $`M`$ in canonical fashion with momentum map $`\mu :M𝔤^{}`$ (cf. Subsection 4.4), one may form another phase space $`M//G=\mu ^1(0)/G`$.<sup>180</sup><sup>180</sup>180Technically, $`M`$ has to be a symplectic manifold, and if $`G`$ acts properly and freely on $`\mu ^1(0)`$, then $`M//G`$ is again a symplectic manifold. Physically, in the case where $`G`$ is a gauge group and $`M`$ is the unconstrained phase space, $`\mu ^1(0)`$ is the constraint hypersurface (i.e. the subspace of $`M`$ on which the constraints defined by the gauge symmetry hold), and $`M//G`$ is the true phase space of the system that only contains physical degrees of freedom.
Unfortunately, the correct way of dealing with constrained quantum systems remains a source of speculation and controversy:<sup>181</sup><sup>181</sup>181Cf. Dirac (1964), Sundermeyer (1982), Gotay (1986), Duval et al. (1991), Govaerts (1991), Henneaux & Teitelboim (1992), and Landsman (1998) for various perspectives on the quantization of constrained systems. practically all rigorous results on quantization (like the ones discussed in the preceding subsections) concern unconstrained systems. Accordingly, one would like to quantize a constrained system by reducing the problem to the unconstrained case. This could be done provided the following scenario applies. One first quantizes the unconstrained phase space $`M`$ (supposedly the easiest part of the problem), ans subsequently imposes a quantum version of symplectic reduction. Finally, one proves by abstract means that the quantum theory thus constructed is equal to the theory defined by first reducing at the classical level and then quantizing the constrained classical phase space (usually an impossible task to perform in practice).
Tragically, sufficiently powerful theorems stating that “quantization commutes with reduction” in this sense remain elusive.<sup>182</sup><sup>182</sup>182 The so-called Guillemin–Sternberg conjecture (Guillemin & Sternberg, 1982) - now a theorem (Meinrenken, 1998, Meinrenken & Sjamaar, 1999) - merely deals with the case of Marsden–Weinstein reduction where $`G`$ and $`M`$ are compact. Mathematically impressive as the “quantization commutes with reduction” theorem already is here, it is a far call from gauge theories, where the groups and spaces are not only noncompact but even infinite-dimensional. So far, this has blocked, for example, a rigorous quantization of Yang–Mills theory in dimension 4; this is one of the Millenium Problems of the Clay Mathematical Institute, rewarded with 1 Million dollars.<sup>183</sup><sup>183</sup>183See http://www.claymath.org/millennium/
On a more spiritual note, the mathematician E. Nelson famously said that ‘First quantization is a mystery, but second quantization is a functor.’ The functoriality of ‘second’ quantization (a construction involving Fock spaces, see Reed & Simon, 1975) being an almost trivial matter, the deep mathematical and conceptual problem lies in the possible functoriality of ‘first’ quantization, which simply means quantization in the sense we have been discussing so far. This was initially taken to mean that canonical transformations $`\alpha `$ of the phase space $`M`$ should be ‘quantized’ by unitary operators $`U(\alpha )`$ on $`(M)`$, in such a way $`U(\alpha )𝒬_{\mathrm{}}(f)U(\alpha )^1=𝒬(L_\alpha f)`$ (cf. (4.10)). This is possible only in special circumstances, e.g., when $`M=^{2n}`$ and $`\alpha `$ is a linear symplectic map, and more generally when $`M=G/H`$ is homogeneous and $`\alpha G`$ (see the end of Subsection 4.2).<sup>184</sup><sup>184</sup>184 Canonical transformations can be quantized in approximate sense that becomes precise as $`\mathrm{}0`$ by means of so-called Fourier integral operators; see Hörmander (1971, 1985b) and Duistermaat (1996). Consequently, the functoriality of quantization is widely taken to be a dead end.<sup>185</sup><sup>185</sup>185See Groenewold (1946), van Hove (1951), Gotay, Grundling, & Tuynman (1996), and Gotay (1999).
However, all no-go theorems establishing this conclusion start from wrong and naive categories, both on the classical and on the quantum side.<sup>186</sup><sup>186</sup>186Typically, one takes the classical category to consist of symplectic manifolds as objects and symplectomorphisms as arrows, and the quantum category to have $`C^{}`$-algebras as objects and automorphisms as arrows. It appears very likely that one may indeed make quantization functorial by a more sophisticated choice of categories, with the additional bonus that deformation quantization and geometric quantization become unified: the former is the object part of the quantization functor, whereas the latter (suitably reinterpreted) is the arrow part. Amazingly, on this formulation the statement that ‘quantization commutes with reduction’ becomes a special case of the functoriality of quantization (Landsman, 2002, 2005a).
To explain the main idea, we return to the geometric quantization of $`M=^2`$ explained in the preceding subsection. The identification of $`pL^2()`$<sup>187</sup><sup>187</sup>187Or the Bargmann–Fock space $`_{BF}`$, see footnote 178. as the correct Hilbert space of the problem may be understood in a completely different way, which paves the way for the powerful reformulation of the geometric quantization program that will eventually define the quantization functor. Namely, $``$ supports a certain linear first-order differential operator $`D\text{/}`$ that is entirely defined by its geometry as a phase space, called the Dirac operator.<sup>188</sup><sup>188</sup>188 Specifically, this is the so-called $`\mathrm{Spin}^c`$ Dirac operator defined by the complex structure of $``$, coupled to the prequantization line bundle. See Guillemin, Ginzburg, & Karshon (2002). This operator is given by<sup>189</sup><sup>189</sup>189Relative to the Dirac matrices $`\gamma ^1=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)`$ and $`\gamma ^2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$.
$$D\text{/}=2\left(\begin{array}{cc}0& \frac{}{z}+\frac{\overline{z}}{4\mathrm{}}\\ \frac{}{\overline{z}}+\frac{z}{4\mathrm{}}& 0\end{array}\right),$$
(4.50)
acting on $`L^2()^2`$ (as a suitably defined unbounded operator). This operator has the generic form
$$D\text{/}=\left(\begin{array}{cc}0& D\text{/}_{}\\ D\text{/}_+& 0\end{array}\right).$$
The index of such an operator is given by
$$\mathrm{index}(D\text{/})=[\mathrm{ker}(D\text{/}_+)][\mathrm{ker}(D\text{/}_{})],$$
(4.51)
where $`[\mathrm{ker}(D\text{/}_\pm )]`$ stand for the (unitary) isomorphism class of $`\mathrm{ker}(D\text{/}_\pm )`$ seen as a representation space of a suitable algebra of operators.<sup>190</sup><sup>190</sup>190 The left-hand side of (4.51) should really be written as $`\mathrm{index}(D\text{/}_+)`$, since $`\mathrm{coker}(D\text{/}_+)=\mathrm{ker}(D\text{/}_+^{})`$ and $`D\text{/}_+^{}=D\text{/}_{}`$, but since the index is naturally associated to $`D\text{/}`$ as a whole, we abuse notation in writing $`\mathrm{index}(D\text{/})`$ for $`\mathrm{index}(D\text{/}_+)`$. The usual index of a linear map $`L:VW`$ between finite-dimensional vector spaces is defined as $`\mathrm{index}(L)=dim(\mathrm{ker}(L))dim(\mathrm{coker}(L))`$, where $`\mathrm{coker}(L)=W/\mathrm{ran}(L)`$. Elementary linear algebra yields $`\mathrm{index}(L)=dim(V)dim(W)`$. This is surprising because it is independent of $`L`$, whereas $`dim(\mathrm{ker}(L))`$ and $`dim(\mathrm{coker}(L))`$ quite sensitively depend on it. For, example, take $`V=W`$ and $`L=\epsilon 1`$. If $`\epsilon 0`$ then $`dim(\mathrm{ker}(\epsilon 1))=dim(\mathrm{coker}(\epsilon 1))=0`$, whereas for $`\epsilon =0`$ one has $`dim(\mathrm{ker}(0))=dim(\mathrm{coker}(0))=dim(V)`$! Similarly, the usual definiton of geometric quantization through (4.47) etc. is unstable against perturbations of the underlying symplectic structure, whereas the refined definition through (4.51) is not. To pass to the latter from the above notion of an index, we first write $`\mathrm{index}(L)=[\mathrm{ker}(L)][\mathrm{coker}(L)]`$, where $`[X]`$ is the isomorphism class of a linear space $`X`$ as a $``$-module. This expression is an element of $`K_0()`$, and we recover the earlier index through the realization that the class $`[X]`$ is entirely determined by $`dim(X)`$, along with and the corresponding isomorphism $`K_0()`$. When a more complicated finite-dimensional $`C^{}`$-algebra $`𝒜`$ acts on $`V`$ and $`W`$ with the property that $`\mathrm{ker}(L)`$ and $`\mathrm{coker}(L)`$ are stable under the $`𝒜`$-action, one may define $`[\mathrm{ker}(L)][\mathrm{coker}(L)]`$ and hence $`\mathrm{index}(L)`$ as an element of the so-called $`C^{}`$-algebraic K-theory group $`K_0(𝒜)`$. Under certain technical conditions, this notion of an index may be generalized to infinite-dimensional Hilbert spaces and $`C^{}`$-algebras; see Baum, Connes & Higson (1994) and Blackadar (1998). The $`K`$-theoretic index is best understood when $`𝒜=C^{}(G)`$ is the group $`C^{}`$-algebra of some locally compact group $`G`$. In the example $`M=^2`$ one might take $`G`$ to be the Heisenberg group $`H_1`$, so that $`\mathrm{index}(D\text{/})K_0(C^{}(H_1))`$. In the case at hand, one has $`\mathrm{ker}(D\text{/}_+)=pL^2()`$ (cf. (4.47) etc.) and $`\mathrm{ker}(D\text{/}_{})=0`$, <sup>191</sup><sup>191</sup>191Since $`(\frac{}{z}+\frac{\overline{z}}{4\mathrm{}})\mathrm{\Phi }=0`$ implies $`\mathrm{\Phi }(z,\overline{z})=\mathrm{exp}(|z^2|/4\mathrm{})f(\overline{z})`$, which lies in $`L^2()`$ iff $`f=0`$. where we regard $`\mathrm{ker}(D\text{/}_+)`$ as a representation space of the Heisenberg group $`H_1`$. Consequently, the geometric quantization of the phase space $``$ is given modulo unitary equivalence by $`\mathrm{index}(D\text{/})`$, seen as a “formal difference” of representation spaces of $`H_1`$.
This procedure may be generalized to arbitrary phase spaces $`M`$, where $`D\text{/}`$ is a certain operator naturally defined by the phase space geometry of $`M`$ and the demands of quantization.<sup>192</sup><sup>192</sup>192 Any symplectic manifold carries an almost complex structure compatible with the symplectic form, leading to a $`\mathrm{Spin}^c`$ Dirac operator as described in footnote 188. See, again, Guillemin, Ginzburg, & Karshon (2002). If $`M=G/H`$, or, more generally, if $`M`$ carries a canonical action of a Lie group $`G`$ with compact quotient $`M/G`$, then $`\mathrm{index}(D\text{/})`$ defines an element of $`K_0(C^{}(G))`$. See footnote 190. In complete generality, $`\mathrm{index}(D\text{/})`$ ought to be an element of $`K_0(𝒜)`$, where $`𝒜`$ is the $`C^{}`$-algebra of observables of the quantum system. This has turned out to be the most promising formulation of geometric quantization - at some cost.<sup>193</sup><sup>193</sup>193On the benefit side, the invariance of the index under continuous deformations of $`D\text{/}`$ seems to obviate the ambiguity of traditional quantization procedures with respect to different ‘operator orderings’ not prescribed by the classical theory. For the original goal of quantizing a phase space by a Hilbert space has now been replaced by a much more abstract procedure, in which the result of quantization is a formal difference of certain isomorphism classes of representation spaces of the quantum algebra of observables. To illustrate the degree of abstraction involved here, suppose we ignore the action of the observables (such as position and momentum in the example just considered). In that case the isomorphism class $`[]`$ of a Hilbert space $``$ is entirely characterized by its dimension $`dim()`$, so that (in case that $`\mathrm{ker}(D\text{/}_{})0`$) quantization (in the guise of $`\mathrm{index}(D\text{/})`$) can even be a negative number! Have we gone mad?
Not quite. The above picture of geometric quantization is indeed quite irrelevant to physics, unless it is supplemented by deformation quantization. It is convenient to work at some fixed value of $`\mathrm{}`$ in this context, so that deformation quantization merely associates some $`C^{}`$-algebra $`𝒜(P)`$ to a given phase space $`P`$.<sup>194</sup><sup>194</sup>194Here $`P`$ is not necessarily symplectic; it may be a Poisson manifold, and to keep Poisson and symplectic manifolds apart we denote the former by $`P`$ from now on, preserving the notation $`M`$ for the latter. Looking for a categorical interpretation of quantization, it is therefore natural to assume that the objects of the classical category $``$ are phase spaces $`P`$,<sup>195</sup><sup>195</sup>195Strictly speaking, to be an object in this category a Poisson manifold $`P`$ must be integrable; see Landsman (2001). whereas the objects of the quantum category $`𝔔`$ are $`C^{}`$-algebras.<sup>196</sup><sup>196</sup>196For technical reasons involving $`K`$-theory these have to be separable. The object part of the hypothetical quantization functor is to be deformation quantization, symbolically written as $`P𝒬(P)`$.
Everything then fits together if geometric quantization is reinterpreted as the arrow part of the conjectural quantization functor. To accomplish this, the arrows in the classical category $``$ should not be taken to be maps between phase spaces, but symplectic bimodules $`P_1MP_2`$.<sup>197</sup><sup>197</sup>197Here $`M`$ is a symplectic manifold and $`P_1`$ and $`P_2`$ are integrable Poisson manifolds; the map $`MP_2`$ is anti-Poisson, whereas the map $`P_1M`$ is Poisson. Such bimodules (often called dual pairs) were introduced by Karasev (1989) and Weinstein (1983). In order to occur as arrows in $``$, symplectic bimodules have to satisfy a number of regularity conditions (Landsman, 2001). More precisely, the arrows in $``$ are suitable isomorphism classes of such bimodules.<sup>198</sup><sup>198</sup>198This is necessary in order to make arrow composition associative; this is given by a generalization of the symplectic reduction procedure. Similarly, the arrows in the quantum category $`𝔔`$ are not morphisms of $`C^{}`$-algebras, as might naively be expected, but certain isomorphism classes of bimodules for $`C^{}`$-algebras, equipped with the additional structure of a generalized Dirac operator.<sup>199</sup><sup>199</sup>199The category $`𝔔`$ is nothing but the category KK introduced by Kasparov, whose objects are separable $`C^{}`$-algebras, and whose arrows are the so-called Kasparov group $`KK(A,B)`$, composed with Kasparov’s product $`KK(A,B)\times KK(B,C)KK(A,C)`$. See Higson (1990) and Blackadar (1998).
Having already defined the object part of the quantization map $`𝒬:𝔔`$ as deformation quantization, we now propose that the arrow part is geometric quantization, in the sense of a suitable generalization of (4.51); see Landsman (2005a) for details. We then conjecture that $`𝒬`$ is a functor; in the cases where this can and has been checked, the functoriality of $`𝒬`$ is precisely the statement that quantization commutes with reduction.<sup>200</sup><sup>200</sup>200A canonical $`G`$-action on a symplectic manifold $`M`$ with momentum map $`\mu :M𝔤^{}`$ gives rise to a dual pair $`ptM𝔤^{}`$, which in $``$ is interpreted as an arrow from the space $`pt`$ with one point to $`𝔤^{}`$. The composition of this arrow with the arrow $`𝔤^{}0pt`$ from $`𝔤^{}`$ to $`pt`$ is $`ptM//Gpt`$. If $`G`$ is connected, functoriality of quantization on these two pairs is equivalent to the Guillemin–Sternberg conjecture (cf. footnote 182); see Landsman (2005a).
Thus Heisenberg’s idea of Umdeutung finds it ultimate realization in the quantization functor.
## 5 The limit $`\mathrm{}0`$
It was recognized at an early stage that the limit $`\mathrm{}0`$ of Planck’s constant going to zero should play a role in the explanation of the classical world from quantum theory. Strictly speaking, $`\mathrm{}`$ is a dimensionful constant, but in practice one studies the semiclassical regime of a given quantum theory by forming a dimensionless combination of $`\mathrm{}`$ and other parameters; this combination then re-enters the theory as if it were a dimensionless version of $`\mathrm{}`$ that can indeed be varied. The oldest example is Planck’s radiation formula (2.1), with temperature $`T`$ as the pertinent variable. Indeed, the observation of Einstein (1905) and Planck (1906) that in the limit $`\mathrm{}\nu /kT0`$ this formula converges to the classical equipartition law $`E_\nu /N_\nu =kT`$ may well be the first use of the $`\mathrm{}0`$ limit of quantum theory.<sup>201</sup><sup>201</sup>201Here Einstein (1905) put $`\mathrm{}\nu /kT0`$ by letting $`\nu 0`$ at fixed $`T`$ and $`\mathrm{}`$, whereas Planck (1906) took $`T\mathrm{}`$ at fixed $`\nu `$ and $`\mathrm{}`$.
Another example is the Schrödinger equation (2.3) with Hamiltonian $`H=\frac{\mathrm{}^2}{2m}\mathrm{\Delta }_x+V(x)`$, where $`m`$ is the mass of the pertinent particle. Here one may pass to dimensionless parameters by introducing an energy scale $`ϵ`$ typical of $`H`$, like $`ϵ=sup_x|V(x)|`$, as well as a typical length scale $`\lambda `$, such as $`\lambda =ϵ/sup_x|V(x)|`$ (if these quantities are finite). In terms of the dimensionless variable $`\stackrel{~}{x}=x/\lambda `$, the rescaled Hamiltonian $`\stackrel{~}{H}=H/ϵ`$ is then dimensionless and equal to $`\stackrel{~}{H}=\stackrel{~}{\mathrm{}}^2\mathrm{\Delta }_{\stackrel{~}{x}}+\stackrel{~}{V}(\stackrel{~}{x})`$, where $`\stackrel{~}{\mathrm{}}=\mathrm{}/\lambda \sqrt{2mϵ}`$ and $`\stackrel{~}{V}(\stackrel{~}{x})=V(\lambda \stackrel{~}{x})/ϵ`$. Here $`\stackrel{~}{\mathrm{}}`$ is dimensionless, and one might study the regime where it is small (Gustafson & Sigal, 2003). Our last example will occur in the theory of large quantum systems, treated in the next Section. In what follows, whenever it is considered variable $`\mathrm{}`$ will denote such a dimensionless version of Planck’s constant.
Although, as we will argue, the limit $`\mathrm{}0`$ cannot by itself explain the classical world, it does give rise to a number of truly pleasing mathematical results. These, in turn, render almost inescapable the conclusion that the limit in question is indeed a relevant one for the recovery of classical physics from quantum theory. Thus the present section is meant to be a catalogue of those pleasantries that might be of direct interest to researchers in the foundations of quantum theory.
There is another, more technical use of the $`\mathrm{}0`$ limit, which is to perform computations in quantum mechanics by approximating the time-evolution of states and observables in terms of associated classical objects. This endeavour is known as semiclassical analysis. Mathematically, this use of the $`\mathrm{}0`$ limit is closely related to the goal of recovering classical mechanics from quantum mechanics, but conceptually the matter is quite different. We will attempt to bring the pertinent differences out in what follows.
### 5.1 Coherent states revisited
As Schrödinger (1926b) foresaw, coherent states play an important role in the limit $`\mathrm{}0`$. We recall from Subsection 4.2 that for some fixed value $`\mathrm{}`$ of Planck’s constant coherent states in a Hilbert space $``$ for a phase space $`M`$ are defined by an injection $`M`$, $`z\mathrm{\Psi }_z^{\mathrm{}}`$, such that (4.16) and (4.17) hold. In what follows, we shall say that $`\mathrm{\Psi }_z^{\mathrm{}}`$ is centered at $`zM`$, a terminology justified by the key example (4.20).
To be relevant to the classical limit, coherent states must satisfy an additional property concerning their dependence on $`\mathrm{}`$, which also largely clarifies their nature (Landsman, 1998). Namely, we require that for each $`fC_c(M)`$ and each $`zM`$ the following function from the set $`I`$ in which $`\mathrm{}`$ takes values (i.e. usually $`I=[0,1]`$, but in any case containing zero as an accumulation point) to $``$ is continuous:
$`\mathrm{}`$ $``$ $`c_{\mathrm{}}{\displaystyle _M}𝑑\mu _L(w)|(\mathrm{\Psi }_w^{\mathrm{}},\mathrm{\Psi }_z^{\mathrm{}})|^2f(w)(\mathrm{}>0);`$ (5.1)
$`0`$ $``$ $`f(z).`$ (5.2)
In view of (4.19), the right-hand side of (5.2) is the same as $`(\mathrm{\Psi }_z^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }_z^{\mathrm{}})`$. In particular, this continuity condition implies
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }_z^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }_z^{\mathrm{}})=f(z).$$
(5.3)
This means that the classical limit of the quantum-mechanical expectation value of the phase space quantization (4.19) of the classical observable $`f`$ in a coherent state centered at $`zM`$ is precisely the classical expectation value of $`f`$ in the state $`z`$. This interpretation rests on the identification of classical states with probability measures on phase space $`M`$, under which points of $`M`$ in the guise of Dirac measures (i.e. delta functions) are pure states. Furthermore, it can be shown (cf. Landsman, 1998) that the continuity of all functions (5.1) - (5.2) implies the property
$$\underset{\mathrm{}0}{lim}|(\mathrm{\Psi }_w^{\mathrm{}},\mathrm{\Psi }_z^{\mathrm{}})|^2=\delta _{wz},$$
(5.4)
where $`\delta _{wz}`$ is the ordinary Kronecker delta (i.e. $`\delta _{wz}=0`$ whenever $`wz`$ and $`\delta _{zz}=1`$ for all $`zM`$). This has a natural physical interpretation as well: the classical limit of the quantum-mechanical transition probability between two coherent states centered at $`w,zM`$ is equal to the classical (and trivial) transition probability between $`w`$ and $`z`$. In other words, when $`\mathrm{}`$ becomes small, coherent states at different values of $`w`$ and $`z`$ become increasingly orthogonal to each other.<sup>202</sup><sup>202</sup>202See Mielnik (1968), Cantoni (1975), Beltrametti & Cassinelli (1984), Landsman (1998), and Subsection 6.3 below for the general meaning of the concept of a transition probability. This has the interesting consequence that
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }_w^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }_z^{\mathrm{}})=0(wz).$$
(5.5)
for all $`fC_c(M)`$. In particular, the following phenomenon of the Schrödinger cat type occurs in the classical limit: if $`wz`$ and one has continuous functions $`\mathrm{}c_w^{\mathrm{}}`$ and $`\mathrm{}c_z^{\mathrm{}}`$ on $`\mathrm{}[0,1]`$ such that
$$\mathrm{\Psi }_{w,z}^{\mathrm{}}=c_w^{\mathrm{}}\mathrm{\Psi }_w^{\mathrm{}}+c_z^{\mathrm{}}\mathrm{\Psi }_z^{\mathrm{}}$$
(5.6)
is a unit vector for $`\mathrm{}0`$ and also $`|c_w^0|^2+|c_z^0|^2=1`$, then
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }_{w,z}^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }_{w,z}^{\mathrm{}})=|c_w^0|^2f(w)+|c_z^0|^2f(z).$$
(5.7)
Hence the family of (typically) pure states $`\psi _{w,z}^{\mathrm{}}`$ (on the $`C^{}`$-algebras $`𝒜_{\mathrm{}}`$ in which the map $`𝒬_{\mathrm{}}^B`$ takes values)<sup>203</sup><sup>203</sup>203For example, for $`M=^{2n}`$ each $`𝒜_{\mathrm{}}`$ is equal to the $`C^{}`$-algebra of compact operators on $`L^2(^n)`$, on which each vector state is certainly pure. defined by the vectors $`\mathrm{\Psi }_{w,z}^{\mathrm{}}`$ in some sense converges to the mixed state on $`C_0(M)`$ defined by the right-hand side of (5.7). This is made precise at the end of this subsection.
It goes without saying that Schrödinger’s coherent states (4.20) satisfy our axioms; one may also verify (5.4) immediately from (4.21). Consequently, by (4.32) one has the same property (5.3) for Weyl quantization (as long as $`f𝒮(^{2n})`$),<sup>204</sup><sup>204</sup>204Here $`𝒮(^{2n})`$ is the usual Schwartz space of smooth test functions with rapid decay at infinity. that is,
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }_z^{\mathrm{}},𝒬_{\mathrm{}}^W(f),\mathrm{\Psi }_z^{\mathrm{}})=f(z).$$
(5.8)
Similarly, (5.5) holds for $`𝒬_{\mathrm{}}^W`$ as well.
In addition, many constructions referred to as coherent states in the literature (cf. the references in footnote 130) satisfy (4.16), (4.17), and (5.4); see Landsman (1998).<sup>205</sup><sup>205</sup>205For example, coherent states of the type introduced by Perelomov (1986) fit into our setting as follows (Simon, 1980). Let $`G`$ be a compact connected Lie group, and $`𝒪_\lambda `$ an integral coadjoint orbit, corresponding to a highest weight $`\lambda `$. (One may think here of $`G=SU(2)`$ and $`\lambda =0,1/2,1,\mathrm{}`$.) Note that $`𝒪_\lambda G/T`$, where $`T`$ is the maximal torus in $`G`$ with respect to which weights are defined. Let $`_\lambda ^{\text{hw}}`$ be the carrier space of the irreducible representation $`U_\lambda (G)`$ with highest weight $`\lambda `$, containing the highest weight vector $`\mathrm{\Omega }_\lambda `$. (For $`G=SU(2)`$ one has $`_j^{\text{hw}}=^{2j+1}`$, the well-known Hilbert space of spin $`j`$, in which $`\mathrm{\Omega }_j`$ is the vector with spin $`j`$ in the $`z`$-direction.) For $`\mathrm{}=1/k`$, $`k`$, define $`_{\mathrm{}}:=_{\lambda /\mathrm{}}^{\text{hw}}`$. Choosing a section $`\sigma :𝒪_\lambda G`$ of the projection $`GG/T`$, one then obtains coherent states $`xU_{\lambda /\mathrm{}}(\sigma (x))\mathrm{\Omega }_{\lambda /\mathrm{}}`$ with respect to the Liouville measure on $`𝒪_\lambda `$ and $`c_{\mathrm{}}=dim(_{\lambda /\mathrm{}}^{\text{hw}})`$. These states are obviously not defined for all values of $`\mathrm{}`$ in $`(0,1]`$, but only for the discrete set $`1/`$. The general picture that emerges is that a coherent state centered at $`zM`$ is the Umdeutung of $`z`$ (seen as a classical pure state, as explained above) as a quantum-mechanical pure state.<sup>206</sup><sup>206</sup>206This idea is also confirmed by the fact that at least Schrödinger’s coherent states are states of minimal uncertainty; cf. the references in footnote 130.
Despite their wide applicability (and some would say beauty), one has to look beyond coherent states for a complete picture of the $`\mathrm{}0`$ limit of quantum mechanics. The appropriate generalization is the concept of a continuous field of states.<sup>207</sup><sup>207</sup>207The use of this concept in various mathematical approaches to quantization is basically folklore. For the $`C^{}`$-algebraic setting see Emch (1984), Rieffel (1989b), Werner (1995), Blanchard (1996), Landsman (1998), and Nagy (2000). This is defined relative to a given deformation quantization of a phase space $`M`$; cf. Subsection 4.3. If one now has a state $`\omega _{\mathrm{}}`$ on $`𝒜_{\mathrm{}}`$ for each $`\mathrm{}[0,1]`$ (or, more generally, for a discrete subset of $`[0,1]`$ containing 0 as an accumulation point), one may call the ensuing family of states a continuous field whenever the function $`\mathrm{}\omega _{\mathrm{}}(𝒬_{\mathrm{}}(f))`$ is continuous on $`[0,1]`$ for each $`fC_c^{\mathrm{}}(M)`$; this notion is actually intrinsically defined by the continuous field of $`C^{}`$-algebras, and is therefore independent of the quantization maps $`𝒬_{\mathrm{}}`$. In particular, one has
$$\underset{\mathrm{}0}{lim}\omega _{\mathrm{}}(𝒬_{\mathrm{}}(f))=\omega _0(f).$$
(5.9)
Eq. (5.3) (or (5.8)) shows that coherent states are indeed examples of continuous fields of states, with the additional property that each $`\omega _{\mathrm{}}`$ is pure. As an example where all states $`\omega _{\mathrm{}}`$ are mixed, we mention the convergence of quantum-mechanical partition functions to their classical counterparts in statistical mechanics along these lines; see Lieb (1973), Simon (1980), Duffield (1990), and Nourrigat & Royer (2004). Finally, one encounters the surprising phenomenon that pure quantum states may coverge to mixed classical ones. The first example of this has just been exhibited in (5.7); other cases in point are energy eigenstates and WKB states (see Subsections 5.4, 5.5, and 5.6 below).
### 5.2 Convergence of quantum dynamics to classical motion
Nonrelativistic quantum mechanics is based on the Schrödinger equation (2.3), which more generally reads
$$H\mathrm{\Psi }(t)=i\mathrm{}\frac{\mathrm{\Psi }}{t}.$$
(5.10)
The formal solution with initial value $`\mathrm{\Psi }(0)=\mathrm{\Psi }`$ is
$$\mathrm{\Psi }(t)=e^{\frac{it}{\mathrm{}}H}\mathrm{\Psi }.$$
(5.11)
Here we have assumed that $`H`$ is a given self-adjoint operator on the Hilbert space $``$ of the system, so that this solution indeed exists and evolves unitarily by Stone’s theorem; cf. Reed & Simon (1972) and Simon (1976). Equivalently, one may transfer the time-evolution from states (Schrödinger picture) to operators (Heisenberg picture) by putting
$$A(t)=e^{\frac{it}{\mathrm{}}H}Ae^{\frac{it}{\mathrm{}}H}.$$
(5.12)
We here restrict ourselves to particle motion in $`^n`$, so that $`=L^2(^n)`$.<sup>208</sup><sup>208</sup>208See Hunziker & Sigal (2000) for a recent survey of $`N`$-body Schrödinger operators. In that case, $`H`$ is typically given by a formal expression like (2.3) (on some specific domain).<sup>209</sup><sup>209</sup>209 One then has to prove self-adjointness (or the lack of it) on a larger domain on which the operator is closed; see the literature cited in footnote 44. Now, the first thing that comes to mind is Ehrenfest’s Theorem (1927), which states that for any (unit) vector $`\mathrm{\Psi }L^2(^n)`$ in the domain of $`𝒬_{\mathrm{}}(q^j)=x^j`$ and $`V(x)/x^j`$ one has
$$m\frac{d^2}{dt^2}x^j(t)=\frac{V(x)}{x^j}(t),$$
(5.13)
with the notation
$`x^j(t)`$ $`=`$ $`(\mathrm{\Psi }(t),x^j\mathrm{\Psi }(t));`$
$`{\displaystyle \frac{V(x)}{x^j}}(t)`$ $`=`$ $`(\mathrm{\Psi }(t),{\displaystyle \frac{V(x)}{x^j}}\mathrm{\Psi }(t)).`$ (5.14)
This looks like Newton’s second law for the expectation value of $`x`$ in the state $`\psi `$, with the tiny but crucial difference that Newton would have liked to see $`(V/x^j)(x(t))`$ on the right-hand side of (5.13). Furthermore, even apart from this point Ehrenfest’s Theorem by no means suffices to have classical behaviour, since it gives no guarantee whatsoever that $`x(t)`$ behaves like a point particle. Much of what follows can be seen as an attempt to sharpen Ehrenfest’s Theorem to the effect that it does indeed yield appropriate classical equations of motion for the expectation values of suitable operators.
We assume that the quantum Hamiltonian has the more general form
$$H=h(𝒬_{\mathrm{}}(p_j),𝒬_{\mathrm{}}(q^j)),$$
(5.15)
where $`h`$ is the classical Hamiltonian (i.e. a function defined on classical phase space $`^{2n}`$) and $`𝒬_{\mathrm{}}(p_j)`$ and $`𝒬_{\mathrm{}}(q^j)`$ are the operators given in (2.2). Whenever this expression is ambiguous (as in cases like $`h(p,q)=pq`$), one has to assume a specific quantization prescription such as Weyl quantization $`𝒬_{\mathrm{}}^W`$ (cf. (4.29)), so that formally one has
$$H=𝒬_{\mathrm{}}^W(h).$$
(5.16)
In fact, in the literature to be cited an even larger class of quantum Hamiltonians is treated by the methods explained here. The quantum Hamiltonian $`H`$ carries an explicit (and rather singular) $`\mathrm{}`$-dependence, and for $`\mathrm{}0`$ one then expects (5.11) or (5.12) to be related in one way or another to the flow of the classical Hamiltonian $`h`$. This relationship was already foreseen by Schrödinger (1926a), and was formalized almost immediately after the birth of quantum mechanics by the well-known WKB approximation (cf. Landau & Lifshitz (1977) and Subsection 5.5 below). A mathematically rigorous understanding of this and analogous approximation methods only emerged much later, when a technique called microlocal analysis was adapted from its original setting of partial differential equations (Hörmander, 1965; Kohn & Nirenberg, 1965; Duistermaat, 1974, 1996; Guillemin & Sternberg, 1977; Howe, 1980; Hörmander, 1979, 1985a, 1985b; Grigis & Sjöstrand, 1994) to the study of the $`\mathrm{}0`$ limit of quantum mechanics. This adaptation (often called semiclassical analysis) and its results have now been explained in various reviews written by the main players, notably Robert (1987, 1998), Helffer (1988), Paul & Uribe (1995), Colin de Verdière (1998), Ivrii (1998), Dimassi & Sjöstrand (1999), and Martinez (2002) (see also the papers in Robert (1992)). More specific references will be given below.<sup>210</sup><sup>210</sup>210For the heuristic theory of semiclassical asymptotics Landau & Lifshitz (1977) is a goldmine.
As mentioned before, the relationship between $`H`$ and $`h`$ provided by semiclassical analysis is double-edged. On the one hand, one obtains approximate solutions of (5.11) or (5.12), or approximate energy eigenvalues and energy eigenfunctions (sometimes called quasi-modes) for small values of $`\mathrm{}`$ in terms of classical data. This is how the results are usually presented; one computes specific properties of quantum theory in a certain regime in terms of an underlying classical theory. On the other hand, however, with some effort the very same results can often be reinterpreted as a partial explanation of the emergence of classical dynamics from quantum mechanics. It is the latter aspect of semiclassical analysis, somewhat understated in the literature, that is of interest to us. In this and the next three subsections we restrict ourselves to the simplest type of results, which nonetheless provide a good flavour of what can be achieved and understood by these means. By the same token, we just work with the usual flat phase space $`M=^{2n}`$ as before.
The simplest of all results relating classical and quantum dynamics is this:<sup>211</sup><sup>211</sup>211More generally, Egorov’s Theorem states that for a large class of Hamiltonians one has $`𝒬_{\mathrm{}}^W(f)(t)=𝒬_{\mathrm{}}^W(f_t)+O(\mathrm{})`$. See, e.g., Robert (1987), Dimassi & Sjöstrand (1999), and Martinez (2002).
> If the classical Hamiltonian $`h(p,q)`$ is at most quadratic in $`p`$ and $`q`$, and the Hamiltonian in (5.12) is given by (5.16), then
$$𝒬_{\mathrm{}}^W(f)(t)=𝒬_{\mathrm{}}^W(f_t).$$
(5.17)
Here $`f_t`$ is the solution of the classical equation of motion $`df_t/dt=\{h,f_t\}`$; equivalently, one may write
$$f_t(p,q)=f(p(t),q(t)),$$
(5.18)
where $`t(p(t),q(t))`$ is the classical Hamiltonian flow of $`h`$ with initial condition $`(p(0),q(0))=(p,q)`$. This holds for all decent $`f`$, e.g., $`f𝒮(^{2n})`$.
This result explains quantum in terms of classical, but the converse may be achieved by combining (5.17) with (5.9). This yields
$$\underset{\mathrm{}0}{lim}\omega _{\mathrm{}}(𝒬_{\mathrm{}}(f)(t))=\omega _0(f_t)$$
(5.19)
for any continuous field of states $`(\omega _{\mathrm{}})`$. In particular, for Schrödinger’s coherent states (4.20) one obtains
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }_{(p,q)}^{\mathrm{}},𝒬_{\mathrm{}}(f)(t)\mathrm{\Psi }_{(p,q)}^{\mathrm{}})=f_t(p,q).$$
(5.20)
Now, whereas (5.17) merely reflects the good symmetry properties of Weyl quantization,<sup>212</sup><sup>212</sup>212Eq. (5.17) is equivalent to the covariance of Weyl quantization under the affine symplectic group; cf. footnote 154. (and is false for $`𝒬_{\mathrm{}}^B`$), eq. (5.20) is actually valid for a large class of realistic Hamiltonians and for any deformation quantization map $`𝒬_{\mathrm{}}`$ that is asymptotically equal to $`𝒬_{\mathrm{}}^W`$ (cf. (4.32)). A result of this type was first established by Hepp (1974); further work in this direction includes Yajima (1979), Hogreve, Potthoff & Schrader (1983), Wang (1986), Robinson (1988a, 1988b), Combescure (1992), Arai (1995), Combescure & Robert (1997), Robert (1998), and Landsman (1998).
Impressive results are available also in the Schrödinger picture. The counterpart of (5.17) is that for any suitably smooth classical Hamiltonian $`h`$ (even a time-dependent one) that is at most quadratic in the canonical coordinates $`p`$ and $`q`$ on phase space $`^{2n}`$ one may construct generalized coherent states $`\mathrm{\Psi }_{(p,q,C)}^{\mathrm{}}`$, labeled by a set $`C`$ of classical parameters dictated by the form of $`h`$, such that
$$e^{\frac{it}{\mathrm{}}𝒬_{\mathrm{}}^W(h)}\mathrm{\Psi }_{(p,q,C)}^{\mathrm{}}=e^{iS(t)/\mathrm{}}\mathrm{\Psi }_{(p(t),q(t),C(t))}^{\mathrm{}}.$$
(5.21)
Here $`S(t)`$ is the action associated with the classical trajectory $`(p(t),q(t))`$ determined by $`h`$, and $`C(t)`$ is a solution of a certain system of differential equations that has a classical interpretation as well (Hagedorn, 1998). Schrödinger’s coherent states (4.20) are a special case for the standard harmonic oscillator Hamiltonian. For more general Hamiltonians one then has an asymptotic result (Hagedorn & Joye, 1999, 2000)<sup>213</sup><sup>213</sup>213See also Paul & Uribe (1995, 1996) as well as the references listed after (5.20) for analogous statements.
$$\underset{\mathrm{}0}{lim}e^{\frac{it}{\mathrm{}}𝒬_{\mathrm{}}^W(h)}\mathrm{\Psi }_{(p,q,C)}^{\mathrm{}}e^{iS(t)/\mathrm{}}\mathrm{\Psi }_{(p(t),q(t),C(t))}^{\mathrm{}}=0.$$
(5.22)
Once again, at first sight such results merely contribute to the understanding of quantum dynamics in terms of classical motion. As mentioned, they may be converted into statements on the emergence of classical motion from quantum mechanics by taking expectation values of suitable $`\mathrm{}`$-dependent obervables of the type $`𝒬_{\mathrm{}}^W(f)`$.
For finite $`\mathrm{}`$, the second term in (5.22) is a good approximation to the first \- the error even being as small as $`𝒪(\mathrm{exp}(\gamma /\mathrm{}))`$ for some $`\gamma >0`$ as $`\mathrm{}0`$ \- whenever $`t`$ is smaller than the so-called Ehrenfest time
$$T_E=\lambda ^^1\mathrm{log}(\mathrm{}^1),$$
(5.23)
where $`\lambda `$ is a typlical inverse time scale of the Hamiltonian (e.g., for chaotic systems it is the largest Lyapunov exponent).<sup>214</sup><sup>214</sup>214Recall that throughout this section we assume that $`\mathrm{}`$ has been made dimensionless through an appropriate rescaling. This is the typical time scale on which semiclassical approximations to wave packet solutions of the time-dependent Schrödinger equation with a general Hamiltonian tend to be valid (Ehrenfest, 1927; Berry et al., 1979; Zaslavsky, 1981; Combescure & Robert, 1997; Bambusi, Graffi, & Paul, 1999; Hagedorn & Joye, 2000).<sup>215</sup><sup>215</sup>215 One should distinguish here between two distinct approximation methods to the time-dependent Schrödinger equation. Firstly, one has the semiclassical propagation of a quantum-mechanical wave packet, i.e. its propagation as computed from the time-dependence of the parameters on which it depends according to the underlying classical equations of motion. It is shown in the references just cited that this approximates the full quantum-mechanical propagation of the wave packet well until $`tT_E`$. Secondly, one has the time-dependent WKB approximation (for integrable systems) and its generalization to chaotic systems (which typically involve tens of thousands of terms instead of a single one). This second approximation is valid on a much longer time scale, typically $`t\mathrm{}^{1/2}`$ (O’Connor, Tomsovic, & Heller, 1992; Heller & Tomsovic, 1993; Tomsovic, & Heller, 1993, 2002; Vanicek & Heller, 2003). Adding to the confusion, Ballentine has claimed over the years that even the semiclassical propagation of a wave packet approximates its quantum-mechanical propagation for times much longer than the Ehrenfest time, typically $`t\mathrm{}^{1/2}`$ (Ballentine, Yang, & Zibin, 1994; Ballentine, 2002, 2003). This claim is based on the criterion that the quantum and classical (i.e. Liouville) probabilities are approximately equal on such time scales, but the validity of this criterion hinges on the “statistical” or “ensemble” interpretation of quantum mechanics. According to this interpretation, a pure state provides a description of certain statistical properties of an ensemble of similarly prepared systems, but need not provide a complete description of an individual system. See Ballentine (1970, 1986). Though once defended by von Neumann, Einstein and Popper, this interpretation has now completely fallen out of fashion. For example, Ehrenfest (1927) himself estimated that for a mass of 1 gram a wave packet would double its width only in about $`10^{13}`$ years under free motion. However, Zurek and Paz (1995) have estimated the Ehrenfest time for Saturn’s moon Hyperion to be of the order of 20 years! This obviously poses a serious problem for the program of deriving (the appearance of) classical behaviour from quantum mechanics, which affects all interpretations of this theory.
Finally, we have not discussed the important problem of combining the limit $`t\mathrm{}`$ with the limit $`\mathrm{}0`$; this should be done in such a way that $`T_E`$ is kept fixed. This double limit is of particular importance for quantum chaos; see Robert (1998) and most of the literature cited in Subsection 5.6.
### 5.3 Wigner functions
The $`\mathrm{}0`$ limit of quantum mechanics is often discussed in terms of the so-called Wigner function, introduced by Wigner (1932).<sup>216</sup><sup>216</sup>216The original context was quantum statistical mechanics; one may write down (5.24) for mixed states as well. See Hillery et al. (1984) for a survey. Each unit vector (i.e. wave function) $`\mathrm{\Psi }L^2(^n)`$ defines such a function $`W_\mathrm{\Psi }^{\mathrm{}}`$ on classical phase space $`M=^{2n}`$ by demanding that for each $`f𝒮(^{2n})`$ one has
$$(\mathrm{\Psi },𝒬_{\mathrm{}}^W(f)\mathrm{\Psi })=_{^{2n}}\frac{d^npd^nq}{(2\pi )^n}W_\mathrm{\Psi }^{\mathrm{}}(p,q)f(p,q).$$
(5.24)
The existence of such a function may be proved by writing it down explicitly as
$$W_\mathrm{\Psi }^{\mathrm{}}(p,q)=_^nd^nve^{ipv}\overline{\mathrm{\Psi }(q+\frac{1}{2}\mathrm{}v)}\mathrm{\Psi }(q\frac{1}{2}\mathrm{}v).$$
(5.25)
In other words, the quantum-mechanical expectation value of the Weyl quantization of the classical observable $`f`$ in a quantum state $`\mathrm{\Psi }`$ formally equals the classical expectation value of $`f`$ with respect to the distribution $`W_\mathrm{\Psi }`$. However, the latter may not be regarded as a probability distribution because it is not necessarily positive definite.<sup>217</sup><sup>217</sup>217Indeed, it may not even be in $`L^1(^{2n})`$, so that its total mass is not necessarily defined, let alone equal to 1. Conditions for the positivity of Wigner functions defined by pure states are given by Hudson (1974); see Bröcker & Werner (1995) for the case of mixed states. Despite this drawback, the Wigner function possesses some attractive properties. For example, one has
$$𝒬_{\mathrm{}}^W(W_\mathrm{\Psi }^{\mathrm{}})=\mathrm{}^n[\mathrm{\Psi }].$$
(5.26)
This somewhat perverse result means that if the Wigner function defined by $`\mathrm{\Psi }`$ is seen as a classical observable (despite its manifest $`\mathrm{}`$-dependence!), then its Weyl quantization is precisely ($`\mathrm{}^n`$ times) the projection operator onto $`\mathrm{\Psi }`$.<sup>218</sup><sup>218</sup>218In other words, $`W_\mathrm{\Psi }`$ is the Weyl symbol of the projection operator $`[\mathrm{\Psi }]`$. Furthermore, one may derive the following formula for the transition probability:<sup>219</sup><sup>219</sup>219This formula is well defined since $`\mathrm{\Psi }L^2(^n)`$ implies $`W_\mathrm{\Psi }^{\mathrm{}}L^2(^{2n})`$.
$$|(\mathrm{\Phi },\mathrm{\Psi })|^2=\mathrm{}^n_{^{2n}}\frac{d^npd^nq}{(2\pi )^n}W_\mathrm{\Psi }^{\mathrm{}}(p,q)W_\mathrm{\Phi }^{\mathrm{}}(p,q).$$
(5.27)
This expression has immediate intuitive appeal, since the integrand on the right-hand side is supported by the area in phase space where the two Wigner functions overlap, which is well in tune with the idea of a transition probability.
The potential lack of positivity of a Wigner function may be remedied by noting that Berezin’s deformation quantization scheme (see (4.28)) analogously defines functions $`B_\mathrm{\Psi }^{\mathrm{}}`$ on phase space by means of
$$(\mathrm{\Psi },𝒬_{\mathrm{}}^B(f)\mathrm{\Psi })=_{^{2n}}\frac{d^npd^nq}{(2\pi )^n}B_\mathrm{\Psi }^{\mathrm{}}(p,q)f(p,q).$$
(5.28)
Formally, (4.28) and (5.28) immediately yield
$$B_\mathrm{\Psi }^{\mathrm{}}(p,q)=|(\mathrm{\Psi }_{(p,q)}^{\mathrm{}},\mathrm{\Psi })|^2$$
(5.29)
in terms of Schrödinger’s coherent states (4.20). This expression is manifestly positive definite. The existence of $`B_\mathrm{\Psi }^{\mathrm{}}`$ may be proved rigorously by recalling that the Berezin quantization map $`f𝒬_{\mathrm{}}^B(f)`$ is positive from $`C_0(^{2n})`$ to $`(L^2(^n))`$. This implies that for each (unit) vector $`\mathrm{\Psi }L^2(^n)`$ the map $`f(\mathrm{\Psi },𝒬_{\mathrm{}}^B(f)\mathrm{\Psi })`$ is positive from $`C_c(^{2n})`$ to $``$, so that (by the Riesz theorem of measure theory) there must be a measure $`\mu _\mathrm{\Psi }`$ on $`^{2n}`$ such that $`(\mathrm{\Psi },𝒬_{\mathrm{}}^B(f)\mathrm{\Psi })=𝑑\mu _\mathrm{\Psi }f`$. This measure, then, is precisely given by $`d\mu _\mathrm{\Psi }(p,q)=(2\pi )^nd^npd^nqB_\mathrm{\Psi }^{\mathrm{}}(p,q)`$. If $`(\mathrm{\Psi },\mathrm{\Psi })=1`$, then $`\mu _\mathrm{\Psi }`$ is a probability measure. Accordingly, despite its $`\mathrm{}`$-dependence, $`B_\mathrm{\Psi }^{\mathrm{}}`$ defines a bona fide classical probability distribution on phase space, in terms of which one might attempt to visualize quantum mechanics to some extent.
For finite values of $`\mathrm{}`$, the Wigner and Berezin distribution functions are different, because the quantization maps $`𝒬_{\mathrm{}}^W`$ and $`𝒬_{\mathrm{}}^B`$ are. The connection between $`B_\mathrm{\Psi }^{\mathrm{}}`$ and $`W_\mathrm{\Psi }^{\mathrm{}}`$ is easily computed to be
$$B_\mathrm{\Psi }^{\mathrm{}}=W_\mathrm{\Psi }^{\mathrm{}}g^{\mathrm{}},$$
(5.30)
where $`g^{\mathrm{}}`$ is the Gaussian function
$$g^{\mathrm{}}(p,q)=(2/\mathrm{})^n\mathrm{exp}((p^2+q^2)/\mathrm{}).$$
(5.31)
This is how physicists look at the Berezin function,<sup>220</sup><sup>220</sup>220 The ‘Berezin’ functions $`B_\mathrm{\Psi }^{\mathrm{}}`$ were introduced by Husimi (1940) from a different point of view, and are therefore actually called Husimi functions by physicists. viz. as a Wigner function smeared with a Gaussian so as to become positive. But since $`g^{\mathrm{}}`$ converges to a Dirac delta function as $`\mathrm{}0`$ (with respect to the measure $`(2\pi )^nd^npd^nq`$ in the sense of distributions), it is clear from (5.30) that as distributions one has<sup>221</sup><sup>221</sup>221 Eq. (5.32) should be interpreted as a limit of the distribution on $`𝒟(^{2n})`$ or $`𝒮(^{2n})`$ defined by $`B_\mathrm{\Psi }^{\mathrm{}}W_\mathrm{\Psi }^{\mathrm{}}`$. Both functions are continuous for $`\mathrm{}>0`$, but lose this property in the limit $`\mathrm{}0`$, generally converging to distributions.
$$\underset{\mathrm{}0}{lim}\left(B_\mathrm{\Psi }^{\mathrm{}}W_\mathrm{\Psi }^{\mathrm{}}\right)=0.$$
(5.32)
See also (4.32). Hence in the study of the limit $`\mathrm{}0`$ there is little advantage in the use of Wigner functions; quite to the contrary, in limiting procedures their generic lack of positivity makes them more difficult to handle than Berezin functions.<sup>222</sup><sup>222</sup>222 See, however, Robinett (1993) and Arai (1995). It should be mentioned that (5.32) expresses the asymptotic equivalence of Wigner and Berezin functions as distributions on $`\mathrm{}`$-independent test functions. Even in the limit $`\mathrm{}0`$ one is sometimes interested in studying $`O(\mathrm{})`$ phenomena, in which case one should make a choice. For example, one would like to write the asymptotic behaviour (5.8) of coherent states in the form $`lim_\mathrm{}0W_{\mathrm{\Psi }_z^{\mathrm{}}}^{\mathrm{}}=\delta _z`$. Although this is indeed true in the sense of distributions, the corresponding limit
$$\underset{\mathrm{}0}{lim}B_{\mathrm{\Psi }_z^{\mathrm{}}}^{\mathrm{}}=\delta _z,$$
(5.33)
exists in the sense of (probability) measures, and is therefore defined on a much larges class of test functions.<sup>223</sup><sup>223</sup>223Namely those in $`C_0(^{2n})`$ rather than in $`𝒟(^{2n})`$ or $`𝒮(^{2n})`$. Here and in what follows, we abuse notation: if $`\mu ^0`$ is some probability measure on $`^{2n}`$ and $`(\mathrm{\Psi }^{\mathrm{}})`$ is a sequence of unit vectors in $`L^2(^n)`$ indexed by $`\mathrm{}`$ (and perhaps other labels), then $`B_\mathrm{\Psi }^{\mathrm{}}^{\mathrm{}}\mu ^0`$ for $`\mathrm{}0`$ by definition means that for any $`fC_c^{\mathrm{}}(^{2n})`$ one has<sup>224</sup><sup>224</sup>224Since $`𝒬_{\mathrm{}}^B`$ may be extended from $`C_c^{\mathrm{}}(^{2n})`$ to $`L^{\mathrm{}}(^{2n})`$, one may omit the stipulation that $`\mu ^0`$ be a probability measure in this definition if one requires convergence for all $`fL^{\mathrm{}}(^{2n})`$, or just for all $`f`$ in the unitization of the $`C^{}`$-algebra $`C_0(^{2n})`$.
$$\underset{\mathrm{}0}{lim}(\mathrm{\Psi }^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }^{\mathrm{}})=_{^{2n}}𝑑\mu ^0f.$$
(5.34)
### 5.4 The classical limit of energy eigenstates
Having dealt with coherent states in (5.33), in this subsection we discuss the much more difficult problem of computing the limit measure $`\mu ^0`$ for eigenstates of the quantum Hamiltonian $`H`$. Thus we assume that $`H`$ has eigenvalues $`E_n^{\mathrm{}}`$ labeled by $`n`$ (defined with or without 0 according to convenience), and also depending on $`\mathrm{}`$ because of the explicit dependence of $`H`$ on this parameter. The associated eigenstates $`\mathrm{\Psi }_𝗇^{\mathrm{}}`$ then by definition satisfy
$$H\mathrm{\Psi }_𝗇^{\mathrm{}}=E_n^{\mathrm{}}\mathrm{\Psi }_𝗇^{\mathrm{}}.$$
(5.35)
Here we incorporate the possibility that the eigenvalue $`E_n^{\mathrm{}}`$ is degenerate, so that the label $`𝗇`$ extends $`n`$. For example, for the one-dimensional harmonic oscillator one has $`E_n^{\mathrm{}}=\mathrm{}\omega (n+\frac{1}{2})`$ ($`n=0,1,2,\mathrm{}`$) without multiplicity, but for the hydrogen atom the Bohrian eigenvalues $`E_n^{\mathrm{}}=m_ee^4/2\mathrm{}^2n^2`$ (where $`m_e`$ is the mass of the electron and $`e`$ is its charge) are degenerate, with the well-known eigenfunctions $`\mathrm{\Psi }_{(n,l,m)}^{\mathrm{}}`$ (Landau & Lifshitz, 1977). Hence in this case one has $`𝗇=(n,l,m)`$ with $`n=1,2,3,\mathrm{}`$, subject to $`l=0,1,\mathrm{},n1`$, and $`m=l,\mathrm{},l`$.
In any case, it makes sense to let $`𝗇\mathrm{}`$; this certainly means $`n\mathrm{}`$, and may in addition involve sending the other labels in $`𝗇`$ to infinity (subject to the appropriate restrictions on $`𝗇\mathrm{}`$, as above). One then expects classical behaviour à la Bohr if one simultaneously lets $`\mathrm{}0`$ whilst $`E_n^{\mathrm{}}E^0`$ converges to some ‘classical’ value $`E^0`$. Depending on how one lets the possible other labels behave in this limit, this may also involve similar asymptotic conditions on the eigenvalues of operators commuting with $`H`$ \- see below for details in the integrable case. We denote the collection of such eigenvalues (including $`E_n^{\mathrm{}}`$) by $`𝖤_𝗇^{\mathrm{}}`$. (Hence in the case where the energy levels $`E_n^{\mathrm{}}`$ are nondegenerate, the label $`𝖤`$ is just $`E`$.) In general, we denote the collective limit of the eigenvalues $`𝖤_𝗇^{\mathrm{}}`$ as $`\mathrm{}0`$ and $`𝗇\mathrm{}`$ by $`𝖤^0`$.
For example, for the hydrogen atom one has the additional operators $`J^2`$ of total angular momentum as well as the operator $`J_3`$ of angular momentum in the $`z`$-direction. The eigenfunction $`\mathrm{\Psi }_{(n,l,m)}^{\mathrm{}}`$ of $`H`$ with eigenvalue $`E_n^{\mathrm{}}`$ is in addition an eigenfunction of $`J^2`$ with eigenvalue $`j_{\mathrm{}}^2=\mathrm{}^2l(l+1)`$ and of $`J_3`$ with eigenvalue $`j_3^{\mathrm{}}=\mathrm{}m`$. Along with $`n\mathrm{}`$ and $`\mathrm{}0`$, one may then send $`l\mathrm{}`$ and $`m\pm \mathrm{}`$ in such a way that $`j_{\mathrm{}}^2`$ and $`j_3^{\mathrm{}}`$ approach specific constants.
The object of interest, then, is the measure on phase space obtained as the limit of the Berezin functions (5.29), i.e.
$$\mu _𝖤^0=\underset{\mathrm{}0,𝗇\mathrm{}}{lim}B_{\mathrm{\Psi }_𝗇^{\mathrm{}}}^{\mathrm{}}.$$
(5.36)
Although the pioneers of quantum mechanics were undoubtedly interested in quantities like this, it was only in the 1970s that rigorous results were obtained. Two cases are well understood: in this subsection we discuss the integrable case, leaving chaotic and more generally ergodic motion to Subsection 5.6.
In the physics literature, it was argued that for an integrable system the limiting measure $`\mu _𝖤^0`$ is concentrated (in the form of a $`\delta `$-function) on the invariant torus associated to $`𝖤^0`$ (Berry, 1977a).<sup>225</sup><sup>225</sup>225This conclusion was, in fact, reached from the Wigner function formalism. See Ozorio de Almeida (1988) for a review of work of Berry and his collaborators on this subject. Independently, mathematicians began to study a quantity very similar to $`\mu _𝖤^0`$, defined by limiting sequences of eigenfunctions of the Laplacian on a Riemannian manifold $`M`$. Here the underlying classical flow is Hamiltonian as well, the corresponding trajectories being the geodesics of the given metric (see, for example, Klingenberg (1982), Abraham & Marsden (1985), Katok & Hasselblatt (1995), or Landsman (1998)).<sup>226</sup><sup>226</sup>226 The simplest examples of integrable geodesic motion are $`n`$-tori, where the geodesics are projections of lines, and the sphere, where the geodesics are great circles (Katok & Hasselblatt, 1995). The ensuing picture largely confirms the folklore of the physicists:
> In the integrable case the limit measure $`\mu _𝖤^0`$ is concentrated on invariant tori.
See Charbonnel (1986, 1988), Zelditch (1990, 1996a), Toth (1996, 1999), Nadirashvili, Toth, & Yakobson (2001), and Toth & Zelditch (2002, 2003a, 2003b).<sup>227</sup><sup>227</sup>227These papers consider the limit $`n\mathrm{}`$ without $`\mathrm{}0`$; in fact, a physicist would say that they put $`\mathrm{}=1`$. In that case $`E_n\mathrm{}`$; in this procedure the physicists’ microscopic $`E𝒪(\mathrm{})`$ and macroscopic $`E𝒪(1)`$ regimes correspond to $`E𝒪(1)`$ and $`E\mathrm{}`$, respectively. Finally, as part of the transformation of microlocal analysis to semiclassical analysis (cf. Subsection 5.2), these results were adapted to quantum mechanics (Paul & Uribe, 1995, 1996).
Let us now give some details for integrable systems (of Liouville type); these include the hydrogen atom as a special case. Integrable systems are defined by the property that on a $`2p`$-dimensional phase space $`M`$ one has $`p`$ independent<sup>228</sup><sup>228</sup>228I.e. $`df_1\mathrm{}df_p0`$ everywhere. At this point we write $`2p`$ instead of $`2n`$ for the dimension of phase space in order to avoid notational confusion. classical observables $`(f_1=h,f_2,\mathrm{},f_p)`$ whose mutual Poisson brackets all vanish (Arnold, 1989). One then hopes that an appropriate quantization scheme $`𝒬_{\mathrm{}}`$ exists under which the corresponding quantum observables $`(𝒬_{\mathrm{}}(f_1)=H,𝒬_{\mathrm{}}(f_2),\mathrm{},𝒬_{\mathrm{}}(f_p))`$ are all self-adjoint and mutually commute (on a common core).<sup>229</sup><sup>229</sup>229There is no general theory of quantum integrable systems. Olshanetsky & Perelomov (1981, 1983) form a good starting point. This is indeed the case for the hydrogen atom, where $`(f_1,f_2,f_3)`$ may be taken to be $`(h,j^2,j_3)`$ (where $`j^2`$ is the total angular momentum and $`j_3`$ is its $`z`$-component),<sup>230</sup><sup>230</sup>230In fact, if $`\mu `$ is the momentum map for the standard $`SO(3)`$-action on $`^3`$, then $`j^2=_{k=1}^3\mu _k^2`$ and $`j_3=\mu _3`$. $`H`$ is given by (5.16), $`J^2=𝒬_{\mathrm{}}^W(j^2)`$, and $`J_3=𝒬_{\mathrm{}}^W(j_3)`$. In general, the energy eigenfunctions $`\mathrm{\Psi }_𝗇^{\mathrm{}}`$ will be joint eigenfunctions of the operators $`(𝒬_{\mathrm{}}(f_1),\mathrm{},𝒬_{\mathrm{}}(f_p))`$, so that $`𝖤_𝗇^{\mathrm{}}=(E_{n_1}^{\mathrm{}},\mathrm{},E_{n_p}^{\mathrm{}})`$, with $`𝒬_{\mathrm{}}(f_k)\mathrm{\Psi }_𝗇^{\mathrm{}}=E_{n_k}^{\mathrm{}}\mathrm{\Psi }_𝗇^{\mathrm{}}`$. We assume that the submanifolds $`_{k=1}^pf_k^1(x_k)`$ are compact and connected for each $`x^p`$, so that they are tori by the Liouville–Arnold Theorem (Abraham & Marsden, 1985, Arnold, 1989).
Letting $`\mathrm{}0`$ and $`𝗇\mathrm{}`$ so that $`E_{n_k}^{\mathrm{}}E_k^0`$ for some point $`E^0=(E_1^0,\mathrm{},E_p^0)^p`$, it follows that the limiting measure $`\mu _𝖤^0`$ as defined in (5.36) is concentrated on the invariant torus $`_{k=1}^pf_k^1(E_k^0)`$. This torus is generically $`p`$-dimensional, but for singular points $`E^0`$ it may be of lower dimension. In particular, in the exceptional circumstance where the invariant torus is one-dimensional, $`\mu _𝖤^0`$ is concentrated on a classical orbit. Of course, for $`p=1`$ (where any Hamiltonian system is integrable) this singular case is generic. Just think of the foliation of $`^2`$ by the ellipses that form the closed orbits of the harmonic oscillator motion.<sup>231</sup><sup>231</sup>231 It may be enlightening to consider geodesic motion on the sphere; this example may be seen as the hydrogen atom without the radial degree of freedom (so that the degeneracy in question occurs in the hydrogen atom as well). If one sends $`l\mathrm{}`$ and $`m\mathrm{}`$ in the spherical harmonics $`Y_l^m`$ (which are eigenfunctions of the Laplacian on the sphere) in such a way that $`limm/l=\mathrm{cos}\phi `$, then the invariant tori are generically two-dimensional, and occur when $`\mathrm{cos}\phi \pm 1`$; an invariant torus labeled by such a value of $`\phi 0,\pi `$ comprises all great circles (regarded as part of phase space by adding to each point of the geodesic a velocity of unit length and direction tangent to the geodesic) whose angle with the $`z`$-axis is $`\phi `$ (more precisely, the angle in question is the one between the normal of the plane through the given great circle and the $`z`$-axis). For $`\mathrm{cos}\phi =\pm 1`$ (i.e. $`m=\pm l`$), however, there is only one great circle with $`\phi =0`$ namely the equator (the case $`\phi =\pi `$ corresponds to the same equator traversed in the opposite direction). Hence in this case the invariant torus is one-dimensional. The reader may be surprised that the invariant tori explicitly depend on the choice of variables, but this feature is typical of so-called degenerate systems; see Arnold (1989), §51.
What remains, then, of Bohr’s picture of the hydrogen atom in this light?<sup>232</sup><sup>232</sup>232We ignore coupling to the electromagnetic field here; see footnote 28. Quite a lot, in fact, confirming his remarkable physical intuition. The energy levels Bohr calculated are those given by the Schrödinger equation, and hence remain correct in mature quantum mechanics. His orbits make literal sense only in the “correspondence principle” limit $`\mathrm{}0`$, $`n\mathrm{}`$, where, however, the situation is even better than one might expect for integrable systems: because of the high degree of symmetry of the Kepler problem (Guillemin & Sternberg, 1990), one may construct energy eigenfunctions whose limit measure $`\mu ^0`$ concentrates on any desired classical orbit (Nauenberg, 1989).<sup>233</sup><sup>233</sup>233Continuing footnote 231, for a given principal quantum number $`n`$ one forms the eigenfunction $`\mathrm{\Psi }_{(n,n1,n1)}^{\mathrm{}}`$ by multiplying the spherical harmonic $`Y_{n1}^{n1}`$ with the appropriate radial wave function. The limiting measure (5.36) as $`n\mathrm{}`$ and $`\mathrm{}0`$ is then concentrated on an orbit (rather than on an invariant torus). Now, beyond what it possible for general integrable systems, one may use the $`SO(4)`$ symmetry of the Kepler problem and the construction in footnote 205 for the group-theoretic coherent states of Perelomov (1986) to find the desired eigenfunctions. See also De Bièvre (1992) and De Bièvre et al. (1993). In order to recover a travelling wave packet, one has to form wave packets from a very large number of energy eigenstates with very high quantum numbers, as explained in Subsection 2.4. For finite $`n`$ and $`\mathrm{}`$ Bohr’s orbits seem to have no meaning, as already recognized by Heisenberg (1969) in his pathfinder days!<sup>234</sup><sup>234</sup>234The later Bohr also conceded this through his idea that causal descriptions are complementary to space-time pictures; see Subsection 3.3.
### 5.5 The WKB approximation
One might have expected a section on the $`\mathrm{}0`$ limit of quantum mechanics to be centered around the WKB approximation, as practically all textbooks base their discussion of the classical limit on this notion. Although the scope of this method is actually rather limited, it is indeed worth saying a few words about it. For simplicity we restrict ourselves to the time-independent case.<sup>235</sup><sup>235</sup>235Cf. Robert (1998) and references therein for the time-dependent case. In its original formulation, the time-independent WKB method involves an attempt to approximate solutions of the time-independent Schrödinger equation $`H\mathrm{\Psi }=E\mathrm{\Psi }`$ by wave functions of the type
$$\mathrm{\Psi }(x)=a_{\mathrm{}}(x)e^{\frac{i}{\mathrm{}}S(x)},$$
(5.37)
where $`a_{\mathrm{}}`$ admits an expansion in $`\mathrm{}`$ as a power series. Assuming the Hamiltonian $`H`$ is of the form (5.15), plugging the Ansatz (5.37) into the Schrödinger equation, and expanding in $`\mathrm{}`$, yields in lowest order the classical (time-independent) Hamilton–Jacobi equation
$$h(\frac{S}{x},x)=E,$$
(5.38)
supplemented by the so-called (homogeneous) transport equation<sup>236</sup><sup>236</sup>236 Only stated here for a classical Hamiltonian $`h(p,q)=p^2/2m+V(q)`$. Higher-order terms in $`\mathrm{}`$ yield further, inhomogeneous transport equations for the expansion coefficients $`a_j(x)`$ in $`a_{\mathrm{}}=_ja_j\mathrm{}^j`$. These can be solved in a recursive way, starting with (5.39).
$$\left(\frac{1}{2}\mathrm{\Delta }S+\underset{k}{}\frac{S}{x^k}\frac{}{x^k}\right)a_0=0.$$
(5.39)
In particular, $`E`$ should be a classically allowed value of the energy. Even when it applies (see below), in most cases of interest the Ansatz (5.37) is only valid locally (in $`x`$), leading to problems with caustics. These problems turn out to be an artefact of the use of the coordinate representation that lies behind the choice of the Hilbert space $`=L^2(^n)`$, and can be avoided (Maslov & Fedoriuk, 1981): the WKB method really comes to its own in a geometric reformulation in terms of symplectic geometry. See Arnold (1989), Bates & Weinstein (1995), and Dimassi & Sjöstrand (1999) for (nicely complementary) introductory treatments, and Guillemin & Sternberg (1977), Hörmander (1985a, 1985b), and Duistermaat (1974, 1996) for advanced accounts.
The basic observation leading to this reformulation is that in the rare cases that $`S`$ is defined globally as a smooth function on the configuration space $`^n`$, it defines a submanifold $``$ of the phase space $`M=^{2n}`$ by $`=\{(p=dS(x),q=x),x^n\}`$. This submanifold is Lagrangian in having two defining properties: firstly, $``$ is $`n`$-dimensional, and secondly, the restriction of the symplectic form (i.e. $`_kdp_kdq^k`$) to $``$ vanishes. The Hamilton–Jacobi equation (5.38) then guarantees that the Lagrangian submanifold $`M`$ is contained in the surface $`\mathrm{\Sigma }_E=h^1(E)`$ of constant energy $`E`$ in $`M`$. Consequently, any solution of the Hamiltonian equations of motion that starts in $``$ remains in $``$.
In general, then, the starting point of the WKB approximation is a Lagrangian submanifold $`\mathrm{\Sigma }_EM`$, rather than some function $`S`$ that defines it locally. By a certain adaptation of the geometric quantization procedure, one may, under suitable conditions, associate a unit vector $`\mathrm{\Psi }_{}`$ in a suitable Hilbert space to $``$, which for small $`\mathrm{}`$ happens to be a good approximation to an eigenfunction of $`H`$ at eigenvalue $`E`$. This strategy is successful in the integrable case, where the nondegenerate tori (i.e. those of maximal dimension $`n`$) provide such Lagrangian submanifolds of $`M`$; the associated unit vector $`\mathrm{\Psi }_{}`$ then turns out to be well defined precisely when $``$ satisfies (generalized) Bohr–Sommerfeld quantization conditions. In fact, this is how the measures $`\mu _𝖤^0`$ in (5.36) are generally computed in the integrable case.
If the underlying classical system is not integrable, it may still be close enough to integrability for invariant tori to be defined. Such systems are called quasi-integrable or perturbations of integrable systems, and are described by the Kolmogorov–Arnold–Moser (KAM) theory; see Gallavotti (1983), Abraham & Marsden (1985), Ozorio de Almeida (1988), Arnold (1989), Lazutkin (1993), Gallavotti, Bonetto & Gentile (2004), and many other books. In such systems the WKB method continues to provide approximations to the energy eigenstates relevant to the surviving invariant tori (Colin de Verdière, 1977; Lazutkin, 1993; Popov, 2000), but already loses some of its appeal.
In general systems, notably chaotic ones, the WKB method is almost useless. Indeed, the following theorem of Werner (1995) shows that the measure $`\mu _𝖤^0`$ defined by a WKB function (5.37) is concentrated on the Lagrangian submanifold $``$ defined by $`S`$:
> Let $`a_{\mathrm{}}`$ be in $`L^2(^n)`$ for each $`\mathrm{}>0`$ with pointwise limit $`a_0=lim_\mathrm{}0a_{\mathrm{}}`$ also in $`L^2(^n)`$,<sup>237</sup><sup>237</sup>237This assumption is not made in Werner (1995), who directly assumes that $`\mathrm{\Psi }=a_0\mathrm{exp}(iS/\mathrm{})`$ in (5.37). and suppose that $`S`$ is almost everywhere differentiable. Then for each $`fC_c^{\mathrm{}}(^{2n})`$:
$$\underset{\mathrm{}0}{lim}(a_{\mathrm{}}e^{\frac{i}{\mathrm{}}S},𝒬_{\mathrm{}}^B(f)a_{\mathrm{}}e^{\frac{i}{\mathrm{}}S})=_^nd^nx|a_0(x)|^2f(\frac{S}{x},x).$$
(5.40)
As we shall see shortly, this behaviour is impossible for ergodic systems, and this is enough to seal the fate of WKB for chaotic systems in general (except perhaps as a hacker’s tool).
Note, however, that for a given energy level $`E`$ the discussion so far has been concerned with properties of the classical trajectories on $`\mathrm{\Sigma }_E`$ (where they are constrained to remain by conservation of energy). Now, it belongs to the essence of quantum mechanics that other parts of phase space than $`\mathrm{\Sigma }_E`$ might be relevant to the spectral properties of $`H`$ as well. For example, for a classical Hamiltonian of the simple form $`h(p,q)=p^2/2m+V(q)`$, this concerns the so-called classically forbidden area $`\{q^nV(q)>E\}`$ (and any value of $`p`$). Here the classical motion can have no properties like integrability or ergodicity, because it does not exist. Nonetheless, and perhaps counterintuitively, it is precisely here that a slight adaptation of the WKB method tends to be most effective. For $`q=x`$ in the classically forbidden area, the Ansatz (5.37) should be replaced by
$$\mathrm{\Psi }(x)=a_{\mathrm{}}(x)e^{\frac{S(x)}{\mathrm{}}},$$
(5.41)
where this time $`S`$ obeys the Hamilton–Jacobi equation ‘for imaginary time’, <sup>238</sup><sup>238</sup>238This terminology comes from the Lagrangian formalism, where the classical action $`S=𝑑tL(t)`$ is replaced by $`iS`$ through the substitution $`t=i\tau `$ with $`\tau `$. i.e.
$$h(i\frac{S}{x},x)=E,$$
(5.42)
and the transport equation (5.39) is unchanged. For example, it follows that in one dimension (with a Hamiltonian of the type (2.3)) the WKB function (5.41) assumes the form
$$\mathrm{\Psi }(x)e^{\frac{\sqrt{2m}}{\mathrm{}}^{|x|}𝑑y\sqrt{V(y)E}}$$
(5.43)
in the forbidden region, which explains both the tunnel effect in quantum mechanics (i.e. the propagation of the wave function into the forbidden region) and the fact that this effect disappears in the limit $`\mathrm{}0`$. However, even here the use of WKB methods has now largely been superseded by techniques developed by Agmon (1982); see, for example, Hislop & Sigal (1996) and Dimassi & Sjöstrand (1999) for reviews.
### 5.6 Epilogue: quantum chaos
Chaos in classical mechanics was probably known to Newton and was famously highlighted by Poincaré (1892–1899),<sup>239</sup><sup>239</sup>239See also Diacu & Holmes (1996) and Barrow-Green (1997) for historical background. but its relevance for (and potential threat to) quantum theory was apparently first recognized by Einstein (1917) in a paper that was ‘completely ignored for 40 years’ (Gutzwiller, 1992).<sup>240</sup><sup>240</sup>240It was the study of the very same Helium atom that led Heisenberg to believe that a fundamentally new ‘quantum’ mechanics was needed to replace the inadequate old quantum theory of Bohr and Sommerfeld. See Mehra & and Rechenberg (1982b) and Cassidy (1992). Another microscopic example of a chaotic system is the hydrogen atom in an external magnetic field. Currently, the study of quantum chaos is one of the most thriving businesses in all of physics, as exemplified by innumerable conference proceedings and monographs on the subject, ranging from the classic by Gutzwiller (1990) to the online opus magnum by Cvitanovic et al. (2005).<sup>241</sup><sup>241</sup>241Other respectable books include, for example, Guhr, Müller-Groeling & Weidenmüller (1998), Haake (2001) and Reichl (2004). Nonetheless, the subject is still not completely understood, and provides a fascinating testing ground for the interplay between classical and quantum mechanics.
One should distinguish between various different goals in the field of quantum chaos. The majority of papers and books on quantum chaos is concerned with the semiclassical analysis of some concretely given quantum system having a chaotic system as its classical limit. This means that one tries to approximate (for small $`\mathrm{}`$) a suitable quantum-mechanical expression in terms of data associated with the underlying classical motion. Michael Berry even described this goal as the “Holy Grail” of quantum chaos. The methods described in Subsection 5.2 contribute to this goal, but are largely independent of the nature of the dynamics. In this subsection we therefore concentrate on techniques and results specific to chaotic motion.
Historically, the first new tool in semiclassical approximation theory that specifically applied to chaotic systems was the so-called Gutzwiller trace formula.<sup>242</sup><sup>242</sup>242This attribution is based on Gutzwiller (1971). A similar result was independently derived by Balian & Bloch (1972, 1974). See also Gutzwiller (1990) and Brack & Bhaduri (2003) for mathematically heuristic but otherwise excellent accounts of semiclassical physics based on the trace formula. Mathematically rigorous discussions and proofs may be found in Colin de Verdière (1973), Duistermaat & Guillemin (1975), Guillemin & Uribe (1989), Paul & Uribe (1995), and Combescure, Ralston, & Robert (1999). Roughly speaking, this formula approximates the eigenvalues of the quantum Hamiltonian in terms of the periodic (i.e. closed) orbits of the underlying classical Hamiltonian.<sup>243</sup><sup>243</sup>243Such orbits are dense but of Liouville measure zero in chaotic classical systems. Their crucial role was first recognized by Poincaré (1892–1899). The Gutzwiller trace formula does not start from the wave function (as the WKB approximation does), but from the propagator $`K(x,y,t)`$. Physicists write this as $`K(x,y,t)=x|\mathrm{exp}(itH/\mathrm{})|y`$, whereas mathematicians see it as the Green’s function in the formula
$$e^{\frac{it}{\mathrm{}}H}\mathrm{\Psi }(x)=d^nyK(x,y,t)\mathrm{\Psi }(y),$$
(5.44)
where $`\mathrm{\Psi }L^2(^n)`$. Its (distributional) Laplace transform
$$G(x,y,E)=\frac{1}{i\mathrm{}}_0^{\mathrm{}}𝑑tK(x,y,t)e^{\frac{itE}{\mathrm{}}}$$
(5.45)
contains information about both the spectrum and the eigenfunctions; for if the former is discrete, one has
$$G(x,y,E)=\underset{j}{}\frac{\mathrm{\Psi }_j(x)\overline{\mathrm{\Psi }_j(y)}}{EE_j}.$$
(5.46)
It is possible to approximate $`K`$ or $`G`$ itself by an expression of the type
$$K(x,y,t)(2\pi i\mathrm{})^{n/2}\underset{P}{}\sqrt{|detV_P|}e^{\frac{i}{\mathrm{}}S_P(x,y,t){\scriptscriptstyle \frac{1}{2}}i\pi \mu _P},$$
(5.47)
where the sum is over all classical paths $`P`$ from $`y`$ to $`x`$ in time $`t`$ (i.e. paths that solve the classical equations of motion). Such a path has an associated action $`S_P`$, Maslov index $`\mu _P`$, and Van Vleck (1928) determinant $`detV_P`$ (Arnold, 1989). For chaotic systems one typically has to include tens of thousands of paths in the sum, but if one does so the ensuing approximation turns out to be remarkably successful (Heller & Tomsovic, 1993; Tomsovic & Heller, 1993). The Gutzwiller trace formula is a semiclassical approximation to
$$g(E)=d^nxG(x,x,E)=\underset{j}{}\frac{1}{EE_j},$$
(5.48)
for a quantum Hamiltonian with discrete spectrum and underlying classical Hamiltonian having chaotic motion. It has the form
$$g(E)g_0(E)+\frac{1}{i\mathrm{}}\underset{P}{}\underset{k=1}{\overset{\mathrm{}}{}}\frac{T_P}{2\mathrm{sinh}(k\chi _P/2)}e^{\frac{ik}{\mathrm{}}S_P(E){\scriptscriptstyle \frac{1}{2}}i\pi \mu _P},$$
(5.49)
where $`g_0`$ is a smooth function giving the mean density of states. This time, the sum is over all (prime) periodic paths $`P`$ of the classical Hamiltonian at energy $`E`$ with associated action $`S_P(E)=p𝑑q`$ (where the momentum $`p`$ is determined by $`P`$, given $`E`$), period $`T_P`$, and stability exponent $`\chi _P`$ (this is a measure of how rapidly neighbouring trajectories drift away from $`P`$). Since the frustration expressed by Einstein (1917), this was the first indication that semiclassical approximations had some bearing on chaotic systems.
Another important development concerning energy levels was the formulation of two key conjectures:<sup>244</sup><sup>244</sup>244Strictly speaking, both conjectures are wrong; for example, the harmonic oscillator yields a counterexamples to the first one. See Zelditch (1996a) for further information. Nonetheless, the conjectures are believed to be true in a deeper sense.
* If the classical dynamics defined by the classical Hamiltonian $`h`$ is integrable, then the spectrum of $`H`$ is “uncorrelated” or “random” (Berry & Tabor, 1977).
* If the classical dynamics defined by $`h`$ is chaotic, then the spectrum of $`H`$ is “correlated” or “regular” (Bohigas, Giannoni, & Schmit, 1984).
The notions of correlation and randomness used here can be made precise using notions like the distribution of level spacings and the pair correlation function of eigenvalues; see Zelditch (1996a) and De Bièvre (2001) for introductory treatments, and most of the literature cited in this subsection for further details.<sup>245</sup><sup>245</sup>245This aspect of quantum chaos has applications to number theory and might even lead to a proof of the Riemann hypothesis; see, for example, Sarnak (1999), Berry & Keating (1999), and many other recent papers. Another relevant connection, related to the one just mentioned, is between energy levels and random matrices; see especially Guhr, Müller-Groeling & Weidenmüller (1998). For the plain relevance of all this to practical physics see Mirlin (2000).
We now consider energy eigenfunctions instead of eigenvalues, and return to the limit measure (5.36). In the non (quasi-) integrable case, the key result is that
> for ergodic classical motion,<sup>246</sup><sup>246</sup>246Ergodicity is the weakest property that any chaotic dynamical system possesses. See Katok & Hasselblatt (1995), Emch & Liu (2002), Gallavotti, Bonetto & Gentile (2004), and countless other books. the limit measure $`\mu _𝖤^0`$ coincides with the (normalized) Liouville measure induced on the constant energy surface $`\mathrm{\Sigma }_Eh^1(E)`$.<sup>247</sup><sup>247</sup>247The unnormalized Liouville measure $`\mu _E^u`$ on $`\mathrm{\Sigma }_E`$ is defined by $`\mu _E^u(B)=_B𝑑S_E(x)(dh(x))^1`$, where $`dS_E`$ is the surface element on $`\mathrm{\Sigma }_E`$ and $`B\mathrm{\Sigma }_E`$ is a Borel set. If $`\mathrm{\Sigma }_E`$ is compact, the normalized Liouville measure $`\mu _E`$ on $`\mathrm{\Sigma }_E`$ is given by $`\mu _E(B)=\mu _E^u(B)/\mu _E^u(\mathrm{\Sigma }_E)`$. It is a probability measure on $`\mathrm{\Sigma }_E`$, reflecting the fact that the eigenvectors $`\mathrm{\Psi }_𝗇^{\mathrm{}}`$ are normalized to unit length so as to define quantum-mechanical states.
This result was first suggested in the mathematical literature for ergodic geodetic motion on compact hyperbolic Riemannian manifolds (Snirelman, 1974), where it was subsequently proved with increasing generality (Colin de Verdière, 1985; Zelditch, 1987).<sup>248</sup><sup>248</sup>248In the Riemannian case with $`\mathrm{}=1`$ the cosphere bundle $`S^{}Q`$ (i.e. the subbundle of the cotangent bundle $`T^{}Q`$ consisting of one-forms of unit length) plays the role of $`\mathrm{\Sigma }_E`$. Low-dimensional examples of ergodic geodesic motion are provided by compact hyperbolic spaces. Also cf. Zelditch (1992a) for the physically important case of a particle moving in an external gauge field. See also the appendix to Lazutkin (1993) by A.I. Shnirelman, and Nadirashvili, Toth, & Yakobson (2001) for reviews. For certain other ergodic systems this property was proved by Zelditch (1991), Gérard & Leichtnam (1993), Zelditch & Zworski (1996), and others; to the best of our knowledge a completely general proof remains to be given.
An analogous version for Schrödinger operators on $`^n`$ was independently stated in the physics literature (Berry, 1977b, Voros, 1979), and was eventually proved under certain assumptions on the potential by Helffer, Martinez & Robert (1987), Charbonnel (1992), and Paul & Uribe (1995). Under suitable assumptions one therefore has
$$\underset{\mathrm{}0,𝗇\mathrm{}}{lim}(\mathrm{\Psi }_𝗇^{\mathrm{}},𝒬_{\mathrm{}}^B(f)\mathrm{\Psi }_𝗇^{\mathrm{}})=_{\mathrm{\Sigma }_E}𝑑\mu _Ef$$
(5.50)
for any $`fC_c^{\mathrm{}}(^{2n})`$, where again $`\mu _E`$ is the (normalized) Liouville measure on $`\mathrm{\Sigma }_E^{2n}`$ (assuming this space to be compact). In particular, in the ergodic case $`\mu _𝖤^0`$ only depends on $`E^0`$ and is the same for (almost) every sequence of energy eigenfunctions $`(\mathrm{\Psi }_𝗇^{\mathrm{}})`$ as long as $`E_n^{\mathrm{}}E^0`$.<sup>249</sup><sup>249</sup>249 The result is not necessarily valid for all sequences $`(\mathrm{\Psi }_𝗇^{\mathrm{}})`$ with the given limiting behaviour, but only for ‘almost all’ such sequences (technically, for a class of sequences of density 1). See, for example, De Bièvre (2001) for a simple explanation of this. Thus the support of the limiting measure is uniformly spread out over the largest part of phase space that is dynamically possible.
The result that for ergodic classical motion $`\mu _𝖤^0`$ is the Liouville measure on $`\mathrm{\Sigma }_E`$ under the stated condition leaves room for the phenomenon of ‘scars’, according to which in chaotic systems the limiting measure is sometimes concentrated on periodic classical orbits. This terminology is used in two somewhat different ways in the literature. ‘Strong’ scars survive in the limit $`\mathrm{}0`$ and concentrate on stable closed orbits;<sup>250</sup><sup>250</sup>250An orbit $`\gamma M`$ is called stable when for each neighbourhood $`U`$ of $`\gamma `$ there is neighbourhood $`VU`$ of $`\gamma `$ such that $`z(t)U`$ for all $`zV`$ and all $`t`$. they may come from ‘exceptional’ sequences of eigenfunctions.<sup>251</sup><sup>251</sup>251Cf. footnote 249. These are mainly considered in the mathematical literature; cf. Nadirashvili, Toth, & Yakobson (2001) and references therein.
In the physics literature, on the other hand, the notion of a scar usually refers to an anomalous concentration of the functions $`B_{\mathrm{\Psi }_𝗇^{\mathrm{}}}^{\mathrm{}}`$ (cf. (5.29)) near unstable closed orbits for finite values of $`\mathrm{}`$; see Heller & Tomsovic (1993), Tomsovic & Heller (1993), Kaplan & Heller (1998a,b), and Kaplan (1999) for surveys. Such scars turn out to be crucial in attempts to explain the energy spectrum of the associated quantum system. The reason why such scars do not survive the (double) limit in (5.36) is that this limit is defined with respect to $`\mathrm{}`$-independent smooth test functions. Physically, this means that one averages over more and more De Broglie wavelengths as $`\mathrm{}0`$, eventually losing information about the single wavelength scale (Kaplan, 1999). Hence to pick them up in a mathematically sound way, one should redefine (5.36) as a pointwise limit (Duclos & Hogreve, 1993, Paul & Uribe, 1996, 1998). In any case, there is no contradiction between the mathematical results cited and what physicists have found.
Another goal of quantum chaos is the identification of chaotic phenomena within a given quantum-mechanical model. Here the slight complication arises that one cannot simply copy the classical definition of chaos in terms of diverging trajectories in phase space, since (by unitarity of time-evolution) in quantum mechanics $`\mathrm{\Psi }(t)\mathrm{\Phi }(t)`$ is constant in time $`t`$ for solutions of the Schrödinger equation. However, this just indicates that should intrinsic quantum chaos exist, it has to be defined differently from classical chaos.<sup>252</sup><sup>252</sup>252As pointed out by Belot & Earman (1997), the Koopman formulation of classical mechanics (cf. footnote 173) excludes classical chaos if this is formulated in terms of trajectories in Hilbert space. The transition from classical to quantum notions of chaos can be smoothened by first reformulating the classical definition of chaos (normally put in terms of properties of trajectories in phase space). This has now been largely accomplished in the algebraic formulation of quantum theory (Benatti, 1993; Emch et al., 1994;, Zelditch, 1996b,c; Belot & Earman, 1997; Alicki & Fannes, 2001; Narnhofer, 2001). The most significant recent development in this direction in the “heuristic” literature has been the study of the quantity
$$M(t)=|(e^{\frac{it}{\mathrm{}}(H+\mathrm{\Sigma })}\mathrm{\Psi },e^{\frac{it}{\mathrm{}}H}\mathrm{\Psi })|^2,$$
(5.51)
where $`\mathrm{\Psi }`$ is a coherent state (or Gaussian wave packet), and $`\mathrm{\Sigma }`$ is some perturbation of the Hamiltonian $`H`$ (Peres, 1984). In what is generally regarded as a breakthrough in the field, Jalabert & Pastawski (2001) discovered that in a certain regime $`M(t)`$ is independent of the detailed form of $`\mathrm{\Sigma }`$ and decays as $`\mathrm{exp}(\lambda t)`$, where $`\lambda `$ is the (largest) Lyapunov exponent of the underlying classical system. See Cucchietti (2004) for a detailed account and further development.
In any case, the possibility that classical chaos appears in the $`\mathrm{}0`$ limit of quantum mechanics is by no means predicated on the existence of intrinsic quantum chaos in the above sense.<sup>253</sup><sup>253</sup>253Arguments by Ford (1988) and others to the effect that quantum mechanics is wrong because it cannot give rise to chaos in its classical limit have to be discarded for the reasons given here. See also Belot & Earman (1997). In fact, using the same argument, such authors could simultaneously have ‘proved’ the opposite statement that any classical dynamics that arises as the classical limit of a quantum theory with non-degenerate spectrum must be ergodic. For the naive definition of quantum ergodic flow clearly is that quantum time-evolution sweeps out all states at some energy $`E`$; but for non-degenerate spectra this is a tautology by definition of an eigenfunction! For even in the unlikely case that quantum dynamics would turn out to be intrinsically non-chaotic, its classical limit is sufficiently singular to admit kinds of classical motion without a qualitative counterpart in quantum theory. This possibility is not only confirmed by most of the literature on quantum chaos (little of which makes any use of notions of intrinsic quantum chaotic motion), but even more so by the possibility of incomplete motion. This is a type of dynamics in which the flow of the Hamiltonian vector field is only defined until a certain time $`t_f<\mathrm{}`$ (or from an initial time $`t_i>\mathrm{}`$), which means that the equations of motion have no solution for $`t>t_f`$ (or $`t<t_i`$).<sup>254</sup><sup>254</sup>254 The simplest examples are incomplete Riemannian manifolds $`Q`$ with geodesic flow; within this class, the case $`Q=(0,1)`$ with flat metric is hard to match in simplicity. Clearly, the particle reaches one of the two boundary points in finite time, and does not know what to do (or even whether its exists) afterwards. Other examples come from potentials $`V`$ on $`Q=^n`$ with the property that the classical dynamics is incomplete; see Reed & Simon (1975) and Gallavotti (1983). On a somewhat different note, the Universe itself has incomplete dynamics because of the Big Bang and possible Big Crunch. The point, then, is that unitary quantum dynamics, though intrinsically complete, may very well have incomplete motion as its classical limit.<sup>255</sup><sup>255</sup>255 The quantization of the Universe is unknown at present, but geodesic motion on Riemannian manifolds, complete or not, is quantized by $`H=\frac{\mathrm{}^2}{2m}\mathrm{\Delta }`$ (perhaps with an additonal term proportional to the Ricci scalar $`R`$, see Landsman (1998)), where $`\mathrm{\Delta }`$ is the Laplacian, and quantization on $`Q=^n`$ is given by the Schrödinger equation (2.3), whether or not the classical dynamics is complete. In these two cases, and probably more generally, the incompleteness of the classical motion is often (but not always) reflected by the lack of essential self-adjointness of the quantum Hamiltonian on its natural initial domain $`C_c^{\mathrm{}}(Q)`$. For example, if $`Q`$ is complete as a Riemannian manifold, then $`\mathrm{\Delta }`$ is essentially self-adjoint on $`C_c^{\mathrm{}}(Q)`$ (Chernoff, 1973, Strichartz, 1983), and if $`Q`$ is incomplete then the Laplacian usually fails to be essentially self-adjoint on this domain (but see Horowitz & Marolf (1995) for counterexamples). One may refer to the latter property as quantum-mechanical incompleteness (Reed & Simon, 1975), although a Hamiltonian that fails to be essentially self-adjoint on $`C_c^{\mathrm{}}(Q)`$ can often be extended (necessarily in a non-unique way) to a self-adjoint operator by a choice of boundary conditions (possibly at infinity). By Stone’s theorem, the quantum dynamics defined by each self-adjoint extension is unitary (and therefore defined for all times). Similarly, although no general statement can be made relating (in)complete classical motion in a potential to (lack of) essential selfadjointness of the corresponding Schrödinger operator, it is usually the case that completeness implies essential selfadjointness, and vice versa. See Reed & Simon (1975), Appendix to §X.1, where the reader may also find examples of classically incomplete but quantum-mechanically complete motion, and vice versa. Now, here is the central point for the present discussion: as probably first noted by Hepp (1974), different self-adjoint extensions have the same classical limit (in the sense of (5.20) or similar criteria), namely the given incomplete classical dynamics. This proves that complete quantum dynamics can have incomplete motion as its classical limit. However, much remains to be understood in this area. See also Earman (2005, 2006).
## 6 The limit $`N\mathrm{}`$
In this section we show to what extent classical physics may approximately emerge from quantum theory when the size of a system becomes large. Strictly classical behaviour would be an idealization reserved for the limit where this size is infinite, which we symbolically denote by “$`limN\mathrm{}`$”. As we shall see, mathematically speaking this limit is a special case of the limit $`\mathrm{}0`$ discussed in the previous chapter. What is more, we shall show that formally the limit $`N\mathrm{}`$ even falls under the heading of continuous fields of $`C^{}`$-algebras and deformation quantization (see Subsection 4.3.) Thus the ‘philosophical’ nature of the idealization involved in assuming that a system is infinite is much the same as that of assuming $`\mathrm{}0`$ in a quantum system of given (finite) size; in particular, the introductory comments in Section 1 apply here as well.
An analogous discussion pertains to the derivation of thermodynamics from statistical mechanics (Emch & Liu, 2002; Batterman, 2005). For example, in theory phase transitions only occur in infinite systems, but in practice one sees them every day. Thus it appears to be valid to approximate a pot of $`10^{23}`$ boiling water molecules by an infinite number of such molecules. The basic point is that the distinction between microscopic and macroscopic regimes is unsharp unless one admits infinite systems as an idealization, so that one can simply say that microscopic systems are finite, whereas macroscopic systems are infinite. This procedure is eventually justified by the results it produces.
Similarly, in the context of quantum theory classical behaviour is simply not found in finite systems (when $`\mathrm{}>0`$ is fixed), whereas, as we shall see, it is found in infinite ones. Given the observed classical nature of the macroscopic world,<sup>256</sup><sup>256</sup>256With the well-known mesoscopic exceptions (Leggett, 2002; Brezger et al., 2002; Chiorescu et al., 2003; Marshall et al., 2003; Devoret et al., 2004). at the end of the day one concludes that the idealization in question is apparently a valid one. One should not be confused by the fact that the error in the number of particles this approximation involves (viz. $`\mathrm{}10^{23}=\mathrm{}`$) is considerably larger than the number of particles in the actual system. If all of the $`10^{23}`$ particles in question were individually tracked down, the approximation is indeed a worthless ones, but the point is rather that the limit $`N\mathrm{}`$ is valid whenever averaging over $`N=10^{23}`$ particles is well approximated by averaging over an arbitrarily larger number $`N`$ (which, then, one might as well let go to infinity). Below we shall give a precise version of this argument.
Despite our opening comments above, the quantum theory of infinite systems has features of its own that deserve a separate section. Our treatment is complementary to texts such as Thirring (1983), Strocchi (1985), Bratteli & Robinson (1987), Haag (1992), Araki (1999), and Sewell (1986, 2002), which should be consulted for further information on infinite quantum systems. The theory in Subsections 6.1 and 6.5 is a reformulation in terms of continuous field of $`C^{}`$-algebras and deformation quantization of the more elementary parts of a remarkable series of papers on so-called quantum mean-field systems by Raggio & Werner (1989, 1991), Duffield & Werner (1992a,b,c), and Duffield, Roos, & Werner (1992). These models have their origin in the treatment of the BCS theory of superconductivity due to Bogoliubov (1958) and Haag (1962), with important further contributions by Thirring & Wehrl (1967), Thirring (1968), Hepp (1972), Hepp & Lieb (1973), Rieckers (1984), Morchio & Strocchi (1987), Duffner & Rieckers (1988), Bona (1988, 1989, 2000), Unnerstall (1990a, 1990b), Bagarello & Morchio (1992), Sewell (2002), and others.
### 6.1 Macroscopic observables
The large quantum systems we are going to study consist of $`N`$ copies of a single quantum system with unital algebra of observables $`𝒜_1`$. Almost all features already emerge in the simplest example $`𝒜_1=M_2()`$ (i.e. the complex $`2\times 2`$ matrices), so there is nothing wrong with having this case in mind as abstraction increases.<sup>257</sup><sup>257</sup>257In the opposite direction of greater generality, it is worth noting that the setting below actually incorporates quantum systems defined on general lattices in $`^n`$ (such as $`^n`$). For one could relabel things so as to make $`𝒜_{1/N}`$ below the algebra of observables of all lattice points $`\mathrm{\Lambda }`$ contained in, say, a sphere of radius $`N`$. The limit $`N\mathrm{}`$ then corresponds to the limit $`\mathrm{\Lambda }^n`$. The aim of what follows is to describe in what precise sense macroscopic observables (i.e. those obtained by averaging over an infinite number of sites) are “classical”.
From the single $`C^{}`$-algebra $`𝒜_1`$, we construct a continuous field of $`C^{}`$-algebras $`𝒜^{(c)}`$ over
$$I=01/=\{0,\mathrm{},1/N,\mathrm{},\frac{1}{3},\frac{1}{2},1\}[0,1],$$
(6.1)
as follows. We put
$`𝒜_0^{(c)}`$ $`=`$ $`C(𝒮(𝒜_1));`$
$`𝒜_{1/N}^{(c)}`$ $`=`$ $`𝒜_1^N,`$ (6.2)
where $`𝒮(𝒜_1)`$ is the state space of $`𝒜_1`$ (equipped with the weak$`^{}`$-topology)<sup>258</sup><sup>258</sup>258In this topology one has $`\omega _\lambda \omega `$ when $`\omega _\lambda (A)\omega (A)`$ for each $`A𝒜_1`$. and $`𝒜_1^N=\widehat{}^N𝒜_1`$ is the (spatial) tensor product of $`N`$ copies of $`𝒜_1`$.<sup>259</sup><sup>259</sup>259When $`𝒜_1`$ is finite-dimensional the tensor product is unique. In general, one needs the projective tensor product at this point. See footnote 91. The point is the same here: any tensor product state $`\omega _1\mathrm{}\omega _N`$ on $`^N𝒜_1`$ \- defined on elementary tensors by $`\omega _1\mathrm{}\omega _N(A_1\mathrm{}A_N)=\omega _1(A_1)\mathrm{}\omega _N(A_N)`$ \- extends to a state on $`\widehat{}^N𝒜_1`$ by continuity. This explains the suffix $`c`$ in $`𝒜^{(c)}`$: it refers to the fact that the limit algebra $`𝒜_0^{(c)}`$ is classical or commutative.
For example, take $`𝒜_1=M_2()`$. Each state is given by a density matrix, which is of the form
$$\rho (x,y,z)=\frac{1}{2}\left(\begin{array}{cc}1+z& xiy\\ x+iy& 1z\end{array}\right),$$
(6.3)
for some $`(x,y,z)^3`$ satisfying $`x^2+y^2+z^21`$. Hence $`𝒮(M_2())`$ is isomorphic (as a compact convex set) to the three-ball $`B^3`$ in $`^3`$. The pure states are precisely the points on the boundary,<sup>260</sup><sup>260</sup>260 The extreme boundary $`_eK`$ of a convex set $`K`$ consists of all $`\omega K`$ for which $`\omega =p\rho +(1p)\sigma `$ for some $`p(0,1)`$ and $`\rho ,\sigma K`$ implies $`\rho =\sigma =\omega `$. If $`K=𝒮(𝒜)`$ is the state space of a $`C^{}`$-algebra $`𝒜`$, the extreme boundary consists of the pure states on $`𝒜`$ (the remainder of $`𝒮(𝒜)`$ consisting of mixed states). If $`K`$ is embedded in a vector space, the extreme boundary $`_eK`$ may or may not coincide with the geometric boundary $`K`$ of $`K`$. In the case $`K=B^3^3`$ it does, but for an equilateral triangle in $`^2`$ it does not, since $`_eK`$ merely consists of the corners of the triangle whereas the geometric boundary includes the sides as well. i.e. the density matrices for which $`x^2+y^2+z^2=1`$ (for these and these alone define one-dimensional projections).<sup>261</sup><sup>261</sup>261Eq. (6.3) has the form $`\rho (x,y,z)=\frac{1}{2}(x\sigma _x+y\sigma _y+z\sigma _z)`$, where the $`\sigma _i`$ are the Pauli matrices. This yields an isomorphism between $`^3`$ and the Lie algebra of $`SO(3)`$ in its spin-$`\frac{1}{2}`$ representation $`𝒟_{1/2}`$ on $`^2`$. This isomorphism intertwines the defining action of $`SO(3)`$ on $`^3`$ with its adjoint action on $`M_2()`$. I.e., for any rotation $`R`$ one has $`\rho (R𝐱)=𝒟_{1/2}(R)\rho (𝐱)𝒟_{1/2}(R)^1`$. This will be used later on (see Subsection 6.5).
In order to define the continuous sections of the field, we introduce the symmetrization maps $`j_{NM}:𝒜_1^M𝒜_1^N`$, defined by
$$j_{NM}(A_M)=S_N(A_M1\mathrm{}1),$$
(6.4)
where one has $`NM`$ copies of the unit $`1𝒜_1`$ so as to obtain an element of $`𝒜_1^N`$. The symmetrization operator $`S_N:𝒜_1^N𝒜_1^N`$ is given by (linear and continuous) extension of
$$S_N(B_1\mathrm{}B_N)=\frac{1}{N!}\underset{\sigma 𝔖_N}{}B_{\sigma (1)}\mathrm{}B_{\sigma (N)},$$
(6.5)
where $`𝔖_N`$ is the permutation group (i.e. symmetric group) on $`N`$ elements and $`B_i𝒜_1`$ for all $`i=1,\mathrm{},N`$. For example, $`j_{N1}:𝒜_1𝒜_1^N`$ is given by
$$j_{N1}(B)=\overline{B}^{(N)}=\frac{1}{N}\underset{k=1}{\overset{N}{}}1\mathrm{}B_{(k)}1\mathrm{}1,$$
(6.6)
where $`B_{(k)}`$ is $`B`$ seen as an element of the $`k`$’th copy of $`𝒜_1`$ in $`𝒜_1^N`$. As our notation $`\overline{B}^{(N)}`$ indicates, this is just the ‘average’ of $`B`$ over all copies of $`𝒜_1`$. More generally, in forming $`j_{NM}(A_M)`$ an operator $`A_M𝒜_1^M`$ that involves $`M`$ sites is averaged over $`NM`$ sites. When $`N\mathrm{}`$ this means that one forms a macroscopic average of an $`M`$-particle operator.
We say that a sequence $`A=(A_1,A_2,\mathrm{})`$ with $`A_N𝒜_1^N`$ is symmetric when
$$A_N=j_{NM}(A_M)$$
(6.7)
for some fixed $`M`$ and all $`NM`$. In other words, the tail of a symmetric sequence entirely consists of ‘averaged’ or ‘intensive’ observables, which become macroscopic in the limit $`N\mathrm{}`$. Such sequences have the important property that they commute in this limit; more precisely, if $`A`$ and $`A^{}`$ are symmetric sequences, then
$$\underset{N\mathrm{}}{lim}A_NA_N^{}A_N^{}A_N=0.$$
(6.8)
As an enlightening special case we take $`A_N=j_{N1}(B)`$ and $`A_N^{}=j_{N1}(C)`$ with $`B,C𝒜_1`$. One immediately obtains from the relation $`[B_{(k)},C_{(l)}]=0`$ for $`kl`$ that
$$[\overline{B}^{(N)},\overline{C}^{(N)}]=\frac{1}{N}\overline{[B,C]}^{(N)}.$$
(6.9)
For example, if $`𝒜_1=M_2()`$ and if for $`B`$ and $`C`$ one takes the spin-$`\frac{1}{2}`$ operators $`S_j=\frac{\mathrm{}}{2}\sigma _j`$ for $`j=1,2,3`$ (where $`\sigma _j`$ are the Pauli matrices), then
$$[\overline{S}_j^{(N)},\overline{S}_k^{(N)}]=i\frac{\mathrm{}}{N}ϵ_{jkl}\overline{S}_l^{(N)}.$$
(6.10)
This shows that averaging one-particle operators leads to commutation relations formally like those of the one-particle operators in question, but with Planck’s constant $`\mathrm{}`$ replaced by a variable $`\mathrm{}/N`$. For constant $`\mathrm{}=1`$ this leads to the interval (6.1) over which our continuous field of $`C^{}`$-algebras is defined; for any other constant value of $`\mathrm{}`$ the field would be defined over $`I=0\mathrm{}/`$, which of course merely changes the labeling of the $`C^{}`$-algebras in question.
We return to the general case, and denote a section of the field with fibers (6.2) by a sequence $`A=(A_0,A_1,A_2,\mathrm{})`$, with $`A_0𝒜_0^{(c)}`$ and $`A_N𝒜_1^N`$ as before (i.e. the corresponding section is $`0A_0`$ and $`1/NA_N`$). We then complete the definition of our continuous field by declaring that a sequence $`A`$ defines a continuous section iff:
* $`(A_1,A_2,\mathrm{})`$ is approximately symmetric, in the sense that for any $`\epsilon >0`$ there is an $`N_\epsilon `$ and a symmetric sequence $`A^{}`$ such that $`A_NA_N^{}<\epsilon `$ for all $`NN_\epsilon `$;<sup>262</sup><sup>262</sup>262A symmetric sequence is evidently approximately symmetric.
* $`A_0(\omega )=lim_N\mathrm{}\omega ^N(A_N)`$, where $`\omega 𝒮(𝒜_1)`$ and $`\omega ^N𝒮(𝒜_1^N)`$ is the tensor product of $`N`$ copies of $`\omega `$, defined by (linear and continuous) extension of
$$\omega ^N(B_1\mathrm{}B_N)=\omega (B_1)\mathrm{}\omega (B_N).$$
(6.11)
This limit exists by definition of an approximately symmetric sequence.<sup>263</sup><sup>263</sup>263If $`(A_1,A_2,\mathrm{})`$ is symmetric with (6.7), one has $`\omega ^N(A_N)=\omega ^M(A_M)`$ for $`N>M`$, so that the tail of the sequence $`(\omega ^N(A_N))`$ is even independent of $`N`$. In the approximately symmetric case one easily proves that $`(\omega ^N(A_N))`$ is a Cauchy sequence.
It is not difficult to prove that this choice of continuous sections indeed defines a continuous field of $`C^{}`$-algebras over $`I=01/`$ with fibers (6.2). The main point is that
$$\underset{N\mathrm{}}{lim}A_N=A_0$$
(6.12)
whenever $`(A_0,A_1,A_2,\mathrm{})`$ satisfies the two conditions above.<sup>264</sup><sup>264</sup>264Given (6.12), the claim follows from Prop. II.1.2.3 in Landsman (1998) and the fact that the set of functions $`A_0`$ on $`𝒮(𝒜_1)`$ arising in the said way are dense in $`C(𝒮(𝒜_1))`$ (equipped with the supremum-norm). This follows from the Stone–Weierstrass theorem, from which one infers that the functions in question even exhaust $`𝒮(𝒜_1)`$. This is easy to show for symmetric sequences,<sup>265</sup><sup>265</sup>265 Assume (6.7), so that $`A_N=j_{NN}(A_N)`$ for $`NM`$. By the $`C^{}`$-axiom $`A^{}A=A^2`$ it suffices to prove (6.12) for $`A_0^{}=A_0`$, which implies $`A_M^{}=A_M`$ and hence $`A_N^{}=A_N`$ for all $`NM`$. One then has $`A_N=sup\{|\rho (A_N)|,\rho 𝒮(𝒜_1^N)\}`$. Because of the special form of $`A_N`$ one may replace the supremum over the set $`𝒮(𝒜_1^N)`$ of all states on $`𝒜_1^N`$ by the supremum over the set $`𝒮^p(𝒜_1^N)`$ of all permutation invariant states, which in turn may be replaced by the supremum over the extreme boundary $`𝒮^p(𝒜_1^N)`$ of $`𝒮^p(𝒜_1^N)`$. It is well known (Størmer, 1969; see also Subsection 6.2) that the latter consists of all states of the form $`\rho =\omega ^N`$, so that $`A_N=sup\{|\omega ^N(A_N)|,\omega 𝒮(𝒜_1)\}`$. This is actually equal to $`A_M=sup\{|\omega ^M(A_M)|\}`$. Now the norm in $`𝒜_0^{(c)}`$ is $`A_0=sup\{|A_0(\omega )|,\omega 𝒮(𝒜_1)\}`$, and by definition of $`A_0`$ one has $`A_0(\omega )=\omega ^M(A_M)`$. Hence (6.12) follows. and follows from this for approximately symmetric ones.
Consistent with (6.8), we conclude that in the limit $`N\mathrm{}`$ the macroscopic observables organize themselves in a commutative $`C^{}`$-algebra isomorphic to $`C(𝒮(𝒜_1))`$.
### 6.2 Quasilocal observables
In the $`C^{}`$-algebraic approach to quantum theory, infinite systems are usually described by means of inductive limit $`C^{}`$-algebras and the associated quasilocal observables (Thirring, 1983; Strocchi, 1985; Bratteli & Robinson, 1981, 1987; Haag, 1992; Araki, 1999; Sewell, 1986, 2002). To arrive at these notions in the case at hand, we proceed as follows (Duffield & Werner, 1992c).
A sequence $`A=(A_1,A_2,\mathrm{})`$ (where $`A_N𝒜_1^N`$, as before) is called local when for some fixed $`M`$ and all $`NM`$ one has $`A_N=A_M1\mathrm{}1`$ (where one has $`NM`$ copies of the unit $`1𝒜_1`$); cf. (6.4). A sequence is said to be quasilocal when for any $`\epsilon >0`$ there is an $`N_\epsilon `$ and a local sequence $`A^{}`$ such that $`A_NA_N^{}<\epsilon `$ for all $`NN_\epsilon `$. On this basis, we define the inductive limit $`C^{}`$-algebra
$$\overline{_N𝒜_1^N}$$
(6.13)
of the family of $`C^{}`$-algebras $`(𝒜_1^N)`$ with respect to the inclusion maps $`𝒜_1^N𝒜_1^{N+1}`$ given by $`A_NA_N1`$. As a set, (6.13) consists of all equivalence classes $`[A]A_0`$ of quasilocal sequences $`A`$ under the equivalence relation $`AB`$ when $`lim_N\mathrm{}A_NB_N=0`$. The norm on $`\overline{_N𝒜_1^N}`$ is
$$A_0=\underset{N\mathrm{}}{lim}A_N,$$
(6.14)
and the rest of the $`C^{}`$-algebraic structure is inherited from the quasilocal sequences in the obvious way (e.g., $`A_0^{}=[A^{}]`$ with $`A^{}=(A_1^{},A_2^{},\mathrm{})`$, etc.). As the notation suggests, each $`𝒜_1^N`$ is contained in $`\overline{_N𝒜_1^N}`$ as a $`C^{}`$-subalgebra by identifying $`A_N𝒜_1^N`$ with the local (and hence quasilocal) sequence $`A=(0,\mathrm{},0,A_N1,A_N11,\mathrm{})`$, and forming its equivalence class $`A_0`$ in $`\overline{_N𝒜_1^N}`$ as just explained.<sup>266</sup><sup>266</sup>266Of course, the entries $`A_1,\mathrm{}A_{N1}`$, which have been put to zero, are arbitrary. The assumption underlying the common idea that (6.13) is “the” algebra of observables of the infinite system under study is that by locality or some other human limitation the infinite tail of the system is not accessible, so that the observables must be arbitrarily close (i.e. in norm) to operators of the form $`A_N11,\mathrm{}`$ for some finite $`N`$.
This leads us to a second continuous field of $`C^{}`$-algebras $`𝒜^{(q)}`$ over $`01/`$, with fibers
$`𝒜_0^{(q)}`$ $`=`$ $`\overline{_N𝒜_1^N};`$
$`𝒜_{1/N}^{(q)}`$ $`=`$ $`𝒜_1^N.`$ (6.15)
Thus the suffix $`q`$ reminds one of that fact that the limit algebra $`𝒜_0^{(q)}`$ consists of quasilocal or quantum-mechanical observables. We equip the collection of $`C^{}`$-algebras (6.15) with the structure of a continuous field of $`C^{}`$-algebras $`𝒜^{(q)}`$ over $`01/`$ by declaring that the continuous sections are of the form $`(A_0,A_1,A_2,\mathrm{})`$ where $`(A_1,A_2,\mathrm{})`$ is quasilocal and $`A_0`$ is defined by this quasilocal sequence as just explained.<sup>267</sup><sup>267</sup>267The fact that this defines a continuous field follows from (6.14) and Prop. II.1.2.3 in Landsman (1998); cf. footnote 264. For $`N<\mathrm{}`$ this field has the same fibers
$$𝒜_{1/N}^{(q)}=𝒜_{1/N}^{(c)}=𝒜_1^N$$
(6.16)
as the continuous field $`𝒜`$ of the previous subsection, but the fiber $`𝒜_0^{(q)}`$ is completely different from $`𝒜_0^{(c)}`$. In particular, if $`𝒜_1`$ is noncommutative then so is $`𝒜_0^{(q)}`$, for it contains all $`𝒜_1^N`$.
The relationship between the continuous fields of $`C^{}`$-algebras $`𝒜^{(q)}`$ and $`𝒜^{(c)}`$ may be studied in two different (but related) ways. First, we may construct concrete representations of all $`C^{}`$-algebras $`𝒜_1^N`$, $`N<\mathrm{}`$, as well as of $`𝒜_0^{(c)}`$ and $`𝒜_0^{(q)}`$ on a single Hilbert space; this approach leads to superselections rules in the traditional sense. This method will be taken up in the next subsection. Second, we may look at those families of states $`(\omega _1,\omega _{1/2},\mathrm{},\omega _{1/N},\mathrm{})`$ (where $`\omega _{1/N}`$ is a state on $`𝒜_1^N`$) that admit limit states $`\omega _0^{(c)}`$ and $`\omega _0^{(q)}`$ on $`𝒜_0^{(c)}`$ and $`𝒜_0^{(q)}`$, respectively, such that the ensuing families of states $`(\omega _0^{(c)},\omega _1,\omega _{1/2},\mathrm{})`$ and $`(\omega _0^{(q)},\omega _1,\omega _{1/2},\mathrm{})`$ are continuous fields of states on $`𝒜^{(c)}`$ and on $`𝒜^{(q)}`$, respectively (cf. the end of Subsection 5.1).
Now, any state $`\omega _0^{(q)}`$ on $`𝒜_0^{(q)}`$ defines a state $`\omega _{0|1/N}^{(q)}`$ on $`𝒜_1^N`$ by restriction, and the ensuing field of states on $`𝒜^{(q)}`$ is clearly continuous. Conversely, any continuous field $`(\omega _0^{(q)},\omega _1,\omega _{1/2},\mathrm{},\omega _{1/N},\mathrm{})`$ of states on $`𝒜^{(q)}`$ becomes arbitrarily close to a field of the above type for $`N`$ large.<sup>268</sup><sup>268</sup>268 For any fixed quasilocal sequence $`(A_1,A_2,\mathrm{})`$ and $`\epsilon >0`$, there is an $`N_\epsilon `$ such that $`|\omega _{1/N}(A_N)\omega _{0|1/N}^{(q)}(A_N)|<\epsilon `$ for all $`N>N_\epsilon `$. However, the restrictions $`\omega _{0|1/N}^{(q)}`$ of a given state $`\omega _0^{(q)}`$ on $`𝒜_0^{(q)}`$ to $`𝒜_1^N`$ may not converge to a state $`\omega _0^{(c)}`$ on $`𝒜_0^{(c)}`$ for $`N\mathrm{}`$.<sup>269</sup><sup>269</sup>269See footnote 289 below for an example. States $`\omega _0^{(q)}`$ on $`\overline{_N𝒜_1^N}`$ that do have this property will here be called classical. In other words, $`\omega _{0|1/N}^{(q)}`$ is classical when there exists a probability measure $`\mu _0`$ on $`𝒮(𝒜_1)`$ such that
$$\underset{N\mathrm{}}{lim}_{𝒮(𝒜_1)}𝑑\mu _0(\rho )(\rho ^N(A_N)\omega _{0|1/N}^{(q)}(A_N))=0$$
(6.17)
for each (approximately) symmetric sequence $`(A_1,A_2,\mathrm{})`$. To analyze this notion we need a brief intermezzo on general $`C^{}`$-algebras and their representations.
* A folium in the state space $`𝒮()`$ of a $`C^{}`$-algebra $``$ is a convex, norm-closed subspace $``$ of $`𝒮()`$ with the property that if $`\omega `$ and $`B`$ such that $`\omega (B^{}B)>0`$, then the “reduced” state $`\omega _B:A\omega (B^{}AB)/\omega (B^{}B)`$ must be in $``$ (Haag, Kadison, & Kastler, 1970).<sup>270</sup><sup>270</sup>270See also Haag (1992). The name ‘folium’ is very badly chosen, since $`𝒮()`$ is by no means foliated by its folia; for example, a folium may contain subfolia. For example, if $`\pi `$ is a representation of $``$ on a Hilbert space $``$, then the set of all density matrices on $``$ (i.e. the $`\pi `$-normal states on $``$)<sup>271</sup><sup>271</sup>271A state $`\omega `$ on $``$ is called $`\pi `$-normal when it is of the form $`\omega (B)=\text{Tr}\rho \pi (B)`$ for some density matrix $`\rho `$. Hence the $`\pi `$-normal states are the normal states on the von Neumann algebra $`\pi ()^{\prime \prime }`$. comprises a folium $`_\pi `$. In particular, each state $`\omega `$ on $``$ defines a folium $`_\omega _{\pi _\omega }`$ through its GNS-representation $`\pi _\omega `$.
* Two representations $`\pi `$ and $`\pi ^{}`$ are called disjoint, written $`\pi \pi ^{}`$, if no subrepresentation of $`\pi `$ is (unitarily) equivalent to a subrepresentation of $`\pi ^{}`$ and vice versa. They are said to be quasi-equivalent, written $`\pi \pi ^{}`$, when $`\pi `$ has no subrepresentation disjoint from $`\pi ^{}`$, and vice versa.<sup>272</sup><sup>272</sup>272Equivalently, two representations $`\pi `$ and $`\pi ^{}`$ are disjoint iff no $`\pi `$-normal state is $`\pi ^{}`$-normal and vice versa, and quasi-equivalent iff each $`\pi `$-normal state is $`\pi ^{}`$-normal and vice versa. Quasi-equivalence is an equivalence relation $``$ on the set of representations. See Kadison & Ringrose (1986), Ch. 10.
* Similarly, two states $`\rho ,\sigma `$ are called either quasi-equivalent ($`\rho \sigma `$) or disjoint ($`\rho \sigma `$) when the corresponding GNS-representations have these properties.
* A state $`\omega `$ is called primary when the corresponding von Neumann algebra $`\pi _\omega ()^{\prime \prime }`$ is a factor.<sup>273</sup><sup>273</sup>273A von Neumann algebra $``$ acting on a Hilbert space is called a factor when its center $`^{}`$ is trivial, i.e. consists of multiples of the identity. Equivalently, $`\omega `$ is primary iff each subrepresentation of $`\pi _\omega ()`$ is quasi-equivalent to $`\pi _\omega ()`$, which is the case iff $`\pi _\omega ()`$ admits no (nontrivial) decomposition as the direct sum of two disjoint subrepresentations.
Now, there is a bijective correspondence between folia in $`𝒮()`$ and quasi-equivalence classes of representations of $``$, in that $`_\pi =_\pi ^{}`$ iff $`\pi \pi ^{}`$. Furthermore (as one sees from the GNS-construction), any folium $`𝒮()`$ is of the form $`=_\pi `$ for some representation $`\pi ()`$. Note that if $`\pi `$ is injective (i.e. faithful), then the corresponding folium is dense in $`𝒮()`$ in the weak$`^{}`$-topology by Fell’s Theorem. So in case that $``$ is simple,<sup>274</sup><sup>274</sup>274In the sense that it has no closed two-sided ideals. For example, the matrix algebra $`M_n()`$ is simple for any $`n`$, as is its infinite-dimensional analogue, the $`C^{}`$-algebra of all compact operators on a Hilbert space. The $`C^{}`$-algebra of quasilocal observables of an infinite quantum systems is typically simple as well. any folium is weak$`^{}`$-dense in the state space.
Two states need not be either disjoint or quasi-equivalent. This dichotomy does apply, however, within the class of primary states. Hence two primary states are either disjoint or quasi-equivalent. If $`\omega `$ is primary, then each state in the folium of $`\pi _\omega `$ is primary as well, and is quasi-equivalent to $`\omega `$. If, on the other hand, $`\rho `$ and $`\sigma `$ are primary and disjoint, then $`_\rho _\sigma =\mathrm{}`$. Pure states are, of course, primary.<sup>275</sup><sup>275</sup>275Since the corresponding GNS-representation $`\pi _\omega `$ is irreducible, $`\pi _\omega ()^{\prime \prime }=(_\omega )`$ is a factor. Furthermore, in thermodynamics pure phases are described by primary KMS states (Emch & Knops, 1970; Bratteli & Robinson, 1981; Haag, 1992; Sewell, 2002). This apparent relationship between primary states and “purity” of some sort is confirmed by our description of macroscopic observables:<sup>276</sup><sup>276</sup>276These claims easily follow from Sewell (2002), §2.6.5, which in turn relies on Hepp (1972).
* If $`\omega _0^{(q)}`$ is a classical primary state on $`𝒜_0^{(q)}=\overline{_N𝒜_1^N}`$, then the corresponding limit state $`\omega _0^{(c)}`$ on $`𝒜_0^{(c)}=C(𝒮(𝒜_1))`$ is pure (and hence given by a point in $`𝒮(𝒜_1)`$).
* If $`\rho _0^{(q)}`$ and $`\sigma _0^{(q)}`$ are classical primary states on $`𝒜_0^{(q)}`$, then
$`\rho _0^{(c)}=\sigma _0^{(c)}`$ $``$ $`\rho _0^{(q)}\sigma _0^{(q)};`$ (6.18)
$`\rho _0^{(c)}\sigma _0^{(c)}`$ $``$ $`\rho _0^{(q)}\sigma _0^{(q)}.`$ (6.19)
As in (6.17), a general classical state $`\omega _0^{(q)}`$ with limit state $`\omega _0^{(c)}`$ on $`C(𝒮(𝒜_1))`$ defines a probability measure $`\mu _0`$ on $`𝒮(𝒜_1)`$ by
$$\omega _0^{(c)}(f)=_{𝒮(𝒜_1)}𝑑\mu _0f,$$
(6.20)
which describes the probability distribution of the macroscopic observables in that state. As we have seen, this distribution is a delta function for primary states. In any case, it is insensitive to the microscopic details of $`\omega _0^{(q)}`$ in the sense that local modifications of $`\omega _0^{(q)}`$ do not affect the limit state $`\omega _0^{(c)}`$ (Sewell, 2002). Namely, it easily follows from (6.8) and the fact that the GNS-representation is cyclic that one can strengthen the second claim above:
> Each state in the folium $`_{\omega _0^{(q)}}`$ of a classical state $`\omega _0^{(q)}`$ is automatically classical and has the same limit state on $`𝒜_0^{(c)}`$ as $`\omega _0^{(q)}`$.
To make this discussion a bit more concrete, we now identify an important class of classical states on $`\overline{_N𝒜_1^N}`$. We say that a state $`\omega `$ on this $`C^{}`$-algebra is permutation-invariant when each of its restrictions to $`𝒜_1^N`$ is invariant under the natural action of the symmetric group $`𝔖_N`$ on $`𝒜_1^N`$ (i.e. $`\sigma 𝔖_N`$ maps an elementary tensor $`A_N=B_1\mathrm{}B_N𝒜_1^N`$ to $`B_{\sigma (1)}\mathrm{}B_{\sigma (N)}`$, cf. (6.5)). The structure of the set $`𝒮^𝔖`$ of all permutation-invariant states in $`𝒮(𝒜_0^{(q)})`$ has been analyzed by Størmer (1969). Like any compact convex set, it is the (weak$`^{}`$-closed) convex hull of its extreme boundary $`_e𝒮^𝔖`$. The latter consists of all infinite product states $`\omega =\rho ^{\mathrm{}}`$, where $`\rho 𝒮(𝒜_1)`$. I.e. if $`A_0𝒜_0^{(q)}`$ is an equivalence class $`[A_1,A_2,\mathrm{}]`$, then
$$\rho ^{\mathrm{}}(A_0)=\underset{N\mathrm{}}{lim}\rho ^N(A_N);$$
(6.21)
cf. (6.11). Equivalently, the restriction of $`\omega `$ to any $`𝒜_1^N𝒜_0^{(q)}`$ is given by $`^N\rho `$. Hence $`_e𝒮^𝔖`$ is isomorphic (as a compact convex set) to $`𝒮(𝒜_1)`$ in the obvious way, and the primary states in $`𝒮^𝔖`$ are precisely the elements of $`_e𝒮^𝔖`$.
A general state $`\omega _0^{(q)}`$ in $`𝒮^𝔖`$ has a unique decomposition<sup>277</sup><sup>277</sup>277This follows because $`𝒮^𝔖`$ is a so-called Bauer simplex (Alfsen, 1970). This is a compact convex set $`K`$ whose extreme boundary $`_eK`$ is closed and for which every $`\omega K`$ has a unique decomposition as a probability measure supported by $`_eK`$, in the sense that $`a(\omega )=_{_eK}𝑑\mu (\rho )a(\rho )`$ for any continuous affine function $`a`$ on $`K`$. For a unital $`C^{}`$-algebra $`𝒜`$ the continuous affine functions on the state space $`K=𝒮(𝒜)`$ are precisely the elements $`A`$ of $`𝒜`$, reinterpreted as functions $`\widehat{A}`$ on $`𝒮(𝒜)`$ by $`\widehat{A}(\omega )=\omega (A)`$. For example, the state space $`𝒮(𝒜)`$ of a commutative unital $`C^{}`$-algebra $`𝒜`$ is a Bauer simplex, which consists of all (regular Borel) probability measures on the pre state space $`𝒫(𝒜)`$.
$$\omega _0^{(q)}(A_0)=_{𝒮(𝒜_1)}𝑑\mu (\rho )\rho ^{\mathrm{}}(A_0),$$
(6.22)
where $`\mu `$ is a probability measure on $`𝒮(𝒜_1)`$ and $`A_0𝒜_0^{(q)}`$.<sup>278</sup><sup>278</sup>278 This is a quantum analogue of De Finetti’s representation theorem in classical probability theory (Heath & Sudderth, 1976; van Fraassen, 1991); see also Hudson & Moody (1975/76) and Caves et al. (2002). The following beautiful illustration of the abstract theory (Unnerstall, 1990a,b) is then clear from (6.17) and (6.22):
> If $`\omega _0^{(q)}`$ is permutation-invariant, then it is classical. The associated limit state $`\omega _0^{(c)}`$ on $`𝒜_0^{(c)}`$ is characterized by the fact that the measure $`\mu _0`$ in (6.20) coincides with the measure $`\mu `$ in (6.22).<sup>279</sup><sup>279</sup>279In fact, each state in the folium $`^𝔖`$ in $`𝒮(𝒜_0^{(q)})`$ corresponding to the (quasi-equivalence class of) the representation $`_{[\omega 𝒮^𝔖]}\pi _\omega `$ is classical.
### 6.3 Superselection rules
Infinite quantum systems are often associated with the notion of a superselection rule (or sector), which was originally introduced by Wick, Wightman, & Wigner (1952) in the setting of standard quantum mechanics on a Hilbert space $``$. The basic idea may be illustrated in the example of the boson/fermion (or “univalence”) superselection rule.<sup>280</sup><sup>280</sup>280See also Giulini (2003) for a modern mathematical treatment. Here one has a projective unitary representation $`𝒟`$ of the rotation group $`SO(3)`$ on $``$, for which $`𝒟(R_{2\pi })=\pm 1`$ for any rotation $`R_{2\pi }`$ of $`2\pi `$ around some axis. Specifically, on bosonic states $`\mathrm{\Psi }_B`$ one has $`𝒟(R_{2\pi })\mathrm{\Psi }_B=\mathrm{\Psi }_B`$, whereas on fermionic states $`\mathrm{\Psi }_F`$ the rule is $`𝒟(R_{2\pi })\mathrm{\Psi }_F=\mathrm{\Psi }_F`$. Now the argument is that a rotation of $`2\pi `$ accomplishes nothing, so that it cannot change the physical state of the system. This requirement evidently holds on the subspace $`_B`$ of bosonic states in $``$, but it is equally well satisfied on the subspace $`_F`$ of fermionic states, since $`\mathrm{\Psi }`$ and $`z\mathrm{\Psi }`$ with $`|z|=1`$ describe the same physical state. However, if $`\mathrm{\Psi }=c_B\mathrm{\Psi }_B+c_F\mathrm{\Psi }_F`$ (with $`|c_B|^2+|c_F|^2=1`$), then $`𝒟(R_{2\pi })\mathrm{\Psi }=c_B\mathrm{\Psi }_Bc_F\mathrm{\Psi }_F`$, which is not proportional to $`\mathrm{\Psi }`$ and apparently describes a genuinely different physical state from $`\mathrm{\Psi }`$.
The way out is to deny this conclusion by declaring that $`𝒟(R_{2\pi })\mathrm{\Psi }`$ and $`\mathrm{\Psi }`$ do describe the same physical state, and this is achieved by postulating that no physical observables $`A`$ (in their usual mathematical guise as operators on $``$) exist for which $`(\mathrm{\Psi }_B,A\mathrm{\Psi }_F)0`$. For in that case one has
$$(c_B\mathrm{\Psi }_B\pm c_F\mathrm{\Psi }_F,A(c_B\mathrm{\Psi }_B\pm c_F\mathrm{\Psi }_F))=|c_B|^2(\mathrm{\Psi }_B,A\mathrm{\Psi }_B)+|c_F|^2(\mathrm{\Psi }_F,A\mathrm{\Psi }_F)$$
(6.23)
for any observable $`A`$, so that $`(𝒟(R_{2\pi })\mathrm{\Psi },A𝒟(R_{2\pi })\mathrm{\Psi })=(\mathrm{\Psi },A\mathrm{\Psi })`$ for any $`\mathrm{\Psi }`$. Since any quantum-mechanical prediction ultimately rests on expectation values $`(\mathrm{\Psi },A\mathrm{\Psi })`$ for physical observables $`A`$, the conclusion is that a rotation of $`2\pi `$ indeed does nothing to the system. This is codified by saying that superpositions of the type $`c_B\mathrm{\Psi }_B+c_F\mathrm{\Psi }_F`$ are incoherent (whereas superpositions $`c_1\mathrm{\Psi }_1+c_2\mathrm{\Psi }_2`$ with $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ both in either $`_B`$ or in $`_F`$ are coherent). Each of the subspaces $`_B`$ and $`_F`$ of $``$ is said to be a superselection sector, and the statement that $`(\mathrm{\Psi }_B,A\mathrm{\Psi }_F)=0`$ for any observbale $`A`$ and $`\mathrm{\Psi }_B_B`$ and $`\mathrm{\Psi }_F_F`$ is called a superselection rule.<sup>281</sup><sup>281</sup>281In an ordinary selection rule between $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ one merely has $`(\mathrm{\Psi },H\mathrm{\Phi })=0`$ for the Hamiltonian $`H`$.
The price one pays for this solution is that states of the form $`c_B\mathrm{\Psi }_B+c_F\mathrm{\Psi }_F`$ with $`c_B0`$ and $`c_F0`$ are mixed, as one sees from (6.23). More generally, if $`=_{\lambda \mathrm{\Lambda }}_\lambda `$ with $`(\mathrm{\Psi },A\mathrm{\Phi })=0`$ whenever $`A`$ is an observable, $`\mathrm{\Psi }_\lambda `$, $`\mathrm{\Phi }_\lambda ^{}`$, and $`\lambda \lambda ^{}`$, and if in addition for each $`\lambda `$ and each pair $`\mathrm{\Psi },\mathrm{\Phi }_\lambda `$ there exists an observable $`A`$ for which $`(\mathrm{\Psi },A\mathrm{\Phi })0`$, then the subspaces $`_\lambda `$ are called superselection sectors in $``$. Again a key consequence of the occurrence of superselection sectors is that unit vectors of the type $`\mathrm{\Psi }=_\lambda c_\lambda \mathrm{\Psi }_\lambda `$ with $`\mathrm{\Psi }_\lambda `$ (and $`c_\lambda 0`$ for at least two $`\lambda `$’s) define mixed states
$$\psi (A)=(\mathrm{\Psi },A\mathrm{\Psi })=\underset{\lambda }{}|c_\lambda |^2(\mathrm{\Psi }_\lambda ,A\mathrm{\Psi }_\lambda )=\underset{\lambda }{}|c_\lambda |^2\psi _\lambda (A).$$
This procedure is rather ad hoc. A much deeper approach to superselection theory was developed by Haag and collaborators; see Roberts & Roepstorff (1969) for an introduction. Here the starting point is the abstract $`C^{}`$-algebra of observables $`𝒜`$ of a given quantum system, and superselection sectors are reinterpreted as equivalence classes (under unitary isomorphism) of irreducible representations of $`𝒜`$ (satisfying a certain selection criterion - see below). The connection between the concrete Hilbert space approach to superselection sectors discussed above and the abstract $`C^{}`$-algebraic approach is given by the following lemma (Hepp, 1972):<sup>282</sup><sup>282</sup>282Hepp proved a more general version of this lemma, in which ‘Two pure states $`\rho ,\sigma `$ on a $`C^{}`$-algebra $``$ define different sectors iff…’ is replaced by ‘Two states $`\rho ,\sigma `$ on a $`C^{}`$-algebra $``$ are disjoint iff…’
> Two pure states $`\rho ,\sigma `$ on a $`C^{}`$-algebra $`𝒜`$ define different sectors iff for each representation $`\pi (𝒜)`$ on a Hilbert space $``$ containing unit vectors $`\mathrm{\Psi }_\rho ,\mathrm{\Psi }_\sigma `$ such that $`\rho (A)=(\mathrm{\Psi }_\rho ,\pi (A)\mathrm{\Psi }_\rho )`$ and $`\sigma (A)=(\mathrm{\Psi }_\sigma ,\pi (A)\mathrm{\Psi }_\sigma )`$ for all $`A𝒜`$, one has $`(\mathrm{\Psi }_\rho ,\pi (A)\mathrm{\Psi }_\sigma )=0`$ for all $`A𝒜`$.
In practice, however, most irreducible representations of a typical $`C^{}`$-algebra $`𝒜`$ used in physics are physically irrelevant mathematical artefacts. Such representations may be excluded from consideration by some selection criterion. What this means depends on the context. For example, in quantum field theory this notion is made precise in the so-called DHR theory (reviewed by Roberts (1990), Haag (1992), Araki (1999), and Halvorson (2005)). In the class of theories discussed in the preceding two subsections, we take the algebra of observables $`𝒜`$ to be $`𝒜_0^{(q)}`$ \- essentially for reasons of human limitation - and for pedagogical reasons define (equivalence classes of) irreducible representations of $`𝒜_0^{(q)}`$ as superselection sectors, henceforth often just called sectors, only when they are equivalent to the GNS-representation given by a permutation-invariant pure state on $`𝒜_0^{(q)}`$. In particular, such a state is classical. On this selection criterion, the results in the preceding subsection trivially imply that there is a bijective correspondence between pure states on $`𝒜_1`$ and sectors of $`𝒜_0^{(q)}`$. The sectors of the commutative $`C^{}`$-algebra $`𝒜_0^{(c)}`$ are just the points of $`𝒮(𝒜_1)`$; note that a mixed state on $`𝒜_1`$ defines a pure state on $`𝒜_0^{(c)}`$! The role of the sectors of $`𝒜_1`$ in connection with those of $`𝒜_0^{(c)}`$ will be clarified in Subsection 6.5.
Whatever the model or the selection criterion, it is enlightening (and to some extent even in accordance with experimental practice) to consider superselection sectors entirely from the perspective of the pure states on the algebra of observables $`𝒜`$, removing $`𝒜`$ itself and its representations from the scene. To do so, we equip the space $`𝒫(𝒜)`$ of pure states on $`𝒜`$ with the structure of a transition probability space (von Neumann, 1981; Mielnik, 1968).<sup>283</sup><sup>283</sup>283See also Beltrametti & Cassinelli (1984) or Landsman (1998) for concise reviews. A transition probability on a set $`𝒫`$ is a function
$$p:𝒫\times 𝒫[0,1]$$
(6.24)
that satisfies
$$p(\rho ,\sigma )=1\rho =\sigma $$
(6.25)
and
$$p(\rho ,\sigma )=0p(\sigma ,\rho )=0.$$
(6.26)
A set with such a transition probability is called a transition probability space. Now, the pure state space $`𝒫(𝒜)`$ of a $`C^{}`$-algebra $`𝒜`$ carries precisely this structure if we define<sup>284</sup><sup>284</sup>284This definition applies to the case that $`𝒜`$ is unital; see Landsman (1998) for the general case. An analogous formula defines a transition probability on the extreme boundary of any compact convex set.
$$p(\rho ,\sigma ):=inf\{\rho (A)A𝒜,0A1,\sigma (A)=1\}.$$
(6.27)
To give a more palatable formula, note that since pure states are primary, two pure states $`\rho ,\sigma `$ are either disjoint ($`\rho \sigma `$) or else (quasi, hence unitarily) equivalent ($`\rho \sigma `$). In the first case, (6.27) yields
$$p(\rho ,\sigma )=0(\rho \sigma ).$$
(6.28)
Ine the second case it follows from Kadison’s transitivity theorem (cf. Thm. 10.2.6 in Kadison & Ringrose (1986)) that the Hilbert space $`_\rho `$ from the GNS-representation $`\pi _\rho (𝒜)`$ defined by $`\rho `$ contains a unit vector $`\mathrm{\Omega }_\sigma `$ (unique up to a phase) such that
$$\sigma (A)=(\mathrm{\Omega }_\sigma ,\pi _\rho (A)\mathrm{\Omega }_\sigma ).$$
(6.29)
Eq. (6.27) then leads to the well-known expression
$$p(\rho ,\sigma )=|(\mathrm{\Omega }_\rho ,\mathrm{\Omega }_\sigma )|^2(\rho \sigma ).$$
(6.30)
In particular, if $`𝒜`$ is commutative, then
$$p(\rho ,\sigma )=\delta _{\rho \sigma }.$$
(6.31)
For $`𝒜=M_2()`$ one obtains
$$p(\rho ,\sigma )=\frac{1}{2}(1+\mathrm{cos}\theta _{\rho \sigma }),$$
(6.32)
where $`\theta _{\rho \sigma }`$ is the angular distance between $`\rho `$ and $`\sigma `$ (seen as points on the two-sphere $`S^2=_eB^3`$, cf. (6.3) etc.), measured along a great circle.
Superselection sectors may now be defined for any transition probability spaces $`𝒫`$. A family of subsets of $`𝒫`$ is called orthogonal if $`p(\rho ,\sigma )=0`$ whenever $`\rho `$ and $`\sigma `$ do not lie in the same subset. The space $`𝒫`$ is called reducible if it is the union of two (nonempty) orthogonal subsets; if not, it is said to be irreducible. A component of $`𝒫`$ is a subset $`𝒞𝒫`$ such that $`𝒞`$ and $`𝒫\backslash 𝒞`$ are orthogonal. An irreducible component of $`𝒫`$ is called a (superselection) sector. Thus $`𝒫`$ is the disjoint union of its sectors. For $`𝒫=𝒫(𝒜)`$ this reproduces the algebraic definition of a superselection sector (modulo the selection criterion) via the correspondence between states and representations given by the GNS-constructions. For example, in the commutative case $`𝒜C(X)`$ each point in $`X𝒫(𝒜)`$ is its own little sector.
### 6.4 A simple example: the infinite spin chain
Let us illustrate the occurrence of superselection sectors in a simple example, where the algebra of observables is $`𝒜_0^{(q)}`$ with $`𝒜_1=M_2()`$. Let $`_1=^2`$, so that $`_1^N=^N^2`$ is the tensor product of $`N`$ copies of $`^2`$. It is clear that $`𝒜_1^N`$ acts on $`_1^N`$ in a natural way (i.e. componentwise). This defines an irreducible representation $`\pi _N`$ of $`𝒜_1^N`$, which is indeed its unique irreducible representation (up to unitary equivalence). In particular, for $`N<\mathrm{}`$ the quantum system whose algebra of observables is $`𝒜_1^N`$ (such as a chain with $`N`$ two-level systems) has no superselection rules. We define the $`N\mathrm{}`$ limit “$`(M_2())^{\mathrm{}}`$” of the $`C^{}`$-algebras $`(M_2())^N`$ as the inductive limit $`𝒜_0^{(q)}`$ for $`𝒜_1=M_2()`$, as introduced in Subsection 6.2; see (6.13). The definition of “$`^{\mathrm{}}^2`$” is slightly more involved, as follows (von Neumann, 1938).
For any Hilbert space $`_1`$, let $`\mathrm{\Psi }`$ be a sequence $`(\mathrm{\Psi }_1,\mathrm{\Psi }_2,\mathrm{})`$ with $`\mathrm{\Psi }_n_1`$. The space $`𝖧_1`$ of such sequences is a vector space in the obvious way. Now let $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ be two such sequences, and write $`(\mathrm{\Psi }_n,\mathrm{\Phi }_n)=\mathrm{exp}(i\alpha _n)|(\mathrm{\Psi }_n,\mathrm{\Phi }_n)|`$. If $`_n|\alpha _n|=\mathrm{}`$, we define the (pre-) inner product $`(\mathrm{\Psi },\mathrm{\Phi })`$ to be zero. If $`_n|\alpha _n|<\mathrm{}`$, we put $`(\mathrm{\Psi },\mathrm{\Phi })=_n(\mathrm{\Psi }_n,\mathrm{\Phi }_n)`$ (which, of course, may still be zero!). The (vector space) quotient of $`𝖧_1`$ by the space of sequences $`\mathrm{\Psi }`$ for which $`(\mathrm{\Psi },\mathrm{\Psi })=0`$ can be completed to a Hilbert space $`_1^{\mathrm{}}`$ in the induced inner product, called the complete infinite tensor product of the Hilbert space $`_1`$ (over the index set $``$).<sup>285</sup><sup>285</sup>285Each fixed $`\mathrm{\Psi }_1`$ defines an incomplete tensor product $`_\mathrm{\Psi }^{\mathrm{}}`$, defined as the closed subspace of $`_1^{\mathrm{}}`$ consisting of all $`\mathrm{\Phi }`$ for which $`_n|(\mathrm{\Psi }_n,\mathrm{\Phi }_n)1|<\mathrm{}`$. If $`_1`$ is separable, then so is $`_\mathrm{\Psi }^{\mathrm{}}`$ (in contrast to $`_1^{\mathrm{}}`$, which is an uncountable direct sum of the $`_\mathrm{\Psi }^{\mathrm{}}`$). We apply this construction with $`_1=^2`$. If $`(e_i)`$ is some basis of $`^2`$, an orthonormal basis of $`_1^{\mathrm{}}`$ then consists of all different infinite strings $`e_{i_1}\mathrm{}e_{i_n}\mathrm{}`$, where $`e_{i_n}`$ is $`e_i`$ regarded as a vector in $`^2`$.<sup>286</sup><sup>286</sup>286The cardinality of the set of all such strings equals that of $``$, so that $`_1^{\mathrm{}}`$ is non-separable, as claimed. We denote the multi-index $`(i_1,\mathrm{},i_n,\mathrm{})`$ simply by $`I`$, and the corresponding basis vector by $`e_I`$.
This Hilbert space $`_1^{\mathrm{}}`$ carries a natural faithful representation $`\pi `$ of $`𝒜_0^{(q)}`$: if $`A_0𝒜_0^{(q)}`$ is an equivalence class $`[A_1,A_2,\mathrm{}]`$, then $`\pi (A_0)e_I=lim_N\mathrm{}A_Ne_i`$, where $`A_N`$ acts on the first $`N`$ components of $`e_I`$ and leaves the remainder unchanged.<sup>287</sup><sup>287</sup>287Indeed, this yields an alternative way of defining $`\overline{_N𝒜_1^N}`$ as the norm closure of the union of all $`𝒜_1^N`$ acting on $`_1^{\mathrm{}}`$ in the stated way. Now the point is that although each $`𝒜_1^N`$ acts irreducibly on $`_1^N`$, the representation $`\pi (𝒜_0^{(q)})`$ on $`_1^{\mathrm{}}`$ thus constructed is highly reducible. The reason for this is that by definition (quasi-) local elements of $`𝒜_0^{(q)}`$ leave the infinite tail of a vector in $`_1^{\mathrm{}}`$ (almost) unaffected, so that vectors with different tails lie in different superselection sectors. Without the quasi-locality condition on the elements of $`𝒜_0^{(q)}`$, no superselection rules would arise. For example, in terms of the usual basis
$$\{=\left(\begin{array}{c}1\\ 0\end{array}\right),=\left(\begin{array}{c}0\\ 1\end{array}\right)\}$$
(6.33)
of $`^2`$, the vectors $`\mathrm{\Psi }_{}=\mathrm{}\mathrm{}`$ (i.e. an infinite product of ‘up’ vectors) and $`\mathrm{\Psi }_{}=\mathrm{}\mathrm{}`$ (i.e. an infinite product of ‘down’ vectors) lie in different sectors. The reason why the inner product $`(\mathrm{\Psi }_{},\pi (A)\mathrm{\Psi }_{})`$ vanishes for any $`A𝒜_0^{(q)}`$ is that for local observables $`A`$ one has $`\pi (A)=A_M1\mathrm{}1\mathrm{}`$ for some $`A_M(_M)`$; the inner product in question therefore involves infinitely many factors $`(,1)=(,)=0`$. For quasilocal $`A`$ the operator $`\pi (A)`$ might have a small nontrivial tail, but the inner product vanishes nonetheless by an approximation argument.
More generally, elementary analysis shows that $`(\mathrm{\Psi }_u,\pi (A)\mathrm{\Psi }_v)=0`$ whenever $`\mathrm{\Psi }_u=^{\mathrm{}}u`$ and $`\mathrm{\Psi }_v=^{\mathrm{}}v`$ for unit vectors $`u,v^2`$ with $`uv`$. The corresponding vector states $`\psi _u`$ and $`\psi _v`$ on $`𝒜_0^{(q)}`$ (i.e. $`\psi _u(A)=(\mathrm{\Psi }_u,\pi (A)\mathrm{\Psi }_u)`$ etc.) are obviously permutation-invariant and hence classical. Identifying $`𝒮(M_2())`$ with $`B^3`$, as in (6.3), the corresponding limit state $`(\psi _u)_0`$ on $`𝒜_0^{(c)}`$ defined by $`\psi _u`$ is given by (evaluation at) the point $`\stackrel{~}{u}=(x,y,z)`$ of $`_eB^3=S^2`$ (i.e. the two-sphere) for which the corresponding density matrix $`\rho (\stackrel{~}{u})`$ is the projection operator onto $`u`$. It follows that $`\psi _u`$ and $`\psi _v`$ are disjoint; cf. (6.19). We conclude that each unit vector $`u^2`$ determines a superselection sector $`\pi _u`$, namely the GNS-representation of the corresponding state $`\psi _u`$, and that each such sector is realized as a subspace $`_u`$ of $`_1^{\mathrm{}}`$ (viz. $`_u=\overline{\pi (𝒜_0^{(q)})\mathrm{\Psi }_u}`$). Moreover, since a permutation-invariant state on $`𝒜_0^{(q)}`$ is pure iff it is of the form $`\psi _u`$, we have found all superselection sectors of our system. Thus in what follows we may concentrate our attention on the subspace (of $`_1^{\mathrm{}}`$) and subrepresentation (of $`\pi `$)
$`_𝔖`$ $`=`$ $`_{\stackrel{~}{u}S^2}_u;`$
$`\pi _𝔖(𝒜_0^{(q)})`$ $`=`$ $`_{\stackrel{~}{u}S^2}\pi _u(𝒜_0^{(q)}),`$ (6.34)
where $`\pi _u`$ is simply the restriction of $`\pi `$ to $`_u_1^{\mathrm{}}`$.
In the presence of superselection sectors one may construct operators that distinguish different sectors whilst being a multiple of the unit in each sector. In quantum field theory these are typically global charges, and in our example the macroscopic observables play this role. To see this, we return to Subsection 6.1. It is not difficult to show that for any approximately symmetric sequence $`(A_1,A_2,\mathrm{})`$ the limit
$$\overline{A}=\underset{N\mathrm{}}{lim}\pi _𝔖(A_N)$$
(6.35)
exists in the strong operator topology on $`(_𝔖)`$ (Bona, 1988). Moreover, if $`A_0𝒜_0^{(c)}=C(𝒮(𝒜_1))`$ is the function defined by the given sequence,<sup>288</sup><sup>288</sup>288Recall that $`A_0(\omega )=lim_N\mathrm{}\omega ^N(A_N)`$. then the map $`A_0\overline{A}`$ defines a faithful representation of $`𝒜_0^{(c)}`$ on $`_𝔖`$, which we call $`\pi _𝔖`$ as well (by abuse of notation). An easy calculation in fact shows that $`\pi _𝔖(A_0)\mathrm{\Psi }=A_0(\stackrel{~}{u})\mathrm{\Psi }`$ for $`\mathrm{\Psi }_u`$, or, in other words,
$$\pi _𝔖(A_0)=_{\stackrel{~}{u}S^2}A_0(\stackrel{~}{u})1__u.$$
(6.36)
Thus the $`\pi _𝔖(A_0)`$ indeed serve as the operators in question.
To illustrate how delicate all this is, it may be interesting to note that even for symmetric sequences the limit $`lim_N\mathrm{}\pi (A_N)`$ does not exist on $`_1^{\mathrm{}}`$, not even in the strong topology.<sup>289</sup><sup>289</sup>289 For example, let us take the sequence $`A_N=j_{N1}(\mathrm{diag}(1,1))`$ and the vector $`\mathrm{\Psi }=\mathrm{},`$ where a sequence of $`2^N`$ factors of $``$ is followed by $`2^{N+1}`$ factors of $``$, etc. Then the sequence $`\{\pi (A_N)\mathrm{\Psi }\}_N`$ in $`_1^{\mathrm{}}`$ diverges: the subsequence where $`N`$ runs over all numbers $`2^n`$ with $`n`$ odd converges to $`\frac{1}{3}\mathrm{\Psi }`$, whereas the subsequence where $`N`$ runs over all $`2^n`$ with $`n`$ even converges to $`\frac{1}{3}\mathrm{\Psi }`$. On the positive side, it can be shown that $`lim_N\mathrm{}\pi (A_N)\mathrm{\Psi }`$ exists as an element of the von Neumann algebra $`\pi (𝒜_0^{(q)})^{\prime \prime }`$ whenever the vector state $`\psi `$ defined by $`\mathrm{\Psi }`$ lies in the folium $`^𝔖`$ generated by all permutation-invariant states (Bona, 1988; Unnerstall, 1990a).
This observation is part of a general theory of macroscopic observables in the setting of von Neumann algebras (Primas, 1983; Rieckers, 1984; Amann, 1986, 1987; Morchio & Strocchi, 1987; Bona, 1988, 1989; Unnerstall, 1990a, 1990b; Breuer, 1994; Atmanspacher, Amann, & Müller-Herold, 1999), which complements the purely $`C^{}`$-algebraic approach of Raggio & Werner (1989, 1991), Duffield & Werner (1992a,b,c), and Duffield, Roos, & Werner (1992) explained so far.<sup>290</sup><sup>290</sup>290Realistic models have been studied in the context of both the $`C^{}`$-algebraic and the von Neumann algebraic approach by Rieckers and his associates. See, for example, Honegger & Rieckers (1994), Gerisch, Münzner, & Rieckers (1999), Gerisch, Honegger, & Rieckers (2003), and many other papers. For altogether different approaches to macroscopic observables see van Kampen (1954, 1988, 1993), Wan & Fountain (1998), Harrison & Wan (1997), Wan et al. (1998), Fröhlich, Tsai, & Yau (2002), and Poulin (2004). In our opinion, the latter has the advantage that conceptually the passage to the limit $`N\mathrm{}`$ (and thereby the idealization of a large system as an infinite one) is very satisfactory, especially in our reformulation in terms of continuous fields of $`C^{}`$-algebras. Here the commutative $`C^{}`$-algebra $`𝒜_0^{(c)}`$ of macroscopic observables of the infinite system is glued to the noncommutative algebras $`𝒜_1^N`$ of the corresponding finite systems in a continuous way, and the continuous sections of the ensuing continuous field of $`C^{}`$-algebras $`𝒜^{(c)}`$ exactly describe how macroscopic quantum observables of the finite systems converge to classical ones. Microscopic quantum observables of the pertinent finite systems, on the other hand, converge to quantum observables of the infinite quantum system, and this convergence is described by the continuous sections of the continuous field of $`C^{}`$-algebras $`𝒜^{(q)}`$. This entirely avoids the language of superselection rules, which rather displays a shocking discontinuity between finite and infinite systems: for superselection rules do not exist in finite systems!<sup>291</sup><sup>291</sup>291We here refer to superselection rules in the traditional sense of inequivalent irreducible representations of simple $`C^{}`$-algebras. For topological reasons certain finite-dimensional systems are described by (non-simple) $`C^{}`$-algebras that do admit inequivalent irreducible representations (Landsman, 1990a,b).
### 6.5 Poisson structure and dynamics
We now pass to the discussion of time-evolution in infinite systems of the type considered so far. We start with the observation that the state space $`𝒮()`$ of a finite-dimensional $`C^{}`$-algebra $``$ (for simplicity assumed unital in what follows) is a Poisson manifold (cf. Subsection 4.3) in a natural way. A similar statement holds in the infinite-dimensional case, and we carry the reader through the necessary adaptations of the main argument by means of footnotes.<sup>292</sup><sup>292</sup>292Of which this is the first. When $``$ is infinite-dimensional, the state space $`𝒮()`$ is no longer a manifold, let alone a Poisson manifold, but a Poisson space (Landsman, 1997, 1998). This is a generalization of a Poisson manifold, which turns a crucial property of the latter into a definition. This property is the foliation of a Poisson manifold by its symplectic leaves (Weinstein, 1983), and the corresponding definition is as follows: A Poisson space $`P`$ is a Hausdorff space of the form $`P=_\alpha S_\alpha `$ (disjoint union), where each $`S_\alpha `$ is a symplectic manifold (possibly infinite-dimensional) and each injection $`\iota _\alpha :S_\alpha P`$ is continuous. Furthermore, one has a linear subspace $`FC(P,)`$ that separates points and has the property that the restriction of each $`fF`$ to each $`S_\alpha `$ is smooth. Finally, if $`f,gF`$ then $`\{f,g\}F`$, where the Poisson bracket is defined by $`\{f,g\}(\iota _\alpha (\sigma ))=\{\iota _\alpha ^{}f,\iota _\alpha ^{}g\}_\alpha (\sigma )`$. Clearly, a Poisson manifold $`M`$ defines a Poisson space if one takes $`P=M`$, $`F=C^{\mathrm{}}(M)`$, and the $`S_\alpha `$ to be the symplectic leaves defined by the given Poisson bracket. Thus we refer to the manifolds $`S_\alpha `$ in the above definition as the symplectic leaves of $`P`$ as well. We write $`K=𝒮()`$.
Firstly, an element $`A`$ defines a linear function $`\widehat{A}`$ on $`^{}`$ and hence on $`K`$ (namely by restriction) through $`\widehat{A}(\omega )=\omega (A)`$. For such functions we define the Poisson bracket by
$$\{\widehat{A},\widehat{B}\}=i\widehat{[A,B]}.$$
(6.37)
Here the factor $`i`$ has been inserted in order to make the Poisson bracket of two real-valed functions real-valued again; for $`\widehat{A}`$ is real-valued on $`K`$ precisely when $`A`$ is self-adjoint, and if $`A^{}=A`$ and $`B^{}=B`$, then $`i[A,B]`$ is self-adjoint (whereras $`[A,B]`$ is skew-adjoint). In general, for $`f,gC^{\mathrm{}}(K)`$ we put
$$\{f,g\}(\omega )=i\omega ([df_\omega ,dg_\omega ]),$$
(6.38)
interpreted as follows.<sup>293</sup><sup>293</sup>293In the infinite-dimensional case $`C^{\mathrm{}}(K)`$ is defined as the intersection of the smooth functions on $`K`$ with respect to its Banach manifold structure and the space $`C(K)`$ of weak$`^{}`$-continuous functions on $`K`$. The differential forms $`df`$ and $`dg`$ in (6.38) also require an appropriate definition; see Duffield & Werner (1992a), Bona (2000), and Odzijewicz & Ratiu (2003) for the technicalities. Let $`_{}`$ be the self-adjoint part of $``$, and interpret $`K`$ as a subspace of $`_{}^{}`$; since a state $`\omega `$ satisfies $`\omega (A^{})=\overline{\omega (A)}`$ for all $`A`$, it is determined by its values on self-adjoint elements. Subsequently, we identify the tangent space at $`\omega `$ with
$$T_\omega K=\{\rho _{}^{}\rho (1)=0\}_{}^{}$$
(6.39)
and the cotangent space at $`\omega `$ with the quotient (of real Banach spaces)
$$T_\omega ^{}K=_{}^{}/1,$$
(6.40)
where the unit $`1`$ is regarded as an element of $`^{}`$ through the canonical embedding $`^{}`$. Consequently, the differential forms $`df`$ and $`dg`$ at $`\omega K`$ define elements of $`_{}^{}/1`$. The commutator in (6.38) is then defined as follows: one lifts $`df_\omega _{}^{}/1`$ to $`_{}^{}`$, and uses the natural isomorphism $`^{}`$ typical of finite-dimensional vector spaces.<sup>294</sup><sup>294</sup>294In the infinite-dimensional case one uses the canonical identification between $`^{}`$ and the enveloping von Neumann algebra of $``$ to define the commutator. The arbitrariness in this lift is a multiple of 1, which drops out of the commutator. Hence $`i[df_\omega ,dg_\omega ]`$ is an element of $`_{}^{}_{}`$, on which the value of the functional $`\omega `$ is defined.<sup>295</sup><sup>295</sup>295If $``$ is infinite-dimensional, one here regards $`^{}`$ as the predual of the von Neumann algebra $`^{}`$. This completes the definition of the Poisson bracket; one easily recovers (6.37) as a special case of (6.38).
The symplectic leaves of the given Poisson structure on $`K`$ have been determined by Duffield & Werner (1992a).<sup>296</sup><sup>296</sup>296See also Bona (2000) for the infinite-dimensional special case where $``$ is the $`C^{}`$-algebra of compact operators. Namely:
> Two states $`\rho `$ and $`\sigma `$ lie in the same symplectic leaf of $`𝒮()`$ iff $`\rho (A)=\sigma (UAU^{})`$ for some unitary $`U`$.
When $`\rho `$ and $`\sigma `$ are pure, this is the case iff the corresponding GNS-representations $`\pi _\rho ()`$ and $`\pi _\sigma ()`$ are unitarily equivalent,<sup>297</sup><sup>297</sup>297Cf. Thm. 10.2.6 in Kadison & Ringrose (1986). but in general the implication holds only in one direction: if $`\rho `$ and $`\sigma `$ lie in the same leaf, then they have unitarily equivalent GNS-representations.<sup>298</sup><sup>298</sup>298An important step of the proof is the observation that the Hamiltonian vector field $`\xi _f(\omega )T_\omega K𝒜_{}^{}`$ of $`fC^{\mathrm{}}(K)`$ is given by $`\xi _f(\omega ),B=i[df_\omega ,B]`$, where $`B_{}_{}^{}`$ and $`df_\omega _{}^{}/1`$. (For example, this gives $`\xi _{\widehat{A}}\widehat{B}=i\widehat{[A,B]}=\{\widehat{A},\widehat{B}\}`$ by (6.37), as it should be.) If $`\phi _t^h`$ denotes the Hamiltonian flow of $`h`$ at time $`t`$, it follows (cf. Duffield, Roos, & Werner (1992), Prop. 6.1 or Duffield & Werner (1992a), Prop. 3.1) that $`\phi _h^t(\omega ),B=\omega ,U_t^hB(U_t^h)^{}`$ for some unitary $`U_t^h`$. For example, if $`h=\widehat{A}`$ then $`U_t^h=\mathrm{exp}(itA)`$.
It follows from this characterization of the symplectic leaves of $`K=𝒮()`$ that the pure state space $`_eK=𝒫()`$ inherits the Poisson bracket from $`K`$, and thereby becomes a Poisson manifold in its own right.<sup>299</sup><sup>299</sup>299More generally, a Poisson space. The structure of $`𝒫()`$ as a Poisson space was introduced by Landsman (1997, 1998) without recourse to the full state space or the work of Duffield & Werner (1992a). This leads to an important connection between the superselection sectors of $``$ and the Poisson structure on $`𝒫()`$ (Landsman, 1997, 1998):
> The sectors of the pure state space $`𝒫()`$ of a $`C^{}`$-algebra $``$ as a transition probability space coincide with its symplectic leaves as a Poisson manifold.
For example, when $`C(X)`$ is commutative, the space $`𝒮(C(X))`$ of all (regular Borel) probability measures on $`X`$ acquires a Poisson bracket that is identically zero, as does its extreme boundary $`X`$. It follows from (6.31) that the sectors in $`X`$ are its points, and so are its symplectic leaves (in view of their definition and the vanishing Poisson bracket). The simplest noncommutative case is $`=M_2()`$, for which the symplectic leaves of the state space $`K=𝒮(M_2())B^3`$ (cf. (6.3)) are the spheres with constant radius.<sup>300</sup><sup>300</sup>300 Equipped with a multiple of the so-called Fubini–Study symplectic structure; see Landsman (1998) or any decent book on differential geometry for this notion. This claim is immediate from footnote 261. More generally, the pure state space of $`M_n()`$ is the projective space $`^n`$, which again becomes equipped with the Fubini–Study symplectic structure. This is even true for $`n=\mathrm{}`$ if one defines $`M_{\mathrm{}}()`$ as the $`C^{}`$-algebra of compact operators on a separable Hilbert space $``$: in that case one has $`𝒫(M_{\mathrm{}}())`$. Cf. Cantoni (1977), Cirelli, Lanzavecchia, & Maniá (1983), Cirelli, Maniá, & Pizzocchero (1990), Landsman (1998), Ashtekar & Schilling (1999), Marmo et al. (2005), etc. The sphere with radius 1 consists of points in $`B^3`$ that correspond to pure states on $`M_2()`$, all interior symplectic leaves of $`K`$ coming from mixed states on $`M_2()`$.
The coincidence of sectors and symplectic leaves of $`𝒫()`$ is a compatibility condition between the transition probability structure and the Poisson structure. It is typical of the specific choices (6.27) and (6.38), respectively, and hence of quantum theory. In classical mechanics one has the freedom of equipping a manifold $`M`$ with an arbitrary Poisson structure, and yet use $`C_0(M)`$ as the commutative $`C^{}`$-algebra of observables. The transition probability (6.31) (which follows from (6.27) in the commutative case) are clearly the correct ones in classical physics, but since the symplectic leaves of $`M`$ can be almost anything, the coincidence in question does not hold.
However, there exists a compatibility condition between the transition probability structure and the Poisson structure, which is shared by classical and quantum theory. This is the property of unitarity of a Hamiltonian flow, which in the present setting we formulate as follows.<sup>301</sup><sup>301</sup>301All this can be boosted into an axiomatic structure into which both classical and quantum theory fit; see Landsman (1997, 1998). First, in quantum theory with algebra of observables $``$ we define time-evolution (in the sense of an automorphic action of the abelian group $``$ on $``$, i.e. a one-parameter group $`\alpha `$ of automorphisms on $``$) to be Hamiltonian when $`A(t)=\alpha _t(A)`$ satisfies the Heisenberg equation $`i\mathrm{}dA/dt=[A,H]`$ for some self-adjoint element $`H`$. The corresponding flow on $`𝒫()`$ \- i.e. $`\omega _t(A)=\omega (A(t))`$ \- is equally well said to be Hamiltonian in that case. In classical mechanics with Poisson manifold $`M`$ we similarly say that a flow on $`M`$ is Hamiltonian when it is the flow of a Hamiltonian vector field $`\xi _h`$ for some $`hC^{\mathrm{}}(M)`$. (Equivalently, the time-evolution of the observables $`fC^{\mathrm{}}(M)`$ is given by $`df/dt=\{h,f\}`$; cf. (5.18) etc.) The point is that in either case the flow is unitary in the sense that
$$p(\rho (t),\sigma (t))=p(\rho ,\sigma )$$
(6.41)
for all $`t`$ and all $`\rho ,\sigma P`$ with $`P=𝒫()`$ (equipped with the transition probabilities (6.27) and the Poisson bracket (6.38)) or $`P=M`$ (equipped with the transition probabilities (6.31) and any Poisson bracket).<sup>302</sup><sup>302</sup>302In quantum theory the flow is defined for any $`t`$. In classical dynamics, (6.41) holds for all $`t`$ for which $`\rho (t)`$ and $`\sigma (t)`$ are defined, cf. footnote 254.
In both cases $`P=𝒫()`$ and $`P=M`$, a Hamiltonian flow has the property (which is immediate from the definition of a symplectic leaf) that for all (finite) times $`t`$ a point $`\omega (t)`$ lies in the same symplectic leaf of $`P`$ as $`\omega =\omega (0)`$. In particular, in quantum theory $`\omega (t)`$ and $`\omega `$ must lie in the same sector. In the quantum theory of infinite systems an automorphic time-evolution is rarely Hamiltonian, but one reaches a similar conclusion under a weaker assumption. Namely, if a given one-parameter group of automorphisms $`\alpha `$ on $``$ is implemented in the GNS-representation $`\pi _\omega ()`$ for some $`\omega 𝒫()`$,<sup>303</sup><sup>303</sup>303This assumption means that there exists a unitary representation $`tU_t`$ of $``$ on $`_\omega `$ such that $`\pi _\omega (\alpha _t(A))=U_t\pi _\omega (A)U_t^{}`$ for all $`A`$ and all $`t`$. then $`\omega (t)`$ and $`\omega `$ lie in the same sector and hence in the same symplectic leaf of $`𝒫()`$.
To illustrate these concepts, let us return to our continuous field of $`C^{}`$-algebras $`𝒜^{(c)}`$; cf. (6.2). It may not come as a great surprise that the canonical $`C^{}`$-algebraic transition probabilities (6.27) on the pure state space of each fiber algebra $`𝒜_{1/N}^{(c)}`$ for $`N<\mathrm{}`$ converge to the classical transition probabilities (6.31) on the commutative limit algebra $`𝒜_0^{(c)}`$. Similarly, the $`C^{}`$-algebraic Poisson structure (6.38) on each $`𝒫(𝒜_{1/N}^{(c)})`$ converges to zero. However, we know from the limit $`\mathrm{}0`$ of quantum mechanics that in generating classical behaviour on the limit algebra of a continuous field of $`C^{}`$-algebras one should rescale the commutators; see Subsection 4.3 and Section 5. Thus we replace the Poisson bracket (6.38) for $`𝒜_{1/N}^{(c)}`$ by
$$\{f,g\}(\omega )=iN\omega ([df_\omega ,dg_\omega ]).$$
(6.42)
Thus rescaled, the Poisson brackets on the spaces $`𝒫(𝒜_{1/N}^{(c)})`$ turn out to converge to the canonical Poisson bracket (6.38) on $`𝒫(𝒜_0^{(c)})=𝒮(𝒜_1)`$, instead of the zero bracket expected from the commutative nature of the limit algebra $`𝒜_0^{(c)}`$. Consequently, the symplectic leaves of the full state space $`𝒮(𝒜_1)`$ of the fiber algebra $`𝒜_1^{(c)}`$ become the symplectic leaves of the pure state space $`𝒮(𝒜_1)`$ of the fiber algebra $`𝒜_0^{(c)}`$. This is undoubtedly indicative of the origin of classical phase spaces and their Poisson structures in quantum theory.
More precisely, we have the following result (Duffield & Werner, 1992a):
> If $`A=(A_0,A_1,A_2,\mathrm{})`$ and $`A^{}=(A_0^{},A_1^{},A_2^{},\mathrm{})`$ are continuous sections of $`𝒜^{(c)}`$ defined by symmetric sequences,<sup>304</sup><sup>304</sup>304The result does not hold for all continuous sections (i.e. for all approximately symmetric sequences), since, for example, the limiting functions $`A_0`$ and $`A_0^{}`$ may not be differentiable, so that their Poisson bracket does not exist. This problem occurs in all examples of deformation quantization. However, the class of sequences for which the claim is valid is larger than the symmetric ones alone. A sufficient condition on $`A`$ and $`B`$ for (6.43) to make sense is that $`A_N=_{MN}j_{NM}(A_M^{(N)})`$ (with $`A_M^{(N)}𝒜_1^M`$), such that $`lim_N\mathrm{}A_M^{(N)}`$ exists (in norm) and $`_{M=1}^{\mathrm{}}Msup_{NM}\{A_M^{(N)}\}<\mathrm{}`$. See Duffield & Werner (1992a). then the sequence
>
> $$(\{A_0,A_0^{}\},i[A_1,A_1^{}],\mathrm{},iN[A_N,A_N^{}],\mathrm{})$$
> (6.43)
> defines a continuous section of $`𝒜^{(c)}`$.
This follows from an easy computation. In other words, although the sequence of commutators $`[A_N,A_N^{}]`$ converges to zero, the rescaled commutators $`iN[A_N,A_N^{}]𝒜_N`$ converge to the macroscopic observable $`\{A_0,A_0^{}\}𝒜_0^{(c)}=C(𝒮(𝒜_1))`$. Although it might seem perverse to reinterpret this result on the classical limit of a large quantum system in terms of quantization (which is the opposite of taking the classical limit), it is formally possible to do so (cf. Section 4.3) if we put
$$\mathrm{}=\frac{1}{N}.$$
(6.44)
Using the axiom of choice if necessary, we devise a procedure that assigns a continuous section $`A=(A_0,A_1,A_2,\mathrm{})`$ of our field to a given function $`A_0𝒜_0^{(c)}`$. We write this as $`A_N=𝒬_{\frac{1}{N}}(A_0)`$, and similarly $`A_N^{}=𝒬_{\frac{1}{N}}(A_0^{})`$. This choice need not be such that the sequence (6.43) is assigned to $`\{A_0,A_0^{}\}`$, but since the latter is the unique limit of (6.43), it must be that
$$\underset{N\mathrm{}}{lim}iN[𝒬_{\frac{1}{N}}(A_0),𝒬_{\frac{1}{N}}(A_0^{})]𝒬_{\frac{1}{N}}(\{A_0,A_0^{}\})=0.$$
(6.45)
Also note that (4.27) is just (6.12). Consequently (cf. (4.25) and surrounding text):
> The continuous field of $`C^{}`$-algebras $`𝒜^{(c)}`$ defined by (6.2) and approximately symmetric sequences (and their limits) as continuous sections yields a deformation quantization of the phase space $`𝒮(𝒜_1)`$ (equipped with the Poisson bracket (6.38)) for any quantization map $`𝒬`$.
For the dynamics this implies:
> Let $`H=(H_0,H_1,H_2,\mathrm{})`$ be a continuous section of $`𝒜^{(c)}`$ defined by a symmetric sequence,<sup>305</sup><sup>305</sup>305Once again, the result in fact holds for a larger class of Hamiltonians, namely the ones satisfying the conditions specified in footnote 304 (Duffield & Werner, 1992a). The assumption that each Hamiltonian $`H_N`$ lies in $`𝒜_1^N`$ and hence is bounded is natural in lattice models, but is undesirable in general. and let $`A=(A_0,A_1,A_2,\mathrm{})`$ be an arbitrary continuous section of $`𝒜^{(c)}`$ (i.e. an approximately symmetric sequence). Then the sequence
>
> $$(A_0(t),e^{iH_1t}A_1e^{iH_1t},\mathrm{}e^{iNH_Nt}A_Ne^{iNH_Nt},\mathrm{}),$$
> (6.46)
> where $`A_0(t)`$ is the solution of the equations of motion with classical Hamiltonian $`H_0`$,<sup>306</sup><sup>306</sup>306See (5.18) and surrounding text. defines a continuous section of $`𝒜^{(c)}`$.
In other words, for bounded symmetric sequences of Hamiltonians $`H_N`$ the quantum dynamics restricted to macroscopic observables converges to the classical dynamics with Hamiltonian $`H_0`$. Compare the positions of $`\mathrm{}`$ and $`N`$ in (5.12) and (6.46), respectively, and rejoice in the reconfirmation of (6.44).
In contrast, the quasilocal observables are not well behaved as far as the $`N\mathrm{}`$ limit of the dynamics defined by such Hamiltonians is concerned. Namely, if $`(A_0,A_1,\mathrm{})`$ is a section of the continuous field $`𝒜^{(q)}`$, and $`(H_1,H_2,\mathrm{})`$ is any bounded symmetric sequence of Hamiltonians, then the sequence
$$(e^{iH_1t}A_1e^{iH_1t},\mathrm{}e^{iNH_Nt}A_Ne^{iNH_Nt},\mathrm{})$$
has no limit for $`N\mathrm{}`$, in that it cannot be extended by some $`A_0(t)`$ to a continuous section of $`𝒜^{(q)}`$. Indeed, this was the very reason why macroscopic observables were originally introduced in this context (Rieckers, 1984; Morchio & Strocchi, 1987; Bona, 1988; Unnerstall, 1990a; Raggio & Werner, 1989; Duffield & Werner, 1992a). Instead, the natural finite-$`N`$ Hamiltonians for which the limit $`N\mathrm{}`$ of the time-evolution on $`𝒜_1^N`$ exists as a one-parameter automorphism group on $`𝒜^{(q)}`$ satisfy an appropriate locality condition, which excludes the global averages defining symmetric sequences.
### 6.6 Epilogue: Macroscopic observables and the measurement problem
In a renowned paper, Hepp (1972) suggested that macroscopic observables and superselection rules should play a role in the solution of the measurement problem of quantum mechanics. He assumed that a macroscopic apparatus may be idealized as an infinite quantum system, whose algebra of observables $`𝒜_A`$ has disjoint pure states. Referring to our discussion in Subsection 2.5 for context and notation, Hepp’s basic idea (for which he claimed no originality) was that as a consequence of the measurement process the initial state vector $`\mathrm{\Omega }_I=_nc_n\mathrm{\Psi }_nI`$ of system plus apparatus evolves into a final state vector $`\mathrm{\Omega }_F=_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n`$, in which each $`\mathrm{\Phi }_n`$ lies in a different superselection sector of the Hilbert space of the apparatus (in other words, the corresponding states $`\phi _n`$ on $`𝒜_A`$ are mutually disjoint). Consequently, although the initial state $`\omega _I`$ is pure, the final state $`\omega _F`$ is mixed. Moreover, because of the disjointness of the $`\omega _n`$ the final state $`\omega _F`$ has a unique decomposition $`\omega _F=_n|c_n|^2\psi _n\phi _n`$ into pure states, and therefore admits a bona fide ignorance interpretation. Hepp therefore claimed with some justification that the measurement “reduces the wave packet”, as desired in quantum measurement theory.
Even apart from the usual conceptual problem of passing from the collective of all terms in the final mixture to one actual measurement outcome, Hepp himself indicated a serious mathematical problem with this program. Namely, if the initial state is pure it must lie in a certain superselection sector (or equivalence class of states); but then the final state must lie in the very same sector if the time-evolution is Hamiltonian, or, more generally, automorphic (as we have seen in the preceding subsection). Alternatively, it follows from a more general lemma Hepp (1972) himself proved:
> If two states $`\rho ,\sigma `$ on a $`C^{}`$-algebra $``$ are disjoint and $`\alpha :`$ is an automorphism of $``$, then $`\rho \alpha `$ and $`\sigma \alpha `$ are disjoint, too.
To reach the negative conclusion above, one takes $``$ to be the algebra of observables of system and apparatus jointly, and computes back in time by choosing $`\alpha =\alpha _{t_Ft_I}^1`$, where $`\alpha _t`$ is the one-parameter automorphism group on $``$ describing the joint time-evolution of system and apparatus (and $`t_I`$ and $`t_F`$ are the initial and final times of the measurement, respectively). However, Hepp pointed out that this conclusion may be circumvented if one admits the possibility that a measurement takes infinitely long to complete. For the limit $`Alim_t\mathrm{}\alpha _t(A)`$ (provided it exists in a suitable sense, e.g., weakly) does not necessarily yield an automorphism of $``$. Hence a state - evolving in the Schrödinger picture by $`\omega _t(A)\omega (\alpha _t(A))`$ \- may leave its sector in infinite time, a possibility Hepp actually demonstrated in a range of models; see also Frigerio (1974), Whitten-Wolfe & Emch (1976), Araki (1980), Bona (1980), Hannabuss (1984), Bub (1988), Landsman (1991), Frasca (2003, 2004), and many other papers.
Despite the criticism that has been raised against the conclusion that a quantum-mechanical measurement requires an infinite apparatus and must take infinite time (Bell, 1975; Robinson, 1994; Landsman, 1995), and despite the fact that this procedure is quite against the spirit of von Neumann (1932), in whose widely accepted description measurements are practically instantaneous, this conclusion resonates well with the modern idea that quantum theory is universally valid and the classical world has no absolute existence; cf. the Introduction. Furthermore, a quantum-mechanical measurement is nothing but a specific interaction, comparable with a scattering process; and it is quite uncontroversial that such a process takes infinite time to complete. Indeed, what would it mean for scattering to be over after some finite time? Which time? As we shall see in the next section, the theory of decoherence requires the limit $`t\mathrm{}`$ as well, and largely for the same mathematical reasons. There as well as in Hepp’s approach, the limiting behaviour actually tends to be approached very quickly (on the pertinent time scale), and one needs to let $`t\mathrm{}`$ merely to make terms $`\mathrm{exp}\gamma t`$ (with $`\gamma >0`$) zero rather than just very small. See also Primas (1997) for a less pragmatic point of view on the significance of this limit.
A more serious problem with Hepp’s approach lies in his assumption that the time-evolution on the quasilocal algebra of observables of the infinite measurement apparatus (which in our class of examples would be $`𝒜_0^{(q)}`$) is automorphic. This, however, is by no means always the case; cf. the references listed near the end of Subsection 6.5. As we have seen, for certain natural Hamiltonian (and hence automorphic) time-evolutions at finite $`N`$ the dynamics has no limit $`N\mathrm{}`$ on the algebra of quasilocal observables \- let alone an automorphic one.
Nonetheless, Hepp’s conclusion remains valid if we use the algebra $`𝒜_0^{(c)}`$ of macroscopic observables, on which (under suitable assumptions - see Subsection 6.5) Hamiltonian time-evolution on $`𝒜_1^N`$ does have a limit as $`N\mathrm{}`$. For, as pointed out in Subsection 6.3, each superselection sector of $`𝒜_0^{(q)}`$ defines and is defined by a pure state on $`𝒜_1`$, which in turn defines a sector of $`𝒜_0^{(c)}`$. Now the latter sector is simply a point in the pure state space $`𝒮(𝒜_1)`$ of the commutative $`C^{}`$-algebra $`𝒜_0^{(c)}`$, so that Hepp’s lemma quoted above boils down to the claim that if $`\rho \sigma `$, then $`\rho \alpha \sigma \alpha `$ for any automorphism $`\alpha `$. This, of course, is a trivial property of any Hamiltonian time-evolution, and it follows once again that a transition from a pure pre-measurement state to a mixed post-measurement state on $`𝒜_0^{(c)}`$ is impossible in finite time. To avoid this conclusion, one should simply avoid the limt $`N\mathrm{}`$, which is the root of the $`t\mathrm{}`$ limit; see Janssens (2004).
What, then, does all this formalism mean for Schrödinger’s cat? In our opinion, it confirms the impression that the appearance of a paradox rests upon an equivocation. Indeed, the problem arises because one oscillates between two mutually exclusive interpretations.<sup>307</sup><sup>307</sup>307Does complementarity re-enter through the back door?
Either one is a bohemian theorist who, in vacant or in pensive mood, puts off his or her glasses and merely contemplates whether the cat is dead or alive. Such a person studies the cat exclusively from the point of view of its macroscopic observables, so that he or she has to use a post-measurement state $`\omega _F^{(c)}`$ on the algebra $`𝒜_0^{(c)}`$. If $`\omega _F^{(c)}`$ is pure, it lies in $`𝒫(𝒜_1)`$ (unless the pre-measurement state was mixed). Such a state corresponds to a single superselection sector $`[\omega _F^{(q)}]`$ of $`𝒜_0^{(q)}`$, so that the cat is dead or alive. If, on the other hand, $`\omega _F^{(c)}`$ is mixed (which is what occurs if Schrödinger has his way), there is no problem in the first place: at the level of macroscopic observables one merely has a statistical description of the cat.
Or one is a hard-working experimental physicist of formidable power, who investigates the detailed microscopic constitution of the cat. For him or her the cat is always in a pure state on $`𝒜_1^N`$ for some large $`N`$. This time the issue of life and death is not a matter of lazy observation and conclusion, but one of sheer endless experimentation and computation. From the point of view of such an observer, nothing is wrong with the cat being in a coherent superposition of two states that are actually quite close to each other microscopically - at least for the time being.
Either way, the riddle does not exist (Wittgenstein, TLP, §6.5).
## 7 Why classical states and observables?
> ‘We have found a strange footprint on the shores of the unknown. We have devised profound theories, one after another, to account for its origins. At last, we have succeeded in reconstructing the creature that made the footprint. And lo! It is our own.’ (Eddington, 1920, pp. 200–201)
The conclusion of Sections 5 and 6 is that quantum theory may give rise to classical behaviour in certain states and with respect to certain observables. For example, we have seen that in the limit $`\mathrm{}0`$ coherent states and operators of the form $`𝒬_{\mathrm{}}(f)`$, respectively, are appropriate, whereas in the limit $`N\mathrm{}`$ one should use classical states (nomen est omen!) as defined in Subsection 6.2 and macroscopic observables. If, instead, one uses superpositions of such states, or observables with the wrong limiting behaviour, no classical physics emerges. Thus the question remains why the world at large should happen to be in such states, and why we turn out to study this world with respect to the observables in question. This question found its original incarnation in the measurement problem (cf. Subsection 2.5), but this problem is really a figure-head for a much wider difficulty.
Over the last 25 years,<sup>308</sup><sup>308</sup>308Though some say the basic idea of decoherence goes back to Heisenberg and Ludwig. two profound and original answers to this question have been proposed.
### 7.1 Decoherence
The first goes under the name of decoherence. Pioneering papers include van Kampen (1954), Zeh (1970), Zurek (1981, 1982),<sup>309</sup><sup>309</sup>309See also Zurek (1991) and the subsequent debate in Physics Today (Zurek, 1993), which drew wide attention to decoherence. and Joos & Zeh (1985), and some recent reviews are Bub (1999), Auletta (2001), Joos et al. (2003), Zurek (2003), Blanchard & Olkiewicz (2003), Bacciagaluppi (2004) and Schlosshauer (2004).<sup>310</sup><sup>310</sup>310The website http://almaak.usc.edu/$``$tbrun/Data/decoherence$`_{}`$list.html contains an extensive list of references on decoherence. More references will be given in due course. The existence (and excellence) of these reviews obviates the need for a detailed treatment of decoherence in this article, all the more so since at the time of writing this approach appears to be in a transitional stage, conceptually as well as mathematically (as will be evident from what follows). Thus we depart from the layout of our earlier chapters and restrict ourselves to a few personal comments.
1. Mathematically, decoherence boils down to the idea of adding one more link to the von Neumann chain (see Subsection 2.5) beyond $`S+A`$ (i.e. the system and the apparatus). Conceptually, however, there is a major difference between decoherence and older approaches that took such a step: whereas previously (e.g., in the hands of von Neumann, London & Bauer, Wigner, etc.)<sup>311</sup><sup>311</sup>311See Wheeler & Zurek (1983). the chain converged towards the observer, in decoherence it diverges away from the observer. Namely, the third and final link is now taken to be the environment (taken in a fairly literal sense in agreement with the intuitive meaning of the word). In particular, in realistic models the environment is treated as an infinite system (necessitating the limit $`N\mathrm{}`$), which has the consequence that (in simple models where the pointer has discrete spectrum) the post-measurement state $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n\chi _n`$ (in which the $`\chi _n`$ are mutually orthogonal) is only reached in the limit $`t\mathrm{}`$. However, as already mentioned in Subsection 6.6, infinite time is only needed mathematically in order to make terms of the type $`\mathrm{exp}\gamma t`$ (with $`\gamma >0`$) zero rather than just very small: in many models the inner products $`(\chi _n,\chi _m)`$ are actually negligible for $`nm`$ within surprisingly short time scales.<sup>312</sup><sup>312</sup>312Cf. Tables 3.1 and 3.2 on pp. 66–67 of Joos et al. (2003).
If only in view of the need for limits of the type $`N\mathrm{}`$ (for the environment) and $`t\mathrm{}`$, in our opinion decoherence is best linked to stance 1 of the Introduction: its goal is to explain the approximate appearance of the classical world from quantum mechanics seen as a universally valid theory. However, decoherence has been claimed to support almost any opinion on the foundations of quantum mechanics; cf. Bacciagaluppi (2004) and Schlosshauer (2004) for a critical overview and also see Point 3 below.
2. Originally, decoherence entered the scene as a proposed solution to the measurement problem (in the precise form stated at the end of Subsection 2.5). For the restriction of the state $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n\chi _n`$ to $`S+A`$ (i.e. its trace over the degrees of freedom of the environment) is mixed in the limit $`t\mathrm{}`$, which means that the quantum-mechanical interference between the states $`\mathrm{\Psi }_n\mathrm{\Phi }_n`$ for different values of $`n`$ has become ‘delocalized’ to the environment, and accordingly is irrelevant if the latter is not observed (i.e. omitted from the description). Unfortunately, the application of the ignorance interpretation of the mixed post-measurement state of $`S+A`$ is illegal even from the point of view of stance 1 of the Introduction. The ignorance interpretation is only valid if the environment is kept within the description and is classical (in having a commutative $`C^{}`$-algebra of observables). The latter assumption (Primas, 1983), however, makes the decoherence solution to the measurement problem circular.<sup>313</sup><sup>313</sup>313On the other hand, treating the environment as if it were classical might be an improvement on the Copenhagen ideology of treating the measurement apparatus as if it were classical (cf. Section 3).
In fact, as quite rightly pointed out by Bacciagaluppi (2004), decoherence actually aggravates the measurement problem. Where previously this problem was believed to be man-made and relevant only to rather unusual laboratory situations (important as these might be for the foundations of physics), it has now become clear that “measurement” of a quantum system by the environment (instead of by an experimental physicist) happens everywhere and all the time: hence it remains even more miraculous than before that there is a single outcome after each such measurement. Thus decoherence as such does not provide a solution to the measurement problem (Leggett, 2002;<sup>314</sup><sup>314</sup>314In fact, Leggett’s argument only applies to strawman 3 of the Introduction and loses its force against stance 1. For his argument is that decoherence just removes the evidence for a given state (of Schrödinger’s cat type) to be a superposition, and accuses those claiming that this solves the measurement problem of committing the logical fallacy that removal of the evidence for a crime would undo the crime. But according to stance 1 the crime is only defined relative to the evidence! Leggett is quite right, however, in insisting on the ‘from “ and” to “or” problem’ mentioned at the end of the Introduction. Adler, 2003; Joos & Zeh, 2003), but is in actual fact parasitic on such a solution.
3. There have been various responses to this insight. The dominant one has been to combine decoherence with some interpretation of quantum mechanics: decoherence then finds a home, while conversely the interpretation in question is usually enhanced by decoherence. In this context, the most popular of these has been the many-worlds interpretation, which, after decades of obscurity and derision, suddenly started to be greeted with a flourish of trumpets in the wake of the popularity of decoherence. See, for example, Saunders (1993, 1995), Joos et al. (2003) and Zurek (2003). In quantum cosmology circles, the consistent histories approach has been a popular partner to decoherence, often in combination with many worlds; see below. The importance of decoherence in the modal interpretation has been emphasized by Dieks (1989b) and Bene & Dieks (2002), and practically all authors on decoherence find the opportunity to pay some lip-service to Bohr in one way or another. See Bacciagaluppi (2004) and Schlosshauer (2004) for a critical assessment of all these combinations.
In our opinion, none of the established interpretations of quantum mechanics will do the job, leaving room for genuinely new ideas. One such idea is the return of the environment: instead of “tracing it out”, as in the original setting of decoherence theory, the environment should not be ignored! The essence of measurement has now been recognized to be the redundancy of the outcome (or “record”) of the measurement in the environment. It is this very redundancy of information about the underlying quantum object that “objectifies” it, in that the information becomes accessible to a large number of observers without necessarily disturbing the object<sup>315</sup><sup>315</sup>315Such objectification is claimed to yield an ‘operational definition of existence’ (Zurek, 2003, p. 749.). (Zurek, 2003; Ollivier, Poulin, & Zurek, 2004; Blume-Kohout & Zurek, 2004, 2005). This insight (called “Quantum Darwinism”) has given rise to the “existential” interpretation of quantum mechanics due to Zurek (2003).
4. Another response to the failure of decoherence (and indeed all other approaches) to solve the measurement problem (in the sense of failing to win a general consensus) has been of a somewhat more pessimistic (or, some would say, pragmatic) kind: all attempts to explain the quantum world are given up, yielding to the point of view that ‘the appropriate aim of physics at the fundamental level then becomes the representation and manipulation of information’ (Bub, 2004). Here ‘measuring instruments ultimately remain black boxes at some level’, and one concludes that all efforts to understand measurement (or, for that matter, epr-correlations) are futile and pointless.<sup>316</sup><sup>316</sup>316 It is indeed in describing the transformation of quantum information (or entropy) to classical information during measurement that decoherence comes to its own and exhibits some of its greatest strength. Perhaps for this reason such thinking pervades also Zurek (2003).
5. Night thoughts of a quantum physicist, then?<sup>317</sup><sup>317</sup>317Kent, 2000. Pun on the title of McCormmach (1982). Not quite. Turning vice into virtue: rather than solving the measurement problem, the true significance of the decoherence program is that it gives conditions under which there is no measurement problem! Namely, foregoing an explanation of the transition from the state $`_nc_n\mathrm{\Psi }_n\mathrm{\Phi }_n\chi _n`$ of $`S+A+`$ to a single one of the states $`\mathrm{\Psi }_n\mathrm{\Phi }_n`$ of $`S+A`$, at the heart of decoherence is the claim that each of the latter states is robust against coupling to the environment (provided the Hamiltonian is such that $`\mathrm{\Psi }_n\mathrm{\Phi }_n`$ tensored with some initial state $`I_{}`$ of the environment indeed evolves into $`\mathrm{\Psi }_n\mathrm{\Phi }_n\chi _n`$, as assumed so far). This implies that each state $`\mathrm{\Psi }_n\mathrm{\Phi }_n`$ remains pure after coupling to the environment and subsequent restriction to the original system plus apparatus, so that at the end of the day the environment has had no influence on it. In other words, the real point of decoherence is the phenomenon of einselection (for environment-induced superselection), where a state is ‘einselected’ precisely when (given some interaction Hamiltonian) it possesses the stability property just mentioned. The claim, then, is that einselected states are often classical, or at least that classical states (in the sense mentioned at the beginning of this section) are classical precisely because they are robust against coupling to the environment. Provided this scenario indeed gives rise to the classical world (which remains to be shown in detail), it gives a dynamical explanation of it. But even short of having achieved this goal, the importance of the notion of einselection cannot be overstated; in our opinion, it is the most important and powerful idea in quantum theory since entanglement (which einselection, of course, attempts to undo!).
6. The measurement problem, and the associated distinction between system and apparatus on the one hand and environment on the other, can now be omitted from decoherence theory. Continuing the discussion in Subsection 3.4, the goal of decoherence should simply be to find the robust or einselected states of a object $`𝒪`$ coupled to an environment $``$, as well as the induced dynamics thereof (given the time-evolution of $`𝒪+`$). This search, however, must include the correct identification of the object $`𝒪`$ within the total $`𝒮+`$, namely as a subsystem that actually has such robust states. Thus the Copenhagen idea that the Heisenberg cut between object and apparatus be movable (cf. Subsection 3.2) will not, in general, extend to the “Primas–Zurek” cut between object and environment. In traditional physics terminology, the problem is to find the right “dressing” of a quantum system so as to make at least some of its states robust against coupling to its environment (Amann & Primas, 1997; Brun & Hartle, 1999; Omnès, 2002). In other words: What is a system? To mark this change in perspective, we now change notation from $`𝒪`$ (for “object”) to $`𝒮`$ (for “system”). Various tools for the solution of this problem within the decoherence program have now been developed - with increasing refinement and also increasing reliance on concepts from information theory (Zurek, 2003) - but the right setting for it seems the formalism of consistent histories, see below.
7. Various dynamical regimes haven been unearthed, each of which leads to a different class of robust states (Joos et al., 2003; Zurek, 2003; Schlosshauer, 2004). Here $`H_𝒮`$ is the system Hamiltonian, $`H_I`$ is the interaction Hamiltonian between system and environment, and $`H_{}`$ is the environment Hamiltonian. As stated, no reference to measurement, object or apparatus need be made here.
* In the regime $`H_𝒮<<H_I`$, for suitable Hamiltonians the robust states are the traditional pointer states of quantum measurement theory. This regime conforms to von Neumann’s (1932) idea that quantum measurements be almost instantaneous. If, moreover, $`H_{}<<H_I`$ as well - with or without a measurement context - then the decoherence mechanism turns out to be universal in being independent of the details of $``$ and $`H_{}`$ (Strunz, Haake, & Braun, 2003).
* If $`H_𝒮H_I`$, then (at least in models of quantum Brownian motion) the robust states are coherent states (either of the traditional Schrödinger type, or of a more general nature as defined in Subsection 5.1); see Zurek, Habib, & Paz (1993) and Zurek (2003). This case is, of course, of supreme importance for the physical relevance of the results quoted in our Section 5 above, and - if only for this reason - decoherence theory would benefit from more interaction with mathematically rigorous results on quantum stochastic analysis.<sup>318</sup><sup>318</sup>318Cf. Davies (1976), Accardi, Frigerio, & Lu (1990), Parthasarathy (1992), Streater (2000), Kümmerer (2002), Maassen (2003), etc.
* Finally, if $`H_𝒮>>H_I`$, then the robust states turn out to be eigenstates of the system Hamiltonian $`H_𝒮`$ (Paz & Zurek, 1999; Ollivier, Poulin & Zurek, 2004). In view of our discussion of such states in Subsections 5.5 and 5.6, this shows that robust states are not necessarily classical. It should be mentioned that in this context decoherence theory largely coincides with standard atomic physics, in which the atom is taken to be the system $`𝒮`$ and the radiation field plays the role of the environment $``$; see Gustafson & Sigal (2003) for a mathematically minded introductory treatment and Bach, Fröhlich, & Sigal (1998, 1999) for a full (mathematical) meal.
8. Further to the above clarification of the role of energy eigenstates, decoherence also has had important things to say about quantum chaos (Zurek, 2003; Joos et al., 2003). Referring to our discussion of wave packet revival in Subsection 2.4, we have seen that in atomic physics wave packets do not behave classically on long time scales. Perhaps surprisingly, this is even true for certain chaotic macroscopic systems: cf. the case of Hyperion mentioned in the Introduction and at the end of Subsection 5.2. Decoherence now replaces the underlying superposition by a classical probability distribution, which reflects the chaotic nature of the limiting classical dynamics. Once again, the transition from the pertinent pure state of system plus environment to a single observed system state remains clouded in mystery. But granted this transition, decoherence sheds new light on classical chaos and circumvents at least the most flagrant clashes with observation.<sup>319</sup><sup>319</sup>319It should be mentioned, though, that any successful mechanism explaining the transition from quantum to classical should have this feature, so that at the end of the day decoherence might turn out to be a red herring here.
9. Robustness and einselection form the state side or Schrödinger picture of decoherence. Of course, there should also be a corresponding observable side or Heisenberg picture of decoherence. But the transition between the two pictures is more subtle than in the quantum mechanics of closed systems. In the Schrödinger picture, the whole point of einselection is that most pure states simply disappear from the scene. This may be beautifully visualized on the example of a two-level system with Hilbert space $`_𝒮=^2`$ (Zurek, 2003). If $``$ and $``$ (cf. (6.33)) happen to be the robust vector states of the system after coupling to an appropriate environment, and if we identify the corresponding density matrices with the north-pole $`(0,0,1)B^3`$ and the south-pole $`(0,0,1)B^3`$, respectively (cf. (6.3)), then following decoherence all other states move towards the axis connecting the north- and south poles (i.e. the intersection of the $`z`$-axis with $`B^3`$) as $`t\mathrm{}`$. In the Heisenberg picture, this disappearance of all pure states except two corresponds to the reduction of the full algebra of observables $`M_2()`$ of the system to its diagonal (and hence commutative) subalgebra $``$ in the same limit. For it is only the latter algebra that contains enough elements to distinguish $``$ and $``$ without containing observables detecting interference terms between these pure states.
10. To understand this in a more abstract and general way, we recall the mathematical relationship between pure states and observables (Landsman, 1998). The passage from a $`C^{}`$-algebra $`𝒜`$ of observables of a given system to its pure states is well known: as a set, the pure state space $`𝒫(𝒜)`$ is the extreme boundary of the total state space $`𝒮(𝒜)`$ (cf. footnote 260). In order to reconstruct $`𝒜`$ from $`𝒫(𝒜)`$, the latter needs to be equipped with the structure of a transition probability space (see Subsection 6.3) through (6.27). Each element $`A𝒜`$ defines a function $`\widehat{A}`$ on $`𝒫(𝒜)`$ by $`\widehat{A}(\omega )=\omega (A)`$. Now, in the simple case that $`𝒜`$ is finite-dimensional (and hence a direct sum of matrix algebras), one can show that each function $`\widehat{A}`$ is a finite linear combination of the form $`\widehat{A}=_ip_{\omega _i}`$, where $`\omega _i𝒫(𝒜)`$ and the elementary functions $`p_\rho `$ on $`𝒫(𝒜)`$ are defined by $`p_\rho (\sigma )=p(\rho ,\sigma )`$. Conversely, each such linear combination defines a function $`\widehat{A}`$ for some $`A𝒜`$. Thus the elements of $`𝒜`$ (seen as functions on the pure state space $`𝒫(𝒜)`$) are just the transition probabilities and linear combinations thereof. The algebraic structure of $`𝒜`$ may then be reconstructed from the structure of $`𝒫(𝒜)`$ as a Poisson space with a transition probability (cf. Subsection 6.5). In this sense $`𝒫(𝒜)`$ uniquely determines the algebra of observables of which it is the pure state space. For example, the space consisting of two points with classical transition probabilities (6.31) leads to the commutative algebra $`𝒜=`$, whereas the unit two-sphere in $`^3`$ with transition probabilities (6.32) yields $`𝒜=M_2()`$.
This reconstruction procedure may be generalized to arbitrary $`C^{}`$-algebras (Landsman, 1998), and defines the precise connection between the Schrödinger picture and the Heisenberg picture that is relevant to decoherence. These pictures are equivalent, but in practice the reconstruction procedure may be difficult to carry through.
11. For this reason it is of interest to have a direct description of decoherence in the Heisenberg picture. Such a description has been developed by Blanchard & Olkiewicz (2003), partly on the basis of earlier results by Olkiewicz (1999a,b, 2000). Mathematically, their approach is more powerful than the Schrödinger picture on which most of the literature on decoherence is based. Let $`𝒜_𝒮=(_𝒮)`$ and $`𝒜_{}=(_{})`$, and assume one has a total Hamiltonian $`H`$ acting on $`_𝒮_{}`$ as well as a fixed state of the environment, represented by a density matrix $`\rho _{}`$ (often taken to be a thermal equilibrium state). If $`\rho _𝒮`$ is a density matrix on $`_𝒮`$ (so that the total state is $`\rho _𝒮\rho _{}`$), the Schrödinger picture approach to decoherence (and more generally to the quantum theory of open systems) is based on the time-evolution
$$\rho _𝒮(t)=\text{Tr}_{_{}}\left(e^{\frac{it}{\mathrm{}}H}\rho _𝒮\rho _{}e^{\frac{it}{\mathrm{}}H}\right).$$
(7.1)
The Heisenberg picture, on the other hand, is based on the associated operator time-evolution for $`A(_𝒮)`$ given by
$$A(t)=\text{Tr}_{_{}}\left(\rho _{}e^{\frac{it}{\mathrm{}}H}A1e^{\frac{it}{\mathrm{}}H}\right),$$
(7.2)
since this yields the equivalence of the Schrödinger and Heisenberg pictures expressed by
$$\text{Tr}__𝒮\left(\rho _𝒮(t)A\right)=\text{Tr}__𝒮\left(\rho _𝒮A(t)\right).$$
(7.3)
More generally, let $`𝒜_𝒮`$ and $`𝒜_{}`$ be unital $`C^{}`$-algebras with spatial tensor product $`𝒜_𝒮𝒜_{}`$, equipped with a time-evolution $`\alpha _t`$ and a fixed state $`\omega _{}`$ on $`𝒜_{}`$. This defines a conditional expectation $`P_{}:𝒜_𝒮𝒜_{}𝒜_𝒮`$ by linear and continuous extension of $`P_{}(AB)=A\omega _{}(B)`$, and consequently a reduced time-evolution $`AA(t)`$ on $`𝒜_𝒮`$ via
$$A(t)=P_{}(\alpha _t(A1)).$$
(7.4)
See, for example, Alicki & Lendi (1987); in our context, this generality is crucial for the potential emergence of continuous classical phase spaces; see below.<sup>320</sup><sup>320</sup>320For technical reasons Blanchard & Olkiewicz (2003) assume $`𝒜_𝒮`$ to be a von Neumann algebra with trivial center. Now the key point is that decoherence is described by a decomposition $`𝒜_𝒮=𝒜_𝒮^{(1)}𝒜_𝒮^{(2)}`$ as a vector space (not as a $`C^{}`$-algebra), where $`𝒜_𝒮^{(1)}`$ is a $`C^{}`$-algebra, with the property that $`lim_t\mathrm{}A(t)=0`$ (weakly) for all $`A𝒜_𝒮^{(2)}`$, whereas $`AA(t)`$ is an automorphism on $`𝒜_𝒮^{(1)}`$ for each finite $`t`$ . Consequently, $`𝒜_𝒮^{(1)}`$ is the effective algebra of observables after decoherence, and it is precisely the pure states on $`𝒜_𝒮^{(1)}`$ that are robust or einselected in the sense discussed before.
12. For example, if $`𝒜_𝒮=M_2()`$ and the states $``$ and $``$ are robust under decoherence, then $`𝒜_𝒮^{(1)}=`$ and $`𝒜_𝒮^{(2)}`$ consists of all $`2\times 2`$ matrices with zeros on the diagonal. In this example $`𝒜_𝒮^{(1)}`$ is commutative hence classical, but this may not be the case in general. But if it is, the automorphic time-evolution on $`𝒜_𝒮^{(1)}`$ induces a classical flow on its structure space, which should be shown to be Hamiltonian using the techniques of Section 6.<sup>321</sup><sup>321</sup>321Since on the assumption in the preceding footnote $`𝒜_𝒮^{(1)}`$ is a commutative von Neumann algebra one should define the structure space in an indirect way; see Blanchard & Olkiewicz (2003). In any case, there will be some sort of classical behaviour of the decohered system whenever $`𝒜_𝒮^{(1)}`$ has a nontrivial center.<sup>322</sup><sup>322</sup>322This is possible even when $`𝒜_𝒮`$ is a factor! If this center is discrete, then the induced time-evolution on it is necessarily trivial, and one has the typical measurement situation where the center in question is generated by the projections on the eigenstates of a pointer observable with discrete spectrum. This is generic for the case where $`𝒜_𝒮`$ is a type i factor. However, type ii and iii factors may give rise to continuous classical systems with nontrivial time-evolution; see Lugiewicz & Olkiewicz (2002, 2003). We cannot do justice here to the full technical details and complications involved here. But we would like to emphasize that further to quantum field theory and the theory of the thermodynamic limit, the present context of decoherence should provide important motivation for specialists in the foundations of quantum theory to learn the theory of operator algebras.<sup>323</sup><sup>323</sup>323See the references in footnote 8.
### 7.2 Consistent histories
Whilst doing so, one is well advised to work even harder and simultaneously familiarize oneself with consistent histories. This approach to quantum theory was pioneered by Griffiths (1984) and was subsequently taken up by Omnès (1992) and others. Independently, Gell-Mann and Hartle (1990, 1993) arrived at analogous ideas. Like decoherence, the consistent histories method has been the subject of lengthy reviews (Hartle, 1995) and even books (Omnès, 1994, 1999; Griffiths, 2002) by the founders. See also the reviews by Kiefer (2003) and Halliwell (2004), the critiques by Dowker & Kent (1996), Kent (1998), Bub (1999), and Bassi & Ghirardi (2000), as well as the various mathematical reformulations and reinterpretations of the consistent histories program (Isham, 1994, 1997; Isham & Linden, 1994, 1995; Isham, Linden & Schreckenberg (1994); Isham & Butterfield, 2000; Rudolph, 1996a,b, 2000; Rudolph & Wright, 1999).
The relationship between consistent histories and decoherence is somewhat peculiar: on the one hand, decoherence is a natural mechanism through which appropriate sets of histories become (approximately) consistent, but on the other hand these approaches appear to have quite different points of departure. Namely, where decoherence starts from the idea that (quantum) systems are naturally coupled to their environments and therefore have to be treated as open systems, the aim of consistent histories is to deal with closed quantum systems such as the Universe, without a priori talking about measurements or observers. However, this distinction is merely historical: as we have seen in item 6 in the previous subsection, the dividing line between a system and its environment should be seen as a dynamical entity to be drawn according to certain stability criteria, so that even in decoherence theory one should really study the system plus its environment as a whole from the outset.<sup>324</sup><sup>324</sup>324This renders the distinction between “open” and “closed” systems a bit of a red herring, as even in decoherence theory the totality of the system plus its environment is treated as a closed system. And this is precisely what consistent historians do.
As in the preceding subsection, and for exactly the same reasons, we format our treatment of consistent histories as a list of items open to discussion.
1. The starting point of the consistent histories formulation of quantum theory is conventional: one has a Hilbert space $``$, a state $`\rho `$, taken to be the initial state of the total system under consideration (realized as a density matrix on $``$) and a Hamiltonian $`H`$ (defined as a self-adjoint operator on $``$). What is unconventional is that this total system may well be the entire Universe. Each property $`\alpha `$ of the total system is mathematically represented by a projection $`P_\alpha `$ on $``$; for example, if $`\alpha `$ is the property that the energy takes some value $`ϵ`$, then the operator $`P_\alpha `$ is the projection onto the associated eigenspace (assuming $`ϵ`$ belongs to the discrete spectrum of $`H`$). In the Heisenberg picture, $`P_\alpha `$ evolves in time as $`P_\alpha (t)`$ according to (5.12); note that $`P_\alpha (t)`$ is once again a projection.
A history $`_A`$ is a chain of properties (or propositions) $`(\alpha _1(t_1),\mathrm{},\alpha _n(t_n))`$ indexed by $`n`$ different times $`t_1<\mathrm{}<t_n`$; here $`A`$ is a multi-label incorporating both the properties $`(\alpha _1,\mathrm{},\alpha _n)`$ and the times $`(t_1,\mathrm{},t_n)`$. Such a history indicates that each property $`\alpha _i`$ holds at time $`t_i`$, $`i=1,\mathrm{},n`$. Such a history may be taken to be a collection $`\{\alpha (t)\}_t`$ defined for all times, but for simplicity one usually assumes that $`\alpha (t)1`$ (where 1 is the trivial property that always holds) only for a finite set of times $`t`$; this set is precisely $`\{t_1,\mathrm{},t_n\}`$. An example suggested by Heisenberg (1927) is to take $`\alpha _i`$ to be the property that a particle moving through a Wilson cloud chamber may be found in a cell $`\mathrm{\Delta }_i^6`$ of its phase space; the history $`(\alpha _1(t_1),\mathrm{},\alpha _n(t_n))`$ then denotes the state of affairs in which the particle is in cell $`\mathrm{\Delta }_1`$ at time $`t_1`$, subsequently is in cell $`\mathrm{\Delta }_2`$ at time $`t_2`$, etcetera. Nothing is stated about the particle’s behaviour at intermediate times. Another example of a history is provided by the double slit experiment, where $`\alpha _1`$ is the particle’s launch at the source at $`t_1`$ (which is usually omitted from the description), $`\alpha _2`$ is the particle passing through (e.g.) the upper slit at $`t_2`$, and $`\alpha _3`$ is the detection of the particle at some location $`L`$ at the screen at $`t_3`$. As we all know, there is a potential problem with this history, which will be clarified below in the present framework.
The fundamental claim of the consistent historians seems to be that quantum theory should do no more (or less) than making predictions about the probabilities that histories occur. What these probabilities actually mean remains obscure (except perhaps when they are close to zero or one, or when reference is made to some measurement context; see Hartle (2005)), but let us first see when and how one can define them. The only potentially meaningful mathematical expression (within quantum mechanics) for the probability of a history $`_A`$ with respect to a state $`\rho `$ is (Groenewold, 1952; Wigner, 1963)
$$p(_A)=\text{Tr}(C_A\rho C_A^{}),$$
(7.5)
where
$$C_A=P_{\alpha _n}(t_n)\mathrm{}P_{\alpha _1}(t_1).$$
(7.6)
Note that $`C_A`$ is generally not a projection (and hence a property) itself (unless all $`P_{\alpha _i}`$ mutually commute). In particular, when $`\rho =[\mathrm{\Psi }]`$ is a pure state (defined by some unit vector $`\mathrm{\Psi }`$), one simply has
$$p(_A)=C_A\mathrm{\Psi }^2=P_{\alpha _n}(t_n)\mathrm{}P_{\alpha _1}(t_1)\mathrm{\Psi }^2.$$
(7.7)
When $`n=1`$ this just yields the Born rule. Conversely, see Isham (1994) for a derivation of (7.5) from the Born rule.<sup>325</sup><sup>325</sup>325See also Zurek (2004) for a novel derivation of the Born rule, as well as the ensuing discussion in Schlosshauer (2004).
2. Whatever one might think about the metaphysics of quantum mechanics, a probability makes no sense whatsoever when it is only attributed to a single history (except when it is exactly zero or one). The least one should have is something like a sample space (or event space) of histories, each (measurable) subset of which is assigned some probability such that the usual (Kolmogorov) rules are satisfied. This is a (well-known) problem even for a single time $`t`$ and a single projection $`P_\alpha `$ (i.e. $`n=1`$). In that case, the problem is solved by finding a self-adjoint operator $`A`$ of which $`P_\alpha `$ is a spectral projection, so that the sample space is taken to be the spectrum $`\sigma (A)`$ of $`A`$, with $`\alpha \sigma (A)`$. Given $`P_\alpha `$, the choice of $`A`$ is by no means unique, of course; different choices may lead to different and incompatible sample spaces. In practice, one usually starts from $`A`$ and derives the $`P_\alpha `$ as its spectral projections $`P_\alpha =_\alpha 𝑑P(\lambda )`$, given that the spectral resolution of $`A`$ is $`A=_{}𝑑P(\lambda )\lambda `$. Subsequently, one may then either coarse-grain or fine-grain this sample space. The former is done by finding a partition $`\sigma (A)=_i\alpha _i`$ (disjoint union), and only admitting elements of the $`\sigma `$-algebra generated by the $`\alpha _i`$ as events (along with the associated spectral projection $`P_{\alpha _i}`$), instead of all (measurable) subsets of $`\sigma (A)`$. To perform fine-graining, one supplements $`A`$ by operators that commute with $`A`$ as well as with each other, so that the new sample space is the joint spectrum of the ensuing family of mutually commuting operators.
In any case, in what follows it turns out to be convenient to work with the projections $`P_\alpha `$ instead of the subsets $`\alpha `$ of the sample space; the above discussion then amounts to extending the given projection on $``$ to some Boolean sublattice of the lattice $`𝒫()`$ of all projections on $``$.<sup>326</sup><sup>326</sup>326This sublattice is supposed to the unit of $`𝒫()`$, i.e. the unit operator on $``$, as well as the zero projection. This comment also applies to the Boolean sublattice of $`𝒫(^N)`$ discussed below. Any state $`\rho `$ then defines a probability measure on this sublattice in the usual way (Beltrametti & Cassinelli, 1984).
3. Generalizing this to the multi-time case is not a trivial task, somewhat facilitated by the following device (Isham, 1994). Put $`^N=^N`$, where $`N`$ is the cardinality of the set of all times $`t_i`$ relevant to the histories in the given collection,<sup>327</sup><sup>327</sup>327See the mathematical references above for the case $`N=\mathrm{}`$. and, for a given history $`_A`$, define
$$_A=P_{\alpha _n}(t_n)\mathrm{}P_{\alpha _1}(t_1).$$
(7.8)
Here $`P_{\alpha _i}(t_i)`$ acts on the copy of $``$ in the tensor product $`^N`$ labeled by $`t_i`$, so to speak. Note that $`_A`$ is a projection on $`^N`$ (whereas $`C_A`$ in (7.6) is generally not a projection on $``$). Furthermore, given a density matrix $`\rho `$ on $``$ as above, define the decoherence functional $`d`$ as a map from pairs of histories into $``$ by
$$d(_A,_B)=\text{Tr}(C_A\rho C_B^{}).$$
(7.9)
The main point of the consistent histories approach may now be summarized as follows: a collection $`\{_A\}_{A𝔸}`$ of histories can be regarded as a sample space on which a state $`\rho `$ defines a probability measure via (7.5), which of course amounts to
$$p(_A)=d(_A,_A),$$
(7.10)
provided that:
1. The operators $`\{_A\}_{A𝔸}`$ form a Boolean sublattice of the lattice $`𝒫(^N)`$ of all projections on $`^N`$;
2. The real part of $`d(_A,_B)`$ vanishes whenever $`_A`$ is disjoint from $`_B`$.<sup>328</sup><sup>328</sup>328This means that $`_A_B=0`$; equivalently, $`P_{\alpha _i}(t_i)P_{\beta _i}(t_i)=0`$ for at least one time $`t_i`$. This condition guarantees that the probability (7.10) is additive on disjoint histories.
In that case, the set $`\{_A\}_{A𝔸}`$ is called consistent. It is important to realize that the possible consistency of a given set of histories depends (trivially) not only on this set, but in addition on the dynamics and on the initial state.
Consistent sets of histories generalize families of commuting projections at a single time. There is no great loss in replacing the second condition by the vanishing of $`d(_A,_B)`$ itself, in which case the histories $`_A`$ and $`_B`$ are said to decohere.<sup>329</sup><sup>329</sup>329Consistent historians use this terminology in a different way from decoherence theorists. By definition, any two histories involving only a single time are consistent (or, indeed, “decohere”) iff condition (a) above holds; condition (b) is trivially satisfied in that case, and becomes relevant only when more than one time is considered. However, in decoherence theory the reduced density matrix at some given time does not trivially “decohere” at all; the whole point of the (original) decoherence program was to provide models in which this happens (if only approximately) because of the coupling of the system with its environment. Having said this, within the context of models there are close links between consistency (or decoherence) of multi-time histories and decoherence of reduced density matrices, as the former is often (approximately) achieved by the same kind of dynamical mechanisms that lead to the latter. For example, in the double slit experiment the pair of histories $`\{_A,_B\}`$ where $`\alpha _1=\beta _1`$ is the particle’s launch at the source at $`t_1`$, $`\alpha _2`$ ($`\beta _2`$) is the particle passing through the upper (lower) slit at $`t_2`$, and $`\alpha _3=\beta _3`$ is the detection of the particle at some location $`L`$ at the screen, is not consistent. It becomes consistent, however, when the particle’s passage through either one of the slits is recorded (or measured) without the recording device being included in the histories (if it is, nothing would be gained). This is reminiscent of the von Neumann chain in quantum measurement theory, which indeed provides an abstract setting for decoherence (cf. item 1 in the preceding subsection). Alternatively, the set can be made consistent by omitting $`\alpha _2`$ and $`\beta _2`$. See Griffiths (2002) for a more extensive discussion of the double slit experiment in the language of consistent histories.
More generally, coarse-graining by simply leaving out certain properties is often a promising attempt to make a given inconsistent set consistent; if the original history was already consistent, it can never become inconsistent by doing so. Fine-graining (by embedding into a larger set), on the other hand, is a dangerous act in that it may render a consistent set inconsistent.
4. What does it all mean? Each choice of a consistent set defines a “universe of discourse” within which one can apply classical probability theory and classical logic (Omnès, 1992). In this sense the consistent historians are quite faithful to the Copenhagen spirit (as most of them acknowledge): in order to understand it, the quantum world has to be looked at through classical glasses. In our opinion, no convincing case has ever been made for the absolute necessity of this Bohrian stance (cf. Subsection 3.1), but accepting it, the consistent histories approach is superior to Copenhagen in not relying on measurement as an a priori ingredient in the interpretation of quantum mechanics.<sup>330</sup><sup>330</sup>330See Hartle (2005) for an analysis of the connection between consistent histories and the Copenhagen interpretation and others. It is also more powerful than the decoherence approach in turning the notion of a system into a dynamical variable: different consistent sets describe different systems (and hence different environments, defined as the rest of the Universe); cf. item 6 in the previous subsection.<sup>331</sup><sup>331</sup>331Technically, as the commutant of the projections occurring in a given history. In other words, the choice of a consistent set boils down to a choice of “relevant variables” against “irrelevant” ones omitted from the description. As indeed stressed in the literature, the act of identification of a certain consistent set as a universe of discourse is itself nothing but a coarse-graining of the Universe as a whole.
5. But these conceptual successes come with a price tag. Firstly, consistent sets turn out not to exist in realistic models (at least if the histories in the set carry more than one time variable). This has been recognized from the beginning of the program, the response being that one has to deal with approximately consistent sets for which (the real part of) $`d(_A,_B)`$ is merely very small. Furthermore, even the definition of a history often cannot be given in terms of projections. For example, in Heisenberg’s cloud chamber example (see item 1 above), because of his very own uncertainty principle it is impossible to write down the corresponding projections $`P_{\alpha _i}`$. A natural candidate would be $`P_\alpha =𝒬_{\mathrm{}}^B(\chi _\mathrm{\Delta })`$, cf. (4.19) and (4.28), but in view of (4.21) this operator fails to satisfy $`P_\alpha ^2=P_\alpha `$, so that it is not a projection (although it does satisfy the second defining property of a projection $`P_\alpha ^{}=P_\alpha `$). This merely reflects the usual property $`𝒬(f)^2𝒬(f^2)`$ of any quantization method, and necessitates the use of approximate projections (Omnès, 1997). Indeed, this point calls for a reformulation of the entire consistent histories approach in terms of positive operators instead of projections (Rudolph, 1996a,b).
These are probably not serious problems; indeed, the recognition that classicality emerges from quantum theory only in an approximate sense (conceptually as well as mathematically) is a profound one (see the Introduction), and it rather should be counted among its blessings that the consistent histories program has so far confirmed it.
6. What is potentially more troubling is that consistency by no means implies classicality beyond the ability (within a given consistent set) to assign classical probabilities and to use classical logic. Quite to the contrary, neither Schrödinger cat states nor histories that look classical at each time but follow utterly unclassical trajectories in time are forbidden by the consistency conditions alone (Dowker & Kent, 1996). But is this a genuine problem, except to those who still believe that the earth is at the centre of the Universe and/or that humans are privileged observers? It just seems to be the case that - at least according to the consistent historians - the ontological landscape laid out by quantum theory is far more “inhuman” (or some would say “obscure”) than the one we inherited from Bohr, in the sense that most consistent sets bear no obvious relationship to the world that we observe. In attempting to make sense of these, no appeal to “complementarity” will do now: for one, the complementary pictures of the quantum world called for by Bohr were classical in a much stronger sense than generic consistent sets are, and on top of that Bohr asked us to only think about two such pictures, as opposed to the innumerable consistent sets offered to us. Our conclusion is that, much as decoherence does not solve the measurement problem but rather aggravates it (see item 314 in the preceding subsection), also consistent histories actually make the problem of interpreting quantum mechanics more difficult than it was thought to be before. In any case, it is beyond doubt that the consistent historians have significantly deepened our understanding of quantum theory - at the very least by providing a good bookkeeping device!
7. Considerable progress has been made in the task of identifying at least some (approximately) consistent sets that display (approximate) classical behaviour in the full sense of the word (Gell-Mann & Hartle, 1993; Omnès, 1992, 1997; Halliwell, 1998, 2000, 2004; Brun & Hartle, 1999; Bosse & Hartle, 2005). Indeed, in our opinion studies of this type form the main concrete outcome of the consistent histories program. The idea is to find a consistent set $`\{_A\}_{A𝔸}`$ with three decisive properties:
1. Its elements (i.e. histories) are strings of propositions with a classical interpretation;
2. Any history in the set that delineates a classical trajectory (i.e. a solution of appropriate classical equations of motion) has probability (7.10) close to unity, and any history following a classically impossible trajectory has probability close to zero;
3. The description is sufficiently coarse-grained to achieve consistency, but is sufficiently fine-grained to turn the deterministic equations of motion following from (b) into a closed system.
When these goals are met, it is in this sense (no more, no less) that the consistent histories program can claim with some justification that it has indicated (or even explained) ‘How the quantum Universe becomes classical’ (Halliwell, 2005).
Examples of propositions with a classical interpretation are quantized classical observables with a recognizable interpretation (such as the operators $`𝒬_{\mathrm{}}^B(\chi _\mathrm{\Delta })`$ mentioned in item 5), macroscopic observables of the kind studied in Subsection 6.1, and hydrodynamic variables (i.e. spatial integrals over conserved currents). These represent three different levels of classicality, which in principle are connected through mutual fine- or coarse-grainings.<sup>332</sup><sup>332</sup>332The study of these connections is relevant to the program laid out in this paper, but really belongs to classical physics per se; think of the derivation of the Navier–Stokes equations from Newton’s equations. The first are sufficiently coarse-grained to achieve consistency only in the limit $`\mathrm{}0`$ (cf. Section 5), whereas the latter two are already coarse-grained by their very nature. Even so, also the initial state will have to be “classical” in some sense in order te achieve the three targets (a) - (c).
All this is quite impressive, but we would like to state our opinion that neither decoherence nor consistent histories can stand on their own in explaining the appearance of the classical world. Promising as these approaches are, they have to be combined at least with limiting techniques of the type described in Sections 5 and 6 \- not to speak of the need for a new metaphysics! For even if it is granted that decoherence yields the disappearance of superpositions of Schrödinger cat type, or that consistent historians give us consistent sets none of whose elements contain such superpositions among their properties, this by no means suffices to explain the emergence of classical phase spaces and flows thereon determined by classical equations of motion. Since so far the approaches cited in Sections 5 and 6 have hardly been combined with the decoherence and/or the consistent histories program, a full explanation of the classical world from quantum theory is still in its infancy. This is not merely true at the technical level, but also conceptually; what has been done so far only represents a modest beginning. On the positive side, here lies an attractive challenge for mathematically minded researchers in the foundations of physics!
## 8 Epilogue
As a sobering closing note, one should not forget that whatever one’s achievements in identifying a “classical realm” in quantum mechanics, the theory continues to incorporate another realm, the pure quantum world, that the young Heisenberg first gained access to, if not through his mathematics, then perhaps through the music of his favourite composer, Beethoven. This world beyond ken has never been better described than by Hoffmann (1810) in his essay on Beethoven’s instrumental music, and we find it appropriate to end this paper by quoting at some length from it:<sup>333</sup><sup>333</sup>333Translation copyright: Ingrid Schwaegermann (2001).
> Should one, whenever music is discussed as an independent art, not always be referred to instrumental music which, refusing the help of any other art (of poetry), expresses the unique essence of art that can only be recognized in it? It is the most romantic of all arts, one would almost want to say, the only truly romantic one, for only the infinite is its source. Orpheus’ lyre opened the gates of the underworld. Music opens to man an unknown realm, a world that has nothing in common with the outer sensual world that surrounds him, a realm in which he leaves behind all of his feelings of certainty, in order to abandon himself to an unspeakable longing. (…)
>
> Beethoven’s instrumental music opens to us the realm of the gigantic and unfathomable. Glowing rays of light shoot through the dark night of this realm, and we see gigantic shadows swaying back and forth, encircling us closer and closer, destroying us (…) Beethoven’s music moves the levers of fear, of shudder, of horror, of pain and thus awakens that infinite longing that is the essence of romanticism. Therefore, he is a purely romantic composer, and may it not be because of it, that to him, vocal music that does not allow for the character of infinite longing - but, through words, achieves certain effects, as they are not present in the realm of the infinite - is harder? (…)
>
> What instrumental work of Beethoven confirms this to a higher degree than his magnificent and profound Symphony in c-Minor. Irresistibly, this wonderful composition leads its listeners in an increasing climax towards the realm of the spirits and the infinite. (…) Only that composer truly penetrates into the secrets of harmony who is able to have an effect on human emotions through them; to him, relationships of numbers, which, to the Grammarian, must remain dead and stiff mathematical examples without genius, are magic potions from which he lets a miraculous world emerge. (…)
>
> Instrumental music, wherever it wants to only work through itself and not perhaps for a certain dramatic purpose, has to avoid all unimportant punning, all dallying. It seeks out the deep mind for premonitions of joy that, more beautiful and wonderful than those of this limited world, have come to us from an unknown country, and spark an inner, wonderful flame in our chests, a higher expression than mere words - that are only of this earth - can spark.
## 9 References
* Abraham, R. & Marsden, J.E. (1985). Foundations of Mechanics, 2nd ed. Addison Wesley, Redwood City.
* Accardi, L., Frigerio, A., & Lu, Y. (1990). The weak coupling limit as a quantum functional central limit. Communications in Mathematical Physics 131, 537–570.
* Adler, S.L. (2003). Why decoherence has not solved the measurement problem: A response to P.W. Anderson. Studies in History and Philosophy of Modern Physics 34B, 135–142.
* Agmon, S. (1982). Lectures on Exponential Decay of Solutions of Second-Order Elliptic Equations. Princeton: Princeton University Press.
* Albeverio, S.A. & Høegh-Krohn, R.J. (1976). Mathematical Theory of Feynman Path Integrals. Berlin: Springer-Verlag.
* Alfsen, E.M. (1970). Compact Convex Sets and Boundary Integrals. Berlin: Springer.
* Ali, S.T., Antoine, J.-P., Gazeau, J.-P. & Mueller, U.A. (1995). Coherent states and their generalizations: a mathematical overview. Reviews in Mathematical Physics 7, 1013–1104.
* Ali, S.T., Antoine, J.-P., & Gazeau, J.-P. (2000). Coherent States, Wavelets and their Generalizations. New York: Springer-Verlag.
* Ali, S.T. & Emch, G.G. (1986). Geometric quantization: modular reduction theory and coherent states. Journal of Mathematical Physics 27, 2936–2943.
* Ali, S.T & Englis, M. (2004). Quantization methods: a guide for physicists and analysts.
arXiv:math-ph/0405065.
* Alicki, A. & Fannes, M. (2001). Quantum Dynamical Systems. Oxford: Oxford University Press.
* Alicki, A. & Lendi, K. (1987). Quantum Dynamical Semigroups and Applications. Berlin: Springer.
* Amann, A. (1986). Observables in $`W^{}`$-algebraic quantum mechanics. Fortschritte der Physik 34, 167–215.
* Amann, A. (1987). Broken symmetry and the generation of classical observables in large systems. Helvetica Physica Acta 60, 384–393.
* Amann, A. & Primas, H. (1997). What is the referent of a non-pure quantum state? Experimental Metaphysics: Quantum Mechanical Studies in Honor of Abner Shimony, S. Cohen, R.S., Horne, M.A., & Stachel, J. (Eds.). Dordrecht: Kluwer Academic Publishers.
* Arai, T. (1995). Some extensions of the semiclassical limit $`\mathrm{}0`$ for Wigner functions on phase space. Journal of Mathematical Physics 36, 622–630.
* Araki, H. (1980). A remark on the Machida-Namiki theory of measurement. Progress in Theoretical Physics 64, 719–730.
* Araki, H. (1999). Mathematical Theory of Quantum Fields. New York: Oxford University Press.
* Arnold, V.I. (1989). Mathematical Methods of Classical Mechanics. Second edition. New York: Springer-Verlag.
* Ashtekar, A. & Schilling, T.A. (1999). Geometrical formulation of quantum mechanics. On Einstein’s Path (New York, 1996), pp. 23–65. New York: Springer.
* Atmanspacher, H., Amann, A., & Müller-Herold, U. (Eds.). (1999). On Quanta, Mind and Matter: Hans Primas in Context. Dordrecht: Kluwer Academic Publishers.
* Auletta, G. (2001). Foundations and Interpretation of Quantum Mechanics. Singapore: World Scientific.
* Bacciagaluppi, G. (1993). Separation theorems and Bell inequalities in algebraic quantum mechanics. Proceedings of the Symposium on the Foundations of Modern Physics (Cologne, 1993), pp. 29–37. Busch, P., Lahti, P.J., & Mittelstaedt, P. (Eds.). Singapore: World Scientific.
* Bacciagaluppi, G. (2004). The Role of Decoherence in Quantum Theory. Stanford Encyclopedia of Philosophy, (Winter 2004 Edition), Zalta, E.N. (Ed.). Online only at
http://plato.stanford.edu/archives/win2004/entries/qm-decoherence/.
* Bach, V., Fröhlich, J., & Sigal, I.M. (1998). Quantum electrodynamics of confined nonrelativistic particles. Advances in Mathematics 137, 299–395.
* Bach, V., Fröhlich, J., & Sigal, I.M. (1999). Spectral analysis for systems of atoms and molecules coupled to the quantized radiation field. Communications in Mathematical Physics 207, 249–290.
* Baez, J. (1987). Bell’s inequality for $`C^{}`$-algebras. Letters in Mathematical Physics 13, 135–136.
* Bagarello, F. & Morchio, G. (1992). Dynamics of mean-field spin models from basic results in abstract differential equations. Journal of Statistical Physics 66, 849–866.
* Ballentine, L.E. (1970). The statistical interpretation of quantum mechanics. Reviews of Modern Physics 42, 358–381.
* Ballentine, L.E. (1986). Probability theory in quantum mechanics. American Journal of Physics 54, 883–889.
* Ballentine, L.E. (2002). Dynamics of quantum-classical differences for chaotic systems. Physical Review A65, 062110-1–6.
* Ballentine, L.E. (2003). The classical limit of quantum mechanics and its implications for the foundations of quantum mechanics. Quantum Theory: Reconsideration of Foundations – 2, pp. 71–82. Khrennikov, A. (Ed.). Växjö: Växjö University Press.
* Ballentine, L.E., Yang, Y. & Zibin, J.P. (1994). Inadequacy of Ehrenfest’s theorem to characterize the classical regime. Physical Review A50, 2854–2859.
* Balian, R. & Bloch, C. (1972). Distribution of eigenfrequencies for the wave equation in a finite domain. III. Eigenfrequency density oscillations. Annals of Physics 69, 76–160.
* Balian, R. & Bloch, C. (1974). Solution of the Schrödinger equation in terms of classical paths. Annals of Physics 85, 514–545.
* Bambusi, D., Graffi, S., & Paul, T. (1999). Long time semiclassical approximation of quantum flows: a proof of the Ehrenfest time. Asymptotic Analysis 21, 149–160.
* Barrow-Green, J. (1997). Poincaré and the Three Body Problem. Providence, RI: (American Mathematical Society.
* Barut, A.O. & Raçka, R. (1977). Theory of Group Representations and Applications. Warszawa: PWN.
* Bassi, A. & Ghirardi, G.C. (2000). Decoherent histories and realism. Journal of Statistical Physics 98, 457–494. Reply by Griffiths, R.B. (2000). ibid. 99, 1409–1425. Reply to this reply by Bassi, A. & Ghirardi, G.C. (2000). ibid. 99, 1427.
* Bates, S. & Weinstein, A. (1995). Lectures on the Geometry of Quantization. Berkeley Mathematics Lecture Notes 8. University of California, Berkeley. Re-issued by the American Mathematical Society.
* Batterman, R.W. (2002). The Devil in the Details: Asymptotic Reasoning in Explanation, Reduction, and Emergence. Oxford: Oxford University Press.
* Batterman, R.W. (2005). Critical phenomena and breaking drops: Infinite idealizations in physics. Studies in History and Philosophy of Modern Physics 36, 225–244.
* Baum, P., Connes, A. & Higson, N. (1994). Classifying space for proper actions and K-theory of group $`C^{}`$-algebras. Contemporary Mathematics 167, 241–291.
* Bayen, F., Flato, M., Fronsdal, C., Lichnerowicz, A. & Sternheimer, D. (1978). Deformation theory and quantization I, II. Annals of Physics 110, 61–110, 111–151.
* Bell, J.S. (1975). On wave packet reduction in the Coleman–Hepp model. Helvetica Physica Acta 48, 93–98.
* Bell, J.S. (1987). Speakable and Unspeakable in Quantum Mechanics. Cambridge: Cambridge University Press.
* Bell, J.S. (2001). John S. Bell on the Foundations of Quantum Mechanics. Singapore: World Scientific.
* Beller, M. (1999). Quantum Dialogue. Chicago: University of Chicago Press.
* Bellissard, J. & Vittot, M. (1990). Heisenberg’s picture and noncommutative geometry of the semiclassical limit in quantum mechanics. Annales de l’ Institut Henri Poincaré - Physique Théorique 52, 175–235.
* Belot, G. (2005). Mechanics: geometrical. This volume.
* Belot, G. & Earman, J. (1997). Chaos out of order: quantum mechanics, the correspondence principle and chaos. Studies in History and Philosophy of Modern Physics 28B, 147–182.
* Beltrametti, E.G. & Cassinelli, G. (1984). The Logic of Quantum Mechanics. Cambridge: Cambridge University Press.
* Benatti, F. (1993). Deterministic Chaos in Infinite Quantum Systems. Berlin: Springer-Verlag.
* Bene, G. & Dieks, D. (2002). A Perspectival Version of the Modal Interpretation of Quantum Mechanics and the Origin of Macroscopic Behavior. Foundations of Physics 32, 645-671.
* Berezin, F.A. (1974). Quantization. Mathematical USSR Izvestia 8, 1109–1163.
* Berezin, F.A. (1975a). Quantization in complex symmetric spaces. Mathematical USSR Izvestia 9, 341–379.
* Berezin, F.A. (1975b). General concept of quantization. Communications in Mathematical Physics 40, 153–174.
* Berry, M.V. (1977a). Semi-classical mechanics in phase space: a study of Wigner’s function. Philosophical Transactions of the Royal Society 287, 237–271.
* Berry, M.V. (1977b). Regular and irregular semi-classical wavefunctions. Journal of Physics A10, 2083–2091.
* Berry, M.V., Balazs, N.L., Tabor, M., & Voros, A. (1979). Quantum maps. Annals of Physics 122, 26–63.
* Berry, M.V. & Tabor, M. (1977). Level clustering in the regular spectrum. Proceedings of the Royal Society A356, 375–394.
* Berry, M.V. & Keating, J.P. (1999). The Riemann zeros and eigenvalue asymptotics. SIAM Review 41, 236–266.
* Binz, E., J., Śniatycki, J. & Fischer, H. (1988). The Geometry of Classical Fields. Amsterdam: North–Holland.
* Birkhoff, G. & von Neumann, J. (1936). The logic of quantum mechanics. Annals of Mathematics (2) 37, 823–843.
* Bitbol, M. (1996). Schrödinger’s Philosophy of Quantum Mechanics. Dordrecht: Kluwer Academic Publishers.
* Bitbol, M. & Darrigol, O. (Eds.) (1992). Erwin Schrödinger: Philosophy and the Birth of Quantum Mechanics. Dordrecht: Kluwer Academic Publishers.
* Blackadar, B. (1998). $`K`$-Theory for Operator Algebras. Second edition. Cambridge: Cambridge University Press.
* Blair Bolles, E. (2004). Einstein Defiant: Genius versus Genius in the Quantum Revolution. Washington: Joseph Henry Press.
* Blanchard, E. (1996). Deformations de $`C^{}`$-algebras de Hopf. Bulletin de la Société mathématique de France 124, 141–215.
* Blanchard, Ph. & Olkiewicz, R. (2003). Decoherence induced transition from quantum to classical dynamics. Reviews in Mathematical Physics 15, 217–243.
* Bohigas, O., Giannoni, M.-J., & Schmit, C. (1984). Characterization of chaotic quantum spectra and universality of level fluctuation laws. Physical Review Letters 52, 1–4.
* Blume-Kohout, R. & Zurek, W.H. (2004). A simple example of ”Quantum Darwinism”: Redundant information storage in many-spin environments. Foundations of Physics, to appear.
arXiv:quant-ph/0408147.
* Blume-Kohout, R. & Zurek, W.H. (2005). Quantum Darwinism: Entanglement, branches, and the emergent classicality of redundantly stored quantum information. Physical Review A, to appear. arXiv:quant-ph/0505031.
* Bohr, N. (1927) The quantum postulate and the recent development of atomic theory. Atti del Congress Internazionale dei Fisici (Como, 1927). Reprinted in Bohr (1934), pp. 52–91.
* Bohr, N. (1934). Atomic Theory and the Description of Nature. Cambridge: Cambridge University Press.
* Bohr, N. (1935). Can quantum-mechanical description of physical reality be considered complete? Physical Review 48, 696–702.
* Bohr, N. (1937). Causality and complementarity. Philosophy of Science 4, 289–298.
* Bohr, N. (1949). Discussion with Einstein on epistemological problems in atomic physics. Albert Einstein: Philosopher-Scientist, pp. 201–241. P.A. Schlipp (Ed.). La Salle: Open Court.
* Bohr, N. (1958). Atomic Physics and Human Knowlegde. New York: Wiley.
* Bohr, N. (1985). Collected Works. Vol. 6: Foundations of Quantum Physics i (1926–1932). Kalckar, J. (Ed.). Amsterdam: North-Holland.
* Bohr, N. (1996). Collected Works. Vol. 7: Foundations of Quantum Physics ii (1933–1958). Kalckar, J. (Ed.). Amsterdam: North-Holland.
* Bogoliubov, N.N. (1958). On a new method in the theory of superconductivity. Nuovo Cimento 7, 794–805.
* Bona, P. (1980). A solvable model of particle detection in quantum theory. Acta Facultatis Rerum Naturalium Universitatis Comenianae Physica XX, 65–94.
* Bona, P. (1988). The dynamics of a class of mean-field theories. Journal of Mathematical Physics 29, 2223–2235.
* Bona, P. (1989). Equilibrium states of a class of mean-field theories. Journal of Mathematical Physics 30, 2994–3007.
* Bona, P. (2000). Extended quantum mechanics. Acta Physica Slovaca 50, 1–198.
* Bonechi, F. & De Bièvre, S. (2000). Exponential mixing and $`\mathrm{ln}\mathrm{}`$ time scales in quantized hyperbolic maps on the torus. Communications in Mathematical Physics 211, 659–686.
* Bosse, A.W. & Hartle, J.B. (2005). Representations of spacetime alternatives and their classical limits. arXiv:quant-ph/0503182.
* Brack, M. & Bhaduri, R.K. Semiclassical Physics. Boulder: Westview Press.
* Bratteli, O. & Robinson, D.W. (1987). Operator Algebras and Quantum Statistical Mechanics. Vol. I: $`C^{}`$\- and $`W^{}`$-Algebras, Symmetry Groups, Decomposition of States. 2nd Ed. Berlin: Springer.
* Bratteli, O. & Robinson, D.W. (1981). Operator Algebras and Quantum Statistical Mechanics. Vol. II: Equilibrium States, Models in Statistical Mechanics. Berlin: Springer.
* Brezger, B., Hackermüller, L., Uttenthaler, S., Petschinka, J., Arndt, M., & Zeilinger, A. (2002). Matter-Wave Interferometer for Large Molecules. Physical Review Letters 88, 100404.
* Breuer T. (1994). Classical Observables, Measurement, and Quantum Mechanics. Ph.D. Thesis, University of Cambridge.
* Bröcker, T. & Werner, R.F. (1995). Mixed states with positive Wigner functions. Journal of Mathematical Physics 36, 62–75.
* Brun, T.A. & Hartle, J.B. (1999). Classical dynamics of the quantum harmonic chain. Physical Review D60, 123503-1–20.
* Brush, S.G. (2002). Cautious revolutionaries: Maxwell, Planck, Hubble. American Journal of Physics 70, 119- 127.
* Bub, J. (1988). How to Solve the Measurement Problem of Quantum Mechanics. Foundations of Physics 18, 701–722.
* Bub, J. (1999). Interpreting the Quantum World. Cambridge: Cambridge University Press.
* Bub, J. (2004). Why the quantum? Studies in History and Philosophy of Modern Physics 35B, 241–266.
* Busch, P., Grabowski, M. & Lahti, P.J. (1998). Operational Quantum Physics, 2nd corrected ed. Berlin: Springer.
* Busch, P., Lahti, P.J., & Mittelstaedt, P. (1991). The Quantum Theory of Measurement. Berlin: Springer.
* Butterfield, J. (2002). Some Worlds of Quantum Theory. R.Russell, J. Polkinghorne et al (Ed.). Quantum Mechanics (Scientific Perspectives on Divine Action vol 5), pp. 111-140. Rome: Vatican Observatory Publications, 2. arXiv:quant-ph/0105052; PITT-PHIL-SCI00000204.
* Butterfield, J. (2005). On symmetry, conserved quantities and symplectic reduction in classical mechanics. This volume.
* Büttner, L., Renn, J., & Schemmel, M. (2003). Exploring the limits of classical physics: Planck, Einstein, and the structure of a scientific revolution. Studies in History and Philosophy of Modern Physics 34B, 37–60.
* Camilleri, K. (2005). Heisenberg and Quantum Mechanics: The Evolution of a Philosophy of Nature. Ph.D. Thesis, University of Melbourne.
* Cantoni, V. (1975). Generalized “transition probability”. Communications in Mathematical Physics 44, 125–128.
* Cantoni, V. (1977). The Riemannian structure on the states of quantum-like systems. Communications in Mathematical Physics 56, 189–193.
* Carson, C. (2000). Continuities and discontinuities in Planck’s Akt der Verzweiflung. Annalen der Physik 9, 851–960.
* Cassidy, D.C. (1992). Uncertainty: the Life and Science of Werner Heisenberg. New York: Freeman.
* Castrigiano, D.P.L. & Henrichs, R.W. (1980). Systems of covariance and subrepresentations of induced representations. Letters in Mathematical Physics 4, 169-175.
* Cattaneo, U. (1979). On Mackey’s imprimitivity theorem. Commentari Mathematici Helvetici 54, 629-641.
* Caves, C.M., Fuchs, C.A., & Schack, R. (2002). Unknown quantum states: the quantum de Finetti representation. Quantum information theory. Journal of Mathematical Physics 43, 4537–4559.
* Charbonnel, A.M. (1986). Localisation et développement asymptotique des éléments du spectre conjoint d’ opérateurs psuedodifférentiels qui commutent. Integral Equations Operator Theory 9, 502–536.
* Charbonnel, A.M. (1988). Comportement semi-classiques du spectre conjoint d’ opérateurs psuedodifférentiels qui commutent. Asymptotic Analysis 1, 227–261.
* Charbonnel, A.M. (1992). Comportement semi-classiques des systèmes ergodiques. Annales de l’ Institut Henri Poincaré - Physique Théorique 56, 187–214.
* Chernoff, P.R. (1973). Essential self–adjointness of powers of generators of hyperbolic equations. Journal of Functional Analysis 12, 401–414.
* Chernoff, P.R. (1995). Irreducible representations of infinite dimensional transformation groups and Lie algebras I. Journal of Functional Analysis 130, 255–282.
* Chevalley, C. (1991). Introduction: Le dessin et la couleur. Niels Bohr: Physique Atomique et Connaissance Humaine. (French translation of Bohr, 1958). Bauer, E. & Omnès, R. (Eds.), pp. 17–140. Paris: Gallimard.
* Chevalley, C. (1999). Why do we find Bohr obscure? Epistemological and Experimental Perspectives on Quantum Physics, pp. 59–74. Greenberger, D., Reiter, W.L., & Zeilinger, A. (Eds.). Dordrecht: Kluwer Academic Publishers.
* Chiorescu, I., Nakamura, Y., Harmans, C.J.P.M., & Mooij, J.E. (2003). Coherent Quantum Dynamics of a Superconducting Flux Qubit. Science 299, Issue 5614, 1869–1871.
* Cirelli, R., Lanzavecchia, P., & Mania, A. (1983). Normal pure states of the von Neumann algebra of bounded operators as a Kähler manifold. Journal of Physics A16, 3829–3835.
* Cirelli, R., Maniá, A., & Pizzocchero, L. (1990). Quantum mechanics as an infinite-dimensional Hamiltonian system with uncertainty structure. I, II. Journal of Mathematical Physics 31, 2891–2897, 2898–2903.
* Colin de Verdière, Y. (1973). Spectre du laplacien et longueurs des géodésiques périodiques. I, II. Compositio Mathematica 27, 83–106, 159–184.
* Colin de Verdière, Y. (1977). Quasi-modes sur les variétés Riemanniennes. Inventiones Mathematicae 43, 15–52.
* Colin de Verdière, Y. (1985). Ergodicité et fonctions propres du Laplacien. Communications in Mathematical Physics 102, 497–502.
* Colin de Verdière, Y. (1998). Une introduction la m canique semi-classique. l’Enseignement Mathematique (2) 44, 23–51.
* Combescure, M. (1992). The squeezed state approach of the semiclassical limit of the time-dependent Schrödinger equation. Journal of Mathematical Physics 33, 3870–3880.
* Combescure, M., Ralston, J., & Robert, D. (1999). A proof of the Gutzwiller semiclassical trace formula using coherent states decomposition. Communications in Mathematical Physics 202, 463–480.
* Combescure, M. & Robert, D. (1997). Semiclassical spreading of quantum wave packets and applications near unstable fixed points of the classical flow. Asymptotic Analysis 14, 377–404.
* Corwin, L. & Greenleaf, F.P. (1989). Representations of Nilpotent Lie Groups and Their Applications, Part I. Cambridge: Cambridge University Press.
* Cucchietti, F.M. (2004). The Loschmidt Echo in Classically Chaotic Systems: Quantum Chaos, Irreversibility and Decoherence. Ph.D. Thesis, Universidad Nacional de Córdobo. arXiv:quant-ph/0410121.
* Cushing, J.T. (1994). Quantum Mechanics: Historical Contingency and the Copenhagen Hegemony. Chicago: University of Chicago Press.
* Cvitanovic, P. et al. (2005). Classical and Quantum Chaos. http://ChaosBook.org.
* Cycon, H. L., Froese, R. G., Kirsch, W., & Simon, B. (1987). Schrödinger Operators with Application to Quantum Mechanics and Global Geometry. Berlin: Springer-Verlag.
* Darrigol, O. (1992). From c-Numbers to q-Numbers. Berkeley: University of California Press.
* Darrigol, O. (2001). The Historians’ Disagreements over the Meaning of Planck’s Quantum. Centaurus 43, 219–239.
* Davidson, D. (2001). Subjective, Intersubjective, Objective. Oxford: Clarendon Press
* Davies, E.B. (1976). Quantum Theory of Open Systems. London: Academic Press.
* De Bièvre, S. (1992). Oscillator eigenstates concentrated on classical trajectories. Journal of Physics A25, 3399-3418.
* De Bièvre, S. (2001). Quantum chaos: a brief first visit. Contemporary Mathematics 289, 161–218.
* De Bièvre, S. (2003). Local states of free bose fields. Lectures given at the Summer School on Large Coulomb Systems, Nordfjordeid.
* De Bièvre, S., Irac-Astaud, M, & Houard, J.C. (1993). Wave packets localised on closed classical trajectories. Differential Equations and Applications in Mathematical Physics, pp. 25–33. Ames, W.F. & Harrell, E.M., & Herod, J.V. (Eds.). New York: Academic Press.
* De Muynck, W.M. (2002) Foundations of Quantum Mechanics: an Empiricist Approach. Dordrecht: Kluwer Academic Publishers.
* Devoret, M.H., Wallraff, A., & Martinis, J.M. (2004). Superconducting Qubits: A Short Review. arXiv:cond-mat/0411174.
* Diacu, F. & Holmes, P. (1996). Celestial Encounters. The Origins of Chaos and Stability. Princeton: Princeton University Press.
* Dickson, M. (2005). Non-relativistic quantum mechanics. This Volume.
* Dieks, D. (1989a). Quantum mechanics without the projection postulate and its realistic interpretation. Foundations of Physics 19, 1397–1423.
* Dieks, D. (1989b). Resolution of the measurement problem through decoherence of the quantum state. Physics Letters 142A, 439–446.
* Dimassi, M. & Sjöstrand, J. (1999). Spectral Asymptotics in the Semi-Classical Limit. Cambridge: Cambridge University Press.
* Dirac, P.A.M. (1926). The fundamental equations of quantum mechanics. Proceedings of the Royal Society A109, 642–653.
* Dirac, P.A.M. (1930). The Principles of Quantum Mechanics. Oxford: Clarendon Press.
* Dirac, P.A.M. (1964). Lectures on Quantum Mechanics. New York: Belfer School of Science, Yeshiva University.
* Dixmier, J. (1977). $`C^{}`$-Algebras. Amsterdam: North–Holland.
* Doebner, H.D. & J. Tolar (1975). Quantum mechanics on homogeneous spaces. Journal of Mathematical Physics 16, 975–984.
* Dowker, F. & Kent, A. (1996). On the Consistent Histories Approach to Quantum Mechanics. Journal of Statistical Physics 82, 1575–1646.
* Dubin, D.A., Hennings, M.A., & Smith, T.B. (2000). Mathematical Aspects of Weyl Quantization and Phase. Singapore: World Scientific.
* Duclos, P. & Hogreve, H. (1993). On the semiclassical localization of the quantum probability. Journal of Mathematical Physics 34, 1681–1691.
* Duffield, N.G. (1990). Classical and thermodynamic limits for generalized quantum spin systems. Communications in Mathematical Physics 127, 27–39.
* Duffield, N.G. & Werner, R.F. (1992a). Classical Hamiltonian dynamics for quantum Hamiltonian mean-field limits. Stochastics and Quantum Mechanics: Swansea, Summer 1990, pp. 115–129. Truman, A. & Davies, I.M. (Eds.). Singapore: World Scientific.
* Duffield, N.G. & Werner, R.F. (1992b). On mean-field dynamical semigroups on $`C^{}`$-algebras. Reviews in Mathematical Physics 4, 383–424.
* Duffield, N.G. & Werner, R.F. (1992c). Local dynamics of mean-field quantum systems. Helvetica Physica Acta 65, 1016–1054.
* Duffield, N.G., Roos, H., & Werner, R.F. (1992). Macroscopic limiting dynamics of a class of inhomogeneous mean field quantum systems. Annales de l’ Institut Henri Poincaré - Physique Théorique 56, 143–186.
* Duffner, E. & Rieckers, A. (1988). On the global quantum dynamics of multilattice systems with nonlinear classical effects. Zeitschrift für Naturforschung A43, 521–532.
* Duistermaat, J.J. (1974). Oscillatory integrals, Lagrange immersions and unfolding of singularities. Communications in Pure and Applied Mathematics 27, 207–281.
* Duistermaat, J.J. & Guillemin, V. (1975). The spectrum of positive elliptic operators and periodic bicharacteristics. Inventiones Mathematicae 29, 39–79.
* Duistermaat, J.J. (1996). Fourier Integral Operators. Original Lecture Notes from 1973. Basel: Birkhäuser.
* Duval, C., Elhadad, J., Gotay, M.J., Śniatycki, J., & Tuynman, G.M. (1991). Quantization and bosonic BRST theory. Annals of Physics 206, 1–26.
* Earman, J. (1986). A Primer on Determinism. Dordrecht: Reidel.
* Earman, J. (2005). Aspects of determinism in modern physics. This volume.
* Earman, J. (2006). Essential self-adjointness: implications for determinism and the classical-quantum correspondence. Synthese, to appear.
* Echterhoff, S., Kaliszewski, S., Quigg, J., & Raeburn, I. (2002). A categorical approach to imprimitivity theorems for C\*-dynamical systems. arXiv:math.OA/0205322.
* Eddington, A.S. (1920). Space, Time, and Gravitation: An Outline of the General Relativity Theory. Cambridge: Cambridge University Press.
* Effros, E.G. & Hahn, F. (1967). Locally compact transformation groups and $`C^{}`$\- algebras. Memoirs of the American Mathematical Society 75.
* Ehrenfest, P. (1927). Bemerkung über die angenäherte Gultigkeit der klassischen Mechanik innerhalb der Quantenmechanik. Zeitschrift für Physik 45, 455–457.
* Einstein, A. (1905). Über einen die Erzeugung und Verwandlung des Lichtes betreffenden heuristischen Gesichtpunkt. Annalen der Physik 17, 132–178.
* Einstein, A. (1917). Zum Quantensatz von Sommerfeld und Epstein. Verhandlungen der deutschen Physikalischen Geselschaft (2) 19, 82–92.
* Einstein, A. (1949). Remarks to the essays appearing in this collective volume. (Reply to criticisms). Albert Einstein: Philosopher-Scientist, pp. 663–688. Schilpp, P.A. (Ed.). La Salle: Open Court.
* Emch, G.G. & Knops, H.J.F. (1970). Pure thermodynamical phases as extremal KMS states. Journal of Mathematical Physics 11, 3008–3018.
* Emch, G.G. (1984) Mathematical and conceptual foundations of 20th-century physics. Amsterdam: North-Holland.
* Emch, G.G. & Liu, C. (2002). The Logic of Thermostatistical Physics. Berlin: Springer-Verlag.
* Emch, G.G., Narnhofer, H., Thirring, W., & Sewell, G. (1994). Anosov actions on noncommutative algebras. Journal of Mathematical Physics 35, 5582–5599.
* d’Espagnat, B. (1995). Veiled Reality: An Analysis of Present-Day Quantum Mechanical Concepts. Reading (MA): Addison-Wesley.
* Esposito, G., Marmo, G., & Sudarshan, G. (2004). From Classical to Quantum Mechanics: An Introduction to the Formalism, Foundations and Applications. Cambridge: Cambridge University Press.
* Enz, C.P. (2002). No Time to be Brief: A Scientific Biography of Wolfgang Pauli. Oxford: Oxford University Press.
* Everett, H. iii (1957). “Relative state” formulation of quantum mechanics. Reviews in Modern Physics 29, 454–462.
* Faye, J. (1991). Niels Bohr: His Heritage and Legacy. An Anti-Realist View of Quantum Mechanics. Dordrecht: Kluwer Academic Publishers.
* Faye, J. (2002). Copenhagen Interpretation of Quantum Mechanics. The Stanford Encyclopedia of Philosophy (Summer 2002 Edition). Zalta, E.N. (Ed.).
http://plato.stanford.edu/archives/sum2002/entries/qm-copenhagen/.
* Faye, J. & Folse, H. (Eds.) (1994). Niels Bohr and Contemporary Philosophy. Dordrecht: Kluwer Academic Publishers.
* Fell, J.M.G. & Doran, R.S. (1988). Representations of $`^{}`$-Algebras, Locally Compact Groups and Banach $`^{}`$-Algebraic Bundles, Vol. 2. Boston: Academic Press.
* Feyerabend, P. (1981). Niels Bohr’s world view. Realism, Rationalism & Scientific Method: Philosophical Papers Vol. 1, pp. 247–297. Cambridge: Cambridge University Press.
* Fleming, G. & Butterfield, J. (2000). Strange positions. From Physics to Philosophy, pp. 108–165. Butterfield, J. & Pagonis, C. (Eds.). Cambridge: Cambridge University Press.
* Folse, H.J. (1985). The Philosophy of Niels Bohr. Amsterdam: North-Holland.
* Ford, J. (1988). Quantum chaos. Is there any? Directions in Chaos, Vol. 2, pp. 128–147. Bai-Lin, H. (Ed.). Singapore: World Scientific.
* Frasca, M. (2003). General theorems on decoherence in the thermodynamic limit. Physics Letters A308, 135–139.
* Frasca, M. (2004). Fully polarized states and decoherence. arXiv:cond-mat/0403678.
* Frigerio, A. (1974). Quasi-local observables and the problem of measurement in quantum mechanics. Annales de l’ Institut Henri Poincaré A3, 259–270.
* Fröhlich, J., Tsai, T.-P., & Yau, H.-T. (2002). On the point-particle (Newtonian) limit of the non-linear Hartree equation. Communications in Mathematical Physics 225, 223–274.
* Gallavotti, G. (1983). The Elements of Mechanics. Berlin: Springer-Verlag.
* Gallavotti, G., Bonetto, F., & Gentile, G. (2004). Aspects of Ergodic, Qualitative and Statistical Theory of Motion. New York: Springer.
* Gell-Mann, M. & Hartle, J.B. (1990). Quantum mechanics in the light of quantum cosmology. Complexity, Entropy, and the Physics of Information, pp. 425–458. Zurek, W.H. (Ed.). Reading, Addison-Wesley.
* Gell-Mann, M. & Hartle, J.B. (1993). Classical equations for quantum systems. Physical Review D47, 3345–3382.
* Gérard, P. & Leichtnam, E. (1993). Ergodic properties of eigenfunctions for the Dirichlet problem. Duke Mathematical Journal 71, 559–607.
* Gerisch, T., Münzner, R., & Rieckers, A. (1999). Global $`C^{}`$-dynamics and its KMS states of weakly inhomogeneous bipolaronic superconductors. Journal of Statistical Physics 97, 751–779.
* Gerisch, T., Honegger, R., & Rieckers, A. (2003). Algebraic quantum theory of the Josephson microwave radiator. Annales Henri Poincaré 4, 1051–1082.
* Geyer, B., Herwig, H., & Rechenberg, H. (Eds.) (1993). Werner Heisenberg: Physiker und Philosoph. Leipzig: Spektrum.
* Giulini, D. (2003). Superselection rules and symmetries. Decoherence and the Appearance of a Classical World in Quantum Theory, pp. 259–316. Joos, E. et al. (Eds.). Berlin: Springer.
* Glimm, J. & Jaffe, A. (1987). Quantum Physics. A Functional Integral Point of View. New York: Springer-Verlag.
* Gotay, M.J. (1986). Constraints, reduction, and quantization. Journal of Mathematical Physics 27, 2051–2066.
* Gotay, M.J. (1999). On the Groenewold-Van Hove problem for $`𝐑^{2n}`$. Journal of Mathematical Physics 40, 2107–2116.
* Gotay, M.J., Grundling, H.B., & Tuynman, G.M. (1996). Obstruction results in quantization theory. Journal of Nonlinear Science 6, 469–498.
* Gotay, M.J., Nester, J.M., & Hinds, G. (1978). Presymplectic manifolds and the Dirac–Bergmann theory of constraints. Journal of Mathematical Physics 19, 2388–2399.
* Götsch, J. (1992). Erwin Schrödinger’s World View: The Dynamics of Knowlegde and Reality. Dordrecht: Kluwer Academic Publishers.
* Govaerts, J. (1991). Hamiltonian Quantization and Constrained Dynamics. Leuven: Leuven University Press.
* Gracia-Bondía, J.M., Várilly, J.C., & Figueroa, H. (2001). Elements of Noncommutative Geometry. Boston: Birkhäuser.
* Griesemer, M., Lieb, E.H., & Loss, M. (2001). Ground states in non-relativistic quantum electrodynamics. Inventiones Mathematicae 145, 557–595.
* Griffiths, R.B. (1984). Consistent histories and the interpretation of quantum mechanics. Journal of Statistical Physics 36, 219–272.
* Griffiths, R.B. (2002). Consistent Quantum Theory. Cambridge: Cambridge University Press.
* Grigis, A. & Sjöstrand, J. (1994). Microlocal Analysis for Differential Operators. Cambridge: Cambridge University Press.
* Groenewold, H.J. (1946). On the principles of elementary quantum mechanics. Physica 12, 405–460.
* Groenewold, H.J. (1952). Information in quantum measurements. Proceedings Koninklijke Nederlandse Akademie van Wetenschappen B55, 219–227.
* Guhr, T., Müller-Groeling, H., & Weidenmüller, H. (1998). Random matrix theories in quantum physics: common concepts. Physics Reports 299, 189–425.
* Guillemin, V., Ginzburg, V. & Karshon, Y. (2002). Moment Maps, Cobordisms, and Hamiltonian Group Actions. Providence (RI): American Mathematical Society.
* Guillemin, V. & Sternberg, S. (1977). Geometric Asymptotics. Providence (RI): American Mathematical Society.
* Guillemin, V. & Sternberg, S. (1990). Variations on a Theme by Kepler. Providence (RI): American Mathematical Society.
* Guillemin, V. & Uribe, A. (1989). Circular symmetry and the trace formula. Inventiones Mathematicae 96, 385–423.
* Gustafson, S.J. & Sigal, I.M. (2003). Mathematical concepts of quantum mechanics. Berlin: Springer.
* Gutzwiller, M.C. (1971). Periodic orbits and classical quantization conditions. Journal of Mathematical Physics 12, 343–358.
* Gutzwiller, M.C. (1990). Chaos in Classical and Quantum Mechanics. New York: Springer-Verlag.
* Gutzwiller, M.C. (1992). Quantum chaos. Scientific American 266, 78–84.
* Gutzwiller, M.C. (1998). Resource letter ICQM-1: The interplay between classical and quantum mechanics. American Journal of Physics 66, 304–324.
* Haag, R. (1962). The mathematical structure of the Bardeen–Cooper–Schrieffer model. Nuovo Cimento 25, 287–298.
* Haag, R., Kadison, R., & Kastler, D. (1970). Nets of $`C^{}`$-algebras and classification of states. Communications in Mathematical Physics 16, 81–104.
* Haag, R. (1992). Local Quantum Physics: Fields, Particles, Algebras. Heidelberg: Springer-Verlag.
* Haake, F. (2001). Quantum Signatures of Chaos. Second Edition. New York: Springer-Verlag.
* Hagedorn, G.A. (1998). Raising and lowering operators for semiclassical wave packets. Annals of Physics 269, 77–104.
* Hagedorn, G.A. & Joye, A. (1999). Semiclassical dynamics with exponentially small error estimates. Communications in Mathematical Physics 207, 439–465.
* Hagedorn, G.A. & Joye, A. (2000). Exponentially accurate semiclassical dynamics: propagation, localization, Ehrenfest times, scattering, and more general states. Annales Henri Poincaré 1, 837–883.
* Halliwell, J.J. (1998). Decoherent histories and hydrodynamic equations. Physical Review D58, 105015-1–12.
* Halliwell, J.J. (2000). The emergence of hydrodynamic equations from quantum theory: a decoherent histories analysis. International Journal of Theoretical Physics 39, 1767–1777.
* Halliwell, J.J. (2004).Some recent developments in the decoherent histories approach to quantum theory. Lecture Notes in Physics 633, 63–83.
* Halliwell, J.J. (2005). How the quantum Universe becomes classical. arXiv:quant-ph/0501119.
* Halvorson, H. (2004). Complementarity of representations in quantum mechanics. Studies in History and Philosophy of Modern Physics B35, 45–56.
* Halvorson, H. (2005). Algebraic quantum field theory. This volume.
* Halvorson, H. & Clifton, R. (1999). Maximal beable subalgebras of quantum-mechanical observables. International Journal of Theoretical Physics 38, 2441–2484.
* Halvorson, H. & Clifton, R. (2002). Reconsidering Bohr’s reply to epr. Non-locality and Modality, pp. 3–18. Placek, T. & Butterfield, J. (Eds.). Dordrecht: Kluwer Academic Publishers.
* Hannabuss, K.C. (1984). Dilations of a quantum measurement. Helvetica Physica Acta 57, 610–620.
* Harrison, F.E. & Wan, K.K. (1997). Macroscopic quantum systems as measuring devices: dc SQUIDs and superselection rules. Journal of Physics A30, 4731–4755.
* Hartle, J.B. (1995). Spacetime quantum mechanics and the quantum mechanics of spacetime. Gravitation et Quantifications (Les Houches, 1992), pp. 285–480. Amsterdam: North-Holland.
* Hartle, J.B. (2005). What connects different interpretations of quantum mechanics? Quo Vadis Quantum Mechanics, pp. 73-82. Elitzur, A., Dolev, S., & Kolenda, N. (Eds.). Heidelberg: Springer-Verlag. arXiv:quant-ph/0305089.
* Heath, D. & Sudderth, W. (1976). De Finetti’s theorem on exchangeable variables. American Statistics 30, 188–189.
* Heelan, P. (1965). Quantum Mechanics and Objectivity: A Study of the Physical Philosophy of Werner Heisenberg. Den Haag: Martinus Nijhoff.
* Heilbron, J. (2000). The Dilemmas of an Upright Man: Max Planck as a Spokesman for German Science. Second Edition. Los Angeles: University of California Press.
* Heisenberg, W. (1925). Über die quantentheoretische Umdeutung kinematischer und mechanischer Beziehungen. Zeitschrift für Physik 33, 879-893.
* Heisenberg, W. (1927). Über den anschaulichen Inhalt der quantentheoretischen Kinematik und Mechanik. Zeitschrift für Physik 43, 172–198.
* Heisenberg, W. (1930). The Physical Principles of the Quantum Theory. Chicago: University of Chicago Press.
* Heisenberg, W. (1942). Ordnung der Wirklichkeit. In Heisenberg (1984a), pp. 217–306. Also available at http://werner-heisenberg.unh.edu/Ordnung.htm.
* Heisenberg, W. (1958). Physics and Philosophy: The Revolution in Modern Science. London: Allen & Unwin.
* Heisenberg, W. (1969). Der Teil und das Ganze: Gespräche im Umkreis der Atomphysik. München: Piper. English translation as Heisenberg (1971).
* Heisenberg, W. (1971). Physics and Beyond. New York: Harper and Row. Translation of Heisenberg (1969).
* Heisenberg, W. (1984a). Gesammelte Werke. Series C: Philosophical and Popular Writings, Vol i: Physik und Erkenntnis 1927–1955. Blum, W., Dürr, H.-P., & Rechenberg, H. (Eds.). München: Piper.
* Heisenberg, W. (1984b). Gesammelte Werke. Series C: Philosophical and Popular Writings, Vol ii: Physik und Erkenntnis 1956–1968. Blum, W., Dürr, H.-P., & Rechenberg, H. (Eds.). München: Piper.
* Heisenberg, W. (1985). Gesammelte Werke. Series C: Philosophical and Popular Writings, Vol iii: Physik und Erkenntnis 1969–1976. Blum, W., Dürr, H.-P., & Rechenberg, H. (Eds.). München: Piper.
* Held, C. (1994). The Meaning of Complementarity. Studies in History and Philosophy of Science 25, 871–893.
* Helffer, B. (1988) Semi-classical Analysis for the Schrödinger Operator and Applications. Lecture Notes in Mathematics 1336. Berlin: Springer-Verlag.
* Heller, E.J. & Tomsovic, S. (1993). Postmodern quantum mechanics. Physics Today July, 38–46.
* Hendry, J. (1984). The Creation of Quantum Mechanics and the Bohr-Pauli Dialogue. Dordrecht: D. Reidel.
* Henneaux, M. & Teitelboim, C. (1992). Quantization of Gauge Systems. Princeton: Princeton University Press.
* Hepp, K. (1972). Quantum theory of measurement and macroscopic observables. Helvetica Physica Acta 45, 237–248.
* Hepp, K. (1974). The classical limit of quantum mechanical correlation functions. Communications in Mathematical Physics 35, 265–277.
* Hepp, K. & Lieb, E. (1974). Phase transitions in reservoir driven open systems with applications to lasers and superconductors. Helvetica Physica Acta 46, 573–602.
* Higson, N. (1990). A primer on $`KK`$-theory. Operator Theory: Operator Algebras and Applications. Proceedings Symposia in Pure Mathematical, 51, Part 1, pp. 239–283. Providence, RI: American Mathematical Society
* Hillery, M., O’Connel, R.F., Scully, M.O., & Wigner, E.P. (1984). Distribution functions in physics – Fundamentals. Physics Reports 106, 121–167.
* Hislop, P. D. & Sigal, I. M. (1996). Introduction to Spectral Theory. With Applications to Schrödinger Operators. New York: Springer-Verlag.
* Hoffmann, E.T.A. (1810). Musikalische Novellen und Aufsätze. Leipzig: Insel-Bücherei.
* Holevo, A.S. (1982). Probabilistic and Statistical Aspects of Quantum Theory. Amsterdam: North-Holland Publishing Co.
* Hogreve, H., Potthoff, J., & Schrader, R. (1983). Classical limits for quantum particles in external Yang–Mills potentials. Communications in Mathematical Physics 91, 573–598.
* Honegger, R. & Rieckers, A. (1994). Quantized radiation states from the infinite Dicke model. Publications of the Research Institute for Mathematical Sciences (Kyoto) 30, 111–138.
* Honner, J. (1987). The Description of Nature: Niels Bohr and the Philosophy of Quantum Physics. Oxford: Oxford University Press.
* Hooker, C.A. (1972). The nature of quantum mechanical reality: Einstein versus Bohr. Paradigms & Paradoxes: The Philosophical Challenges of the Quantum Domain, pp. 67–302. Colodny, J. (Ed.). Pittsburgh: University of Pittsburgh Press.
* Hörmander, L. (1965). Pseudo-differential operators. Communications in Pure Applied Mathematical 18, 501–517
* Hörmander, L. (1979). The Weyl calculus of pseudo-differential operators. Communications in Pure Applied Mathematical 32, 359–443.
* Hörmander, L. (1985a). The Analysis of Linear Partial Differential Operators, Vol. III. Berlin: Springer-Verlag.
* Hörmander, L. (1985b). The Analysis of Linear Partial Differential Operators, Vol. IV. Berlin: Springer-Verlag.
* Horowitz, G.T. & Marolf, D. (1995). Quantum probes of spacetime singularities. Physical Review D52, 5670–5675.
* Hörz, H. (1968). Werner Heisenberg und die Philosophie. Berin: VEB Deutscher Verlag der Wissenschaften.
* Howard, D. (1990). ‘Nicht sein kann was nicht sein darf’, or the Prehistory of epr, 1909-1935: Einstein’s early worries about the quantum mechanics of composite systems. Sixty-Two Years of Uncertainty, pp. 61–11. Miller, A.I. (Ed.). New York: Plenum.
* Howard, D. (1994). What makes a classical concept classical? Towards a reconstruction of Niels Bohr’s philosophy of physics. Niels Bohr and Contemporary Philosophy, pp. 201–229. Faye, J. & Folse, H. (Eds.). Dordrecht: Kluwer Academic Publishers.
* Howard, D. (2004). Who Invented the Copenhagen Interpretation? Philosophy of Science 71, 669-682.
* Howe, R. (1980). Quantum mechanics and partial differential equations. Journal of Functional Analysis 38, 188–254
* Hudson, R.L. (1974). When is the Wigner quasi-probability density non-negative? Reports of Mathematical Physics 6, 249–252.
* Hudson, R.L. & Moody, G.R. (1975/76). Locally normal symmetric states and an analogue of de Finetti’s theorem. Z. Wahrscheinlichkeitstheorie und Verw. Gebiete 33, 343–351
* Hunziker, W. & Sigal, I. M. (2000). The quantum $`N`$-body problem. Journal of Mathematical Physics 41, 3448–3510.
* Husimi, K. (1940). Some formal properties of the density matrix. Progress of the Physical and Mathematical Society of Japan 22, 264–314.
* Isham, C.J. (1984). Topological and global aspects of quantum theory. Relativity, Groups and Topology, II (Les Houches, 1983), 1059–1290. Amsterdam: North-Holland.
* Isham, C.J. (1994). Quantum logic and the histories approach to quantum theory. Journal of Mathematical Physics 35, 2157–2185.
* Isham, C.J. (1997). Topos theory and consistent histories: the internal logic of the set of all consistent sets. International Journal of Theoretical Physics 36, 785–814.
* Isham, C.J. & Butterfield, J. (2000). Some possible roles for topos theory in quantum theory and quantum gravity. Foundations of Physics 30, 1707–1735.
* Isham, C.J. & Linden, N. (1994). Quantum temporal logic and decoherence functionals in the histories approach to generalized quantum theory. Journal of Mathematical Physics 35, 5452–5476.
* Isham, C.J. & Linden, N. (1995). Continuous histories and the history group in generalized quantum theory. Journal of Mathematical Physics 36, 5392–5408.
* Isham, C.J., Linden, N., & Schreckenberg, S. (1994). The classification of decoherence functionals: an analog of Gleason’s theorem. Journal of Mathematical Physics 35, 6360–6370.
* Ivrii, V. (1998). Microlocal Analysis and Precise Spectral Asymptotics. New York: Springer-Verlag.
* Jalabert, R.O. & Pastawski, H.M. (2001). Environment-Independent Decoherence Rate in Classically Chaotic Systems. Physical Review Letters 86, 2490–2493.
* Jammer, M. (1966). The Conceptual Development of Quantum Mechanics. New York: McGraw-Hill.
* Jammer, M. (1974). The Philosophy of Quantum Mechanics. New York: Wiley.
* Janssens, B. (2004). Quantum Measurement: A Coherent Description. M.Sc. Thesis, Radboud Universiteit Nijmegen.
* Jauch, J.M. (1968). Foundations of Quantum Mechanics. Reading (MA): Addison-Wesley.
* Joos, E. & Zeh, H.D. (1985). The emergence of classical properties through interaction with the environment. Zeitschrift für Physik B59, 223–243.
* Joos, E., Zeh, H.D., Kiefer, C., Giulini, D., Kupsch, J., & Stamatescu, I.-O. (2003). Decoherence and the Appearance of a Classical World in Quantum Theory. Berlin: Springer-Verlag.
* Jørgensen, P.E.T. & Moore, R.T. (1984). Operator Commutation Relations. Dordrecht: Reidel.
* Kaplan, L. & Heller, E.J. (1998a). Linear and nonlinear theory of eigenfunction scars. Annals of Physics 264, 171–206.
* Kaplan, L. & Heller, E.J. (1998b). Weak quantum ergodicity. Physica D 121, 1–18.
* Kaplan, L. (1999). Scars in quantum chaotic wavefunctions. Nonlinearity 12, R1–R40.
* Katok, A. & Hasselblatt, B. (1995). Introduction to the Modern Theory of Dynamical Systems. Cambridge: Cambridge University Press.
* Kadison, R.V. & Ringrose, J.R. (1983). Fundamentals of the theory of operator algebras. Vol. 1: Elementary Theory. New York: Academic Press.
* Kadison, R.V. & Ringrose, J.R. (1986). Fundamentals of the theory of operator algebras. Vol. 2: Advanced Theory. New York: Academic Press.
* Karasev, M.V. (1989). The Maslov quantization conditions in higher cohomology and analogs of notions developed in Lie theory for canonical fibre bundles of symplectic manifolds. I, II. Selecta Mathematica Formerly Sovietica 8, 213–234, 235–258.
* Kent, A. (1990). Against Many-Worlds Interpretations. International Journal of Modern Physics A5 (1990) 1745.
* Kent, A. (1997). Consistent sets yield contrary inferences in quantum theory. Physical Review Letters 78, 2874–2877. Reply by Griffiths, R.B. & Hartle, J.B. (1998). ibid. 81, 1981. Reply to this reply by Kent, A. (1998). ibid. 81, 1982.
* Kent, A. (1998). Quantum histories. Physica Scripta T76, 78–84.
* Kent, A. (2000). Night thoughts of a quantum physicist. Philosophical Transactions of the Royal Society of London A 358, 75–88.
* Kiefer, C. (2003). Consistent histories and decoherence. Decoherence and the Appearance of a Classical World in Quantum Theory, pp. 227–258. Joos, E. et al. (Eds.). Berlin: Springer-Verlag.
* Kirchberg, E. & Wassermann, S. (1995). Operations on continuous bundles of $`C^{}`$-algebras. Mathematische Annalen 303, 677–697.
* Kirillov, A.A. (1990). Geometric Quantization. Dynamical Systems IV, pp. 137–172. Arnold, V.I. & S.P. Novikov (Eds.). Berlin: Springer-Verlag.
* Kirillov, A.A. (2004). Lectures on the Orbit Method. Providence, RI: American Mathematical Society.
* Klauder, J.R. & B.-S. Skagerstam (Eds.). (1985). Coherent States. Singapore: World Scientific.
* Klingenberg, W. (1982). Riemannian Geometry. de Gruyter, Berlin.
* Kohn, J.J. & Nirenberg, L. (1965). An algebra of pseudo-differential operators. Communications in Pure and Applied Mathematics 18 269–305.
* Koopman, B.O. (1931). Hamiltonian systems and transformations in Hilbert space. Proceedings of the National Academy of Sciences 18, 315–318.
* Kostant, B. (1970). Quantization and unitary representations. Lecture Notes in Mathematics 170, 87–208.
* Krishnaprasad, P. S. & Marsden, J. E. (1987). Hamiltonian structures and stability for rigid bodies with flexible attachments. Archive of Rational Mechanics and Analysis 98, 71–93.
* Kuhn, T. S. (1978). Black-body Theory and the Quantum Discontinuity: 1894 1912. New York: Oxford University Press.
* Kümmerer, B. (2002). Quantum Markov processes. Coherent Evolution in Noisy Environment (Lecture Notes in Physics Vol. 611), Buchleitner, A. & Hornberger, K. (Eds.), pp. 139-198. Berlin: Springer-Verlag.
* Lahti, P. & Mittelstaedt, P. (Eds.) (1987). The Copenhagen Interpretation 60 Years After the Como Lecture. Singapore: World Scientific.
* Landau, L.D. & Lifshitz, E.M. (1977). Quantum Mechanics: Non-relativistic Theory. 3d Ed. Oxford: Pergamon Press.
* Landsman, N.P. (1990a). Quantization and superselection sectors I. Transformation group $`C^{}`$-algebras. Reviews in Mathematical Physics 2, 45–72.
* Landsman, N.P. (1990b). Quantization and superselection sectors II. Dirac Monopole and Aharonov–Bohm effect. Reviews in Mathematical Physics 2, 73–104.
* Landsman, N.P. (1991). Algebraic theory of superselection sectors and the measurement problem in quantum mechanics. International Journal of Modern Physics A30, 5349–5371.
* Landsman, N.P. (1992) Induced representations, gauge fields, and quantization on homogeneous spaces. Reviews in Mathematical Physics 4, 503–528.
* Landsman, N.P. (1993). Deformations of algebras of observables and the classical limit of quantum mechanics. Reviews in Mathematical Physics 5, 775–806.
* Landsman, N.P. (1995). Observation and superselection in quantum mechanics. Studies in History and Philosophy of Modern Physics 26B, 45–73.
* Landsman, N.P. (1997). Poisson spaces with a transition probability. Reviews in Mathematical Physics 9, 29–57.
* Landsman, N.P. (1998). Mathematical Topics Between Classical and Quantum Mechanics. New York: Springer-Verlag.
* Landsman, N.P. (1999a). Quantum Mechanics on phase Space. Studies in History and Philosophy of Modern Physics 30B, 287–305.
* Landsman, N.P. (1999b). Lie groupoid $`C^{}`$-algebras and Weyl quantization. Communications in Mathematical Physics 206, 367–381.
* Landsman, N.P. (2001). Quantized reduction as a tensor product. Quantization of Singular Symplectic Quotients, pp. 137–180. Landsman, N.P., Pflaum, M.J., & Schlichenmaier, M. (Eds.). Basel: Birkhäuser.
* Landsman, N.P. (2002). Quantization as a functor. Contemporary Mathematics 315, 9–24.
* Landsman, N.P. (2005a). Functorial quantization and the Guillemin–Sternberg conjecture. Twenty Years of Bialowieza: A Mathematical Anthology, pp. 23–45. Ali, S.T., Emch, G.G., Odzijewicz, A., Schlichenmaier, M., & Woronowicz, S.L. (Eds). Singapore: World Scientific. arXiv:math-ph/0307059.
* Landsman, N.P. (2005b). Lie Groupoids and Lie algebroids in physics and noncommutative geometry. Journal of Geom. Physics , to appear.
* Landsman, N.P. (2006). When champions meet: Rethinking the Bohr–Einstein debate. Studies in History and Philosophy of Modern Physics, to appear. arXiv:quant-ph/0507220.
* Landsman, N.P. & Ramazan, B. (2001). Quantization of Poisson algebras associated to Lie algebroids. Contemporary Mathematics 282, 159–192.
* Laurikainen, K.V. (1988). Beyond the Atom: The Philosophical Thought of Wolfgang Pauli. Berlin: Springer-Verlag.
* Lazutkin, V.F. (1993). KAM Theory and Semiclassical Approximations to Eigenfunctions. Berlin: Springer-Verlag.
* Leggett, A.J. (2002). Testing the limits of quantum mechanics: motivation, state of play, prospects. Journal of Physics: Condensed Matter 14, R415–R451.
* Liboff, R.L. (1984). The correspondence principle revisited. Physics Today February, 50–55.
* Lieb, E.H. (1973). The classical limit of quantum spin systems. Communications in Mathematical Physics 31, 327–340.
* Littlejohn, R.G. (1986). The semiclassical evolution of wave packets. Physics Reports 138, 193–291.
* Ludwig, G. (1985). An Axiomatic Basis for Quantum Mechanics. Volume 1: Derivation of Hilbert Space Structure. Berlin: Springer-Verlag.
* Lugiewicz, P. & Olkiewicz, R. (2002). Decoherence in infinite quantum systems. Journal of Physics A35, 6695–6712.
* Lugiewicz, P. & Olkiewicz, R. (2003). Classical properties of infinite quantum open systems. Communications in Mathematical Physics 239, 241–259.
* Maassen, H. (2003). Quantum probability applied to the damped harmonic oscillator. Quantum Probability Communications, Vol. XII (Grenoble, 1998), pp. 23–58. River Edge, NJ: World Scientific Publishing.
* Mackey, G.W. (1962). The Mathematical Foundations of Quantum Mechanics. New York: Benjamin.
* Mackey, G.W. (1968). Induced Representations of Groups and Quantum Mechanics. New York: W. A. Benjamin; Turin: Editore Boringhieri.
* Mackey, G.W. (1978). Unitary Group Representations in Physics, Probability, and Number Theory. Reading, Mass.: Benjamin/Cummings Publishing Co.
* Mackey, G.W. (1992). The Scope and History of Commutative and Noncommutative Harmonic Analysis. Providence, RI: American Mathematical Society.
* Majid, S. (1988). Hopf algebras for physics at the Planck scale. Classical & Quantum Gravity 5, 1587–1606.
* Majid, S. (1990). Physics for algebraists: noncommutative and noncocommutative Hopf algebras by a bicrossproduct construction. Journal of Algebra 130, 17–64.
* Marmo, G., Scolarici, G., Simoni, A., & Ventriglia, F. (2005). The quantum-classical transition: the fate of the complex structure. International Journal of Geometric Methods in Physics 2, 127–145.
* Marsden, J.E. (1992). Lectures on Mechanics. Cambridge: Cambridge University Press.
* Marsden, J.E. & T.S. Ratiu (1994). Introduction to Mechanics and Symmetry. New York: Springer-Verlag.
* Marsden, J.E., Raţiu, T., & Weinstein, A. (1984). Semidirect products and reduction in mechanics. Transactions of the American Mathematical Society 281, 147–177.
* Marshall, W., Simon, C., Penrose, R., & Bouwmeester, D. (2003). Towards quantum superpositions of a mirror. Physical Review Letters 91, 130401-1–4.
* Martinez, A. (2002). An Introduction to Semiclassical and Microlocal Analysis. New York: Springer-Verlag.
* Maslov, V.P. & Fedoriuk, M.V. (1981). Semi-Classical Approximation in Quantum Mechanics. Dordrecht: Reidel.
* McCormmach, R. (1982). Night Thoughts of a Classical Physicist. Cambridge (MA): Harvard University Press.
* Mehra, J. & and Rechenberg, H. (1982a). The Historical Development of Quantum Theory. Vol. 1: The Quantum Theory of Planck, Einstein, Bohr, and Sommerfeld: Its Foundation and the Rise of Its Difficulties. New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (1982b). The Historical Development of Quantum Theory. Vol. 2: The Discovery of Quantum Mechanics . New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (1982c). The Historical Development of Quantum Theory. Vol. 3: The formulation of matrix mechanics and its modifications, 1925-1926. New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (1982d). The Historical Development of Quantum Theory. Vol. 4: The fundamental equations of quantum mechanics 1925-1926. The reception of the new quantum mechanics. New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (1987). The Historical Development of Quantum Theory. Vol. 5: Erwin Schrödinger and the Rise of Wave Mechanics. New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (2000). The Historical Development of Quantum Theory. Vol. 6: The Completion of Quantum Mechanics 1926–1941. Part 1: The probabilistic Interpretation and the Empirical and Mathematical Foundation of Quantum Mechanics, 1926-1936. New York: Springer-Verlag.
* Mehra, J. & and Rechenberg, H. (2001). The Historical Development of Quantum Theory. Vol. 6: The Completion of Quantum Mechanics 1926–1941. Part 2: The Conceptual Completion of Quantum Mechanics. New York: Springer-Verlag.
* Meinrenken, E. (1998). Symplectic surgery and the $`\mathrm{Spin}^c`$-Dirac operator. Adv. Mathematical 134, 240–277.
* Meinrenken, E. & Sjamaar, R. (1999). Singular reduction and quantization. Topology 38, 699–762.
* Mermin, N.D. (2004). What’s wrong with this quantum world? Physics Today 57 (2), 10.
* Mielnik, B. (1968). Geometry of quantum states. Communications in Mathematical Physics 9, 55–80.
* Miller, A.I. (1984). Imagery in Scientific Thought: Creating 20th-Century Physics. Boston: Birkhäuser.
* Mirlin, A.D. (2000). Statistics of energy levels and eigenfunctions in disordered systems. Physics Reports 326, 259–382.
* Mittelstaedt, P. (2004). The Interpretation of Quantum Mechanics and the Measurement Process. Cambridge: Cambridge University Press.
* Moore, G.E. (1939). Proof of an external world. Proceedings of the British Academy 25, 273–300. Reprinted in Philosophical Papers (George, Allen and Unwin, London, 1959) and in Selected Writings (Routledge, London, 1993).
* Moore, W. (1989). Schrödinger: Life and Thought. Cambridge: Cambridge University Press.
* Morchio, G. & Strocchi, F. (1987). Mathematical structures for long-range dynamics and symmetry breaking. Journal of Mathematical Physics 28, 622–635.
* Muller, F.A. (1997). The equivalence myth of quantum mechanics I, II. Studies in History and Philosophy of Modern Physics 28 35–61, 219–247.
* Murdoch, D. (1987). Niels Bohr s Philosophy of Physics. Cambridge: Cambridge University Press.
* Nadirashvili, N., Toth, J., & Yakobson, D. (2001). Geometric properties of eigenfunctions. Russian Mathematical Surveys 56, 1085–1105.
* Nagy, G. (2000). A deformation quantization procedure for $`C^{}`$-algebras. Journal of Operator Theory 44, 369–411.
* Narnhofer, H. (2001). Quantum K-systems and their abelian models. Foundations of Probability and Physics, pp. 274–302. River Edge, NJ: World Scientific.
* Natsume, T. & Nest, R. (1999). Topological approach to quantum surfaces. Communications in Mathematical Physics 202, 65–87.
* Natsume, T., Nest, R., & Ingo, P. (2003). Strict quantizations of symplectic manifolds. Letters in Mathematical Physics 66, 73–89.
* Nauenberg, M. (1989). Quantum wave packets on Kepler elliptic orbits. Physical Review A40, 1133–1136.
* Nauenberg, M., Stroud, C., & Yeazell, J. (1994). The classical limit of an atom. Scientific American June, 24–29.
* Neumann, J. von (1931). Die Eindeutigkeit der Schrödingerschen Operatoren. Mathematische Annalen 104, 570–578.
* Neumann, J. von (1932). Mathematische Grundlagen der Quantenmechanik. Berlin: Springer–Verlag. English translation (1955): Mathematical Foundations of Quantum Mechanics. Princeton: Princeton University Press.
* Neumann, J. von (1938). On infinite direct products. Compositio Mathematica 6, 1–77.
* Neumann, J. von (1981). Continuous geometries with a transition probability. Memoirs of the American Mathematical Society 252, 1–210. (Edited by I.S. Halperin. MS from 1937).
* Neumann, H. (1972). Transformation properties of observables. Helvetica Physica Acta 25, 811-819.
* Nourrigat, J. & Royer, C. (2004). Thermodynamic limits for Hamiltonians defined as pseudodifferential operators. Communications in Partial Differential Equations 29, 383–417.
* O’Connor, P.W., Tomsovic, S., & Heller, E.J. (1992). Semiclassical dynamics in the strongly chaotic regime: breaking the log time barrier. Physica D55, 340–357.
* Odzijewicz, A. (1992). Coherent states and geometric quantization. Communications in Mathematical Physics 150, 385–413.
* Odzijewicz, A. & Ratiu, T.S. (2003). Banach Lie-Poisson spaces and reduction. Communications in Mathematical Physics 243, 1–54.
* Olkiewicz, R. (1999a). Dynamical semigroups for interacting quantum and classical systems. Journal of Mathematical Physics 40, 1300–1316.
* Olkiewicz, R. (1999b). Environment-induced superselection rules in Markovian regime. Communications in in Mathematical Physics 208, 245–265.
* Olkiewicz, R. (2000). Structure of the algebra of effective observables in quantum mechanics. Annals of Physics 286, 10–22.
* Ollivier, H., Poulin, D., & Zurek, W.H. (2004). Environment as witness: selective proliferation of information and emergence of objectivity. arXiv: quant-ph/0408125.
* Olshanetsky, M.A. & Perelomov, A.M. (1981). Classical integrable finite-dimensional systems related to Lie algebras. Physics Reports 71, 313–400.
* Olshanetsky, M.A. & Perelomov, A.M. (1983). Quantum integrable systems related to Lie algebras. Physics Reports 94, 313–404.
* Omnès, R. (1992). Consistent interpretations of quantum mechanics. Reviews of Modern Physics 64, 339–382.
* Omnès, R. (1994). The Interpretation of Quantum Mechanics. Princeton: Princeton University Press.
* Omnès, R. (1997). Quantum-classical correspondence using projection operators. Journal of Mathematical Physics 38, 697–707.
* Omnès, R. (1999). Understanding Quantum Mechanics. Princeton: Princeton University Press.
* Ørsted, B. (1979). Induced representations and a new proof of the imprimitivity theorem. Journal of Functional Analysis 31, 355–359.
* Ozorio de Almeida, A.M. (1988). Hamiltonian Systems: Chaos and Quantization. Cambridge: Cambridge University Press.
* Pais, A. (1982). Subtle is the Lord: The Science and Life of Albert Einstein. Oxford: Oxford University Press.
* Pais, A. (1991). Niels Bohr s Times: In Physics, Philosophy, and Polity. Oxford: Oxford University Press.
* Pais, A. (1997). A Tale of Two Continents: A Physicist’s Life in a Turbulent World. Princeton: Princeton University Press.
* Pais, A. (2000). The Genius of Science. Oxford: Oxford University Press.
* Parthasarathy, K.R. (1992). An Introduction to Quantum Stochastic Calculus. Basel: Birkhäuser.
* Paul, T. & Uribe, A. (1995). The semi-classical trace formula and propagation of wave packets. Journal of Functional Analysis 132, 192–249.
* Paul, T. & Uribe, A. (1996). On the pointwise behavior of semi-classical measures. Communications in Mathematical Physics 175, 229–258.
* Paul, T. & Uribe, A. (1998). A. Weighted trace formula near a hyperbolic trajectory and complex orbits. Journal of Mathematical Physics 39, 4009–4015.
* Pauli, W. (1925). Über den Einfluß der Geschwindigkeitsabhängigkeit der Elektronenmasse auf den Zeemaneffekt. Zeitschrift für Physik 31, 373–385.
* Pauli, W. (1933). Die allgemeinen Prinzipien der Wellenmechanik. Flügge, S. (Ed.). Handbuch der Physik, Vol. V, Part I. Translated as Pauli, W. (1980). General Principles of Quantum Mechanics. Berlin: Springer-Verlag.
* Pauli, W. (1949). Die philosophische Bedeutung der Idee der Komplementarität. Reprinted in von Meyenn, K. (Ed.) (1984). Wolfang Pauli: Physik und Erkenntnistheorie, pp. 10–23. Braunschweig: Vieweg Verlag. English translation in Pauli (1994).
* Pauli, W. (1979). Wissenschaftlicher Briefwechsel mit Bohr, Einstein, Heisenberg. Vol 1: 1919–1929. Hermann, A., von Meyenn, K., & Weisskopf, V. (Eds.). New York: Springer-Verlag.
* Pauli, W. (1985). Wissenschaftlicher Briefwechsel mit Bohr, Einstein, Heisenberg. Vol 2: 1930–1939. von Meyenn, K. (Ed.). New York: Springer-Verlag.
* Pauli, W. (1994). Writings on Physics and Philosophy. Enz, C.P. & von Meyenn, K. (Eds.). Berlin: Springer-Verlag.
* Paz, J.P. & Zurek, W.H. (1999). Quantum limit of decoherence: environment induced superselection of energy eigenstates. Physical Review Letters 82, 5181–5185.
* Pedersen, G.K. (1979.) $`C^{}`$-algebras and their Automorphism Groups. London: Academic Press.
* Pedersen, G.K. (1989). Analysis Now. New York: Springer-Verlag.
* Perelomov, A. (1986.) Generalized Coherent States and their Applications. Berlin: Springer-Verlag.
* Peres, A. (1984). Stability of quantum motion in chaotic and regular systems. Physical Review A30, 1610–1615.
* Peres, A. (1995). Quantum Theory: Concepts and Methods. Dordrecht: Kluwer Academic Publishers.
* Petersen, A. (1963). The Philosophy of Niels Bohr. Bulletin of the Atomic Scientists 19, 8–14.
* Pitowsky, I. (1989). Quantum Probability - Quantum Logic. Berlin: Springer-Verlag.
* Planck, M. (1906). Vorlesungen Über die Theorie der Wärmestrahlung Leipzig: J.A. Barth.
* Poincaré, H. (1892–1899). Les Méthodes Nouvelles de la Méchanique Céleste. Paris: Gauthier-Villars.
* Popov, G. (2000). Invariant tori, effective stability, and quasimodes with exponentially small error terms. I & II. Annales Henri Poincaré 1, 223–248 & 249–279.
* Poulin, D. (2004). Macroscopic observables. arXiv:quant-ph/0403212.
* Poulsen, N.S. (1970). Regularity Aspects of the Theory of Infinite-Dimensional Representations of Lie Groups. Ph.D Thesis, MIT.
* Primas, H. (1983). Chemistry, Quantum Mechanics and Reductionism. Second Edition. Berlin: Springer-Verlag.
* Primas, H. (1997). The representation of facts in physical theories. Time, Temporality, Now, pp. 241-263. Atmanspacher, H. & Ruhnau, E. (Eds.). Berlin: Springer-Verlag.
* Prugovecki, E. (1971). Quantum Mechanics in Hilbert Space. New York: Academic Press.
* Puta, M. (1993). Hamiltonian Dynamical Systems and Geometric Quantization. Dordrecht: D. Reidel.
* Raggio, G.A. (1981). States and Composite Systems in $`W^{}`$-algebras Quantum Mechanics. Ph.D Thesis, ETH Zürich.
* Raggio, G.A. (1988). A remark on Bell’s inequality and decomposable normal states. Letters in Mathematical Physics 15, 27–29.
* Raggio, G.A. & Werner, R.F. (1989). Quantum statistical mechanics of general mean field systems. Helvetica Physica Acta 62, 980–1003.
* Raggio, G.A. & Werner, R.F. (1991). The Gibbs variational principle for inhomogeneous mean field systems. Helvetica Physica Acta 64, 633–667.
* Raimond, R.M., Brune, M., & Haroche, S. (2001). Manipulating quantum entanglement with atoms and photons in a cavity. Reviews of Modern Physics 73, 565–582.
* Rédei, M. (1998). Quantum logic in algebraic approach. Dordrecht: Kluwer Academic Publishers.
* Rédei, M. & Stöltzner, M. (Eds.). (2001). John von Neumann and the Foundations of Modern Physics. Dordrecht: Kluwer Academic Publishers.
* Reed, M. & Simon, B. (1972). Methods of Modern Mathematical Physics. Vol I. Functional Analysis. New York: Academic Press.
* Reed, M. & Simon, B. (1975). Methods of Modern Mathematical Physics. Vol II. Fourier Analysis, Self-adjointness. New York: Academic Press.
* Reed, M. & Simon, B. (1979). Methods of Modern Mathematical Physics. Vol III. Scattering Theory. New York: Academic Press.
* Reed, M. & Simon, B. (1978). Methods of Modern Mathematical Physics. Vol IV. Analysis of Operators. New York: Academic Press.
* Reichl, L.E. (2004). The Transition to Chaos in Conservative Classical Systems: Quantum Manifestations. Second Edition. New York: Springer-Verlag.
* Rieckers, A. (1984). On the classical part of the mean field dynamics for quantum lattice systems in grand canonical representations. Journal of Mathematical Physics 25, 2593–2601.
* Rieffel, M.A. (1989a). Deformation quantization of Heisenberg manifolds. Communications in Mathematical Physics 122, 531–562.
* Rieffel, M.A. (1989b). Continuous fields of $`C^{}`$-algebras coming from group cocycles and actions. Mathematical Annals 283, 631–643.
* Rieffel, M.A. (1994). Quantization and $`C^{}`$-algebras. Contemporary Mathematics 167, 66–97.
* Riesz, F. & Sz.-Nagy, B. (1990). Functional Analysis. New York: Dover.
* Robert, D. (1987) Autour de l’Approximation Semi-Classique. Basel: Birkhäuser.
* Robert, D. (Ed.). (1992) Méthodes Semi-Classiques. Astérisque 207, 1–212, ibid. 210, 1–384.
* Robert, D. (1998). Semi-classical approximation in quantum mechanics. A survey of old and recent mathematical results. Helvetica Physica Acta 71, 44–116.
* Roberts, J.E. (1990). Lectures on algebraic quantum field theory. The Algebraic Theory of Superselection Sectors. Introduction and Recent Results, pp. 1–112. Kastler, D. (Ed.). River Edge, NJ: World Scientific Publishing Co.
* Roberts, J.E. & Roepstorff, G. (1969). Some basic concepts of algebraic quantum theory. Communications in Mathematical Physics 11, 321–338.
* Robinett, R.W. (2004). Quantum wave packet revival. Physics Reports 392, 1-119.
* Robinson, D. (1994). Can Superselection Rules Solve the Measurement Problem? British Journal for the Philosophy of Science 45, 79-93.
* Robinson, S.L. (1988a). The semiclassical limit of quantum mechanics. I. Time evolution. Journal of Mathematical Physics 29, 412–419.
* Robinson, S.L. (1988b). The semiclassical limit of quantum mechanics. II. Scattering theory. Annales de l’ Institut Henri Poincaré A48, 281–296.
* Robinson, S.L. (1993). Semiclassical mechanics for time-dependent Wigner functions. Journal of Mathematical Physics 34, 2185–2205.
* Robson, M.A. (1996). Geometric quantization of reduced cotangent bundles. Journal of Geometry and Physics 19, 207–245.
* Rosenfeld, L. (1967). Niels Bohr in the Thirties. Consolidation and extension of the conception of complementarity. Niels Bohr: His Life and Work as Seen by His Friends and Colleagues, pp. 114–136. Rozental, S. (Ed.). Amsterdam: North-Holland.
* Rudolph, O. (1996a). Consistent histories and operational quantum theory. International Journal of Theoretical Physics 35, 1581–1636.
* Rudolph, O. (1996b). On the consistent effect histories approach to quantum mechanics. Journal of Mathematical Physics 37, 5368–5379.
* Rudolph, O. (2000). The representation theory of decoherence functionals in history quantum theories. International Journal of Theoretical Physics 39, 871–884.
* Rudolph, O. & Wright, J.D. M. (1999). Homogeneous decoherence functionals in standard and history quantum mechanics. Communications in Mathematical Physics 204, 249–267.
* Sarnak, P. (1999). Quantum chaos, symmetry and zeta functions. I. & II. Quantum Chaos. Current Developments in Mathematics, 1997 (Cambridge, MA), pp. 127–144 & 145–159. Boston: International Press.
* Saunders, S. (1993). Decoherence, relative states, and evolutionary adaptation. Foundations of Physics 23, 1553–1585.
* Saunders, S. (1995) Time, quantum mechanics, and decoherence. Synthese 102, 235–266.
* Saunders, S. (2004). Complementarity and Scientific Rationality. arXiv:quant-ph/0412195.
* Scheibe, E. (1973). The Logical Analysis of Quantum Mechanics. Oxford: Pergamon Press.
* Scheibe, E. (1991). J. v. Neumanns und J. S. Bells Theorem. Ein Vergleich. Philosophia Naturalis 28, 35–53. English translation in Scheibe, E. (2001). Between Rationalism and Empiricism: Selected Papers in the Philosophy of Physics. New York: Springer-Verlag.
* Scheibe, E. (1999). Die Reduktion Physikalischer Theorien. Ein Beitrag zur Einheit der Physik. Teil II: Inkommensurabilität und Grenzfallreduktion. Berlin: Springer-Verlag.
* Schlosshauer, M. (2004). Decoherence, the measurement problem, and interpretations of quantum mechanics. Reviews of Modern Physics 76, 1267–1306.
* Schmüdgen, K. (1990). Unbounded Operator Algebras and Representation Theory. Basel: Birkhäuser Verlag.
* Schrödinger, E. (1926a). Quantisierung als Eigenwertproblem. I.-IV. Annalen der Physik 79, 361–76, 489–527, ibid. 80, 437–90, ibid. 81, 109–39. English translation in Schrödinger (1928).
* Schrödinger, E. (1926b). Der stetige Übergang von der Mikro-zur Makromekanik. Die Naturwissenschaften 14, 664–668. English translation in Schrödinger (1928).
* Schrödinger, E. (1926c). Über das Verhaltnis der Heisenberg–Born–Jordanschen Quantenmechanik zu der meinen. Annalen der Physik 79, 734–56. English translation in Schrödinger (1928).
* Schrödinger, E. (1928). Collected Papers on Wave Mechanics. London: Blackie and Son.
* Schroeck, F.E., Jr. (1996). Quantum Mechanics on Phase Space. Dordrecht: Kluwer Academic Publishers.
* Scutaru, H. (1977). Coherent states and induced representations. Letters in Mathematical Physics 2, 101-107.
* Segal, I.E. (1960). Quantization of nonlinear systems. Journal of Mathematical Physics 1, 468–488.
* Sewell, G. L. (1986). Quantum Theory of Collective Phenomena. New York: Oxford University Press.
* Sewell, G. L. (2002). Quantum Mechanics and its Emergent Macrophysics. Princeton: Princeton University Press.
* Simon, B. (1976). Quantum dynamics: from automorphism to Hamiltonian. Studies in Mathematical Phyiscs: Essays in Honour of Valentine Bargmann, pp. 327–350. Lieb, E.H., Simon, B. & Wightman, A.S. (Eds.). Princeton: Princeton University Press.
* Simon, B. (1980). The classical limit of quantum partition functions. Communications in Mathematical Physics 71, 247–276.
* Simon, B. (2000) Schrödinger operators in the twentieth century. Journal of Mathematical Physics 41, 3523–3555.
* Śniatycki, J. (1980). Geometric Quantization and Quantum Mechanics. Berlin: Springer-Verlag.
* Snirelman, A.I. (1974). Ergodic properties of eigenfunctions. Uspekhi Mathematical Nauk 29, 181–182.
* Souriau, J.-M. (1969). Structure des systèmes dynamiques. Paris: Dunod. Translated as Souriau, J.-M. (1997).
* Souriau, J.-M. (1997). Structure of Dynamical Systems. A Symplectic View of Physics. Boston: Birkhäuser.
* Stapp, H.P. (1972). The Copenhagen Interpretation. American Journal of Physics 40, 1098–1116.
* Steiner, M. (1998). The Applicability of Mathematics as a Philosophical Problem. Cambridge (MA): Harvard University Press.
* Stinespring, W. (1955). Positive functions on $`C^{}`$-algebras. Proceedings of the American Mathematical Society 6, 211–216.
* Størmer, E. (1969). Symmetric states of infinite tensor products of $`C^{}`$-algebras. Jornal of Functional Analysis 3, 48–68.
* Strawson, P.F. (1959). Individuals: An Essay in Descriptive Metaphysics. London: Methuen.
* Streater, R.F. (2000). Classical and quantum probability. Journal of Mathematical Physics 41, 3556–3603.
* Strichartz, R.S. (1983). Analysis of the Laplacian on a complete Riemannian manifold. Journal of Functional Analysis 52, 48–79.
* Strocchi, F. (1985). Elements of Quantum mechanics of Infinite Systems. Singapore: World Scientific.
* Strunz, W.T., Haake, F., & Braun, D. (2003). Universality of decoherence for macroscopic quantum superpositions. Physical Review A 67, 022101–022114.
* Summers, S.J. & Werner, R. (1987). Bell s inequalities and quantum field theory, I, II. Journal of Mathematical Physics 28, 2440–2447, 2448–2456.
* Sundermeyer, K. (1982). Constrained Dynamics. Berlin: Springer-Verlag.
* Takesaki, M. (2003). Theory of Operator Algebras. Vols. I-III. New York: Springer-Verlag.
* Thirring, W. & Wehrl, A. (1967). On the mathematical structure of the BCS model. I. Communications in Mathematical Physics 4, 303–314.
* Thirring, W. (1968). On the mathematical structure of the BCS model. II. Communications in Mathematical Physics 7, 181–199.
* Thirring, W. (1981). A Course in Mathematical Physics. Vol. 3: Quantum Mechanics of Atoms and Molecules. New York: Springer-Verlag.
* Thirring, W. (1983). A Course in Mathematical Physics. Vol. 4: Quantum Mechanics of Large Systems. New York: Springer-Verlag.
* Tomsovic, S. & Heller, E.J. (1993). Long-time semiclassical dynamics of chaos: The stadium billiard. Physical Review E47, 282–299.
* Tomsovic, S. & Heller, E.J. (2002). Comment on” Ehrenfest times for classically chaotic systems”. Physical Review E65, 035208-1–2.
* Toth, J.A. (1996). Eigenfunction localization in the quantized rigid body. Journal of Differential Geometry 43, 844–858.
* Toth, J.A. (1999). On the small-scale mass concentration of modes. Communications in Mathematical Physics 206, 409–428.
* Toth, J.A. & Zelditch, S. (2002). Riemannian manifolds with uniformly bounded eigenfunctions. Duke Mathematical Journal 111, 97–132.
* Toth, J.A. & Zelditch, S. (2003a). $`L^p`$ norms of eigenfunctions in the completely integrable case. Annales Henri Poincaré 4, 343–368.
* Toth, J.A. & Zelditch, S. (2003b). Norms of modes and quasi-modes revisited. Contemporary Mathematics 320, 435–458.
* Tuynman, G.M. (1987). Quantization: towards a comparison between methods. Journal of Mathematical Physics 28, 2829–2840.
* Tuynman, G.M. (1998). Prequantization is irreducible. Indagationes Mathematicae (New Series) 9, 607–618.
* Unnerstall, T. (1990a). Phase-spaces and dynamical descriptions of infinite mean-field quantum systems. Journal of Mathematical Physics 31, 680–688.
* Unnerstall, T. (1990b). Schrödinger dynamics and physical folia of infinite mean-field quantum systems. Communications in Mathematical Physics 130, 237–255.
* Vaisman, I. (1991). On the geometric quantization of Poisson manifolds. Journal of Mathematical Physics 32, 3339–3345.
* van Fraassen, B.C. (1991). Quantum Mechanics: An Empiricist View. Oxford: Oxford University Press.
* van Hove, L. (1951). Sur certaines représentations unitaires d’un groupe infini de transformations. Memoires de l’Académie Royale de Belgique, Classe des Sciences 26, 61–102.
* Van Vleck, J.H. (1928). The Correspondence Principle in the Statistical Interpretation of Quantum Mechanics. Proceedings of the National Academy of Sciences 14, 178–188.
* van der Waerden, B.L. (Ed.). (1967). Sources of Quantum Mechanics. Amsterdam: North-Holland.
* van Kampen, N. (1954). Quantum statistics of irreversible processes. Physica 20, 603–622.
* van Kampen, N. (1988). Ten theorems about quantum mechanical measurements. Physica A153, 97–113.
* van Kampen, N. (1993). Macroscopic systems in quantum mechanics. Physica A194, 542–550.
* Vanicek, J. & Heller, E.J. (2003). Semiclassical evaluation of quantum fidelity Physical Review E68, 056208-1–5.
* Vergne, M. (1994). Geometric quantization and equivariant cohomology. First European Congress in Mathematics, Vol. 1, pp. 249–295. Boston: Birkhäuser.
* Vermaas, P. (2000). A Philosopher’s Understanding of Quantum Mechanics: Possibilities and Impossibilities of a Modal Interpretation. Cambridge: Cambridge University Press.
* Vey, J. (1975). Déformation du crochet de Poisson sur une variété symplectique. Commentarii Mathematici Helvetici 50, 421–454.
* Voros, A. (1979). Semi-classical ergodicity of quantum eigenstates in the Wigner representation. Stochastic Behaviour in Classical and Quantum Hamiltonian Systems. Lecture Notes in Physics 93, 326–333.
* Wallace, D. (2002). Worlds in the Everett interpretation. Studies in History and Philosophy of Modern Physics 33B, 637–661.
* Wallace, D. (2003). Everett and structure. Studies in History and Philosophy of Modern Physics 34B, 87–105.
* Wan, K.K. & Fountain, R.H. (1998). Quantization by parts, maximal symmetric operators, and quantum circuits. International Journal of Theoretical Physics 37, 2153–2186.
* Wan, K.K., Bradshaw, J., Trueman, C., & Harrison, F.E. (1998). Classical systems, standard quantum systems, and mixed quantum systems in Hilbert space. Foundations of Physics 28, 1739–1783.
* Wang, X.-P. (1986). Approximation semi-classique de l’equation de Heisenberg. Communications in Mathematical Physics 104, 77–86.
* Wegge-Olsen, N.E. (1993). K-theory and $`C^{}`$-algebras. Oxford: Oxford University Press.
* Weinstein, A. (1983). The local structure of Poisson manifolds. Journal of Differential Geometry 18, 523–557.
* Werner, R.F. (1995). The classical limit of quantum theory. arXiv:quant-ph/9504016.
* Weyl, H. (1931). The Theory of Groups and Quantum Mechanics. New York: Dover.
* Wheeler, J.A. & and Zurek, W.H. (Eds.) (1983). Quantum Theory and Measurement. Princeton: Princeton University Press.
* Whitten-Wolfe, B. & Emch, G.G. (1976). A mechanical quantum measuring process. Helvetica Physica Acta 49, 45-55.
* Wick, C.G., Wightman, A.S., & Wigner, E.P. (1952). The intrinsic parity of elementary particles. Physical Review 88, 101–105.
* Wightman, A.S. (1962) On the localizability of quantum mechanical systems. Reviews of Modern Physics 34, 845-872.
* Wigner, E.P. (1932). On the quantum correction for thermodynamic equilibrium. Physical Review 40, 749–759.
* Wigner, E.P. (1939) Unitary representations of the inhomogeneous Lorentz group. Annals of Mathematics 40, 149-204.
* Wigner, E.P. (1963). The problem of measurement. American Journal of Physics 31, 6–15.
* Williams, F.L. (2003). Topics in Quantum Mechanics. Boston: Birkhäuser.
* Woodhouse, N. M. J. (1992). Geometric Quantization. Second edition. Oxford: The Clarendon Press.
* Yajima, K. (1979). The quasi-classical limit of quantum scattering theory. Communications in Mathematical Physics 69, 101–129.
* Zaslavsky, G.M. (1981). Stochasticity in quantum systems. Physics Reports 80, 157–250.
* Zeh, H.D. (1970). On the interpretation of measurement in quantum theory. Foundations of Physics 1, 69–76.
* Zelditch, S. (1987). Uniform distribution of eigenfunctions on compact hyperbolic surfaces. Duke Mathematical J. 55, 919–941.
* Zelditch, S. (1990). Quantum transition amplitudes for ergodic and for completely integrable systems. Journal of Functional Analysis 94, 415–436.
* Zelditch, S. (1991). Mean Lindelöf hypothesis and equidistribution of cusp forms and Eisenstein series. Journal of Functional Analysis 97, 1–49.
* Zelditch, S. (1992a). On a “quantum chaos” theorem of R. Schrader and M. Taylor. Journal of Functional Analysis 109, 1–21.
* Zelditch, S. (1992b). Quantum ergodicity on the sphere. Communications in Mathematical Physics 146, 61–71.
* Zelditch, S. (1996a). Quantum dynamics from the semiclassical viewpoint. Lectures at the Centre E. Borel. Available at http://mathnt.mat.jhu.edu/zelditch/Preprints/preprints.html.
* Zelditch, S. (1996b). Quantum mixing. Journal of Functional Analysis 140, 68–86.
* Zelditch, S. (1996c). Quantum ergodicity of $`C^{}`$ dynamical systems. Communications in Mathematical Physics 177, 507–528.
* Zelditch, S. & Zworski, M. (1996). Ergodicity of eigenfunctions for ergodic billiards. Communications in Mathematical Physics 175, 673–682.
* Zurek, W.H. (1981). Pointer basis of quantum apparatus: into what mixture does the wave packet collapse? Physical Review D24, 1516–1525.
* Zurek, W.H. (1982) Environment-induced superselections rules. Physical Review D26, 1862–1880.
* Zurek, W.H. (1991). Decoherence and the transition from quantum to classical. Physics Today 44 (10), 36–44.
* Zurek, W.H. (1993). Negotiating the tricky border between quantum and classical. Physics Today 46 (4), 13–15, 81–90.
* Zurek, W.H. (2003). Decoherence, einselection, and the quantum origins of the classical. Reviews of Modern Physics 75, 715–775.
* Zurek, W.H. (2004). Probabilities from entanglement, Born’s rule from envariance.
arXiv:quant-ph/0405161.
* Zurek, W.H., Habib, S., & Paz, J.P. (1993). Coherent states via decoherence. Physical Review Letters 70, 1187–1190.
* Zurek, W.H. & Paz, J.P. (1995). Why We Don’t Need Quantum Planetary Dynamics: Decoherence and the Correspondence Principle for Chaotic Systems. Proceedings of the Fourth Drexel Meeting. Feng, D.H. et al. (Eds.). New York: Plenum. arXiv:quant-ph/9612037.
|
warning/0506/cond-mat0506676.html
|
ar5iv
|
text
|
# Spin transverse force on spin current in an electric field
(Revised on July 27, 2005)
## Abstract
As a relativistic quantum mechanical effect, it is shown that the electric field exerting a transverse force on an electron spin 1/2 only if the electron is moving. The spin force, analogue to the Lorentz for on an electron charge in a magnetic field, is perpendicular to the electric field and the spin current which spin polarization is projected anong the electric field. This spin-dependent force can be used to understand the zitterbewegung of electron wave packet with spin-orbit coupling and is relevant to the generation of the charge Hall effect driven by the spin current in semiconductors.
Recent years spintronics has became an emerging field because of its potential application to semiconductor industry, and more and more attentions are focused on the effect of spin-orbit coupling in metals and semiconductors.Prinz98Science The spin-orbit coupling is a relativistic effect describing the interaction of the electron spin, momentum and electric field, and provides a route to manipulate and to control the quantum spin states via external fields.Datta90 ; Dyakonov71 It is desirable to understand the motion of electron spin with spin-orbit coupling in an electromagnetic field. In electrodynamics it is known that a magnetic field would exert a Lorentz force on an electric charge if it is moving. This Lorentz force can generate a lot of fundamental phenomena such as the Hall effect in solid.Hall The interaction of the spin in the electromagnetic field behaves as if the spin is a gauge charge and the interaction is due to the SU(2) gauge field.Anandan89 It is essential that the electron spin is an intrinsic quantum variable, not a just classical tiny magnetic moment. The physical meaning of the interaction is closely associated to the Aharonov-Casher effect and the Berry phase.
In this paper it is found that an electric field exerts a transverse force on an electron spin if it is moving and the spin is projected along the electric field. The force is proportional to the square of electric field, and the spin current projected along the field, and its direction is always perpendicular to the electric field and the speed of the spin. The force stems from the spin-orbital coupling which can be derived from the Dirac equation for an electron in a potential in the non-relativistic limit or the Kane model with the $`𝐤𝐩`$ coupling between the conduction band and valence band. From an exact solution of a single electron with the spin-orbit coupling it is heuristic to understand that the zitterbewegung of electronic wave packet is driven by the spin transverse force on a moving spin. The role of spin transverse force is also discussed in the spin Hall effect and its reciprocal effect driven by the pure spin current in semiconductors.
We start with the Dirac equation of an electron in a confining potential $`V`$ and a vector potential $`𝐀`$ for a magnetic field, $`𝐁=\times 𝐀`$,
$$i\mathrm{}\frac{}{t}\mathrm{\Psi }=\left[c\alpha \left(𝐩+\frac{e}{c}𝐀\right)+mc^2\gamma +V\right]\mathrm{\Psi },$$
(1)
where $`\alpha `$ and $`\gamma `$ are the $`4\times 4`$ Dirac matrices. $`m`$ and $`e`$ are the electron mass and charge, respectively, and $`c`$ is the speed of light. We let $`\mathrm{\Psi }=\left(\begin{array}{c}\phi \\ \chi \end{array}\right)e^{imc^2t/\mathrm{}}`$ such that the rest mass energy of electron is removed from the energy eigenvalue of the electron. In the nonrelativistic limit, $`\chi `$ is a very small component, $`\chi \frac{1}{mc^2V}c\sigma \left(𝐩+\frac{e}{c}𝐀\right)\phi `$ where $`\sigma `$ are the Pauli matrices. Thus the component $`\phi `$ of the wavefunction satisfies the following equation, $`i\mathrm{}\frac{}{t}\phi H\phi `$, where
$$H\frac{\left(𝐩+\frac{e}{c}𝐀\right)^2}{2m}+V_{eff}+\mu _B\sigma B+\frac{\mathrm{}\left(𝐩+\frac{e}{c}𝐀\right)}{4m^2c^2}\left(\sigma \times V\right)\text{,}$$
(2)
where $`\mu _B=e\mathrm{}/2mc`$ and $`V_{eff}=V+\frac{\mathrm{}^2}{8m^2c^2}^2V.`$ In the last step we neglect the higher order terms of expansion. Thus the Dirac equation is reduced to the Schrödinger equation with the spin-orbit coupling. The same form of effective spin-orbit coupling and Zeeman splitting can be also derived from the $`8\times 8`$ Kane model that takes into account only the $`𝐤𝐩`$ coupling between the $`\mathrm{\Gamma }_6^c`$ conduction band and the $`\mathrm{\Gamma }_8^v`$ and $`\mathrm{\Gamma }_7^v`$ valence bands, although the effective mass, the effective Lande g-factor, and the effective coupling coefficients have to be introduced as material-specific parameters such as $`mm^{}`$, $`\mu _Bg\mu _B/2`$, and $`\mathrm{}^2/(4m^2c^2)r_{41}^{6c6c}`$. Winkler03
In the Heisenberg picture the kinetic velocity is
$$𝐯=\frac{1}{i\mathrm{}}[𝐫,H]=\frac{1}{m}\left[𝐩+\frac{e}{c}\left(𝐀+𝒜\right)\right]$$
(3)
where $`𝒜=\frac{\mathrm{}}{4mce}\sigma \times V`$ and comes from the spin-orbit coupling. It indicates clearly that $`𝒜`$ plays a role of a SU(2) gauge vector potential. The spin dependence of the gauge field can separates the charged carriers with different spins in cyclotron motion experimentally.Rokhinson04prl Even though we have $`[p_\alpha ,p_\beta ]=0`$ for canonical momentum, the analogous commutators do not vanish for the kinetic velocity
$$[v_\alpha ,v_\beta ]=i\frac{\mathrm{}e}{m^2c}ϵ_{\alpha \beta \gamma }_\gamma +\frac{e^2}{m^2c^2}[𝒜_\alpha ,𝒜_\beta ]$$
(4)
where the total magnetic field $`=𝐁+\times 𝒜,`$ and $`\times 𝒜=\mathrm{}\left[\sigma (^2V)(\sigma )V\right]/\left(4mce\right).`$ Notice that $`[𝒜_\alpha ,𝒜_\beta ]=2\mathrm{}^2i(\sigma V)ϵ_{\alpha \beta \gamma }_\gamma V/(4mce)^2.`$ We can derive the quantum mechanical version of the force,
$$m\frac{d𝐯}{dt}=F_h+F_g+F_f$$
(5)
with
$`F_h`$ $`=`$ $`{\displaystyle \frac{e}{c}}𝐯\times \left(V_{eff}+\mu _B\sigma B\right)`$ (6a)
$`F_g`$ $`=`$ $`{\displaystyle \frac{\mu _B}{2mc^2}}[\sigma (𝐁V)𝐁(\sigma V)],`$ (6b)
$`F_f`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{8m^2c^4}}(\sigma V)(𝐯\times V)`$ (6c)
This is the quantum mechanical analogue of Newton’s second law. Of course we should notice that this is just an operator equation. The uncertainty relationship tells that the position and momentum cannot be measured simultaneously, and there is no concept of force in quantum mechanics. To see the physical meaning of the equation, we take the expectation values of both sides with respect to a Heisenberg state $`|\mathrm{\Phi }`$ which does not varies with time. The expectation values of the observable describe the motion of the center of the wave package of electrons. In this sense we have an equation of the force experienced by the moving electron. Actually the first term $`F_h`$ in Eq. (5) is the Lorentz force for a charged particle in a magnetic field $``$ which contains the contribution from the SU(2) gauge field $`𝒜`$ as well as the conventional electromagnetic field. We have recovered the Ehrenfest theorem as one of the examples of the corresponding principle in quantum mechanics. The term, $`(\mu _B\sigma B),`$ results from the non-uniform magnetic field. Its role was first realized in the Stern-Gerlach experiment, where a shaped magnet is used to generate a non-uniform magnetic field to split the beam of silver atoms. In the classical limit it is written as the interaction between the magnetic momentum $`\mu =\mu _B\sigma `$ and magnetic field. This spin force depends on the spin. Recently it is proposed that the force can generate a pure spin current if we assume $`B_z`$ is a constant.Zhang04xxx In the term $`F_g`$ we can also use $`\mu `$ to replace the spin. It is non-zero only when the electric and magnetic fields coexist, as suggested by Anandan and others.Anandan89 This term will play an essential role in generating spin Hall current in two-dimensional Rashba system, which we will discuss it later.
The last term, $`F_f`$, comes from the SU(2) gauge potential or spin-orbital coupling. As the force is related to the Planck constant it has no counterpart in classic mechanics and is purely quantum mechanic effect. The force is irrelevant of the magnetic field. In the classical limit we cannot simply use the magnetic momentum $`\mu `$ to replace the spin $`\sigma `$ in the potential $`𝒜`$. Otherwise the force vanishes. To see the physical meaning of the force, we write the spin force for a single electron on a quantum state in a compact form,
$$F_f=\frac{e^2\left|\right|}{4m^2c^4}𝐣_s^{}\times $$
(7)
where the spin current is defined conventionally, $`𝐣_s^{}=\frac{\mathrm{}}{4}\{𝐯,\sigma /\left|\right|\},`$ and in the last step the relation $`\{𝒜,\sigma \}=0`$ has been used. This is the main result in this paper. The force is proportional to the square of electric field $``$ and the spin current which polarization is projected along the field. It is important to note that an electron in a spin state perpendicular to the electric field will not experience any force. Comparing with a charged particle in a magnetic field, $`𝐣_c\times 𝐁,`$ where $`𝐣_c`$ is a charge current density, the spin force is nonlinear to the electric field and depends on the spin state of electron.
We discuss several examples relevant to the spin force. Though the gauge field $`𝒜`$ provides a spin-dependent magnetic field $`\times 𝒜`$, the Lorentz force caused by the field on the charge will vanish in a uniform electric field $`V=e.`$ The spin depend force $`F_g`$ also vanishes in the absence of magnetic field. Here we consider the motion of an electron confining in a two-dimensional plane subjected to a perpendicular electric field,
$$H=\frac{𝐩^2}{2m}+\lambda (𝐩_x\sigma _y𝐩_y\sigma _x)$$
(8)
where $`\lambda =\mathrm{}e/(4m^2c^2)`$ from Eq.(2). This can be regarded as counterpart of a charged particle in a magnetic field. On the other hand it has the same form of the Rashba coupling in a semiconductor heterojuction with the structural inversion asymmetry,Rashba60 where the spin-orbit coupling is induced by the offsets of valence bands at the interfaces and the structure inversion asymmetry.Winkler03 A typical value of this coefficient $`\lambda `$ is of order $`10^4c`$ ($`c`$ the speed of light), and can be adjusted by an external field.Nitta97apl Because of the spin-orbit coupling the electron spin will precess with time,
$$\frac{d\sigma (t)}{dt}=\frac{2\lambda }{\mathrm{}}\sigma (t)\times \left(𝐩\times \widehat{z}\right).$$
(9)
Since the momentum $`𝐩`$ is a good quantum number, without loss of generality, we take $`p=p_x,`$ just along the x direction. Correspondingly the wave function in the position space $`r|\mathrm{\Phi }=\mathrm{exp}(ip_xx)\chi _s/\sqrt{L}`$ where $`\chi _s`$ is the initial spin state. Equivalently the spin-orbit coupling provides an effective magnetic field along the $`y`$ direction, $`𝐁_{eff}=\lambda 𝐩_x\widehat{y}`$. This problem can be solved analytically, and the electron spin precesses in the spin x-z plane,Shen04prb $`\sigma _x(t)=\sigma _x\mathrm{cos}\omega _ct\sigma _z\mathrm{sin}\omega _ct,`$ $`\sigma _z(t)=\sigma _z\mathrm{cos}\omega _ct+\sigma _x\mathrm{sin}\omega _ct,`$ and $`\sigma _y(t)=\sigma _y`$ where the Larmor frequency $`\omega _c=2p_x\lambda /\mathrm{}.`$ The spin $`\sigma _z(t)`$ varies with time and the spin current is always along the x direction, $`𝐣_s^z=\frac{\mathrm{}}{2}\frac{p_x}{m}\left(\sigma _z\mathrm{cos}\omega _ct+\sigma _x\mathrm{sin}\omega _ct\right)\widehat{x}`$ where $`\mathrm{}`$ means the expectation value over an initial state $`r|\mathrm{\Phi }`$. As a result the spin transverse force on the spin is always perpendicular to the x direction. Correspondingly the kinetic velocity$`v_x`$ and $`v_y`$ at a time $`t`$ are
$`v_x_t`$ $`=`$ $`{\displaystyle \frac{p_x}{m}}+\lambda \sigma _y;`$ (10a)
$`v_y_t`$ $`=`$ $`\lambda \left(\sigma _x\mathrm{cos}\omega _ct\sigma _z\mathrm{sin}\omega _ct\right),`$ (10b)
respectively. Though $`p_y=0`$ the kinetic velocity $`v_y_t`$ oscillates with the frequency $`\omega _c`$ while $`v_x_t`$ remains constant. The y-component of the position is
$$y_t=y_{t=0}\frac{\mathrm{}}{p_x}\mathrm{sin}\frac{\omega _ct}{2}\left(\sigma _x\mathrm{cos}\frac{\omega _ct}{2}\sigma _z\mathrm{sin}\frac{\omega _ct}{2}\right).$$
(11)
If the initial state is polarized along y direction the electron spin does not vary with time as it is an energy eigen state of the system, as discussed by Datta and Das.Datta90 In this case the spin current $`𝐣_s^z`$ carried by the electron is always zero and the spin transverse force is zero. Thus $`v_y_t=0.`$ If the initial state is along the spin z direction at $`t=0`$, i.e., $`\sigma _z=s=\pm 1,`$ it is found that $`v_y_{t,s}=s\lambda \mathrm{sin}\omega _ct.`$ Different spins will move in opposite directions. It can be understood that the spin precession makes the spin current which polarization is projected along the electric field changes with time such that the spin force along the y direction also oscillates with the frequency $`\omega _c`$. This force will generate a non-zero velocity of electron oscillating along the y direction. Though $`v_y_{t,s=1}v_y_{t,s=1}=2\lambda \mathrm{sin}\omega _ct,`$ the velocity does not contribute to the spin current along the y direction, i.e. $`\{v_y,\sigma _z\}_{t,s}=0.`$ The trajectory oscillates with the time. The amplitude of the oscillation is $`\mathrm{}/p_x`$ and the frequency is $`\omega _c=2p_x\lambda /\mathrm{}.`$ For a typical two dimensional electron gas the electron density is $`n_e=10^{11}10^{12}/`$cm<sup>2</sup> and the wave length near the Fermi surface $`\mathrm{}/p_x310`$nm. For a typical Rashba coupling $`\lambda =10^4c`$, $`\omega _c=0.31.0\times 10^{14}s.`$ The rapid oscillation of the electron wave packet is known in literatures as the zitterbewegung of electron as a relativistic effect, which is usually regarded as a result of admixture of the positron state in electron wave packet as a relativistic effect.Schrodinger30 In semiconductors the Rashba coupling reflects the admixture of the particle and hole states in the conduction and valence bands. Recently Schliemann et al obtained the solution of the trajectory and proposed that this effect can be observed in a III-V zinc-blende semiconductor quantum wells.Schliemann05prl In the p-doped semiconductors described by the Luttinger model there exists also a spin force, and will generate the zitterbewegung as calculated by Jiang et al.Jiang04xxx Though the spin transverse force on a moving spin is very analogous to the Lorentz force on a moving charge, because of spin precession, its effect is completely different with the motion of a charged particle in a magnetic field, where the amplitude of the Lorentz force is constant and the charged particle will move in a circle. The zitterbewegung of the electronic wave package near the boundary will cause some edge effect as shown in recent numerical calculations.edge The edge effect is determined by the electron momentum. The smaller the momentum, the larger the edge effect. The amplitude and frequency of the zitterbewegung satisfy a relation that $`\left(\mathrm{}/2p_x\right)(2p_x\lambda /\mathrm{})=\lambda `$, which is the amplitude of oscillation of the velocity $`v_y`$. In Fig.1 it is illustrated that for two electrons with different spins experience opposite forces in an electric field.
Furthermore we consider a two-dimensional electron gas lacking both the bulk and structural inversion symmetries. A Dresselhaus term $`\beta (𝐩_x\sigma _x𝐩_y\sigma _y)`$ will be included in the total Hamiltonian in Eq. (8).Dresselhaus55 In this model the spin force formula givesLi05prb
$$F_f=\frac{4m^2}{\mathrm{}^2}(\lambda ^2\beta ^2)(j_s^z\times \widehat{z})$$
(12)
for each moving electron. First of all the force disappears at the symmetric point of $`\lambda =\pm \beta .`$ At this point the operator $`\sigma _x\pm \sigma _y`$ is a good quantum number an there is no spin flip in the system. For $`\lambda \pm \beta ,`$ the moving electron will experience a spin-dependent force and the force will change its sign near $`\lambda =\beta `$. A heuristic picture from this formula is that when a non-zero spin current $`j_s^z`$ goes through this system the spin-orbit coupling exerts the spin transverse force on the spin current, and drives electrons to form a charge Hall current perpendicular to the spin current. The injected spin current can be generated in various ways, such as by the spin force $`(\mu 𝐁)`$ Zhang04xxx and circularly or linearly polarized light injection.Sipe ; Hankiewicz05xxx For instance we assume the spin current $`j_s^z`$ be generated by the linear polarized light injection which is proportional to the transition rate of from the valence band to the conduction band with finite momentum and the life time of electrons at the excited states. In the relaxation time $`\tau `$ approximation in a steady state the drift velocity orthogonal to the spin current is $`v_y=`$ $`\frac{4m}{\mathrm{}^2}(\lambda ^2\beta ^2)j_s^z\tau `$ if $`\tau `$ is not so long, i.e. $`2p_x\sqrt{\lambda ^2+\beta ^2}\tau <<\mathrm{}.`$ This non-zero drift velocity will form a Hall current orthogonal to the the spin current. This is the charge Hall effect driven by the spin current. In ferromagnetic metal or diluted magnetic semiconductors the charge current is spin polarized it can generate the spin polarized Hall current via the spin transverse force. Thus the spin transverse force can be also regarded a driven force of an anomalous Hall effectAHE and the spin-resolved Hall effect.Bulgakov99prl ; Li05prb A detailed calculation for this Hall conductance is given by the Kubo formula as a linear response to the field $`B=\mathrm{\Delta }Bx\widehat{z}`$. This field will generate a spin force, $`(g\mu _BB\sigma )=g\mu _B\mathrm{\Delta }B\sigma _z\widehat{x},`$ which will circulate a spin current along the x direction, and furthermore the spin-orbit coupling provides a driving force to generate a transverse charge current, $`j_{c,y}`$. In the clear limit the Hall conductance $`\sigma _{xy}=j_{c,y}/(g\mu _B\mathrm{\Delta }B)=0`$ for $`\lambda =\pm \beta `$ and $`(e/4\pi \mathrm{})(\lambda ^2\beta ^2)/\left|\lambda ^2\beta ^2\right|`$. However following Inoue et al and Mishchenko et alInoue05prb the inclusion of impurities scattering will suppress the Hall conductance completely just like the spin Hall effect. On the other hand numerical calculation in mesoscopic systems shows the existence of the effect.Hankiewicz05xxx Another example is the two-dimensional p-doped system with the cubic Rashba coupling,Schliemann05prb
$$H=\frac{𝐩^2}{2m}+i\alpha (𝐩_+^3\sigma _{}𝐩_{}^3\sigma _+)$$
(13)
where $`\sigma _\pm `$ are spin increasing and decreasing operators, and $`𝐩_\pm =𝐩_x\pm i𝐩_y`$. The spin force on the moving electron in this system is $`F_f=(2m𝐩\alpha /\mathrm{})^2(𝐣_s^z\times \widehat{z}).`$ The linear response theory gives the Hall conductance $`9e/(2\pi \mathrm{})`$ which is robust against the vertex correction from impurities scattering. Calculations by means of the Green-Keldysh function technique and linear response theory Zhang04xxx ; Hankiewicz05xxx show the existence of charge Hall effect driven by the spin current, and the Onsager relation between the charge Hall effect and its reciprocal. The key features of the numerical results are in good agreement with the picture of spin force qualitatively.
In conclusion, an electric field exerts a transverse force on a moving spin just like a magnetic field exerts a Lorentz force on a moving charge. This force is proportional to the square of electric field and the spin current with spin projected along the field. This is a purely relativistic quantum mechanical effect. As the origin of the force the spin current should be also observable physically. From the solution of the motion of a single electron in an electric field, the zitterbewegung of electronic wave packet in the spin-orbit coupling can be regarded as an explicit consequence of this force. Due to the similarity of this spin transverse force and the Lorentz force, the spin transverse force plays a similar role in the formation of the charge Hall effect driven by the spin current and the spin Hall effect driven by the charge current as the Lorentz force does in the Hall effect in a magnetic field.
The author thanks F. C. Zhang for helpful discussion. This work was supported by the Research Grant Council of Hong Kong (No.: HKU 7039/05P), and a CRCG grant of The University of Hong Kong.
|
warning/0506/nucl-th0506048.html
|
ar5iv
|
text
|
# Single 𝛼-particle orbits and Bose-Einstein condensation in 12C
## I Introduction
Four-nucleon correlations of the $`\alpha `$-cluster-type play an important role in nuclei. The microscopic $`\alpha `$-cluster model Wildermuth77 ; Brink66 ; Bertsch71 ; Fujiwara80 has succeeded in describing the structure of many states in light nuclei, in particular, around the threshold energy of breakup into constituent clusters. As for <sup>12</sup>C, detailed analyses were made by several authors with the microscopic $`3\alpha `$ cluster model about 25 years ago. The $`3\alpha `$ GCM (generator coordinate method) Uegaki77 and $`3\alpha `$ RGM (resonating group method) Fukushima77 calculations showed that the second $`0^+`$ state of <sup>12</sup>C ($`E_x`$=7.65 MeV), located at $`E_{3\alpha }`$=0.38 MeV above the $`3\alpha `$ threshold, has a loosely coupled $`3\alpha `$ structure, although the ground state is a shell-model-like compact state. On the other hand, special attention has been paid to an $`\alpha `$-type condensation in symmetric nuclear matter, analogue to the Bose-Einstein condensation for finite number of dilute bosonic atoms such as <sup>87</sup>Rb or <sup>23</sup>Na at very low temperature where all atoms occupy the lowest $`S`$ orbit Dalfovo99 . Several authors Ropke98 ; Beyer00 showed the possibility of such $`\alpha `$-particle condensation in low-density nuclear matter, although the ordinary pairing correlation can prevail at higher density. This result suggests that dilute $`\alpha `$ condensate states in finite nuclear system may exist in excited states as a weakly interacting gas of $`\alpha `$ particles.
Recently, a new $`\alpha `$-cluster wave function was proposed which is of the $`N\alpha `$ particle condensate type Tohsaki01
$`|\mathrm{\Phi }_{N\alpha }=(C_\alpha ^+)^N|\mathrm{vac},`$ (1)
$`𝒓_1\mathrm{}𝒓_N|\mathrm{\Phi }_{N\alpha }𝒜\left\{e^{\nu \left(𝒓_1^2+\mathrm{}+𝒓_N^2\right)}\varphi (\alpha _1)\mathrm{}\varphi (\alpha _N)\right\},`$ (2)
where $`C_\alpha ^+`$ is the $`\alpha `$-particle creation operator, $`\varphi (\alpha )`$ denotes the internal wave function of the $`\alpha `$ cluster, $`𝒓_i`$ is the center-of-mass coordinate of the i-th $`\alpha `$ cluster, and $`𝒜`$ presents the antisymmetrizer among the nucleons belonging to different $`\alpha `$ clusters. The important characteristic of the wave function is that the center-of-mass motion of each $`\alpha `$ cluster is of $`S`$-wave type. Applications of the condensate-type wave function to <sup>12</sup>C and <sup>16</sup>Tohsaki01 indicated that the second $`0^+`$ state of <sup>12</sup>C ($`E_x`$=7.65 MeV) and fifth $`0^+`$ state of <sup>16</sup>O ($`E_x`$=14.0 MeV), around the $`3\alpha `$ and $`4\alpha `$ threshold, respectively, are conjectured to be dilute $`N\alpha `$ condensate states, which are quite similar to the Bose-Einstein condensation of bosonic atoms at very low temperature. The calculated nuclear radii for both of those states are about 4 fm, significantly larger than that for the ground state (about 2.5 fm). As for <sup>12</sup>C, a detailed analysis with a deformed alpha condensate wave function, slightly different from the spherical one in Eq. (2), was performed to investigate the structure of the $`0_1^+`$ and $`0_2^+`$ states. Funaki02 It was found that each of the $`0_2^+`$ wave functions obtained by the $`3\alpha `$ GCM and RGM calculations has a large squared overlap value of more than $`90\%`$ with the single $`3\alpha `$ condensate wave function.
The above-mentioned results for <sup>12</sup>C and <sup>16</sup>O lead us to the further intriguing issue that dilute $`\alpha `$-cluster states with $`J^\pi =0^+`$ near the $`N\alpha `$ threshold may exist in other heavier $`4N`$ self-conjugate nuclei. The Gross-Pitaevskii-equation approach Yamada04 is useful to explore such dilute multi $`\alpha `$ systems, because this equation Pitaevskii61 , based on mean field theory, has succeeded in describing the structure of the Bose-Einstein condensation for dilute neutral atomic systems, for example, <sup>87</sup>Rb or <sup>23</sup>Na, at very low temperature, trapped by an external magnetic field. Dalfovo99 The present authors Yamada04 applied the Gross-Pitaevskii equation to self-conjugate $`4N`$ nuclei. They found that 1) there exists a critical number of $`\alpha `$ bosons that the dilute $`N\alpha `$ system can sustain as a self-bound nucleus, and 2) the Coulomb-potential barrier plays an important role to confine such dilute $`N\alpha `$-particle condensate states.
It is interesting to explore also the possibility of the $`\alpha `$ condensate states with non-zero angular momentum in <sup>12</sup>C. The old theoretical calculations based on the microscopic $`3\alpha `$ cluster model Fujiwara80 ; Uegaki77 ; Fukushima77 suggested the existence of a $`2_2^+`$ state of <sup>12</sup>C at around $`E_{3\alpha }3`$ MeV above the $`3\alpha `$ threshold, the structure of which is similar to the $`0_2^+`$ state except for the angular momentum. Quite recently the $`2_2^+`$ state was observed at $`E_{3\alpha }=2.6\pm 0.3`$ MeV with the alpha decay width $`\mathrm{\Gamma }_\alpha =1.0\pm 0.3`$ MeV. Itoh04 The $`\alpha `$ condensate-type wave function with axially deformation Funaki04 was applied to study the structure of the $`2_2^+`$ state with help of the method of ACCC (analytic continuation in the coupling constant) Kukulin77 . They found that the $`2_2^+`$ state has a large overlap with the single condensate wave function of $`3\alpha `$ gas-like structure, the squared value of which amounts to about $`88\%`$. This result implies that the $`2_2^+`$ state has a similar structure as the $`0_2^+`$ state, namely, dilute $`3\alpha `$ condensation.
Here, it is an intriguing problem to discuss whether the $`0_2^+`$ and $`2_2^+`$ states of <sup>12</sup>C are ideal dilute $`\alpha `$ condensates or not. The condensate-type $`\alpha `$-cluster wave function in Eq. (2) has succeeded in describing the $`0_2^+`$ state of <sup>12</sup>C. This result, however, does not necessarily mean that the $`0_2^+`$ state of <sup>12</sup>C is an ideal $`\alpha `$-condensate state. If the $`0_2^+`$ state of <sup>12</sup>C is an ideal dilute $`\alpha `$-condensate, the single $`\alpha `$-particle orbit in the state should be of the zero-node long-ranged $`S`$-wave type with an occupation probability of $`100\%`$, as suggested from the Gross-Pitaevskii-equation approach Yamada04 . The antisymmetrizer $`𝒜`$ in Eq. (2) generally perturbs the single $`\alpha `$ motion in the nucleus, and one should remember that the condensate-type wave function can also describe the shell-model-like compact structure of the ground state of <sup>12</sup>C. The effect of the antisymmetrizer should have a close relation to the rms radius of the nucleus or the distance between $`2\alpha `$ clusters in a nucleus. The ideal $`3\alpha `$ condensate state is expected to be realized if the distance between two arbitrary $`\alpha `$ clusters is large enough so that the effect of the antisymmetrizer can be neglected. The calculated nuclear radius for the $`0_2^+`$ state of <sup>12</sup>C, about 4 fm Tohsaki01 , suggests that the action of the antisymmetrizer is weakened significantly in that state. In order to give more decisive theoretical evidence that the $`0_2^+`$ state of <sup>12</sup>C as well as the $`2_2^+`$ state has dilute $`3\alpha `$ condensation structure, it is needed to study quantitatively the bosonic properties such as single $`\alpha `$-particle orbits and corresponding occupation probabilities, starting from the microscopic wave function.
The first attempt to derive the $`\alpha `$-boson properties for $`0^+`$ states in <sup>12</sup>C from a microscopic model was performed in Ref. Matsuura04 , where the $`3\alpha `$ RGM equation was solved in terms of the correlated Gaussian basis with the stochastic variational method. Although they formulated a derivation of the $`3\alpha `$ boson wave function starting from the microscopic $`3\alpha `$ wave function, the $`\alpha `$ bosonic properties of <sup>12</sup>C were studied not with the $`3\alpha `$ bosonic wave function but with the normalized spectroscopic amplitude, because the derivation of the $`3\alpha `$ boson wave function is numerically difficult due to the non-local properties of the norm kernel. Although the normalized spectroscopic amplitude seems to be a good approximation for the boson wave function in the region where the effect of the antisymmetrizer is negligible, the approximation becomes worse when the spatial overlap of the $`3\alpha `$ clusters becomes larger. It is requested to demonstrate quantitatively how good the approximation is for the $`0_2^+`$ state within their framework.
The purposes in the present paper are twofold. First we study the bosonic properties such as single $`\alpha `$-particle orbits and occupation probabilities for the $`0^+`$ and $`2^+`$ states in <sup>12</sup>C with direct use of the wave function obtained by the $`3\alpha `$ OCM (orthogonality condition model) Saito68 . The OCM is a semi-microscopic model and a simple version of the RGM, taking into account properly the antisymmetrization among nucleons, which successfully describes the structure of <sup>12</sup>C. Fujiwara80 ; Horiuchi74 ; Kato89 ; Kurokawa04 The second purpose is to explore the possibility of the dilute $`3\alpha `$ condensation with negative parity within the present framework. The $`3_1^{}`$ ($`1_1^{}`$) state of <sup>12</sup>C at $`E_{3\alpha }=2.37`$ (3.57) MeV above the $`3\alpha `$ threshold appears at the same energy region as that for the $`0_2^+`$ ($`E_x`$=0.38 MeV) and $`2_2^+`$ (2.6 MeV) states. According to the old theoretical study based on the $`3\alpha `$ GCM and RGM calculations Uegaki77 ; Fukushima77 , the nuclear radius of the $`3^{}`$ state is intermediate between a compact shell-model-like state ($`0_1^+`$) and a loosely coupled $`3\alpha `$ cluster state ($`0_2^+`$), while the $`1^{}`$ state has a radius only a little smaller than that of the $`0_2^+`$ state. Thus, it is quite interesting to study the bosonic properties for the negative parity states, as well.
In order to clarify the relationship between the $`N\alpha `$ boson wave function and the one of $`N\alpha `$ OCM, we first outline a way of mapping of the microscopic $`N\alpha `$ wave function onto the $`N\alpha `$ boson wave function, and derive the equation of motion for $`N\alpha `$ bosons from the $`N\alpha `$ RGM equation in Sec. II. The $`N\alpha `$ OCM equation is illustrated as an approximation of the $`N\alpha `$ boson equation. The $`N\alpha `$ OCM wave function, thus, has bosonic properties. The $`3\alpha `$ OCM equation is solved properly with modern numerical techniques. The calculated single $`\alpha `$-particle orbits and occupation probabilities in <sup>12</sup>C are discussed in Sec. III. Finally, we give the summary in Sec. IV.
## II Formulation
A way of mapping of the microscopic $`N\alpha `$ wave function onto the $`N\alpha `$ boson wave function is illustrated in order to derive the equation of motion for the $`N\alpha `$ bosons. The $`N\alpha `$ OCM equation is given as an approximation of the equation of motion for $`N\alpha `$ bosons. We formulate the evaluation of the single-$`\alpha `$ orbits and occupation numbers from the $`N\alpha `$ OCM wave function together with other physical quantities. Finally, we give an outline of how to solve the $`3\alpha `$ OCM equation for <sup>12</sup>C with a phenomenological $`\alpha \alpha `$ potential.
### II.1 Mapping of the fermionic $`N\alpha `$ wave function onto a $`N\alpha `$ boson wave function
In the microscopic $`N\alpha `$ cluster model, the total wave function, $`\mathrm{\Psi }_J^{(F)}`$, with the total angular momentum $`J`$ is given as
$`\mathrm{\Psi }_J^{(F)}`$ $`=`$ $`𝒜\left\{{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi _{\alpha _i}^{(int)}\chi _J\left(𝒓\right)\right\},`$ (3)
$`=`$ $`{\displaystyle 𝑑𝒂\mathrm{\Psi }_J^{(F)}(𝒂)\chi (𝒂)},`$ (4)
where $`\varphi _\alpha ^{(int)}`$ denotes the intrinsic wave function of the $`\alpha `$ particle with the simple $`(0s)^4`$ shell-model configuration, and $`\chi _J`$ represents the relative wave function with a set of the relative (Jacobi) coordinates, $`𝒓=\{𝒓_1,𝒓_2,\mathrm{},𝒓_{N1}\}`$, with respect to the c.m. of $`\alpha `$ clusters. The antisymmetrization among $`4N`$ nucleons is properly taken into account in terms of the operator $`𝒜`$. The function $`\mathrm{\Psi }_J^{(F)}(𝒂)`$ in Eq. (4) is defined as
$`\mathrm{\Psi }_J^{(F)}(𝒂)=𝒜\left\{{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi _{\alpha _i}^{(int)}{\displaystyle \underset{j=1}{\overset{N1}{}}}\delta (𝒓_j𝒂_j)\right\},`$ (5)
which describes the $`\alpha `$-cluster state located at the relative positions specified by the Jacobi parameter coordinate $`𝒂=\{𝒂_1,𝒂_2,\mathrm{},𝒂_{N1}\}`$.
The total Hamiltonian for $`4N`$ fermions is given as
$`H={\displaystyle \underset{i=1}{\overset{4N}{}}}t_iT_{cm}+{\displaystyle \underset{i<j=1}{\overset{4N}{}}}\left(\upsilon _{ij}+\upsilon _{ij}^{Coul}\right),`$ (6)
where $`t_i`$ and $`T_{cm}`$ denote, respectively, the kinetic energy operator of the i-th nucleon and of the center-of-mass of the total system. The nucleon-nucleon interaction (Coulomb interaction) between the i-th and j-th nucleons is expressed as $`\upsilon _{ij}`$ ($`\upsilon _{ij}^{Coul}`$).
The Schrödinger equation for the fermionic $`N\alpha `$ system is
$`H\mathrm{\Psi }_J^{(F)}=E\mathrm{\Psi }_J^{(F)}.`$ (7)
Substituting the total wave function of Eq. (4) into Eq. (7), we obtain the equation of motion for the relative wave function $`\chi _J`$,
$`{\displaystyle 𝑑𝒂^{}\left\{H(𝒂,𝒂^{})EN(𝒂,𝒂^{})\right\}\chi _J(𝒂^{})}=0,\mathrm{or}\left(𝒩\right)\chi _J=0,`$ (8)
which is called the RGM (resonating group method) equation Horiuchi86 . The Hamiltonian and norm kernels, $`H(𝒂,𝒂^{})`$ and $`N(𝒂,𝒂^{})`$, are defined as
$`\left\{\begin{array}{c}H(𝒂,𝒂^{})\\ N(𝒂,𝒂^{})\end{array}\right\}=\mathrm{\Psi }_J^{(F)}(𝒂)\left\{\begin{array}{c}H\\ 1\end{array}\right\}\mathrm{\Psi }_J^{(F)}(𝒂^{}).`$ (13)
Recalling the normalization condition $`\mathrm{\Psi }_J^{(F)}\mathrm{\Psi }_J^{(F)}=1`$ for the total wave function in Eq. (3), the normalization of $`\chi _J`$ in Eq. (8) is given by
$`{\displaystyle 𝑑𝒂𝑑𝒂^{}\chi _J^{}(𝒂^{})N(𝒂,𝒂^{})\chi _J(𝒂^{})}=1.`$ (14)
This suggests that an $`N\alpha `$ boson wave function $`\mathrm{\Phi }_J^{(B)}`$ corresponding to the fermionic wave function $`\mathrm{\Psi }_J^{(F)}`$ in Eq. (3) should be taken to be
$`\mathrm{\Phi }_J^{(B)}(𝒂)=𝒩^{1/2}\chi _J={\displaystyle 𝑑𝒂^{}N^{1/2}(𝒂,𝒂^{})\chi _J(𝒂^{})},{\displaystyle 𝑑𝒂\left|\mathrm{\Phi }_J^{(B)}(𝒂)\right|^2}=1,`$ (15)
where $`N^{1/2}`$ is defined as
$`{\displaystyle 𝑑𝒂^{\prime \prime }N^{1/2}(𝒂,𝒂^{\prime \prime })N^{1/2}(𝒂^{\prime \prime },𝒂^{})}=N(𝒂,𝒂^{}).`$ (16)
It is noted that the boson wave function $`\mathrm{\Phi }_J^{(B)}(𝒂)`$ has only the Jacobi coordinates $`𝒂=\{𝒂_1,𝒂_2,\mathrm{},𝒂_{N1}\}`$ of the $`N\alpha `$ system and the internal coordinates in the $`\alpha `$ cluster are integrated out completely. In addition, $`\mathrm{\Phi }_J^{(B)}(𝒂)`$ is totally symmetric for any two-$`\alpha `$-cluster exchange. Thus, it has bosonic property. From the RGM equation (8), $`\mathrm{\Phi }_J^{(B)}`$ should satisfy the following equation
$`\left(𝒩^{1/2}𝒩^{1/2}E\right)\mathrm{\Phi }_J^{(B)}=0,`$ (17)
Thus, we can interpret $`𝒩^{1/2}𝒩^{1/2}`$ as the $`N\alpha `$ boson Hamiltonian, and Equation (17) is the equation of motion for the $`N\alpha `$ boson wave function.
If the eigenvalue problem for the norm kernel
$`𝒩u_\lambda ={\displaystyle 𝑑𝒂^{}N(𝒂,𝒂^{})u_\lambda (𝒂^{})}=\lambda u_\lambda (𝒂)`$ (18)
is solved, the boson wave function $`\mathrm{\Phi }_J^{(B)}(𝒂)`$ in Eq. (15), then, is obtained as
$`\mathrm{\Phi }_J^{(B)}(𝒂)={\displaystyle \underset{\lambda }{}}\sqrt{\lambda }u_\lambda (𝒂)u_\lambda \chi _J.`$ (19)
The eigenvalue of the norm kernel, $`\lambda `$, is non-negative, and the eigenfunction $`u_\lambda `$ with $`\lambda =0`$ is called the Pauli-forbidden state, which satisfies the condition $`𝒜\{_{i=1}^N\varphi _{\alpha _i}^{(int)}u_\lambda \}=0`$. The boson wave function $`\mathrm{\Phi }_J^{(B)}`$, thus, has no component of the Pauli-forbidden state,
$`u_\lambda \mathrm{\Phi }_J^{(B)}=0\mathrm{for}u_\lambda \mathrm{with}\lambda =0.`$ (20)
### II.2 $`N\alpha `$ orthogonality condition model (OCM) with bosonic properties
In the previous section, we mapped the fermionic wave function onto the $`N\alpha `$ boson wave function $`\mathrm{\Phi }_J^{(B)}`$ within the framework of the resonating group method (RGM). The boson wave function has the following properties: 1) $`\mathrm{\Phi }_J^{(B)}`$ is totally symmetric for any 2$`\alpha `$-particle exchange, 2) $`\mathrm{\Phi }_J^{(B)}`$ satisfies the equation motion in Eq. (17), and 3) $`\mathrm{\Phi }_J^{(B)}`$ is orthogonal to the Pauli forbidden state (see Eq. (20)). In order to obtain the boson wave function, we need to solve the RGM equation (8) and the eigenvalue equation of the norm kernel (18) or to solve directly the equation of motion (17). Solving the eigenvalue equation of the norm kernel, however, is difficult in general even for the 3$`\alpha `$ case. In addition, computational problems are encountered for solving the $`N\alpha `$ RGM equation (8) for $`N4`$. Thus, it is requested to use more feasible frameworks for the study of the bosonic properties and the amount of $`\alpha `$ condensation for the $`N\alpha `$ system. In the present study, we take the orthogonality condition model (OCM) Saito68 as one of the more feasible frameworks. The OCM scheme, which is an approximation to the RGM, is known to describe nicely the structure of low-lying states in light nuclei Fujiwara80 ; Saito68 ; Horiuchi74 ; Kato89 ; Kurokawa04 . The essential properties of the $`N\alpha `$ boson wave function $`\mathrm{\Phi }_J^{(B)}`$, as mentioned above, can be taken into account in OCM in a simple manner. We will demonstrate this briefly.
In OCM, the $`\alpha `$ cluster is treated as a point-like particle. We approximate the $`N\alpha `$ boson Hamiltonian (non-local) in Eq. (17) by an effective (local) one $`H^{(\mathrm{OCM})}`$,
$`𝒩^{1/2}𝒩^{1/2}H^{(\mathrm{OCM})}`$ (21)
$`H^{(\mathrm{OCM})}{\displaystyle \underset{i=1}{\overset{N}{}}}T_iT_{cm}+{\displaystyle \underset{i<j=1}{\overset{N}{}}}V_{2\alpha }^{eff}(i,j)+{\displaystyle \underset{i<j<k=1}{\overset{N}{}}}V_{3\alpha }^{eff}(i,j,k),`$ (22)
where $`T_i`$ denotes the kinetic energy of the i-th $`\alpha `$ cluster, and the center-of-mass kinetic energy $`T_{cm}`$ is subtracted from the Hamiltonian. The effective 2$`\alpha `$ and 3$`\alpha `$ potentials are presented as $`V_{2\alpha }^{eff}`$ and $`V_{3\alpha }^{eff}`$, respectively. Referring to the RGM framework in Eqs. (17)$``$(20), the equation of the relative motions for the $`N\alpha `$ particles with $`H^{(\mathrm{OCM})}`$, called the OCM equation, is expressed as
$`\left\{H^{(\mathrm{OCM})}E\right\}\mathrm{\Phi }_J=0,`$ (23)
$`u_F\mathrm{\Phi }_J=0,`$ (24)
where $`u_F`$ denotes the Pauli-forbidden state, satisfying the following condition
$`𝒜\{{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi _{\alpha _i}^{(int)}u_F\}=0.`$ (25)
The orthogonality condition in Eq. (24) corresponds to Eq. (20) in Sec. IIA. The bosonic property of the wave function $`\mathrm{\Phi }_J`$ can be taken into account by symmetrizing the wave function with respect to any 2$`\alpha `$-particle exchange,
$`\mathrm{\Phi }_J=𝒮\mathrm{\Phi }_J(1,2,\mathrm{},N),`$ (26)
where $`𝒮`$ denotes the symmetrization operator, $`𝒮=(1/\sqrt{N!})_𝒫𝒫`$, where the sum runs over all permutations $`𝒫`$ for the $`N\alpha `$ particles. It is noted that the complete overlapped state of the $`3\alpha `$ particles is forbidden within the present framework because of the Pauli-blocking effect in Eq. (24), although we take into account the bosonic statistics for the constituent $`\alpha `$ particles. In the next section, we will demonstrate i) how to solve the OCM equation and ii) what kind of effective $`\alpha `$-$`\alpha `$ potential we should choose in $`H^{(\mathrm{OCM})}`$ for the $`3\alpha `$ case of <sup>12</sup>C.
Here, it is useful to define various quantities characterizing the structure of the $`N\alpha `$ system with use of the $`N\alpha `$ boson wave function $`\mathrm{\Phi }_J`$ obtained by solving the OCM equation in Eqs. (23) and (24). The single $`\alpha `$-particle density is defined as
$`\rho (𝒓)=\mathrm{\Phi }_J{\displaystyle \underset{i=1}{\overset{N}{}}}\delta (𝒓𝒓_i^{(cm)})\mathrm{\Phi }_J,`$ (27)
where $`𝒓_i^{(cm)}`$ is the coordinate of the i-th $`\alpha `$ particle measured from the center-of-mass coordinate of the total system. The nuclear root-mean-square radius is given as
$`\sqrt{r_N^2}=\sqrt{r_\alpha ^2+1.71^2},`$ (28)
$`\sqrt{r_\alpha ^2}={\displaystyle 𝑑𝒓r^2\rho (𝒓)},`$ (29)
where we take into account the finite size effect of the $`\alpha `$ particle. The correlation functions with respect to the $`\alpha `$-$`\alpha `$ relative coordinate $`𝒓_{\alpha \alpha }`$ as well as the relative coordinate between one of the $`\alpha `$ particles and the remaining $`(N1)\alpha `$ system $`𝒓_{\alpha (N1)\alpha }`$ are given as
$`f_{\alpha \alpha }(r)=\mathrm{\Phi }_J\delta (𝒓𝒓_{\alpha \alpha })\mathrm{\Phi }_J,`$ (30)
$`f_{\alpha (N1)\alpha }(r)=\mathrm{\Phi }_J\delta (𝒓𝒓_{\alpha (N1)\alpha })\mathrm{\Phi }_J,`$ (31)
where the way of choosing the coordinates, $`𝒓_{\alpha \alpha }`$ and $`𝒓_{\alpha (N1)\alpha }`$, is arbitrary in the set of Jacobi coordinates of the $`N\alpha `$ particles because of the totally symmetrization for $`\mathrm{\Phi }_J`$. The root-mean-square (rms) distances with respect to $`𝒓_{\alpha \alpha }`$ and $`𝒓_{\alpha (N1)\alpha }`$ are, respectively, given by
$`\sqrt{r_{\alpha \alpha }^2}=\left[\mathrm{\Phi }_J𝒓_{\alpha \alpha }^2\mathrm{\Phi }_J\right]^{1/2},`$ (32)
$`\sqrt{r_{\alpha (N1)\alpha }^2}=\left[\mathrm{\Phi }_J𝒓_{\alpha (N1)\alpha }^2\mathrm{\Phi }_J\right]^{1/2}.`$ (33)
The reduced width amplitude for the $`\alpha `$-$`(N1)\alpha `$ part is defined as
$`𝒴_{\mathrm{}LJ}(r)=r\times \left[Y_L(𝒓)\varphi _{\mathrm{}}\left((N1)\alpha \right)\right]_J\mathrm{\Phi }_J,`$ (34)
where $`𝒓`$ denotes the relative coordinate between the $`\alpha `$ particle and the $`(N1)\alpha `$ nucleus, and $`\varphi _{\mathrm{}}\left((N1)\alpha \right)`$ represents the wave function of the $`(N1)\alpha `$ nucleus with total angular momentum $`\mathrm{}`$ which is obtained by solving the OCM equation for the $`(N1)\alpha `$ system. The integration in Eq. (34) is done over all of the relative (Jacobi) coordinates for the $`N\alpha `$ system except for the radial part of $`𝒓`$.
In order to discuss the bosonic properties such as the degree of $`\alpha `$ condensation in a nucleus, it is needed to extract information on the single $`\alpha `$-particle orbits and its occupation probabilities in the nucleus from the total wave function $`\mathrm{\Phi }_J`$. The one-particle density matrix for the $`N\alpha `$ system is very useful for this Matsuura04 . Defining the one-particle density operator as
$`𝒟(𝒓,𝒓^{})={\displaystyle \underset{i=1}{\overset{N}{}}}\delta (𝒓_𝒊^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓^{})\delta (𝒓_𝒊^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓),`$ (35)
then, the single $`\alpha `$-particle density matrix is given as
$`\rho (𝒓,𝒓^{})`$ $`=`$ $`\mathrm{\Phi }_J𝒟(𝒓,𝒓^{})\mathrm{\Phi }_J,`$ (36)
$`=`$ $`N\times \mathrm{\Phi }_J\delta (𝒓_\mathrm{𝟏}^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓^{})\delta (𝒓_\mathrm{𝟏}^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓)\mathrm{\Phi }_J,`$ (37)
where $`𝒓_i^{(cm)}`$ is the same as that in Eq. (27). It is noted that the diagonal matrix element reduces to the single $`\alpha `$-particle density in Eq. (27): $`\rho (𝒓,𝒓^{}=𝒓)=\rho (𝒓)`$. The single $`\alpha `$-particle orbit and its occupation number in the nucleus can be evaluated by solving the eigenvalue equation of the one-particle density matrix
$`{\displaystyle 𝑑𝒓^{}\rho (𝒓,𝒓^{})\phi _\mu (𝒓^{})}=\mu \phi _\mu (𝒓),`$ (38)
where the eigenvalue $`\mu `$ presents the occupation number. The eigenfunction $`\phi _\mu `$ denotes the single-$`\alpha `$ orbital wave function in a nucleus with the argument of the intrinsic coordinate ($`𝒓_\alpha ^{(cm)}`$) of an arbitrary $`\alpha `$ particle in a nucleus measured from the center-of-mass coordinate. The ratio $`\mu /N`$ represents the occupation probability of an $`\alpha `$ particle in the orbit $`\phi _\mu `$. The spectrum of the occupation probabilities offers important information about the occupancy of the single $`\alpha `$-particle orbit in a nucleus. If each of the $`N\alpha `$ particles in an $`N\alpha `$-boson state is occupied by only one orbit, the occupation probability for this orbit becomes $`100\%`$.
The <sup>8</sup>Be ($`2\alpha `$) system is a good example to demonstrate the characteristic of the single-$`\alpha `$ orbital wave function. From the definition of Eqs. (37) and (38), the $`L_\alpha `$-wave single-$`\alpha `$ orbit in the <sup>8</sup>Be($`J^\pi `$) state corresponds to the relative wave function (which is obtained by solving the $`2\alpha `$ OCM equation with $`J`$ ($`=L_\alpha `$) in Eqs. (23) and (24)), scaling to $`1/2`$ with respect to the relative coordinate between the $`2\alpha `$ clusters. Then, the occupation probability becomes exactly (mathematically) $`100\%`$ for any $`L_\alpha `$-value.
The radial behavior of the $`L_\alpha `$-wave single-$`\alpha `$ orbit, $`\phi _\mu (r_\alpha ^{(cm)})`$, in Eq. (38) generally has a close relationship with that of the reduced width amplitude, $`𝒴_{\mathrm{}LJ}(r_{\alpha (N1)\alpha })`$, in Eq. (34). This is due to the fact that both represent the behavior of the single-$`\alpha `$-particle motion in a nucleus in which all degrees of freedom of the other $`\alpha `$ particles are integrated out, and $`𝒓_\alpha ^{(cm)}`$ is given as $`𝒓_\alpha ^{(cm)}=\frac{N1}{N}\times 𝒓_{\alpha (N1)\alpha }`$.
The momentum distribution of the single $`\alpha `$ particle is also helpful for the study of $`\alpha `$ condensation in a nucleus Matsuura04 . It is defined as a double Fourier transformation of the one-particle density matrix
$`\rho (k)={\displaystyle 𝑑𝒓^{}𝑑𝒓\frac{e^{i𝒌𝒓^{}}}{(2\pi )^{3/2}}\rho (𝒓,𝒓^{})\frac{e^{i𝒌𝒓}}{(2\pi )^{3/2}}},`$ (39)
$`{\displaystyle 𝑑𝒌\rho (k)}=1,`$ (40)
It is reminded that $`\rho (k)`$ would have a $`\delta `$-function like peak around $`k=0`$ for an ideal dilute condensed state of infinite size.
### II.3 3$`\alpha `$ OCM for <sup>12</sup>C
In the previous section, we outlined the $`N\alpha `$ orthogonality condition model (OCM) and discussed how to extract the properties and the amount of $`\alpha `$ condensation in the $`N\alpha `$ system. Here, we will apply the OCM framework to the $`3\alpha `$ system of <sup>12</sup>C.
The total wave function of <sup>12</sup>C (the total angular momentum $`J`$) within the frame of the $`3\alpha `$ OCM is presented as
$`\mathrm{\Phi }_J({}_{}{}^{12}\mathrm{C})=\mathrm{\Phi }_J^{(12,3)}+\mathrm{\Phi }_J^{(23,1)}+\mathrm{\Phi }_J^{(31,2)},`$ (41)
where $`\mathrm{\Phi }_J^{(12,3)}`$ denotes the relative wave function of the $`3\alpha `$ system with the Jacobi-coordinate system shown in Fig. 1(a), and others are self-explanatory. In the present study, $`\mathrm{\Phi }_J({}_{}{}^{12}\mathrm{C})`$ is expanded in terms of the Gaussian basis Kamimura88 ,
$`\mathrm{\Phi }_J({}_{}{}^{12}\mathrm{C})={\displaystyle \underset{c}{}}{\displaystyle \underset{\nu ,\mu }{}}A_c(\nu ,\mu )\mathrm{\Phi }_c^{3\alpha }(\nu ,\mu ),`$ (42)
$`\mathrm{\Phi }_c^{3\alpha }(\nu ,\mu )=\mathrm{\Phi }_c^{(12,3)}(\nu ,\mu )+\mathrm{\Phi }_c^{(23,1)}(\nu ,\mu )+\mathrm{\Phi }_c^{(31,2)}(\nu ,\mu ),`$ (43)
$`\mathrm{\Phi }_c^{(ij,k)}(\nu ,\mu )=\left[\phi _{\mathrm{}}(𝒓_{ij},\nu )\phi _L(𝒓_k,\mu )\right]_J,`$ (44)
$`\phi _{\mathrm{}}(𝒓,\nu )=N_{\mathrm{}}(\nu )r^{\mathrm{}}\mathrm{exp}(\nu r^2)Y_{\mathrm{}}(\widehat{𝒓}),`$ (45)
where $`N_{\mathrm{}}`$ is the normalization factor, and $`𝒓_{ij}`$ ($`𝒓_k`$) denotes the relative coordinate between the i\- and j-th $`\alpha `$ particle (the k-th $`\alpha `$ particle and the center-of-mass coordinate of the i-th and j-th $`\alpha `$ particle). The angular momentum channel is presented as $`c=(\mathrm{},L)_J`$, where $`\mathrm{}`$ ($`L`$) denotes the relative orbital angular momentum between 2$`\alpha `$ clusters (the center-of-mass for the $`2\alpha `$ clusters and the third $`\alpha `$). The Gaussian parameter $`\nu `$ is taken to be of geometrical progression,
$$\nu _n=1/b_n^2,b_n=b_{\mathrm{min}}a^{n1},n=1n_{\mathrm{max}}.$$
(46)
It is noted that the prescription is found to be very useful in optimizing the ranges with a small number of free parameters with high accuracy Kamimura88 .
The total Hamiltonian for the $`3\alpha `$ system is presented as
$`={\displaystyle \underset{i=1}{\overset{3}{}}}T_iT_{cm}+{\displaystyle \underset{i<j=1}{\overset{3}{}}}\left[V_{2\alpha }(r_{ij})+V_{2\alpha }^{Coul}(r_{ij})\right]+V_{3\alpha }(r_{12},r_{23},r_{31})+V_{\mathrm{Pauli}},`$ (47)
where $`T_i`$, $`V_{2\alpha }`$ and $`V_{3\alpha }`$ stand for the kinetic energy operator for the i-th $`\alpha `$ particle, phenomenological $`2\alpha `$ and $`3\alpha `$ potentials, respectively, and $`V_{2\alpha }^{Coul}`$ is the Coulomb potential between 2$`\alpha `$ particles. The center-of-mass kinetic energy is subtracted from the Hamiltonian. The Pauli-blocking operator $`V_{\mathrm{Pauli}}`$ Kukulin84 is represented as
$`V_{\mathrm{Pauli}}=\underset{\lambda \mathrm{}}{lim}\lambda \widehat{O}_{\mathrm{Pauli}},`$ (48)
$`\widehat{O}_{\mathrm{Pauli}}={\displaystyle \underset{2n+\mathrm{}<4,\mathrm{}=even}{}}{\displaystyle \underset{(ij)=(12),(23),(31)}{}}|u_n\mathrm{}(𝒓_{ij}u_n\mathrm{}(𝒓_{ij})|,`$ (49)
which removes the Pauli forbidden states, $`u_{00}`$, $`u_{10}`$ and $`u_{20}`$, between any two $`\alpha `$ particles from the 3$`\alpha `$ model space. The Gaussian size parameter of the nucleon $`(0s)`$ orbit in the $`\alpha `$ cluster is taken to be $`b_N=1.358`$ fm, which reproduces the size of the $`\alpha `$ particle in free space. In the present study, we take the harmonic oscillator wave functions as the Pauli forbidden states. The eigenenergy $`E`$ of <sup>12</sup>C and coefficients $`A_c`$ in Eq. (42) are obtained in terms of the variational principle,
$`\delta \left[\mathrm{\Phi }_JE\mathrm{\Phi }_J\right]=0.`$ (50)
We use an effective $`2\alpha `$ potential which reproduces the observed $`\alpha \alpha `$ scattering phase shifts ($`S`$-, $`D`$\- and $`G`$-waves) and the resonant ground-state energy within the $`2\alpha `$ OCM framework. The effective $`2\alpha `$ potential and Coulomb potential, $`V_{2\alpha }`$ and $`V_{2\alpha }^{Coul}`$, are constructed with the folding procedure, where we fold the modified Hasegawa-Nagata effective $`NN`$ interaction (MHN) and the $`pp`$ Coulomb potential with the $`\alpha `$-cluster density. Also the strength of the second-range triplet-odd part in MHN is modified so as to reproduce the $`2\alpha `$ scattering phase shifts.
Only using the effective $`2\alpha `$ potential leads to a significant overbinding energy for the ground state of <sup>12</sup>C within the frame of the $`3\alpha `$ OCM. Thus, we introduce an effective, repulsive, $`3\alpha `$ potential, $`V_{3\alpha }`$, in addition to the $`2\alpha `$ potential,
$`V_{3\alpha }=V_0\mathrm{exp}\left[\beta \left(𝒓_{12}^2+𝒓_{23}^2+𝒓_{31}^2\right)\right],`$ (51)
where $`𝒓_{ij}`$ denotes the relative coordinate between the i\- and j-th $`\alpha `$ particles, and $`V_0`$ and $`\beta `$ are taken to be $`V_0=87.5`$ MeV and $`\beta =0.15`$ fm<sup>-2</sup>. Including the $`3\alpha `$ potential, the energy of the ground state of <sup>12</sup>C is reproduced with respect to the $`3\alpha `$ threshold, together with the nuclear radius (see Sec. III).
Single-$`\alpha `$ orbits and corresponding occupation probabilities for $`0^+`$, $`2^+`$, $`1^{}`$, and $`3^{}`$ states of <sup>12</sup>C are investigated by solving the eigenvalue equation of the single $`\alpha `$-particle density matrix in Eqs. (37) and (38) \[see Sec. II(b)\]. They will lead to a deep understanding about the structure of <sup>12</sup>C.
In the present investigation, we make a further structure study for the $`0^+`$ states of <sup>12</sup>C, because they have very intriguing features. According to Ref. Tohsaki01 , the $`0_1^+`$ state has a compact shell-model-like state, while the $`0_2^+`$ one is conjectured to have a dilute $`3\alpha `$ condensate structure, the nuclear radius of which is 4.3 fm, much larger than that of the ground $`0_1^+`$ state (2.48 fm). Thus, it is interesting to see in detail the structure change of the $`0^+`$ state of <sup>12</sup>C by taking the nuclear radius as a parameter. We investigate the dependence of the occupation probabilities and radial behaviors of the single $`\alpha `$-particle orbits in the $`0^+`$ state on its rms radius within the $`3\alpha `$ OCM framework. The results will give us helpful understanding about the structure of <sup>12</sup>C. The procedure of evaluating them is formulated hereafter.
First, we consider a Pauli-principle respecting $`3\alpha `$ OCM basis wave function. For the purpose, the eigenvalue problem for the Pauli operator in Eq. (49) is solved to obtain the Pauli forbidden state in the $`3\alpha `$ OCM model space
$`\widehat{O}_{\mathrm{Pauli}}|G_P^{3\alpha }=\lambda _P|G_P^{3\alpha },`$ (52)
where $`\lambda _P`$ denotes the eigenvalue for the eigenfunction $`|G_P^{3\alpha }`$. The Pauli operator, then, is expressed as
$`\widehat{O}_{\mathrm{Pauli}}={\displaystyle \underset{P}{}}|G_P^{3\alpha }\lambda _PG_P^{3\alpha }|.`$ (53)
If $`\lambda _P`$ is non-zero, its eigenfunction corresponds to the Pauli forbidden state. In the present study, the eigenvalue problem is solved with use of the $`3\alpha `$ OCM basis in Eq. (43). Then, the Pauli-principle respecting OCM basis wave function is given by
$`\stackrel{~}{\mathrm{\Phi }}_c^{3\alpha }(\nu ,\mu )=N_c(\nu ,\mu )\left[\mathrm{\Phi }_c^{3\alpha }(\nu ,\mu ){\displaystyle \underset{\lambda _P0}{}}|G_P^{3\alpha }G_P^{3\alpha }|\mathrm{\Phi }_c^{3\alpha }(\nu ,\mu )\right],`$ (54)
where $`N_c`$ denotes the normalization factor with the angular momentum channel $`c=(\mathrm{},L)_J`$, and $`\mathrm{\Phi }_c^{3\alpha }(\nu ,\mu )`$ is given in Eq. (43).
According to the results of the $`3\alpha `$ OCM calculation (see Sec. IIIA), the ground state ($`0_1^+`$) and second $`0_2^+`$ states of <sup>12</sup>C have the equilateral triangle configuration of the $`3\alpha `$ clusters. In addition, it is found that the single-angular-momentum-channel calculation with $`c=(\mathrm{}L)_J=(00)_0`$ gives a good approximation to the results of the full coupled-channel calculation. Thus, only the single angular momentum channel $`c=(\mathrm{}L)_J=(00)_0`$ is taken in the present calculation, and the equilateral triangle configuration is assumed for the basis wave function. The latter can be realized easily by putting the condition $`\nu =\mu `$ in Eq. (54). Then, the rms radius of the Pauli-principle respecting $`3\alpha `$ OCM basis wave function is evaluated as
$`\sqrt{r^2_\nu }=\left[\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu )|{\displaystyle \underset{i=1}{\overset{3}{}}}𝒓_{i}^{(cm)}{}_{}{}^{2}|\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu )+1.71^2\right]^{1/2},`$ (55)
which depends on $`\nu `$, and where we take into account the finite size of the $`\alpha `$ cluster. The energy of the $`3\alpha `$ system is given by
$`E_{3\alpha }(\nu )=\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu )|\stackrel{~}{}|\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu ).`$ (56)
where $`\stackrel{~}{}`$ denotes the total Hamiltonian of the $`3\alpha `$ system in which we subtract the Pauli-blocking operator $`V_{\mathrm{Pauli}}`$ from the $`3\alpha `$ OCM Hamiltonian $``$ in Eq. (47). The single-$`\alpha `$ density matrix is given as
$`\rho _\nu (𝒓,𝒓^{})=\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu ){\displaystyle \underset{i=1}{\overset{3}{}}}\delta (𝒓_𝒊^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓^{})\delta (𝒓_𝒊^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓)\stackrel{~}{\mathrm{\Phi }}_{(00)_0}^{3\alpha }(\nu ,\mu =\nu ).`$ (57)
The single $`\alpha `$-particle orbit and its occupation number in the basis wave function are obtained by solving the eigenvalue equation of the single-$`\alpha `$ density matrix,
$`{\displaystyle 𝑑𝒓^{}\rho _\nu (𝒓,𝒓^{})\phi _\eta (𝒓^{})}=\eta \phi _\eta (𝒓),`$ (58)
where the eigenfunction $`\phi _\eta `$ denotes the $`\alpha `$-particle orbit with the occupation number $`\eta `$ (eigenvalue). Thus, we can study the dependence of the occupancy of the single $`\alpha `$-particle orbits in the $`0^+`$ state of <sup>12</sup>C on its nuclear radius by choosing the parameter value $`\nu `$ so as to reproduce a given nuclear radius.
## III Results and discussion
### III.1 $`0_1^+`$ and $`0_2^+`$ states
Table 1 presents the results for the energy, measured from the $`3\alpha `$ threshold, and nuclear radii for the ground ($`0_1^+`$) and excited states ($`0_2^+`$) of <sup>12</sup>C. The energy for the ground state is reproduced well, and the corresponding nuclear radius, 2.44 fm, is in good agreement with the experimental charge radius ($`2.4829\pm 0.019`$ fm) with an error of about $`2\%`$. On the other hand, the rms distance between $`2\alpha `$ clusters in the $`0_1^+`$ state is $`\sqrt{r^2_{\alpha \alpha }}`$=3.02 fm (see Table 1), which is larger than that between the center-of-mass of the $`2\alpha `$ clusters and the third $`\alpha `$ cluster, $`\sqrt{r^2_{\alpha 2\alpha }}`$=2.61 fm. Then, the square of the ratio, $`\left[\sqrt{r^2_{\alpha 2\alpha }}/\sqrt{r^2_{\alpha \alpha }}\right]^2`$, is about 3/4. The results mean that the ground state has an equilateral-triangle-like intrinsic shape. Figure 2 shows the density distribution of the $`\alpha `$ particle for the $`0_1^+`$ state of <sup>12</sup>C. We see a prominent peak at $`r2`$ fm, which demonstrates clearly the shell-model-like compact structure of the ground state of <sup>12</sup>C.
As for the $`0_2^+`$ state, the energy measured from the $`3\alpha `$ threshold is $`E_{3\alpha }`$=0.86 MeV ($`E_x`$=8.13 MeV), which agrees well with the experimental data $`E_{3\alpha }^{exp}`$=0.38 MeV ($`E_x^{exp}`$=7.65 MeV). The calculated nuclear radius for the $`0_2^+`$ state is as large as 4.31 fm (see Table 1). This means that the state has a dilute $`3\alpha `$ structure, although our nuclear radius is a little larger than that in Ref. Tohsaki01 . The density distribution of the $`\alpha `$ particle for the $`0_2^+`$ state is illustrated in Fig. 2. In comparison with that for the ground state, we can easily recognize the dilute structure of the $`0_2^+`$ state, which is in contrast with the compact structure of the ground state.
The difference between the structures of the $`0_1^+`$ and $`0_2^+`$ states can be also seen in the radial behavior of the correlation functions, $`f_{\alpha \alpha }`$ and $`f_{\alpha 2\alpha }`$, with respect to the $`\alpha `$-$`\alpha `$ and $`\alpha `$-$`2\alpha `$ relative coordinates, respectively, of Eqs. (30) and (31). They are illustrated in Fig. 3. Reflecting the compact structure of the $`0_1^+`$ state, both of $`f_{\alpha \alpha }`$ and $`f_{\alpha 2\alpha }`$ have prominent peaks at $`r2.6`$ fm and $`2.5`$ fm, respectively, and extend to $`r5`$ fm, while those for the $`0_2^+`$ state show bump structures with peaks at $`r`$4 fm and $`r`$5 fm, respectively, and have a long tail up to $`r15`$ fm.
It is instructive to study the single $`\alpha `$-particle orbits (eigenfunctions) and occupation numbers (eigenvalues) of the one-body density matrix in Eq. (37). The results of the diagonalization of the latter are shown in Table 2 together with the occupation probability defined as the occupation number divided by the number of $`\alpha `$ particles. As for the ground state, the occupation probabilities spread out over $`S`$, $`D`$ and $`G`$ waves, but they are concentrated to the first orbits, $`S_1`$, $`D_1`$ and $`G_1`$ orbits, respectively, where $`L_k`$ denotes the $`k`$-th orbit for the $`L`$ wave. The occupation probabilities are about $`30\%`$ for all orbits. This result is expected from the fact that the ground-state wave function is of nuclear SU(3)-like character, SU(3)$`[f](\lambda \mu )_{J^\pi }=[444](04)_{0^+}`$ with quanta $`Q`$=8, where the SU(3) bases with $`Q<8`$ correspond to the Pauli-forbidden states. Since the SU(3)$`[444](04)_{0^+}`$ state is the eigenfunction of the $`3\alpha `$ RGM norm kernel in Eq. (18), it can be regarded as the $`3\alpha `$ boson wave function with $`Q=8`$, see Eq. (19). The state is described as
$`|[444](04)_{0^+}`$ $`=`$ $`{\displaystyle \underset{n\mathrm{}NL}{}}a_{nlNL}|(n\mathrm{})(NL),`$ (59)
$`=`$ $`\sqrt{{\displaystyle \frac{64}{225}}}|2s2S\sqrt{{\displaystyle \frac{80}{225}}}|1d1D)+\sqrt{{\displaystyle \frac{81}{225}}}|0g0G,`$
where $`|(n\mathrm{})(NL)`$ presents the basis function $`|u_n\mathrm{}(𝒓_{2\alpha })u_{NL}(𝒓_{\alpha 2\alpha })`$ with $`2n+\mathrm{}+2N+L=8`$, and $`u_n\mathrm{}`$ ($`u_{NL}`$) denotes the harmonic oscillator wave function with the number of nodes $`n`$ ($`N`$) and orbital angular momentum $`\mathrm{}`$ ($`L`$) referring to the coordinate vector $`𝒓_{2\alpha }`$ ($`𝒓_{\alpha 2\alpha }`$) between $`2\alpha `$ clusters (between the center-of-mass for the $`2\alpha `$ clusters and the third $`\alpha `$ cluster). Let us define $`L_\alpha `$ as the orbital angular momentum of a single-$`\alpha `$ orbit. Then, $`L`$ in Eq. (59) corresponds to $`L_\alpha `$, because $`L_\alpha `$ is defined as the orbital angular momentum with respect to $`𝒓_\alpha ^{(cm)}`$, coordinate vector of the $`\alpha `$ particle measured from the center-of-mass coordinate of <sup>12</sup>C (see Eq. (27)), which is parallel to $`𝒓_{\alpha 2\alpha }`$ ($`𝒓_\alpha ^{(cm)}=\frac{2}{3}𝒓_{\alpha 2\alpha }`$). From the definition of the one-body density matrix in Eq. (37), the single-$`\alpha `$ orbits and occupation probabilities for the SU(3) state in Eq. (59) are given as follows: $`64/25528\%`$ for $`S`$-orbit, $`80/22536\%`$ for $`D`$-orbit, and $`81/22536\%`$ for $`G`$-orbit. Thus, we can understand the reason why the $`S_1`$, $`D_1`$ and $`G_1`$ orbits in Table 2 have about $`30\%`$ occupation probabilities each.
Figure 4(a) demonstrates the radial parts for the $`S_1`$-, $`D_1`$\- and $`G_1`$-orbits, the number of nodes of which are two, one and zero, respectively. Reflecting the SU(3) character, the radial behaviors of the three orbits are similar to those of the harmonic oscillator wave functions ($`u_{NL}`$) with $`Q=4`$, $`u_{02}`$, $`u_{21}`$ and $`u_{40}`$, respectively, where $`N`$ ($`L`$) denotes the number of nodes (orbital angular momentum). We see that the radial parts of the single $`\alpha `$-particle orbits oscillate widely in the inside region ($`r<4`$ fm). This is due to the strong Pauli blocking effect for the ground state with the compact shell-model-like structure. The large oscillation can also be seen in the reduced width amplitude of the $`\alpha `$+<sup>8</sup>Be($`0^+`$) channel for the ground state shown in Fig. 5.
Concerning the $`0_2^+`$ state, the occupation probabilities are shown in Table 2. We see a strong concentration on a single orbit: the occupation probability of the $`S_1`$ orbit is largest, amounting to about $`70\%`$, and those for other orbits are very small. This means that each of the three $`\alpha `$ particles in the $`0_2^+`$ state is in the $`S_1`$ orbit with occupation probability as large as about $`70\%`$. The radial behavior of the $`S_1`$ orbit is illustrated with the solid line in Fig. 4(b). We see no nodal behavior but small oscillations in the inner region ($`r<4`$ fm) and a long tail up to $`r`$10 fm. For reference, the radial behavior of the $`S`$-wave Gaussian function, $`\phi _{0s}(r)=N_{0s}(a)\mathrm{exp}(ar^2)`$, is drawn with the dashed line in Fig. 4(b), where the size parameter $`a`$ is chosen to be 0.038 fm<sup>-2</sup>, and $`N_{0s}(a)`$ denotes the normalization factor. The radial behavior of the $`S_1`$ orbit is similar to that of the $`S`$-wave Gaussian function, in particular, in the outer region ($`r>4`$ fm), whereas a slight oscillation of the former around the latter can be seen in the inner region ($`r<4`$ fm).
The small oscillation of the $`S_1`$ orbit in the inner region can also be seen in the reduced width amplitude of the $`0_2^+`$ state for the $`\alpha `$-<sup>8</sup>Be($`0^+`$) channel in Fig. 6(a). In order to study the origin of the small oscillation, we show in Fig. 6(b) the results of the reduced width amplitudes of the $`0_2^+`$ state for the $`\alpha `$-<sup>8</sup>Be($`0^+`$) channel, fixing the distance between the $`2\alpha `$ clusters in <sup>8</sup>Be to $`r_{\alpha \alpha }`$=0.5, 2.5, 4.5 and 6.5 fm. In the case of $`r_{\alpha \alpha }<4`$ fm, we see the nodal behavior with the large oscillation in the inner region, coming from the strong Pauli blocking effect among the $`3\alpha `$ clusters, while the nodal behavior is disappearing and the oscillations are getting weaker for the larger $`r_{\alpha \alpha }`$ ($`4`$ fm), reflecting the weaker Pauli blocking effect. Thus, the small oscillations in the radial behavior of the $`S_1`$ orbit is evidence for the weak Pauli blocking effect for the $`0_2^+`$ state with the dilute structure.
The momentum distributions of the $`\alpha `$ particles, $`\rho (k)`$ and $`k^2\times \rho (k)`$, are shown for the $`0_1^+`$ and $`0_2^+`$ states in Fig. 7. Reflecting the dilute structure of the $`0_2^+`$ state, we see a strong concentration of the momentum distribution in the $`k<1`$ fm<sup>-1</sup> region, and the behavior of $`\rho (k)`$ is of the $`\delta `$-function type, similar to the momentum distribution of the dilute neutral atomic condensate states at very low temperature trapped by the external magnetic field Dalfovo99 . On the other hand, the ground state has higher momentum component up to $`k6`$ fm<sup>-1</sup> as seen from the behavior of $`k^2\times \rho (k)`$ reflecting the compact structure. The above results for the radial behavior of the $`S_1`$ orbit, occupation probability and momentum distribution for the $`0_2^+`$ state leads us to conclude that this state is similar to an ideal dilute $`3\alpha `$ condensate with as much as about $`70\%`$ occupation probability.
Let us make some remarks on the calculated energy ($`E_{3\alpha }`$=0.85 MeV) and wave function of the $`0_2^+`$ state. They were evaluated under the bound state approximation in the present study (see Sec. II). The quite small experimental width for $`0_2^+`$ ($`\mathrm{\Gamma }`$=8.5 eV) Ajzenberg90 means that the bound state approximation is very good to describe the wave function. The complex scaling method Kuruppa88 is powerful to search for resonant states and calculate the exact energies and widths, and is applied easily to the $`3\alpha `$ system by slightly modifying the present framework. The detailed framework is skipped here and referred to Ref. Kuruppa88 . In the present study, we investigated the energy of the $`0_2^+`$ state with the complex scaling method. It was found that a resonant state, corresponding to the $`0_2^+`$ state, appears at $`E_{3\alpha }`$=0.85 MeV with a width less than the numerical uncertainty ($``$100 keV in the present calculation). The results confirm that the bound state approximation is good to describe the $`0_2^+`$ resonant state.
It is interesting to compare our results with those given by Matsuura et al. Matsuura04 , who used the normalized spectroscopic amplitude to obtain the bosonic quantities such as the single-$`\alpha `$ orbits and occupation probabilities for the $`0_2^+`$ state in place of the $`3\alpha `$ boson wave function. According to their results, the occupation probability of the $`S_1`$ orbit ($`0S`$ orbit in Ref. Matsuura04 ) for the $`0_2^+`$ state is about $`70\%`$, the value of which is quite similar to ours in Table 2. However, the radial behavior of the $`S_1`$ orbit for the $`0_2^+`$ state as well as the one of the $`0_1^+`$ state given by Matsuura et al. are quite different from our results, and seems unnatural. For example, the $`S_1`$ orbit for the $`0_2^+`$ state has as much as $`68`$ nodes and shows a behavior similar to that for the $`0_1^+`$ state, in spite of the fact that the $`0_2^+`$ state has a dilute $`3\alpha `$ condensate structure (see Fig. 6 in Ref. Matsuura04 ). In addition, the $`G`$-orbit for $`0_1^+`$ state has a prominent peak at $`r13`$ fm, although the state has a shell-model-like compact structure. Also the radial behavior of the single-$`\alpha `$ orbits given in Ref. Matsuura04 is hard to understand. This may be due to the fact that those authors used the normalized spectroscopic amplitude in place of the $`3\alpha `$ boson wave function.
### III.2 $`2_1^+`$ and $`2_2^+`$ states
The $`2_1^+`$ state at $`E_{3\alpha }^{exp}=2.83`$ MeV ($`E_x`$=4.44 MeV) belongs to the rotational band of the ground state starting from the $`0_1^+`$ state at $`E_{3\alpha }^{exp}=7.27`$ MeV. The calculated energy and nuclear radius for $`2_1^+`$ in the present study are shown in Table 1$`E_{3\alpha }=5.28`$ MeV and 2.45 fm, respectively. The nuclear radius is almost the same as the one for the ground state, although the calculated excitation energy is underestimated in comparison with the experimental one, in line with what is discussed in other papers with the microscopic or semi-microscopic $`3\alpha `$ cluster model Fujiwara80 ; Saito68 ; Horiuchi74 ; Kato89 ; Kurokawa04 .
The occupation probabilities of the single-$`\alpha `$ orbits for $`2_1^+`$ are demonstrated in Table 2. The occupation numbers are concentrated to the first $`D_1`$ and $`G_1`$ orbits with about $`50\%`$. Comparing with those for the $`0_1^+`$ state, we notice the smallness of the occupation number for the $`S_1`$ orbit. This feature can be understood from the fact that the $`2_1^+`$ state is of the SU(3)$`[f](\lambda \mu )_J=[444](04)_{2^+}`$ type with $`Q`$=8. The SU(3) state is described as
$`|[444](04)_{2^+}`$ $`=`$ $`\sqrt{0.07111}|1d2S,`$ (60)
$`+`$ $`\sqrt{0.07111}|2s1D\sqrt{0.43900}|1d1D\sqrt{0.00735}|0g1D,`$
$``$ $`\sqrt{0.00735}|1d0G+\sqrt{0.40408}|0g0G,`$
where $`|(n\mathrm{})(NL)`$ presents the basis function, $`|u_n\mathrm{}(𝒓_{2\alpha })u_{NL}(𝒓_{\alpha 2\alpha })`$ (with $`2n+\mathrm{}+2N+L=8`$), with the harmonic oscillator wave function $`u_n\mathrm{}`$. From the definition of the one-body density matrix in Eq. (37), the occupation probabilities for the SU(3) state in Eq. (60) are given as 0.07111 for $`S`$ orbit, 0.07111+0.43900+0.00735=0.51746 for $`D`$ orbit, and 0.00735+0.40408=0.41143 for $`G`$ orbit. Reflecting the character of the SU(3) structure, the occupation probability for the $`S_1`$ orbit in Table 2 is as small as $`8.5\%`$. The radial behavior of the single-$`\alpha `$ orbits, $`S_1`$, $`D_1`$ and $`G_1`$ ones, is shown in Fig. 8(a). They are similar to those for the $`0_1^+`$ state shown in Fig. 4(a).
The structure study of <sup>12</sup>C based on the $`3\alpha `$-condensate type wave function Funaki04 indicated that the $`2_2^+`$ state at $`E_{3\alpha }^{exp}=2.6\pm 0.3`$ MeV with the width of $`\mathrm{\Gamma }=1.0\pm 0.3`$ MeV Itoh04 has a structure similar to the $`0_2^+`$ state at $`E_{3\alpha }=0.38`$ MeV with the dilute $`3\alpha `$ condensation Funaki04 . The conclusion stems from the result that the $`2_2^+`$ state has a large overlap with the single condensate wave function of a $`3\alpha `$ gas-like structure, the squared value of which amounts to about $`88\%`$. Thus, it is interesting to study the structure of the $`2_2^+`$ state in the present framework. Since the $`2_2^+`$ state is a resonant state with non negligible width, a continuum treatment is requested to estimate exactly the resonant energy and width.
In order to study the resonant properties of the $`2_2^+`$ state, we take the complex-scaling method Kuruppa88 , which can be applied easily to the present $`3\alpha `$ system by slightly modifying the framework given in Sec. II. The method is powerful to evaluate not only the resonant energy and width but also the nuclear radius. The details are again skipped here and we refer to Ref. Kuruppa88 . The calculated results are as follows: 1) the $`2_2^+`$ resonant state is located at $`E_{3\alpha }`$=2.3 MeV with $`\mathrm{\Gamma }`$=1.0 MeV, results which are in good agreement with the experimental data Itoh04 , and 2) the calculated nuclear radius is 4.3 fm, almost the same as that of the $`0_2^+`$ state. Thus, the $`2_2^+`$ state has a dilute $`3\alpha `$ structure.
It is interesting to study the single-$`\alpha `$ orbits and occupation probabilities in the $`2_2^+`$ state. For this purpose, we need to have the wave function of the $`2_2^+`$ state. Since the calculated width is not so large in comparison with the resonance energy, the bound state approximation is rather good to describe the resonant wave function. The bound state approximation of the wave function is obtained within the framework of Sec. II, although the wave function gives a large nuclear radius, about 6 fm (see Table 1). Table 2 illustrates the occupation probabilities of the single-$`\alpha `$ orbits ($`S`$-, $`D`$\- and $`G`$-waves) for the $`2_2^+`$ state. We see that the occupation probability concentrates on only one orbit, the $`D_1`$ orbit, with occupancy as large as $`83\%`$, and the radial behavior of the orbit is likely to be of the $`D`$-wave Gaussian-function-type with a long tail \[see Fig. 8(b)\], reflecting a dilute structure. These characteristics are quite similar to those for the $`0_2^+`$ state. Thus, we conclude that the $`2_2^+`$ state belongs to the $`3\alpha `$-condensate structure.
According to the results in Ref. Funaki04 , it was found that the $`2_2^+`$ state has dominant $`S`$-wave between $`2\alpha `$ particles and a $`D`$-wave between the center-of-mass of the $`2\alpha `$ particles and the third $`\alpha `$,
$`\mathrm{\Phi }(2_2^+)|u_{\mathrm{}=0}(𝒓_{2\alpha })U_{L=2}(𝒓_{\alpha 2\alpha }).`$ (61)
This interpretation is consistent with the preset result. The reason is as the follows. From the definition of the single-$`\alpha `$ density matrix in Eq. (37), the single-$`\alpha `$ density of the $`2_2^+`$ state is presented as
$`\rho (𝒓,𝒓^{})`$ $`=`$ $`3\times \mathrm{\Phi }(2_2^+)\delta (𝒓_\mathrm{𝟏}^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓^{})\delta (𝒓_\mathrm{𝟏}^{\mathbf{(}𝒄𝒎\mathbf{)}}𝒓)\mathrm{\Phi }(2_2^+),`$ (62)
$``$ $`3\times N_{2\alpha }\times U_{L=2}(𝒓)U_{L=2}^{}{}_{}{}^{}(𝒓^{}),`$ (63)
where $`N_{2\alpha }=𝑑𝒓_{2\alpha }u_{0}^{}{}_{}{}^{}(𝒓_{2\alpha })u_0(𝒓_{2\alpha })1`$. Thus, the $`2_2^+`$-state wave function, Eq. (61), has a dominant occupation probability of the $`D`$-orbit, $`U_{L=2}`$. The results are in good agreement with the present study.
### III.3 $`3_1^{}`$ state
The $`3^{}`$ state at $`E_{3\alpha }^{exp}=`$2.37 MeV is an interesting one from the point of view of the dilute $`\alpha `$ condensation. If the state is a condensate with all of the $`3\alpha `$ particles in the $`P`$ orbit, there is the possibility of a superfuid with vortex lines, similar to the rotating dilute atomic condensate at very low temperature Dalfovo99 . Thus, it is an intriguing problem to study the structure in the present framework.
The calculated energy of the $`3^{}`$ state is in good agreement with the experimental data (see Table 1). The very small width ($`\mathrm{\Gamma }^{exp}=3.4`$ keV) Ajzenberg90 indicates that the bound state approximation is very good to describe the state. In fact, we checked it theoretically with the complex-scaling method, and found that the calculated resonant energy (width) is almost the same as the one with the bound state approximation (less than 100 keV, which is the numerical uncertainty in the present calculation). Thus, we use the $`3^{}`$ wave function under the bound state approximation to study the characteristics of the state.
The calculated nuclear radius for the $`3^{}`$ state is 2.95 fm, the value of which is larger than that for the ground state ($`0_1^+`$), while it is smaller than that for the $`0_2^+`$ state (see Table 1). This suggests that the structure of the $`3^{}`$ state is intermediate between the shell-model-like compact structure ($`0_1^+`$) and the dilute $`3\alpha `$ structure ($`0_2^+`$). The occupation probability of the single-$`\alpha `$ orbits for the state are shown in Table 3$`44.7\%`$ for $`P_1`$-orbit and $`27.9\%`$ for $`F_1`$-orbit. Although the concentration of the single orbit $`P_1`$ amounts to about $`50\%`$, the radial behavior of the single-$`\alpha `$ orbit in Fig. 9 has two nodes in the inner region. However, the amplitude of the inner oscillations is significantly smaller than that for the ground state in Fig. 3(a). The small oscillations indicate the weak Pauli-blocking effect, and thus, we can see the precursor of the $`3\alpha `$ condensate state, although the $`3^{}`$ state is not an ideal rotating dilute $`3\alpha `$ condensate.
### III.4 $`1_1^{}`$ state
The experimental width of the $`1_1^{}`$ state at $`E_{3\alpha }^{exp}`$=3.57 MeV is as small as $`\mathrm{\Gamma }`$=315 keV. Ajzenberg90 This means that the bound state approximation is good to describe the state. In fact, the calculated energy of the $`1_1^{}`$ state under the bound state approximation is $`E_{3\alpha }`$=3.11 MeV, which is quite similar to that with the complex scaling method ($`E_{3\alpha }`$=3.1 MeV and $`\mathrm{\Gamma }`$=0.1 MeV) and in good agreement with the experimental value.
The calculated nuclear radius, 3.32 fm, is larger than that of the ground state (2.44 fm) and the $`3_1^{}`$ state (2.95 fm) but is smaller than that of the $`0_2^+`$ one (4.3 fm). The occupation probabilities of the $`\alpha `$ particles in the $`1_1^{}`$ state are shown in Table 3$`35\%`$ for $`P_1`$ orbit and $`16\%`$ for $`F_1`$ orbit. Thus, there is no concentration of the occupation probability to a single orbit like the $`0_2^+`$ and $`2_2^+`$ states. Since the $`\alpha `$ particles in the $`1_1^{}`$ state are distributed over in several orbits, the state is not of the dilute $`\alpha `$-condensate type. Figure 10 shows the radial behavior of the $`P_1`$ and $`F_1`$ orbits in the $`1_1^{}`$ state. The $`P_1`$ orbit has two nodes in the inner region, the behavior of which is rather similar to the $`2P`$ harmonic oscillator wave function. However, the $`F_1`$ orbit has a $`F`$-wave Gaussian-type behavior. (Exactly speaking, the orbit has one node at the vicinity of the origin, which can not be seen in Fig. 10.) Also we see the oscillatory behavior of the $`F_1`$ orbit for $`0<r<2`$ fm, similar to the one of the $`S_1`$ orbit in the $`0_2^+`$ state in Fig. 4(b). These interesting behaviors of the $`F_1`$ orbit indicate some signal of the dilute $`\alpha `$ condensation, reflecting the relatively large nuclear radius (3.32 fm) for the $`1_1^{}`$ state.
### III.5 Structure change of the $`0^+`$ state with nuclear radius
In Sec. IIIA, we found that the $`0_2^+`$ state has a dilute $`3\alpha `$ structure characterized by the nuclear radius as large as about 4.3 fm, in which the $`\alpha `$ particle occupies in the single orbit ($`S_1`$) with about $`70\%`$ probability, and the radial behavior of the $`S_1`$ orbit is similar to the $`S`$-wave Gaussian wave function with a very long tail. On the other hand, the $`0_1^+`$ state has a compact structure with a nuclear radius of 2.44 fm, where the occupation probabilities of the $`\alpha `$ particles spread out over the $`S`$, $`D`$ and $`G`$ orbits, amounting to about $`30\%`$, each. The feature is much in contrast with that of the $`0_2^+`$ state. The nuclear radius or density of <sup>12</sup>C seems to have a close relation with making the compact structure and the dilute $`3\alpha `$ structure in the <sup>12</sup>C $`0^+`$ state. Thus, it is very interesting to see the structure change of the $`0^+`$ state of <sup>12</sup>C by taking the nuclear radius (or density) as parameter. The dependence of the occupation probabilities and radial behaviors of the single $`\alpha `$-particle orbits in the $`0^+`$ state on its nuclear radius is investigated with the use of the simple framework given in the latter part of Sec. IIC (see Eqs. (52)$``$(58)).
Figure 11 shows the dependence of the energy of <sup>12</sup>C measured from the $`3\alpha `$ threshold on the nuclear radius $`R_N`$, 2.20 fm $`R_N4.86`$ fm, corresponding to a nuclear density 0.15 $`\rho /\rho _0`$ 1.6 ($`\rho _0`$ denotes the normal density). The energy minimum point appears around $`R_N`$2.4 fm, corresponding to the normal density region. We see the strong repulsion in the region of $`R_N<2.2`$ fm, due to the kinetic-energy effect and Pauli-blocking effect, while the almost flat region appears at $`R_N>4`$ fm and the energy is positive and small, less than 1 MeV with respect to the $`3\alpha `$ threshold.
The occupation probabilities of the single-$`\alpha `$ orbits ($`S_1`$-, $`D_1`$-, and $`G_1`$-orbits) are shown in Fig. 12 with respect to the nuclear radius, where $`L_k`$ denotes the $`k`$-th orbit for the $`L`$ wave. In the region of $`R_N=2.22.4`$ fm (normal density region), the occupation probabilities of the $`\alpha `$ particles spread out over the $`S`$, $`D`$ and $`G`$ orbits, amounting to about $`30\%`$, each. This feature is almost the same as that of the $`0_1^+`$ state obtained by the $`3\alpha `$ OCM calculation, the nuclear radius of which is 2.43 fm (see Sec. IIIA). Figure 13 shows the radial behavior of the single-$`\alpha `$ orbit, $`S_1`$, with respect to the nuclear radius. The $`S_1`$-orbit at $`R_N2.42`$ fm (Fig. 13(a)) has two nodes and the radial behavior is of the $`2S`$ harmonic oscillator wave function (howf) type, the result of which is almost the same as that of the $`0_1^+`$ state obtained by the $`3\alpha `$ OCM calculation (see Fig. 4(a)). Thus, the wave function with $`R_N2.4`$ fm has the SU(3)$`[f](\lambda \nu )=[444](04)`$ character (see Eq. (59)).
Increasing the nuclear radius from $`R_N=2.42`$ fm, the occupation probability of the single-$`\alpha `$ orbits concentrates gradually on a single orbit ($`S_1`$), and it amounts to be about 90 % at $`R_N`$=4.84 fm ($`\rho /\rho _0`$=0.14) in the present calculation (see Fig. 12). The radial behaviors of the $`S_1`$ orbit with $`R_N`$=2.42, 2.70, 3.11 and 4.84 fm are demonstrated in Figs. 13(a), (b), (c) and (d), respectively. We can see that increasing the nuclear radius, the internal oscillation observed in the $`S_1`$ orbit with $`R_N`$=2.42 fm is gradually disappearing and, finally, the $`2S`$-type radial behavior transits to the zero-node long-ranged $`S`$-wave type (Gaussian) with the occupation probability of about 90 %, approaching an ideal dilute $`\alpha `$ condensate. The reason of why only the $`S`$ wave survives in the case of increasing the nuclear radius is due to the fact that the centrifugal barrier is not at work for the $`S`$-wave $`\alpha `$ orbit. The $`S`$-wave $`\alpha `$ particles, thus, can move in a nucleus with a given large nuclear radius, although they are confined by the Coulomb potential barrier produced self-consistently Yamada04 . According to the results of the $`3\alpha `$ OCM calculation (see Sec. IIIA), the $`\alpha `$ particle in the $`0_2^+`$ state ($`R_N`$=4.3 fm) is occupied in the single orbit ($`S_1`$) with about 70 % probability, the radial behavior of which is similar to the $`S`$-wave Gaussian wave function with a very long tail. Their results are consistent with those in Figs. 12 and 13.
## IV Summary
In this work we have investigated the bosonic properties such as single-$`\alpha `$ particle orbits and occupation numbers in the $`J^\pi `$=$`0^+`$, $`2^+`$, $`1^{}`$, and $`3^{}`$ states of <sup>12</sup>C around the $`3\alpha `$ threshold within the framework of the $`3\alpha `$ OCM (orthogonality condition model). The $`3\alpha `$ OCM equation is based on the equation of motion for the $`N\alpha `$ bosons derived from the microscopic $`N\alpha `$ cluster model theory. The experimental energy spectra for $`0_1^+`$, $`0_2^+`$, $`2_2^+`$, $`1_1^{}`$, and $`3_1^{}`$ are reproduced well with the $`3\alpha `$ OCM.
The main results to be emphasized here are as follows.
(1) The $`0_2^+`$ state at $`E_{3\alpha }^{exp}`$=0.38 MeV has a dilute $`3\alpha `$ structure characterized by the nuclear radius as large as about 4.3 fm. The analysis of the single-$`\alpha `$ orbits and occupation probabilities for the dilute state shows that the $`\alpha `$ particle is occupied in a single orbit ($`S_1`$) with about $`70\%`$ probability, and the radial behavior of the $`S_1`$ orbit is similar to the $`S`$-wave Gaussian wave function with a very long tail. The momentum distribution of the $`\alpha `$ particle illustrates the $`\delta `$-function like behavior, similar to the momentum distribution of a dilute neutral atomic condensate states at very low temperature, a feature which eventually can be measured experimentally. These results give more theoretical evidence that the $`0_2^+`$ state is a dilute $`3\alpha `$ condensate. On the other hand, the $`0_1^+`$ state has a compact structure with a nuclear radius of 2.44 fm. The occupation probabilities of the $`\alpha `$ particles spread out over the $`S`$, $`D`$ and $`G`$ orbits, amounting to about $`30\%`$, each, the results of which comes from the fact that the $`0_1^+`$ state is characterized by the nuclear SU(3) wave function, $`[f](\lambda \mu )=[444](04)`$. The feature is much in contrast with that of the $`0_2^+`$ state.
(2) In order to understand further the characteristic structure of the two $`0^+`$ states, we have studied the single-$`\alpha `$ orbital behavior in the <sup>12</sup>C($`0^+`$) state by taking the nuclear radius $`R_N`$ (or density $`\rho /\rho _0`$) as parameter, 2.42 $`R_N`$ 4.84 fm, (0.15 $`\rho /\rho _0`$ 1.2), where $`\rho _0`$ denotes the normal density). We found that the single-$`\alpha `$ orbits in the <sup>12</sup>C($`0^+`$) state are smoothly changed with the nuclear radius $`R_N`$, and their behavior is classified into the following three types, depending on the value of $`R_N`$: i) at $`R_N2.4`$ fm ($`\rho /\rho _0`$ 1), we have two-nodal $`S`$-orbit ($`2S`$), one-nodal $`D`$-orbit ($`1D`$) and zero-nodal $`G`$-orbit ($`0G`$) with about $`30\%`$ occupation probability, each, characterized by a nuclear SU(3) wave function, ii) increasing the nuclear radius from $`R_N2.4`$ fm, the occupation probability of the single-$`\alpha `$ orbits concentrates gradually on a single $`S`$-orbit, in which the two-nodal behavior is disapperaing, and then, iii) at $`R_N4`$ fm ($`\rho /\rho _0`$ 0.2), there appears a dominant zero-nodal Gaussian ($`0S`$-type) orbit with a very long tail, the radial behavior of which is similar to that of the $`0_2^+`$ state in <sup>12</sup>C as mentioned above. The structure change is caused mainly by the Pauli-blocking effect, the strength of which depends dominantly on the nuclear radius $`R_N`$ in the present framework.
(3) The structure of the $`2_2^+`$ state at $`E_{3\alpha }^{exp}`$=$`2.6\pm 0.3`$ MeV with $`\mathrm{\Gamma }=1.0\pm 0.3`$ MeV was studied with the present $`3\alpha `$ OCM and the complex-scaling method. We found that the $`2_2^+`$ resonant state appears at $`E_{3\alpha }`$=2.3 MeV with $`\mathrm{\Gamma }=`$1.0 MeV, in agreement with the experimental data, and the calculated nuclear radius is 4.3 fm, similar to that of the $`0_2^+`$ state. The $`2_2^+`$ wave function obtained with the $`3\alpha `$ OCM was used to study the bosonic properties of the state. It was found that the occupation probability of the $`\alpha `$ particle concentrates only on the $`D_1`$ orbit, amounting to be as large as about $`80\%`$, and the radial behavior is of the $`D`$-wave Gaussian type with long tail. The characteristics of the boson properties in $`2_2^+`$ is quite similar to those in $`0_2^+`$ at $`E_{3\alpha }^{exp}`$=0.38 MeV. Thus, the $`2_2^+`$ state has the dilute $`3\alpha `$-condensate-like structure. On the other hand, the $`2_1^+`$ state has a compact structure with the nuclear radius, 2.44 fm, like the ground state. The occupation probabilities of the $`\alpha `$ particles spread out over the $`D`$ and $`G`$ orbits, amounting to about $`56\%`$ and $`33\%`$, respectively, reflecting the SU(3) character of the $`2_1^+`$ state.
(4) We investigated the $`\alpha `$ bosonic properties of the negative parity states, $`1_1^{}`$ at $`E_{3\alpha }^{exp}`$=3.57 MeV and $`3_1^{}`$ at $`E_{3\alpha }^{exp}`$=2.37 MeV. Their nuclear radii are 3.32 and 2.95 fm, respectively, which are larger than that of the ground state ($`0_1^+`$) but smaller than that of $`0_2^+`$. The calculated occupation probabilities of the $`\alpha `$ particles in those states show that there is no concentration on a single $`\alpha `$ orbit like in the $`0_2^+`$ and $`2_2^+`$ states. The results indicates that the $`1^{}`$ and $`3_1^{}`$ states are not of the dilute $`3\alpha `$ condensate. The radial behavior of the $`P`$\- and $`F`$-wave single $`\alpha `$ orbits, however, suggests that small components of the $`3\alpha `$ condensation exist even in the negative parity states, which is reflected by their relatively large nuclear radii.
## Acknowledgments
We acknowledge helpful discussions with H. Horiuchi, K. Ikeda, G. Röpke, Y. Suzuki, and A. Tohsaki.
|
warning/0506/hep-ph0506213.html
|
ar5iv
|
text
|
# Analytic evaluation of the amplitudes for orthopositronium decay to three photons to one-loop order
## I Introduction
Positronium, the electron-positron bound state, is well-suited for probing many fundamental aspects of particle physics. Karshenboim04 The physics of positronium is governed almost exclusively by the electromagnetic force—weak interaction effects are negligible compared to present experimental and theoretical uncertainties Bernreuther81 ; Alcorta94 ; Govaerts96 ; Czarnecki99 . As a consequence, positronium is an ideal system for testing QED through high precision comparison between experimental and theoretical results for energy levels and decay rates. The states of positronium are eigenstates of the charge conjugation and parity operators C and P, so positronium can be used to test the discrete symmetries C, P, and T and combinations thereof. Vetter04 Positronium has been the focus of many past and ongoing attempts to observe physics beyond the standard model. Gninenko02 ; Rubbia04 In this work we focus on the decay of spin-1 orthopositronium to three photons.
The orthopositronium decay rate has been the subject of continuing experimental and theoretical work since the first measurement by Deutsch in 1951. Deutsch51 A summary of all experimental and theoretical results has been given by Adkins, Fell, and Sapirstein Adkins02 and updated with commentary by Rubbia Rubbia04 and by Sillou Sillou04 . By 1990 it was apparent that there was an “orthopositronium lifetime puzzle”, as the most precise experimental determinations (gas Westbrook89 and vacuum Nico90 results from the Michigan group) were in disagreement with theory Caswell77 ; Caswell79 by several standard deviations. Many experiments were mounted to look for exotic decays of orthopositronium in an attempt to resolve the discrepancy. Rubbia04 ; Dobroliubov93 ; Dvoeglazov93 ; Skalsey97 ; Vetter04a Newer, somewhat less precise powder results from the Tokyo group in 1995 Asai95 and 2000 Jinnouchi00 were consistent with theory and inconsistent with the earlier Michigan results. In 2000 the calculation of all $`O(\alpha ^2)`$ corrections to the decay rate were completed. Adkins00 ; Adkins02 Including yet higher order logarithmic corrections as well, Hill00 ; Kniehl00 ; Melnikov00 the theoretical prediction is Adkins02
$$\mathrm{\Gamma }(\mathrm{theory})=7.039979(11)\mu s^1.$$
(1)
The $`O(\alpha ^2)`$ correction was found to be not unusually large, leaving the discrepancy with the Michigan results intact. Finally, in 2003 the lifetime puzzle was resolved by two new high-precision results from the Tokyo Jinnouchi03 and Michigan Vallery03 groups:
$`\mathrm{\Gamma }(\mathrm{Tokyo})`$ $`=`$ $`7.0396(12\mathrm{stat}.)(11\mathrm{syst}.)\mu s^1`$ (2a)
$`\mathrm{\Gamma }(\mathrm{Michigan})`$ $`=`$ $`7.0404(10\mathrm{stat}.)(8\mathrm{syst}.)\mu s^1,`$ (2b)
consistent with each other and with theory.
The resolution of the o-Ps lifetime puzzle does not decrease the long-term usefulness of positronium decay as a probe of fundamental physics. Ongoing and proposed experiments involving positronium decay include those of Refs. Vallery03 ; Crivelli04 ; Rubbia04 ; Gninenko04 ; Vetter04 . One challenge is to improve the experimental precision of the o-Ps decay rate (currently about 200 ppm) to a level closer to the present theoretical value (about 2 ppm). The $`O(\alpha ^2)`$ contribution to that rate is 250 ppm, so improved experimental precision will be required in order to test the $`O(\alpha ^2)`$ calculated result.
In this work we describe an analytic evaluation of the one-loop o-Ps$`3\gamma `$ decay amplitudes. We use these amplitudes to obtain a precise value for the $`O(\alpha )`$ decay rate contribution, and also to calculate the part of the $`O(\alpha ^2)`$ correction coming from the square of the one-loop amplitudes. These results have been reported already Adkins96 —here we give further details. We also supply an explicit analytic expression for the $`O(\alpha )`$ differential decay rate in terms of photon energy variables. From the differential decay rate it is easy to obtain the $`O(\alpha )`$ corrected one-photon energy spectrum. (This energy spectrum, calculated more laboriously by numerical methods, has been useful in developing simulations of experimental arrangements. Asai95 ; Jinnouchi03 )
We adapt the formalism of covariant decay amplitudes, originally developed for the study of Z boson decay to three photons, Glover93 to the case of o-Ps$`3\gamma `$. In Sec. II we use the extensive symmetries of the decay tensor to show that there are only three independent amplitudes for the o-Ps$`3\gamma `$ decay. In Sec. III we express the decay amplitudes in terms of helicity variables since the spin sums are most convenient in this form. In Sec. IV the integral for the decay rate is reduced to its minimal two-dimensional form. In Sec. V the preceeding formalism is applied to the lowest-order decay process and the lowest-order decay rate of Ore and Powell Ore49 is reproduced. In Sec. VI the method of Passarino and Veltman Passarino79 for evaluating one-loop integrals is developed. In Sec. VII the one-loop calculation is described. Finally, in Sec. VIII our results for the $`O(\alpha )`$ and part of the $`O(\alpha ^2)`$ decay rates are given. The Appendix contains our explicit form for the one-loop decay distribution.
## II Symmetries of the decay tensor
The decay of the massive vector particle orthopositronium to three photons is described by the matrix element conventions
$$M=ϵ_{1\mu _1}^{}ϵ_{2\mu _2}^{}ϵ_{3\mu _3}^{}ϵ_\alpha M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3),$$
(3)
where the three photons have momenta $`k_i`$ and polarizations $`ϵ_i`$, and the positronium atom has momentum $`P=k_1+k_2+k_3`$ and polarization $`ϵ`$. The decay tensor is a linear combination of terms like $`k_a^{\mu _1}k_b^{\mu _2}k_c^{\mu _3}k_d^\alpha `$, $`k_a^{\mu _1}k_b^{\mu _2}g^{\mu _3\alpha }`$, and $`g^{\mu _1\mu _2}g^{\mu _3\alpha }`$. The most general such tensor has 81 terms of the first type, 54 of the second, and 3 of the third. However, gauge invariance and Bose symmetry reduce the number of independent contributions to only three Glover93 . We review the argument below.
Because the decay tensor is always contracted with physical polarization vectors of on-shell photons, which satisfy $`ϵ_{a\mu }k_a^\mu =0`$ (for $`a=1`$, 2, or 3), we can drop terms containing factors of $`k_1^{\mu _1}`$, $`k_2^{\mu _2}`$, and $`k_3^{\mu _3}`$. This leaves only 24 terms of the first type, 30 of the second, and still 3 of the third.
By Bose symmetry, the tensor $`M`$ is totally symmetric under photon interchange. This means, for the interchange of photons 1 and 2, that
$$M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)=M^{\mu _2\mu _1\mu _3\alpha }(k_2,k_1,k_3).$$
(4)
This symmetry leaves only four independent terms of the first type, six of the second, and one of the third. The decay tensor can be written in the manifestly symmetric way
$$M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)=\underset{S_3}{}^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3),$$
(5)
where the sum is over the six photon permutations, and the tensor $``$ has the form
$`^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)`$ $`=`$ $`a_1(k_1,k_2,k_3)k_3^{\mu _1}k_1^{\mu _2}k_1^{\mu _3}k_1^\alpha +a_2(k_1,k_2,k_3)k_3^{\mu _1}k_3^{\mu _2}k_1^{\mu _3}k_1^\alpha `$ (6)
$`+`$ $`a_3(k_1,k_2,k_3)k_3^{\mu _1}k_3^{\mu _2}k_2^{\mu _3}k_1^\alpha +a_4(k_1,k_2,k_3)k_3^{\mu _1}k_1^{\mu _2}k_2^{\mu _3}k_1^\alpha `$ (7)
$`+`$ $`b_1(k_1,k_2,k_3)k_1^{\mu _2}k_1^\alpha g^{\mu _1\mu _3}+b_2(k_1,k_2,k_3)k_3^{\mu _2}k_1^\alpha g^{\mu _1\mu _3}`$ (8)
$`+`$ $`b_3(k_1,k_2,k_3)k_3^{\mu _1}k_1^\alpha g^{\mu _2\mu _3}+b_4(k_1,k_2,k_3)k_1^{\mu _2}k_2^{\mu _3}g^{\mu _1\alpha }`$ (9)
$`+`$ $`b_5(k_1,k_2,k_3)k_1^{\mu _2}k_1^{\mu _3}g^{\mu _1\alpha }+b_6(k_1,k_2,k_3)k_3^{\mu _2}k_2^{\mu _3}g^{\mu _1\alpha }`$ (10)
$`+`$ $`c(k_1,k_2,k_3)g^{\mu _1\alpha }g^{\mu _2\mu _3}.`$ (11)
The quantities $`a_i(k_1,k_2,k_3)`$, $`b_i(k_1,k_2,k_3)`$, and $`c(k_1,k_2,k_3)`$ are scalar functions of their arguments, and $`b_5`$, $`b_6`$, and $`c`$ are symmetric under the interchange $`k_2k_3`$.
Gauge invariance requires that the tensor $`M`$ be transverse
$$k_{1\mu _1}M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)=0,$$
(12)
with similar relations holding for contractions with $`k_{2\mu _2}`$ and $`k_{3\mu _3}`$. The condition of Eq. (12) provides 13 independent relations among the 19 variables $`a_i(k_1,k_2,k_3)`$, $`a_i(k_1,k_3,k_2)`$, $`b_i(k_1,k_2,k_3)`$, $`b_i(k_1,k_3,k_2)`$, $`b_5(k_1,k_2,k_3)`$, $`b_6(k_1,k_2,k_3)`$, and $`c(k_1,k_2,k_3)`$ (with $`i=1,2,3,4`$), which lead (using permutation symmetry) to three independent solutions. So, the tensor $``$ can be expressed in terms of three scalar functions $`A_1`$, $`A_2`$, and $`A_3`$ as
$`^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)`$ $`=`$ $`A_1(k_1,k_2,k_3){\displaystyle \frac{1}{k_1k_3}}\left({\displaystyle \frac{k_3^{\mu _1}k_1^{\mu _3}}{k_1k_3}}g^{\mu _1\mu _3}\right)k_1^\alpha \left({\displaystyle \frac{k_3^{\mu _2}}{k_2k_3}}{\displaystyle \frac{k_1^{\mu _2}}{k_1k_2}}\right)`$ (13)
$`+`$ $`A_2(k_1,k_2,k_3)\{{\displaystyle \frac{1}{k_2k_3}}({\displaystyle \frac{k_1^\alpha k_3^{\mu _1}}{k_1k_3}}g^{\alpha \mu _1})({\displaystyle \frac{k_1^{\mu _2}k_2^{\mu _3}}{k_1k_2}}g^{\mu _2\mu _3})`$
$`+{\displaystyle \frac{1}{k_1k_3}}({\displaystyle \frac{k_1^{\mu _2}}{k_1k_2}}{\displaystyle \frac{k_3^{\mu _2}}{k_2k_3}})(k_1^{\mu _3}g^{\alpha \mu _1}k_1^\alpha g^{\mu _1\mu _3})\}`$
$`+`$ $`A_3(k_1,k_2,k_3){\displaystyle \frac{1}{k_1k_3}}\left({\displaystyle \frac{k_1^\alpha k_3^{\mu _1}}{k_1k_3}}g^{\alpha \mu _1}\right)\left({\displaystyle \frac{k_3^{\mu _2}k_2^{\mu _3}}{k_2k_3}}g^{\mu _2\mu _3}\right).`$
The amplitudes $`A_1`$, $`A_2`$, and $`A_3`$ can be identified by writing the decay tensor as
$`M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)=`$ $``$ $`A_1(k_1,k_2,k_3)k_3^{\mu _1}k_1^{\mu _2}k_1^{\mu _3}k_1^\alpha \left[(k_1k_3)^2(k_1k_2)\right]^1`$ (14)
$`+`$ $`A_2(k_1,k_2,k_3)k_3^{\mu _1}k_1^{\mu _2}k_2^{\mu _3}k_1^\alpha \left[(k_1k_2)(k_2k_3)(k_3k_1)\right]^1`$
$`+`$ $`A_3(k_1,k_2,k_3)k_3^{\mu _1}k_3^{\mu _2}k_2^{\mu _3}k_1^\alpha \left[(k_1k_3)^2(k_2k_3)\right]^1`$
$`+`$ $`\mathrm{}.`$
One finds $`A_1`$, $`A_2`$, and $`A_3`$ by taking the coefficients of $`k_3^{\mu _1}k_1^{\mu _2}k_1^{\mu _3}k_1^\alpha `$, $`k_3^{\mu _1}k_1^{\mu _2}k_2^{\mu _3}k_1^\alpha `$, and $`k_3^{\mu _1}k_3^{\mu _2}k_2^{\mu _3}k_1^\alpha `$.
## III Helicity Amplitudes
The formula for the $`\text{o-Ps}3\gamma `$ decay rate involves the absolute square of the decay matrix element summed over final state spins and averaged over the initial state spin:
$$\overline{|M|^2}=\underset{ϵ_1,ϵ_2,ϵ_3}{}\frac{1}{3}\underset{ϵ}{}|M|^2.$$
(15)
This is a Lorentz invariant quantity, and can be calculated in any frame. It is convenient to calculate it in a two-photon rest frame.
Since we will use the orthopositronium center of mass frame for the decay rate integration, it is useful to express our results in terms of invariant variables. A convenient set is given by the Mandelstam variables, which are defined by
$$s_{ij}=s_{ji}=(k_i+k_j)^2=2k_ik_j,$$
(16)
and satisfy
$$s_{ij}+s_{jk}+s_{ki}=M_{\mathrm{Ps}}^2,$$
(17)
where $`M_{\mathrm{Ps}}`$ here is the orthopositronium mass and $`\{i,j,k\}`$ is any permutation of $`\{1,2,3\}`$. Bar variables are defined by
$$\overline{s}_{ij}=M_{\mathrm{Ps}}^2s_{ij}=s_{ik}+s_{jk}.$$
(18)
They satisfy
$$\overline{s}_{ij}+\overline{s}_{jk}+\overline{s}_{ki}=2M_{\mathrm{Ps}}^2.$$
(19)
We note that each $`s_{ij}`$ and $`\overline{s}_{ij}`$ is non-negative.
We calculate $`\overline{|M|^2}`$ in the $`k_1k_2`$ rest frame. The photon and positronium momentum vectors in $`(E,p_x,p_y,p_z)`$ notation are given by
$`k_1`$ $`=`$ $`(k,0,0,k)`$ (20)
$`k_2`$ $`=`$ $`(k,0,0,k)`$ (21)
$`k_3`$ $`=`$ $`(q,q\mathrm{sin}\theta ,0,q\mathrm{cos}\theta )`$ (22)
$`P`$ $`=`$ $`(E,q\mathrm{sin}\theta ,0,q\mathrm{cos}\theta ),`$ (23)
where $`E^2=q^2+M_{\mathrm{Ps}}^2`$. The $`k_1k_2`$ rest frame kinematic variables are given in terms of invariants by
$`k`$ $`=`$ $`{\displaystyle \frac{\sqrt{s_{12}}}{2}},`$ (24)
$`q`$ $`=`$ $`{\displaystyle \frac{\overline{s}_{12}}{2\sqrt{s_{12}}}},`$ (25)
$`E`$ $`=`$ $`{\displaystyle \frac{M_{\mathrm{Ps}}^2+s_{12}}{\overline{s}_{12}}},`$ (26)
$`\mathrm{sin}\theta `$ $`=`$ $`{\displaystyle \frac{2\sqrt{s_{13}s_{23}}}{\overline{s}_{12}}}.`$ (27)
The helicity vectors for photon 1 are
$`\widehat{e}_1^+`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(0,1,i,0),`$ (28)
$`\widehat{e}_1^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(0,1,i,0).`$ (29)
For photon 2 we rotate these by $`180^{}`$ around the $`y`$ axis using
$$R_2=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)$$
(30)
to find
$$\widehat{e}_2^\pm =R_2\widehat{e}_1^\pm =\widehat{e}_1^{}.$$
(31)
For photon 3 we rotate using
$$R_3=\left(\begin{array}{ccc}\mathrm{cos}\theta & 0& \mathrm{sin}\theta \\ 0& 1& 0\\ \mathrm{sin}\theta & 0& \mathrm{cos}\theta \end{array}\right)$$
(32)
to find
$`\widehat{e}_3^+`$ $`=`$ $`R_3\widehat{e}_1^+={\displaystyle \frac{1}{\sqrt{2}}}(0,\mathrm{cos}\theta ,i,\mathrm{sin}\theta ),`$ (33)
$`\widehat{e}_3^{}`$ $`=`$ $`R_3\widehat{e}_1^{}={\displaystyle \frac{1}{\sqrt{2}}}(0,\mathrm{cos}\theta ,i,\mathrm{sin}\theta ).`$ (34)
The positronium spin $`\pm 1`$ helicity vectors are the same as those for photon 3:
$$\widehat{e}_{Ps}^\pm =\widehat{e}_3^\pm ,$$
(35)
while the positronium spin 0 helicity vector is
$$\widehat{e}_{Ps}^0=\frac{1}{M_{\mathrm{Ps}}}(q,E\mathrm{sin}\theta ,0,E\mathrm{cos}\theta ).$$
(36)
The helicity amplitudes are defined by
$$M_{\lambda _1\lambda _2\lambda _3\lambda }=\widehat{e}_{1\mu _1}^{\lambda _1}\widehat{e}_{2\mu _2}^{\lambda _2}\widehat{e}_{3\mu _3}^{\lambda _3}e_{Ps\alpha }^\lambda M^{\mu _1\mu _2\mu _3\alpha }(k_1,k_2,k_3)$$
(37)
There are nine independent helicity amplitudes with $`\lambda _1=+`$. They are
$`M_{++++}`$ $`=`$ $`2\{{\displaystyle \frac{A_1(123)+A_1(132)}{\overline{s}_{12}}}+{\displaystyle \frac{A_2(123)A_2(132)}{\overline{s}_{12}}}`$ (38b)
$`{\displaystyle \frac{s_{23}A_3(132)}{s_{12}\overline{s}_{12}}}{\displaystyle \frac{A_3(312)}{s_{23}}}+(12)\},`$
$`M_{+++}`$ $`=`$ $`2\{{\displaystyle \frac{A_1(123)A_1(132)}{\overline{s}_{12}}}+{\displaystyle \frac{A_2(123)+A_2(132)}{\overline{s}_{12}}}`$ (38d)
$`{\displaystyle \frac{A_3(123)}{s_{13}}}{\displaystyle \frac{s_{13}A_3(132)}{s_{12}\overline{s}_{12}}}+(12)\},`$
$`M_{+++}`$ $`=`$ $`2\{{\displaystyle \frac{A_1(132)}{\overline{s}_{12}}}+(12)\},`$ (38e)
$`M_{++}`$ $`=`$ $`2\{{\displaystyle \frac{A_1(132)}{\overline{s}_{12}}}+{\displaystyle \frac{A_2(123)A_2(132)}{s_{23}}}{\displaystyle \frac{A_3(312)}{s_{23}}}+(12)\},`$ (38f)
$`M_{+++}`$ $`=`$ $`2\{{\displaystyle \frac{A_1(123)}{\overline{s}_{12}}}+{\displaystyle \frac{A_2(123)A_2(132)}{\overline{s}_{12}}}+{\displaystyle \frac{A_2(312)A_2(321)}{s_{12}}}`$ (38h)
$`{\displaystyle \frac{A_3(213)}{s_{23}}}{\displaystyle \frac{s_{23}A_3(231)}{s_{12}\overline{s}_{12}}}\},`$
$`M_{++}`$ $`=`$ $`2\left\{{\displaystyle \frac{A_1(123)}{\overline{s}_{12}}}+{\displaystyle \frac{s_{13}\left(A_2(123)A_2(132)\right)}{s_{23}\overline{s}_{12}}}{\displaystyle \frac{s_{13}A_3(231)}{s_{12}\overline{s}_{12}}}\right\},`$ (38i)
$`M_{+++0}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Delta }}{M_{\mathrm{Ps}}^2}}\{{\displaystyle \frac{r_{13}\left(A_1(123)A_1(132)\right)}{\overline{s}_{12}}}\overline{s}_{12}A_1(312)+{\displaystyle \frac{r_{13}\left(A_2(123)+A_2(132)\right)}{\overline{s}_{12}}}`$ (38l)
$`+\overline{s}_{12}A_2(312)+{\displaystyle \frac{s_{12}s_{23}A_3(123)}{s_{13}}}+{\displaystyle \frac{(M_{\mathrm{Ps}}^2+s_{12})s_{13}s_{23}A_3(132)}{s_{12}\overline{s}_{12}}}`$
$`+{\displaystyle \frac{s_{12}s_{13}A_3(312)}{s_{23}}}(12)\},`$
$`M_{++0}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Delta }}{M_{\mathrm{Ps}}^2}}\{{\displaystyle \frac{r_{13}A_1(132)r_{23}A_1(231)}{\overline{s}_{12}}}+s_{12}(A_2(123)+A_2(132)A_2(213)A_2(231))`$ (38n)
$`{\displaystyle \frac{s_{12}s_{13}A_3(312)}{s_{23}}}+{\displaystyle \frac{s_{12}s_{23}A_3(321)}{s_{13}}}\},`$
$`M_{++0}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Delta }}{M_{\mathrm{Ps}}^2}}\{{\displaystyle \frac{r_{13}A_1(123)}{\overline{s}_{12}}}\overline{s}_{12}A_1(321)`$ (38q)
$`+{\displaystyle \frac{(M_{\mathrm{Ps}}^2+s_{12})s_{13}\left(A_2(123)+A_2(132)\right)}{\overline{s}_{12}}}+s_{13}\left(A_2(312)+A_2(321)\right)`$
$`+{\displaystyle \frac{s_{12}s_{13}A_3(213)}{s_{23}}}+{\displaystyle \frac{(M_{\mathrm{Ps}}^2+s_{12})s_{13}s_{23}A_3(231)}{s_{12}\overline{s}_{12}}}\},`$
where we have used the abbreviated notation $`A_i(abc)=A_i(k_a,k_b,k_c)`$ and the definitions
$`r_{ij}`$ $`=`$ $`M_{\mathrm{Ps}}^2s_{ij}s_{ik}s_{jk},`$ (39)
$`\mathrm{\Delta }`$ $`=`$ $`\sqrt{{\displaystyle \frac{M_{\mathrm{Ps}}^2}{2s_{12}s_{13}s_{23}}}}.`$ (40)
The other three $`\lambda _1=+`$ amplitudes are related to the previous ones by
$`M_{++}`$ $`=`$ $`M_{++}(12),`$ (41)
$`M_+`$ $`=`$ $`M_{+++}(12),`$ (42)
$`M_{+0}`$ $`=`$ $`M_{++0}(12).`$ (43)
The $`\lambda _1=`$ amplitudes are given by the parity relations
$`M_{\lambda _2\lambda _3\pm }`$ $`=`$ $`M_{+\lambda _2\lambda _3},`$ (44)
$`M_{\lambda _2\lambda _30}`$ $`=`$ $`M_{+\lambda _2\lambda _30}.`$ (45)
The squared decay matrix element can be written as
$$\overline{|M|^2}=\frac{2}{3}\underset{\lambda _2,\lambda _3,\lambda }{}|M_{+\lambda _2\lambda _3\lambda }|^2.$$
(46)
## IV The Decay Rate Integral
We will calculate the $`\text{o-Ps}3\gamma `$ decay rate integral in the positronium center of mass frame. The decay rate integral is given by
$$\mathrm{\Gamma }=\frac{1}{3!}\frac{1}{2M_{\mathrm{Ps}}}\frac{d^3k_1}{(2\pi )^32\omega _1}\frac{d^3k_2}{(2\pi )^32\omega _2}\frac{d^3k_3}{(2\pi )^32\omega _3}(2\pi )^4\delta (Pk_1k_2k_3)\overline{|M|^2}$$
(47)
where $`w_i=|\stackrel{}{k}_i|`$ are the photon energies. Of the nine variables in $`\stackrel{}{k}_1`$, $`\stackrel{}{k}_2`$, $`\stackrel{}{k}_3`$, four are determined in terms of the others by energy-momentum conservation
$`\omega _1+\omega _2+\omega _3`$ $`=`$ $`M_{\mathrm{Ps}},`$ (48)
$`\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3`$ $`=`$ $`0.`$ (49)
Three variables describe the orientation in space of the decay plane. The remaining two variables describe the relative orientation of the photons in the decay plane. We will use the energies of two of the photons for this last pair of variables. Each photon can have any energy between $`0`$ and $`W=M_{\mathrm{Ps}}/2`$. We find it convenient to introduce dimensionless variables $`x_i=\omega _i/W`$ which satisfy $`0x_i1`$, $`x_1+x_2+x_3=2`$ and are given in terms of invariants by $`x_i=\overline{s}_{jk}/M_{\mathrm{Ps}}^2`$. In terms of the $`x`$’s, one has
$$\mathrm{\Gamma }=\frac{W}{768\pi ^3}_0^1𝑑x_1_{1x_1}^1𝑑x_2\overline{|M|^2}.$$
(50)
## V The Lowest Order Decay Rate
The lowest order decay amplitude is given by
$$M_{\mathrm{LO}}=\underset{S_3}{}\mathrm{tr}\left[(ie\gamma ϵ_3^{})\frac{i}{\gamma (P/2+k_3)m}(ie\gamma ϵ_2^{})\frac{i}{\gamma (P/2k_1)m}(ie\gamma ϵ_1^{})\mathrm{\Psi }\right],$$
(51)
where the sum is over the six permutations of the final state photons. The wave function factor is given by
$$\mathrm{\Psi }=\sqrt{2M_{\mathrm{Ps}}}\left(\begin{array}{cc}0& \stackrel{}{\sigma }\widehat{ϵ}/\sqrt{2}\\ 0& 0\end{array}\right)\varphi _0,$$
(52)
which contains the spin-1 spin factor, a normalization factor, and the wave function at contact
$$\varphi _0=\sqrt{\frac{m^3\alpha ^3}{8\pi }}.$$
(53)
We write
$$\left(\begin{array}{cc}0& \stackrel{}{\sigma }\widehat{ϵ}\\ 0& 0\end{array}\right)=\frac{1}{4}(\gamma N+1)\gamma ϵ(\gamma N1)$$
(54)
for the positronium spin factor where $`N=P/(2W)`$. The lowest order decay amplitude (see Fig. 1) becomes
$$M_{\mathrm{LO}}=\frac{i\pi \alpha ^3}{4}\underset{S_3}{}\frac{1}{x_1x_3}\mathrm{tr}\left[\gamma ϵ_3^{}(\gamma R_3+1)\gamma ϵ_2^{}(\gamma R_1+1)\gamma ϵ_1^{}(\gamma N+1)\gamma ϵ(\gamma N1)\right]$$
(55)
where $`R_i=NK_i`$, $`K_i=k_i/W`$ and $`Wm`$. The lowest order decay tensor has the corresponding form
$$M_{\mathrm{LO}}^{\mu _1\mu _2\mu _3\alpha }=\frac{i\pi \alpha ^3}{4}\underset{S_3}{}\frac{1}{x_1x_3}\mathrm{tr}\left[\gamma ^{\mu _3}(\gamma R_3+1)\gamma ^{\mu _2}(\gamma R_1+1)\gamma ^{\mu _1}(\gamma N+1)\gamma ^\alpha (\gamma N1)\right].$$
(56)
We replace $`N`$ by $`(K_1+K_2+K_3)/2`$ and expand this out, and identify the $`A_i`$ functions by use of Eq. (14). The lowest order functions are
$`A_1^{\mathrm{LO}}(x_1,x_2,x_3)`$ $`=`$ $`0,`$ (57)
$`A_2^{\mathrm{LO}}(x_1,x_2,x_3)`$ $`=`$ $`16i\pi m^2\alpha ^3{\displaystyle \frac{\overline{x}_1\overline{x}_2\overline{x}_3}{x_1x_2x_3}},`$ (58)
$`A_3^{\mathrm{LO}}(x_1,x_2,x_3)`$ $`=`$ $`0,`$ (59)
where $`x_i=\overline{s}_{jk}/M_{\mathrm{Ps}}^2`$ and $`\overline{x}_i=1x_i=s_{jk}/M_{\mathrm{Ps}}^2`$. Clearly, $`A_2^{\mathrm{LO}}`$ is a factor in each helicity amplitude. One has
$`M_{++}^{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{x_3}{\overline{x}_1\overline{x}_2}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (60)
$`M_{+++}^{\mathrm{LO}}`$ $`=`$ $`M_+^{\mathrm{LO}}={\displaystyle \frac{1}{x_3\overline{x}_3}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (61)
$`M_{++}^{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{\overline{x}_2}{\overline{x}_1x_3}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (62)
$`M_{++}^{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{\overline{x}_1}{\overline{x}_2x_3}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (63)
$`M_{++0}^{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{\overline{x}_2}{x_3}}\sqrt{{\displaystyle \frac{2}{\overline{x}_1\overline{x}_2\overline{x}_3}}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (64)
$`M_{+0}^{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{\overline{x}_1}{x_3}}\sqrt{{\displaystyle \frac{2}{\overline{x}_1\overline{x}_2\overline{x}_3}}}{\displaystyle \frac{A_2^{\mathrm{LO}}}{m^2}},`$ (65)
with all other $`M_{+\lambda _2\lambda _3\lambda }^{\mathrm{LO}}`$ amplitudes equal to zero. One finds that
$$\overline{|M_{\mathrm{LO}}|^2}=\frac{512}{3}\pi ^2\alpha ^6\left\{\left(\frac{\overline{x}_1}{x_2x_3}\right)^2+\left(\frac{\overline{x}_2}{x_1x_3}\right)^2+\left(\frac{\overline{x}_3}{x_1x_2}\right)^2\right\}.$$
(66)
The lowest order decay rate is the Ore and Powell result Ore49
$`\mathrm{\Gamma }_{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{2}{9\pi }}m\alpha ^6{\displaystyle _0^1}𝑑x_1{\displaystyle _{1x_1}^1}𝑑x_2\left\{\left({\displaystyle \frac{\overline{x}_1}{x_2x_3}}\right)^2+\left({\displaystyle \frac{\overline{x}_2}{x_1x_3}}\right)^2+\left({\displaystyle \frac{\overline{x}_3}{x_1x_2}}\right)^2\right\}`$ (67)
$`=`$ $`{\displaystyle \frac{2}{9\pi }}(\pi ^29)m\alpha ^6.`$ (68)
## VI One-Loop Integrals
We used the method of Passarino and Veltman Passarino79 to evaluate the one-loop integrals. Since this approach is widely used, and lengthy to describe in detail, we will just list the one-loop integrals that are required but only work the scalar integrals out in detail. The Passarino-Veltman method will be illustrated in the case of the three-point functions.
The general definition of the one-loop form factors is through
$`\{X_0,X_\mu ,X_{\mu \nu },\mathrm{}\}`$ $`=`$ $`\mu ^{2ϵ}{\displaystyle (dq)_n^{\prime \prime }\{1,q_\mu ,q_\mu q_\nu ,\mathrm{}\}}`$ (69)
$`\times `$ $`\left[(q^2+m_1^2)((q+p_1)^2+m_2^2)((q+p_1+p_2)^2+m_3^2)\mathrm{}\right]^1.`$ (70)
Ultraviolet divergences are controlled through dimensional regularization with $`n=42ϵ`$ the dimensionality of spacetime. We define $`(dq)_n^{\prime \prime }=d^nq/(i\pi ^{n/2})`$. The quantity $`\mu `$ is a reference mass introduced with the regularization which we take to be equal to the electron mass $`m`$. Functional dependences on the masses and momenta are indicated by $`X(m_1,m_2,m_3,\mathrm{};p_1,p_2,\mathrm{})`$.
The one-point function is trivially evaluated:
$`A(m_1)`$ $`=`$ $`m^{2ϵ}{\displaystyle (dq)_n^{\prime \prime }\frac{1}{(q^2+m_1^2)}}`$ (71)
$`=`$ $`m_1^2\mathrm{\Gamma }(ϵ)+\overline{A}(m_1)+O(ϵ)`$ (72)
where
$$\overline{A}(m_1)=m_1^2\left[1\mathrm{ln}(m_1^2/m^2)\right].$$
(73)
The two-point functions are defined by
$$\{B_0,B_\mu ,B_{\mu \nu }\}=m^{2ϵ}(dq)_n^{\prime \prime }\{1,q_\mu ,q_\mu q_\nu \}\left[(q^2+m_1^2)((q+p)^2+m_2^2)\right]^1.$$
(74)
The scalar function $`B_0`$ is
$`B_0(m_1,m_2;p)`$ $`=`$ $`m^{2ϵ}{\displaystyle 𝑑x\frac{\mathrm{\Gamma }(ϵ)}{\mathrm{\Delta }^ϵ}}`$ (75)
$`=`$ $`\mathrm{\Gamma }(ϵ)+\overline{B}_0(m_1,m_2;p)`$ (76)
where $`\mathrm{\Delta }=(1x)m_1^2+xm_2^2x(1x)p^2`$ and
$$\overline{B}_0(m_1,m_2;p)=𝑑x\mathrm{ln}\left(\mathrm{\Delta }/m^2\right).$$
(77)
All parametric integrals will be taken between the limits $`0`$ and $`1`$. The cases of interest are
$$\overline{B}_0(0,m;p)=2+\frac{1\rho }{\rho }\mathrm{ln}(1\rho )$$
(78)
where $`\rho =p^2/m^2`$, and
$$\overline{B}_0(m,m;p)=2\left\{1\sqrt{\frac{4\rho }{\rho }}\mathrm{arctan}\sqrt{\frac{\rho }{4\rho }}\right\},$$
(79)
valid for $`0\rho 4`$.
The three-point functions are defined by
$`\{C_0,C_\mu ,C_{\mu \nu },C_{\mu \nu \alpha }\}`$ $`=`$ $`m^{2ϵ}{\displaystyle (dq)_n^{\prime \prime }\{1,q_\mu ,q_\mu q_\nu ,q_\mu q_\nu q_\alpha \}}`$ (80)
$`\times `$ $`\left[(q^2+m_1^2)((q+p_1)^2+m_2^2)((q+p_1+p_2)^2+m_3^2)\right]^1.`$ (81)
The general forms for $`C_\mu `$, $`C_{\mu \nu }`$, and $`C_{\mu \nu \alpha }`$ are
$`C_\mu `$ $`=`$ $`p_{1\mu }C_{11}+p_{2\mu }C_{12},`$ (82a)
$`C_{\mu \nu }`$ $`=`$ $`p_{1\mu }p_{1\nu }C_{21}+p_{2\mu }p_{2\nu }C_{22}+\{p_1p_2\}_{\mu \nu }C_{23}+g_{\mu \nu }C_{24},`$ (82b)
$`C_{\mu \nu \alpha }`$ $`=`$ $`p_{1\mu }p_{1\nu }p_{1\alpha }C_{31}+p_{2\mu }p_{2\nu }p_{2\alpha }C_{32}+\{p_1p_1p_2\}_{\mu \nu \alpha }C_{33}`$ (82d)
$`+\{p_1p_2p_2\}_{\mu \nu \alpha }C_{34}+\{p_1g\}_{\mu \nu \alpha }C_{35}+\{p_2g\}_{\mu \nu \alpha }C_{36},`$
where
$`\{pk\}_{\mu \nu }`$ $`=`$ $`p_\mu k_\nu +k_\mu p_\nu ,`$ (83)
$`\{ppk\}_{\mu \nu \alpha }`$ $`=`$ $`p_\mu p_\nu k_\alpha +p_\mu k_\nu p_\alpha +k_\mu p_\nu p_\alpha ,`$ (84)
$`\{pg\}_{\mu \nu \alpha }`$ $`=`$ $`p_\mu g_{\nu \alpha }+p_\nu g_{\mu \alpha }+p_\alpha g_{\mu \nu }.`$ (85)
The only divergent terms here are $`C_{24}`$, $`C_{35}`$, and $`C_{36}`$.
We illustrate the Passarino-Veltman procedure by describing the evaluation of $`C_{11}`$ and $`C_{12}`$ in terms of $`C_0`$ and the $`B`$ functions. We start by multiplying Eq. (82a) by $`p_1^\mu `$ and $`p_2^\mu `$:
$`p_1^2C_{11}+p_{12}C_{12}`$ $`=`$ $`qp_1_CR_1,`$ (86a)
$`p_{12}C_{11}+p_2^2C_{12}`$ $`=`$ $`qp_2_CR_2,`$ (86b)
where $`p_{ij}=p_ip_j`$ and $`_C`$ is the integral operator on the RHS of Eq. (80) (so that for example $`C_0=1_C`$). We write Eqs. (86) as
$$X\left(\begin{array}{c}C_{11}\\ C_{12}\end{array}\right)=\left(\begin{array}{c}R_1\\ R_2\end{array}\right),X=\left(\begin{array}{cc}p_1^2& p_{12}\\ p_{12}& p_2^2\end{array}\right)$$
(87)
with the solution
$$\left(\begin{array}{c}C_{11}\\ C_{12}\end{array}\right)=X^1\left(\begin{array}{c}R_1\\ R_2\end{array}\right).$$
(88)
We find $`R_1`$ and $`R_2`$ by noting that
$`qp_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{(q^2+m_1^2)+((q+p_1)^2+m_2^2)+f_1\right\},`$ (89a)
$`qp_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{((q+p_1)^2+m_2^2)+((q+p_1+p_2)^2+m_3^2)+f_2\right\},`$ (89b)
where
$`f_1`$ $`=`$ $`m_1^2m_2^2+p_1^2,`$ (90a)
$`f_2`$ $`=`$ $`m_2^2m_3^2+(p_1+p_2)^2p_1^2.`$ (90b)
Then we see that
$`R_1`$ $`=`$ $`qp_1_C={\displaystyle \frac{1}{2}}\{B_0(m_2,m_3;p_2)+B_0(m_1,m_3;p_1+p_2)`$ (91b)
$`+f_1C_0(m_1,m_2,m_3;p_1,p_2)\},`$
$`R_2`$ $`=`$ $`qp_2_C={\displaystyle \frac{1}{2}}\{B_0(m_1,m_3;p_1+p_2)+B_0(m_1,m_2;p_1)`$ (91d)
$`+f_2C_0(m_1,m_2,m_3;p_1,p_2)\}.`$
Since the $`B`$ functions are already known, only the scalar $`C_0`$ function remains to be computed. Similarly, the $`C_{2i}`$ functions can be evaluated in terms of the $`C_{1i}`$’s and $`B`$’s, etc. At each level in the ladder, only the scalar functions are new.
The general three-point scalar integral is
$`C_0(m_1,m_2,m_3;p_1,p_2)`$ $`=`$ $`{\displaystyle }(dq)^{\prime \prime }[(q^2+m_1^2)((q+p_1)^2+m_2^2)`$ (93)
$`\times ((q+p_1+p_2)^2+m_3^2)]^1`$
$`=`$ $`{\displaystyle 𝑑z𝑑x\frac{z}{\mathrm{\Delta }}},`$ (94)
where the limit $`n4`$ has been taken since $`C_0`$ is ultraviolet finite, and $`(dq)^{\prime \prime }(dq)_4^{\prime \prime }=d^4q/(i\pi ^2)`$. For $`\mathrm{\Delta }`$ one finds
$$\mathrm{\Delta }=(1z)m_1^2+z(1x)m_2^2+zxm_3^2z(1z)p_1^2xz(1z)2p_{12}xz(1xz)p_2^2.$$
(95)
The cases of interest here are
$$C_0(0,m,m;p_1,p_2)=\frac{1}{2p_{12}}\left\{\mathrm{Li}_2\left(\frac{p_1^2+2p_{12}}{m^2}\right)\mathrm{Li}_2\left(\frac{p_1^2}{m^2}\right)\right\}$$
(96)
which holds when $`p_2^2=0`$;
$$C_0(0,m,m;p_1,p_2)=\frac{1}{2m^2\alpha }\left\{\mathrm{Li}_2(12\alpha )+2\zeta (2)2\left(\mathrm{arctan}\sqrt{\frac{1\alpha }{\alpha }}\right)^2\right\},$$
(97)
which holds when $`p_1^2=m^2`$ and $`(2p_1+p_2)^2=0`$ and where $`\alpha =2+p_{12}/m^2`$; and
$$C_0(m,m,m;p_1,p_2)=\frac{1}{2p_{12}}\left\{\mathrm{L}\left(\frac{(p_1+p_2)^2}{m^2}\right)\mathrm{L}\left(\frac{p_1^2}{m^2}\right)\right\}$$
(98)
where $`p_2^2=0`$ and
$$\mathrm{L}(s)=𝑑z\frac{\mathrm{ln}(1z(1z)s)}{(1z)}=2\left(\mathrm{arctan}\sqrt{\frac{s}{4s}}\right)^2.$$
(99)
The dilogarithm function $`\mathrm{Li}_2(x)`$ is discussed in detail by Lewin. Lewin81
The four-point functions are defined by
$`\{D_0,D_\mu ,D_{\mu \nu },D_{\mu \nu \alpha },D_{\mu \nu \alpha \beta }\}`$ $`=`$ $`{\displaystyle (dq)^{\prime \prime }\{1,q_\mu ,q_\mu q_\nu ,q_\mu q_\nu q_\alpha ,q_\mu q_\nu q_\alpha q_\beta \}}`$ (100)
$`\times `$ $`[(q^2+m_1^2)((q+p_1)^2+m_2^2)((q+p_1+p_2)^2+m_3^2)`$ (101)
$`\times `$ $`((q+p_1+p_2+p_3)^2+m_4^2)]^1.`$ (102)
We have dispensed with the regularization here since all of the $`D`$ functions needed for our calculation are ultraviolet finite. While general expressions for $`D_0`$ exist, we need only a few special cases. In particular, we find that
$$D_0(0,m,m,m;mN,k_1,k_2)=\frac{8}{\sqrt{s_{12}\overline{s}_{12}s_{23}\overline{s}_{23}}}\left\{\mathrm{Li}_2(\frac{x_+}{\sqrt{D}},\theta )\mathrm{Li}_2(\frac{x_{}}{\sqrt{D}},\theta )\right\},$$
(103)
where
$`D`$ $`=`$ $`{\displaystyle \frac{m^2\overline{s}_{12}}{s_{12}\overline{s}_{23}}},`$ (104a)
$`x_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1\pm \sqrt{{\displaystyle \frac{s_{23}\overline{s}_{12}}{s_{12}\overline{s}_{23}}}}\right),`$ (104b)
$`\mathrm{tan}\theta `$ $`=`$ $`\sqrt{{\displaystyle \frac{\overline{s}_{23}}{s_{23}}}},`$ (104c)
and $`\mathrm{Li}_2(r,\theta )`$ is the dilogarithm of complex argument. Lewin81 By some transformations among the momentum vectors, one can show that
$$D_0(0,m,m,m;mNk_1,k_2,k_3)=D_0(0,m,m,m;mN,k_3,k_2).$$
(105)
Finally, we also have
$`D_0(m,m,m,m;k_1,k_2,k_3)`$ $`=`$ $`{\displaystyle \frac{4}{\sqrt{s_{12}\overline{s}_{12}s_{23}\overline{s}_{23}}}}\{\mathrm{Li}_2({\displaystyle \frac{y_+}{\sqrt{D_1}}},\theta _1)\mathrm{Li}_2({\displaystyle \frac{y_{}}{\sqrt{D_1}}},\theta _1)`$ (108)
$`+\mathrm{Li}_2({\displaystyle \frac{y_+}{\sqrt{D_3}}},\theta _3)\mathrm{Li}_2({\displaystyle \frac{y_{}}{\sqrt{D_3}}},\theta _3)`$
$`\mathrm{Li}_2({\displaystyle \frac{y_+}{\sqrt{D_0}}},0)+\mathrm{Li}_2({\displaystyle \frac{y_{}}{\sqrt{D_0}}},0)\},`$
where
$`D_0`$ $`=`$ $`{\displaystyle \frac{\overline{s}_{12}\overline{s}_{23}}{4s_{12}s_{23}}},`$ (109a)
$`D_1`$ $`=`$ $`{\displaystyle \frac{m^2\overline{s}_{23}}{s_{12}s_{23}}},`$ (109b)
$`D_3`$ $`=`$ $`{\displaystyle \frac{m^2\overline{s}_{12}}{s_{12}s_{23}}},`$ (109c)
$`y_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1\pm 2\sqrt{D_0}\right),`$ (109d)
$`\theta _1`$ $`=`$ $`\mathrm{arctan}\sqrt{{\displaystyle \frac{s_{12}}{\overline{s}_{12}}}},`$ (109e)
$`\theta _3`$ $`=`$ $`\mathrm{arctan}\sqrt{{\displaystyle \frac{s_{23}}{\overline{s}_{23}}}}.`$ (109f)
All of these integrals were done directly by way of Feynman parameters.
The five-point functions are required for the ladder diagram. The five-point functions are very difficult to evaluate in general. We require only a special case, where $`m_1=0`$, $`m_2=m_3=m_4=m_5=m`$, $`p_1=mN`$, $`p_2=k_1`$, $`p_3=k_2`$, $`p_4=k_3`$. One feature of this special case is that there is a binding singularity: the scalar five-point function diverges, so we will have to base our implementation of the Passarino-Veltman formalism on the integral of the vector $`q_\mu `$, which is finite, instead of on the divergent scalar integral. Also, we have not yet evaluated the three- and four-point functions with the necessary momenta. We give the three-, four-, and five-point functions with the special case mass and momenta values the names $`E`$, $`F`$, and $`G`$:
$`f_E`$ $`=`$ $`m^{2ϵ}{\displaystyle (dq)_n^{\prime \prime }f\left[(q^2)((q+p_1)^2+m^2)((qp_1)^2+m^2)\right]^1},`$ (110a)
$`f_{F(p_2)}`$ $`=`$ $`m^{2ϵ}{\displaystyle }(dq)_n^{\prime \prime }f[(q^2)((q+p_1)^2+m^2)((qp_1)^2+m^2)`$ (110c)
$`\times ((q+p_1+p_2)^2+m^2)]^1,`$
$`f_{G(p_2,p_3)}`$ $`=`$ $`m^{2ϵ}{\displaystyle }(dq)_n^{\prime \prime }f[(q^2)((q+p_1)^2+m^2)((qp_1)^2+m^2)`$ (110e)
$`\times ((q+p_1+p_2)^2+m^2)((q+p_1+p_2+p_3)^2+m^2)]^1,`$
where $`p_1=p=mN`$. The first two of these are special cases of the three- and four-point functions. Because of the binding singularity, $`E_0=1_E`$, $`F_0=1_F`$, and $`G_0=1_G`$ all diverge. We start our analysis with the vector integrals $`E_\mu =q_\mu _E`$, $`F_\mu =q_\mu _F`$, and $`G_\mu =q_\mu _G`$.
The three-point special case vector integral has the general form
$$E_\mu =q_\mu _E=p_\mu E_1.$$
(111)
It is not hard to show (by use of symmetric integration) that $`E_1=0`$, so that
$$E_\mu =0.$$
(112)
The four-point special case vector integral has the general form
$$F_\mu (p_2)=q_\mu _F=p_{1\mu }F_{11}(p_2)+p_{2\mu }F_{12}(p_2).$$
(113)
The necessary vector integrals for $`p_2=k_1`$ are
$`F_{11}(k_1)`$ $`=`$ $`{\displaystyle \frac{1}{4x_1^2}}\left\{\mathrm{Li}_2(12x_1)2x_1\mathrm{ln}(2x_1)2\theta ^24\sqrt{x_1\overline{x}_1}\theta +2\zeta (2)\right\},`$ (114a)
$`F_{12}(k_1)`$ $`=`$ $`{\displaystyle \frac{1}{x_1}}F_{11}(k_1)+{\displaystyle \frac{1}{8x_1^2}}\left[2\mathrm{L}\mathrm{i}_2(12x_1)+\zeta (2)2\theta ^2\right]`$ (114b)
where
$$\theta =\mathrm{arctan}\sqrt{\frac{\overline{x}_1}{x_1}}.$$
(115)
When $`p_2=k_1k_2`$ one finds
$$F_\mu (k_1k_2)=F_\mu (k_3),$$
(116)
which implies that
$`F_{11}(k_1k_2)`$ $`=`$ $`F_{11}(k_3)+2F_{12}(k_3),`$ (117)
$`F_{12}(k_2k_2)`$ $`=`$ $`F_{12}(k_3).`$ (118)
The five-point special case vector integral has the general form
$$G_\mu (k_1,k_2)=q_\mu _G=p_\mu G_{11}+k_{1\mu }G_{12}+k_{3\mu }G_{13}.$$
(119)
The $`G_{1i}`$ functions are given by
$`G_{11}(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{8\overline{x}_1}}\left[I_0(x_1,x_3)+I_1(x_1,x_3)\right]{\displaystyle \frac{1}{8\overline{x}_3}}\left[I_0(x_3,x_1)+I_1(x_3,x_1)\right],`$ (120a)
$`G_{12}(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{16x_1\overline{x}_1}}\left[(12x_1)I_0(x_1,x_3)I_1(x_1,x_3)\right]`$ (120c)
$`+{\displaystyle \frac{1}{16x_1\overline{x}_3}}\left[I_0(x_3,x_1)+I_1(x_3,x_1)\right],`$
$`G_{13}(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{16\overline{x}_1x_3}}\left[I_0(x_1,x_3)+I_1(x_1,x_3)\right]`$ (120e)
$`{\displaystyle \frac{1}{16x_3\overline{x}_3}}\left[(12x_3)I_0(x_3,x_1)I_1(x_3,x_1)\right],`$
where
$`I_0(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{x_1\overline{x}_1x_3\overline{x}_3}}}\left[\mathrm{Li}_2(r_+,\theta )\mathrm{Li}_2(r_{},\theta )\right],`$ (121a)
$`I_1(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{(x_1x_3)}}\mathrm{ln}\left({\displaystyle \frac{x_1}{x_3}}\right){\displaystyle \frac{2}{\sqrt{x_3\overline{x}_3}}}\mathrm{arctan}\left(\sqrt{{\displaystyle \frac{\overline{x}_3}{x_3}}}\right),`$ (121b)
with $`r_\pm =\sqrt{\overline{x}_1}\pm \sqrt{x_1\overline{x}_3/x_3}`$ and $`\theta =\mathrm{arctan}\sqrt{x_1/\overline{x}_1}`$.
## VII Analysis of the one-loop decay diagrams
The decay amplitudes can be written as
$$A_i=A_i^{(0)}+A_i^{(1)}+A_i^{(2)}+\mathrm{}$$
(122)
for $`i=1,2,3`$, where the superscript indicates the power of $`\alpha `$ above that of the lowest order amplitudes $`A_i^{(0)}=A_i^{LO}`$. (Terms of order $`A_i^{(2)}`$ and higher also contain factors of $`\mathrm{ln}(1/\alpha )`$.) The expressions for the squares $`|M_{\lambda _1,\lambda _2,\lambda _3;m}|^2`$ contain parts of the form
$$A_i^{}A_j=A_i^{(0)}A_j^{(0)}+\left[A_i^{(0)}A_j^{(1)}+A_i^{(1)}A_j^{(0)}\right]+A_i^{(1)}A_j^{(1)}+\left[A_i^{(0)}A_j^{(2)}+A_i^{(2)}A_j^{(0)}\right]+\mathrm{}$$
(123)
for various combinations of $`i`$ and $`j`$. The $`A_i^{(0)}A_j^{(0)}`$ terms give the lowest-order differential decay distribution. The $`A_i^{(0)}A_j^{(1)}+A_i^{(1)}A_j^{(0)}`$ terms give the order-$`\alpha `$ correction, and the $`A_i^{(1)}A_j^{(1)}`$ and $`A_i^{(0)}A_j^{(2)}+A_i^{(2)}A_j^{(0)}`$ terms give the order-$`\alpha ^2`$ corrections.
The graphs contributing to the order-$`\alpha `$ corrected decay amplitudes $`A_i^{(1)}`$ in the renormalized Feynman gauge are shown in Fig. 2. The infrared divergence induced by mass-shell renormalization is regulated by use of a photon mass $`\lambda `$. The self-energy (Fig. 2a) and vertex graphs (Figs. 2b, 2c) contain infrared divergences of the form $`\mathrm{ln}\lambda M_{LO}`$. The ladder graph Fig. 2e requires special care in its evaluation since it contains an infrared binding singularity. This divergence can be identified and subtracted out, as discussed in detail in Ref. Adkins02 . The result is that
$$M_L=\left\{\frac{\pi }{\lambda }+\mathrm{ln}\lambda 1+O(\lambda )\right\}\left(\frac{\alpha }{\pi }\right)M_{LO}+M_{LS}.$$
(124)
The subtracted ladder graph is
$$M_{LS}=i\alpha ^4m^2\underset{S_3}{}(d\mathrm{})^{\prime \prime }\left[\mathrm{}^2(\mathrm{}^22\mathrm{}p)(\mathrm{}^2+2\mathrm{}p)Z(\mathrm{})\right]^1\left\{\left(\mathrm{tr}(\mathrm{})\mathrm{tr}(0)\right)\frac{\mathrm{tr}(0)}{Z(0)}\left(Z(\mathrm{})Z(0)\right)\right\},$$
(125)
with $`p=mN`$,
$`\mathrm{tr}(\mathrm{})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{tr}[\gamma ^\mu (\gamma (\mathrm{}p)+m)\gamma ϵ_3^{}(\gamma (\mathrm{}p+k_3)+m)\gamma ϵ_2^{}(\gamma (\mathrm{}+pk_1)+m)`$ (127)
$`\times \gamma ϵ_1^{}(\gamma (\mathrm{}+p)+m)\gamma _\mu (\gamma N+1)\gamma ϵ(\gamma N1)],`$
and
$$Z(\mathrm{})=((\mathrm{}p+k_3)^2m^2)((\mathrm{}+pk_1)^2m^2).$$
(128)
The subtraction in Eq.(125) takes away the $`\mathrm{}`$-independent part of $`\mathrm{tr}(\mathrm{})/Z(\mathrm{})`$, which would have had an infrared singularity. This binding singularity, regulated by the photon mass, is displayed in Eq. (124). The contributions of the order-$`\alpha `$ decay graphs were evaluated one by one and summed. The $`1/\lambda `$ binding singularity was removed according to the usual procedure of NRQED Caswell86 ; Adkins02 . The $`\mathrm{ln}\lambda `$ terms cancel between the self-energy, vertex, and ladder graphs. The remaining expressions are a finite sums of rational functions of the $`x_i`$ times logarithms, dilogarithms, and inverse tangent functions.
## VIII Results and Conclusions
We use our analytic results for the order-$`\alpha `$ decay amplitudes $`A_i^{(1)}`$ to calculate the order-$`\alpha `$ correction to the o-Ps $`3\gamma `$ decay rate and a part of the order-$`\alpha ^2`$ correction. The individual amplitudes are quite lengthy and will not be displayed. A simplified form for the complete order-$`\alpha `$ decay rate contribution is given in the Appendix. The result for the order-$`\alpha `$ decay rate is integration\_ref\_1
$$\mathrm{\Gamma }_1=10.286606(10)\frac{\alpha }{\pi }\mathrm{\Gamma }_{LO}.$$
(129)
This represents a 60-fold improvement in precision over the previous best result $`10.2866(6)`$ Adkins92 done using a higher dimensional integration. The two-dimensional integral for the part of the order-$`\alpha ^2`$ correction to the decay rate coming from the $`A_i^{(1)}A_j^{(1)}`$ terms gives integration\_ref\_2
$$\mathrm{\Gamma }_2(\mathrm{square})=28.860(2)\left(\frac{\alpha }{\pi }\right)^2\mathrm{\Gamma }_{LO}.$$
(130)
The previous result for this contribution was $`28.8(2)`$. Burichenko93
In this work we obtained analytic expressions for the o-Ps$`3\gamma `$ decay amplitudes. We used these expressions to obtain precise results for the one-loop and “square” decay rate contributions, which were incorporated into the full calculation of two-loop corrections to the o-Ps$`3\gamma `$ decay rate. Adkins00 ; Adkins02 We also give an explicit form for the one-loop decay distribution (see the Appendix) which can be used to obtain the one-loop energy spectrum in a convenient form.
###### Acknowledgements.
I am grateful for the assistance of Kunal Das in an early stage of this work, and to Zvi Bern, Richard Fell, Russell Kauffman, Andrew Morgan, and Jonathan Sapirstein for useful conversations. I thank Aditya Narayanan, Sharmini Ilankovan, and Aba Mensah-Brown for helping to check formulas. I appreciate the hospitality of the Physics Department at UCLA, where part of this work was done, and acknowledge the support of the National Science Foundation (through grant No. PHY-9408215) and of the Franklin and Marshall College Grants Committee.
*
## Appendix A The one-loop correction
In the appendix we present the integral for the one-loop correction to the decay rate in compact form. From this integral the one-loop phase-space distribution and energy spectrum can be obtained. We note that for very soft photons additional effects must be taken into account in order to obtain accurate results for the phase-space distribution and energy spectrum. Pestieau02 ; Manohar04 ; Voloshin04 ; Ruiz04
The one-loop correction to the decay rate is
$$\mathrm{\Gamma }_1=\frac{m\alpha ^7}{36\pi ^2}_0^1𝑑x_1_{1x_1}^1𝑑x_2\frac{1}{x_1x_2x_3}\left\{F(x_1,x_3)+\mathrm{permutations}\right\},$$
(131)
where $`x_1+x_2+x_3=2`$ and the “permutations” are the six permutations of the variables $`x_1`$, $`x_2`$, $`x_3`$. The one-loop phase space distribution is just the integrand. The energy spectrum is found by integrating over $`x_2`$ but not $`x_1=E_1/m`$. (The corresponding lowest-order expression is given in Eq. (68).) The function $`F(x_1,x_3)`$ is given by
$$F(x_1,x_3)=g_0(x_1,x_3)+\underset{i=1}{\overset{5}{}}g_i(x_1,x_3)h_i(x_1)+\underset{i=6}{\overset{7}{}}g_i(x_1,x_3)h_i(x_1,x_3).$$
(132)
The $`h`$ functions are given by
$`h_1(x_1)`$ $`=`$ $`\mathrm{ln}(2x_1),`$ (133a)
$`h_2(x_1)`$ $`=`$ $`\sqrt{{\displaystyle \frac{x_1}{\overline{x}_1}}}\theta _1,`$ (133b)
$`h_3(x_1)`$ $`=`$ $`{\displaystyle \frac{1}{2x_1}}\left\{\zeta (2)\mathrm{Li}_2(12x_1)\right\},`$ (133c)
$`h_4(x_1)`$ $`=`$ $`{\displaystyle \frac{1}{2x_1}}\left\{\left({\displaystyle \frac{\pi }{2}}\right)^2\theta _1^2\right\},`$ (133d)
$`h_5(x_1)`$ $`=`$ $`{\displaystyle \frac{1}{2\overline{x}_1}}\theta _1^2,`$ (133e)
$`h_6(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{x_1\overline{x}_1x_3\overline{x}_3}}}\left\{\mathrm{Li}_2(r_A^+,\overline{\theta }_1)\mathrm{Li}_2(r_A^{},\overline{\theta }_1)\right\},`$ (133f)
$`h_7(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{x_1\overline{x}_1x_3\overline{x}_3}}}\{\mathrm{Li}_2(r_B^+,\theta _1)\mathrm{Li}_2(r_B^{},\theta _1)`$ (133g)
$`{\displaystyle \frac{1}{2}}\mathrm{Li}_2(r_C^+,0)+{\displaystyle \frac{1}{2}}\mathrm{Li}_2(r_C^{},0)\},`$
where $`\overline{x}_i=1x_i`$ and
$`\theta _1`$ $`=`$ $`\mathrm{arctan}\left(\sqrt{\overline{x}_1/x_1}\right),`$ (134a)
$`\overline{\theta }_1`$ $`=`$ $`\mathrm{arctan}\left(\sqrt{x_1/\overline{x}_1}\right),`$ (134b)
$`r_A^\pm `$ $`=`$ $`\sqrt{\overline{x}_1}\left(1\pm \sqrt{{\displaystyle \frac{x_1\overline{x}_3}{\overline{x}_1x_3}}}\right),`$ (134c)
$`r_B^\pm `$ $`=`$ $`\sqrt{x_1}\left(1\pm \sqrt{{\displaystyle \frac{\overline{x}_1\overline{x}_3}{x_1x_3}}}\right),`$ (134d)
$`r_C^\pm `$ $`=`$ $`r_B^\pm /\sqrt{x_1}.`$ (134e)
The $`g`$ functions are given in terms of $`x^{mn}=x_1^mx_3^n`$ and $`\overline{x}_2=x_1+x_31`$ as
$`g_0(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{9x_1\overline{x}_1(12x_1)x_3\overline{x}_3(12x_3)}}\{180+2196x^{10}4968x^{20}`$ (135d)
$`+5292x^{30}2664x^{40}+504x^{50}5848x^{11}+22639x^{21}20280x^{31}`$
$`+8405x^{41}1240x^{51}24x^{61}17551x^{22}+22982x^{32}5857x^{42}`$
$`+264x^{52}+48x^{62}3776x^{33}878x^{43}+400x^{53}+536x^{44}\},`$
$`g_1(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{4}{x_1^2(12x_1)^2(x_1x_3)x_3}}\{2x^{20}13x^{30}+35x^{40}36x^{50}+8x^{60}`$ (135h)
$`+4x^{70}+9x^{11}59x^{21}+149x^{31}210x^{41}+162x^{51}51x^{61}+x^{71}`$
$`4x^{02}+3x^{12}+55x^{22}126x^{32}+104x^{42}39x^{52}+x^{62}+8x^{03}`$
$`26x^{13}+7x^{23}+22x^{33}+2x^{43}2x^{53}4x^{04}+14x^{14}8x^{24}8x^{34}\},`$
$`g_2(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{2}{3x_1^3\overline{x}_1x_3}}\{48x^{10}+180x^{20}276x^{30}+228x^{40}108x^{50}+24x^{60}`$ (135l)
$`+48x^{01}48x^{11}144x^{31}+244x^{41}106x^{51}+2x^{61}+4x^{71}96x^{02}`$
$`+156x^{12}108x^{22}+168x^{32}132x^{42}+7x^{52}+6x^{62}+48x^{03}60x^{13}`$
$`36x^{23}+42x^{33}+9x^{43}6x^{53}+6x^{34}4x^{44}\},`$
$`g_3(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{4}{x_1^2(x_1x_3)x_3}}\{2x^{20}2x^{40}4x^{60}+5x^{11}6x^{21}+14x^{31}4x^{41}`$ (135o)
$`+18x^{51}x^{61}4x^{02}2x^{12}+4x^{22}2x^{32}26x^{42}x^{52}+8x^{03}`$
$`7x^{13}2x^{23}+12x^{33}+2x^{43}4x^{04}+4x^{14}\},`$
$`g_4(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{8}{x_1^2}}\{4+7x^{10}7x^{20}+12x^{30}10x^{40}+2x^{50}+8x^{01}10x^{11}`$ (135q)
$`+3x^{21}3x^{31}+2x^{41}4x^{02}+3x^{12}+2x^{22}+x^{32}\},`$
$`g_5(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{2\overline{x}_1}{x_1}}\left\{834x^{10}+29x^{20}4x^{30}+6x^{11}+8x^{02}4x^{12}\right\},`$ (135r)
$`g_6(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{x_1\overline{x}_2x_3}}\{1676x^{10}+136x^{20}124x^{30}+64x^{40}16x^{50}60x^{01}`$ (135v)
$`+272x^{11}424x^{21}+294x^{31}104x^{41}+22x^{51}+92x^{02}392x^{12}`$
$`+484x^{22}187x^{32}+13x^{42}+2x^{52}76x^{03}+294x^{13}259x^{23}+30x^{33}`$
$`+3x^{43}+36x^{04}120x^{14}+61x^{24}+3x^{34}8x^{05}+22x^{15}+2x^{25}\},`$
$`g_7(x_1,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{\overline{x}_2}}\{1648x^{10}+46x^{20}12x^{30}2x^{40}48x^{01}+60x^{11}+9x^{21}`$ (135x)
$`31x^{31}+10x^{41}+46x^{02}+9x^{12}42x^{22}+11x^{32}12x^{03}31x^{13}`$
$`+11x^{23}2x^{04}+10x^{14}\}.`$
|
warning/0506/physics0506027.html
|
ar5iv
|
text
|
# Is There a Real-Estate Bubble in the US?
## 1 Is there a real-estate bubble in the US? Lessons from the past UK bubble
In the aftermath of the burst of the “new economy” bubble in 2000, the Federal Reserve aggressively reduced short-term rate yields in less than two years from 6<sup>1/2</sup> % to 1 % in June 2003 in an attempt to coax forth a stronger recovery of the US economy. In March 2003, we released a paper <sup>1</sup><sup>1</sup>1See, W.-X. Zhou and D. Sornette, http://arXiv.org/abs/physics/0303028 published a few months later addressing the growing apprehension at the time (see for instance ) that this loosening of the US monetary policy could lead to a new bubble in real estate, as strong housing demand was being fueled by historically low mortgage rates. As of March 2003, we concluded that, “while there is undoubtedly a strong growth rate, there is no evidence of a super-exponential growth in the latest six years,” giving “ no evidence whatsoever of a bubble in the US real estate market” .
More than two years have passed. During that period, the historically low Fed rate of 1% remained stable from June 2003 to June 2004. Since June 2004, the Fed (specifically, the Federal Open Market Committee (FOMC)) has increased its discount rate by increments of 0.25% at each of its successive meetings (the FOMC holds 8 meetings per year): at the time of writing (end of May 2005), the last 0.25% increase occurred on May 3rd, 2005 to yield a short-term rate of 3%, the next meeting of the FOMC being scheduled on 29/30 June 2005 <sup>2</sup><sup>2</sup>2Federal Open Market Committee, $`http://www.federalreserve.gov/FOMC/\mathrm{\#}calendars`$. While the short-term interest rates are following a steady upward trend of 2% per year since June 2004, long-term rates have not followed, some going down while other long-term rates increasing only slightly. Thus, long-term mortgage interest rates have remained extremely low by historical standard. The Office of Federal Housing Enterprise Oversight (OFHEO), the government unit tasked with regulating Fannie Mae and Freddie Mac <sup>3</sup><sup>3</sup>3Fannie Mae (resp. Freddie Mac) is a stockholder-owned corporation chartered by the Federal Government in 1938 (resp. by Congress in 1970) to keep money flowing to mortgage lenders in support of homeownership and rental housing., recently published a research paper <sup>4</sup><sup>4</sup>4Office of Federal Housing Enterprise Oversight (OFHEO), Mortgage markets and the enterprises (October 2004), prepared by V.L. Smith and L.R. Bowes (http://www.ofheo.gov/Research.asp) stating that “The housing market achieved record levels of activity and contributed significantly to the economic recovery … Falling mortgage rates stimulated housing starts and sales, and many refinancing borrowers took out loans that were larger than those they paid off, providing additional funds for consumption expenditures… According to Freddie Mac, homeowners who refinanced in 2003 converted almost $139 billion in home equity into cash, up from $105 billion in 2002.” This has led to renewed worries that a real-estate bubble is on its way.
The purpose of this paper is to revisit this question, using the more than two additional years of data since our previous analysis . As explained in our previous paper on real-estate bubbles , our analysis relies on a general theory of financial crashes and of stock market instabilities developed in a series of works (see and references therein). The main ingredient of the theory is the existence of positive feedbacks in stock markets as well as in the economy. Positive feedbacks, i.e., self-reinforcement, refer to the fact that, conditioned on the observation that the market has recently moved up (respectively down), this makes it more probable to keep it moving up (respectively down), so that a large cumulative move may ensue. The concept of “positive feedbacks” has a long history in economics. It can occur for instance in the form of “increasing returns”– which says that goods become cheaper the more of them are produced (and the closely related idea that some products, like fax machines, become more useful the more people use them). Positive feedbacks, when unchecked, can produce runaways until the deviation from equilibrium is so large that other effects can be abruptly triggered and lead to rupture or crashes. Alternatively, it can give prolonged depressive bearish markets. There are many mechanisms leading to positive feedbacks including investors’ over-confidence, imitative behavior and herding between investors, refinancing releasing new cash re-invested in houses, lower requirement margins due to uprising prices, and so on. Such positive feedbacks provide the fuel for the development of speculative bubbles, by the mechanism of cooperativity, that is, the interactions and imitation between investors may lead to collective behaviors similar to crowd phenomena. Different types of collective regimes are separated by so-called critical points which, in physics, are widely considered to be one of the most interesting properties of complex systems. A system goes critical when local influences propagate over long distances and the average state of the system becomes exquisitely sensitive to a small perturbation, i.e. different parts of the system become highly correlated. Another characteristic is that critical systems are self-similar across scales: at the critical point, an ocean of traders who are mostly bearish may have within it several continents of traders who are mostly bullish, each of which in turns surrounds seas of bearish traders with islands of bullish traders; the progression continues all the way down to the smallest possible scale: a single trader . Intuitively speaking, critical self-similarity is why local imitation cascades through the scales into global coordination. Critical points are described in mathematical parlance as singularities associated with bifurcation and catastrophe theory. At critical points, scale invariance holds and its signature is the power law behavior of observables.
Mathematically, these ideas are captured by the power law
$$\mathrm{ln}[p(t)]=A+B(t_ct)^m,$$
(1)
where $`p(t)`$ is the house price or index, $`t_c`$ is an estimate of the end of a bubble so that $`t<t_c`$ and $`A,B,m`$ are coefficients. If the exponent $`m`$ is negative, $`\mathrm{ln}[p(t)]`$ is singular when $`tt_c^{}`$ and $`B>0`$ ensuring that $`\mathrm{ln}[p(t)]`$ increases. If $`0<m<1`$, $`\mathrm{ln}[p(t)]`$ is finite but its first derivative $`d\mathrm{ln}[p(t)]/dt`$ is singular at $`t_c`$ and $`B<0`$ ensuring that $`\mathrm{ln}[p(t)]`$ increases. Extension of this power law (1) takes the form of log-periodic power law (LPPL) for the logarithm of the price
$$\mathrm{ln}[p(t)]=A+B(t_ct)^m+C(t_ct)^m\mathrm{cos}\left[\omega \mathrm{log}(t_ct)\varphi \right],$$
(2)
where $`\varphi `$ is a phase constant and $`\omega `$ is the angular log-frequency. This first version (2) amounts to assume that the potential correction or crash at the end of the bubble is proportional to the total price . In contrast, a second version assumes that the potential correction or crash at the end of the bubble is proportional to the bubble part of the total price, that is to the total price minus the fundamental price . This gives the following price evolution:
$$p(t)=A+B(t_ct)^m+C(t_ct)^m\mathrm{cos}\left[\omega \mathrm{log}(t_ct)\varphi \right].$$
(3)
As explained in , we diagnose a bubble using these models by demonstrating a faster-than-exponential increase of $`p(t)`$, possibly decorated by log-periodic oscillations.
Before presenting the result of our analysis using these models on the US real-estate bubble, it is appropriate to discuss how our detection of a bubble in the UK real-estate market fared since March 2003. In , we reported “unmistakable signatures (log-periodicity and power law super-exponential acceleration) of a strong unsustainable bubble” for the UK real-estate market. We identified two potential turning points in the UK bubble reported in Tables 2 and 3 of : end of 2003 and mid-2004. The former (resp. later) was based on the use of formula (2) (resp. (3)). These predictions were performed in Feb. 2003 (again our paper was released in early March 2003 on an electronic archive <sup>5</sup><sup>5</sup>5W.-X. Zhou and D. Sornette, http://arXiv.org/abs/physics/0303028). We stress that these turning points can be either crashes or changes of regimes according to the theory coupling rational expectation bubbles with collective herding behavior described in . In other words, the theory describes bubbles and their end but not the crash itself: the end of a bubble is the most probable time for a crash, but a crash can occur earlier (with low probability) or not at all; the possibility that no crash occurs is necessary for the bubble to exist, otherwise, rational investors would anticipate the crash and, by backward reasoning, would make it impossible to develop.
Figure 1 plots the UK Halifax house price index (HPI) <sup>6</sup><sup>6</sup>6The Halifax house price index has been used extensively by government departments, the media and businesses as an authoritative indicator of house price movements in the United Kingdom. This index is based on the largest sample of housing data and provides the longest unbroken series of any similar UK index. The monthly house price index data are retrieved from the web site of HBOS http://www.hbosplc.com/view/housepriceindex/housepriceindex.asp.. The six time series are the following. AllMon: All Houses (All Buyers); AllMonSA: All Houses (All Buyers) (seasonally adjusted); Existing: Existing Houses (All Buyers); New: New Houses (All Buyers); FOO: Former Owner Occupiers (All Houses); FTB: First Time Buyers (All Houses). from 1993 to April 2005 (the latest available quote at the time of writing). The two groups of vertical lines correspond to the two predicted turning points mentioned above. The first set of predicted turning points (dashed lines in Figure 1) anticipated by half-a-year the turning point which occurred mid-2004 as predicted by the second set.
Our analysis presented below uses three data sets: (1) the regional data (Northeast, Midwest, west, south and USA as a whole) of the quarterly average sale prices of new houses up to the fourth quarter of 2004 (the latest data available); (2) the house price index of individual states (50 states $`+`$ DC), up to Q1 of 2005, quarterly data; and (3) daily data of the S&P 500 Home Index, up to May 6, 2005. We first present in section 2 a broad-brush analysis using the exponential versus power law models of house price appreciation for the whole continental US and then by regions. We then turn to a state-by-state analysis which leads to a partition into three classes: (i) non-bubbling states, (ii) recent-bubbling states and (iii) clearly-bubbling states. For the states for which a bubble seems to be clearly established according to our criterion, we provide a first estimation of the critical time of the end of the bubble. We then turn to the more elaborate LPPL models (2) and (3) and a nonlinear extension, using the daily data of the S&P 500 Home Index up to May 6, 2005.
## 2 Evidence of a US real-estate bubble by faster-than-exponential growth
Figures 2 show the quarterly average sale prices of new houses sold in all the states of the USA as well as in the four main regions, Northeast, Midwest, South and West, from 1993 to the fourth quarter of 2004 as a function of time $`t`$. The smooth curve is the power-law fit (1) to the data. Except for the midwest and south regions, one can observe a strong upward curvature in these linear-logarithmic plots, which characterize a faster-than-exponential price growth (recall that an exponential growth would qualify as a straight line in such linear-logarithmic plots). The existence of a strong upward curvature characterizing a faster-than-exponential growth is quantified by the relatively small values of the exponent $`m`$ ($`=0.55`$ for all states, $`=0.64`$ for the Northeast region, $`=0.18`$ for the West region).
To have a closer look, we examined quarterly data of House Price Index (HPI) for each individual state. Rather than following a formal procedure and developing sophisticated statistical tests, the obvious differences between the price trajectories in the different states led us to prefer a more intuitive approach consisting of classifying the different states according to how strongly they depart from a steady exponential growth. We found three families, shown in figures 3, 4 and 5. Figure 3 shows the quarterly HPI in the 21 states which have an approximately constant exponential growth, qualified by a linear trend in a linear-logarithmic scale. The thick straight line at the bottom of the figure is the average over all 21 states corresponding to an annual growth rate of 4.6% over the last 13 years (we did not use the data prior to 1992 to avoid contamination by the turning point of the previous bubble in 1991). Figure 4 shows the quarterly HPI in the 8 states exhibiting a recent upward acceleration following an approximately constant exponential growth rate. Figure 5 shows the quarterly HPI in the 22 states exhibiting a clear upward faster-than-exponential growth. These 22 states thus exhibit the hallmark of a real-estate bubble. Figure 6 provides a geographical synopsis of this classification in three families: the first family of figure 3 is green, the second family of figure 4 is magenta, and the third family of figure 5 is red. As often discussed by commentators, prices have accelerated mostly to the Northeast and West regions, which is consistent with Fig. 2.
Consider the third family of 22 states where we diagnose a bubble, as shown in figure 5. Can a power law fit with (1) reveal the end of the bubble? Such turning point is in principle measured by the time $`t_c`$ in expression (1), which gives the time at which the bubble should end. In order to get less noisy data, we averaged over the 22 price trajectories of figure 5 and then fitted the obtained average with (1) (with the modification that $`t_ct`$ is replaced by $`|t_ct|`$ to allow for a more robust estimation) over a time interval from $`t_{\mathrm{start}}`$ to the last available data point (2005Q1). Figure 7 shows the obtained critical time $`t_c`$ and exponent $`m`$ as a function of $`t_{\mathrm{start}}`$. Varying $`t_{\mathrm{start}}`$ allows us to test for sensitivity with respect to the different time periods and assess the robustness of the results. Not surprisingly, we find that the fitted $`t_c`$ is close to the last data points for some $`t_{\mathrm{start}}`$, a result which has been found to systematically characterize power law behaviors . Therefore, the power law fit is not very reliable to determine the end of the real-estate bubble. However, the relative stability of $`m`$ in the range $`0.10.5`$ as a function of $`t_{\mathrm{start}}`$, which characterizes a faster-than-exponential growth, indicates that the simple power law (1) is already a good model.
## 3 Extension to the LPPL model and discussion
The previous tests performed in (and references therein) show that the problem with the too-large sensitivity of $`t_c`$ in the simple power law model with respect to the few last data points are alleviated by using the more sophisticated LPPL models (2) and (3). Here, we use both LPPL models (2) and (3) as well as the so-called 2nd-order Landau LPPL introduced in <sup>7</sup><sup>7</sup>7See also $`http://www.ess.ucla.edu/faculty/sornette/prediction/index.asp\mathrm{\#}prediction`$ for a recent application to the US stock market. In a nutshell, the 2nd-order Landau LPPL extends the LPPL model by allowing for a first nonlinear correction which amounts to combining two log-frequencies $`\omega _0`$ (close to $`t_c`$) and $`\omega _{\mathrm{}}`$ (far from $`t_c`$).
We fit the daily data of the S&P 500 Home Index to the LPPL and 2nd-order Landau LPPL models in a time interval from $`t_{\mathrm{start}}`$ to the last available data point (April 2005). Figure 8 presents the dependence of $`\omega `$ for the LPPL model and of $`\omega _0`$ for the 2nd-order Landau LPPL model as a function of $`t_{\mathrm{start}}`$. We find $`\omega 10`$ and $`\omega _0`$ oscillating between $`\omega 10`$ and $`\omega /2`$, as we expect from the generic existence of harmonics (see our previous extended discussions in ). The stability of $`\omega _0`$ and the compatibility between the two descriptions is a signature of the robustness of the signal. Figure 9 shows the predicted critical time $`t_c`$ as a function of $`t_{\mathrm{start}}`$ obtained from the fits with the LPPL and the 2nd-order Landau LPPL models. The large spreads of values for $`t_{\mathrm{start}}`$ earlier than 1993 reflects the fact that the bubble has really started only after 1993.
We observe a good stability of the predicted $`t_c`$ mid-2006 for the two LPPL models (2) and (3). The spread of $`t_c`$ is larger for the second-order LPPL fits but brackets mid-2006. As mentioned before, the power law fits are not reliable. We conclude that the turning point of the bubble will probably occur around mid-2006.
|
warning/0506/astro-ph0506433.html
|
ar5iv
|
text
|
# Imaging Bright Spots in the Accretion Flow Near the Black Hole Horizon of Sgr A*
## 1 Introduction
Despite its successes, there has not yet been a satisfactory test of general relativity in the strong gravity limit. By their very nature, studies of black holes are likely to provide the best opportunity for constraining strong field relativity. Unfortunately, current attempts to do this rely upon modelling the accretion flow and are thus indirect (see, e.g., Narayan & Heyl, 2002; Tanaka et al., 1995; Reynolds & Nowak, 2003; Pariev et al., 2001).
Of these, the most direct method involves observations of line-like features in the X-ray spectra of black hole candidates, typically interpreted as the Fe K$`\alpha `$ fluorescence line, broadened as a result of the Keplerian motion of the disk and frame dragging due to the rotation of the black hole (see, e.g., Reynolds & Nowak, 2003; Pariev et al., 2001). The lack of a correlation between the variability in the soft X-ray continuum and the line emission implies that the simplest model used for the interpretation of the observations is incomplete (see, e.g., Weaver et al., 2001; Wang et al., 2001, 1999; Chiang et al., 2000; Lee et al., 2000), though attempts to rectify this situation with the inclusion of gravitational lensing have been made (Matt et al., 1997). The case is further complicated by the existence of alternative interpretations (see, e.g., Elvis, 2000; You et al., 2003).
General relativistic effects can also play a substantial role in the polarisation properties of Thomson thick disks (see, e.g. Connors et al., 1980; Laor et al., 1990; Bao et al., 1997). Therefore, detailed spectropolarimetric observations may shed light on both the physics of the accretion flow and the curvature of space-time. However, these necessarily require an accretion model and consequently suffer from considerable uncertainties associated with the accretion physics.
In contrast, it may be possible to directly image a black hole, and thus measure the space-time curvature down to the photon orbit, the radius at which photons execute a circular orbit, $`3GM/c^2`$ for a Schwarzschild black hole and down to $`GM/c^2`$ for a maximally rotating Kerr black hole<sup>1</sup><sup>1</sup>1Note that the coincidence between the photon orbit and the horizon in the maximally rotating Kerr case is an artifact of Boyer-Lindquist coordinates. At all values of black hole spin, a non-vanishing radial proper distance exists between the horizon and the photon orbit (see, e.g., Chandrasekhar, 1992). (Falcke et al., 2000). The black hole in the Galactic centre provides the best candidate for such an observation as it possesses the largest apparent size on the sky, with $`GM/c^2`$ corresponding to an angular scale of $`5\mu \mathrm{as}`$. Within the next decade it is expected that a Very Large Baseline Array (VLBA) at sub-millimetre wavelengths (where scattering is no longer the limiting factor) will exist which is expected to provide $`20\mu \mathrm{as}`$ resolution (Greenhill, 2005; Miyoshi et al., 2004).
Earlier theoretical work has focused on optically thin, azimuthally symmetric accretion flows, which are generically expected to show a shadow around the black hole (e.g., Falcke et al., 2000; Takahashi, 2005, 2004; Beckwith & Done, 2005). However, numerical general-relativistic magnetohydrodynamic simulations of accretion flows suggest that the region near the innermost stable circular orbit (ISCO) may be strongly inhomogeneous (see, e.g., De Villiers et al., 2003). Indeed, recent observations of Sgr A\* in the infrared and X-ray bands revealed flaring activity on short time scales (Ghez et al., 2004; Eckart et al., 2004; Genzel et al., 2003; Porquet et al., 2003; Aschenbach et al., 2004; Goldwurm et al., 2003; Baganoff et al., 2001) indicating strong inhomogeneities in the emission close to the black hole horizon. Computations of the light curves of inhomogeneous accretion flows have been performed in the context of quasars and X-ray binaries more generally (Bao et al., 1998, 1997), but without providing images to be compared with future observations and without studying the dependence of the light curve on the black hole parameters.
In this paper, the unpolarised and polarised light curves and images are computed for an optically thick emitting sphere at a number of radii, viewing angles, and black hole spins. While this may strictly apply for a star orbiting a supermassive black hole, it should be understood that our treatment is a proxy for any mechanism which enhances the emission in a compact region of space (e.g., due to reconnection events, over densities in mass or magnetic field strength).
The method by which rays are traced, the radiative transfer performed, and the optically thick sphere is modelled are discussed in section 2. Section 3 compares the expected light curves (both polarised and unpolarised) for a number of orbits, viewing angles, and black holes spins. Images and centroid motions are presented in sections 4 and 5, respectively. Finally, concluding remarks are summarised in section 6.
In what follows, the metric signature is taken to be $`(+++)`$, and geometrised unites are used ($`G=c=1`$).
## 2 Computational Methods
Due to the complexities introduced by the Kerr metric, in the strong field limit, images of the accretion flow are most readily obtained via numerical methods. This computation may be succinctly segregated into the problems of tracing null geodesics, performing the radiative transfer along these rays, and modelling the inhomogeneity (in this case an optically thick sphere). The first two of these are treated in subsection 2.1, and the third in subsection 2.2.
### 2.1 Ray Tracing & Radiative Transfer
The light rays are traced in curved space following the scheme of Broderick & Blandford (2003) (for which tracing null geodesics is a limiting case), as well as the radiative transfer method of Broderick & Blandford (2004). Below, only the essential aspects are summarised.
Null geodesics are constructed by integrating the equations
$`{\displaystyle \frac{dx^\mu }{d\lambda }}`$ $`=f(r)k^\mu `$
$`{\displaystyle \frac{dk_\mu }{d\lambda }}`$ $`=f(r)\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{k^\nu k_\nu }{x^\mu }}\right)_{k_\alpha },`$ (1)
where the partial differentiation is taken holding the covariant components of the wave four-vector, $`k_\mu `$, constant, and
$$f(r)=r^2\sqrt{1\frac{r_h}{r}},$$
(2)
(where $`r_h`$ is the horizon radius) simply reparameterises the affine parameter, $`\lambda `$, along the ray to avoid singular behaviour near the horizon. An explicit demonstration that $`x^\mu (\lambda )`$ reproduce the null geodesics can be found in Broderick & Blandford (2003).
Typically, polarised radiative transfer is performed using the Stokes parameters. However, since these are not Lorentz scalars, in relativistic environments they are not ideal. In contrast, the photon distribution function ($`I_\nu /\nu ^3`$) is a Lorentz scalar, and thus much simpler to evolve along the null geodesics. As shown in Broderick & Blandford (2004), it is possible to define analogues of the photon distribution function for the rest of the Stokes parameters using an orthonormal tetrad propagated along the null geodesics.
While the numerical scheme utilised here is not optimized for this particular problem, it has the virtue of already being implemented and tested in the context of imaging accreting black holes.
### 2.2 Optically Thick Sphere
An inhomogeneity in the accretion flow is modelled as an optically thick sphere orbiting around the central black hole. The surface of the inhomogeneity is given implicitly by
$$\mathrm{\Delta }r^\mu \mathrm{\Delta }r_\mu +\left(u_S^\mu \mathrm{\Delta }r_\mu \right)^2=R_S^2$$
(3)
where $`\mathrm{\Delta }r^\mu r^\mu r_S^\mu `$ is the displacement from the sphere centre, located at $`r_S^\mu `$, moving with velocity $`u_S^\mu `$, and $`R_S`$ is the radius of the sphere. The second term accounts for the length contraction<sup>2</sup><sup>2</sup>2Despite the length contraction, it is long been known that in the far field, relativistic aberration combined with time of flight effects result in the relativistically moving sphere appearing rotated but otherwise unchanged (Penrose, 1958; Terrell, 1959). This may indeed be seen explicitly in the images presented in section 4.. In the limit that $`\mathrm{\Delta }r^\mu \mathrm{\Delta }r_\mu r`$ this definition indeed produces a sphere in the comoving frame. For $`\mathrm{\Delta }r^\mu \mathrm{\Delta }r_\mu r`$ the surface defined by equation (3) necessarily departs from sphericity (since the $`\mathrm{\Delta }r^\mu \mathrm{\Delta }r_\mu `$ is only the differential line element), and thus in reality is only quasi-spherical. However, considering the level of approximation inherent in treating the inhomogeneity as a sphere in the first place, the added simplicity outweighs a more detailed effort to produce a true sphere. Here a value of $`R_S=1.5M`$ is used, which is sufficiently compact to reveal the strong field effects of interest in this study.
Since the sphere is an extended body, it cannot have a uniform velocity if it is to remain coherent. Here all points upon the surface of the sphere are taken to have the same angular velocity as measured by the Boyer-Lindquist observer, i.e., the same $`u_S^\varphi /u_S^t`$, ensuring that all parts of the sphere execute an orbit in the same amount of observer time. An explicit proof that this ensures that the sphere does not shear can be found in Broderick & Blandford (2004). The velocity of the centre of the sphere was taken to be that associated with the stable circular orbit at $`r_S^\mu `$.
Images are produced by tracing a bundle of parallel rays back from a screen far from the black hole. Rays that impinged upon the sphere were given an intensity (subsequently transformed into a photon distribution function) governed by the Rayleigh-Jeans law,
$$I_\nu =2\frac{\nu ^2}{c^2}kT,$$
(4)
where $`\nu =u_\mu k^\mu /2\pi `$ and the temperature was assumed to be proportional to its virial value:
$$T\left(u_t\frac{1}{\sqrt{g^{tt}}}\right),$$
(5)
where $`g^{\mu \nu }`$ are the contravariant components of the metric.
Polarisation measurements place an additional observational constraint upon the emission and propagation physics. Because the polarisation is parallel–propagated along the ray, it is especially sensitive to the effects of strong gravity. However, any discussion of polarisation must be prefaced with the caveat that there will be considerable uncertainty in the emitted polarisation, e.g., due to the emission process (synchrotron or Compton scattering), and the geometry of the emitting region (tangled magnetic fields, thick accretion disks, etc). Because of this large uncertainty, a simple fiducial polarisation model is adopted. For the purpose of investigating the generic effects of strong lensing upon the polarisation, the emitted polarisation fraction is set to be constant and polarised orthogonally to the spin axis of the hole (also the orbital axis). Such a geometry might be expected, e.g., when the primary emission mechanism is synchrotron, and the magnetic field is vertically aligned. The polarisation may be substantially reduced in the presence of realistic field geometries, alternate polarising mechanisms, or radiative transfer effects such as Faraday rotation or depolarisation. Nonetheless, this provides some insight into the possible complexity that may arise in the polarised spectrum.
## 3 Light Curves
Although this paper is primarily concerned with resolved images, considerable information can be extracted from the unresolved unpolarised and polarised light curves. Since obtaining light curves doees not require imaging capabilities, and thus is likely to be technically easier, they are discussed first.
### 3.1 Unpolarised Flux
In general, the shape of the light curve in its entirety is required to make statements regarding the parameters of the black hole and the orbit. Nevertheless, some of these parameters leave their strongest signatures on particular portions of the light curves.
Any gravitational lensing transient is characterised by a time scale and a magnification. For a compact bright spot on a circular orbit about a black hole, the former is set by the period of the orbit,
$$P=2\pi \left(r_S^{3/2}+a\right),$$
(6)
where $`r_S`$ is the orbital radius of the spot and a positive Kerr spin parameter, $`a`$, corresponds to prograde orbits. For a given $`a`$, it is straightforward to determine the radius of the orbit. This is readily apparent in Figure 1 in which the magnification (i.e., the integrated flux normalised by its time-averaged value) is plotted as a function of time for orbital radii of $`6M`$, $`8M`$, and $`10M`$ around a Schwarzschild black hole.
While the magnification does vary with orbital radius, it is a stronger indicator of the inclination of the orbit relative to the line of sight (not to be confused with orbits lying out of the equatorial plane of a Kerr black hole). This is shown in Figure 2 in which the magnification is shown for viewing angles ranging from edge on ($`0^{}`$) to nearly face on ($`89^{}`$). In addition to the strong magnification for orbits which pass directly behind the black hole, there is a second feature near $`t30M`$ resulting from those null geodesics which make a complete orbit before escaping to infinity. The significance of these in the context of the Fe K$`\alpha `$ lines has been previously noted by Beckwith & Done (2005). Here it provides a second signature of the strong bending of light, though only for orbits that are viewed nearly edge on.
The dependence of the magnification light curve upon $`a`$ is presented in Figures 3 and 4. In Figure 3, different spin parameters ($`a=0`$, $`0.5`$, $`0.998`$) are compared at a fixed radius in Boyer-Lindquist coordinates, namely $`6M`$. In this case, higher black hole spins tend to demagnify the source slightly. However, this may be of secondary significance when compared to the fact that for higher $`a`$, stable circular orbits extend closer to the horizon. Indeed, as seen in Figure 4, placing the orbits at the ISCO results in a net increase in the maximum magnification (and a substantially reduced orbital period!).
For rapidly rotating black holes one would expect a difference in the lensing signature of prograde and retrograde orbits. This will be mitigated somewhat by the fact that the retrograde ISCO moves out considerably, reaching $`9M`$ for the maximally rotating case. Nonetheless, as seen in Figure 5, there are quantitative differences in the magnification similar to those seen when the spin is varied (cf. Figure 3) as expected.
### 3.2 Polarised Flux
In addition to the magnification, Figures 14 also show the fractional change in the emitted polarisation and the polarisation angle for the fiducial polarisation model described in section 2.2 (and subject to the caveats presented at the end of that section). In general, the polarisation fraction and polarisation-angle light curves show considerable structure. The variability results primarily from special relativistic aberration associated with the rapid orbital motion, coupled with the choice of the emitted polarisation. This effect is substantially amplified by both gravitational lensing (which allows many viewing angles to be sampled) and general relativistic transfer effects (which rotate the polarisation direction according to the rules of parallel propagation).
The two primary features in the polarisation light curves are the decrease in polarisation fraction (referred to in what follows as the primary minimum) and rotation of the polarisation angle immediately preceding the maximum magnification. Both phenomena result from the development of an Einstein ring/arc, and are thus expected to be generic features. For the configuration considered here, these are a strong function of radius (see, e.g., Figure 1) and viewing angle (Figure 2). The dependence upon viewing angle reaches a maximum when a substantial portion of the rays are incident along the orbital axis (note that due to gravitational lensing this does not occur when the orbit is face on).
A third notable feature is the location of a second minimum due to the development of a secondary ring/arc resulting from rays which complete an orbit before escaping to infinity. This second minimum is considerably shallower than (and lags substantially behind) the primary minimum and should not be confused with the substructure within the primary minimum. The decrease in the polarisation fraction results for reasons similar to those that produce the primary minimum, and is similarly expected to be a generic feature. It is present in Figures 1 and 2, where it leads the primary minimum by $`50M`$. However, of greater interest is the strong dependence of this feature upon the structure of the space-time near the photon orbit, and thus its sensitivity to strong-field general-relativity. This is explicitly demonstrated in Figures 3 and 5, in which the lag between the primary and secondary minima is a strong function of black hole spin.
As the orbit becomes more relativistic, a permanent partial ring develops, qualitatively changing the polarisation features. This can be seen in Figure 4. However, it can be more clearly seen in the images, \[e.g., Figure 6, panel (e)\], and thus is discussed in the following section.
## 4 Images
In addition to measuring the light curves, future sub-millimetre observations promise to image the inner regions of the Galactic centre at $`20\mu \mathrm{as}`$ resolution (Greenhill, 2005; Miyoshi et al., 2004). Illustrative images of the orbits are presented in this section. These may be compared with calculations which assume a uniform and steady accretion flow (e.g., Falcke et al., 2000; Takahashi, 2004, 2005).
Figure 6 shows a sequence of six images distributed roughly evenly throughout an orbit for three of the cases discussed in the previous section. In particular the upper panels show an orbit around a Schwarzschild black hole at $`r_S=6M`$ viewed edge on, the middle panels show the same orbit viewed from $`45^{}`$ above the orbital plane, and the bottom panels show the prograde ISCO of a Kerr black hole with $`a=0.9`$ viewed from $`45^{}`$ above the orbital plane. Panels on the left show the computed resolution, while panels on the right are smoothed by a Gaussian filter with full width $`4M`$ (corresponding to $`20\mu \mathrm{as}`$ for Sgr A\*) and thus are comparable to those that may be observed within the next decade. Each sequence of images is separately normalised.
While all three cases appear qualitatively different, they share three generic features: brightening due to relativistic beaming, a primary Einstein ring/arc, and a secondary Einstein ring/arc due to photons which complete at least one orbit around the black hole (also complete only on the top panel). The first two features are responsible for the magnification and primary minimum in the unpolarised and polarised light curves, respectively. The third is largely responsible for the secondary minimum in the polarisation fraction. Since the secondary ring/arc is formed by photons which pass near the photon orbit ($`3M`$ for Schwarzschild, $`M`$ for maximally rotating Kerr; but see the footnote in section 1), it provides a sensitive test of strong field relativity. As a direct result, the position of the secondary minimum is a diagnostic of the black hole spin, as found in section 3.
Despite the commonalities, there are considerable differences between the three cases shown. The most obvious is that the primary and secondary Einstein rings are incomplete for orbits viewed significantly out of the orbital plane. This results in reduced magnification and polarisation deviations in the lightcurves. In the high spin case, a permanent incomplete ring develops resulting in a qualitative change in the polarised light curves.
As seen in the right-hand panels, these features remain in the smoothed images. Therefore, imaging may be a diagnostic of the orbital and black hole parameters.
The use of the polarisation as a diagnostic of the space-time structure is illustrated in Figure 7, in which images of the polarisation angles are overlayed upon the intensity maps. For the purpose of comparison, these are shown when the emitting sphere is $`180^{}`$ out of phase with the maximum magnification (and thus typical of the unlensed polarisation, top) and for a typical point within the primary polarisation minimum (bottom). As seen in the bottom panels, the development of a primary Einstein ring/arc generically rotates the polarisation angle and reduces the integrated polarisation flux. The development of a secondary ring/arc yields similar results, and hence polarisation is sensitive to those image features which probe strong gravity most effectively. Therefore, polarisation can be use to probe the spacetime at the photon orbit, as long as its emission and transfer are understood.
The images in Figures 6 and 7 are instantaneous snapshots, and are only observable through exposures that are considerably shorter than the orbital period. Typically, observations will average these images over the instrument integration time. To reflect more realistic circumstances, Figures 8-10 show the intensity and polarisation maps averaged over a full orbit (for the same orbits shown in Figures 6 and 7). These are produced by summing the images used to compute the light curves in section 3. All of these have similar structures: a ring/arc associated with the direct image of the sphere, and a second ring/arc associated with the secondary Einstein ring/arc (which closely follows the photon capture impact parameter).
The remarkable symmetry is an artifact of the emission scheme used; for thermal emission in the Rayleigh-Jeans limit, the special relativistic red-shift is precisely offset by the apparent motion on the sky<sup>3</sup><sup>3</sup>3This can be explicitly seen in the images. Indeed the left side of the images have a lower resolution, resulting from the poorer quality of the averages than the right side. This is a direct result of the emitting sphere spending less time on the left (where it is brighter) than on the right (where it is brighter).. This can be explicitly demonstrated for a Schwarzschild black hole by considering a pair of spots moving at a given angle $`\theta `$ relative to the line of sight with velocities $`\beta `$ and $`\beta `$ and time averaged intensities $`I_\nu ^+`$ and $`I_\nu ^{}`$, respectively. For an emitted spectrum with spectral slope $`\alpha `$ (i.e., $`I_\nu \nu ^\alpha `$) the time averaged intensities are related by
$$\frac{I_\nu ^+}{I_\nu ^{}}=\frac{I_\nu ^+}{I_\nu ^{}}\frac{v_a^{}}{v_a^+}=\left(\frac{1+\beta \mathrm{cos}\theta }{1\beta \mathrm{cos}\theta }\right)^{\alpha +2},$$
(7)
where $`v_a^\pm =\beta \mathrm{sin}\theta /(1\beta \mathrm{cos}\theta )`$ is the apparent velocity on the sky. Thus, in the Rayleigh-Jeans approximation, $`I_\nu ^+=I_\nu ^{}`$ and the images are indeed expected to be symmetric. However, for $`\alpha >2`$, as is the case for optically thin synchrotron emission, equation (7) implies a substantial brightness asymmetry. This can be seen explicitly in Figure 11, which is identical to Figure 9 except that $`\alpha `$ was taken to be $`1.3`$ (as suggested by infrared and X-ray flare observations, Eckart et al., 2004). However, despite the considerable asymmetry, the morphology of the image remains unchanged. Therefore, measuring the brightness asymmetry provides a method to probe the spectrum of the bright spot. Note that in both cases the symmetry is clearly broken in the polarisation map, and thus polarisation data may be used to infer the direction of the source motion.
The images do not show a clear black hole shadow. This is not a result of approximating an inhomogeneity by an optically thick sphere, but rather due to emission geometry. In particular, the black hole is not everywhere back lit, and hence at positions on the sky where the primary emitting region lies in front of the black hole no shadow is present. Thus, even when the emission is optically thin, it is expected that the image of inhomogeneous accretion flows will qualitatively differ from that produced by a quasi-spherical accretion flow.
It is significant that the average images of the $`a=0`$ and $`a=0.9`$ case differ substantially (cf. Figures 9 and 10). In particular, the relative positions of the orbital and secondary ring are shifted. Therefore, the spin of the black hole leads to a relative change in position, which is considerably simpler to measure through sub-millimetre interferometry.
## 5 Centroids
It should be possible to measure the image centroid with greater precision than the resolution of an image. As seen in Figure 6, the motion of the centroid contains a considerable amount of the information in the image, and thus the path of the centroid provides a diagnostic of the orbital and black hole parameters.
Figures 12-15 compare the nearly elliptical centroid paths for various orbital and black hole parameters. The major-axis of the elliptical path is indicative of the orbital radius in a similar way as it is for Newtonian orbits (Figure 12). In general, gravitational lensing will increase the minor-axis. Nevertheless, the minor-axis may be used to infer the orbital inclination (Figure 13). In contrast with variations of the orbital radii, gravitational lensing results in changes that are substantially different from the Newtonian case. Nonetheless, these results suggest that it is indeed possible to place significant constraints upon the orbital parameters from the centroid motion alone.
Of more interest is the possibility of measuring the black hole spin. Given the orbital radius and its period, equation (6) may be used to deduce the black hole spin. Figures 14 and 15 show that this may be constrained by the centroid paths as well. However, except for a feature present only near the maximally rotating case, the paths are very similar, and thus it is likely that imaging will be required to distinguish the different cases. Nevertheless, it is not clear that this is the appropriate comparison because the prograde ISCO moves substantially inwards with increasing black hole spin. Therefore, it may be more suitable to compare different black hole spins at their respective prograde ISCOs, as shown in Figure 15. In this case, there is a morphological change in the centroid paths, in addition to the difference that results trivially from the decrease in orbital radius. As a result, when combined with period measurements, an accurate determination of the centroid path provides a method with which to measure $`a`$.
As suggested by Figure 11, for emission models with spectral index $`>2`$ the centroid position will be dominated by the portion of the orbit in which the bright spot is moving towards the observer. This is shown in Figure 16, in which the centroid paths of the Rayliegh-Jeans and $`\alpha =1.3`$ emission models are compared. Thus, estimates of the orbital parameters from the motion of the centroid may require additional information about the spectrum of the emission. However, as mentioned before, even in the absence of spectroscopy, the spectral index may be obtained from brightness asymmetry in the the time-averaged images (see, e.g., equation 7).
## 6 Conclusions
The orbit-integrated image of a bright sphere (hot spot) in motion near the horizon of Sgr A\* is qualitatively different from the image produced by a quasi-spherical accretion flow. In particular, the time-averaged image of a compact hot spot does not reveal a clear black hole shadow, as found for spherical accretion flows (Falcke et al., 2000).
The evidence for short-term flaring activity in the Galactic centre (Ghez et al., 2004; Eckart et al., 2004; Genzel et al., 2003; Porquet et al., 2003; Aschenbach et al., 2004; Goldwurm et al., 2003; Baganoff et al., 2001), implies that the accretion flow around the horizon of Sgr A\* is clumpy or unsteady. In the likely event that the accretion flow has bright spots, imaging of these spots could be utilised as a method for inferring the black hole parameters. In particular, the signature of the secondary Einstein ring/arc in the light curves and the images could potentially be more sensitive to strong gravity effects than imaging the black hole shadow alone. In the time averaged images, black hole spin produces relative offsets between features (see, Figure 10), which are amenable to interferometric imaging in a way that measurements of the shadow are not.
Imaging a bright spot may also independently constrain the mass and distance of Sgr A\*. Timing data combined with the projection of the orbit on the sky (or the apparent velocity) allows the measurement of the ratio of Sgr A\*’s mass and distance. If the line-of-sight velocity of the inhomogeneity can be determined as well (e.g., from spectral measurements combined with the magnitude of the brightness asymmetry), this degeneracy can be broken.
If the emission is polarised (and the nature of the intrinsic polarisation is understood), light curves and images of the polarisation are strongly dependent upon those features in the image that are produced near the photon orbit. Hence, polarisation in general may be diagnostic of the black hole parameters.
Lastly, even if high precision images are unavailable, measurements of the centroid path coupled with observations of the light curve provide an alternative method by which to determine the black hole spin.
In reality, images of the Galactic centre may be expected to result from inhomogeneities during flaring events superimposed upon a homogeneous background. Thus, the results of this paper are complimentary to previous work. It is therefore likely that observations of both, a black hole shadow between flaring events and the dynamical and/or averaged properties of the images during flaring events, can be coupled to reduce the uncertainty inherent to each.
## Acknowledgements
This work was supported in part by NASA grant NAG 5-13292, and by NSF grants AST-0071019, AST-0204514 (for A.L.). A.E.B. gratefully acknowledges the support of an ITC Fellowship from Harvard College Observatory.
|
warning/0506/quant-ph0506043.html
|
ar5iv
|
text
|
# Segre variety, conifold, Hopf fibration, and separable multi-qubit states
## 1 Introduction
Quantum entanglement is one of the most interesting features of quantum theory. In quantum mechanics, the space of pure states in is an $`N+1`$-dimensional Hilbert space can be described by the complex projective space $`\mathrm{𝐂𝐏}^N`$. For bipartite, pure states, the entanglement of formation can be written in terms of concurrence . The connection between concurrence and geometry is found in a map called Segre embedding, see D. C. Brody and L. P. Hughston . They illustrate this map for a pair of qubits, and point. The Segre embedding has also been discussed in . There is also another geometrical description to describe pure state called Hopf fibration. The relation between Hopf fibration and single qubit and two-qubit states is discussed by R. Mosseri and R. Dandoloff . They have shown that $`𝐒^2`$ base space of a suitably oriented $`𝐒^3`$ Hopf fibration is nothing but the Bloch sphere, while the circular fibres represent the qubit overall phase degree of freedom. For two-qubit states, the Hilbert space is a seven-dimensional sphere $`𝐒^7`$, which also allows for a second Hopf fibration which is entanglement sensitive, with $`𝐒^3`$ fibres and a $`𝐒^4`$ base. Moreover, a generalization of Hopf fibration to three-qubit state has been presented in Ref. , where the Hilbert space of the three-qubit state is the fifteen-dimensional sphere $`𝐒^{15}`$, which allows for the third Hopf fibration with $`𝐒^8`$ as base and $`𝐒^7`$ as fiber. In this paper we will describe the Segre variety, which is a quadric space in algebraic geometry , by giving a complete and explicit formula for it. We will compare the Segre variety with the concurrence of pure, two-qubit states. The vanishing of the concurrence of a pure two-qubit state coincides with the Segre variety. Moreover, we will establish relations between Segre variety, conifold and Hopf fibration. In algebraic geometry, a conifold is a generalization of the notion of a manifold. Unlike manifolds, a conifold can contain conical singularities, i.e., points whose neighborhood look like a cone with a certain base. The base is usually a five-dimensional manifold. Conifold are very important in string theory, i.e., in the process of compactification of Calabi-Yau manifolds. A Calabi-Yau manifold is a compact Kähler manifold with a vanishing first Chern class. A Calabi-Yau manifold can also be defined as a compact Ricci-flat Kähler manifold. Finally, we will discuss the geometry and topology of pure multi-qubit states based on some mathematical tools from algebraic geometry and algebraic topology, namely the multi-projective Segre variety and higher-order Hopf fibration. Let us start by denoting a general, pure, composite quantum system with $`m`$ subsystems $`𝒬=𝒬_m^p(N_1,N_2,\mathrm{},N_m)=𝒬_1𝒬_2\mathrm{}𝒬_m`$, consisting of a pure state $`|\mathrm{\Psi }=_{i_1=1}^{N_1}_{i_2=1}^{N_2}\mathrm{}_{i_m=1}^{N_m}\alpha _{i_1,i_2,\mathrm{},i_m}|i_1,i_2,\mathrm{},i_m`$ and corresponding Hilbert space as $`_𝒬=_{𝒬_1}_{𝒬_2}\mathrm{}_{𝒬_m}`$, where the dimension of the $`j`$th Hilbert space is given by $`N_j=dim(_{𝒬_j})`$. We are going to use this notation throughout this paper, i.e., we denote a pure two-qubit states by $`𝒬_2^p(2,2)`$. Next, let $`\rho _𝒬`$ denotes a density operator acting on $`_𝒬`$. The density operator $`\rho _𝒬`$ is said to be fully separable, which we will denote by $`\rho _𝒬^{sep}`$, with respect to the Hilbert space decomposition, if it can be written as $`\rho _𝒬^{sep}=_{k=1}^\mathrm{N}p_k_{j=1}^m\rho _{𝒬_j}^k,_{k=1}^Np_k=1`$ for some positive integer $`\mathrm{N}`$, where $`p_k`$ are positive real numbers and $`\rho _{𝒬_j}^k`$ denotes a density operator on Hilbert space $`_{𝒬_j}`$. If $`\rho _𝒬^p`$ represents a pure state, then the quantum system is fully separable if $`\rho _𝒬^p`$ can be written as $`\rho _𝒬^{sep}=_{j=1}^m\rho _{𝒬_j}`$, where $`\rho _{𝒬_j}`$ is a density operator on $`_{𝒬_j}`$. If a state is not separable, then it is said to be entangled state.
## 2 Complex projective variety
In this section we will review basic definition of complex projective variety. Let $`\{f_1,f_2,\mathrm{},f_q\}`$ be continuous functions $`𝐊^n𝐊`$, where $`𝐊`$ is field of real $`𝐑`$ or complex number $`𝐂`$. Then we define real (complex) space as the set of simultaneous zeroes of the functions
$$𝒱_𝐊(f_1,f_2,\mathrm{},f_q)=\{(z_1,z_2,\mathrm{},z_n)𝐊^n:f_i(z_1,z_2,\mathrm{},z_n)=01iq\}.$$
(1)
These real (complex) spaces become a topological spaces by giving them the induced topology from $`𝐊^n`$. Now, if all $`f_i`$ are polynomial functions in coordinate functions, then the real (complex) space is called a real (complex) affine variety. A complex projective space $`\mathrm{𝐂𝐏}^n`$ which is defined to be the set of lines through the origin in $`𝐂^{n+1}`$, that is, $`\mathrm{𝐂𝐏}^n=(𝐂^{n+1}0)/`$, where $``$ is an equivalence relation define by $`(x_1,\mathrm{},x_{n+1})(y_1,\mathrm{},y_{n+1})\lambda 𝐂0`$ such that $`\lambda x_i=y_i0in`$. For $`n=1`$ we have a one dimensional complex manifold $`\mathrm{𝐂𝐏}^1`$ which is very important one, since as a real manifold it is homeomorphic to the 2-sphere $`𝐒^2`$. Moreover every complex compact manifold can be embedded in some $`\mathrm{𝐂𝐏}^n`$. In particular, we can embed a product of two projective spaces into a third one. Let $`\{f_1,f_2,\mathrm{},f_q\}`$ be a set of homogeneous polynomials in the coordinates $`\{\alpha _1,\alpha _2,\mathrm{},\alpha _{n+1}\}`$ of $`𝐂^{n+1}`$. Then the projective variety is defined to be the subset
$$𝒱(f_1,f_2,\mathrm{},f_q)=\{[\alpha _1,\mathrm{},\alpha _{n+1}]\mathrm{𝐂𝐏}^n:f_i(\alpha _1,\mathrm{},\alpha _{n+1})=01iq\}.$$
(2)
We can view the complex affine variety $`𝒱_𝐂(f_1,f_2,\mathrm{},f_q)𝐂^{n+1}`$ as complex cone over projective variety $`𝒱(f_1,f_2,\mathrm{},f_q)`$. We can also view $`\mathrm{𝐂𝐏}^n`$ as a quotient of the unit $`2n+1`$ sphere in $`𝐂^{n+1}`$ under the action of $`U(1)=𝐒^1`$, that is $`\mathrm{𝐂𝐏}^n=𝐒^{2n+1}/U(1)=𝐒^{2n+1}/𝐒^1`$, since every line in $`𝐂^{n+1}`$ intersects the unit sphere in a circle.
## 3 Hopf fibration and two- and three-qubit states
For a pure one-qubit state $`𝒬_1^p(2)`$ with $`|\mathrm{\Psi }=\alpha _1|1+\alpha _2|2`$, where $`\alpha _1,\alpha _2𝐂`$, and $`|\alpha _1|^2+|\alpha _2|^2=1`$, we can parameterize this state as
$$\left(\begin{array}{c}\alpha _1\\ \alpha _2\end{array}\right)=\left(\begin{array}{c}\mathrm{cos}(\frac{\vartheta }{2})\mathrm{exp}(i(\frac{\phi }{2}+\frac{\chi }{2}))\\ \mathrm{cos}(\frac{\vartheta }{2})\mathrm{exp}(i(\frac{\phi }{2}\frac{\chi }{2}))\end{array}\right)$$
(3)
where $`\vartheta [0,\pi ]`$, $`\phi [0,2\pi ]`$ and $`\chi [0,2\pi ]`$. The Hilbert space $`_𝒬`$ of a single qubit is the unit 3-dimensional sphere $`𝐒^3𝐑^4=𝐂^2`$. But since quantum mechanics is $`U(1)`$ projective, the projective Hilbert space is defined up to a phase $`\mathrm{exp}(i\phi )`$, so we have $`\mathrm{𝐂𝐏}^1=𝐒^3/U(1)=𝐒^3/𝐒^1=𝐒^2`$. Now, the first Hopf map, as an $`𝐒^1`$ fibration over a base space $`𝐒^2`$. For a pure two-qubit state $`𝒬_2^p(2,2)`$ with $`|\mathrm{\Psi }=\alpha _{1,1}|1,1+\alpha _{1,2}|1,2+\alpha _{2,1}|2,1+\alpha _{2,2}|2,2`$, where $`\alpha _{1,1},\alpha _{1,2},\alpha _{2,1},\alpha _{2,2}𝐂`$ and $`_{k,l}^2|\alpha _{k,l}|^2=1`$. The normalization condition identifies the Hilbert space $`_𝒬`$ to be the seven dimensional sphere $`𝐒^7𝐑^8=𝐂^4`$ and the projective Hilbert space to be $`\mathrm{𝐂𝐏}^3=𝐒^7/\mathrm{U}(1).`$ Thus we can parameterized the sphere $`𝐒^7`$ as a $`𝐒^3`$ fiber over $`𝐒^4`$, that is which is called the Hopf second fibration. This Hopf map is entanglement sensitive and the separable states satisfy $`\alpha _{1,1}\alpha _{2,2}=\alpha _{1,2}\alpha _{2,1}`$, see Ref. .
## 4 Segre variety for a general bipartite state and concurrence
For given general pure bipartite state $`𝒬_2^p(N_1,N_2)`$ we want make $`\mathrm{𝐂𝐏}^{N_11}\times \mathrm{𝐂𝐏}^{N_21}`$ into a projective variety by its Segre embedding which we construct as follows. Let $`(\alpha _1,\alpha _2,\mathrm{},\alpha _{N_1})`$ and $`(\alpha _1,\alpha _2,\mathrm{},\alpha _{N_2})`$ be two points defined on $`\mathrm{𝐂𝐏}^{N_11}`$ and $`\mathrm{𝐂𝐏}^{N_21}`$, respectively, then the Segre map
$$\begin{array}{cc}𝒮_{N_1,N_2}:\mathrm{𝐂𝐏}^{N_11}\times \mathrm{𝐂𝐏}^{N_21}& \mathrm{𝐂𝐏}^{N_1N_21}\end{array}$$
(4)
$$\begin{array}{c}((\alpha _1,\mathrm{},\alpha _{N_1}),(\alpha _1,\mathrm{},\alpha _{N_2}))(\alpha _{1,1},\mathrm{},\alpha _{1,N_1},\mathrm{},\alpha _{N_1,1},\mathrm{},\alpha _{N_1,N_2})\end{array}$$
is well defined. Next, let $`\alpha _{i,j}`$ be the homogeneous coordinate function on $`\mathrm{𝐂𝐏}^{N_1N_21}`$. Then the image of the Segre embedding is an intersection of a family of quadric hypersurfaces in $`\mathrm{𝐂𝐏}^{N_1N_21}`$, that is
$`\mathrm{Im}(𝒮_{N_1,N_2})`$ $`=`$ $`<\alpha _{i,k}\alpha _{j,l}\alpha _{i,l}\alpha _{j,k}>=𝒱\left(\alpha _{i,k}\alpha _{j,l}\alpha _{i,l}\alpha _{j,k}\right).`$ (5)
This quadric space is the space of separable states and it coincides with the definition of general concurrence $`𝒞(𝒬_2^p(N_1,N_2))`$ of a pure bipartite state because
$`𝒞(𝒬_2^p(N_1,N_2))`$ $`=`$ $`\left(𝒩{\displaystyle \underset{j,i=1}{\overset{N_1}{}}}{\displaystyle \underset{l,k=1}{\overset{N_2}{}}}\left|\alpha _{i,k}\alpha _{j,l}\alpha _{i,l}\alpha _{j,k}\right|^2\right)^{\frac{1}{2}},`$ (6)
where $`𝒩`$ is a somewhat arbitrary normalization constant. The separable set is defined by $`\alpha _{i,k}\alpha _{j,l}=\alpha _{i,l}\alpha _{j,k}`$ for all $`i,j`$ and $`k,l`$. I.e., for a two qubit state we have $`𝒮_{2,2}:\mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1\mathrm{𝐂𝐏}^3`$ and
$$\mathrm{Im}(𝒮_{2,2})=𝒱\left(\alpha _{1,1}\alpha _{2,2}\alpha _{1,2}\alpha _{2,1}\right)\alpha _{1,1}\alpha _{2,2}=\alpha _{1,2}\alpha _{2,1}$$
(7)
is a quadric surface in $`\mathrm{𝐂𝐏}^3`$ which coincides with the space of separable set of pairs of qubits. In following section comeback to this result.
## 5 Conifold
In this section we will give a short review of conifold. An example of real (complex) affine variety is conifold which is defined by
$$𝒱_𝐂(\underset{i=1}{\overset{4}{}}z_i^2)=\{(z_1,z_2,z_3,z_4)𝐂^4:\underset{i=1}{\overset{4}{}}z_i^2=0\}.$$
(8)
and conifold as a real affine variety is define by
$$𝒱_𝐑(f_1,f_2)=\{(x_1,\mathrm{},x_4,y_1,\mathrm{},y_4)𝐑^8:\underset{i=1}{\overset{4}{}}x_i^2=\underset{j=1}{\overset{4}{}}y_j^2,\underset{i=1}{\overset{4}{}}x_iy_i=0\}.$$
(9)
where $`f_1=_{i=1}^4(x_i^2y_i^2)`$ and $`f_2=_{i=1}^4x_iy_i`$. This can be seen by defining $`z=x+iy`$ and identifying imaginary and real part of equation $`_{i=1}^4z_i^2=0`$. As a real topological space $`𝒱_𝐑(f_1,\mathrm{},f_n)𝐑^n`$, $`x𝒱_𝐑(f_1,\mathrm{},f_n)`$ is a smooth point of $`𝒱_𝐑(f_1,\mathrm{},f_n)`$ if there is a neighborhood $`V`$ of $`x`$ such that $`V`$ is homeomorphic to $`𝐑^d`$ for some $`d`$ which is usually called the local dimension of $`𝒱_𝐑(f_1,\mathrm{},f_n)`$ in $`x`$. If there is no such neighborhood $`V`$, then $`x`$ is said to be a singular point of $`𝒱_𝐑(f_1,\mathrm{},f_n)`$. Now, we can call $`𝒱_𝐑(f_1,\mathrm{},f_n)`$ a topological manifold if all points $`x𝒱_𝐑(f_1,\mathrm{},f_n)`$ are smooth. $`𝐒^n`$ is compact, since it is a closed and bounded subset of $`𝐑^{n+1}`$. Now, let us define a cone as a real space $`𝒱_𝐑(f_1,\mathrm{},f_n)𝐑^n`$ with a specified point $`s`$ such that for all $`x𝒱_𝐑(f_1,\mathrm{},f_n)`$ we have that the line $`sx𝒱_𝐑(f_1,\mathrm{},f_n)`$. But every line $`s𝐑^n`$ intersect any sphere $`𝐒^{n1}`$ with center $`s`$, the cone $`𝒱_𝐑(f_1,\mathrm{},f_n)`$ can be determined by a compact space $`=𝒱_𝐑(f_1,\mathrm{},f_n)𝐒^{n1}`$ called the base space of the cone. As a real space, the conifold is cone in $`𝐑^8`$ with top the origin and base space the compact manifold $`𝐒^2\times 𝐒^3`$. One can reformulate this relation in term of a theorem. The conifold $`𝒱_𝐂(_{i=1}^4z_i^2)`$ is the complex cone over the Segre variety $`\mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1𝐒^2\times 𝐒^2`$. To see this let us make a complex linear change of coordinate $`\alpha _{1,1}^{^{}}=z_1+iz_2`$, $`\alpha _{1,2}^{^{}}=z_4+iz_3`$, $`\alpha _{2,1}^{^{}}=z_4+iz_3`$, and $`\alpha _{2,2}^{^{}}=z_1iz_2`$. Thus after this linear coordinate transformation we have
$$𝒱_𝐂(\alpha _{1,1}^{^{}}\alpha _{2,2}^{^{}}\alpha _{1,2}^{^{}}\alpha _{2,1}^{^{}})=𝒱_𝐂(\underset{i=1}{\overset{4}{}}z_i^2)𝐂^4.$$
(10)
We will comeback to this result in section 6 where we establish a relation between these varieties, Hopf fibration and two-qubit state. Moreover, removal of singularity of a conifold leads to a Segre variety which also describes the separable two-qubit states. We will investigate this connection in the following section. We can also define a metric on conifold as $`dS_6^2=dr^2+r^2dS_{T^{1,1}}^2`$, where
$$dS_{T^{1,1}}^2=\frac{1}{9}\left(d\psi +\underset{i=1}{\overset{2}{}}\mathrm{cos}\theta _id\varphi _i\right)^2+\frac{1}{6}\underset{i=1}{\overset{2}{}}\left(d\varphi _i^2+\mathrm{sin}^2\theta _id\varphi _i^2\right)^2,$$
(11)
is the metric on the Einstein manifold $`T^{1,1}=\frac{SU(2)\times SU(2)}{U(1)}`$, with $`U(1)`$ being a diagonal subgroup of the maximal torus of $`SU(2)\times SU(2)`$. Moreover, $`T^{1,1}`$ is a $`U(1)`$ bundle over $`𝐒^2\times 𝐒^2`$, where $`0\psi 4`$ is an angular coordinate and $`(\theta _i,\varphi _i)`$ for all $`i=1,2`$ parameterize the two $`𝐒^2`$, see Ref. . One can even relate these angular coordinate to the $`\alpha _{k,l}^{^{}}`$ for all $`k,l=1,2`$ as follows
$$\begin{array}{cc}\alpha _{1,1}^{^{}}=r^{3/2}e^{\frac{i}{2}(\psi \varphi _1\varphi _2)}\mathrm{sin}\frac{\theta _1}{2}\mathrm{sin}\frac{\theta _2}{2}& \alpha _{1,2}^{^{}}=r^{3/2}e^{\frac{i}{2}(\psi +\varphi _1\varphi _2)}\mathrm{cos}\frac{\theta _1}{2}\mathrm{sin}\frac{\theta _2}{2}\\ \alpha _{2,1}^{^{}}=r^{3/2}e^{\frac{i}{2}(\psi \varphi _1+\varphi _2)}\mathrm{sin}\frac{\theta _1}{2}\mathrm{cos}\frac{\theta _2}{2}& \alpha _{2,2}^{^{}}=r^{3/2}e^{\frac{i}{2}(\psi +\varphi _1+\varphi _2)}\mathrm{cos}\frac{\theta _1}{2}\mathrm{cos}\frac{\theta _2}{2}\end{array}.$$
Moreover, if we define the conifold as $`𝒱_𝐂(_{i=1}^4z_i^2)`$, then we identify the Einstein manifold $`T^{1,1}`$ as the intersection of conifold with the variety $`𝒱_𝐂(_{i=1}^4|z_i^2|r^3)`$ and $`T^{1,1}`$ is invariant under rotations $`SO(4)=SU(2)\times SU(2)`$ of $`z_i`$ coordinate and under an overall phase rotation.
## 6 Conifold, Segre variety, and a pure two-qubit state
In this section we will investigate relations between pure two-qubit states, Segre variety, and conifold. For a pure two-qubit state the Segre variety is given by $`𝒮_{2,2}:\mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1\mathrm{𝐂𝐏}^3`$ and
$`\mathrm{Im}(𝒮_{2,3})`$ $`=`$ $`𝒱\left(\alpha _{1,1}\alpha _{2,2}\alpha _{1,2}\alpha _{2,1}\right)`$
$`=`$ $`𝒱(\alpha _{1,1}^{}_{}{}^{}2+\alpha _{2,2}^{}_{}{}^{}2+\alpha _{1,2}^{}_{}{}^{}2+\alpha _{2,1}^{}_{}{}^{}2)`$
$`=`$ $`\mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1𝐒^2\times 𝐒^2`$
$``$ $`\text{}.`$ (15)
where we have performed a coordinate transformation on ideal of Segre variety $`\mathrm{Im}(𝒮_{2,2})`$. Moreover, we have the following commutative diagram
where $`\mathrm{𝐇𝐏}^1`$ denotes projective space over quaternion number field and we have the second Hopf fibration . Thus we have established a direct relation between two-qubit state, Segre variety, conic variety and Hopf fibration. Thus the result from algebraic geometry and algebraic topology give a unified picture of two-qubit state. Now, let us investigate what happens to our state, when we do the coordinate transformation to establish relation between conic variety and Segre variety. By the coordinate transformation $`\alpha _{1,1}^{^{}}=\alpha _{1,1}+i\alpha _{1,2}`$, $`\alpha _{1,2}^{^{}}=\alpha _{2,2}+i\alpha _{2,1}`$, $`\alpha _{2,1}^{^{}}=\alpha _{2,2}+i\alpha _{2,1}`$, and $`\alpha _{2,2}^{^{}}=\alpha _{1,1}i\alpha _{1,2}`$ we perform the following map $`|\mathrm{\Psi }=\alpha _{1,1}|1,1+\alpha _{1,2}|1,2+\alpha _{2,1}|2,1+\alpha _{2,2}|2,2|\mathrm{\Psi }^{^{}}`$ which is given by
$`|\mathrm{\Psi }^{^{}}`$ $`=`$ $`\alpha _{1,1}^{^{}}|1,1+\alpha _{1,2}^{^{}}|1,2+\alpha _{2,1}^{^{}}|2,1+\alpha _{2,2}^{^{}}|2,2`$
$`=`$ $`\alpha _{1,1}(|1,1+|2,2)+i\alpha _{1,2}(|1,1|2,2)`$
$`+i\alpha _{2,1}(|1,2+|2,1)\alpha _{2,2}(|1,2|2,1)`$
$`=`$ $`\sqrt{2}\left(\alpha _{1,1}|\mathrm{\Psi }^++i\alpha _{1,2}|\mathrm{\Psi }^{}+i\alpha _{2,1}|\mathrm{\Phi }^+\alpha _{2,2}|\mathrm{\Phi }^{}\right).`$
Thus the equality between Segre variety, conic variety means that we rewrite a pure two-qubit state in terms of Bell’s basis. For higher dimensional space we have Segre variety but we couldn’t find any relation between these two variety.
## 7 Segre variety, Hopf fibration, and multi-qubit states
In this section, we will generalize the Segre variety to a multi-projective space and then we will establish connections between Segre variety for multi-qubit state and Hopf fibration. As in the previous section, we can make $`\mathrm{𝐂𝐏}^{N_11}\times \mathrm{𝐂𝐏}^{N_21}\times \mathrm{}\times \mathrm{𝐂𝐏}^{N_m1}`$ into a projective variety by its Segre embedding following almost the same procedure. Let $`(\alpha _1,\alpha _2,\mathrm{},\alpha _{N_j})`$ be points defined on $`\mathrm{𝐂𝐏}^{N_j1}`$. Then the Segre map
$$\begin{array}{ccc}𝒮_{N_1,\mathrm{},N_m}:\mathrm{𝐂𝐏}^{N_11}\times \mathrm{𝐂𝐏}^{N_21}\times \mathrm{}\times \mathrm{𝐂𝐏}^{N_m1}& & \mathrm{𝐂𝐏}^{N_1N_2\mathrm{}N_m1}\\ ((\alpha _1,\alpha _2,\mathrm{},\alpha _{N_1}),\mathrm{},(\alpha _1,\alpha _2,\mathrm{},\alpha _{N_m}))& & (\mathrm{},\alpha _{i_1,i_2,\mathrm{},i_m},\mathrm{}).\end{array}$$
(17)
is well defined for $`\alpha _{i_1,i_2,\mathrm{},i_m}`$,$`1i_1N_1,1i_2N_2,\mathrm{},1i_mN_m`$ as a homogeneous coordinate-function on $`\mathrm{𝐂𝐏}^{N_1N_2\mathrm{}N_m1}`$. Now, let us consider the composite quantum system $`𝒬_m^p(N_1,N_2,\mathrm{},N_m)`$ and let the coefficients of $`|\mathrm{\Psi }`$, namely $`\alpha _{i_1,i_2,\mathrm{},i_m}`$, make an array as follows
$$𝒜=\left(\alpha _{i_1,i_2,\mathrm{},i_m}\right)_{1i_jN_j},$$
(18)
for all $`j=1,2,\mathrm{},m`$. $`𝒜`$ can be realized as the following set $`\{(i_1,i_2,\mathrm{},i_m):1i_jN_j,j\}`$, in which each point $`(i_1,i_2,\mathrm{},i_m)`$ is assigned the value $`\alpha _{i_1,i_2,\mathrm{},i_m}`$. Then $`𝒜`$ and it’s realization is called an $`m`$-dimensional box-shape matrix of size $`N_1\times N_2\times \mathrm{}\times N_m`$, where we associate to each such matrix a sub-ring $`\mathrm{S}_𝒜=𝐂[𝒜]\mathrm{S}`$, where $`\mathrm{S}`$ is a commutative ring over the complex number field. For each $`j=1,2,\mathrm{},m`$, a two-by-two minor about the $`j`$-th coordinate of $`𝒜`$ is given by
$`𝒞_{k_1,l_1;k_2,l_2;\mathrm{};k_m,l_m}`$ $`=`$ $`\alpha _{k_1,k_2,\mathrm{},k_m}\alpha _{l_1,l_2,\mathrm{},l_m}`$
$`\alpha _{k_1,k_2,\mathrm{},k_{j1},l_j,k_{j+1},\mathrm{},k_m}\alpha _{l_1,l_2,\mathrm{},l_{j1},k_j,l_{j+1},\mathrm{},l_m}\mathrm{S}_𝒜.`$
Then the ideal $`_𝒜^m`$ of $`\mathrm{S}_𝒜`$ is generated by $`𝒞_{k_1,l_1;k_2,l_2;\mathrm{};k_m,l_m}`$ and describes the separable states in $`\mathrm{𝐂𝐏}^{N_1N_2\mathrm{}N_m1}`$. The image of the Segre embedding $`\mathrm{Im}(𝒮_{N_1,N_2,\mathrm{},N_m})`$ which again is an intersection of families of quadric hypersurfaces in $`\mathrm{𝐂𝐏}^{N_1N_2\mathrm{}N_m1}`$ is given by
$`\mathrm{Im}(𝒮_{N_1,N_2,\mathrm{},N_m})`$ $`=`$ $`<𝒞_{k_1,l_1;k_2,l_2;\mathrm{};k_m,l_m}>`$
$`=`$ $`𝒱\left(𝒞_{k_1,l_1;k_2,l_2;\mathrm{};k_m,l_m}\right).`$
In our paper , we showed that the Segre variety defines the completely separable states of a general multipartite state. Furthermore, based on this sub-determinant, we define an entanglement measure for general pure bipartite and three-partite states which coincide with generalized concurrence. Let us consider a general multi-qubit state $`𝒬_m^p(2,\mathrm{},2)`$. For this state the Segre variety is given by equation (7) and
$`\mathrm{Im}(𝒮_{2,\mathrm{},2})`$ $`=`$ $`𝒱\left(𝒞_{1,2;1,2;\mathrm{};1,2}\right)`$
$`=`$ $`\stackrel{m\text{times}}{\stackrel{}{\mathrm{𝐂𝐏}^1\times \mathrm{}\times \mathrm{𝐂𝐏}^1}}𝐒^2\times \mathrm{}\times 𝐒^2`$
$``$ $`\text{}.`$ (24)
We can parameterized the sphere $`𝐒^{2^{m+1}1}`$ as a $`𝐒^{2^m1}`$ fiber over $`𝐒^{2^m}`$, that is which are higher order Hopf fibration. Moreover, we have the following commutative diagram
Thus we have established relations between Segre variety and higher order Hopf fibration and separable set of a multi-qubit state. As an example, let us look at a pure three-qubit state. For such state we have
$`\mathrm{Im}(𝒮_{2,2,2})`$ $`=`$ $`𝒱\left(𝒞_{1,2;1,2;1,2}\right)`$
$`=`$ $`\alpha _{1,1,1}\alpha _{2,1,2}\alpha _{1,1,2}\alpha _{2,1,1},\alpha _{1,1,1}\alpha _{2,2,1}\alpha _{1,2,1}\alpha _{2,1,1}`$
$`,\alpha _{1,1,1}\alpha _{2,2,2}\alpha _{1,2,2}\alpha _{2,1,1},\alpha _{1,1,2}\alpha _{2,2,1}\alpha _{1,2,1}\alpha _{2,1,2}`$
$`,\alpha _{1,1,2}\alpha _{2,2,2}\alpha _{1,2,2}\alpha _{2,1,2},\alpha _{1,2,1}\alpha _{2,2,2}\alpha _{1,2,2}\alpha _{2,2,1}`$
$`,\alpha _{1,1,1}\alpha _{1,2,2}\alpha _{1,1,2}\alpha _{1,2,1},\alpha _{1,1,1}\alpha _{2,2,2}\alpha _{1,2,1}\alpha _{2,1,2}`$
$`,\alpha _{1,1,2}\alpha _{2,2,1}\alpha _{1,2,2}\alpha _{2,1,1},,\alpha _{2,1,1}\alpha _{2,2,2}\alpha _{2,1,2}\alpha _{2,2,1}`$
$`,\alpha _{1,1,1}\alpha _{2,2,2}\alpha _{1,1,2}\alpha _{2,2,1},\alpha _{1,2,1}\alpha _{2,1,2}\alpha _{1,2,2}\alpha _{2,1,1}`$
$`=`$ $`\mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1\times \mathrm{𝐂𝐏}^1𝐒^2\times 𝐒^2\times 𝐒^2`$
$``$ $`\text{}.`$ (28)
This is what we have expected to see. Moreover, we have the following commutative diagram
where we have the third Hopf fibration for three-qubit state which has been discussed in Ref. .
## 8 Conclusion
In this paper, we have discussed a geometric picture of the separable pure two-qubit states based on Segre variety, conifold, and Hopf fibration. We have shown that these varieties and mappings give a unified picture of two-qubit states. Moreover, we have discussed the geometry and topology of pure multi-qubit states based on multi-projective Segre variety and higher-order Hopf fibration. Thus we have established relations between algebraic geometry, algebraic topology and fundamental quantum theory of entanglement. Perhaps, these geometrical and topological visualization puts entanglement in a broader perspective and hopefully gives some hint about how we can solve the problem of quantify entanglement.
Acknowledgments: This work was supported by the Wenner-Gren Foundation
|
warning/0506/math0506192.html
|
ar5iv
|
text
|
# Une base symétrique de l’algèbre des coinvariants quasi-symétriques
## 0 Introduction
L’algèbre des coinvariants est un objet classique associé à chaque groupe de Coxeter fini $`W`$ . Cette algèbre est définie comme le quotient de l’algèbre des polynômes sur l’espace vectoriel sur lequel $`W`$ agit par réflexions, par l’idéal homogène engendré par les polynômes invariants homogènes non constants. Le quotient est une algèbre graduée de dimension finie donnée par l’ordre de $`W`$.
Dans le cas du groupe symétrique sur $`n+1`$ lettres, on peut expliciter cette construction comme le quotient de l’algèbre des polynômes en $`𝗑_1,\mathrm{},𝗑_{n+1}`$ par l’idéal homogène engendré par les fonctions symétriques élémentaires. Plus récemment, une notion plus faible que celle de polynôme symétrique est apparue . Un polynôme est dit quasi-symétrique si, pour toute suite d’exposants $`(m_1,\mathrm{},m_k)`$ fixée, tous les monômes $`𝗑_{i_1}^{m_1}\mathrm{}𝗑_{i_k}^{m_k}`$ pour une suite *croissante* d’indices $`i_1<i_2<\mathrm{}<i_k`$ ont le même coefficient. En particulier, les polynômes symétriques sont aussi quasi-symétriques.
Dans , la notion d’algèbre coinvariante quasi-symétrique a été introduite et étudiée. Elle est définie comme le quotient de l’algèbre des polynômes en $`𝗑_1,\mathrm{},𝗑_{n+1}`$ par l’idéal homogène engendré par les polynômes quasi-symétriques homogènes non constants. C’est une algèbre graduée. Il est démontré dans que cette algèbre est de dimension finie, donnée par le nombre de Catalan $`c_{n+1}`$. La preuve est la construction explicite d’un ensemble de monômes indexés par les chemins de Dyck de longueur $`2n+2`$, dont les images dans le quotient forment une base de l’algèbre des coinvariants quasi-symétriques.
Dans le cas des coinvariants usuels, le groupe de Coxeter $`W`$ agit par automorphismes sur le quotient et on obtient une décomposition intéressante du module régulier. Dans le cas des coinvariants quasi-symétriques, le seul automorphisme de la situation est le renversement qui envoie $`𝗑_i`$ sur $`𝗑_{n+2i}`$. Cette involution préserve l’idéal des fonctions quasi-symétriques sans terme constant et passe donc au quotient.
La motivation initiale de cet article est le fait que l’action de cette involution semble difficile à décrire dans la base des monômes associés aux chemins de Dyck. On construit donc une nouvelle base, dans laquelle l’involution agit par permutation. Cette base est formée de polynômes dont le terme dominant pour l’ordre naturel sur les variables $`𝗑_1,\mathrm{},𝗑_{n+1}`$ redonne les monômes associés aux chemins de Dyck.
Il apparaît que l’ensemble naturel d’indexation de cette nouvelle base est non pas l’ensemble des chemins de Dyck, mais celui des triangulations d’un polygone régulier. Cet ensemble joue un rôle primordial dans la théorie des algèbres à grappes de Fomin et Zelevinsky . Cet article donne donc un premier rapprochement entre les algèbres à grappes et les fonctions quasi-symétriques. Il se trouve que la construction de la base indexée par les triangulations passe par le choix d’une triangulation de base en forme d’éventail. Dans le cadre de la théorie des algèbres à grappes, ce choix correspond au carquois équi-orienté de type $`A_n`$, voir par exemple .
L’article est organisé comme suit. On commence par définir une bijection ad hoc entre triangulations et chemins de Dyck. Ensuite on montre que, par cette bijection, le monôme dominant du polynôme associé à une triangulation est le monôme associé au chemin de Dyck correspondant, ce qui entraîne immédiatement le résultat principal.
## 1 Bijection
Soit $`n`$ un entier positif ou nul. On définit dans cette section une bijection entre
1. les triangulations d’un polygone régulier à $`n+3`$ cotés,
2. les chemins de Dyck de longueur $`2n+2`$.
Il est bien connu que ces deux ensembles ont pour cardinal le nombre de Catalan
$$c_{n+1}=\frac{1}{n+2}\left(\genfrac{}{}{0pt}{}{2n+2}{n+1}\right).$$
(1)
Par définition, un chemin de Dyck est une suite de pas verticaux (“montées”) et horizontaux (“descentes”) qui reste au dessus de la diagonale, voir la partie droite de la figure 1.
La bijection est illustrée par un exemple dans la figure 1.
Avant toute chose, on fixe une triangulation de base en forme d’éventail, c’est-à-dire formée par toutes les diagonales contenant un sommet choisi, noté $`\mathrm{\#}`$. On dessine cette triangulation avec le sommet commun à toutes les diagonales placé en bas. Les diagonales de cette triangulation de base seront dites “négatives” et numérotées de $`1`$ à $`n`$ de gauche à droite. Les diagonales qui n’interviennent pas dans la triangulation de base sont dites “positives”. On numérote aussi de $`1`$ à $`n`$ les sommets aux extrémités des diagonales négatives.
On associe alors un chemin de Dyck $`D(T)`$ à chaque triangulation $`T`$, par récurrence sur $`n`$. Pour $`n=0`$, à la seule triangulation du polygone à trois cotés est associée le seul chemin de Dyck de longueur $`2`$.
Si $`n`$ est non nul, on regarde le sommet $``$ du polygone placé à droite du sommet $`\mathrm{\#}`$ dans le sens trigonométrique. On distingue deux cas.
Si le sommet $``$ participe à un seul triangle de la triangulation $`T`$ i.e. n’est contenu dans aucune diagonale de $`T`$, on lui associe le chemin de Dyck obtenu en encadrant par une montée et une descente le chemin de Dyck $`D(T^{})`$ associé à la triangulation $`T^{}`$ du polygone à $`n+2`$ cotés qui est définie comme $`T`$ moins le triangle adjacent à $``$. Le sommet distingué $`\mathrm{\#}`$ de $`T^{}`$ est celui de $`T`$.
Si le sommet $``$ participe à plusieurs triangles, on découpe la triangulation en autant de morceaux (le long des diagonales contenant $``$), voir la figure 2. Le sommet $``$ donne un sommet dans chacun de ces morceaux. On prend dans chacun des morceaux le sommet à gauche de $``$ comme sommet distingué $`\mathrm{\#}`$. Par récurrence, on associe un chemin de Dyck à chacun des morceaux et on les concatène dans l’ordre des morceaux induit par l’ordre de gauche à droite au voisinage du sommet $``$ dans $`T`$, voir les figures 1 et 2.
C’est clairement une bijection. La bijection inverse est aussi définie par récurrence sur $`n`$. On décompose un chemin de Dyck réductible pour la concaténation en ses composantes irréductibles et on recompose une triangulation par juxtaposition. Pour les chemins de Dyck irréductibles, on enlève une montée et une descente, on obtient une triangulation par récurrence et on rajoute un triangle.
###### Lemme 1.1
Le nombre de pas verticaux initiaux du chemin de Dyck $`D(T)`$ est le nombre de diagonales négatives dans $`T`$ plus $`1`$.
* Preuve. La preuve se fait par récurrence. L’énoncé est vrai pour $`n=0`$. On distingue deux cas comme dans la définition de la bijection. Dans le cas où $``$ est dans une seule diagonale, les deux quantités augmentent de $`1`$. Dans l’autre cas, les deux quantités sont inchangées.
## 2 Polynômes
On associe à chaque diagonale un polynôme en les variables $`\{𝗑_1,\mathrm{},𝗑_{n+1}\}`$ comme suit. On associe la constante $`1`$ aux diagonales négatives. Chaque diagonale positive coupe un ensemble de diagonales négatives consécutives de $`i`$ à $`j`$. En fait, ceci donne une bijection entre les diagonales positives et les segments de $`\{1,\mathrm{},n\}`$. On peut donc parler de la diagonale positive $`(i,j)`$, à qui on associe alors la somme des $`𝗑_k𝗑_{k+1}`$ pour $`k=i,\mathrm{},j`$ soit $`𝗑_i𝗑_{j+1}`$.
On associe alors à chaque triangulation $`T`$ le produit $`𝖡_T`$ des polynômes associés à ses diagonales. Dans l’exemple de la figure 1, on obtient
$$𝖡_T=(𝗑_1𝗑_2)(1)(𝗑_3𝗑_4)(𝗑_3𝗑_6)(𝗑_5𝗑_6).$$
(2)
Par ailleurs, comme dans , on associe un monôme $`M_D`$ en $`\{𝗑_1,\mathrm{},𝗑_n\}`$ à chaque chemin de Dyck $`D`$. On représente un chemin de Dyck par une suite de pas d’une unité vers le haut (“montée”) ou vers la droite (“descente”) dans une grille. On numérote les colonnes internes de la grille de $`1`$ à $`n`$, voir la partie droite de la figure 1. On convient que chaque pas vertical d’indice $`i`$ correspond à la variable $`𝗑_i`$. Le monôme $`M_D`$ est alors le produit des contributions des pas verticaux. Dans l’exemple de la figure 1, on obtient
$$M_{D(T)}=𝗑_1𝗑_3𝗑_3𝗑_5.$$
(3)
On définit un ordre sur les monômes en ordonnant les variables par
$$𝗑_1𝗑_2\mathrm{}𝗑_{n+1}.$$
(4)
Le monôme dominant d’un polynôme pour cet ordre est celui où intervient la plus grande puissance de $`𝗑_1`$, puis en cas d’ambiguïté la plus grand puissance de $`𝗑_2`$ et ainsi de suite.
###### Proposition 2.1
Le monôme dominant du polynôme $`𝖡_T`$ associé à une triangulation $`T`$ est le monôme $`M_{D(T)}`$ associé au chemin de Dyck $`D(T)`$ correspondant à $`T`$ via la bijection ci-dessus.
* Preuve. Par récurrence sur $`n`$. La proposition est vraie pour $`n=0`$. Soit donc $`n`$ non nul. On distingue deux cas.
Supposons d’abord que le sommet $``$ participe à un seul triangle de la triangulation $`T`$. Alors la triangulation $`T`$ contient la diagonale négative $`n`$. Le polynôme $`𝖡_T`$ ne fait donc pas intervenir $`𝗑_{n+1}`$ et est égal au polynôme $`𝖡_T^{}`$ associé à la triangulation raccourcie en $``$. De même, le chemin de Dyck $`D(T)`$ est obtenu par concaténation d’une montée, du chemin de Dyck $`D(T^{})`$ et d’une descente. Donc le monôme associé à $`D(T)`$ est le même que celui associé au chemin $`D(T^{})`$. On conclut par hypothèse de récurrence que le monôme dominant de $`𝖡_T`$ est $`M_{D(T)}`$.
Supposons maintenant que le sommet $``$ participe à plusieurs triangles de $`T`$. Soit $`\mathrm{Ext}(T)`$ l’ensemble des nombres $`k`$ dans $`\{1,\mathrm{},n\}`$ tels que la diagonale négative $`k`$ partage un sommet avec un diagonale de $`T`$ contenant $``$. On va numéroter les diagonales de $`T`$ contenant $``$ par les éléments de $`\mathrm{Ext}(T)`$.
Dans la définition de $`𝖡_T`$ comme produit sur les diagonales de $`T`$, on peut séparer les contributions des diagonales strictement contenues dans les différents morceaux et la contribution des diagonales de $`T`$ séparant les morceaux. On va traiter séparément le morceau le plus à gauche et les autres morceaux. Ces autres morceaux sont numérotés par l’élément de $`\mathrm{Ext}(T)`$ qui les borde sur leur gauche.
La contribution des diagonales entre les morceaux est
$$\underset{k\mathrm{Ext}(T)}{}\left(𝗑_{k+1}𝗑_{n+1}\right).$$
(5)
Considérons le premier morceau et soit $`k_{\mathrm{min}}`$ le plus petit élément de $`\mathrm{Ext}(T)`$. La contribution du premier morceau est
$$\underset{\genfrac{}{}{0pt}{}{1ij<k_{\mathrm{min}}}{(i,j)T}}{}\left(𝗑_i𝗑_{j+1}\right).$$
(6)
Considérons maintenant $`k\mathrm{Ext}(T)`$ et le morceau correspondant, situé à droite de $`k`$. Soit $`k^{}`$ l’élément suivant de $`\mathrm{Ext}(T)`$ ou bien posons $`k^{}=n+1`$ si $`k`$ est le plus grand élément de $`\mathrm{Ext}(T)`$. La contribution du morceau $`k`$ est alors
$$\underset{\genfrac{}{}{0pt}{}{k+1i<k^{}}{(k+1,i)T}}{}\left(𝗑_{k+1}𝗑_{i+1}\right)\underset{\genfrac{}{}{0pt}{}{k+1<ij<k^{}}{(i,j)T}}{}\left(𝗑_i𝗑_{j+1}\right),$$
(7)
où le premier facteur est associé aux diagonales du morceau $`k`$ qui contiennent le sommet $`k`$.
On a donc montré que $`𝖡_T`$ est le produit de facteurs associés à chaque morceau : pour le premier morceau,
$$\underset{\genfrac{}{}{0pt}{}{1ij<k_{\mathrm{min}}}{(i,j)T}}{}\left(𝗑_i𝗑_{j+1}\right)$$
(8)
et, pour le morceau à droite de $`k`$ dans $`\mathrm{Ext}(T)`$,
$$\left(𝗑_{k+1}𝗑_{n+1}\right)\underset{\genfrac{}{}{0pt}{}{k+1i<k^{}}{(k+1,i)T}}{}\left(𝗑_{k+1}𝗑_{i+1}\right)\underset{\genfrac{}{}{0pt}{}{k+1<ij<k^{}}{(i,j)T}}{}\left(𝗑_i𝗑_{j+1}\right).$$
(9)
Regardons maintenant l’image $`D(T)`$ de $`T`$ par la bijection. C’est la concaténation des images des morceaux de $`T`$. Par définition du monôme associé, celui-ci est le produit des contributions de chaque morceau avec un décalage des indices convenable et des contributions des pas verticaux initiaux des morceaux (sauf le premier).
Par hypothèse de récurrence, la contribution du premier morceau est
$$\underset{\genfrac{}{}{0pt}{}{1ij<k_{\mathrm{min}}}{(i,j)T}}{}𝗑_i.$$
(10)
La contribution du morceau entre $`k\mathrm{Ext}(T)`$ et l’élément suivant $`k^{}`$ de $`\mathrm{Ext}(T)`$ est donnée, par hypothèse de récurrence, par
$$𝗑_{k+1}^\mathrm{}_k\underset{\genfrac{}{}{0pt}{}{k+1<ij<k^{}}{(i,j)T}}{}𝗑_i,$$
(11)
$`\mathrm{}_k`$ est le nombre de pas verticaux initiaux du morceau $`k`$.
Par le lemme 1.1 appliqué au morceau $`k`$, on sait que le nombre $`\mathrm{}_k`$ de pas verticaux initiaux dans le morceau $`k`$ de $`D(T)`$ est égal à $`1`$ plus le nombre de diagonales dans le morceau $`k`$ de $`T`$ qui contiennent le sommet $`k`$. La contribution du morceau $`k`$ au monôme $`M_{D(T)}`$ est donc
$$𝗑_{k+1}\underset{\genfrac{}{}{0pt}{}{k+1i<k^{}}{(k+1,i)T}}{}𝗑_{k+1}\underset{\genfrac{}{}{0pt}{}{k+1<ij<k^{}}{(i,j)T}}{}𝗑_i.$$
(12)
On vérifie que le terme dominant de la contribution de chaque morceau à $`𝖡_T`$ est bien égal à la contribution de chaque morceau à $`M_{D(T)}`$. En prenant le produit des contributions des morceaux, on obtient l’égalité voulue.
###### Théorème 2.2
Les polynômes $`𝖡_T`$ associés aux triangulations forment une base de l’algèbre des coinvariants quasi-symétriques. Cette base est stable par le renversement des variables $`𝗑_i𝗑_{n+2i}`$. Les deux choix naturels d’ordre total sur les variables donnent deux bases monomiales, en prenant les monômes dominants des polynômes $`𝖡_T`$.
* Preuve. Dans , il est démontré que les classes des monômes $`M_D`$ associés aux chemins de Dyck forment une base de l’anneau des coinvariants quasi-symétriques. On déduit alors de la proposition 2.1 que les classes des polynômes $`𝖡_T`$ forment aussi une base. Le fait que cette base soit stable par le renversement est immédiat : l’image de $`𝖡_T`$ est $`𝖡_T^{}`$ où la triangulation $`T^{}`$ est obtenue par renversement de $`T`$. Enfin la dernière assertion est juste une reformulation de la proposition 2.1 et son image par le renversement.
|
warning/0506/hep-ex0506045.html
|
ar5iv
|
text
|
# Observation of 𝜒_{𝑐𝐽}→𝜔𝜔 decays
## 1 Introduction
Exclusive quarkoninum decays provide an important laboratory for investigating perturbative quantum chromodynamics. Compared with $`J/\psi `$ and $`\psi (2S)`$ decays, one has much less knowledge on $`PC=++\chi _{cJ}`$ decays. While a few exclusive decays of $`\chi _{cJ}`$ have been measured, many decay modes remain unknown. Current theoretical analyses of $`\chi _{cJ}`$ decays provide only a rough treatment of the color-octet wave function. For $`\chi _{cJ}vectorvector`$ mode, so far only measurements of $`\chi _{cJ}\varphi \varphi `$ and $`\chi _{cJ}K^{}(892)^0\overline{K}^{}(892)^0`$ are available with low statistics. Precise measurements for more channels will help in better understanding the various mechanism of $`\chi _{cJ}`$ decays and the nature of $`{}_{}{}^{3}P_{J}^{}c\overline{c}`$ bound states.
Further, the decays of $`\chi _{cJ}`$, especially $`\chi _{c0}`$ and $`\chi _{c2}`$, provide a direct window on glueball dynamics in the $`0^{++}`$ and $`2^{++}`$ channels since the hadronic decays may proceed via $`c\overline{c}ggq\overline{q}q\overline{q}`$.
Recently, the branching ratio for $`\chi _{c0}f_0(980)f_0(980)`$ has been measured by the BES collaboration. In the present analysis, a search for $`\chi _{c0,2}`$ decaying into $`\pi ^+\pi ^{}\pi ^0\pi ^+\pi ^{}\pi ^0`$ final states is carried out using 14 million $`\psi (2S)`$ events accumulated at the upgraded BES detector (BESII). Signals of $`\chi _{c0}`$ and $`\chi _{c2}`$ decaying to $`\omega `$ pairs in $`\psi (2S)`$ radiative decays are observed for the first time.
## 2 The BES detector
The Beijing Spectrometer (BES) is a conventional solenoidal magnet detector that is described in detail in Ref. ; BESII is the upgraded version of the BES detector . A 12-layer vertex chamber (VC) surrounding the beam pipe provides trigger and position information. A forty-layer main drift chamber (MDC), located radially outside the VC, provides trajectory and energy loss ($`dE/dx`$) information for charged tracks over $`85\%`$ of the total solid angle. The momentum resolution is $`\sigma _p/p=0.017\sqrt{1+p^2}`$ ($`p`$ in GeV/c), and the $`dE/dx`$ resolution for hadron tracks is $`8\%`$. An array of 48 scintillation counters surrounding the MDC measures the time-of-flight (TOF) of charged tracks with a resolution of $`200`$ ps for hadrons. Outside of the TOF counters is a 12-radiation-length barrel shower counter (BSC) comprised of gas proportional tubes interleaved with lead sheets. This measures the energies of electrons and photons over $`80\%`$ of the total solid angle with an energy resolution of $`\sigma _E/E=22\%/\sqrt{E}`$ ($`E`$ in GeV). Outside of the solenoidal coil, which provides a 0.4 Tesla magnetic field over the tracking volume, is an iron flux return that is instrumented with three double layers of counters that identify muons of momentum greater than 0.5 GeV/c.
A GEANT3 based Monte Carlo (MC) program with detailed consideration of the detector performance (such as dead electronic channels) is used to simulate the BESII detector. The consistency between data and Monte Carlo has been carefully checked in many high purity physics channels, and the agreement is quite reasonable .
## 3 Event selection
### 3.1 $`𝝎𝝎`$ signal
In this analysis, $`\chi _{cJ}\omega \omega \pi ^+\pi ^{}\pi ^0\pi ^+\pi ^{}\pi ^0`$ channels are investigated using $`\psi (2S)`$ radiative decays to $`\chi _{cJ}`$. Events with four charged tracks and five or six photons are selected. Each charged track is required to be well fit by a three-dimensional helix and to have a polar angle, $`\theta `$, within the fiducial region $`|\mathrm{cos}\theta |<0.8`$. To ensure tracks originate from the interaction region, we require $`V_{xy}=\sqrt{V_x^2+V_y^2}<2`$ cm and $`|V_z|<20`$ cm, where $`V_x`$, $`V_y`$, and $`V_z`$ are the $`x,y`$ and $`z`$ coordinates of the point of closest approach of each track to the beam axis.
A neutral cluster is considered to be a photon candidate if it is located within the BSC fiducial region ($`|\mathrm{cos}\theta |<0.8`$), the energy deposited in the BSC is greater than 40 MeV, the first hit appears in the first 10 radiation lengths, and the angle between the cluster and the nearest charged track is greater than $`6^{}`$.
A six constraint (6-C) kinematic fit to the hypothesis $`\psi (2S)\gamma \pi ^+\pi ^{}\pi ^0\pi ^+\pi ^{}\pi ^0`$ with the invariant mass of the two photon pairs constrained to the $`\pi ^0`$ mass is performed, and the $`\chi ^2`$ of the 6-C fit is required to be less than 15. For events with six photons candidates, the combination having the minimum $`\chi ^2`$ is chosen, and the probability of the 6-C fit is required to be larger than that of the 7-C fit to the hypothesis $`\psi (2S)2\pi ^+2\pi ^{}3\pi ^0`$ to suppress potential background from $`\psi (2S)\omega \pi ^+\pi ^{}\pi ^0\pi ^02\pi ^+2\pi ^{}3\pi ^0`$.
Since there are four $`\omega `$ pair combinations from $`\pi ^+\pi ^{}\pi ^0\pi ^+\pi ^{}\pi ^0`$, the $`\omega `$ pair with the minimum $`R`$, which is defined as
$$R=\sqrt{(M_{\pi ^+\pi ^{}\pi ^0}^10.783)^2+(M_{\pi ^+\pi ^{}\pi ^0}^20.783)^2},$$
is chosen for further analysis. Here, $`M_{\pi ^+\pi ^{}\pi ^0}`$ is the invariant mass of three pions and superscript 1, 2 denote different pion combinations. Therefore, there is only one entry for each event.
Figures 1 and 2 show mass distributions for candidate events in the high mass ($`M_{6\pi }>3.2`$ GeV/$`c^2`$) and low mass regions ($`M_{6\pi }<3.2`$ GeV/$`c^2`$), respectively. Here (a) is the scatter plot of $`M_{\pi ^+\pi ^{}\pi ^0}`$ versus $`M_{\pi ^+\pi ^{}\pi ^0}`$, (b) is the $`M_{\pi ^+\pi ^{}\pi ^0}`$ distribution recoiling against the opposite $`\omega `$, selected by requiring $`|M_{\pi ^+\pi ^{}\pi ^0}783|<50`$ MeV/$`c^2`$, and (c) is the $`M_{\omega \omega }`$ invariant mass distribution for events in the $`\omega `$ pair signal region, defined by $`R<50`$ MeV/$`c^2`$. In Fig. 1, clear $`\omega `$ signal can be seen in (b), and clear $`\chi _{c0}`$ and $`\chi _{c2}`$ signals in (c), indicating the existence of $`\chi _{c0,2}\omega \omega `$ decays. By contrast, in the low $`M_{6\pi }`$ mass region, shown in Fig. 2, the $`\omega `$ pair signal is less significant than in the high mass region. Here, only $`\omega `$ pair events in high mass region are studied.
In order to test if the selection criteria in this analysis will give ‘fake’ $`\omega `$ pair events from non-$`\omega `$ pair events, 300000 MC simulated $`\psi \gamma \chi _{c0}\gamma 6\pi `$ events are generated in which $`\chi _{c0}6\pi `$ decays according to the phase space. Fig. 3 shows the $`M_{\pi ^+\pi ^{}\pi ^0}`$ distributions of the surviving MC phase space events after requiring the same selection criteria as for the real data. No peak around the $`\omega `$ mass is seen, which shows that the $`\omega `$ pair selection criteria in this analysis does not generate fake $`\omega `$ pair signals.
The annular region around the $`\omega `$ pair signal circle, shown in Fig 1(a), is taken as the sideband region. Fig. 4 shows the $`M_{6\pi }`$ sideband distributions defined using the radius R to be (a) $`150<R<300`$ MeV$`/c^2`$ and (b) $`100<R<200`$ MeV/$`c^2`$. No obvious $`\chi _{cJ}`$ signals seen in these sideband distributions.
### 3.2 MC simulation
A MC simulation of $`\psi (2S)\gamma \chi _{cJ},\chi _{cJ}\omega \omega `$ is used to determine the detection efficiency. The proper angular distributions of the photon emitted in $`\psi (2S)\gamma \chi _{cJ}`$ are used . Fig. 5 shows the distributions, identical to those in Fig. 1 for MC simulated $`\psi (2S)\gamma \chi _{c0},\chi _{c0}\omega \omega `$ events passing the same selection criteria as for the real data. MC simulated $`\psi (2S)\gamma \chi _{c2},\chi _{c2}\omega \omega `$ events have similar distributions.
### 3.3 Mass spectrum fit
The Maximum Likelihood (ML) method is used to fit the $`M_{\omega \omega }`$ mass spectrum of events in the $`\omega `$ pair signal region (Fig. 1(c)). The $`\chi _{0,2}`$ signal functions are determined from MC simulation, as shown in Fig. 5(c) for $`\chi _{c0}`$, while the background function is taken from the sideband distribution, shown in Fig 4(a). The fit result is represented by the solid curve in Fig. 6, and the fit yields
$$N_{\chi _{c0}}=38.1\pm 9.6,N_{\chi _{c2}}=27.7\pm 7.4.$$
The statistical significances of $`\chi _{c0}`$ and $`\chi _{c2}`$ are $`4.4\sigma `$ and $`4.7\sigma `$, respectively, which are estimated from $`\sqrt{2\mathrm{\Delta }ln}`$, where $`\mathrm{\Delta }ln`$ is the difference between the logarithmic ML values of the fit with and without the corresponding signal function.
## 4 Systematic error
The systematic error in this branching ratio measurement includes the uncertainties in the MDC tracking efficiency, photon efficiency, kinematic fit, background shape, number of $`\psi (2S)`$ events, etc.
### 4.1 MDC tracking efficiency and photon efficiency
For charged tracks, the uncertainty of the tracking efficiency is determined by comparing data and MC for some very clean $`J/\psi `$ decay channels , and an error of 2% is found for each track. A similar comparison has also been performed for photons , and the difference is also about 2% for a single photon.
### 4.2 Kinematic fit
The systematic error associated with the kinematic fit is due to differences between data and MC simulation in the determination of the track momentum, the track fitting error matrix, and the photon energy and direction. The effect is studied for charged tracks and neutral tracks separately. By comparing the number of events before and after the kinematic fit for very clean event samples for data and MC simulated data, the difference is determined to be 8.4%, which is taken as the systematic error.
### 4.3 Background shape
Two different sideband $`M_{6\pi }`$ spectrum shapes, shown in Fig. 4, are used as the background function. The difference in the number of $`\chi _{c0,2}`$ events obtained with the two different shapes is taken as a systematic error.
### 4.4 Binning, fit range, and signal region
The differences caused by different binning and fit ranges in the $`\omega \omega `$ mass spectrum fit are 1.2% and 3.4% for $`\chi _{c0}`$ and $`\chi _{c2}`$, respectively. Different sized signal regions yield differences of 3.1% and 2.4% for $`\chi _{c0}`$ and $`\chi _{c2}`$, respectively, which are taken as a systematic error.
### 4.5 Angular distribution of $`𝝌_{𝒄𝑱}\mathbf{}𝝎𝝎`$
In the estimation of the efficiency, a phase space generator with only the angular distribution of the radiative photon is considered. While this is correct for $`\chi _{c0}`$ decays, it may introduce bias for $`\chi _{c2}`$ decays. The effect is estimated by generating different angular distributions of the omega in the $`\chi _{c2}`$ rest frame. The efficiency difference between these tests and the phase space generator is estimated to be 9.4%, which is put into the systematic error.
### 4.6 Branching ratios of intermediate states
The errors on intermediate state branching ratios are obtained from the PDG except for $`(\psi (2S)\gamma \chi _{cJ})`$, where recent CLEO results are used. Table 1 summarizes all contributions to the systematic errors, and the total systematic error is determined by the quadratic sum of all terms.
## 5 Results
The branching ratio of $`(\chi _{cJ}\omega \omega )`$ is determined from
$$(\chi _{cJ}\omega \omega )=\frac{N_{\chi _{cJ}}^{obs}}{N_{\psi (2S)}f_1f_2^2f_3^2ϵ}$$
where $`N_{\chi _{cJ}}^{obs}`$ is the number of events selected, $`N_{\psi (2S)}`$ the total number of $`\psi (2S)`$ events, $`ϵ`$ is the detection efficiency for the investigated channel, and $`f_1,f_2`$ and $`f_3`$ are the branching ratios of $`\psi (2S)\gamma \chi _{cJ},\omega 3\pi `$, and $`\pi ^0\gamma \gamma `$, respectively. Table 2 lists the $`\chi _{c0,2}\omega \omega `$ branching ratio results, together with numbers used in the branching ratio calculation.
In summary, $`\omega \omega `$ signals in the decay of $`\chi _{c0,2}`$ are observed, and their branching ratios measured for the first time. $`\chi _{c0}`$ and $`\chi _{c2}`$ decays to $`\omega \omega `$ have similar decay branching ratios, which is different from other $`\chi _{cJ}VV`$ decays $`(\chi _{cJ}\varphi \varphi ,\overline{K}^{}(892)^0K^{}(892)^0)`$. This measurement, together with previous measurements of $`\chi _{cJ}VV`$, will be helpful in understanding the nature of $`\chi _{cJ}`$ states.
## 6 Acknowlegements
The BES collaboration thanks the staff of BEPC for their hard efforts. This work is supported in part by the National Natural Science Foundation of China under contracts Nos. 10491300, 10225524, 10225525, 10425523, the Chinese Academy of Sciences under contract No. KJ 95T-03, the 100 Talents Program of CAS under Contract Nos. U-11, U-24, U-25, and the Knowledge Innovation Project of CAS under Contract Nos. U-602, U-34 (IHEP), the National Natural Science Foundation of China under Contract No. 10225522 (Tsinghua University), and the Department of Energy under Contract No.DE-FG02-04ER41291 (U Hawaii).
|
warning/0506/hep-ph0506190.html
|
ar5iv
|
text
|
# Muon anomalous magnetic moment due to the brane-stretching effect
## I Introduction
The BNL E821 group recently reported the precision measurement of the muon anomalous magnetic moment. Based on their result, a new world average recorded $`a_\mu ^{(\text{exp})}=11659208(6)\times 10^{10}(\pm 0.7\text{ppm})`$ Brown:2000sj , whereas Höcker et al. obtained the Standard Model (SM) prediction $`a_\mu ^{(\text{SM})}=11659182(6)\times 10^{10}(\pm 0.7\text{ppm})`$ Hocker:2004xc . The difference in values, $`\mathrm{\Delta }a_\mu a_\mu ^{(\text{exp})}a_\mu ^{(\text{SM})}=(26\pm 9.4)\times 10^{10}`$, suggested that the SM does not strictly hold in the low energy region when the difference exceeds the calculation uncertainties of the Hadronic process and measurement error. This difference has been extensively analyzed by various approaches, such as supersymmetry supersymmetric model , lepton flavor violation nonconservation of lepton , extra dimensions Graesser:1999yg , etc. etc . However, there is no conclusive explanation for this observed deviation.
In this context, we attempt to estimate the order of the muon anomalous magnetic moment by using a braneworld model (see Rubakov:2001kp for recent reviews). The general formalism using higher dimensional physics was constructed by R. Sundrum Sundrum:1998sj . This formalism suggests that the SM particles are constrained to live on the world volume of a (3+1)-dimensional hypersurface or a “$`3`$-brane,” while only gravity freely propagates in bulk space-time.
The most important aspect of this theory is that the metric and the vierbein are replaced by the induced metric and the induced vierbein, respectively Sundrum:1998sj ; Akama:1982jy . Hence, higher dimensional gravity can be discussed apart from the usual Kaluza-Klein (KK) theory. In general, the KK modes require periodic limits such as the torus structure in extra dimensions. However, the braneworld scenarios need not have these limits because the configuration of extra dimensions is determined by gravity, the position of branes, the cosmological constant, etc. In this paper, we adopt the factorizable $`6`$-dimensional braneworld model of the ADD type (the name is derived from the paper by N. Arkani-Hamed, S. Dimopoulos, and G. Dvali Arkani-Hamed:1998rs ), which has a theoretical motivation that explains the gauge hierarchy problem, i.e., the reason for the scale of electroweak symmetry breaking being so much smaller than the scales of quantum gravity or grand unification. This model has extra compact spaces, and it finds a simple exact solution to the Einstein equation, including explicit brane sources. Applying Gauss’s law to this model, we have
$`M_{Pl}^2=M_f^4V_{(2)},`$ (1)
where $`M_{Pl}10^{19}`$ $`[\text{GeV}]`$, $`M_f`$, and $`V_{(2)}`$ denote the $`4`$-dimensional Planck mass, $`6`$-dimensional fundamental Planck mass, and volume of the two extra dimensions, respectively. This relation will be crucial in our discussion because it can give $`M_f1`$ \[TeV\] for $`V_{(2)}^{1/2}r=0.1`$ \[mm\]<sup>1</sup><sup>1</sup>1The conversion factor is $`1[\text{GeV}^1]=210^{13}[\text{mm}]`$.. This implies that gravity would be unified with other forces on a TeV scale. On the other hand, there exists an established higher dimensional model, the RS-model Randall:1999ee , which could possibly resolve the hierarchy problem. However, we shall specifically concentrate on the $`6`$-dimensional ADD-type model Sundrum:1998ns ; Carroll:2003db . This involves the two $`3`$-branes (i.e., our world and another world) and the $`U(1)`$ gauge field in the bulk. The model can realize a mechanism that does not require any fine-tuning between the brane tension and bulk parameters; this implies that the brane tension can be freely changed. This type of model is referred to as a self-tuning model Carroll:2003db ; Nilles:2003km , which is one of the simplest models for exploring the braneworld phenomenology and the effects of extra dimensions.
This paper is organized as follows. Section II introduces some basic notations. Section III comprises a brief explanation of the 6-dimensional model and the scaling property of $`4`$-dimensional physics, section IV focuses on the brane-stretching effect and the estimation of the order of muon (g-2), and section V is the conclusion.
## II Setup
The effective theory has presented a picture of the low-energy dynamics of a 3-brane universe, i.e., the SM particle is confined to the braneworld-volume topology as $`𝐌_4`$. Further, only gravity is free to move in bulk space-time with $`d>4`$ dimensions, $`𝐌_4\times 𝐒^{d4}`$ topology, where $`𝐒^d`$ denotes the $`d`$-sphere. The coordinates of bulk space are denoted by $`X^M`$, where the ones on the brane are denoted by $`x^\mu `$ and the extra dimensions by $`y^m`$. The curved bulk coordinate indices, which run over all dimensions, are denoted by uppercase Roman letters beginning from the middle: $`M,N\mathrm{}=0,\mathrm{}d1`$; the indices denoted by Greek letters run over the first four dimensions: $`\mu ,\nu ,\mathrm{}=0,1,2,3`$; and the indices denoted by lowercase Roman letters run over the remaining $`d4`$ dimensions: $`m,n,\mathrm{}=4,\mathrm{}d1`$. The local Lorentz indices in the bulk are similarly denoted: the indices denoted by uppercase Roman letters run over all dimensions: $`A,B\mathrm{}=0,\mathrm{}d1`$; Greek letters run over the first four dimensions: $`\alpha ,\beta \mathrm{}=0,1,2,3`$; the indices denoted by lowercase Roman letters run over the remaining $`d4`$ dimensions: $`a,b\mathrm{}=4,\mathrm{}d1`$ (see Table 1).
The bulk metric $`G_{MN}`$ describes the fundamental gravitational degrees of freedom. The Lorentz metric in the bulk is $`\eta _{AB}`$, and vielbein is $`E_M^A(X)`$. The bulk and Lorentz metrics are related by the following equation:
$`E_M^A(X)\eta _{AB}E_N^B(X)`$ $`=G_{MN}(X)`$ (2)
$`E_M^A(X)G^{MN}(X)E_N^B(X)`$ $`=\eta ^{AB}.`$ (3)
The bulk coordinates occupied by a point $`x`$ on the brane are denoted by $`Y^M(x)`$. However, since the theory has reparametrization invariance, a different parametrization of the surface describing the brane $`xx^{}(x)`$ would lead to the same physics. Therefore, it is necessary to identify the coordinates spanned by the brane with the first four bulk components in order to eliminate the non-physical components from $`Y^M(x)`$. Hence, we choose the gauge fixing condition
$`Y^\mu (x)=x^\mu .`$ (4)
## III Scaling property
We review the $`6`$-dimensional model with two brane sources in two extra dimensions, where the brane and the extra space have an $`𝐌_4`$ and an $`𝐒^2`$ topology, respectively Carroll:2003db ; Aghababaie:2003wz ; Nilles:2003km ; Chen:2000at ; Navarro:2003vw ; Vinet:2004bk ; Garriga:2004tq ; Lee:2004vn ; Mukohyama:2005yw . The total action consists of the $`6`$-dimensional Einstein-Maxwell action and the two brane actions with negative tension. In this model, the stability of bulk geometry and brane fluctuations requires the negative tension brane.
First, in order to obtain a background solution, we discuss the description that does not consider the localized fields on brane and brane fluctuations. The effective action is shown to be as follows:
$`S_{\text{total}}`$ $`=`$ $`S_{\text{branes}}+S_6`$ (5)
$`=`$ $`T_0{\displaystyle d^4x\sqrt{g}}T_1{\displaystyle d^4x\sqrt{g}}+{\displaystyle d^6x\sqrt{G}\left[M_f^4R_6\mathrm{\Lambda }_6\frac{1}{4}F_{MN}^2\right]}.`$
Here, $`T_i`$ $`(i=0,1)`$ denotes the brane tension, $`M_f`$ is the $`6`$-dimensional Planck mass, $`\mathrm{\Lambda }_6`$ is the $`6`$-dimensional cosmological constant, and $`F_{MN}`$ is the $`6`$-dimensional $`2`$-form field strength. We can obtain the $`6`$-dimensional Einstein equation including brane sources by varying the action with respect to the $`6`$-dimensional metric. We consider that a brane is located on a conical singularity in the extra dimensions. Fortunately, in this scenario, we can easily obtain a solution that maintains a $`4`$-dimensional Minkowski space-time, because the equation can split into $`4`$\- and $`2`$-dimensional components. Thus, we obtain the solution<sup>2</sup><sup>2</sup>2The metric signature is diag $`(+,,,,,)`$.
$`ds_6^2=\eta _{\mu \nu }dx^\mu dx^\nu +\gamma _{mn}(y)dy^mdy^n,`$ (6)
and the solution for the equation of motion for $`F_{MN}`$ as
$`F_{mn}=\sqrt{\gamma }B_0ϵ_{mn},`$ (7)
where $`B_0`$ is a constant, $`\gamma `$ is the determinant of $`\gamma _{mn}`$, and $`ϵ_{mn}`$ is a completely antisymmetric tensor, i.e., $`ϵ_{45}=ϵ_{54}=1`$. Solution (7) denotes a magnetic flux through the compactified two extra dimensions.
In this background, the simplest technique to realize the stabilized bulk geometry would be to locate two fixed branes having identical tensions, $`T_0=T_1`$, at opposite poles of the spherical two extra dimensions. This condition can be ensured by imposing a $`𝐙_\mathrm{𝟐}`$ symmetry at the equator Carroll:2003db . Further, using the conformal symmetry, we can then obtain the solution
$`ds_6^2=\eta _{\mu \nu }dx^\mu dx^\nu +a_0^2(d\theta ^2+\alpha ^2\mathrm{sin}^2\theta d\varphi ^2),`$ (8)
if the parameters $`B_0`$ and $`\lambda _6`$ satisfy
$`{\displaystyle \frac{1}{a_0^2}}={\displaystyle \frac{B_0^2}{2M_f^4}}`$ $`,\lambda _6={\displaystyle \frac{B_0^2}{2}}.`$ (9)
These relations are necessary to maintain a $`4`$-dimensional Minkowski space-time and spherical two extra dimensions. The solution has the following relation on the conical singularity;
$`\delta =2\pi (1\alpha )={\displaystyle \frac{T_0}{2M_f^4}},`$ (10)
where $`\delta `$ is the deficit angle of the two extra dimensional sphere, and $`\alpha `$ is a dimensionless fixed parameter, $`0<\alpha <1`$. On the basis of a property in $`2`$-dimensional gravity Deser:tn , the Einstein equation for the extra dimensional component presents a solution that removes a wedge from the sphere and was identified with opposite sides of the wedge. Thus, the $`4`$-dimensional component remains exactly Lorentz invariant because the change in the tension affects only the geometry of the extra dimensions. This means that the tension can be freely changed since there is no fine tuning between bulk parameters and brane tension. This type of model is referred to as a self-tuning model.
In the following, we will briefly describe the manner in which this mechanism affects $`4`$-dimensional physics. The change in $`T_0`$ retains the regular part of the geometry and modifies only the singular part of the geometry, i.e., the deficit angle $`\delta `$ given by (10). This results in a change in the bulk volume related to $`M_{Pl}`$ by (1). Hence, the change in $`T_0`$ signifies a change in $`M_{Pl}`$. Interestingly, a self-tuning model of this type can be constructed only in six dimensions Nilles:2003km .
Subsequently, we focus on fermion $`\psi (x)`$ and gauge field $`A_\mu (x)`$ and ignore scalar field on brane. However, prior to the discussing these behaviors, we should elaborate on a covariant derivative for the fermion. It behaves as a spin 1/2-spinor under the local Lorentz group. Lorentz generators of $`n`$-dimensional spinor representation are usually denoted as:
$`\sigma _{(\alpha \beta )}={\displaystyle \frac{1}{4}}[\gamma _\alpha ,\gamma _\beta ],`$ (11)
where $`\gamma _\alpha `$ represents a set of Dirac matrices satisfying the following condition:
$`\{\gamma _\alpha ,\gamma _\beta \}=2\eta _{\alpha \beta }.`$ (12)
The local Lorentz group on the brane is regarded as an internal $`SO(3,1)`$ group, which connects the Minkowski space with the curved space through the vielbein that satisfy Eqs. (2) and (3). The covariant derivative that maintains the Lorentz and gauge symmetry for $`\psi `$ is
$`D_\mu `$ $`=`$ $`_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{\alpha \beta }\sigma _{(\alpha \beta )},`$ (13)
$`\omega _\mu ^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\alpha \nu }(_\mu e_\nu ^\beta _\nu e_\mu ^\beta )+{\displaystyle \frac{1}{4}}e^{\alpha \nu }e^{\beta \sigma }(_\sigma e_\nu ^\gamma _\nu e_\sigma ^\gamma )e_{\gamma \mu }(\alpha \beta )`$ (14)
Veltman . Thus, the effective brane action is as follows:
$`S_{\text{brane}}`$ $`=`$ $`{\displaystyle d^4x\sqrt{g}\left[T_0+i\overline{\psi }e_\alpha ^\mu \gamma ^\alpha D_\mu \psi m_f\overline{\psi }\psi \frac{1}{4}g^{\mu \rho }g^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }+\mathrm{}\right]},`$ (15)
where the ellipsis represents the higher dimensional interactions that can be constructed with coefficients given by powers of 1/$`M_f`$, and $`m_f`$ is the mass parameter of the fermion in fundamental gravity.
In the following, we discuss only the effect of the brane tension on $`4`$-dimensional physics (see Nilles:2003km for details). The higher dimensional theory that results in a change in $`M_{Pl}`$ generates an effective theory depending on the change in $`M_{Pl}`$. Thus, the $`4`$-dimensional effective action consists of
$`S_{\text{eff}}=M_{Pl}^2(T_0){\displaystyle d^4x\sqrt{g}R_4}+{\displaystyle d^4x\sqrt{g}_4},`$ (16)
where $`M_{Pl}`$ is dependent on $`T_0`$ as follows:
$`M_{Pl}^2(T_0)=\left[1{\displaystyle \frac{T_0}{4\pi M_f^4}}\right]M_{Pl}^2(0),`$ (17)
where $`M_{pl}(0)`$ represents the Planck mass in the absence of branes. When we rescale $`g_{\mu \nu }=\stackrel{~}{g}_{\mu \nu }/\alpha `$, we obtain
$`S_{\text{eff}}=M_{Pl}^2(0){\displaystyle d^4x\sqrt{\stackrel{~}{g}}\stackrel{~}{R}_4}+{\displaystyle d^4x\sqrt{\stackrel{~}{g}}\stackrel{~}{}_4(\alpha )}.`$ (18)
It is obvious that the $`\alpha `$ dependence shifts from the Planck mass to the fields localized on the brane. Hence, after rescaling the fermions as $`\psi =\alpha ^{\frac{3}{4}}\stackrel{~}{\psi }`$ on the basis of (15), we obtain
$`S_4`$ $`=`$ $`{\displaystyle }d^4x\sqrt{\stackrel{~}{g}}[i\overline{\stackrel{~}{\psi }}e_\alpha ^\mu \gamma ^\alpha (_\mu ieA_\mu +{\displaystyle \frac{1}{2}}\omega _\mu ^{\beta \gamma }\sigma _{(\beta \gamma )})\stackrel{~}{\psi }`$ (19)
$`{\displaystyle \frac{m_f}{\sqrt{\alpha }}}\overline{\stackrel{~}{\psi }}\stackrel{~}{\psi }{\displaystyle \frac{1}{4}}g^{\mu \rho }g^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }].`$
Based on the redefinition $`\psi =\stackrel{~}{\psi }`$ and $`g_{\mu \nu }=\stackrel{~}{g}_{\mu \nu }`$, we recognize the action as invariant, except for the mass term. Thus, since $`\alpha `$ is the fixed parameter, we can regard $`m_f/\sqrt{\alpha }`$ as a physical mass $`m`$. As a result, the effect of bulk gravity does not become apparent in the $`4`$-dimensional world. This implies that if fermion is massless, the action becomes scale-invariant, i.e., the scale invariance is broken by fermion mass. The usual field theory also maintains this property. In the next section, we consider the effect of the brane tension and brane fluctuations. The $`4`$-dimensional field theory should be extended to a brane world that maintains this property. Further, we show that the scale transformation is instrumental in restricting the form of the induced metric.
## IV Application to Muon (g-2)
We estimate the muon (g-2) deviation by assuming that brane fluctuations are static in time. The new compensation terms occur through the induced vierbein. This would lead to the possibility of compensating the magnetic moment which has a static property. Under gauge fixing condition (4), the induced metric is as follows:
$`g_{\mu \nu }=\eta _{\mu \nu }+\gamma _{mn}_\mu Y^m_\nu Y^n.`$ (20)
For simplicity, we suppose that off-diagonal components of the $`6`$-bein are zero, as shown below:
$`E_M^A(X)=\left(\begin{array}{cc}\delta _\mu ^\alpha & 0\\ 0& E_m^a\end{array}\right).`$ (23)
In order to obtain the induced vierbein on the brane, we use the following definition Sundrum:1998sj :
$`e_\mu ^\alpha R_A^\alpha E_M^A(X)_\mu Y^M.`$ (24)
Thus, the induced vierbein obtains, up to the second order;
$`e_\mu ^\alpha =\delta _\mu ^\alpha +{\displaystyle \frac{1}{2}}\gamma _{mn}^\alpha Y^m_\mu Y^n+O(ϵ^4).`$ (25)
The expansion of $`\sqrt{g}`$ of induced metric (20) becomes
$`\sqrt{g}`$ $`=`$ $`1{\displaystyle \frac{1}{2}}^\mu Y^m_\mu Y^m+\mathrm{}.`$ (26)
The ellipsis consists of higher dimension terms of $`_\mu Y^m`$ in pairs. When the above expansion is substituted into the minimal brane action
$`S_{\text{brane}}`$ $`=`$ $`{\displaystyle d^4x\sqrt{g}\left[T_0+_\text{4}(g_{\mu \nu })\right]},`$ (27)
we obtain
$`S_{\text{brane}}`$ $`=`$ $`S_{\text{eff}}^{(0)}+S_{\text{eff}}^{(2)}+\mathrm{},`$ (28)
$`S_{\text{eff}}^{(0)}`$ $`=`$ $`{\displaystyle d^4x\left[T_0+_\text{4}(\eta _{\mu \nu })\right]},`$ (29)
$`S_{\text{eff}}^{(2)}`$ $`=`$ $`{\displaystyle d^4x\left[\frac{T_0}{2}_\mu Y^m^\mu Y^m+\frac{1}{2}_\mu Y^m_\nu Y^mT_\text{4}^{\mu \nu }\right]},`$ (30)
where $`T_\text{4}^{\mu \nu }`$ is the conserved energy-momentum tensor of matter fields evaluated in the $`4`$-dimensional Minkowski space-time. Considering the canonically normalized condition for $`_\mu Y^m`$ in (30), we can put
$`_\mu Y^m^\mu Y^m{\displaystyle \frac{1}{T_0}}_\mu Y^m^\mu Y^m.`$ (31)
$`Y^m`$ is considered as the Nambu-Goldstone mode associated with the spontaneous isometry breaking due to the presence of the brane in bulk Sundrum:1998sj Sundrum:1998ns Hisano:1999bn ; Bando:1999di ; Kugo:1999mf ; Dobado:2000gr . Before discussing muon (g-$`2`$), the relation between the brane fluctuations and the scaling property mentioned in section III should be noted. On the assumption that the change in $`Y^m`$ is static in time, the induced metric becomes
$`g_{\mu \nu }=\left(\begin{array}{cc}1& 0\\ 0& \eta _{ij}+{\displaystyle \frac{1}{T_0}}\gamma _{mn}_iY^m_jY^n\end{array}\right),`$ (34)
where $`i,j=1,2,3`$: indices are raised and lowered by the Euclidean metric $`\delta _{ij}=\eta _{ij}`$. The $`4`$-dimensional field theory is scale-invariant for massless fermions and gauge fields, but not for massive fermions. We consider that the braneworld would preserve this property. Thus, induced metric (34) requires the following rescaling for $`\eta _{\mu \nu }\eta _{\mu \nu }/\alpha `$:
$`g_{\mu \nu }`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ 0& \eta _{ij}+{\displaystyle \frac{1}{T_0}}\gamma _{mn}_iY^m_jY^n\end{array}\right)`$ (39)
$``$ $`{\displaystyle \frac{1}{\alpha }}\left(\begin{array}{cc}1& 0\\ 0& \eta _{ij}+{\displaystyle \frac{1}{T_0}}\gamma _{mn}_iY^m_jY^n\end{array}\right).`$
However, the existence of the brane breaks the isometry symmetry. It denotes that the $`6`$-dimensional bulk is separated into the $`4`$-dimensional branes and $`2`$-dimensional extra dimensions. This implies that $`\gamma _{mn}`$ does not have the abovementioned transformation because we can rescale $`G_{MN}G_{MN}/\alpha `$ if and only if no brane exists in the bulk. Therefore, in order to recover the scaling property, we restrict the form of $`g_{ij}`$:
$`g_{ij}=\eta _{ij}+\eta _{ij}{\displaystyle \frac{1}{T_0}}H^2M_f^2,`$ (40)
where $`H`$ has a mass dimension of $`1`$. This form guarantees that the $`\alpha `$ dependence changes from $`M_{Pl}(\alpha )`$ into the fermion mass in the same way as action (19). Subsequently, we present a solution $`Y^m`$ that satisfies (40). The $`Y^m`$ equation of motion derived from effective action (30) is written as
$`_\mu \left[^\mu Y^m+{\displaystyle \frac{1}{T_0}}_\nu Y^mT_\text{4}^{\mu \nu }\right]=0.`$ (41)
When we introduce the dimensionless coordinate $`Ex^i`$ that characterizes the physical process at energy $`E`$, we parametrize $`Y^m`$ as follows (see Fig.1):
$`Y^m(\text{x})=Y_0^m+M_f\stackrel{~}{e}_i^mEx^i,`$ (42)
where $`Y_0^m`$ is a constant and the basis vectors
$`{\displaystyle \frac{Y^m}{x^i}}=M_fE\stackrel{~}{e}_i^m`$ (43)
satisfy the completeness relation
$`\gamma _{mn}{\displaystyle \frac{Y^m}{x^i}}{\displaystyle \frac{Y^n}{x^j}}=M_f^2E^2\eta _{ij},`$ (44)
i.e.,
$`\gamma _{mn}\stackrel{~}{e}_i^m\stackrel{~}{e}_j^n=\eta _{ij}.`$ (45)
The coordinate $`Y^m`$ (42) satisfies (41) and maintains (40) as $`H=E`$. This implies that the spatial part of the brane is stretched due to brane fluctuations, whose magnitude depends on the energy scale of the physical process. This is physically plausible because under general relativity, space-time is not rigid but dynamical. In addition, solution (42) is consistent with the general covariance of general relativity. Substituting (42) into spin connection (14) via induced vierbein (25), we can directly obtain
$`\omega _\mu ^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu ^\beta Y^m^\alpha Y^m{\displaystyle \frac{1}{2}}_\mu ^\alpha Y^m^\beta Y^m`$ (46)
$`=`$ $`0.`$
The vanishing of the spin connection denotes that the equation of motion for a fermion agrees with laws of special relativity. Therefore, solution (42) supports Lorentz symmetry.
In the following, we will see that the brane-stretching effect generates the suitable order for muon (g-$`2`$). The variation of action (27) with respect to $`\overline{\psi }`$ yields the equation of motion:
$`\left[ie_\alpha ^\mu \gamma ^\alpha \left(_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{\beta \gamma }\sigma _{(\beta \gamma )}\right)m\right]\psi =0.`$ (47)
Then, we perform a nonrelativistic approximation, i.e., the Schrödinger approximation. This is demonstrated in Appendix A by using the static induced metric (34) and solution (42). Since the two extra dimensions give the relation
$`M_{pl}^2=4\pi a_0^2\alpha M_f^4`$ (48)
by (1), we obtain the anomalous magnetic moment:
$`a_\mu `$ $`={\displaystyle \frac{1}{T_0}}E^2M_f^2`$ (49)
$`={\displaystyle \frac{1}{4\pi M_f^4(1\alpha )}}E^2M_f^2`$ (50)
$`={\displaystyle \frac{a_0^2\alpha }{M_{pl}^2(1\alpha )}}E^2M_f^2.`$ (51)
Finally, since we are interested in the physics at the muon scale ($`E106[\text{MeV}]`$) and $`M_f1[\text{TeV}]`$ for $`a_00.1`$ $`[\text{mm}]`$, we obtain the following:
$`a_\mu `$ $`{\displaystyle \frac{\alpha }{1\alpha }}\left({\displaystyle \frac{0.1[\text{mm}]}{10^{19}[\text{GeV}]}}\right)^2\times \left(106[\text{MeV}]1[\text{TeV}]\right)^2`$ (52)
$`={\displaystyle \frac{\alpha }{1\alpha }}10^{10}.`$ (53)
This result almost reproduces the deviation of the muon (g-2) measurement, except for the previous dimensionless factor. $`\alpha `$ may be determined by future studies on the self-tuning mechanism Aghababaie:2003wz ; Nilles:2003km ; Chen:2000at ; Navarro:2003vw ; Vinet:2004bk ; Garriga:2004tq ; Lee:2004vn ; Mukohyama:2005yw . However, it is important that we consider its behavior in the bound $`0<\alpha <1`$ because it is possible that $`\alpha `$ has an extreme value. If $`\alpha 1`$, muon (g-2) has a value greater than the experimental result. On the contrary, if $`\alpha 0`$, muon (g-2) has a small value. Moreover, it generates a large hierarchy between the fundamental parameter $`M_f`$ and $`m_f`$. Consequently, when $`\alpha `$ has a moderate value, this model would be capable of reproducing $`\mathrm{\Delta }a_\mu a_\mu ^{(\text{exp})}a_\mu ^{(\text{SM})}=(26\pm 9.4)\times 10^{10}`$.
As a side remark, from recent astrophysical research, it is known that the bounds on the mass of KK-gravitons Hall:1999mk impose much tighter constraints on the radius of Large extra dimension. These suggest the exclusion of the TeV scale gravity. This indicates that we need to consider much more than the TeV scale. However, even in this case, if the order of $`a_0`$ is smaller than $`0.1`$ $`[\text{mm}]`$, the region $`\alpha 1`$ can give the appropriate (g-2) value if $`\alpha `$ is suitably selected.
## V Conclusion
This paper has presented a new approach according to which brane fluctuations compensate for the muon anomalous magnetic moment. The most important fact to be considered is that we have obtained a new potential term for the magnetic moment based on the assumption that brane fluctuations are static in time. This method reflects the effect of a novel classical contribution, namely, brane-stretching effect due to brane fluctuations, which is not based on the previously studied KK-gravitons Graesser:1999yg . In particular, we would obtain a suitable order for $`a_\mu `$ in the $`6`$-dimensional model. This implies that the SM is consistently extended to the braneworld model that maintains the usual scaling property and Lorentz invariance for fermion. In future research, we should promote the investigation of $`a_\mu `$ by using the metric constructed by other higher dimensional models. Moreover, we can expect that the brane-stretching effect will evolve into different configurations in a very high energy. This may be related to the Lorentz violation Kostelecky:2000mm . Since our study leaves a lot of issues to be discussed further, we are confident that this will be a crucial subject on which further research should be conducted.
###### Acknowledgements.
This study was partly supported by Iwanami Fūjyukai.
*
## Appendix A The schödinger approximation
In this appendix, we demonstrate the non-relativistic approximation for fermion in the action (27), and drive the magnetic moment. Varying the action with respect to $`\overline{\psi }`$, we obtain the equation of motion:
$`\left[ie_\alpha ^\mu \gamma ^\alpha \left(_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{\beta \gamma }\sigma _{(\beta \gamma )}\right)m\right]\psi =0,`$ (54)
where $`e_\alpha ^\mu `$ is represented by (25). Operating on
$`\left[ie_\alpha ^\mu \gamma ^\alpha \left(_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{\beta \gamma }\sigma _{(\beta \gamma )}\right)+m\right]`$ (55)
from the left, we get
$`[g^{\mu \nu }(_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{(\stackrel{´}{\beta }\stackrel{´}{\gamma })}\sigma _{(\stackrel{´}{\beta }\stackrel{´}{\gamma })})(_\nu ieA_\nu {\displaystyle \frac{1}{2}}\omega _\nu ^{(\beta \gamma )}\sigma _{(\beta \gamma )})+ie\sigma ^{(\stackrel{´}{\alpha }\alpha )}e_{\stackrel{´}{\alpha }}^{\stackrel{´}{\mu }}e_\alpha ^\mu F_{\mu \stackrel{´}{\mu }}`$
$`+{\displaystyle \frac{1}{2}}\sigma ^{(\stackrel{´}{\alpha }\alpha )}e_{\stackrel{´}{\alpha }}^{\stackrel{´}{\mu }}e_\alpha ^\mu \mathrm{\Omega }_{\mu \stackrel{´}{\mu }}^{\beta \gamma }\sigma _{\beta \gamma }+\gamma ^{\stackrel{´}{\alpha }}\gamma ^\alpha e_{\stackrel{´}{\alpha }}^{\stackrel{´}{\mu }}_{\stackrel{´}{\mu }}e_\alpha ^\mu (_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{(\beta \gamma )}\sigma _{(\beta \gamma )})+m^2]\psi =0`$ (56)
by using the formula
$`\gamma ^{\stackrel{´}{\alpha }}\gamma ^\alpha =\eta ^{\stackrel{´}{\alpha }\alpha }+2\sigma ^{(\stackrel{´}{\alpha }\alpha )},`$ (57)
and
$`\sigma ^{(\stackrel{´}{\alpha }\alpha )}e_{\stackrel{´}{\alpha }}^{\stackrel{´}{\mu }}e_\alpha ^\mu \left(_{\stackrel{´}{\mu }}ieA_{\stackrel{´}{\mu }}{\displaystyle \frac{1}{2}}\omega _\mu ^{(\stackrel{´}{\beta }\stackrel{´}{\gamma })}\sigma _{(\stackrel{´}{\beta }\stackrel{´}{\gamma })}\right)\left(_\mu ieA_\mu {\displaystyle \frac{1}{2}}\omega _\mu ^{\beta \gamma }\sigma _{(\beta \gamma )}\right)`$
$`={\displaystyle \frac{1}{2}}\sigma ^{(\stackrel{´}{\alpha }\alpha )}e_{\stackrel{´}{\alpha }}^{\stackrel{´}{\mu }}e_\alpha ^\mu \left(ieF_{\mu \stackrel{´}{\mu }}+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_{\mu \stackrel{´}{\mu }}^{\beta \gamma }\sigma _{\beta \gamma }\right)`$ (58)
where $`F_{\mu \nu }_{[\mu }A_{\nu ]}`$ and $`\mathrm{\Omega }_{\mu \nu }^{\beta \gamma }_{[\mu }\omega _{\nu ]}^{\beta \gamma }`$. Further, given the assumption that the change in $`Y^m(x)`$ is static in time, we obtain the induced metric
$`g_{\mu \nu }=\left(\begin{array}{cc}1& 0\\ 0& \eta _{ij}+\gamma _{mn}_iY^m_jY^n\end{array}\right).`$ (61)
Thus, rewriting
$`p^\mu =i^\mu ,A_\mu =(\varphi ,\stackrel{}{A}),`$ (62)
the (56) transforms into
$`\left(iE+ie\varphi {\displaystyle \frac{1}{2}}\omega _0^{\beta \gamma }\sigma _{(\beta \gamma )}\right)^2\psi `$ $`=`$ $`[g^{ij}(ip_iieA_i{\displaystyle \frac{1}{2}}\omega _i^{(\stackrel{´}{\beta }\stackrel{´}{\gamma })}\sigma _{(\stackrel{´}{\beta }\stackrel{´}{\gamma })})(ip_jieA_j{\displaystyle \frac{1}{2}}\omega _j^{(\beta \gamma )}\sigma _{(\beta \gamma )})`$ (63)
$`ie\sigma ^{ij}e_i^ke_j^lF_{kl}{\displaystyle \frac{1}{2}}\sigma ^{ij}e_i^ke_j^l\mathrm{\Omega }_{kl}^{\beta \gamma }\sigma _{\beta \gamma }`$
$`+\gamma ^i\gamma ^je_i^k_ke_j^l(ip_lieA_l{\displaystyle \frac{1}{2}}\omega _i^{(\beta \gamma )}\sigma _{\beta \gamma })+m^2]\psi `$
where $`i,j,k,l=1,2,3`$ and $`E`$ represents the energy eigenvalue. Putting $`E=m+W`$ where $`m`$ is the rest energy, the L.H.S of (63) is as follows:
$`\text{L}.\text{H}.\text{S}`$ $`=`$ $`[m^2+2m(W+e\varphi +{\displaystyle \frac{i}{2}}\omega _0^{\beta \gamma }\sigma _{(\beta \gamma )})`$ (64)
$`+(W+e\varphi +{\displaystyle \frac{i}{2}}\omega _0^{\beta \gamma }\sigma _{(\beta \gamma )})^2]\psi .`$
In addition, we assume that $`Wm`$, i.e., the energy due to a magnetic field is extremely small. In this case, dividing both the L.H.S and R.H.S of (63) by $`2m`$ so as to ignore the last term in (64), we obtain
$`W\psi `$ $`=`$ $`{\displaystyle \frac{1}{2m}}[g^{ij}(ip_iieA_i{\displaystyle \frac{1}{2}}\omega _i^{(\stackrel{´}{\beta }\stackrel{´}{\gamma })}\sigma _{(\beta \gamma )})(ip_jieA_j{\displaystyle \frac{1}{2}}\omega _j^{\beta \gamma }\sigma _{(\beta \gamma )})`$ (65)
$`ie\sigma ^{ij}e_i^ke_j^lF_{kl}{\displaystyle \frac{1}{2}}\sigma ^{ij}e_i^ke_j^l\mathrm{\Omega }_{kl}^{\beta \gamma }\sigma _{\beta \gamma }+\gamma ^i\gamma ^je_i^k_ke_j^l(ip_lieA_l+{\displaystyle \frac{1}{2}}\omega _j^{\beta \gamma }\sigma _{(\beta \gamma )})]\psi `$
$`\left(e\varphi +{\displaystyle \frac{i}{2}}\omega _0^{\beta \gamma }\sigma _{(\beta \gamma )}\right)\psi .`$
This is the eigenvalue equation for a charged particle in a magnetic field and gravity. From this equation, we can ascertain the energy shift term, which is produced by the following interaction:
$`{\displaystyle \frac{W}{H_i}}H_i={\displaystyle \frac{ie}{2m}}\sigma ^{ij}e_i^ke_j^lF_{kl}.`$ (66)
Hence, when evaluating Eq. (66) by using $`e_i^k`$ which is the inverse of Eq. (25) and the solution (42), we obtain
$`{\displaystyle \frac{W}{H_i}}H_i={\displaystyle \frac{e}{2m}}\left[\left(2+2{\displaystyle \frac{E^2M_f^2}{T_0}}\right){\displaystyle \frac{\stackrel{}{\sigma }}{2}}\stackrel{}{H}\right]`$ (67)
where $`F_{23}=F_{32}=H_1`$, $`F_{31}=F_{13}=H_2`$, and $`F_{12}=F_{21}=H_3`$. The parenthesis of the term proportional to $`\stackrel{}{\sigma }/2\stackrel{}{H}`$ represents the magnetic moment.
|
warning/0506/hep-ph0506192.html
|
ar5iv
|
text
|
# Large Lepton Mixing and Nonsymmetric Mass Matrices with Flavor 2 ↔ 3 Symmetry
## I Introduction
In order to explain a nearly bimaximal lepton mixing $`(\mathrm{sin}^22\theta _{12}1`$, $`\mathrm{sin}^22\theta _{23}1)`$ observed from neutrino oscillation experiments skamioka , mass matrices with various structures have been investigated in the literature. For example, mass matrices with texture zeros fritzsch -Ramond , with a flavor $`23`$ symmetry Fukuyama Kaneko , and so on have been proposed. Recently we have proposed Matsuda3 a following nonsymmetric mass matrix model for all quarks and leptons based on an extended flavor $`23`$ symmetry with one phase:
$$M_f=\left(\begin{array}{ccc}0\hfill & a_fe^{i\varphi _f}\hfill & a_f\hfill \\ a_f^{}e^{i\varphi _f}\hfill & b_fe^{2i\varphi _f}\hfill & (1\xi _f)b_f\hfill \\ a_f^{}\hfill & (1\xi _f)b_f\hfill & b_f\hfill \end{array}\right),\text{(}f=u,d,\nu \text{}e\text{, (}D\text{, and }R\text{))}$$
(1)
where $`a_f`$, $`b_f`$, $`\xi _f`$, and $`a_f^{}`$ are real parameters and $`\varphi _f`$ is a phase parameter. Here, $`M_u`$, $`M_d`$, $`M_\nu `$, and $`M_e`$ are mass matrices for up quarks ($`u,c,t`$), down quarks ($`d,s,b`$), neutrinos ($`\nu _e,\nu _\mu ,\nu _\tau `$) and charged leptons ($`e,\mu ,\tau `$), respectively. The mass matrices $`M_D`$ and $`M_R`$ are, respectively, the Dirac and the right-handed Majorana neutrino mass matrices, which are included in the model if we assume the seesaw mechanism Yanagida for neutrino masses. In this model, we assume that all the mass matrices for quarks and leptons have this common structure, which is against the conventional picture that the mass matrix forms in the quark sector take somewhat different structures from those in the lepton sector. In our previous works Matsuda3 , we have pointed out that this structure leads to reasonable values for the Cabibbo–Kobayashi–Maskawa (CKM) CKM quark mixing, if we use a specific assignment of the quark masses.
In this paper, we shall discuss the Maki-Nakagawa-Sakata-Pontecorv (MNSP) MNSP lepton mixing of the model by assuming that neutrinos are the Majorana particles. Under this assumption, the neutrino mass matrix $`M_\nu `$ should be symmetric. Namely, we further assume
$$a_\nu =a_\nu ^{}.$$
(2)
In the scenario that the neutrino mass matrix is constructed via the seesaw mechanism, i.e. $`M_\nu =M_D^TM_R^1M_D`$, the structure of $`M_\nu `$ mentioned above is alternatively realized by using the following two assumptions: (i) The mass matrices $`M_D`$ and $`M_R`$ have the same extended flavor 2 $``$ 3 symmetry in Eq.(1) with identical phase parameters, i.e. $`\varphi _D=\varphi _R\varphi _\nu `$. (ii) $`M_D`$ and $`M_R`$ are proportional to each other, except for their (2,1) and (3,1) elements.
On the other hand, $`M_e`$ is assumed to have the above nonsymmetric structure given in Eq. (1). Namely, in this paper, the mass matrices $`M_e`$ and $`M_\nu `$ are assumed to have the following forms:
$`M_e`$ $`=`$ $`\left(\begin{array}{ccc}0\hfill & a_ee^{i\varphi _e}\hfill & a_e\hfill \\ a_e^{}e^{i\varphi _e}\hfill & b_ee^{2i\varphi _e}\hfill & (1\xi _e)b_e\hfill \\ a_e^{}\hfill & (1\xi _e)b_e\hfill & b_e\hfill \end{array}\right),`$ (6)
$`M_\nu `$ $`=`$ $`\left(\begin{array}{ccc}0\hfill & a_\nu e^{i\varphi _\nu }\hfill & a_\nu \hfill \\ a_\nu e^{i\varphi _\nu }\hfill & b_\nu e^{2i\varphi _\nu }\hfill & (1\xi _\nu )b_\nu \hfill \\ a_\nu \hfill & (1\xi _\nu )b_\nu \hfill & b_\nu \hfill \end{array}\right),`$ (10)
where $`\varphi _e`$ and $`\varphi _\nu `$ are phase parameters.
This article is organized as follows. In Sec. II, we discuss the diagonalization of mass matrix of our model. The analytical expressions of the lepton mixings and phases of the model are given in Sec. III. Sec. IV is devoted to a summary.
## II Diagonalization of Mass matrix
### II.1 mass matrix for charged leptons
The diagonalization of the mass matrix for the charged leptons $`M_e`$ can be done similarly to the case of the mass matrices for up and down quarks. These mass matrices have common structure given by $`M_f`$ in Eq. (1). Thus let us present a method for diagonalization of $`M_f`$, which is treated in Ref Matsuda3 in details.
First, let us mention free parameters in the mass matrix. There are five parameters, $`a_f`$, $`a_f^{}`$, $`b_f`$, $`\xi _f`$, and $`\varphi _f`$ in $`M_f`$. Even if we fix three eigenvalues, $`m_{if}`$, of $`M_f`$ by the observed fermion masses, there still exist two free parameters. As the free parameters, let us choose a parameter $`\alpha _f`$ defined by
$`\alpha _f{\displaystyle \frac{a_f^{}}{a_f}},`$ (11)
and a phase parameter $`\eta _f`$ defined in Fig. 1 independently of mass eigenvalues $`m_{if}`$. Note that for the neutrino mass matrix $`M_\nu `$ we assume $`\alpha _\nu =1`$ as mentioned above.
Second, let us discuss an unitary matrix $`U_{Lf}`$ which diagonalizes $`M_fM_f^{}`$. The explicit expression of $`U_{Lf}`$ depends on the following three types of assignment for $`m_{if}`$:
(i) Type A:
In this type, the mass eigenvalues $`|m_{1f}|`$, $`m_{2f}`$, and $`m_{3f}`$ of $`M_f`$ are allocated to the masses of the first, second, and third generations, respectively. (i.e. $`|m_{1f}|`$$``$$`m_{2f}`$$``$$`m_{3f}`$.) In this type, the $`M_fM_f^{}`$ is diagonalized as
$$U_{Lf}^{}M_fM_f^{}U_{Lf}=\text{diag}(m_{1f}^2,m_{2f}^2,m_{3f}^2),$$
(12)
by an unitary matrix $`U_{Lf}`$ given by
$$U_{Lf}=P_f^{}O_f.$$
(13)
Here $`P_f`$ is the diagonal phase matrix expressed as
$$P_f=\text{diag}(1,e^{i(\varphi _f\phi _f)},e^{i\phi _f}),$$
(14)
where $`\phi _f`$ and $`\varphi _f`$ are given by
$`\mathrm{cos}\phi _f`$ $`=`$ $`{\displaystyle \frac{|X_f|m_{3f}\mathrm{cos}\eta _f}{\sqrt{|X_f|^2+m_{3f}^22m_{3f}|X_f|\mathrm{cos}\eta _f}}},`$ (15)
$`\mathrm{cos}\varphi _f`$ $`=`$ $`{\displaystyle \frac{|X_f|^2m_{3f}^2}{\sqrt{\left(|X_f|^2+m_{3f}^2\right)^24m_{3f}^2|X_f|^2\mathrm{cos}^2\eta _f}}}.`$ (16)
Here $`X_f`$ is defined by $`X_fb_f+(1\xi _f)b_fe^{i\varphi _f}|X_f|e^{i\phi _f}`$, and $`|X_f|`$ is expressed in term of $`\alpha _f`$ and $`m_{if}`$ as
$$|X_f|^2=m_{1f}^2+m_{2f}^2|m_{1f}|m_{2f}\left(\frac{1+\alpha _f^2}{\alpha _f}\right).$$
(17)
In Eq. (13), $`O_f`$ is the orthogonal matrix given by
$$O_f\left(\begin{array}{ccc}c_f& s_f& 0\\ \frac{s_f}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ \frac{s_f}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right),$$
(18)
where
$$c_f=\sqrt{\frac{m_{2f}^2\frac{|m_{1f}|m_{2f}}{\alpha _f}}{m_{2f}^2m_{1f}^2}},s_f=\sqrt{\frac{\frac{|m_{1f}|m_{2f}}{\alpha _f}m_{1f}^2}{m_{2f}^2m_{1f}^2}}.$$
(19)
It should be noted that the mixing angles are functions of only $`\alpha _f`$, since the $`m_{if}`$ is fixed by the experimental fermion mass values. We find from Eq. (16) that $`\varphi _f\pm \pi `$ for $`m_{1f}^2m_{2f}^2m_{3f}^2`$ in this type A assignment.
(ii) Type B:
In this type, the mass eigenvalues $`|m_{1f}|`$, $`m_{3f}`$, and $`m_{2f}`$ are allocated to the masses of the first, second, and third generations, respectively. (i.e. $`|m_{1f}|`$$``$$`m_{3f}`$$``$$`m_{2f}`$.) The $`M_fM_f^{}`$ is diagonalized as
$$U_{Lf}^{}M_fM_f^{}U_{Lf}=\text{diag}(m_{1f}^2,m_{3f}^2,m_{2f}^2).$$
(20)
by an unitary matrix $`U_{Lf}`$ given by
$$U_{Lf}=P_f^{}O_f^{}.$$
(21)
Here $`O_f^{}`$ is obtained from $`O_f`$ by exchanging the second row for the third one as
$$O_f^{}\left(\begin{array}{ccc}c_f& 0& s_f\\ \frac{s_f}{\sqrt{2}}& \frac{1}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}\\ \frac{s_f}{\sqrt{2}}& \frac{1}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}\end{array}\right).$$
(22)
(iii) Type C:
In this type, the mass eigenvalues $`m_{3f}`$, $`|m_{1f}|`$, and $`m_{2f}`$ are allocated to the masses of the first, second, and third generations, respectively. (i.e. $`m_{3f}`$$``$$`|m_{1f}|`$$``$$`m_{2f}`$.) In this type, we have
$$U_{Lf}^{}M_fM_f^{}U_{Lf}=\text{diag}(m_{3f}^2,m_{1f}^2,m_{2f}^2),$$
(23)
where
$$U_{Lf}=P_f^{}O_f^{\prime \prime }.$$
(24)
Here, the orthogonal matrix $`O_f^{\prime \prime }`$ is given by
$$O_f^{\prime \prime }\left(\begin{array}{ccc}0& c_f& s_f\\ \frac{1}{\sqrt{2}}& \frac{s_f}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}& \frac{s_f}{\sqrt{2}}& \frac{c_f}{\sqrt{2}}\end{array}\right).$$
(25)
This type is not so useful to get the reasonable lepton mixing values.
### II.2 mass matrix for neutrinos
Since the mass matrix for the Majorana neutrinos $`M_\nu `$ is symmetric, $`M_\nu `$ is diagonalized as follows depending on the following three types of assignments for the neutrino mass $`m_i`$:
(i) Type A:
In this type, the mass eigenvalues $`m_1`$, $`m_2`$, and $`m_3`$ of $`M_\nu `$ are allocated to the masses of the first, second, and third generations, respectively. In this type, $`M_\nu `$ is diagonalized as
$$U_\nu ^{}M_\nu U_\nu ^{}=\text{diag}(m_1,m_2,m_3),$$
(26)
where $`m_i(i=1,2,\text{and}3)`$ are real positive neutrino masses. The unitary matrix $`U_\nu `$ is described as
$$U_\nu =P_\nu ^{}O_\nu Q_\nu .$$
(27)
Here, in order to make the neutrino masses $`m_i`$ to be real positive, we introduce a diagonal phase matrix $`Q_\nu `$ defined by
$$Q_\nu \text{diag}(e^{i(\phi _\nu \pi )/2},e^{i(\phi _\nu )/2},e^{i(\phi _\nu \eta _\nu +\pi )/2}).$$
(28)
The diagonal phase matrix $`P_\nu `$ and the orthogonal matrix $`O_\nu `$ are obtained from Eqs. (14) – (17) and (18) – (19) with $`f=\nu `$ by replacing $`|m_{1f}|`$, $`m_{2f}`$, and $`m_{3f}`$ with $`m_1`$, $`m_2`$, and $`m_3`$, respectively and by setting $`\alpha _\nu =1`$.
(ii) Type B:
In this type, the mass eigenvalues $`m_1`$, $`m_3`$, and $`m_2`$ are allocated to the masses of the first, second, and third generations, respectively. In this type, $`M_\nu `$ is diagonalized as
$$U_\nu ^{}M_\nu U_\nu ^{}=\text{diag}(m_1,m_3,m_2).$$
(29)
The unitary matrix $`U_\nu `$ is described as
$$U_\nu =P_\nu ^{}O_\nu ^{}Q_\nu ^{}.$$
(30)
Here the diagonal phase matrix $`Q_\nu ^{}`$ is defined by
$$Q_\nu ^{}\text{diag}(e^{i(\phi _\nu \pi )/2},e^{i(\phi _\nu \eta _\nu +\pi )/2},e^{i(\phi _\nu )/2}).$$
(31)
The orthogonal matrix $`O_\nu ^{}`$ is obtained from Eqs. (18) and (19) with $`f=\nu `$ by replacing $`|m_{1f}|`$, $`m_{2f}`$, and $`m_{3f}`$ with $`m_1`$, $`m_3`$, and $`m_2`$, respectively and by setting $`\alpha _\nu =1`$.
(iii) Type C:
In this type, the mass eigenvalues $`m_3`$, $`m_1`$, and $`m_2`$ are allocated to the masses of the first, second, and third generations, respectively. In this type, $`M_\nu `$ is diagonalized as
$$U_\nu ^{}M_\nu U_\nu ^{}=\text{diag}(m_3,m_1,m_2).$$
(32)
The unitary matrix $`U_\nu `$ is described as
$$U_\nu =P_\nu ^{}O_\nu ^{\prime \prime }Q_\nu ^{\prime \prime }.$$
(33)
Here the diagonal phase matrix $`Q_\nu ^{\prime \prime }`$ is defined by
$$Q_\nu ^{\prime \prime }\text{diag}(e^{i(\phi _\nu \eta _\nu +\pi )/2},e^{i(\phi _\nu \pi )/2},e^{i(\phi _\nu )/2}).$$
(34)
The orthogonal matrix $`O_\nu ^{\prime \prime }`$ is obtained from Eqs. (18) and (19) with $`f=\nu `$ by replacing $`|m_{1f}|`$, $`m_{2f}`$, and $`m_{3f}`$ with $`m_3`$, $`m_1`$, and $`m_2`$, respectively and by setting $`\alpha _\nu =1`$.
These types B and C are not so useful to get the reasonable lepton mixing values.
## III MNSP lepton mixing matrix
Now let us discuss the MNSP lepton mixing matrix of the model by taking the type A, the type B, and the type C assignments for charged leptons and neutrinos. We find that the assignment that is consistent with the present experimental data is only one case. Namely, the case in which type B assignment for charged leptons and type A for neutrinos are taken. The other possible cases fail to reproduce consistent lepton mixing. In this case, we obtain the MNSP lepton mixing matrix $`U`$ as follows.
$`U`$ $`=`$ $`U_{Le}^{}U_\nu =O_e^TP_eP_\nu ^{}O_\nu Q_\nu =O_e^TPO_\nu Q_\nu `$ (38)
$`=`$ $`\left(\begin{array}{ccc}c_e^{}c_\nu +\rho _\nu s_e^{}s_\nu & c_e^{}s_\nu \rho _\nu s_e^{}c_\nu & \sigma _\nu s_e^{}\\ \sigma _\nu s_\nu & \sigma _\nu c_\nu & \rho _\nu \\ s_e^{}c_\nu \rho _\nu c_e^{}s_\nu & s_e^{}s_\nu +\rho _\nu c_e^{}c_\nu & \sigma _\nu c_e^{}\end{array}\right)Q_\nu ,`$
where
$`s_e^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\frac{|m_e|m_\tau }{\alpha _e}m_e^2}{m_\tau ^2m_e^2}}},c_e^{}=\sqrt{{\displaystyle \frac{m_\tau ^2\frac{|m_e|m_\tau }{\alpha _e}}{m_\tau ^2m_e^2}}},`$
$`s_\nu `$ $`=`$ $`\sqrt{{\displaystyle \frac{|m_1|}{m_2+|m_1|}}},c_\nu =\sqrt{{\displaystyle \frac{m_2}{m_2+|m_1|}}}.`$ (39)
Here the phase matrix $`Q_\nu `$ is shown in Eq. (28), and we have put
$$PP_eP_\nu ^{}\text{diag}(1,e^{i\delta _{\nu 2}},e^{i\delta _{\nu 3}}).$$
(40)
The orthogonal matrices $`O_\nu `$ and $`O_e^{}`$ are obtained from Eq. (18) and Eq. (22), respectively. Here we denote the lepton masses $`(m_{1f},m_{3f},m_{2f})`$ as $`(m_e,m_\tau ,m_\mu )`$ for $`f=e`$, and as $`(m_1,m_2,m_3)`$ for $`f=\nu `$. Note also that $`\alpha _\nu =1`$.
The parameters $`\rho _\nu `$ and $`\sigma _\nu `$ in Eq. (38) are defined by
$`\rho _\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}(e^{i\delta _{\nu 3}}+e^{i\delta _{\nu 2}})=\mathrm{cos}\left({\displaystyle \frac{\delta _{\nu 3}\delta _{\nu 2}}{2}}\right)\mathrm{exp}i\left({\displaystyle \frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}}\right),`$ (41)
$`\sigma _\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}(e^{i\delta _{\nu 3}}e^{i\delta _{\nu 2}})=\mathrm{sin}\left({\displaystyle \frac{\delta _{\nu 3}\delta _{\nu 2}}{2}}\right)\mathrm{exp}i\left({\displaystyle \frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}}+{\displaystyle \frac{\pi }{2}}\right).`$ (42)
Note that the phases of $`\rho _\nu `$ and $`\sigma _\nu `$ are
$`\text{arg}\rho _\nu `$ $`=`$ $`\{\begin{array}{cc}\frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}\hfill & \text{ for }\mathrm{cos}\left(\frac{\delta _{\nu 3}\delta _{\nu 2}}{2}\right)>0\hfill \\ \frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}+\pi \hfill & \text{ for }\mathrm{cos}\left(\frac{\delta _{\nu 3}\delta _{\nu 2}}{2}\right)<0\hfill \end{array},`$ (45)
$`\text{arg}\sigma _\nu `$ $`=`$ $`\{\begin{array}{cc}\frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}+\frac{\pi }{2}\hfill & \text{ for }\mathrm{sin}\left(\frac{\delta _{\nu 3}\delta _{\nu 2}}{2}\right)>0\hfill \\ \frac{\delta _{\nu 3}+\delta _{\nu 2}}{2}\frac{\pi }{2}.\hfill & \text{ for }\mathrm{sin}\left(\frac{\delta _{\nu 3}\delta _{\nu 2}}{2}\right)<0\hfill \end{array}.`$ (48)
By using Eqs. (40) and (14), the phases $`\delta _{\nu 2}`$ and $`\delta _{\nu 3}`$ in our model are given by
$`\delta _{\nu 2}`$ $`=`$ $`\phi _\nu \phi _e(\varphi _\nu \varphi _e),`$ (49)
$`\delta _{\nu 3}`$ $`=`$ $`\phi _\nu \phi _e.`$ (50)
Here the phases $`\varphi _e`$, $`\phi _e`$, $`\varphi _\nu `$, and $`\phi _\nu `$ are expressed as
$`\mathrm{cos}\varphi _e`$ $`=`$ $`{\displaystyle \frac{|X_e|^2m_\mu ^2}{\sqrt{\left(|X_e|^2+m_\mu ^2\right)^24m_\mu ^2|X_e|^2\mathrm{cos}^2\eta _e}}},`$ (51)
$`\mathrm{cos}\varphi _\nu `$ $`=`$ $`{\displaystyle \frac{|X_\nu |^2m_3^2}{\sqrt{\left(|X_\nu |^2+m_3^2\right)^24m_3^2|X_\nu |^2\mathrm{cos}^2\eta _\nu }}},`$ (52)
$`\mathrm{cos}\phi _e`$ $`=`$ $`{\displaystyle \frac{|X_e|m_\mu \mathrm{cos}\eta _e}{\sqrt{|X_e|^2+m_\mu ^22m_\mu |X_e|\mathrm{cos}\eta _e}}},`$ (53)
$`\mathrm{cos}\phi _\nu `$ $`=`$ $`{\displaystyle \frac{|X_\nu |m_3\mathrm{cos}\eta _\nu }{\sqrt{|X_\nu |^2+m_3^22m_3|X_\nu |\mathrm{cos}\eta _\nu }}},`$ (54)
where
$`|X_e|^2`$ $`=`$ $`m_e^2+m_\tau ^2|m_e|m_\tau \left({\displaystyle \frac{1+\alpha _e^2}{\alpha _e}}\right),`$ (55)
$`|X_\nu |^2`$ $`=`$ $`m_1^2+m_2^22|m_1|m_2.`$ (56)
In the following discussions we consider the normal mass hierarchy $`|m_1|<m_2m_3`$ for the neutrino masses. Then the evolution effects can be ignored. Scenarios in which the neutrino masses have the quasi degenerate or the inverse hierarchy are denied from Eqs. (64) and (69).
In order to reproduce the large lepton mixing between the second and third generation, we now choose specific values of the parameters $`\alpha _e`$ and $`\eta _e`$ such that
$`\alpha _e={\displaystyle \frac{m_e^2m_\mu ^2+m_\tau ^2+\sqrt{(m_e^2m_\mu ^2+m_\tau ^2)^24m_e^2m_\tau ^2}}{2m_em_\tau }}{\displaystyle \frac{m_\tau }{|m_e|}}\left[1\left({\displaystyle \frac{m_\mu }{m_\tau }}\right)^2\right],`$ (57)
$`\mathrm{cos}^2\eta _e1.`$ (58)
Then, we obtain
$`\phi _e`$ $`{\displaystyle \frac{\pi \eta _e}{2}},`$ $`\phi _\nu `$ $`\pi \eta _\nu ,`$ $`\varphi _e`$ $`\{\begin{array}{cc}\frac{\pi }{2}\hfill & \text{for }0<\eta _e<\pi \hfill \\ \frac{3\pi }{2}\hfill & \text{for }\pi <\eta _e<2\pi \hfill \end{array},`$ $`\varphi _\nu `$ $`\pi ,`$ (61)
and
$`|\rho _\nu |`$ $`{\displaystyle \frac{1}{\sqrt{2}}},`$ $`|\sigma _\nu |`$ $`{\displaystyle \frac{1}{\sqrt{2}}},`$ $`s_e^{}`$ $`{\displaystyle \frac{|m_e|m_\mu }{m_\tau ^2}}=1.63\times 10^5,`$ $`c_e^{}`$ $`1.`$ (62)
Thus, the explicit magnitudes of the components of $`|U_{ij}|`$ are obtained as
$`\left|U_{11}\right|`$ $`\sqrt{{\displaystyle \frac{m_2}{m_2+m_1}}},`$ $`\left|U_{12}\right|`$ $`\sqrt{{\displaystyle \frac{m_1}{m_2+m_1}}},`$ $`\left|U_{13}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{|m_e|m_\mu }{m_\tau ^2}},`$
$`\left|U_{21}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{{\displaystyle \frac{m_1}{m_2+m_1}}},`$ $`\left|U_{22}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{{\displaystyle \frac{m_2}{m_2+m_1}}},`$ $`\left|U_{23}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}},`$
$`\left|U_{31}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{{\displaystyle \frac{m_1}{m_2+m_1}}},`$ $`\left|U_{32}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{{\displaystyle \frac{m_2}{m_2+m_1}}},`$ $`\left|U_{33}\right|`$ $`{\displaystyle \frac{1}{\sqrt{2}}}.`$ (63)
In Fig. 2, we present more detailed predicted values for $`|U_{23}|`$ in the $`\eta _e\eta _\nu `$ parameter space, by taking the value for $`\alpha _e`$ given in Eq. (57). It can be seen from Fig. 2 that the large mixing angle between the second and third generation is well realized in the model if we use the specific values of $`\alpha _e`$ given in Eq. (57).
As seen from Eq. (63), the neutrino oscillation angles of the model are related to the lepton masses as follows:
$`\mathrm{tan}^2\theta _{\text{solar}}`$ $`=`$ $`{\displaystyle \frac{|U_{12}|^2}{|U_{11}|^2}}{\displaystyle \frac{m_1}{m_2}},`$ (64)
$`\mathrm{sin}^22\theta _{\text{atm}}`$ $`=`$ $`4|U_{23}|^2|U_{33}|^21,`$ (65)
$`|U_{13}|^2`$ $``$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{m_em_\mu }{m_\tau ^2}}\right)^2.`$ (66)
It should be noted that the present model leads to the same results for $`\theta _{\text{solar}}`$ and $`\theta _{\text{atm}}`$ as the model in RefMatsuda2 , while a different feature for $`|U_{13}|^2`$ is derived.
On the other hand, we haveGarcia a experimental bound for $`|U_{13}|_{\text{exp}}^2`$ from the CHOOZchooz , solarsno , and atmospheric neutrino experimentsskamioka . From the global analysis of the SNO solar neutrino experimentsno ; Garcia , we have $`\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{tan}^2\theta _{12}`$ for the large mixing angle (LMA) Mikheyev-Smirnov-Wolfenstein (MSW) solution. From the atmospheric neutrino experimentskamioka ; Garcia , we also have $`\mathrm{\Delta }m_{23}^2`$ and $`\mathrm{tan}^2\theta _{23}`$. These experimental data with $`3\sigma `$ range are given by
$`|U_{13}|_{\text{exp}}^2<0.054,`$ (67)
$`\mathrm{\Delta }m_{12}^2=m_2^2m_1^2=\mathrm{\Delta }m_{\text{sol}}^2=(5.29.8)\times 10^5\text{eV}^2,`$ (68)
$`\mathrm{tan}^2\theta _{12}=\mathrm{tan}^2\theta _{\text{sol}}=0.290.64,`$ (69)
$`\mathrm{\Delta }m_{23}^2=m_3^2m_2^2\mathrm{\Delta }m_{\text{atm}}^2=(1.43.4)\times 10^3\text{eV}^2,`$ (70)
$`\mathrm{tan}^2\theta _{23}\mathrm{tan}^2\theta _{\text{atm}}=0.492.2.`$ (71)
Hereafter, for simplicity, we take $`\mathrm{tan}^2\theta _{\text{atm}}1`$. Thus, by combining the present model with the mixing angle $`\theta _{\text{sol}}`$, we have
$$\frac{m_1}{m_2}\mathrm{tan}^2\theta _{\text{sol}}=0.290.64.$$
(72)
Therefore we predict the neutrino masses as follows.
$`m_1^2`$ $`=`$ $`(0.486.8)\times 10^5\mathrm{eV}^2,`$
$`m_2^2`$ $`=`$ $`(5.716.6)\times 10^5\mathrm{eV}^2,`$ (73)
$`m_3^2`$ $`=`$ $`(1.43.4)\times 10^3\mathrm{eV}^2.`$
Let us mention a specific feature of the model. Our model imposes a stringent restriction on $`|U_{13}|`$ as
$$|U_{13}|^2\frac{1}{2}\left(\frac{m_em_\mu }{m_\tau ^2}\right)^2=1.3\times 10^{10}.$$
(74)
Here we have used the running charged lepton masses at the unification scale $`\mu =\mathrm{\Lambda }_X`$ Fusaoka : $`m_e(\mathrm{\Lambda }_X)=0.325\text{MeV}`$, $`m_\mu (\mathrm{\Lambda }_X)=68.6\text{MeV}`$, and $`m_\tau (\mathrm{\Lambda }_X)=1171.4\pm 0.2\text{MeV}`$. The value in Eq.(74) is consistent with the present experimental constraints Eq.(67), however it is too small to be checked in neutrino factories in future. The very small predicted value for $`|U_{13}|`$ is in contrast to previously proposed modelKoide Matsuda2 .
Next let us discuss the CP-violation phases in the lepton mixing matrix. The Majorana neutrino fields do not have the freedom of rephasing invariance, so that we can use only the rephasing freedom of $`M_e`$ to transform Eq. (38) to the standard form
$`U_{\mathrm{std}}=\text{diag}(e^{i\alpha _1^e},e^{i\alpha _2^e},e^{i\alpha _2^e})U`$
$`=\left(\begin{array}{ccc}c_{\nu 13}c_{\nu 12}& c_{\nu 13}s_{\nu 12}e^{i\beta }& s_{\nu 13}e^{i(\gamma \delta _\nu )}\\ (c_{\nu 23}s_{\nu 12}s_{\nu 23}c_{\nu 23}s_{\nu 13}e^{i\delta _\nu })e^{i\beta }& c_{\nu 23}c_{\nu 12}s_{\nu 23}s_{\nu 12}s_{\nu 13}e^{i\delta _\nu }& s_{\nu 23}c_{\nu 13}e^{i(\gamma \beta )}\\ (s_{\nu 23}s_{\nu 12}c_{\nu 23}c_{\nu 12}s_{\nu 13}e^{i\delta _\nu })e^{i\gamma }& (s_{\nu 23}c_{\nu 12}c_{\nu 23}s_{\nu 12}s_{\nu 13}e^{i\delta _\nu })e^{i(\gamma \beta )}& c_{\nu 23}c_{\nu 13}\end{array}\right).`$ (78)
(79)
Here, $`\alpha _i^e`$ comes from the rephasing in the charged lepton fields to make the choice of phase convention. The CP-violating phase $`\delta _\nu `$, the additional Majorana phase $`\beta `$ and $`\gamma `$ bilenky ; Doi in the representation Eq. (79) are calculable and obtained as
$`\delta _\nu `$ $`=`$ $`\text{arg}\left[{\displaystyle \frac{U_{12}U_{22}^{}}{U_{13}U_{23}^{}}}+{\displaystyle \frac{|U_{12}|^2}{1|U_{13}|^2}}\right]\text{arg}\left({\displaystyle \frac{U_{12}U_{22}^{}}{U_{13}U_{23}^{}}}\right)`$ (82)
$``$ $`\phi _e\phi _\nu {\displaystyle \frac{1}{2}}(\varphi _e\varphi _\nu ),\{\begin{array}{cc}\frac{\eta _e}{2}+\eta _\nu \frac{\pi }{4}\hfill & \text{for }0<\eta _e<\pi \hfill \\ \frac{\eta _e}{2}+\eta _\nu \frac{3\pi }{4}\hfill & \text{for }\pi <\eta _e<2\pi \hfill \end{array},`$
$`\beta `$ $`=`$ $`\text{arg}\left({\displaystyle \frac{U_{12}}{U_{11}}}\right){\displaystyle \frac{3\pi }{2}},`$ (83)
$`\gamma `$ $`=`$ $`\text{arg}\left({\displaystyle \frac{U_{13}}{U_{11}}}e^{i\delta _\nu }\right)\{\begin{array}{cc}\frac{\eta _\nu }{2}+\frac{\pi }{2}\hfill & \text{for }0<\eta _e<\pi \hfill \\ \frac{\eta _\nu }{2}\frac{\pi }{2}\hfill & \text{for }\pi <\eta _e<2\pi \hfill \end{array},`$ (86)
by using the relation $`m_em_\tau `$. Hence, we also predict the averaged neutrino mass $`m_\nu `$ which appears in the neutrinoless double beta decayDoi as follows:
$`m_\nu `$ $``$ $`\left|m_1U_{11}^2+m_2U_{12}^2+m_3U_{13}^2\right|`$ (87)
$``$ $`\left|m_1{\displaystyle \frac{m_2}{m_2+m_1}}+m_2{\displaystyle \frac{m_1}{m_2+m_1}}e^{3\pi i}\right|m_1.`$
This value of $`m_\nu `$ is too small to be observed in near future experiments.
In Fig. 2, we present more detailed predicted values for the lepton mixing matrix elements ($`|U_{12}|`$, $`|U_{23}|`$, and $`|U_{13}|`$), and phases ($`\mathrm{sin}\delta _\nu `$, $`\mathrm{sin}\beta `$, and $`\mathrm{sin}\gamma `$) in the $`\eta _e`$ \- $`\eta _\nu `$ parameter plane. Here we take a value given in Eq. (57) for the parameter $`\alpha _e`$ and center values given in Eq. (73) for neutrino masses $`m_i`$.
## IV conclusion
We have analyzed the lepton mixing matrix of a recently proposed nonsymmetric mass matrix model. The model gives a universal description of quark and lepton with the same texture form (1) based on an extended flavor $`23`$ symmetry including a phase $`\varphi `$. By using the charged lepton masses as inputs, the present model has six adjustable parameters, $`\alpha _e`$, $`\eta _e`$, $`\eta _\nu `$, $`m_1`$, $`m_2`$, and $`m_3`$ to reproduce the observed MNSP lepton mixing matrix parameters and neutrino-mass-squared differences. We have shown that only the case where the type B assignment for charged leptons and the type A for neutrinos are taken can lead to consistent values with neutrino oscillation experiments. In this case, we find that the observed large lepton mixing between the second and third generation is realized by a fine tuning of the parameter $`\alpha _e`$ as given in Eq. (57). It is also shown that the model predicts very small value for $`|U_{13}|`$, which is in contrast to previously proposed modelKoide Matsuda2 . The $`CP`$ violating phases $`\delta _\nu `$, $`\beta `$, and $`\gamma `$ in the lepton mixing matrix are obtained. The decay rate of the neutrinoless double beta decay is also predicted to be almost suppressed.
###### Acknowledgements.
This work of K.M. was supported by the JSPS, No. 3700.
|
warning/0506/hep-lat0506020.html
|
ar5iv
|
text
|
# QED3 on a space-time lattice: compact versus noncompact formulation
## I Introduction
Quantum electrodynamics in 2+1 dimensions (QED<sub>3</sub>) is interesting as a toy model for investigating the mechanism of confinement in gauge theories Po , and as an effective description of low-dimensional, correlated, electronic condensed matter systems, like spin systems us ; senthil , or high-$`T_c`$ superconductors MaPa . While the compact formulation of QED<sub>3</sub> appears to be more suitable for studying the mechanism of confinement, both compact AM and noncompact formulations arise in condensed matter systems. Our paper aims to elucidate some aspects of the relationship between these two formulations of QED<sub>3</sub> on the lattice.
Polyakov showed that compact QED<sub>3</sub> without fermion degrees of freedom is always confining Po . Any pair of test electric charge and anti-charge is confined by a linear potential, as an effect of proliferation of instantons, which are magnetic monopole solutions in three dimensions. The plasma of such monopoles is what is responsible for confinement of electrically charged particles. If compact QED<sub>3</sub> is coupled to matter fields it has been argued deconf that the interaction between monopoles could turn from $`1/x`$ to $`\mathrm{ln}(x)`$ at large distances $`x`$, so that the deconfined phase may become stable at low temperature. The issue of the existence of a confinement-deconfinement transition in QED<sub>3</sub> at $`T=0`$ is still controversial, as it has also been proposed that compact QED<sub>3</sub> with massless fermions is always in the confined phase HeSe ; AzLu ; also, in the limit of large flavor number, it has been argued that monopoles should not play any role in the confinement mechanism Herm . At finite temperature, parity invariant QED<sub>3</sub> coupled with fermionic matter undergoes a Berezinsky-Kosterlitz-Thouless transition to a deconfined phase griseso .
The issue of charge confinement in $`2+1`$ dimensional gauge models comes out to be relevant in the context of quantum phase transitions, as well. Indeed, recently it has been proposed that phenomena similar to deconfinement in high energy physics might appear in planar correlated systems, driven to a quantum (that is, zero-temperature) phase transition between an antiferromagnetically ordered (Neél) phase, and a phase with no order by continuous symmetry breaking us ; senthil . The most suitable candidate for a theoretical description of the system near the quantum critical point is a planar gauge theory, either with Fermionic matter us , or with Bosonic matter senthil .
At finite $`T`$ noncompact QED<sub>3</sub> comes about to be relevant in the analysis of the pseudo gap phase pseudo of cuprates. This phase arises from the fact that, upon doping the cuprate, a gap opens at some temperature $`T^{}`$ which is quite larger than the critical temperature $`T_C`$ for the onset of superconductivity. Both temperatures $`T^{}`$ and $`T_C`$ are doping dependent quantities and the gap is strongly dependent upon the direction in momentum space, since it exhibits $`d`$-wave symmetry TsuKi .
In Fig. 1 we report the phase diagram of high-$`T_c`$ cuprates. For small-$`x`$ phase is characterized HaTh by an insulating antiferromagnet (AF); by increasing $`x`$, this phase evolves into a spin density wave (SDW), that is a weak antiferromagnet. The pseudo gap phase is located between this phase and the $`d`$-wave superconducting (dSC) one.
The effective theory of the pseudo gap phase pseudo turns out to be QED<sub>3</sub> MaPa ; Fra ; Her , with spatial anisotropies in the covariant derivatives, that is with different values for the Fermi and the Gap velocities HaTh , and with Fermionic matter given by spin-$`1/2`$ chargeless excitations of the superconducting state (spinons). These excitations are minimally coupled to a massless gauge field, which arises from the fluctuating topological defects in the superconducting phase. The SDW order parameter is identified with the order parameter for chiral symmetry breaking (CSB) in the gauge theory, that is, $`\overline{\psi }\psi `$ Her . There can be two possibilities; if $`\overline{\psi }\psi `$ is different from zero, then the $`d`$-wave superconducting phase is connected to the spin density wave one (see Fig. 1 case b); otherwise the two phases are separated at $`T=0`$ by the pseudo gap phase (see Fig. 1 case a).
Confinement and chiral symmetry breaking go essentially together as strong coupling phenomena in gauge theories; while confinement is an observed property of the strong interactions and it is an unproven, but widely believed feature of non-abelian gauge theories in four space-time dimensions, chiral symmetry is only an approximate symmetry of particle physics, since the up and down quarks are light but not massless. Central to our understanding of CSB is the existence of a critical coupling: when fermions have a sufficiently strong attractive interaction there is a pairing instability and the ensuing condensate breaks some of the flavor symmetries, generate quark masses, and represents chiral symmetry in the Nambu-Goldstone mode namgot ; modernnamgot . The issue of a critical coupling has been widely investigated in 2+1 dimensional gauge theories cc1 ; cc2 ; cc3 . Typically, the dimensionless expansion parameter is $`1/N_f`$. Using the Schwinger-Dyson equations cc1 or a current algebra approach diaseso for QED<sub>3</sub> and QCD<sub>3</sub> one finds that there is a critical number of flavors, $`N_{f,c}`$, such that only for $`N_f`$ lesser than $`N_{f,c}`$ chiral symmetry is broken; for $`N_f`$ bigger than $`N_{f,c}`$ chirality is unbroken and quarks remain massless. For QED<sub>3</sub> this result has been the subject of some debate cc1 ; cc2 ; cc4 ; cc5 ; ApWi ; ApCoSc ; Fischer:2004nq ; there are, however, numerical simulations simulflav ; HaKoSt ; HaKoScSt of QED<sub>3</sub>, which find an $`N_{f,c}`$ remarkably close to the results reported in Ref. cc1 .
Even if far from the scaling regime, strong coupling gauge theories on the lattice provide interesting clues on the issue of CSB. In fact, one can show that, in the strong coupling limit, a Hamiltonian with $`N_c`$ colors of fermions and $`N_f/2`$ lattice flavors of staggered fermions is effectively a $`U(N_f/2)`$ quantum antiferromagnet with representations determined by $`N_c`$ and $`N_f`$ strongcoupling . CSB is then associated strongcoupling either to the formation of a $`U(1)`$ commensurate charge density wave or of a $`SU(N_f/2)`$ spin density wave, i.e. to the formation of Neél order. Quantum antiferromagnets with the representations considered in Ref. strongcoupling have been analyzed in Ref. reasac where it was found that, for small enough $`N_f`$, the ground state is ordered. Also, when $`N_f`$ is increased there is a phase transition, for $`N_fN_c`$, to a disordered state. In this picture, the large $`N_c`$ limit is the classical limit where Neél order is favored and the small $`N_c`$ and large $`N_f`$ limit are where fluctuations are large and disordered ground states are favored.
We shall not try to ascertain in this paper the critical number of flavours $`N_{f,c}`$. Here, we shall analyze the relationship between monopole density and fermion mass and compare the results obtained for the compact and noncompact lattice formulation of this gauge model. In particular, we revisit the analysis of Fiebig and Woloshyn of Refs. FiWo1 ; FiWo2 , where the dynamic equivalence between the two formulations of (isotropic) QED<sub>3</sub> is claimed to be valid in the finite lattice regime. In this paper we shall extend the comparison to the continuum limit, following the same approach as in Refs. FiWo1 ; FiWo2 : namely we shall analyze the behavior of the chiral condensate and of the monopole density as the continuum limit is reached.
In Section II we describe the model and its properties both in the continuum and on the lattice. Moreover, the method for detecting monopoles on the lattice is illustrated.
In Section III a description of both compact and noncompact formulations of QED<sub>3</sub> is given.
In Section IV we present our numerical result for the chiral condensate and the monopole density in the region in which the continuum limit is reached. Then, we compare our results with those of Fiebig and Woloshyn FiWo1 ; FiWo2 .
Section V is devoted to conclusions.
## II The model and its properties
The continuum Lagrangian density describing QED<sub>3</sub> is given in Minkowski metric Ap by
$$=\frac{1}{4}F_{\mu \nu }^2+\overline{\psi }_iiD_\mu \gamma ^\mu \psi _im_0\overline{\psi }_i\psi _i,$$
(1)
where $`D_\mu =_\mu ieA_\mu `$, $`F_{\mu \nu }`$ is the field strength and the fermions $`\psi _i`$ ($`i=1,\mathrm{},N_f`$) are 4-component spinors. Since QED<sub>3</sub> is a super-renormalizable theory, dim$`[e]=+1/2`$, the coupling does not display any energy dependence. One may define three $`4\times 4`$ Dirac matrices
$$\gamma ^0=\left(\begin{array}{cc}\sigma _3& 0\\ 0& \sigma _3\end{array}\right),\gamma ^1=\left(\begin{array}{cc}i\sigma _1& 0\\ 0& i\sigma _1\end{array}\right),$$
$$\gamma ^2=\left(\begin{array}{cc}i\sigma _2& 0\\ 0& i\sigma _2\end{array}\right),$$
(2)
and two more matrices anticommuting with them: namely
$$\gamma ^3=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma ^5=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(3)
The massless theory will therefore be invariant under the chiral transformations
$$\psi e^{i\alpha \gamma ^3}\psi ,\psi e^{i\beta \gamma ^5}\psi .$$
(4)
If one writes a 4-component spinor as $`\psi =\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right),`$ the mass term becomes
$$m\overline{\psi }\psi =m\psi _1^{}\sigma _3\psi _1m\psi _2^{}\sigma _3\psi _2.$$
Since in three dimensions the parity transformation reads
$`\psi _1(x_0,x_1,x_2)`$ $``$ $`\sigma _1\psi _2(x_0,x_1,x_2),`$
$`\psi _2(x_0,x_1,x_2)`$ $``$ $`\sigma _1\psi _1(x_0,x_1,x_2),`$ (5)
then $`m\overline{\psi }\psi `$ is parity conserving.
The lattice Euclidean action AlFaHaKoMo ; HaKoSt using staggered fermion fields $`\overline{\chi },\chi `$, is given by
$$S=S_G+\underset{i=1}{\overset{N}{}}\underset{n,m}{}\overline{\chi }_i(n)M_{n,m}\chi _i(m),$$
(6)
where $`S_G`$ is the gauge field action and the fermion matrix is given by
$`M_{n,m}[U]`$ (7)
$`={\displaystyle \underset{\nu =1,2,3}{}}{\displaystyle \frac{\eta _\nu (n)}{2}}\left\{[U_\nu (n)]\delta _{m,n+\widehat{\nu }}[U_\nu ^{}(m)]\delta _{m,n\widehat{\nu }}\right\}.`$
The action (6) allows to simulate $`N=1,2`$ flavours of staggered fermions corresponding to $`N_f=2,4`$ flavours of 4-component fermions $`\psi `$ BuBu . $`S_G`$ is different for the compact and noncompact formulation of QED<sub>3</sub>.
For the compact formulation one has
$$S_G[U]=\beta \underset{n,\mu <\nu }{}\left[1\frac{1}{2}\left(U_{\mu \nu }(n)+U_{\mu \nu }^{}(n)\right)\right],$$
(8)
where $`U_{\mu \nu }(n)`$ is the “plaquette variable” and $`\beta =1/(e^2a)`$, $`a`$ being the lattice spacing. Instead, in the noncompact formulation one has
$$S_G[\alpha ]=\frac{\beta }{2}\underset{n,\mu <\nu }{}F_{\mu \nu }(n)F_{\mu \nu }(n),$$
(9)
where
$$F_{\mu \nu }(n)=\{\alpha _\nu (n+\widehat{\mu })\alpha _\nu (n)\}\{\alpha _\mu (n+\widehat{\nu })\alpha _\mu (n)\}$$
(10)
and $`\alpha _\mu (n)`$ is the phase of the “link variable” $`U_\mu (n)=e^{i\alpha _\mu (n)}`$, related to gauge field by $`\alpha _\mu (n)=aeA_\mu (n)`$.
Monopoles are detected in the lattice using the method given by DeGrand and Toussaint DeTo : due to the Gauss’s law, the total magnetic flux emanating from a closed surface allows to determine if the surface encloses a monopole. The monopole density is defined by half of the total number of monopoles and antimonopoles divided by the number of elementary cubes in the lattice. We apply this definition for both the compact and the noncompact formulations of the theory, although some caution should be used in this respect. Indeed, monopoles are classical solutions of the theory with finite action only for compact QED<sub>3</sub>, where they are known to play a relevant role. In the noncompact formulation of QED<sub>3</sub> they are not classical solutions, but they could give a contribution to the Feynman path integral owing to the periodic structure of the fermionic sector HaWe .
## III Compact versus noncompact formulation
In order to investigate the onset of the continuum physics, it is convenient to consider a dimensionless observable and to evaluate it from the lattice for increasing $`\beta `$ until it reaches a plateau. Such an observable can be taken to be $`\beta ^2\overline{\chi }\chi `$, which is expected to become constant in the continuum ($`\beta \mathrm{}`$) limit simulflav ; DaKoKo2 . Numerical simulations show two regimes: for $`\beta `$ larger than a certain value, the theory is in the continuum limit (flat dependence of a dimensionless observable from $`\beta `$), otherwise the system is in a phase with finite lattice spacing. In the former regime, the theory describes continuum physics, in the latter one it is appropriate to describe a lattice condensed-matter-like system.
There are a couple of papers by Fiebig and Woloshyn in which the two formulations are compared in the finite lattice regime FiWo1 ; FiWo2 . In these papers the $`\beta `$-dependence of the chiral condensate and of the monopole density for lattice QED<sub>3</sub> with $`N_f=0`$ and $`N_f=2`$ are analyzed for both compact and noncompact formulations in the finite lattice regime.
It is shown there that, when $`\overline{\chi }\chi `$ is plotted versus the monopole density $`\rho _m`$, data points for both theories fall on the same curve to a good approximation (see Fig. 2). This led the authors of Refs. FiWo1 ; FiWo2 to the conclusion that the physics of the chiral symmetry breaking is the same in the two theories.
Our program is to study if the conclusion reached by Fiebig and Woloshyn can be extended to the continuum limit, by looking at the same observables they considered: namely the chiral condensate and the monopole density.
## IV Numerical results
Since QED<sub>3</sub> is a super-renormalizable theory, the coupling constant does not display any lattice space dependence. The continuum limit is approached by merely sending $`\beta =1/(e^2a)`$ to infinity. In this limit all physical quantities can be expressed in units of the scale set by the coupling $`e`$. Therefore, it is natural to work in terms of dimensionless variables such as $`\beta m`$, $`L/\beta `$ or $`\beta ^2\overline{\chi }\chi `$, which depend on $`e`$ ($`L`$ is the lattice size).
The signature that the continuum limit is approached is that data taken at different $`\beta `$ should overlap on a single curve when plotted in dimensionless units HaKoSt .
In practice, numerical results will not describe the correct physics of the system even in the continuum limit because of finite volume effects which are particularly significant in our case, due to the presence of a massless particle, the photon. In principle one should get rid of these effects by taking $`L/\beta \mathrm{}`$. In practice, this ratio is taken to be large, but finite. In Ref. GuRe the authors conclude that in order to find chiral symmetry breaking for $`N_f=2`$ at least a ratio $`L/\beta 5\times 10^3`$ is required. In our simulations the largest value for the $`L/\beta `$ ratio has been 20.
Our Monte Carlo simulation code was based on the hybrid updating algorithm, with a microcanonical time step set to $`dt=0.02`$. We simulated one flavour of staggered fermions corresponding to two flavours of 4-component fermions. Most simulations were performed on a $`12^3`$ lattice, for bare quark mass ranging in the interval $`am=0.01÷0.05`$. We made refreshments of the gauge (pseudofermion) fields every 7 (13) steps of the molecular dynamics. In order to reduce autocorrelation effects, “measurements” were taken every 50 steps. Data were analyzed by the jackknife method combined with binning.
As a first step, we have reproduced the results by Fiebig and Woloshyn which are shown in Fig. 2. We find that also in our case data points from the two formulations nicely overlap (see Fig. 3). It should be noticed that data of Fig. 2 were obtained using a linear fit with two masses ($`am`$=0.025, 0.05) whilst those of Fig. 3 have been obtained by a quadratic fit with four masses ($`am`$=0.02, 0.03, 0.04, 0.05), nevertheless the conclusion is the same in both cases. We have verified that if we perform a linear fit on the subset of our data with masses $`am`$=0.02 and 0.05 and on the subset with masses $`am`$=0.03 and 0.05, our results nicely compare with those plotted in Fig. 2.
Then, in Fig. 4 we plot data for $`\beta ^2\overline{\chi }\chi `$ obtained in the compact formulation versus $`\beta m`$. We restrict our attention to the subset of $`\beta `$ values for which data points fall approximately on the same curve, which in the present case means $`\beta =1.9,2.0,2.1`$, corresponding to $`L/\beta =6.31,6.00,5.71`$. A linear fit of these data points gives $`\chi ^2`$/d.o.f. $`8.4`$ and the extrapolated value for $`\beta m0`$ turns out to be $`\beta ^2\overline{\chi }\chi =(1.54\pm 0.25)\times 10^3`$. Restricting the sample to the data at $`\beta =2.1`$, the $`\chi ^2/`$d.o.f. lowers to $`1.3`$ and the extrapolated value becomes $`\beta ^2\overline{\chi }\chi =(0.94\pm 0.28)\times 10^3`$, thus showing that there is a strong instability in the determination of the chiral limit. If instead a quadratic fit is used for the points obtained with $`\beta =1.9,2.0,2.1`$, we get $`\beta ^2\overline{\chi }\chi =(0.91\pm 0.45)\times 10^3`$ with $`\chi ^2`$/d.o.f. $`8.7`$. Owing to the large uncertainty, this determination turns out to be compatible with both the previous ones.
In Fig. 5 we plot data for $`\beta ^2\overline{\chi }\chi `$ obtained in the noncompact formulation versus $`\beta m`$. Following the same strategy outlined before, we restrict our analysis to the data obtained with $`\beta =0.7,0.75,0.8`$, which correspond to $`L/\beta =17.14,16,15`$.
If we consider a linear fit of these data and extrapolate to $`\beta m0`$, we get $`\beta ^2\overline{\chi }\chi =(0.45\pm 0.03)\times 10^3`$ with $`\chi ^2`$/d.o.f. $`17`$. Performing the fit only on the data obtained with $`\beta =0.8`$, for which a linear fit gives the best $`\chi ^2/`$d.o.f. value $`16`$, we obtain the extrapolated value $`\beta ^2\overline{\chi }\chi =(0.66\pm 0.07)\times 10^3`$. Therefore, also in the noncompact formulation the chiral extrapolation resulting from a linear fit is largely unstable. A quadratic fit in this case gives instead a negative value for $`\beta ^2\overline{\chi }\chi `$.
The comparison of the extrapolated value for $`\beta ^2\overline{\chi }\chi `$ in the two formulations is difficult owing to the instabilities of the fits and to the low reliability of the linear fits, as suggested by the large values of the $`\chi ^2/`$d.o.f. Taking an optimistic point of view, one could say that the extrapolated $`\beta ^2\overline{\chi }\chi `$ for $`\beta =2.1`$ in the compact formulation is compatible with the extrapolated value obtained in the noncompact formulation for $`\beta =0.8`$.
It is worth mentioning that our results in the noncompact formulation are consistent with known results: indeed, if we carry out a linear fit of the data for $`\beta =0.6,0.7,0.8`$ and $`am`$=0.02, 0.03, 0.04, 0.05 and extrapolate, we get $`\beta ^2\overline{\chi }\chi =(1.30\pm 0.07)\times 10^3`$ with an admittedly large $`\chi ^2/`$d.o.f. $`20`$, but very much in agreement with the value $`\beta ^2\overline{\chi }\chi =(1.40\pm 0.16)\times 10^3`$ obtained in Ref. AlFaHaKoMo .
We stress again that our results are plagued by strong finite volume effects, therefore our conclusions on the extrapolated values of $`\beta ^2\overline{\chi }\chi `$ are significant only in the compact versus noncompact comparison we are interested in. We do not even try to draw any conclusion from our data on the critical number of the flavours. As a matter of fact a recent paper HaKoSt shows that, if effects are carefully monitored and large lattices, up to $`50^3`$, are used, it is possible to establish that $`\beta ^2\overline{\chi }\chi 5\times 10^5`$. For the comparison between compact and noncompact QED<sub>3</sub> it is pertinent to carry out the numerical analysis with an (approximately) constant value of the ratio $`L/\beta `$. This condition is indeed verified even if we performed simulations on lattices with fixed ($`L=12`$) size, since the range of allowed values for $`\beta `$ corresponding to the continuum limit is narrow ($`\beta =1.8÷2.2`$ in the compact case, $`\beta =0.6÷0.9`$ in the noncompact case). Finite volume effects play a “second order” role in our work, since they probably only affect the extension of the continuum limit window of $`\beta `$ values.
In Fig. 6 we plot $`\overline{\chi }\chi `$ versus the monopole density $`\rho _m`$. Differently from Figs. 2-3, it is not evident with the present results that the two formulations are equivalent also in the continuum limit, although such an equivalence cannot yet be excluded.
In Fig. 7 we plot again $`\overline{\chi }\chi `$ versus the monopole density $`\rho _m`$, but now on a $`32^3`$ lattice. In this case the chiral condensate is extrapolated to zero mass by a quadratic fit. In spite of the negative value taken by $`\overline{\chi }\chi `$ for large $`\beta `$, in this case data for both formulations seem to fall on the same curve.
In Fig. 8 and Fig. 9 we plot $`\beta ^3\rho _m`$ versus $`\beta m`$ for the two formulations; the former quantity is dimensionless, therefore, in analogy with the previous cases, we expect that data at different $`\beta `$ values should fall on a single curve in the continuum limit. Our results show that this is not the case, this suggesting that the continuum limit has not been reached for the monopole density.
Simulations on the $`32^3`$ lattice give practically the same results for $`\beta ^3\rho _m`$, indicating that this observable, unlike $`\beta ^2\overline{\chi }\chi `$, is volume independent.
It is important to observe, however, that the monopole density is independent of the fermion mass. Since the mechanism of confinement in the theory with infinitely massive fermions, i.e. in the pure gauge theory, is based on monopoles and since the monopole density is not affected by the fermion mass, we may conjecture that this same mechanism holds also in the chiral limit. This supports the arguments by Herbut about the confinement in the presence of massless fermion HeSe ; AzLu .
## V Conclusions
In this paper we have compared the compact and the noncompact formulations of QED<sub>3</sub> by looking at the behavior of the chiral condensate and the monopole density.
Numerical results for $`\beta ^2\overline{\chi }\chi `$ are compatible with those obtained by other groups, although it is still questionable if the continuum limit has been reached and if the chiral limit is stable. The biggest difficulty for this observable is that the chiral extrapolation is rough when a linear fit is performed, but gives a negative value when instead a quadratic fit is considered. Massive calculations on larger lattices are needed to further reduce the finite volume effects and to stabilize the chiral limit.
As far as monopoles are concerned, they appear in smaller and smaller numbers for large $`\beta `$, this making the determination of the continuum limit for $`\beta ^3\rho _m`$ rather problematic. Our results show, however, a very weak volume dependence.
We have analyzed also the relationship between the monopole density and the fermion mass, both in compact and noncompact QED<sub>3</sub>. The weak dependence observed leads us to conclude that the Polyakov mechanism for confinement holds not only in the pure gauge theory, but also in presence of massless fermions.
Finally, we note that, although the chiral condensate and monopole density approach the continuum limit in two different ranges of $`\beta `$, the analysis à la “Fiebig and Woloshin” does not allow to exclude the equivalence of the compact and noncompact lattice formulations of QED<sub>3</sub>.
|
warning/0506/cond-mat0506182.html
|
ar5iv
|
text
|
# Flow behavior of colloidal rod-like viruses in the nematic phase
## I Introduction
When subjected to shear flow liquid crystals can exhibit a variety of surprising phenomena, which arise because of the anisotropic shape of the constituent rods. Theoretically the behavior of a suspension of hard rods during shear flow can be described by the equation that governs the time development of their probability distribution function, as derived by Hess Hess76 and by Doi and Edwards Doi86 . In the absence of a flow, the Doi-Edwards-Hess (DEH) theory reduces to the Onsager description of equilibrium nematic liquid crystals and can be used to describe the isotropic-nematic (I-N) phase transition of a hard rod suspension Onsager49 . The rheological properties are predicted to be highly non-linear functions of the Péclet number (Pe), which is the ratio of shear rate $`\dot{\gamma }`$ over rotational diffusion constant $`D_r`$. This is not surprising as the Pe number can be much larger than unity when the rod-like molecules have large aspect ratios.
The nonlinear response of the rheological properties indicates that the shear flow distorts the equilibrium distribution of macromolecules or rods. The spatiotemporal microstructural changes during flow are even more complex. At low shear rates, the DEH theory predicts that the pseudo vector describing the average alignment of the rods, i.e. the “director”, undergoes a continuous “tumbling” motion in the plane defined by the velocity and the velocity gradient vectors. At high shear rates the director is predicted to align with the flowMarrucci89 ; Larson90 . At intermediate shear rates it is possible to obtain multiple solutions to the Doi-Edwards-Hess equation, which are dependent on the initial orientation of the director Faraoni99 ; Forest03 . For one stable solution called ”wagging” the nematic director oscillates between two angles in the plane defined by the flow and the gradient of the flow. Other solutions such as kayaking and log-rolling are also possible, in which the director oscillates out of the flow-gradient plane at these intermediate shear rates Grosso03 .
Experiments on polymeric liquid crystals have confirmed several predictions of the Doi-Edwards equation. Using a combination of rheological and rheo-optical measurements it was shown that nematic solutions of poly(benzyl-glutamate) tumble (PBG) at low shear rate and become flow aligning at high shear rates Burghardt91 . The existence of a wagging regime and a potential coexistence of wagging and log-rolling regimes at intermediate flow rates have also been revealed in experiments Mewis97 ; Grosso03 . However, there remain significant difficulties when comparing experiments on polymeric liquid crystals (PLC) to theoretical predictions. One problem is that different levels of the microstructure may lead to different contributions to the stress tensor Larson96 . In addition to the molecular contribution to the stress tensor, textural aspects contribute. The latter include Frank elasticity contributions due the presence of spatial distortions of nematic director, and viscous interactions between ’domains’. In addition, there is an indirect effect to the stress tensor as the defects disturb the orientation distribution function. These textural contributions to the total stress dominates the behavior at high concentrations and low shear rates Walker96b , making it difficult to accurately extract information about the concentration dependence of different flow transitions. The textural portion of the stress typically displays scaling of the transient rheological response with strain rather than with Pe number Moldenaers86 . The strain scaling is a typical feature of materials where the time response is determined by an inherent length scale which in the case of PLC’s is set by the size of the large non-Brownian nematic domainsLarson91 .
The DEH theory describes the flow behavior of a homogeneous ensemble of rods but does not consider any polydomain effects. Therefore an ideal system for testing DEH theory should have small textural contributions. In this paper we use rod-like *fd* virus suspensions to access the concentration dependence of the transition of tumbling to wagging, and wagging to flow aligning. We show that the contribution of textural stress is very low, although the spatial distribution of directors still has to be accounted for. The main motive for using *fd* virus is the thorough understanding of its equilibrium behavior, which has been quantitatively described using the Onsager theory extended to take into account the the semi-flexible nature of *fd* as well as its surface charge Purdy03 . Moreover, *fd* has already successfully been used for (micro-)rheology experiments in the isotropic phase Graf93 ; Schmidt00 . The aim of the present paper is to make a comparison between the *dynamic* flow behavior of *fd* suspensions and the available microscopic theoretical predictions of the DEH theory for a homogeneous system of colloidal rods under shear.
The paper is organized as follows. In section II we discuss the equation of motion of the orientational distribution function and the numerical method we use to solve it. The experimental details about sample preparation and measurements are given in section III. The results are discussed in five parts: the stationary viscosity of *fd* suspensions, the concentration and shear rate dependence of the oscillatory response to a flow reversal, the relaxation after cessation of flow at high concentration and *in situ* microscopy under shear. In section IV the textural contribution to the stress tensor is investigated more detail. Finally we present a non-equilibrium phase diagram of shear and concentration dependence of different flow transitions.
## II Theory
The distribution of an ensemble of rods can be described by the probability density function $`P(\widehat{u}_1,..,\widehat{u}_N,\stackrel{}{r}_1,..,\stackrel{}{r}_n)`$ of the positions $`\{\stackrel{}{r}_i\}`$ and orientations $`\{\widehat{u}_i\}`$ of the rods. Ignoring any spatial correlations, i.e. restricting to a monodomain, we have $`P(\widehat{u}_1,..,\widehat{u}_N,\stackrel{}{r}_1,..,\stackrel{}{r}_n)=\overline{\rho }P(\widehat{u}_1,..,\widehat{u}_N)`$, where $`\overline{\rho }=N/V`$ is the particle density. Therefore, the orientational probability density function, or orientational distribution function (ODF), fully characterizes the system. The time evolution of the ODF for a suspension of rods during flow is obtained by solving the equation of motion for the ODF, given by the $`N`$-particle Smoluchowski equation.
$`{\displaystyle \frac{P(\widehat{𝐮},t)}{t}}`$ $`=`$ $`D_r\widehat{}\{\widehat{}P(\widehat{𝐮},t)`$ (1)
$`+DL^2\overline{\rho }P(\widehat{𝐮},t)\widehat{}{\displaystyle }d\widehat{𝐮}^{}P(\widehat{𝐮}^{},t)|\widehat{𝐮}^{}\times \widehat{𝐮}|\}`$
$`\widehat{}P(\widehat{𝐮},t)\widehat{𝐮}\times \left(\widehat{\mathrm{\Gamma }}\widehat{𝐮}\right),`$
where $`(\mathrm{})=\widehat{u}\times _{\widehat{u}}(\mathrm{})`$ is the rotation operator with respect to the orientation $`\widehat{u}`$ of a rod. $`D_r`$ is the rotational diffusion of a rod at infinite dilution. Furthermore, $`D`$ is the thickness of the rods and $`L`$ is their length. $`\mathrm{\Gamma }=\dot{\gamma }\widehat{\mathrm{\Gamma }}`$ is the velocity-gradient tensor with $`\dot{\gamma }`$ the shear rate. Here we choose
$`\widehat{\mathrm{\Gamma }}=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right),`$ (5)
which corresponds to a flow $`𝐯`$ in the $`x`$-direction and its gradient $`𝐯`$ in the $`y`$-direction.
The concentration where the isotropic phase becomes unstable *in the absence of shear flow* can be calculated by solving the Smoluchowski equation at zero shear rate. This equation agrees with Onsagers approach to the I-N transition. Often the Maier-Saupe potential is used instead of the exact potential, which in fact corresponds to the first term of the Ginzburg-Landau expansion of the outer product in the exact potential given between the brackets in Eq. 1 Dhont03c . Under *flow conditions*, a rich dynamics phase behavior is found as a function of shear rate and rod concentration. Marrucci and Maffettone were the first to solve the equation of motion of the ODF numerically, restricting themselves to two dimensions in order to reduce the computational effortMarrucci89 . They found that the director undergoes a tumbling motion with respect to the flow direction, resulting in a negative normal stress $`N_1`$. Larson expanded the ODF in three dimensions using spherical harmonics and truncated the expansion after checking for convergence Larson90 . This treatment predicts a transition from tumbling to ”wagging” and finally to flow aligning state with increasing shear rates. A closure relation is frequently used for the interaction term on the right side of Eq. 1. This can greatly bias the results, see i.e. Feng et al.Feng98 . The location of the flow transitions in the flow-concentration diagram is very sensitive to the choice of the closure, and no satisfactory closure has been found up till now.
In this paper we use a finite element method to numerically solve the equation of motion for the ODF, thus avoiding the use of any specific closure relation. As a typical diffusion-convection equation, this equation describes the diffusive-convective transport dynamics of an orientation of a homogeneous ensemble of thin rigid rods. A surface of a sphere is constructed on which a tip of the rod moves with respect to its center of mass. The equation for the probability of finding the tip of a rod in an area is determined by the transport fluxes on its boundaries due to (1) the Brownian diffusion (the first term in the brace brackets of Eq. 1), (2) the convection induced by the interparticle forces (the second term in the brace brackets of Eq. 1) and (3) the convection due to the imposed shear flow (the third term of Eq. 1).
To solve Eq. 1 numerically, a discretization scheme is used, and meshes on the surface of a unit sphere are constructed. For those operators inside the brace brackets which represent the transport fluxes we apply the central differences approximations. However, the rotation operator outside of the brace brackets needs to be discretized using the concept of transport fluxes through the boundaries of the mesh. In other words, the integral form of the Eq.(1) is invoked and applied to each of the mesh elements. To do this the identity, $`(\widehat{u}\times _{\widehat{u}})𝐅=_{\widehat{u}}(𝐅\times \widehat{u})`$ is used in order to transform the angular transport flux of a rod to the translational transport flux of one tip of that rod. It differs from the conventional method of discretizing a differential equation where the operators are written explicitly into the sum of the first- and the second order derivatives and then the latter are approximated by selected difference schemes. The advantage of the current method is that, since neighboring meshes share boundaries, the fluxes leaving one mesh are always absorbed by the surrounding meshes and vice versa. Therefore, there is no loss and generation in the total amount of the ODF’s as the computation proceeds (see Fig.1). In practice a $`40\times 80`$ mesh was used on the surface of a unit sphere with $`40`$ equi-spaced grids in the polar angle and $`80`$ equi-spaced grids in the azimuthal angle in a spherical coordinates. The right hand side of Eq.1 is discritized on the meshes according to the flux-conservative method mentioned above. A fourth order Adams’ predictor-corrector method Korn68 was invoked to follow the time evolution of the ODF. More details will be published in a forthcoming paper.
The time-dependent ODF is now used to calculate the time-dependence of three parameters characterizing the flow behavior of a nematic phase: (1) $`\theta `$ describing the angle between the nematic director and flow direction, (2) the scalar magnitude of the director of defined by the order parameter $`P_2`$ and (3) the total stress of an ensemble of flowing rods. The angle and magnitude of the order parameter are obtained from the order parameter tensor
$`𝐒={\displaystyle 𝑑\widehat{𝐮}\widehat{𝐮}\widehat{𝐮}P(\widehat{𝐮},t)}.`$ (6)
The largest eigenvalue of the order parameter tensor $`\lambda `$, characterizes the degree of alignment of rods with respect to the director given by the corresponding eigenvector $`\widehat{𝐧}`$. The largest eigenvalue of $`𝐒`$ is $`1/3`$ in the isotropic phase and to 1 for a perfectly aligned nematic phase. Scalar order parameter $`P_2`$ is defined as $`P_2=(3\lambda 1)/2`$.
The stress $`\sigma _{12}`$ is obtained from the deviatoric part of the stress tensor derived by Dhont and Briels Dhont03c :
$`𝚺_D=\eta _0\dot{\gamma }+3\overline{\rho }k_BT\{𝐒{\displaystyle \frac{1}{3}}\widehat{𝐈}+{\displaystyle \frac{L}{D}}\varphi 𝚺_I^D+`$ (7)
$`{\displaystyle \frac{1}{6}}Pe_r[𝐒^{(4)}:\widehat{𝐄}{\displaystyle \frac{1}{3}}\widehat{𝐈}𝐒:\widehat{𝐄}]\},`$
where
$`𝚺_I^D={\displaystyle \frac{8}{3\pi }}{\displaystyle 𝑑\widehat{𝐮}𝑑\widehat{𝐮}^{}\widehat{𝐮}\widehat{𝐮}\times \frac{\widehat{𝐮}\times \widehat{𝐮}^{}}{|\widehat{𝐮}\times \widehat{𝐮}^{}|}\widehat{𝐮}\widehat{𝐮}^{}P(\widehat{𝐮},t)P(\widehat{𝐮}^{},t)}`$ (8)
and
$`𝐒^{(4)}={\displaystyle 𝑑\widehat{𝐮}\widehat{𝐮}\widehat{𝐮}\widehat{𝐮}\widehat{𝐮}P(\widehat{𝐮},t)}.`$ (9)
Here, $`\varphi =\frac{\pi }{4}D^2L\overline{\rho }`$ is the volume fraction of rods, and $`Pe_r=\dot{\gamma }/D_r`$ the rotational Péclet number which is defined as the shear rate scaled with the rotational diffusion of a rod at infinite dilution. The first term between the brackets, $`𝐒\frac{1}{3}\widehat{𝐈}`$, stems from the Brownian contribution to the stress. The second term stems from the direct interaction between rods and describes the elastic contribution to the total stress. The proportionality constant $`\varphi \frac{L}{D}`$ is the dimensionless rod concentration and is also called the nematic strength. The terms proportional to $`Pe_r`$ stem from the flow of the suspension and described the viscous contribution to the total stress. This term disappears instantaneously when the shear is switched off.
In Fig. 2 we plot the evolution of the three parameters (angle $`\theta `$, order parameter $`P_2`$ and stress $`\sigma _{12}`$) as a function of strain for different shear rates at a dimensionless rod concentration of $`\varphi \frac{L}{D}=4.5`$. For this calculation we used an initial rod orientation in the flow-gradient plane. The flow behavior between Péclet numbers of 4.5 and 5.0 exhibits a sharp transition from tumbling behavior, where the director continuously rotates, to wagging behavior where the director hops back and forth between two well defined angles. At higher shear rates the director is found to be flow aligning. The order parameter at low shear rates remains unchanged, but is significantly reduced at the point of the tumbling to wagging flow transition.
## III materials and methods
The viscosity and stress response was measured using an ARES strain controlled rheometer (TA instruments, Delaware). A double wall Couette geometry was used because of fairly low viscosity of the samples. Polarized light microscopy images of fd under shear flow were taken using a Linkam CSS450 plate-plate shear cell.
The physical characteristics of the bacteriophage *fd* are its length $`L=880nm`$, diameter $`D=6.6nm`$, persistence length of 2200 $`nm`$ and a charge per unit length of around 10 $`e^{}/nm`$ at pH 8.2 Fraden95 . When in solution, *fd* exhibits isotropic, cholesteric, and smectic phases with increasing concentration Dogic97 ; Dogic01 . Fd forms a cholesteric phase while the DEH theory is valid for nematic structures. In practice nematic and cholesteric phase are locally almost identical and the free energy difference between these phases is very small Dogic00 . In this paper we do not distinguish between these two phases. The fd virus was prepared according to standard biological protocols using XL1-Blue strain of E. coli as the host bacteria Sambrook85 . The standard yields are approximately $`50mg`$ of fd per liter of infected bacteria, and virus is typically grown in 6 liter batches. The virus is purified by repetitive centrifugation (108 000 g for 5 hours) and re-dispersed in a 20 mM Tris-HCl buffer at pH 8.2.
### III.1 fd as a model hard rod system
The Onsager theory for hard rod dispersions predicts a first order phase transition between a disordered, isotropic phase and an orientationally ordered, nematic phase. Due to hard core athermal interactions considered in the Onsager model, the phase diagram is temperature independent and the rod concentration is the only parameter that determines the location of the I-N phase transition. The two points spanning the region of isotropic-nematic coexistence are called the binodal points. The spinodal point is located at a rod concentration higher then the lower binodal point and is determined by the following condition $`\varphi \frac{L}{D}=4`$. *Fd* viruses are not true hard rods, due to surface charge and limited flexibility. As a consequence, their equilibrium phase behavior differs from the ideal hard rod case described by Onsager based theory, e.g. DEH. The finite flexibility of *fd* viruses drives the concentration of the binodal points to a 30% higher value when compared to equivalent but perfectly stiff hard rods. In addition, flexibility also reduces the value of the order parameter of the coexisting nematic phase. For *fd* the order parameter of the coexisting nematic is about 0.65 while Onsager theory for hard rods in equilibrium predicts the order parameter of 0.8 Purdy03 . The effect of surface charge is to increase the effective diameter of the rod $`D_{\text{eff}}`$ and therefore the excluded volume interaction between charged rods. As a consequence the charge reduces the real concentration of the phase transition Tang95 .
For the fd suspension used, the binodal point at high rod concentration $`c_{IN}`$ occurs at $`11mg/ml`$. After taking the effects of flexibility and charge into account it was shown that the order parameter of the nematic solution of *fd* is quantitatively described by the extensions of the Onsager theory to the semi-flexible case Purdy03 . Hence, even though *fd* is flexible and charged, it can be used to quantitatively test predictions of the DEH theory. It is, however, a very difficult and until now unfulfilled task to incorporate charge and flexibility into a non-equilibrium equation of motion such as Eq. 1. Therefore in this paper we use data from reference Purdy03 to convert the measured concentration of *fd* to the nematic order parameter of the sample. After that we compare experiments and theory at the same values of the order parameter.
## IV Results
### IV.1 Stationary viscosity
The measurements of a stationary viscosity as a function of the shear rate for different fd concentrations are shown in Fig. 3. For the lowest concentrations of fd the viscosity decreases continuously with shear rate except for a small hesitation at a shear rate of 10 s<sup>-1</sup>. This hesitation is similar to what is observed for solutions of PBG at low concentration in solvent m-cresol Kiss78 ; Moldenaers86 . For fd at intermediate concentrations, shear thinning becomes less pronounced, the hesitation shifts to higher shear rates and turns into a local maximum. For the highest fd concentration almost no shear thinning is observed, only a pronounced peak in the viscosity. This shear thickening behavior has not been previously reported.
A hesitation in the shear rate dependence of the viscosity was predicted theoretically by Larson Larson90 . It was argued that the transition from the tumbling regime to the wagging regime implies a broadening of the ODF which leads to higher dissipative stresses. The broadening of the ODF is illustrated in Fig. 2b. As can be seen in Fig. 2c it is not straightforward that ODF broadening really has an effect on the stress. We calculated the time-dependent viscosity by numerically solving the equation of motion of the ODF for 20 different initial orientations of the director. From the time-dependent ODFs we calculated the viscosity using either only the elastic term or both elastic and viscous terms. The viscosity is averaged over all 20 traces and a tumbling period after the transient start up flows have died out. The results are scaled to the experiment using a typical concentration of $`16;mg/ml`$ for $`\overline{\rho }`$ in eq. 7 and the value of $`D_{\text{inf}}^0=20s^1`$, taken for the rotational diffusion at infinite dilution Purdy03 ; Kramer92 . Fig. 4 the stationary viscosity decreases continuously with increasing shear rate and only shows a hesitation when the viscous contribution to the stress is not included. The shear rate where this hesitation occurs corresponds with the shear rate where the system nematic ordering is significantly reduced and the transition from tumbling to wagging takes place, as can be concluded from Fig. 2. Comparing the model predictions to the experiments it should be noted that the experimentally observed features are much more pronounced. Moreover, there is no real reason to leave out the viscous contribution although it does obscure the behavior we see in the experiment. Still, the maximum in the viscosity is interpreted as a signature of the transition from tumbling to a wagging state.
There are three observations to keep in mind when considering *fd* in the nematic phase under shear flow, which all point to very low stresses in such systems when compared to polymeric liquid crystals. First, the viscosity of *fd* in the nematic phase is two to three orders of magnitude lower than the viscosity of typical polymeric liquid crystals such as poly(benzyl glutamate) (PBG) Vermant94b , although the difference in solvent viscosity is only one order of magnitude. Second, the range over which the viscosity of *fd* suspension varies is more limited with changing shear rate and rod concentration: the viscosity lies between 70 times the solvent viscosity for low shear rate and low rod concentration and 20 times the solvent viscosity for high shear rate and rod concentration. Moreover, the viscosity as calculated from the equation of motion of the ODF is of the same order as the measured viscosity. Third, polymer nematics exhibit negative first normal stress differences for certain shear rates as was first observed for PBG solution Kiss78 . This is a direct consequence of the tumbling of the nematic director. Attempts have been made to measure the first normal stress difference for nematic fd solutions but due to very low force the signals were too small to be measured.
### IV.2 Flow reversal experiments
In flow reversal experiments, the sample is first sheared at a constant shear rate in one direction until the steady state condition is reached. Subsequently, the direction of flow is suddenly reversed while keeping the magnitude of shear rate constant. Such experiments have been very useful in characterizing and understanding the dynamics of sheared liquid crystalline polymers Moldenaers86 . In the present work, flow reversal experiments were performed covering a wide range of shear rates and fd concentrations. Typical flow reversal experiments are depicted in Fig. 5 for a fd concentration of 11.5 mg/ml which corresponds to $`c/c_{IN}=1.05`$. At the lowest shear rates a damped oscillatory response is obtained which decays within few oscillations (Fig. 5a). Increasing the shear rate results in a more pronounced oscillatory response, which damps out relatively slowly. The oscillatory response in Fig. 5b is most pronounced at a shear rates of 12 s<sup>-1</sup>. At even higher shear rates, the damping again increases (Fig. 5c). In order to quantitatively characterize the response to a nematic to a flow reversal, the data is fitted to a damped sinusoidal superimposed onto a asymptotically decaying function of the following form:
$`\eta (t)=\eta _{stat}\{1+Ae^{\frac{\dot{\gamma }t}{\tau _d}}\mathrm{sin}(2\pi {\displaystyle \frac{\dot{\gamma }t\phi }{P}})\}(1bg^{\dot{\gamma }t}.)`$ (10)
This is an empirical choice, but each variable in the fit contains important information about the behavior of rods in shear flow. Fig. 6 shows the behavior of fit parameters as a function of the shear rate at few selected concentrations of fd virus. In this figure we indicate with vertical dashed lines the shear rates at which the steady state viscosity exhibits a local maximum for four different concentrations. Interestingly, these are exactly the same shear rates at which the damping constant $`\tau _d`$ as well as the tumbling period $`P`$ show a sharp increase. The asymptotic constant $`b`$, on the contrary, shows a decrease. These features disappear for the highest fd concentration. Presumably the three regions showing different flow reversal behavior correspond to tumbling, wagging and flow aligning regime. This will be discussed in more detail in section V.2. In the next section we first discuss the concentration and shear rate dependence of the tumbling period in the regime where rods exhibit tumbling flow behavior.
### IV.3 Tumbling period as a function of shear rate and rod concentration
DEH theory predicts that as the De (or Pe) number is increased, that the ’molecular’ period of oscillation decreases with increasing shear rate in the tumbling regime Larson90 . This feature was never fully explored, since in most polymeric liquid crystals it was found that the tumbling period was strain scaling, implying that the response overlaps when the period is scaled with the applied shear rate and the stress is normalized by its steady state value. The strain scaling arises as a consequence of the presence of a large, non-Brownian, length scale in the sample that determines the time response, even at relatively high De (or Pe) numbers. This most probably is the domain size characterizing the nematic texture. The log-log plot of the tumbling period (T=P/$`\dot{\gamma }`$) as a function of the shear rate is shown in Fig. 7. Here the data are only shown for a low shear rate region which is associated with the tumbling region. Strain scaling, if present, would give a slope of -1. However, as can be seen in the inset of Fig.7, the reciprocal indicating strain scaling is only approached and not reached at the highest rod concentration studied here.
The shear rate dependence of the tumbling period is compared to the theoretical prediction for the same rod concentration as well as the same order parameter, see Fig. 8. The reason for using the order parameter to assess the theoretical predictions was discussed at length in section III.1. For purposes of comparison the order parameter was obtained form x-ray experiments and the value of $`D_{\text{inf}}^0=20s^1`$. We emphasize that DEH theory is microscopic and that there are no adjustable parameters in the comparison between theory and experiments. Clearly there is a qualitative correspondence between theory and experiment, both showing a continuous decrease of the period. The quantitative correspondence, on the other hand, is limited. This is probably due to fact that texture, although not dominating the response, is still present. It will be shown later in section V.2 that the shear rate and rod concentration dependence of a tumbling to wagging and wagging to flow-aligning transition agree much better with DEH theory.
The concentration dependence of the tumbling period is shown in Fig. 9. Here, theory and experiments are compared at a fixed shear rate at which the tumbling to wagging flow transition occurs. The tumbling period increases with increasing rod concentration (Fig. 9a) or, equivalently, increasing order parameter of the nematic phase (Fig. 9b). The increase of the tumbling period with increasing order parameter was already predicted using a linearized version of the DEH theoryKuzuu84 .
In conclusion, the absence of strain scaling of the tumbling period and the qualitative agreement between theory and experiment the tumbling period indicates that the response of the suspension of *fd*-virus is dominated by the molecular elasticity arising from the distortion of the ODF of particles.
### IV.4 Relaxation at high concentration
In order to measure the relative magnitude of the elastic texture contribution to the overall stress, relaxation experiments were performed. For polymeric liquid crystals like polybenzylglutamate (PBG) solutions in m-cresol, Walker et al. Walker96b showed that there are three different regimes of relaxation behavior, each of which is related to a distinct structural relaxation. There is a ”fast” relaxation of the nematic fluid; a ”slower” relaxation that exhibits scaling with the shear rate before the cessation of flow, which is due to the indirect contribution of the texture to the overall stress; and a ”long-time” relaxation due to the reorganization of the texture on a supra-molecular level which will not be addressed here.
Stress relaxation experiments were performed in the low shear rate “tumbling” region, at shear rates smaller than those corresponding to the maximum in the viscosity. The sample used had a relatively high *fd* concentration of $`25mg/ml`$, corresponding with $`c/c_{IN}=2.3`$. Some typical responses to the cessation of flow are depicted in Fig. 10. The stress is normalized to its value before the cessation of flow, and the time axis is scaled by the shear rate. The fast component of the decay takes place at less than a tenth of a second, which is comparable to the response of the force re-balanced transducer and therefore not shown. The slow component of the stress relaxation scales when time is multiplied with the previous shear rate, but only from the point that the stress has decayed to less than 30 % of its original value, or less for higher initial shear rates. From Fig. 10 it can be concluded that the contribution to the stress for the highest concentration used and for low shear rates is 30 %. This is the absolute upper limit for the samples used in this paper. It should be noted that for PBG solutions 30 % it was found to be the lower limit Walker96b .
### IV.5 *In situ* microscopy
The flow-induced changes of the liquid crystalline texture during steady state shear flow were studied using a plate-plate geometry in combination with a polarization microscope. Measurements were performed for *fd* concentrations of 14 mg/ml and 25 mg/ml. Typical images are shown in Fig. 11 for different shear rates. Interestingly the characteristic size of the ”domains” was very large. Birefringent regions of up to half a millimeter were observed under static conditions. When the sample is subjected to shear flow, these domains will elongate and eventually disappear, at values of the shear rate which correspond to the maximum in the viscosity (see Fig.6a). An important difference between the two concentrations is that the elongated domains merge into bands for high rod concentration, whereas for the low concentration the structure disappears before such bands are formed. Interestingly, this transition to a banded structure in the high concentration fluid takes place at a shear rate which is higher than the shear rate where the low concentration fluid loses its features.
## V Discussion
When comparing the flow behavior of the polymeric nematic phase and the colloidal nematic phase of the dispersed *fd* viruses, the most striking observation is the qualitative agreement between the two systems, despite the fact that *fd* is an order of magnitude larger. The viscosity of the fd nematic is much smaller and the rotational diffusion of *fd* is much slower when compared to polymeric liquid crystals. Flow reversal experiments reveal typical transitions in the transient rheological behavior: damped oscillations occur at low shear rates changing to undamped oscillations at intermediate shear rate, which disappear if the shear rate is increased even further; the time scale of the oscillations of the stress transients is comparable. Also other well known phenomena like the formation of very large bands upon cessation of flow along the vorticity direction which have been studied in detail in polymeric systems Vermant94 , can also be observed here (data not shown).
Having established that *fd* virus dispersions indeed undergo a tumbling motion under flow, the *dynamic* behavior of *fd* suspensions can be rationalized on the basis of the microscopic theoretical predictions for a homogeneous system of rods under shear. Doing so, one important prerequisite needs to be fulfilled, namely that the dominating contribution to the stress is coming from the nematic fluid and not from the texture. It will be argued here that this indeed is the case. Having done so, we will be able to map out a phase diagram of the dynamic transitions from tumbling to wagging to flow aligning.
### V.1 Textural evolution during flow
The word “texture” refers to disclination points and lines where the director of the nematic phase changes discontinuously, marking domains in the sample. When a system containing these domains and disclinations is subjected to shear flow, part of the dissipated energy is used to destroy these structures. Fig. 11 shows that the domains tend to elongate and align with the flow. Disclinations can also cause a direct contribution to the total stress resulting in a high viscosity and a very pronounced shear thinning behavior, typically referred to as region I Walker94 . Experiments on polymeric liquid crystals have revealed several features of the flow behavior of nematic liquid crystals which are attributed to the presence of texture in the nematic phase. Tumbling induces distortions in the director field and the defects arrest the tumbling, thereby inducing an elastic stress. The length scale over which this distortion occurs, i.e. the ’domain’ length scale, is an inherent non-Brownian length scale, see ref. Burghardt90 . As a consequence, stress patterns during flow reversal will display strain scaling. Also the damping of the oscillations is explained on the base of the presence of the polydomain structure, where e.g. the ”friction” between the domains would lead to a damping of the oscillations Larson91 ; Kawaguchi99 . The scaling of the stress relaxation process after the flow is stopped with shear rate has been explained using the same arguments. From such an experiment the relative contribution to the total stress of a homogeneous nematic phase and the polydomain texture can be estimated since the relaxation dynamics of the nematic phase is much faster than that of polydomain structure Walker96b .
The micrographs in Fig. 11 clearly reveal that texture under flow exists in nematic *fd* dispersions. Their contribution to the rheology is, however, far less prominent when compared to polymeric liquid crystals such as PBG. This we can infer from several observations. First, very moderate shear thinning is observed in the low shear rate regime for the *low* concentrations, which gradually disappears with increasing concentration (Fig. 6b). This is very similar to theoretical predictions for a homogeneous nematic phase (Fig. 2b in Ref. Dhont03c ). Also, the calculated and measured viscosity are of the same order of magnitude. In contrast, shear thinning can be fairly strong in the low shear rate region (Region I) where texture dominates the response and it will increase with increasing concentration Marrucci93 , although also other microstructural features can contribute here BurghardtRI . Second, the tumbling period is not strain scaling (Fig. 7), which could be due to either a smaller relative magnitude of the textural stress or due to the fact that we are not in a low enough Pe regime. Third, ’strain’ scaling is recovered for the slow ’textural’ relaxation process after the flow has been stopped. This experiment shows that at the highest rod concentrations used and at low shear rates the distortional textural contribution is about 30 %. For most experiments done this value is probably significantly lower. So, where texture is important, even dominating the stress response for molecular LCPs, molecular elasticity is far more dominating the *fd*-virus. Though we just argued that the texture does not dominate the shear response of the system, this does not mean that the shear response is not influenced by texture. For one the oscillations we observe are still strongly damped, and the damping only decreases when the transition to the flow aligning state is reached (see the behavior of $`\tau _d`$ in Fig. 6). Moreover, the presence of texture might explain the discrepancy in the behavior of the period of the oscillations between experiment and theory (Fig. 9). Most importantly, we know from microscopy that texture is present under shear (see Fig. 11). It should be noted, however, that the size of the polydomain structure of the *fd* dispersions is one order of magnitude bigger as compared to PBG Vermant94b , so that the density of disclination lines and points is about three orders of magnitude lower for *fd*. Note that the length scale of the texture during flow is still small compared to the dimension of the flow cell. Since the contribution of texture scales with the density of the disclinations Marrucci93 , texture will be far more dominating for e.g. PBG than for *fd*, even when elastic constants are almost the same for the two systems Dogic00 ,taratuta85 ).
### V.2 Phase diagram of dynamical flow transitions
In this section the experimental results or combined and a non-equilibrium phase diagram of fd rods under shear flow is presented. The results for the four fit parameters plotted in Fig. 6 show clear transitions at well defined shear rates for all fd concentration. Although they only give an indirect proof of the transitions, they can be used to infer information about the flow transitions. For all fd concentrations (except for the highest one) the shear rate where the maximum viscosity is reached is identical with the shear rate where the period as well as the damping constant start to increase (indicated by the vertical dashed lines in Fig. 6). The microscopic observations are in fairly good agreement with the transitions inferred from the rheology. Upon approaching the tumbling to wagging transition from tumbling to flow aligning, the texture becomes to faint to resolve in the microscope and texture subsequently disappears upon reaching the FA region. For the high fd concentration, i.e. the sample showing shear banding (Fig. 11 last), one can identify a sharp transition from a structured to an unstructured region in the same micrograph. Since this picture was taken in the plate-plate geometry, there is a shear rate distribution across the image: the shear rate is increasing going from the left side to the right. A sharp spatial transition therefore also represents a sharp transition at a given shear rate. Although, due to the method of zero gap-setting, the value of the shear rate is not exactly known ($`\pm 20\%`$), one can still identify the shear rate where structure disappears as the shear rate where the viscosity reaches its local maximum (the down pointing triangles in Fig. 6a). For low fd concentration of ($`14mg/ml`$) the structure disappears around the point where the viscosity reaches its local maximum, although the morphological transition for the lower concentration is less abrupt.
Fig. 12 shows the behavior of flow transitions as a function of shear rate for various fd concentrations. For the experiment we plotted the Péclet numbers where the viscosity shows a local maximum and where the damping constant reaches a maximum. The theoretical predictions for the tumbling to wagging and wagging to flow aligning transitions are obtained from the plots the angle of the nematic director $`\theta `$ under flow, see Fig. 2. Similar to the method used in Fig. 9, the experimental concentration is scaled to the theoretical concentration in two different ways: effective concentration (Fig. 12 a) and the order parameter $`P_2`$ (Fig. 12 b). This figure was shown in a preliminary paper without a detailed explanation Lettinga04a . The shear rate is rescaled to the Péclet number by using the rotational diffusion coefficient at infinite dilution. Fig. 12 allows us to draw some important conclusions. First, it is clear that scaling the concentration with the equilibrium order parameter gives better agreement when compared to the scaling by the dimensionless concentration. The fact that theory and experiment agree without using any fitting parameters ($`P_2`$ was obtained in a separate experiments Purdy03 ) leads to the conclusion that the DEH theory describes the flow behavior of the fd nematics quite well, as long as the effects of flexibility and charge of the experimental rods are included in the calculation of the order parameter. Less convincing agreement is obtained when comparing the experimental and theoretically calculated periods (Fig. 9). The reason for this could be the remaining textural contribution to the overall stress which, although small, cannot be neglected. Since we deduce from Fig. 12b, that a dimensionless concentration of $`\varphi \frac{L}{D}=4`$ corresponds with a *fd* concentration of $`16mg/ml`$, we used this number *a posteriori* to scale the calculated molecular viscosity in Fig. 4. The pure elastic contribution shows a very nice quantitative correspondence with the experimental data. Interestingly, when the viscous term is added, the theoretical viscosity is higher than the experimental viscosity, despite of the fact that no hydrodynamics is incorporated.
In the previous subsection it was argued that the influence of textural contribution to the stress tensor of fd are relatively small, as compared to PLC’s. There are however strong indications that the dynamic behavior is influenced by the macroscopic bands which are formed for the samples at the highest concentrations used (see Fig. 11 end). As can be seen in Fig. 6, the typical features for the transition to wagging disappear: there is no increase in the damping constant, nor in the period of the oscillations. Moreover, the theory shows only a moderate hesitation of the stationary viscosity (which even disappears then the viscous term is added, Fig. 4), whereas in experiments a local peak is observed which is more pronounced with increasing concentration. The microscopy pictures results that at high concentrations the systems finds another way to handle the distortion of the particle distribution at high shear rates by forming shear bands where the overall orientational distribution is alternating, as was already observed and partially explained for the polymeric systems Larson93a ; Larson93b ; Vermant94b . In the present work, the concentration dependence of the phenomenon at hand suggests that this merits further experimental as well as theoretical work. In this context one should not forget that we compar experiments on charged and semi-flexible fd with theory for hard and stiff rods. It could well be that these factors also play an important role. It will be a major challenge especially to take the semi-flexibility into account in the equation of motion.
## VI Conclusions
Colloidal suspensions of rod-like *fd* viruses are an ideal model system to study the behavior of the nematic liquid crystalline phase under shear flow. Flow reversal experiments show signatures for tumbling, wagging, and flow aligning behavior, very similar to the behavior found in polymeric liquid crystals. The rigid rod nature of the *fd* suspension, possibly combined with a smaller relative textural contributions to the overall stress tensor make *fd*-virus a suitable model system for the DEH theory. Important in this respect is that the overall viscosity is only one to two orders of magnitude higher than the solvent viscosity. Also it is important to note that stress relaxation experiments combined with the absence of strain scaling in flow reversal experiments suggest that there is only a limited contribution of textural aspects to the overall stress, even for the highest fd concentration used in this work. The shear thickening of the viscosity observed for a range of fd concentrations is as yet, unexplained. The maximum in the viscosity occurs at the critical shear rate where the tumbling to wagging transition takes place. Microscopic observations show that at this shear rate the morphological features disappear, suggesting a strong connection between the dynamic transitions and structure formation.
The experimental results have been compared to a microscopic theory for rod like molecules subjected to shear flow. A non-equilibrium phase diagram is constructed, describing the transitions from tumbling to wagging and from wagging to flow-aligning as a function of rod concentration and applied shear stress. When scaling the results to the concentration where the isotropic-nematic transition takes place, the experiment and theory show only a qualitative agreement, possibly due to the fact that the real rods are are both semi-flexible and charged. However, when scaling the results using the order parameter, which is determined by the interactions between the rods, theory and experiment show an excellent agreement without using any fit parameters. Thus, it can be concluded that the Doi-Edward-Hess theory accurately captures the dynamic features of a hard rod system. *Fd* dispersions constitute such a hard rod system as long as flexibility and charge are properly taken into account, which can be simply achieved by using the order parameter to scale the data. More theoretical work is needed, however, to explain the clear connection between the observed band formation at high concentrations and the dynamic transitions, and to incorporate the effect of flexibility of the rods.
## Acknowledgement
We thank Jan Dhont for many discussions and critical reading of the manuscript. Pier-Luca Maffetonne is acknowledged for stimulating remarks. MPL is supported by the Transregio SFB TR6, ”Physics of colloidal dispersions in external fields”. ZD is supported by Junior Fellowship at Rowland Institute at Harvard. The authors acknowledge support of the EU (6th FP) in the framework of the Network of Excelence ”SOFTCOMP”.
|
warning/0506/cond-mat0506001.html
|
ar5iv
|
text
|
# Selective epitaxial growth of sub-micron complex oxide structures by amorphous SrTiO3
## Abstract
A chemical-free technique for fabricating submicron complex oxide structures has been developed based on selective epitaxial growth. The crystallinity and hence the conductivity of the complex oxide is inhibited by amorphous SrTiO<sub>3</sub> (STO). Using a combination of pulsed laser deposition and electron-beam lithography, amorphous STO barriers are first deposited on a single crystal substrate. A thin film is then deposited on the patterned substrate with the amorphous STO barriers acting to electrically and physically isolate different regions of the film. Since no chemical or physical etchants come in contact with the deposited film, its integrity and stability are preserved. This technique has successfully produced sub-micron YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> and La<sub>2/3</sub>Ca<sub>1/3</sub>MnO<sub>3</sub> structures.
High-Tc, Superconductivity, Nano, Nanostructures, Fabrication, Oxides
Fabrication of complex oxide nanostructures is important from both a physical and technological standpoint. Many complex oxides exhibit interesting electrical behavior including high-*T<sub>c</sub>* superconductivity, colossal magnetoresistance (CMR) and ferroelectricity. Potential device applications including high sensitivity sensors,Braginski-FED1:1990 flux flow transistors,Martens-IEEE1:1991 ferroelectric field effect transistorsRamesh-SCI252:1991 and magnetic memoryGrishen-APL74:1999 require suitable microfabrication techniques. However, due to inherent stoichiometric and structural complexities associated with complex oxides, producing high quality sub-micron structures has proven difficult because conventional methods tend to require etching of the complex oxide.
In this letter we present a technique for fabricating sub-micron complex oxide structures which requires no chemical or physical etching. Based on selective epitaxial growth, a single crystal substrate is patterned both vertically and laterally by the deposition of amorphous SrTiO<sub>3</sub> (STO) barriers. The amorphous nature of the STO barriers acts to ensure that any material deposited on top of the barriers is electrically insulating. This technique has the advantage that complex structures can be fabricated without any degradation due to chemical or physical etching. The minimum dimensions of the microstructure are well defined by the amorphous STO barriers and are not determined by diffusive processes. This technique also allows for the deposition of passivation layers which can improve the stability of the complex oxide.Copetti-APL61:1992
Conventional techniques for fabricating sub-micron complex oxide structures involve either post-deposition etching of the oxide or patterning of the substrate before deposition. Post-deposition patterning has commonly been obtained using wet-chemical etching of the oxide in solutions of bromine in ethanol,Vasquez-APL53:1988 ethylenediaminetetraacetic acid (EDTA),EDTA phosphoric acidLyons-IEEE27:1991 or hydrofluoric acid.Eidelloth-APL59:1991 However, wet chemical etching of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (YBCO) thin films has shown to cause a significant increase in the high-frequency surface resistance as well as to cause insulating dead layers and a change in the surface morphology at the exposed surfaces. Roshko-IEEE5:1995 Physical etching methods such as reactive ion etching RIE , focused ion beam etching FIB , and pulsed laser etching Inam-APL51:1987 have also been used to pattern complex oxide microstructures. However, these methods can cause physical damage to the exposed surfaces of the thin film. Also, heat generated from these physical processes can be sufficiently large to alter the doping and hence their electrical transport properties of the oxide.
Sub-micron structures can also be fabricated by patterning a single crystal substrate prior to deposition. The substrate is patterned such that regions of the deposited film are physically separated through the creation of ridge or trench structures Mohanty-PHYC408:2004 or electrically isolated through the inhibition of the conductivity of select regions. The inhibition of conductivity can be achieved by destroying the local crystallinity of select regions. This can be achieved by diffusion of Si Ma-APL55:1989 , SiO<sub>2</sub> Copetti-APL61:1992 , Si<sub>x</sub>N<sub>y</sub> Kern-JVSB9:1991 or Ti Rossel-PHYC185:1991 or by selectively determining where epitaxial growth can occur through the deposition of a Ti or W layer on selective regions of the substrate. Damen-SST11:1998
Our chemical-free procedure for fabricating complex sub-micron structures is outlined schematically in figure 1. First, a polymethylmethacrylate (PMMA) mask is defined on a single crystal STO substrate by electron-beam lithography (EBL). A $``$700 nm layer of 7% PMMA in anisole with a molecular weight of 450 amu is spun onto a single crystal STO substrate at 4500 rpm for 45 s. The resist is then baked on a hotplate at 180<sup>o</sup>C for 5 minutes. The resist is then exposed by EBL using a JEOL IC-848A, tungsten filament scanning electron microscope at an accelerating voltage of 30 kV, probe current of 8 pA and a line dosage of 0.9 nC/cm. The low probe current was used to allow sufficient time for any accumulated charge to dissipate, in order to reduce any electron charging effects that could readily occur when imaging an insulating material. The resist is then developed in a 3:1 solution of isopropyl alcohol (IPA) to methyl isobutyl ketone (MIBK) for 30 seconds and then rinsed in IPA.
After the mask has been defined, a layer of amorphous STO is then deposited by pulsed laser deposition (PLD) using a 248 nm KrF excimer laser at a laser energy density of $``$1.2 J/cm<sup>2</sup>. The STO is pulsed laser ablated at an oxygen partial pressure of 20 mTorr and at ambient temperature to ensure no evaporation or diffusion of the PMMA mask occurs. An ablation time of 30 min using a pulsed laser repetition rate of 10 Hz results in the deposition of $``$350 nm of amorphous STO. The remaining PMMA is then removed by acetone in an ultrasonic bath for 45 s, leaving behind an inhibitive pattern of amorphous STO. A complex oxide film, such as YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (YBCO) or La<sub>2/3</sub>Ca<sub>1/3</sub>MnO<sub>3</sub> (LCMO) can then be deposited by PLD using typical ablation conditions. Since the deposition of the complex oxide is the last step in the process, the integrity of the resulting structure is preserved as it does not undergo subsequent heating due to physical bombardment or come in contact with chemical etchants.
The electrical transport properties of our fabricated microstructures were characterized by standard four-point, phase sensitive, AC resistance versus temperature measurements, as well as synchronous pulsed current versus voltage ($`I`$-$`V`$) and voltage versus current ($`V`$-$`I`$) measurements. In the latter technique, pulses of 200 $`\mu `$s with a duty cycle of 5% were used in order to minimize any effects of Joule heating.
Sub-micron YBCO strips fabricated using our technique exhibit superconducting critical temperatures, $`T_c`$, from 80-92 K, with transition widths of 1-10 K. YBCO strips show room temperature resistivities of tens to several hundreds of $`\mu \mathrm{\Omega }`$cm and critical current densities, j<sub>c</sub>, on the order of 10<sup>7</sup> A/cm<sup>2</sup>, which approach those of unpatterned YBCO thin films. To ensure that an applied current was indeed confined to the selective regions of the film, we also measured the YBCO film that was deposited on the amorphous STO barriers. These control samples show no superconducting transition and exhibit room temperature resistivities on the order of 10<sup>6</sup> $`\mu \mathrm{\Omega }`$cm. An atomic force micrograph (AFM) shown in figure 1 confirms the amorphous nature of the deposited YBCO.
Previous transport studies on chemically-patterned cuprate samples ranging from 100 nm to 2 $`\mu `$m in width have been reported. Highly non-linear $`I`$-$`V`$ characteristics and anomalous resistance versus temperature behavior have been attributed to either collective flux-flow, phase slips or mesoscopic domains Rogalla ; Dmitriev ; Jelila ; AbdelhadiJung ; Bonetti . The $`V`$-$`I`$ characteristics of YBCO microstrips fabricated using our technique are shown in Figure 2. For wider strips, the $`V`$-$`I`$ characteristic can be described by a power law relationship caused by thermally activated flux creep and flux flow. However, as the applied current becomes more confined in a narrower strip, the power law relation is only valid below some threshold value of applied current. At this threshold value, a discontinuity in the voltage occurs and the $`VI`$ relationship becomes linear. An $`I`$-$`V`$ measurement of the same microstrip shows that voltage discontinuity coincides with an $`s`$-shaped, negative conductance region (see inset Figure 2). Both the discontinuity of the voltage in the $`V`$-$`I`$ characteristic and the $`s`$-shaped $`I`$-$`V`$ characteristic of the microstrip are consistent with phase-slip behavior due to the lateral confinement of the superconductor. SBT ; Vodolazov An investigation of the phase-slip behavior in our YBCO microstrips as well as information derived from such behavior on the dynamics of the superconducting order parameter is the focus of separate papers. Morales-NanoYBCO ; Morales-Tau\_Psi\_Phi
Sub-micron LCMO strips were also fabricated using our technique. Sub-micron LCMO strips are of interest because their magnetoresistive properties can be tuned by externally applied strain. Koo-APL71:1997 Resistance versus temperature measurements of strips of differing widths at zero field are shown in figure 3. Fabricated sub-micron LCMO strips show a maximum in resistance, $`T_m`$, ranging from 220 K to 300 K. The position of $`T_m`$ has been shown to be highly dependant on stress relaxation and improved crystallinity resulting from grain growth. Thomas-JAP84:1998 A field dependent study of the CMR properties of sub-micron LCMO strips fabricated using our technique will be presented in a separate paper. Morales-LCMO
In summary, we have developed a novel chemical-free technique for fabricating sub-micron complex oxide structures. The technique is based on selective epitaxial growth of a complex oxide thin film. The crystallinity and hence the conductivity of the complex oxide thin film are inhibited by amorphous SrTiO<sub>3</sub> barriers deposited upon a single crystal substrate. This technique has been successfully applied to fabricate strips of sub-micron high-*T<sub>c</sub>* superconducting YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> and colossal magnetoresistive La<sub>2/3</sub>Ca<sub>1/3</sub>MnO<sub>3</sub>.
The authors acknowledge assistance by Stephanie Chiu and Eugenia Tam and funding from: NSERC, CFI/OIT, MMO/EMK and the Canadian Institute for Advanced Research under the Quantum Materials Program.
|
warning/0506/quant-ph0506269.html
|
ar5iv
|
text
|
# Linear Optics C-Phase gate made simple
## Abstract
Linear optics quantum logic gates are the best tool to generate multi-photon entanglement. Simplifying a recent approach Ralph\_PG we were able to implement the conditional phase gate with only one second order interference at a polarization dependent beam splitter, thereby significantly increasing its stability. The improved quality of the gate is evaluated by analysing its entangling capability and by performing full process tomography. The achieved results ensure that this device is well suited for implementation in various multi photon quantum information protocols.
The quantum computer is one of the most promising and desirable goals in quantum information science. Its implementation relies strongly on the capability to engineer entanglement in the quantum system of choice. For qubits it was shown that entangling gates (like the C-phase gate or the CNOT) together with single qubit operations are sufficient to create any kind of quantum network.
Photons are well suited for quantum information tasks, as their interaction with the environment is small guaranteeing low decoherence. While the creation of entangled photon pairs via spontaneous parametric down conversion became a standard technique, its control is still a major challenge, mainly due to low nonlinear interaction efficiencies. One solution to this problem is using linear optics components and introducing the nonlinearity via ancillary single photons and photon number resolving detectors KLM . Initial demonstrations showed that such gates can be implemented, once the necessary sources and detectors become available on a larger scale KLM-exp . Another solution, requiring much less resources becomes possible if one focuses on performing only a limited number of quantum logic operations. Then one can control the action of the gate by conditioning it to the detection of one photon in each of the output ports. This will occur only with a certain probability, which, however, might be larger than the one achieveable with the first method and equivalent resources. In particular, a controlled phase (C-Phase) gate was introduced recently Ralph\_PG , which uses a combination of first and second order interference to obtain C-Phase action in 1/9 of the cases. Yet, since first-order interference requires stability of the setup on the order of less than the photon’s wavelength, for multi-photon experiments schmid more simple and stable implementations surely are desirable.
Here we introduce a linear optics C-Phase gate, which uses only a single two-photon interference at a polarization dependent beam splitter. The stability requirements are thereby relaxed to the coherence length of the detected photons ($`150\mu m`$) and can easily be fulfilled without additional stabilization equipment. To characterize the C-phase gate we first study the entangling capability of the gate by determining the fidelity and negativity of the output for four different input states. Second, we use linear quantum process tomography (QPT) linQPT ; Steinberg to analyze the gate operation. As imperfect interference reduces the quality of the gate and induces state-selective incoherence we had to account for the non-trace-preserving character of the gate. Instead of the usual maximum likelihood approach, we use prior knowledge of the intrinsic features of our setup, in order to obtain physical and easily understandable parameters for characterizing the gate and estimating its performance.
The ideal C-Phase gate acts on two-qubit input states
$`\psi _{in}`$ $`=`$ $`(c_{HH}HH+c_{HV}HV`$
$`+`$ $`c_{VH}VH+c_{VV}VV),`$
and applies a relative $`\pi `$-phaseshift to the contribution VV only, such that
$`\psi _{out}`$ $`=`$ $`(c_{HH}HH+c_{HV}HV`$
$`+`$ $`c_{VH}VHc_{VV}VV).`$
Here we encode the logical 0 (1) in linear horizontal $`H`$ (vertical $`V`$) polarization of single photons. $`c_{HH}`$ denotes the amplitude of the $`HH`$-term (for the other terms accordingly).
The application of the phase shift relies on second-order interference of indistinguishable photons at a polarization dependent beam splitter (PDBS) (Fig.1) HOM ; CST . Two input modes a and b are overlapped at $`\mathrm{PDBS}_\mathrm{O}`$. The transmission of $`1/3`$ for vertical polarization results in a total amplitude of $`1/3`$ for the $`VV`$ output terms, as can be seen by adding the amplitudes for a coincident detection:
$$(t_V^at_V^b)+(ir_V^air_V^b)=\sqrt{\frac{1}{3}}\sqrt{\frac{1}{3}}\sqrt{\frac{2}{3}}\sqrt{\frac{2}{3}}=1/3$$
(3)
where $`t_i^x`$ ($`r_i^x`$) is the amplitude for transmission (reflection) of state $`i`$ in mode $`x`$. Perfect transmission of horizontal polarization causes that no interference happens on the contributions $`HH,HV`$ and $`VH`$. As the absolute values of the amplitudes need to be equal for any input we still need to attenuate the contributions that include horizontal polarization. This is achieved by $`\mathrm{PDBS}_{\mathrm{a}/\mathrm{b}}`$ with the transmission $`1/3`$ for horizontal polarization and perfect transmission for vertical polarization in both output modes. All together we find a probability of $`1/9`$ to obtain a coincidence in the outputs and thus a gate operation with perfect fidelity.
Working with real components results in deviations from the theoretical derivation. A detailed calculation with arbitrary transmission and reflection amplitudes at $`\mathrm{PDBS}_\mathrm{O}`$ and $`\mathrm{PDBS}_{\mathrm{a},\mathrm{b}}`$ shows how their parameters influence the gate operation. In general we obtain from $`\psi _{in}`$
$$\begin{array}{c}\psi _{out}=(c_{HH}t_H^at_H^ba_Hb_Hc_{HH}r_H^ar_H^ba_Hb_H)HH\hfill \\ +(c_{HV}t_H^at_V^ba_Hb_Vc_{VH}r_V^ar_H^ba_Hb_V)HV\hfill \\ +(c_{VH}t_V^at_H^ba_Vb_Hc_{HV}r_H^ar_V^ba_Vb_H)VH\hfill \\ +(c_{VV}t_V^at_V^ba_Vb_Vc_{VV}r_V^ar_V^ba_Vb_V)VV),\hfill \end{array}$$
(4)
where $`a_i`$ ($`b_i`$) are the transmission amplitudes of $`i`$ in mode a (b). To obtain the expected C-phase gate operation one has to fulfill several conditions, which give an insight in how the setup has to be built. First, $`(r_V^ar_V^b)/(t_V^at_V^b)=2`$, which is approximately achieved by slightly varying the angle of incidence at $`\mathrm{PDBS}_\mathrm{O}`$. Experimentally we reach a value of $`2.018\pm 0.003`$. Second, $`r_H^a=0=r_H^b`$, which requires the reflection of the horizontal polarization at the overlap beam splitter to be zero. The third condition, $`t_H^aa_H=t_V^aa_V`$, and $`t_H^bb_H=t_V^bb_V`$, respectively, determines the setting for the attenuation at $`\mathrm{PDBS}_{\mathrm{a},\mathrm{b}}`$.
To experimentally test the gate operation we used photon pairs emitted from spontaneous parametric down conversion. A 2 mm thick BBO ($`\beta `$-Barium Borate) crystal was pumped by UV pump pulses with a central wavelength of 390 nm and an average power of 700 mW from a frequency-doubled mode-locked Ti:sapphire laser (pulse length 130 fs). The pulsed operation is not necessary when working only with photon pairs, but as the gate is intended to work in multi-photon applications we preferred to characterize it for this mode of operation. The emission is filtered with polarizers to prepare input product states with high quality. We couple the photon pairs into single mode fibers for selection of the spatial modes. This guarantees identical beam modes which eases the alignment of spatial mode matching at $`\mathrm{PDBS}_\mathrm{O}`$. The spectral mode selection is improved via 2nm bandwidth filters behind the gate.
To ensure the same optical path length between the crystal and the overlap beam splitter for both photons, one of the output couplers of the single mode fibers is mounted on a translation stage. The position of zero delay is determined from the minimum of the coincidence rate for $`VV`$-input (Hong-Ou-Mandel HOM , ”HOM”, Dip Fig. 2). In each output mode of the C-Phase gate the polarization is analyzed via quarter and half waveplates and a polarizing beam splitter with single photon avalanche photo diodes. For the analysis of the final two-photon states the coincidence count rates for each of the four contributions have to be corrected for the different detector efficiencies. The errors on all quantities are deduced from propagated Poissonian counting statistics of the raw detection events and efficiencies.
The HOM-measurement shown in Fig. 2 also gives information about the indistinguishability of the photons at the PDBS. For large delay, the two photons are completely distinguishable due to their time of arrival. The probability to get a coincidence from a $`VV`$-input is then $`5/9`$. In case of perfectly indistinguishable photons at zero delay, the probability drops to $`1/9`$. The Dip-Visibility is defined by $`𝒱=\left(c_{\mathrm{}}c_0\right)/c_{\mathrm{}}`$, where $`c_0`$ is the count rate at zero delay and $`c_{\mathrm{}}`$ at positions with big delay. From the above considerations we obtain a theoretical value of $`𝒱_{th}=80\%`$, and experimentally, applying least-square fit, we find $`𝒱_{exp}=72.8\%\pm 0.7\%`$. We call $`𝒬=𝒱_{exp}/𝒱_{th}=91.0\%\pm 0.9\%`$ the overlap quality. One can conclude that the amount of additional $`VV`$$`VV`$-noise depends on the input, but is $`9\%`$ at maximum.
As a first step in the analysis of the performance of our gate we look at its capability to entangle. We choose $`++,+L,L+`$, and $`LL`$, with $`+=1/\sqrt{2}(H+V)`$ and $`L=1/\sqrt{2}(H+iV)`$, as input product states and perform state tomography on the output states MOQ . We use linear tomography as the resulting matrices all have eigenvalues greater or equal -0.02, i.e., are almost physical without corrections. For an ideal C-phase gate one would obtain a maximally entangled output for these input states, for example
$`L+`$ $`=`$ $`1/2(HH+iVH+HV+iVV)`$
$`\overline{PG}L+`$ $`=`$ $`1/2(HH+iVH+HViVV)`$ (5)
$`=`$ $`1/\sqrt{2}(LH+RV).`$
The experimentally observed fidelities relative to the expected output states are all better $`F_{exp}80.5\%\pm 0.6\%`$. Fig. 3 exemplarily shows the experimental result for $`L+`$-input ($`F_{exp}^{L+}=87.8\%\pm 0.6\%`$). Note that for states with a fidelity larger than $`(2+3\sqrt{2})/8=0.78`$ CHSH-inequalities are violated (CHSH ), which is the case for all of our examples. To quantify the entanglement we also calculated the logarithmic negativity — for all output states we find $`𝒩_{exp}0.73\pm 0.02`$ ($`𝒩_{exp}^{L+}=0.75\pm 0.02`$).
For a complete characterization of an arbitrary unknown process one can use quantum process tomography (QPT). For QPT the process is represented by a superoperator $``$ which is decomposed in a linear combination of a basis of unitary transformations $`E_i`$:
$$(\rho _{in})=\underset{i,j}{}\chi _{ij}E_i\rho _{in}E_j^{}$$
(6)
The matrix $`\chi `$ completely describes the process. In order to obtain all components $`\chi _{ij}`$, the normalized output density matrices $`\rho _{out}^k`$ for a tomographic set of, usually separable, input states are measured, in our case for the inputs ($`HH`$ , $`HV`$ , $`H+`$ , $`H`$ , $`VH`$ , $`VV`$ , $`V+`$ , $`VL`$ , $`+H`$ , $`+V`$ , $`++`$ , $`+L`$ , $`LH`$ , $`LV`$ , $`L+`$ , $`LL`$). As the contribution of the $`VV`$$`VV`$-noise is input state dependent our process is non-trace-preserving. This means that the outcomes occur with different probabilities $`p_k`$ linQPT for the different input states $`\rho _{in}^k`$ $`\mathrm{Tr}((\rho _{in}^k))`$:
$$\begin{array}{c}\rho _{out}^k=\frac{(\rho _{in}^k)}{\mathrm{Tr}((\rho _{in}^k))}\hfill \\ \hfill (\rho _{in}^k)=\mathrm{Tr}((\rho _{in}^k))\rho _{out}^k=p_k\rho _{out}^k.\end{array}$$
(7)
We determine these probabilities from the diagonal entries of all measured output density matrices. The normalized density matrices together with the corresponding probabilities can be used to evaluate $`\chi _{ij}`$ via Eq. (7) and (6). To account for the probabilistic nature of the gate an overall normalization is performed such that $`p_{HH}=1/9`$.
Fig. 4a shows the process matrix $`\chi _{th}`$ of the ideal linear optics phase gate. It represents the decomposition of the C-Phase gate into unitary operations, for our choice of $`E_i`$ resulting in
$$\overline{PG}_{ideal}=\left(1𝐥1𝐥+\sigma _z1𝐥+1𝐥\sigma _z\sigma _z\sigma _z\right)/3.$$
(8)
The four peaks in the diagonal of $`\chi _{th}`$ show the equal weights of the contributions, while the negative entries at the edges represent the negative sign at $`\sigma _z\sigma _z`$. This matrix now can be compared with the experimentally obtained one (Fig. 4b). Only the real parts are shown since the imaginary parts are close to zero (average $`0.0\pm 0.002`$). As the introduced noise is not too big, the experimentally measured process matrix still demonstrates nicely the features of the gate operation. The main differences are the lower non-diagonal terms indicating reduced coherence in the system. From the estimated process tomography matrix we calculated a process fidelity of $`F_p=Tr(\chi _{th}.\chi _{exp})/(\mathrm{Tr}(\chi _{th}).\mathrm{Tr}(\chi _{exp}))=81.8\%`$. Still, due to Poissonian counting statistics $`\chi _{exp}`$ has non-physical, negative eigenvalues, and the above value has to be treated with care. To circumvent this problem one can use the maximum likelihood approach, where a physical process matrix is fitted to the observed data. Yet, the process is not really unknown to us and we can try to describe it via a theoretical model according to equation 4.
The transformation of the phase gate consists of interference between both photons transmitted or both photons reflected $`\overline{PG}_{gen}=M_{tt}+M_{rr}`$, where both $`M_{tt}`$ and $`M_{rr}`$ are matrices with components given by the coefficients of Eq. (4). For simplicity we assume $`t_V^a=t_V^b`$ and $`r_H^a=r_H^b=0`$. $`M_{rr}=|r_V|^2VV`$ reduces then to only one nonvanishing matrix element. The state dependent noise originates from the fact that interference occurs only with a probability according to the quality parameter $`𝒬^{}`$ and is incoherent otherwise, which finally yields
$`\overline{PG}_{mod}\rho \overline{PG}_{mod}^{}`$ $`=`$ $`𝒬^{}(M_{tt}+M_{rr})\rho (M_{tt}+M_{rr})^{}`$
$`+`$ $`(1𝒬^{})(M_{tt}\rho M_{tt}^{}+M_{rr}\rho M_{rr}^{})`$
From this ansatz we obtain a model QPT-matrix $`\chi _{mod}`$ by minimizing the sum of the absolute squared values of all the matrix elements of $`\chi _{mod}\chi _{exp}`$ numerically (see Fig. 4c). The obtained quality value $`𝒬^{}=0.904`$ is in very good agreement with $`𝒬`$ obtained from the fit to the HOM-dip. This indicates that it is indeed mainly imperfect overlap at the beam splitter which causes the state dependent noise. In order to compare the model with the real setup we calculate the fidelities between the predicted and the experimentally measured output density matrices obtaining an average value of F$`{}_{mod𝒬^{}}{}^{exp}=96.6\%\pm 1.7\%`$. An alternative model including depolarization in the gate did not significantly change the figure, the resulting white noise was negligible.
In conclusion we have presented a C-Phase gate acting on the polarization degree of photons. The gate relies only on one second order interference at a polarization dependent beam splitter and thus significantly simplifies previous approaches. We have demonstrated the entangling quality of the gate for various input states. A linear quantum process tomography allowed us to match a model of the gate to the experimental data. The resulting fit proofs the assumption, that the main deviation from optimal performance is due to distinguishability of incident photons at the overlap beam splitter. By means of further filtering this can be improved on the cost of count rate. The results ensure that this gate is ready to be used in various quantum information processing tasks such as generating multi photon entanglement or for complete Bell state analysis in quantum teleportation and entanglement swapping experiments.
We acknowledge stimulating discussions with A. Zeilinger. This work was supported by the Deutsche Forschungsgemeinschaft and the European Commission through the EU Project RamboQ (IST-2001-38864)
|
warning/0506/cond-mat0506063.html
|
ar5iv
|
text
|
# Survival Probabilities of History-Dependent Random Walks
\[
## Abstract
We analyze the dynamics of random walks with long-term memory (binary chains with long-range correlations) in the presence of an absorbing boundary. An analytically solvable model is presented, in which a dynamical phase-transition occurs when the correlation strength parameter $`\mu `$ reaches a critical value $`\mu _c`$. For strong positive correlations, $`\mu >\mu _c`$, the survival probability is asymptotically finite, whereas for $`\mu <\mu _c`$ it decays as a power-law in time (chain length).
\]
Dynamical systems with long-range spatial or temporal correlations are attracting considerable interest across many disciplines, lending applications to physical, biological, social, and economical sciences (see, e.g., and references therein). Such systems are often analyzed using random walks, a fundamental concept of statistical physics. Random walks lend applications to numerous scientific fields (see, e.g., and references therein). In particular, random walks in the presence of absorbing traps (boundaries) have been studied in recent years as models for various systems such as absorbing-state phase transitions , polymer adsorption , granular segregation , the spreading of an epidemic , and in the context of complex adaptive systems . In this work we analytically study random walks with an absorbing boundary, in which the jump probability is history-dependent, resulting in long-range correlations.
The statistical properties of data such as DNA strings, written texts and financial data (e.g., stock market quotes), are known to significantly deviate from those of purely random sequences . Such systems may be studied by mapping them onto a correlated sequence of symbols. Although the nature of the resulting sequence may depend upon the choice of mapping (see, e.g., ), the essential statistical properties of the original system are often preserved. By choosing a mapping of these systems onto two symbols , the problem is reduced to the exploration of correlated binary chains, which are equivalent to one dimensional random walks with a constant step size. These binary chains have long-range correlations and often exhibit a super-diffusive nature, in which the variance grows asymptotically faster than the string length.
The preceding discussion motivates a study of random walks with a history-dependent jump probability, both with and without an absorbing boundary, as means of facilitating our understanding of systems with long-range correlation. In Ref. we presented a model for a history-dependent random walk. Although simple, this model features a dynamical phase transition between normal diffusion and super-diffusion as the correlation strength parameter reaches a critical value. In this work we analyze random walks with long-range correlations in the presence of an absorbing boundary.
We begin by introducing a simple model that incorporates long-range correlations into an otherwise random sequence. Consider a discrete string of binary symbols, $`a_i\{0,1\}`$, in which the conditional probability of a given symbol (say, zero) occurring at the position $`L+1`$ is history-dependent, and is given by
$$p(k,L)=\frac{1}{2}\left(1\mu \frac{L2k}{L+L_0}\right),$$
(1)
where $`k`$ is the number of such symbols (zeros) appearing in the preceding $`L`$ bits. The correlation parameter $`\mu `$, where $`1<\mu <1`$, determines the strength of correlations in the system. The persistence condition $`\mu >0`$ implies that a given symbol in the preceding sequence promotes the birth of a new identical symbol. In the anti-persistence regime $`\mu <0`$, on the other hand, each symbol inhibits the appearance of a new identical symbol. The parameter $`L_0>1`$ is a constant transient time. For $`LL_0`$ the sequence is approximately random (uncorrelated), whereas for $`LL_0`$ the effect of correlations takes over .
In this model, the conditional probability $`p(k,L;\mu ,L_0)`$ depends on the number of zeros (or unities) in the preceding bits, and is independent of their arrangement. This allows one to obtain an analytical description of the system’s dynamical behavior. As demonstrated in , this two-parameter model provides a good description of the observed statistical properties of various systems such as coarse-grained DNA strings, written texts, and financial data, when mapped onto a binary chain.
The probability $`P(k,L+1)`$ of finding $`k`$ identical symbols (say, zeroes) in a sequence of length $`L+1`$ follows from the evolution equation
$`P(k,L+1)`$ $`=`$ $`p(k1,L)P(k1,L)`$ (3)
$`+[1p(k,L)]P(k,L).`$
Crossing to the continuous limit, one obtains the Fokker-Planck diffusion equation for the correlated process,
$$\frac{P}{L}=\frac{1}{2}\frac{^2P}{x^2}\frac{\mu }{L+L_0}\frac{(xP)}{x},$$
(4)
where $`x2kL`$ is the distance from the origin in the corresponding random walk, and we have neglected high order terms which are irrelevant for $`1L_0L`$. Solutions of Eq. (4) under the initial condition $`P(x,L=0)=\delta (x)`$ were given in .
We introduce an absorbing boundary at $`x=0`$, by imposing the boundary condition $`P(x0,L>0)=0`$. The evolution equation (4) along with this boundary condition and the initial condition $`P(x,L=0)=\delta (x)`$, has a solution in the form
$$P(x,L)\frac{(L+L_0)^\mu }{D(L)^{3/2}}x\mathrm{\Theta }(x)\mathrm{exp}\left[\frac{x^2}{2D(L)}\right],$$
(5)
where $`\mathrm{\Theta }`$ is the Heaviside step function, and $`D(L)`$ is given by
$$D(L;\mu ,L_0)\frac{L+L_0}{12\mu }\left[1\left(\frac{L_0}{L+L_0}\right)^{12\mu }\right].$$
(6)
Equation (6) breaks down in the special case $`\mu =\mu _c`$, where $`\mu _c1/2`$. In this case, $`D(L)`$ is given by
$$D(L;\mu =\mu _c,L_0)(L+L_0)\mathrm{ln}\left(\frac{L+L_0}{L_0}\right).$$
(7)
The first two moments of the distribution function of the survived walkers are given by $`x^2=\pi D(L)/2`$ and $`x^2=2D(L)`$. The variance of the probability distribution $`P(x,L)`$ thus equals
$$V(L;\mu ,L_0)=\frac{4\pi }{2}D(L).$$
(8)
This result implies that for $`\mu <\mu _c`$, the asymptotic variance scales linearly with the string length, whereas for $`\mu >\mu _c`$ it scales as $`VL^{2\mu }`$. Hence, a history-dependent sequence with strong positive correlations ($`\mu >\mu _c`$) is characterized by a super-diffusion phase in which the variance grows asymptotically faster than $`L`$, both without and with an absorbing boundary.
The survival probability $`S(L)_0^{\mathrm{}}P(x,t)𝑑x`$ of the walkers is given, for $`\mu \mu _c`$, by
$$S(L;\mu \mu _c)\left|\left(\frac{L+L_0}{L_0}\right)^{12\mu }1\right|^{1/2},$$
(9)
whereas for $`\mu =\mu _c`$ we find
$$S(L;\mu _c)\mathrm{ln}^{\frac{1}{2}}\left(\frac{L+L_0}{L_0}\right).$$
(10)
The survival probability thus changes its asymptotic ($`LL_0`$) behavior at the phase transition value $`\mu _c=1/2`$, and one identifies three qualitatively different regimes,
$$S(LL_0)\{\begin{array}{cc}L^{\frac{1}{2}+\mu }\hfill & \mu <\mu _c\text{ ;}\hfill \\ \mathrm{ln}^{\frac{1}{2}}(L/L_0)\hfill & \mu =\mu _c\text{ ;}\hfill \\ const.\hfill & \mu >\mu _c\text{ .}\hfill \end{array}$$
(11)
The normalization of the survival probability is sensitive to the discrete initial conditions. Since the chain is nearly random for $`LL_0`$, the normalization may be approximated by equating S(L), for $`1LL_0`$, to the survival probability of a purely random, continuous walk. We thus find
$$S(1LL_0)\text{erf}\left(|x_0|/\sqrt{2L}\right),$$
(12)
where $`x_0`$ is the distance between the absorbing boundary and the origin.
In order to confirm the analytical results, we perform numerical simulations of (discrete) binary sequences. Figure 1 displays the resulting survival probability $`S(L)`$ of correlated strings with various values of the correlation parameter $`\mu `$. We find an excellent agreement between the analytical results \[Eqs. (9), (10) and (12)\] and the numerical ones.
The preceding results, starting from the distribution function in Eq. (5), are valid only when the absorbing boundary is placed at the origin. One would like to generalize the results for boundaries located at arbitrary locations $`x_0`$, under the boundary condition $`P(xx_0,L>0)=0`$. Note that the diffusion equation (4) is not invariant under the translation $`xx+d`$, where $`d`$ is a constant. The generalization of our solution for $`x_00`$ is therefore non-trivial. Nevertheless, we have verified numerically that the asymptotic behavior of the survival probability given by Eq. (11) remains valid for arbitrary values of $`x_0`$, in all three regimes. The similar asymptotic behavior for different choices of $`x_0`$ is demonstrated numerically in Fig. 2.
The above results can be readily generalized for a biased random walk with a moving boundary. For example, for the biased jump probability
$$p(k,L)=\frac{1}{2}\left(1+q\mu \frac{L2k}{L+L_0}\right)$$
(13)
the above results will hold if we apply the transformation $`xxx_c(L)`$, with
$$x_c(L)q\frac{L+L_0}{1\mu }\left[1\left(\frac{L_0}{L+L_0}\right)^{1\mu }\right].$$
(14)
In summary, in this work we have analyzed the dynamics of random walks in which the jump probabilities are history-dependent, in the presence of an absorbing boundary. Using an analytically solvable model, we identify a dynamical phase transition characterizing the system’s global behavior. For small values of the correlation strength ($`\mu <\mu _c`$) the survival probability decays as $`SL^{(\mu _c\mu )}`$, whereas for $`\mu >\mu _c`$ the system is characterized by finite asymptotic survival probabilities.
ACKNOWLEDGMENTS
The research of SH was supported by G.I.F. Foundation.
|
warning/0506/astro-ph0506119.html
|
ar5iv
|
text
|
# Investigating ionized disc models of the variable narrow-line Seyfert 1 PG 1404+226
## 1 Introduction
Narrow Line Seyfert 1 (NLS1) galaxies are a class of Active Galactic Nuclei (AGN) first studied by Osterbrock & Pogge (1985). The class is defined by its optical properties: the Balmer lines such as H$`\beta `$ have a FWHM $`<`$ 2000 km/s, the ratio of the \[O III\] $`\lambda `$5007 Å forbidden line to the H$`\beta `$ lines is $`<`$ 3, and emission lines from Fe II or higher ionization states are often present. These properties are discussed in, e.g. Boroson & Green (1992). In the X-ray band Seyfert galaxies are usually well described by a power-law and a ‘soft excess’, more emission below $``$1 keV than expected from an extrapolation of the power-law spectrum observed at higher energies. Modelling the soft excess as a black-body usually results in a good fit. The soft excess is notable for having similar temperatures over a wide range of objects, and is the subject of some debate: if it is thermal in origin it may be due to a slim disc (e.g Abramowicz et al. 1988), Gierliński & Done (2004) show that soft excess temperature in a sample of radio-quiet PG quasars is 0.1 – 0.2 keV over a large range of luminosity, and put forward the idea that the soft excess is an illusion caused by relativistically blurred strong absorption. This paper is based on the idea that it is due to photoionized emission blurred by relativistic motion in an accretion disc.
Seyferts are highly variable, in the optical on time-scales of weeks to years, and in the X-ray on time-scales of hours or even less. NLS1s are more variable in the X-ray band than typical Seyfert 1s, implying they are more compact (Leighly 1999a). X-rays are generated only in a small energetic area of an AGN, i.e. near the central engine. The detection of iron lines broadened by relativistic effects in some sources shows that the central energy source is matter accreting onto a black hole (e.g. Tanaka 1995; Vaughan & Fabian 2004).
This concept is extended in relativistically blurred photoionized disc models, the most recent of which (Ross & Fabian 2005) is used in this paper. This model is intended to simulate radiation from the inner regions of a black hole accretion disc, and it includes atomic physics and general relativistic effects. The underlying spectrum is due to a slab of optically thick gas of constant density illuminated by a power-law, which produces fluorescence lines with a continuum of reflected photons. A power-law with the spectral index of the illuminating photons is added to include light from the illuminating object (e.g. from the base of a jet or a hot corona). The power-law component dominates the observed spectrum in many NLS1s but in PG 1404+226 the soft excess is the most powerful component (Fig. 1). This spectrum is then relativistically blurred as calculated for an accretion disc around a Kerr black hole. The Laor (1991) line profile (standard in xspec) is used as a convolution kernel to do this.
This model is more physically motivated than the ad hoc models usually adopted (i.e. a power-law and a black-body to model the soft excess). We show that this new model is a better fit to the source than previous models and adequately accounts for variability.
## 2 The source
PG 1404+226 is a Narrow Line Seyfert 1 galaxy, with a redshift of 0.098. Wang & Lu (2001), using the relation with \[O III\] width, find the central black hole mass to be $`1.0\times 10^7`$ M, with an error of $``$0.5 dex. During the observation we observe a flux of $`1.1\times 10^{12}`$ ergs cm<sup>-2</sup> s<sup>-1</sup> in the 0.3 – 12.0 keV range, which we convert (using the inverse square law with a Hubble parameter of 70 km s<sup>-1</sup> Mpc<sup>-1</sup>) to a luminosity of $`2.3\times 10^{43}`$ erg s<sup>-1</sup>.
The source has previously been observed with the ASCA satellite (e.g. Vaughan et al. 1999), and is included in the two-part review by Leighly (1999a,b). The ASCA observation showed PG 1404+226 to be strongly variable, with a clear soft excess. The source showed possible apparent absorption around 1 keV, which is discussed in detail in e.g. Leighly et al. (1997). PG 1404+226 did not show an iron line. As this article will show, the 1 keV feature is a natural consequence of an ionized disc.
## 3 Observations and data reduction
PG 1404+226 was observed on-axis with XMM-Newton during revolution 0279 (2001-06-17 – 2001-06-18). All EPIC cameras were operated in PrimeFullWindow mode with the filter Thin1. The RGS instruments were used in the HighEventRateWithSES mode. The OM did not take data. During the observation the EPIC MOS took 20.8 ks of data, the EPIC pn 18.3 ks and the RGS instruments 21.4 ks.
The Observation Data Files (ODFs) were reduced in the standard way using sas 5.4.1 to produce event lists. Light curves and spectra were created for the sources by using a circular, source-centred extraction region of 20 arcsec radius for the pn and 30 arcsec for the MOS. The smaller than standard extraction radii were chosen to exclude as much as possible of the strong background. The background light curves and spectra were extracted using regions away from any sources and chip-gaps. The presence of strong background flaring in the observation was noted and compensated for using a Good Time Interval (GTI) file, leaving 16.9 ks for the pn and 20.1 ks for the MOS.
The pn images were also examined for Out-Of-Time (OOT) events and pile-up. Any OOT events in PG 1404+226 are unobservable due to a low signal to noise ratio, so no OOT correction was used. Use of the sas task epatplot showed no pile-up.
The MOS cameras are not accurately calibrated against the pn (Kirsch et al. 2004), and the pn has a larger effective area and spectral range. Therefore all fits were done on pn data alone, and checked against the MOS for consistency. The MOS gives a flatter spectral index than the pn, but no other serious differences were noted. The RGS data were reduced and found to be consistent with the other instruments, however, the statistics are dominated by the pn and MOS so RGS data were not included in any of the fits. The limits in accurate calibration were taken as 0.3 keV and 12.0 keV for the pn cameras, and the MOS was used over the 0.3 keV – 10.0 keV range. The data reduction yielded 7807 (pn) source photons (including an estimated 232 background photons).
The spectral analysis was performed using spectra grouped so that each bin contained at least 20 source counts, to ensure that $`\chi ^2`$ statistics would give reasonable results. Response files generated by sas were used, and the spectra were examined using xspec v11.3 (Arnaud 1996). A Galactic hydrogen absorption column of $`2.14\times 10^{20}`$ cm<sup>-2</sup> was taken from the nh ftool (Dickey & Lockman 1990). Fits were also performed allowing for free hydrogen columns at the source redshift, no evidence for any extra absorption was found. All quoted errors are 90% limits on one parameter.
## 4 Analysis
The data were reduced and spectra were created as described in Section 3. The spectrum is shown in Fig. 2. We fit the data with the standard black-body and power-law. Adding a red-shifted edge fixed at a source-frame energy of 0.87 keV (the O VIII K absorption edge) improves the fit ($`\mathrm{\Delta }\chi ^2`$ of 18 for 1 degree of freedom). Adding a second edge at 0.74 keV (O VII) does not improve the fit. Allowing the energy of the first edge to vary improves $`\chi ^2`$ by 4 for 1 degree of freedom. The results are given in Table 1. To model cold iron reflection from distant gas (e.g. a torus) we added a narrow red-shifted Gaussian at 6.4 keV (in the source frame). This improves the fit further, to a $`\chi _\nu ^2`$ of 0.916 for 190 degrees of freedom (with all other parameters consistent with those reported for the previous fit); allowing the width or energy of the line to vary does not improve the fit. An F-test gives an 89% probability that the addition of the line is valid. Including MOS data does not increase the significance, an F-test gives an 84% probability that the line is valid when all three EPIC observations are fit simultaneously. The use of the F-test on hypotheses on the boundary of the parameter space, such as testing for the existence of an emission line, is of questionable statistical value (Protassov et al. 2002). We therefore calibrate the F-test using simulated observations (from the xspec ‘fakeit’ command), by generating data from a model without an iron line. In this data we measure an iron line at the same or higher significance (measured by an F-test) as detected in our real data 12% of the time, therefore we conclude that the F-test is reasonably calibrated in this case and that the line is 88% significant. The statistics are poor at these energies, so the equivalent width of the line is barely constrained; xspec reports the equivalent width of the line as $`0.5_{0.5}^{+1.6}`$ keV. Given that the line is just on the boundary of 90% significance, we choose not to include the line in any of our further fits. We therefore consider the free edge energy model to be our best simple fit, with a $`\chi _\nu ^2`$ of 0.934 for 191 degrees of freedom.
We also fit PG 1404+226 using the ionized disc model. The free parameters of the model are the iron abundance (Fe, in units of solar metallicity), ionization parameter ($`\xi `$, the number of incident ionizing photons cm<sup>-2</sup>s<sup>-1</sup>), and the spectral index of the illuminating continuum ($`\mathrm{\Gamma }`$, this is also the index of the power-law component we add). The spectrum is blurred with a Laor line profile, which is standard in xspec, and has as free parameters the emissivity index of the disc (Index), the inclination of the disc to the line of sight ($`\theta `$, in degrees), the inner radius of the disc ($`R_{in}`$, in units of gravitational radii) and outer radius (which we fix at 100 gravitational radii as it is not strongly constrained). This is intended to reproduce the spectrum from a photoionized disc around a maximally-rotating (Kerr) black hole. Fits made using relativistic blurring from a non-rotating (Schwarzschild) black hole were not found to be satisfactory; there were strong residuals below 1 keV where the blurring is insufficient to smooth out the emission and absorption features. A more rapidly rotating black hole allows for stronger blurring.
The parameters obtained in fitting PG 1404+226 with the ionized disc model are given in Table 2, and the model is plotted showing the two components in Fig. 1. We find no evidence for absorption edges with this model. Our best ionized disc fit has a $`\chi _\nu ^2`$ of 0.892 for 189 degrees of freedom.
### 4.1 Variability
PG 1404+226 varied by a factor of $``$3 over the course of the observation, as shown in Fig. 3. The Figure shows soft (0.3 – 1.0 keV) and hard (1.0 – 8.0 keV) count rates. Counts above 8 keV are not included as the source is extremely weak at those energies. The hard count rate does not strongly vary whereas the soft count rate increases during the observation; this indicates that the spectral shape changes. To investigate this change we divide the observation into four different periods suggested by the shape of the light-curve, as indicated by the vertical lines on the Figure. We fit these using our best simple model, testing to see if adding the absorption edge (at a fixed energy of 0.9 keV in the source frame) improves the fit. The results are given in Table 3. Adding an edge results in an improvement to the fit in all time periods but the first one. If an edge is included in the first period fit xspec reports an optical depth of $`0.05_{0.05}^{+0.41}`$, with a negligible change in $`\chi ^2`$ for the loss of a degree of freedom. To further investigate this we perform simulated observations to evaluate whether we would detect an absorption edge with a $`\tau `$ of 0.9. We measure an optical depth $`>`$ 0.05 in $``$99.9% of our simulated observations, therefore we are confident that the absence of an absorption edge in the first time period is not due to insufficient data quality.
We also fit the same time periods using the ionized disc model, with the results in Table 2. The ionized disc model and simple model fit the data about equally well when analysing the shorter time periods individually.
It has been suggested that the variability in some sources is due to changes in the normalization of the underlying continuum, with a fairly constant reflection component from an ionized disc (Fabian et al. 2002, Vaughan & Fabian 2004). To investigate this we simultaneously fit all four time periods, allowing each to have independent reflection and power-law component normalizations and tying all the other parameters together. We investigate possible changes in spectral index of the underlying continuum by repeating the previous fit with a free spectral index ($`\mathrm{\Gamma }`$). The results from both of these fits are reported in Table 2. We further attempt to fit the data with no power-law component at all. This did not improve the fits so we do not report the results.
We similarly fit the four periods simultaneously with the simple model, tying the black-body temperature, power-law index and edge energy together but allowing the normalizations and optical depths to vary independently. The results are given in Table 3. Our best fit to the data is with the ionized disc model, simultaneously fitting all parameters except the normalizations and $`\mathrm{\Gamma }`$; this has a $`\chi _\nu ^2`$ of 0.918 for 335 degrees of freedom.
## 5 Discussion
### 5.1 Simple model
The simple model is an acceptable fit to the data. This model includes absorption at $``$0.9 keV, which is likely to be due to the O VIII K absorption edge at 0.87 keV. The $`>0.87`$ keV energy might imply the absorption originates in a slightly relativistic outflow, although the data do not constrain this (see Leighly et al. 1997).
We learn more when we consider spectral variability over the observation. The absorption appears to ‘switch on’ some way into the observation, as shown in Table 3. The absorption increases at the same time as the luminosity of the source increases (at $``$7.5 ks, see Fig. 3), and we have eliminated the possibility that this is a statistical effect. If the absorption is due to O VIII, the lack of strong O VII absorption implies the absorbing gas has a high ionization parameter; this means it is close to the central black hole and the rapid change implied by the variability of the feature is plausible. This variation in absorption has been previously investigated by Dasgupta et al. (2004), who understand it within the model of absorption due to a disc wind - as the source brightens more matter is driven off in the wind and the absorption intensifies. We have shown that the ionized disc model naturally explains this apparent increase in absorption.
### 5.2 Ionized disc model
The ionized disc model is the best fit to the entire observation. This model has the benefit of being self-consistent and not ad hoc. It does not require any additional absorption features to give a good fit, the model itself reproduces the dip in the spectrum. When the spectral variability of the source is investigated using the ionized disc model, it accounts for it naturally; the spectral index of the illuminating powerlaw varies over the course of the observation. This fit (at the bottom of Table 2) has a $`\chi ^2`$ 20 lower than the equivalent simple fit (Table 3) for two fewer degrees of freedom. We conclude that the ionized disc model is the most satisfactory explanation of the data.
The fit parameters give us information about the central source. Firstly, it is evident from Table 2 that the inclination of the source is quite high; around $`60^{}`$. This does not fit very well with the unified model of AGN (Antonucci & Miller 1985, Antonucci 1993), in which Seyfert 1 galaxies are those visible at low inclinations, with no molecular torus in the line of sight. However, PG 1404+226 is fairly luminous, and there have been suggestions that brighter quasars tend to show less absorption from a torus (Lawrence 1991, Brandt & Hasinger 2005). The high inclination also implies that the brightness of the source is partly due to relativistic Doppler beaming from the rotation of the disc. We investigated the effect of inclination angle on flux within our model, and find that the brightness is decreased by a factor of $``$5 if the inclination is changed to 0 (all other parameters unchanged), and increased by a factor of $``$3 at 90. The measured luminosity is therefore a reasonable order-of-magnitude estimate to the intrinsic luminosity of the source.
Information on the rotation can also be obtained. Note that fits were performed assuming a Kerr black hole, and that Schwarzschild fits were unacceptable (Section 4). The measured inner disc radius provides further evidence that the black hole is rotating - it is at the last stable orbit radius for a maximally rotating black hole; well inside the plunging region for a non-rotating black hole. The black hole rotation is not a free parameter in this analysis so we must be careful in drawing conclusions, however it seems that a rapidly rotating black hole is likely. It is also a possibility that radiation from the plunging region is important (see Krolik & Hawley 2002).
The power-law component of this source is weak compared to the reflection component, i.e. the reflection fraction is very high. This may be due to light bending effects on the primary continuum. If the continuum source is compact and close to the black hole almost all the radiation it emits is bent onto the disc rather than escaping to the observer; this reduces the observed power-law flux while enhancing the disc illumination, and therefore the reflection component (Miniutti et al. 2003; Miniutti & Fabian 2004). To produce a reflection dominated spectrum in this model the primary source in PG 1404+226 must be within 3 – 4 gravitational radii of the black hole. Reflected radiation escapes the black hole more easily than the illuminating continuum as it is emitted by rapidly moving matter and tends to be beamed along the plane of the disc to the observer. The high inclination therefore also contributes to the high reflection fraction.
This model raises the question of how energy is transported through the system. The original source of energy is the matter falling into the gravitational potential of the central black hole. This heats the electrons that provide the illuminating continuum, and the energy reflects off the disc and escapes. The question of how the in-falling matter heats the source of the illuminating continuum is an unsolved one, magnetic reconnection is a possibility, as is emission from the base of a weak jet, or shocks in a failed jet (Ghisellini, Haardt & Matt 2004). In the case of a jet the spin of the black hole may also provide energy.
Given the small disc inner radius we calculate for the source, most of the X-ray emission is coming from within a few gravitational radii of the black hole, so strong gravity effects are expected. This small size also explains the rapid variability. If the illuminating continuum is, e.g. UV radiation being Compton up-scattered by hot electrons, the change in spectral index could be due to a change in temperature of the electrons, and the extra flux could be due to an increase in emitting area.
## 6 Conclusions
We have shown that:
* Physically plausible models based on relativistically blurred reflection from photoionized discs fit observations to PG 1404+226 better than the typical ad hoc simple model of a black-body and power-law. In particular, a ‘soft excess’ at low X-ray energies is a natural consequence of a blurred ionized disc; the ‘power-law’ is largely made up of broadened iron line emission and the ‘soft excess’ of other blurred lines, plus the Compton reflection component. The ionized disc model also leads naturally to apparent absorption features at low X-ray energies.
* The X-ray emitting accretion disc in PG 1404+226 has an inclination of $`60^{}`$.
* The central black hole is most likely strongly rotating, alternatively, radiation from the plunging region may be important.
* The source has strong and complex spectral variability, which is well described within the ionized disc model by a variation in the spectral index of the illuminating continuum, as well as variation in the flux of the two components.
* If the alternative simple model is adopted, the spectral variation is reasonably fit with a rapid increase in absorption at $``$0.9 keV coincident with the increase in luminosity of the source, plus some variation in the other model parameters.
The new relativistically blurred photoionized disc model of Ross & Fabian (2005) is clearly a valuable tool in investigating NLS1s, and in sources where an ionized disc fit is appropriate information about the central conditions can be obtained. The application of these models to more, brighter sources (currently in progress) should lead to advances in our knowledge of AGN.
## Acknowledgments
J.C. is a PPARC funded PhD student. A.C.F. thanks the Royal Society for support. W.N.B. acknowledges support from NASA grant NAG5-9924 and NASA LTSA grant NAG5-13035. The XMM-Newton satellite is an ESA science mission (with instruments and contributions from NASA and ESA member states). This work made use of the NASA/IPAC Extragalactic Database. J.C. would like to thank Luigi Gallo and Giovanni Miniutti for many helpful comments, and Randy Ross, Simon Vaughan, and Roderick Johnstone for contributing xspec models and code.
## References
Abramowicz M. A., Czerny B., Lasota J. P., Szuszkiewicz E., 1988, ApJ, 332, 646
Antonucci R., 1993, ARA&A, 31, 473
Antonucci R.R.J., Miller J.S., 1985, ApJ, 297, 621
Arnaud K.A., 1996, ADASS 5, 17A
Boroson T.A., Green R.F., 1992, ApJS, 80, 109
Brandt W.N., Hasinger G., 2005, astro-ph/0501058
Dasgupta S., Rao, A.R., Dewangan G.C., Agrawal V.K., 2005, ApJ, 618, L87
Dickey J.M., Lockman F.J., 1990, ARA&A, 28, 215
Fabian A.C., et al., 2002, MNRAS, 335, L1
Ghisellini G., Haardt F., Matt G., 2004, A&A, 413, 535
Gierliński M., Done C., 2004, MNRAS, 349, L7
Kirsch M.G.F., et al, 2004, Proc. SPIE, 5488, 103
Krolik J.H., Hawley J.F., 2002, ApJ, 573, 754
Laor A., 1991, ApJ, 376, 90
Lawrence A., 1991, MNRAS, 252, 586
Leighly K., 1999a, ApJSS, 125, 297
Leighly K., 1999b, ApJSS, 125, 317
Leighly K., Mushotzky R., Nandra K., Forster K., 1997, ApJ, 489, L25
Miniutti G., Fabian A.C., 2004, MNRAS, 349, 1435
Miniutti G., Fabian A.C., Goyder R., Lasenby A.N., 2003, MNRAS, 344, 22
Osterbrock D.E., Pogge R.W., 1985, ApJ, 297, 166
Protassov R., van Dyk D.A., Connors A., Kashyap V.L., Siemiginowska A., 2002, 571, 545
Ross R.R., Fabian A.C., 2005, MNRAS, 358, 211
Tanaka Y., et al., 1995, Nature, 375, 659
Vaughan S., Fabian A.C., 2004, MNRAS, 348, 1415
Vaughan S., Reeves J., Warwick R., Edelson R., 1999, 309, 113
Wang J.-M., Netzer H., 2003, A&A, 398, 927
Wang T., Lu Y., 2001, A&A, 377, 52
|
warning/0506/cond-mat0506005.html
|
ar5iv
|
text
|
# All-electronic coherent population trapping in quantum dots
<sup>1</sup><sup>1</sup>institutetext: Instituut-Lorentz, Universiteit Leiden, P.O. Box 9506, 2300 RA Leiden, The Netherlands
Department of Physics, University of California San Diego, La Jolla, California 92093–0319, USA
Decoherence; open systems; quantum statistical methods Quantum dots Coulomb blockade; single-electron tunneling
## Abstract
We present a fully electronic analogue of coherent population trapping in quantum optics, based on destructive interference of single-electron tunneling between three quantum dots. A large bias voltage plays the role of the laser illumination. The trapped state is a coherent superposition of the electronic charge in two of these quantum dots, so it is destabilized as a result of decoherence by coupling to external charges. The resulting current $`I`$ through the device depends on the ratio of the decoherence rate $`\mathrm{\Gamma }_\varphi `$ and the tunneling rates. For $`\mathrm{\Gamma }_\varphi 0`$ one has simply $`I=e\mathrm{\Gamma }_\varphi `$. With increasing $`\mathrm{\Gamma }_\varphi `$ the current peaks at the inverse trapping time. The direct relation between $`I`$ and $`\mathrm{\Gamma }_\varphi `$ can serve as a means of measuring the coherence time of a charge qubit in a transport experiment.
Coherent population trapping is a quantum optical phenomenon in which the laser illumination of an atom drives an atomic electron into a coherent superposition of orbital states and traps it there . Such superpositions can be “dark”, in that they are further decoupled from the optical fields. Brandes and Renzoni have shown how such states can also be formed in artificial atoms (quantum dots) through the use of laser illumination . In this paper we present an all-electronic analogue, i.e. without laser illumination, of coherent population trapping in quantum dots. (For an analogy in superconducting Josephson junctions, see Ref. ; for an analogy in single benzene molecules, see Ref. .) We illustrate this effect by considering a system of three tunnel-coupled quantum dots and show that, under proper bias and resonance conditions, an electron can become trapped in a coherent superposition of states in different dots. This state is “dark” in the sense that, due to the Coulomb blockade, no further electrons can pass through the dots and current flow is blocked in the absence of decoherence.
The trapping effect provides a novel mechanism for current rectification, since the blocking is effective for one sign of the bias voltage only. This quantum mechanical mechanism is distinct from mechanisms discussed previously. In particular, the classical rectification mechanism of Stopa and collaborators traps the electron in a single quantum dot, rather than in a coherent superposition of spatially separated states. Experiments by Ono and collaborators on rectification in double quantum dots likewise trap an electron in a single dot. The three-dot configuration requires no Aharonov-Bohm phase to trap an electron, in contrast to the two-dot configuration of Marquardt and Bruder . Because of the phase coherent origin of the effect discussed here, the current that leaks through the device when it is blocked provides a method by which one can determine the coherence time of a charge qubit.
The three-dot trap is shown schematically in Fig. 1. Three quantum dots and three electron reservoirs are connected by reversible or by irreversible transitions. Reversible transitions between the quantum dots are described by the tunnel Hamiltonian
$$H=T|CA|+T|CB|+\mathrm{H}.\mathrm{c}.=2^{1/2}T|C\mathrm{\Phi }_+|+\mathrm{H}.\mathrm{c}.$$
(1)
We have defined the states
$$|\mathrm{\Phi }_\pm =2^{1/2}(|A\pm |B).$$
(2)
We consider the case that the energies of the single-particle levels $`|A`$, $`|B`$, $`|C`$ in the three dots are all the same (set at zero), so that inelastic transitions between these levels do not play a role. To minimize the number of free parameters, all tunnel rates are put equal to $`T`$. (The more general case of unequal tunnel rates will be considered at the end of the paper.) We assume time-reversal symmetry, hence $`T`$ is a real number. (Since results do not depend on the sign of $`T`$, we will take $`T`$ positive for ease of notation.) We furthermore assume that the electrostatic charging energy of the combined three-dot system is sufficiently large that the total number of electrons does not exceed one. Many-electron states are projected out and hence we may ignore spin.
For a bias voltage $`|V|T/e`$, and at zero temperature, the transitions from the source reservoirs into dots $`A`$ and $`B`$ and from dot $`C`$ into the drain reservoir are irreversible. (Because of this restriction, the rectification provided by our device does not apply to the range $`|V|T/e`$ around zero bias.) The tunnel rates between dots and reservoirs are all set equal to $`\mathrm{\Gamma }`$. The quantum jump operators are
$$L_A=\sqrt{\mathrm{\Gamma }}|A0|,L_B=\sqrt{\mathrm{\Gamma }}|B0|,L_C=\sqrt{\mathrm{\Gamma }}|0C|,$$
(3)
where $`|0`$ is the state with all three dots empty.
We study the dynamics of this device by means of the master equation approach to single-electron tunneling , which describes not only the populations of the dot levels, but also accounts for quantum coherences between them. The master equation gives the time evolution of the three-dot density matrix $`\rho (t)`$ in the Lindblad form
$$\frac{d\rho }{dt}=i[H,\rho ]+\underset{X=A,B,C}{}\left(L_X^{}\rho L_X^{}\frac{1}{2}L_X^{}L_X^{}\rho \frac{1}{2}\rho L_X^{}L_X^{}\right).$$
(4)
(We have set $`\mathrm{}1`$.) As initial condition we take $`\rho (0)=|00|`$.
We use as a basis for the density matrix the four states
$$|e_1=|\mathrm{\Phi }_+,|e_2=|\mathrm{\Phi }_{},|e_3=|C,|e_4=|0.$$
(5)
This four-dimensional space may be reduced to a three-dimensional subspace by noting that the master equation (4) couples only $`\rho _{44}`$ and $`\rho _{ij}`$ with $`i,j3`$. The matrix elements $`\rho _{ij}`$ with $`i=4,j4`$ or $`j=4,i4`$ remain zero. We may therefore seek a solution of the form
$$\rho (t)=\stackrel{~}{\rho }(t)+\left[1\mathrm{Tr}\stackrel{~}{\rho }(t)\right]|00|,$$
(6)
where $`\stackrel{~}{\rho }`$ is restricted to the three-dimensional subspace spanned by the states $`|e_i`$ with $`i3`$. The evolution equation for $`\stackrel{~}{\rho }`$ reads
$`{\displaystyle \frac{d\stackrel{~}{\rho }}{dt}}=M\stackrel{~}{\rho }+\stackrel{~}{\rho }M^{}+Q,\stackrel{~}{\rho }(0)=0,`$ (7)
$`M=\left(\begin{array}{ccc}0& 0& 2^{1/2}iT\\ 0& 0& 0\\ 2^{1/2}iT& 0& \mathrm{\Gamma }/2\end{array}\right),`$ (8)
$`Q=\mathrm{\Gamma }(1\mathrm{Tr}\stackrel{~}{\rho })\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right).`$ (9)
All off-diagonal elements of $`\stackrel{~}{\rho }`$ vanish, except the purely imaginary $`\stackrel{~}{\rho }_{13}=\stackrel{~}{\rho }_{31}`$. Four real independent variables remain, which we collect in a vector $`v=(\mathrm{Tr}\stackrel{~}{\rho },\stackrel{~}{\rho }_{11},\stackrel{~}{\rho }_{33},\mathrm{Im}\stackrel{~}{\rho }_{13})`$ satisfying
$`{\displaystyle \frac{dv}{dt}}=X(vv_{\mathrm{}}),v(0)=0,`$ (10)
$`X=\left(\begin{array}{cccc}2\mathrm{\Gamma }& 0& \mathrm{\Gamma }& 0\\ \mathrm{\Gamma }& 0& 0& 2^{3/2}T\\ 0& 0& \mathrm{\Gamma }& 2^{3/2}T\\ 0& 2^{1/2}T& 2^{1/2}T& \mathrm{\Gamma }/2\end{array}\right),v_{\mathrm{}}=\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right).`$
The solution is
$$v(t)=v_{\mathrm{}}e^{Xt}v_{\mathrm{}}.$$
(12)
All four eigenvalues $`\lambda _n`$ of $`X`$ have a negative real part, so $`v(t)v_{\mathrm{}}`$ for $`t\mathrm{}`$ and hence
$$\underset{t\mathrm{}}{lim}\rho (t)=|\mathrm{\Phi }_{}\mathrm{\Phi }_{}|.$$
(13)
This is the trapped state: it does not decay because it is an eigenstate of $`H`$. For large times $`|v(t)v_{\mathrm{}}|e^{\alpha t}`$, with trapping rate $`\alpha =\mathrm{min}(|\mathrm{Re}\lambda _1|,|\mathrm{Re}\lambda _2|,|\mathrm{Re}\lambda _3|,|\mathrm{Re}\lambda _4|)`$. The full expression for $`\alpha `$ is lengthy, but the two asymptotic limits have a compact form,
$$\alpha =\{\begin{array}{cc}4T^2/\mathrm{\Gamma }& \mathrm{if}T\mathrm{\Gamma },\\ \frac{1}{4}(5\sqrt{17})\mathrm{\Gamma }0.22\mathrm{\Gamma }& \mathrm{if}\mathrm{\Gamma }T.\end{array}$$
(14)
If the coupling of the quantum dots to the reservoirs is weaker than between themselves, then the trapping time is of order $`1/\mathrm{\Gamma }`$. One might have guessed the trapping time to be of order $`1/T`$ in the opposite regime $`T\mathrm{\Gamma }`$, but this guess underestimates the correct answer, which is larger by a factor $`\mathrm{\Gamma }/T`$. The fact that $`\alpha 0`$ when $`\mathrm{\Gamma }\mathrm{}`$ can be understood as a decoherence of the inter-dot dynamics induced by a strong coupling to the electron reservoirs.
The full dependence of $`\alpha `$ on $`\mathrm{\Gamma }`$ and $`T`$ is shown in Fig. 2. If $`\mathrm{\Gamma }`$ is increased at constant $`T`$, the trapping rate has a maximum of $`\alpha _{\mathrm{max}}=0.58T`$ at $`\mathrm{\Gamma }=4.35T`$.
The trapping effect requires the coherent superposition of spatially separated electronic states in quantum dots $`A`$ and $`B`$. Such a charge qubit is sensitive to decoherence by coupling to other charges in the environment, which effectively project the qubit on one of the three localized states $`|A`$, $`|B`$, $`|C`$. We model this decoherence by including into the master equation the three quantum jump operators
$$L_{\varphi _X}=\mathrm{\Gamma }_\varphi ^{1/2}|XX|,X=A,B,C.$$
(15)
The decoherence rate $`\mathrm{\Gamma }_\varphi `$ parameterizes the strength of the charge noise and is taken to be dot independent. For a microscopic foundation of the charge noise model we refer to Ref. . We also note that charge noise causes phase as well as energy relaxation.<sup>1</sup><sup>1</sup>1 To calculate the energy relaxation, we decouple the three quantum dots from the electron reservoirs (setting $`\mathrm{\Gamma }0`$) and calculate $`dE/dt=(d/dt)\mathrm{Tr}\rho H`$ from Eq. (16). One finds $`dE/dt=\mathrm{\Gamma }_\varphi E`$, so the energy of the three-dot system relaxes to zero with rate $`\mathrm{\Gamma }_\varphi `$.
The master equation reads
$$\frac{d\rho }{dt}=i[H,\rho ]+\underset{X=A,B,C,\varphi _A,\varphi _B,\varphi _C}{}\left(L_X^{}\rho L_X^{}\frac{1}{2}L_X^{}L_X^{}\rho \frac{1}{2}\rho L_X^{}L_X^{}\right).$$
(16)
The steady-state current,
$$I=\underset{t\mathrm{}}{lim}e\mathrm{\Gamma }C|\rho (t)|C,$$
(17)
is obtained by solving Eq. (16) with the left-hand-side set to zero. We find
$`I`$ $`=`$ $`{\displaystyle \frac{4e\mathrm{\Gamma }T^2}{\mathrm{\Gamma }^2+14T^2+2\mathrm{\Gamma }\mathrm{\Gamma }_\varphi (1+2T^2/\mathrm{\Gamma }_\varphi ^2)}}`$ (20)
$``$ $`\{\begin{array}{cc}e\mathrm{\Gamma }_\varphi \hfill & \mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{\Gamma },T,\hfill \\ 2eT^2/\mathrm{\Gamma }_\varphi \hfill & \mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{\Gamma },T.\hfill \end{array}`$
As illustrated in Fig. 3, the current vanishes both in the limit $`\mathrm{\Gamma }_\varphi 0`$, because of the trapping effect, and in the limit $`\mathrm{\Gamma }_\varphi \mathrm{}`$, because of the quantum Zeno effect . The maximal current is reached at $`\mathrm{\Gamma }_\varphi =2^{1/2}T`$ and is equal to
$`I_{\mathrm{max}}`$ $`=`$ $`{\displaystyle \frac{4e\mathrm{\Gamma }T^2}{\mathrm{\Gamma }^2+14T^2+4\sqrt{2}\mathrm{\Gamma }T}}`$ (23)
$``$ $`\{\begin{array}{cc}4eT^2/\mathrm{\Gamma }\hfill & \mathrm{if}T\mathrm{\Gamma },\hfill \\ \frac{2}{7}e\mathrm{\Gamma }0.29e\mathrm{\Gamma }\hfill & \mathrm{if}\mathrm{\Gamma }T.\hfill \end{array}`$
Comparison of Eqs. (14) and (23) shows that the maximal current $`I_{\mathrm{max}}`$ in the presence of decoherence is set by the trapping rate $`\alpha `$ in the absence of decoherence. For $`T\mathrm{\Gamma }`$ one has exactly $`I_{\mathrm{max}}=e\alpha `$, while for $`\mathrm{\Gamma }T`$ the two quantities differ by a numerical coefficient of order unity. In Fig. 2 both $`I_{\mathrm{max}}/e`$ and $`\alpha `$ are plotted together, and are seen to differ by less than a factor of two over the whole $`\mathrm{\Gamma },T`$ range.
The trapping effect does not happen if the bias is inverted, so that the drain reservoir becomes the source and vice versa. In that case we find for the steady-state current the expression
$`I`$ $`=`$ $`{\displaystyle \frac{4e\mathrm{\Gamma }T^2(2\mathrm{\Gamma }+\mathrm{\Gamma }_\varphi )}{\mathrm{\Gamma }(\mathrm{\Gamma }+\mathrm{\Gamma }_\varphi )(\mathrm{\Gamma }+2\mathrm{\Gamma }_\varphi )+4T^2(6\mathrm{\Gamma }+5\mathrm{\Gamma }_\varphi )}}`$ (26)
$``$ $`\{\begin{array}{cc}8e\mathrm{\Gamma }T^2(\mathrm{\Gamma }^2+24T^2)^1\hfill & \mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{\Gamma },T,\hfill \\ 2eT^2/\mathrm{\Gamma }_\varphi \hfill & \mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{\Gamma },T.\hfill \end{array}`$
For strong decoherence the current is the same in both bias directions, but for weak decoherence the current in the case of inverted bias does not drop to zero but saturates at a finite value. The two cases are compared in Fig. 3. We see that the appearance of a maximum current as a function of $`\mathrm{\Gamma }_\varphi `$ is characteristic for the trapping effect.
We have for simplicity assumed that all three dots have the same tunnel rates and decoherence rates, but this assumption may be easily relaxed. Let us consider first the case that the three-dot structure still has a reflection symmetry, so that dots $`A`$ and $`B`$ are equivalent, but that dot $`C`$ has a different tunnel rate $`\mathrm{\Gamma }^{}`$ into the reservoir and a different decoherence rate $`\mathrm{\Gamma }_\varphi ^{}`$. We denote $`\overline{\mathrm{\Gamma }}_\varphi =(\mathrm{\Gamma }_\varphi +\mathrm{\Gamma }_\varphi ^{})/2`$. The result (20) for the steady-state current generalizes to
$`I`$ $`=`$ $`{\displaystyle \frac{4e\mathrm{\Gamma }^{}T^2}{\mathrm{\Gamma }^2+2T^2(6+\mathrm{\Gamma }^{}/\mathrm{\Gamma })+2\mathrm{\Gamma }^{}\mathrm{\Gamma }_\varphi (\overline{\mathrm{\Gamma }}_\varphi /\mathrm{\Gamma }_\varphi +2T^2/\mathrm{\Gamma }_\varphi ^2)}}`$ (29)
$``$ $`\{\begin{array}{cc}e\mathrm{\Gamma }_\varphi \hfill & \mathrm{if}\mathrm{\Gamma }_\varphi 0,\hfill \\ 2eT^2/\overline{\mathrm{\Gamma }}_\varphi \hfill & \mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{}.\hfill \end{array}`$
The steady-state current still contains the desired information on the rates of decoherence, with the regimes of weak and strong decoherence governed by $`\mathrm{\Gamma }_\varphi `$ and $`\overline{\mathrm{\Gamma }}_\varphi `$, respectively.
In the most general case of arbitrarily different tunnel rates $`T_A,T_B,\mathrm{\Gamma }_A,\mathrm{\Gamma }_B,\mathrm{\Gamma }_C`$ and decoherence rates $`\mathrm{\Gamma }_{\varphi _A},\mathrm{\Gamma }_{\varphi _B},\mathrm{\Gamma }_{\varphi _C}`$, the steady state current in the limit of weak and strong decoherence takes the form
$`I`$ $``$ $`{\displaystyle \frac{w_0e(\mathrm{\Gamma }_{\varphi _A}+\mathrm{\Gamma }_{\varphi _B})}{w_A+w_B}}\mathrm{if}\mathrm{\Gamma }_\varphi 0,`$ (30a)
$`I`$ $``$ $`{\displaystyle \frac{4eT_AT_B}{w_A\mathrm{\Gamma }_{\varphi _A}+w_B\mathrm{\Gamma }_{\varphi ,B}+(w_A+w_B)\mathrm{\Gamma }_{\varphi _C}}}\mathrm{if}\mathrm{\Gamma }_\varphi \mathrm{},`$
with weight factors
$$w_0=\frac{T_AT_B}{T_A^2+T_B^2},w_A=\frac{\mathrm{\Gamma }_AT_B/T_A}{\mathrm{\Gamma }_A+\mathrm{\Gamma }_B},w_B=\frac{\mathrm{\Gamma }_BT_A/T_B}{\mathrm{\Gamma }_A+\mathrm{\Gamma }_B},$$
(31)
that are functions of the tunnel rates — but independent of the decoherence rates. Notice that in this asymmetric case the trapped state $`\sqrt{w_0T_B/T_A}|A\sqrt{w_0T_A/T_B}|B`$ has unequal weights on the two dots $`A`$ and $`B`$.
In conclusion, we have demonstrated how the well known concept of coherent population trapping in atoms may be transferred to a purely electronic system. A large voltage bias plays the role of the laser illumination and single-electron tunneling between quantum dots plays the role of intra-atomic transitions. Because the quantum dots are charged, the trapped electronic state is sensitive to decoherence by coupling to charges in the environment. This decoherence destabilizes the trapped state, causing a leakage current $`I`$ to flow through the quantum dots. We have found that the maximal $`I`$ in the presence of decoherence is set by the trapping rate $`\alpha `$, with $`I_{\mathrm{max}}e\alpha `$ within a factor of two over the whole parameter range. For small decoherence rate $`\mathrm{\Gamma }_\varphi `$ we find $`I=e\mathrm{\Gamma }_\varphi `$, which provides a way to measure the coherence time of a charge qubit in a transport experiment. We finally note that extensions to many-electron trapping can serve as a source for the formation of entangled electron pairs .
###### Acknowledgements.
We have benefited from discussions with B. Trauzettel. This work was supported by the Dutch Science Foundation NWO/FOM and the US NSF project DMR 0403465.
|
warning/0506/quant-ph0506119.html
|
ar5iv
|
text
|
# Weak cloning of an unknown quantum state
## Abstract
The impossibility to clone an unknown quantum state is a powerful principle to understand the nature of quantum mechanics, especially within the context of quantum computing and quantum information. This principle has been generalized to quantitative statements as to what extent imperfect cloning is possible. We delineate an aspect of the border between the possible and the impossible concerning quantum cloning, by putting forward an entanglement-assisted scheme for simulating perfect cloning in the context of weak measurements. This phenomenon we call weak cloning of an unknown quantum state.
Although the no-cloning principle wootters82 ; dieks82 forbids perfect cloning of an unknown quantum state, approximate cloning whose quality is independent of the input state is known to be possible buzek96 ; linares02 . The prize which is paid in such approximate cloning procedures is that no state is cloned perfectly; all states are cloned by a maximum fidelity less than one bruss98 . In contrast, we here demonstrate a scheme that we shall call ‘weak cloning’, which allows independent local weak measurements aharonov90 as if performed on a perfectly cloned unknown input state.
Weak measurements on post-selected quantum ensembles leads to the concept of weak values aharonov90 . Although conceptually similar to expectation values, these weak values may show curious behavior. For instance, they can lie far outside the range of allowed values of the measured observables aharonov87 ; aharonov88 or they can lie in classically forbidden regions of space aharonov93a ; aharonov93b . Weak cloning is a phenomenon that arises from weak measurements by post-selection.
To describe the idea of weak cloning, let us assume that two observers, Alice and Bob, can do experiments on qubits. Alice is given a qubit prepared in the state $`|\varphi _1=\alpha |0_1+\beta |1_1`$, where $`\alpha `$ and $`\beta `$ are unknown complex numbers with $`|\alpha |^2+|\beta |^2=1`$. Alice shares also a Bell state, such as $`|\phi _+_{23}=\left(|0_2|0_3+|1_2|1_3\right)/\sqrt{2}`$, with Bob. Thus, initially the total state reads
$`|\mathrm{\Phi }=|\varphi _1|\phi _+_{23}.`$ (1)
Alice and Bob perform weak measurements on qubits $`1`$ and $`3`$, respectively. These measurements can be modeled by the interaction Hamiltonian
$`H(t)`$ $`=`$ $`g_a(t)AI_2q_aI_3I_b`$ (2)
$`+g_b(t)I_1I_2I_aBq_b,`$
where $`q_a`$ and $`q_b`$ are Alice’s and Bob’s pointer position variable, respectively, $`A`$ and $`B`$ are the corresponding measured observables. The identity operators $`I_1,I_2,I_a`$ pertain to Alice’s two qubits and her pointer system, and $`I_3,I_b`$ are the identity operators of Bob’s qubit and pointer system. Note that the observable $`A`$ belongs to the first of Alice’s qubits, i.e., the one prepared in the unknown state $`|\varphi `$. The Hamiltonian in Eq. (2) results in an evolution that is local with respect to Alice’s and Bob’s locations, and leaves Alice’s second qubit unaffected. The time-dependent coupling parameters $`g_a`$ and $`g_b`$ turn on and off the interaction between the measurement devices and the measured qubits. With the measuring devices initially in the product state $`|m_a|m_b`$, we obtain the final total state as ($`\mathrm{}=1`$)
$`|\mathrm{\Gamma }`$ $`=`$ $`e^{i\gamma _aAq_a}I_2e^{i\gamma _bBq_b}|\mathrm{\Phi }|m_a|m_b`$ (3)
with the measurement strengths
$`\gamma _x={\displaystyle g_x(t)𝑑t},x=a,b.`$ (4)
The measurements are assumed to be weak upon the fulfillment of a certain weakness condition aharonov90 , which essentially entails that the measurement strengths should be much smaller, in some appropriate units, than the width of the corresponding (Gaussian) pointer wave functions. It can be shown aharonov90 that it is only necessary to keep terms to first order in the measurement strengths, in this weak measurement limit.
After the pointers have interacted weakly with the qubits Alice performs a post-selection of the state $`|\phi _+_{12}`$ on her qubit pair. Conditioned on the post-selection, one finds to first order in $`\gamma _a`$ and $`\gamma _b`$, i.e., in the weak measurement limit, the unnormalized state vector of the system and the pointers to be
$`|\mathrm{\Gamma }_{ps}`$ $``$ $`{\displaystyle \frac{1}{2}}|\phi _+_{12}|\varphi _3e^{i\gamma _aA_wq_a}|m_ae^{i\gamma _bB_wq_b}|m_b`$ (5)
$`{\displaystyle \frac{i}{2}}|\phi _+_{12}|\varphi ^{}_3(\gamma _a\varphi ^{}|A|\varphi q_aI_b`$
$`+\gamma _b\varphi ^{}|B|\varphi I_aq_b)|m_a|m_b,`$
where $`X_w=\mathrm{\Psi }|X|\mathrm{\Phi }/\mathrm{\Psi }|\mathrm{\Phi }`$, $`X=A,B`$, are the weak values with $`|\mathrm{\Psi }=|\phi _+_{12}|\varphi _3`$, and $`|\varphi ^{}`$ is orthogonal to $`|\varphi `$. Using the explicit form of $`|\mathrm{\Phi }`$ and $`|\mathrm{\Psi }`$ we obtain
$`X_w=\varphi |X|\varphi .`$ (6)
By taking the partial trace over the qubits, the state $`\rho _{ab}`$ of the measuring devices becomes
$`\rho _{ab}`$ $``$ $`{\displaystyle \frac{1}{4}}e^{i\gamma _aA_wq_a}|m_am_a|e^{i\gamma _aA_wq_a}`$ (7)
$`e^{i\gamma _bB_wq_b}|m_bm_b|e^{i\gamma _bB_wq_b}.`$
Thus, at the expense of entangled particle pairs, Alice and Bob can measure expectation values with respect to the same unknown state $`|\varphi `$ as shifts in the pointer momentum, in the weak measurement limit. This is identical to the result obtained if Alice and Bob had performed their local weak measurements on the state $`|\varphi |\varphi `$. In this sense, the above scheme simulates perfect cloning of quantum states.
Note that Bob may perform his measurement as soon as he has received his part in the entangled pair of qubits. He might even do this measurement before Alice receives the state $`|\varphi `$. After his measurement Bob awaits the message from Alice which details whether the measurement result should be kept or not. By repeating this procedure Bob can reconstruct his expectation value. This is in contrast with a scheme based on teleportation. In such a scheme, Alice first makes a weak measurement on the unknown state $`|\varphi `$ followed by a teleportation to Bob who makes his weak measurement on the teleported state. In this way, Alice and Bob also obtain the weak values $`A_w=\varphi |A|\varphi `$ and $`B_w=\varphi |B|\varphi `$, respectively. The difference is that the weak measurements have to be sequential since Bob has to await the arrival of the signal from Alice before he can do his measurement. In the weak cloning scheme, however, the weak measurements are ‘simultaneous’ in the sense that there is no causal ordering between the two measurement events.
The independence of the two measurement events can perhaps be discerned more clearly in a slightly modified version of the weak cloning scheme. As described above, Alice and Bob initially share an entangled state over systems $`2`$ and $`3`$. Alice performs the post-selection locally on her systems $`1`$ and $`2`$, and transmits the result to Bob. One can consider an alternative with Bob in sole possession of the entangled pair $`2`$ and $`3`$, whereas Alice only possesses system $`1`$. As a consequence, the post-selection on system $`1`$ and $`2`$ cannot be performed locally, as it directly involves both Alice and Bob. In this alternative scheme, Alice and Bob do not even share any correlation at the time of their weak measurements.
By sharing another Bell state with a third party (Charlie), Bob can post-select a Bell state conditioned on Alice’s post-selection. By performing a weak measurement of an observable $`C`$, say, Charlie can reconstruct the weak value $`C_w=\varphi |C|\varphi `$ upon receiving the results of Bob’s post-selection. Continuing in this way, the unknown state can be weakly cloned arbitrary many times, at the expense of an intensity loss that is exponential in the number of post-selection steps.
The weak cloning scheme can be seen as a parallelization of consecutive weak measurements. By a similar post-selection procedure one can obtain parallelization of other operations, like other types of measurements, or unitary transformations. It is important to note, though, that the post-selection procedure results in that Bob operates on the output of Alice, and similarly that Charlie operates on the output of Bob. That Bob operates, or measures, on the output of Alice clarifies the role of the weak measurement in the weak cloning scheme. Since in the weak limit Alice’s measurement does not perturb the state, Bob can measure as if on the original input state. This would not be the case if Alice made an ordinary type of measurement, and similarly for Charlie if Bob made an ordinary measurement.
Generalization of the above scheme to systems described by $`N3`$ dimensional Hilbert spaces can be achieved by replacing the Bell state by a maximally entangled state of the type $`|\psi =_j|j|j/\sqrt{N}`$. Alice then performs a post-selection by measuring an observable with $`|\psi `$ as a nondegenerate eigenvector. She communicates which of her systems pass the test to Bob, who in turn uses the corresponding data to obtain his weak value.
Experimental realization of the weak cloning scheme is most likely to be found for photonic systems. We note that two-photon systems have already been used to demonstrate some properties of weak values pryde04 . In view of the rapid development of techniques for quantum control in general and for photons in particular, we believe experimental test of weak cloning should be feasible in the near future.
|
warning/0506/quant-ph0506204.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Study of periodic potentials has evoked renewed interest in the literature in light of their appearance in Bose-Einstein condensates , and photonic crystals , , , . It has been possible to experimentally change structure of the potential, so as to produce superfluid - insulator transition , the former having delocalized states and the latter localized ones. In this context, traditional Kronig-Penny model is used for illustrative purposes, wherein the known wave functions lead to transcendental equations involving energy and momentum, when appropriate boundary conditions are implemented. The possibility of investigating superfluid-insulator type transitions mentioned above does not arise here, due to the lack of any control parameter. Quite sometime back, Scarf showed that a solvable model exists, which exhibits both discrete bound states and band spectra , as a function of the coupling parameter. The group theoretical aspects of this problem have recently been investigated . The fact that Scarf potential yields both bound states and band structure, as a function of a coupling parameter, makes this model an ideal one to study the interplay of these two types of distinct behavior in a given quantal problem.
The goal of this paper is to first map the Scarf eigenvalue problem into the zero energy sector of another quasi-exactly solvable (QES) problem. We then use the quantum Hamilton-Jacobi (QHJ) approach , which naturally takes advantage of the singularities of the new potential, to isolate the domains corresponding to discrete and band spectra. The subtle aspects of the boundary conditions in quantum mechanics, which lead to the existence of both bound states and band structure in the Scarf potential, come out naturally in this approach. We then proceed to obtain the eigenvalues and eigenfunctions, for both the cases. In this procedure, the energy eigenvalues can be obtained, without finding the eigenfunctions. The QHJ formalism, being formulated in the complex domain where the non-linear Riccati equation replaces the Schrödinger equation, makes use of powerful theorems in complex variable theory to obtain the solutions.
Apart from the fact that QHJ formalism is relatively new and requires detailed study, this approach may provide a different perturbative treatment for the traditional problems. As will be clear form the text, WKB approximation scheme is close to this method . In the following section, we briefly describe the working principles of the QHJ formalism, which is then used for the analysis of the Scarf potential in Section 3. The origin of the bound and the band spectra is then illustrated, without getting into the explicit computation of the eigenvalues, whose details are given in Section 3.2. We obtain the solutions pertaining to both the spectra. We conclude in the final section after pointing out various directions for future investigations.
## 2 Quantum Hamilton - Jacobi formalism
The QHJ formalism, formulated as a theory analogous to the classical canonical transformation theory , , , was proposed by Leacock and Padgett in 1983. It has been applied to one dimensional bound state problems and separable problems in higher dimensions . In our earlier studies, we have shown that one could use the QHJ formalism to analyze one dimensional exactly solvable (ES), quasi - exactly solvable models, consisting of both periodic and aperiodic potentials , , , , , and the recently discovered PT symmetric potentials . The advantage of this method lies in the fact that it requires a modest understanding of basic quantum mechanics and complex analysis as a prerequisite.
In this formalism, the logarithmic derivative of the wave function $`\psi (x)`$, given by
$$p=i\mathrm{}\frac{d}{dx}\mathrm{ln}\psi (x),$$
(1)
plays an important role. This is referred to as the quantum momentum function (QMF), since it is defined analogous to the classical momentum function as, $`p=\frac{dS}{dx}`$. Here, $`S`$ is the Hamilton’s characteristic function which is related to the wave function by $`\psi (x)=\mathrm{exp}(iS/\mathrm{})`$. Substituting $`\psi (x)`$ in terms of $`S`$ in the Schrödinger equation,
$$\frac{\mathrm{}^2}{2m}\frac{d^2\psi (x)}{dx^2}+V(x)\psi (x)=E\psi (x),$$
(2)
and using the relation between $`p`$ and $`\psi `$, one obtains the non-linear Riccati equation:
$$p^2i\mathrm{}p^{}=2m(EV(x)).$$
(3)
The above equation is known as the QHJ equation; here $`x`$ is treated as a complex variable, thereby extending the definition of $`p`$ to the complex plane. We show that one can arrive at the required results by studying the singularity structure of the QMF.
Singularity structure
The QMF has two types of singularities, the moving and the fixed singularities. From (1) one can see that, the $`n`$ nodes of the $`n^{th}`$ excited state, whose locations depend on the initial conditions and energy, correspond to the singularities of $`p`$. These are known as the moving singularities. It is a fact that only poles can appear as moving singularities in the solutions of the Riccati equation. One can calculate the residue at a moving pole $`x_0`$, where $`V(x)`$ is analytic, by doing a Laurent expansion of $`p`$ around $`x_0`$ as,
$$p=\underset{k=1}{\overset{l}{}}(xx_0)^k+\underset{k=0}{\overset{\mathrm{}}{}}(xx_0)^k.$$
(4)
Substituting this in (3) and comparing individually the coefficients of different powers of $`xx_0`$, one obtains $`l=1`$, with the corresponding residue equalling $`i\mathrm{}`$.
The fixed singularities originate from the potential and are present in all the solutions of the Riccati equation. One can calculate the residue at the fixed poles in the same way, as is done for the moving poles. Owing to the quadratic nature of the QHJ equation, one obtains two solutions. In order to arrive at the right solution, one needs to choose the residue that gives the correct physical behavior. The right value of the residue is chosen by applying the appropriate boundary conditions, details of which will be given in the text, as and when required. Thus, knowing the singularity structure of QMF and the behavior of $`p`$ at infinity, one gets the complete form of the QMF. In all the models studied so far, including the periodic potentials, the assumption that, the QMF has finite number of singularities, is equivalent to saying that the point at infinity is an isolated singular point, has been found to be true. We expect it to be valid for the present case also.
For most exactly solvable models, the QMF has been found to be a rational function. As is known, for a rational function the sum of all residues including that at infinity is zero. This result has been used to obtain the energy eigenvalues for all the models studied in the QHJ approach. It should be pointed out that, this condition is equivalent to the exact quantization condition satisfied by the action $`J`$ :
$$J=_Cp𝑑x=n\mathrm{}.$$
(5)
Hence, for the case of Scarf potential, one first tries to bring the QMF into a rational form through a suitable change of variable, as discussed in the next section. It is interesting to note that in the classical limit,
$$pp_c=\sqrt{2m(EV(x))},$$
(6)
where $`p_c`$ is the classical momentum. The QHJ quantization condition then leads to the WKB approximation scheme. The boundary condition (6) was originally used by Leacock and Padgett to obtain the constraints on the residues.
## 3 The Scarf Potential
The Scarf potential is given by
$$V(x)=\left(\frac{(\frac{1}{4}s^2)\pi ^2}{2ma^2\mathrm{sin}^2(\frac{\pi x}{a})}\right),$$
(7)
where, $`a`$ is the potential period. One finds that in the range $`s>1/2`$, the potential is an array of infinite potential wells as shown in Fig 1.
A quantum particle is then confined to only one well, implying that the wave function should vanish at $`x=\pm a`$. Thus, in the above range, the potential exhibits bound state spectra. As shown in Fig.2, in the range $`0<s<1/2`$, the potential is similar to that of a potential in a crystal lattice,
leading to the possibility of energy bands. In this scenario, a particle can escape to infinity. Therefore the wave function need not vanish at $`x=\pm a`$. However, $`\psi (x)`$ should not diverge anywhere, on physical grounds.
The QHJ equation for the Scarf potential, with $`p=iq`$ and $`\mathrm{}=1`$ in (3), is given by,
$$q^2+q^{}+\frac{\pi ^2}{a^2}\left(\lambda ^2+\frac{(\frac{1}{4}s^2)}{\mathrm{sin}^2(\frac{\pi x}{a})}\right)=0,$$
(8)
where $`\lambda ^2=2mEa^2/\pi ^2`$. We perform a change of variable using
$$y=\mathrm{cot}\left(\frac{\pi x}{a}\right),$$
(9)
which transforms (8) to
$$q^2(y)\frac{\pi }{a}(1+y^2)\frac{dq}{dy}+\frac{\pi ^2}{a^2}\left(\lambda ^2+(\frac{1}{4}s^2)(1+y^2)\right)=0.$$
(10)
In order to get all the coefficients in the above equation to a rational form, which in turn will easily yield the singularities of the QMF, we use the transformation equations
$$q=\frac{\pi \mathrm{}}{a}(1+y^2)\varphi ;\varphi =\chi \frac{y}{1+y^2}.$$
(11)
This leads to the QHJ equation in terms of $`\chi `$ as,
$$\chi ^2+\frac{d\chi }{dy}+\frac{\lambda ^21}{(y^2+1)^2}+\frac{(\frac{1}{4}s^2)}{y^2+1}=0.$$
(12)
Henceforth, the above equation will be treated as the QHJ equation and $`\chi `$ as the QMF. It is interesting to note that, substituting $`\chi =\frac{d}{dy}(\mathrm{ln}(\stackrel{~}{\psi }(y))`$ in the above equation, one gets a Schrödinger equation, which describes the zero energy sector of the potential, $`(\lambda ^21)/(y^2+1)^2+(\frac{1}{4}s^2)/(y^2+1)`$. By analyzing the singularity structure of this quasi-exactly solvable problem, we obtain the required results for the solvable Scarf potential, as shown below.
### 3.1 Form of the QMF $`\chi `$
The QMF has $`n`$ moving poles with residue one on the real line, as is clear from the Riccati equation. From (12), one can see that $`\chi `$ has fixed poles at $`y=\pm i`$. Making use of the assumption that the QMF has finite number of moving poles, one can write $`\chi `$ in the rational form, separating its analytical and singular parts as
$$\chi =\frac{b_1}{yi}+\frac{b_1^{}}{y+i}+\underset{k=0}{\overset{n}{}}\frac{1}{yy_k}+Q.$$
(13)
Here, $`b_1`$ and $`b_1^{}`$ are the residues at $`y=i`$ and $`y=i`$ respectively and the summation term describes the sum of all the singular parts coming from the moving poles. Note that, $`_{k=0}^n\frac{1}{yy_k}=\frac{P_n^{}(y)}{P_n(y)}`$, where $`P_n(y)`$ is an $`n^{th}`$ degree polynomial. The quantity $`Q`$ represents the analytic part of $`\chi `$ and from (12) one can see that $`\chi `$ is bounded for large $`y`$. Thus, from Liouville’s theorem, $`Q`$ is a constant; denoting it as $`C`$, (13) can be written as,
$$\chi =\frac{b_1}{yi}+\frac{b_1^{}}{y+i}+\frac{P_n^{}(y)}{P_n(y)}+C.$$
(14)
One can calculate the residues at the fixed poles $`y=\pm i`$, by making a Laurent expansion of $`\chi `$ around the pole. For example, to calculate the residue at $`y=i`$, we expand $`\chi `$ as,
$$\chi =\frac{b_1}{yi}+a_0+a_1(yi)+\mathrm{}.$$
(15)
Comparing the coefficients of different powers of $`(yi)`$ individually, one obtains
$$b_1=\frac{1\pm \lambda }{2}.$$
(16)
Similarly the other residue at $`y=i`$ is found to be
$$b_1^{}=\frac{1\pm \lambda }{2}.$$
(17)
To find the eigenvalues, we now make use of the fact that, for a rational function, the sum of all the residues equals zero. As noted earlier, this is equivalent to the quantization condition (5) of Leacock and Padgett. Thus, we obtain
$$b_1+b_1^{}+n=d_1,$$
(18)
where $`d_1`$ is the residue at infinity, which is calculated by taking Laurent expansion of $`\chi `$ around the point at infinity:
$$\chi =d_0+\frac{d_1}{y}+\frac{d_2}{y^2}+\mathrm{}.$$
(19)
Substitution of the above in the QHJ equation yields,
$$d_1^2d_1+(\frac{1}{4}s^2)=0,$$
(20)
from which, the values of $`d_1`$ can be deduced:
$$d_1=\frac{1\pm 2s}{2}.$$
(21)
Substituting the values of the residues in (18), one obtains
$$n=\frac{1}{2}\pm s\lambda ,$$
(22)
which gives the degree of the polynomial $`P_n(y)`$ in (14). From the definition of $`\lambda `$, one can see that if $`E<0`$, $`\lambda `$ becomes imaginary, in which case (22) will not be satisfied. Thus, from the above equation, we have the condition $`E>0`$, which in turn implies $`\lambda >0`$ and real. Hence, for any range of $`s`$, the energy eigenvalues are greater than zero. With this condition on $`\lambda `$, we now proceed to select the values of the residues at the fixed poles and at infinity, which will give us the physically acceptable results.
### 3.2 Choice of the residues
One needs to use the boundary conditions obeyed by the QMF to choose the right value of residues. Although, there are several ways of implementing the boundary conditions in the QHJ formalism, we have chosen the one closest to the conventional approach for clarity. First, we shall fix the value of the residue at infinity. From the prior discussion of the potential, we know that the wave functions should not become infinite anywhere, in particular, for $`x=\pm a`$. From (1), one can obtain $`\psi (x)`$ in terms of the QMF. Writing $`p=iq`$ and doing the change of variable, one obtains the wave function:
$$\psi (y)=\mathrm{exp}\left(\frac{a}{\pi }\left(\frac{q}{1+y^2}\right)𝑑y\right).$$
(23)
Using the transformation equations in (11), the above expression for the wave function becomes
$$\psi (y)=\mathrm{exp}\left(\left(\chi \frac{y}{1+y^2}\right)𝑑y\right).$$
(24)
For large $`y`$, the leading behavior of $`\chi `$ is obtained as $`\chi \frac{d_1}{y}`$, which when substituted in (24), yields,
$$\psi (y)\mathrm{exp}\left(\left(\frac{d_1}{y}\frac{y}{1+y^2}\right)𝑑y\right)$$
(25)
$$\frac{y^{d_1}}{(y^2+1)^{1/2}}.$$
(26)
Using the value of $`d_1`$ from (21) in the above equation, one obtains
$$\psi (y)\frac{y^{\frac{1}{2}\pm s}}{(y^2+1)^{1/2}}.$$
(27)
For $`0<s<\frac{1}{2}`$, one can see that $`\psi 0`$, in the limit $`y\mathrm{},xma`$, with $`m`$ being an integer, for both the values of $`d_1`$. This range corresponds to the case where the potential exhibits band structure.
For $`s>1/2`$, $`\psi (y)0`$, in the limit $`y\mathrm{},xma`$, with $`m`$ being an integer, only if $`d_1`$ takes the value $`1/2s`$. In this way, the two different ranges of the potential parameter $`s`$ emerge simultaneously, while fixing the values of $`d_1`$. In order to select the values of $`b_1`$ and $`b_1^{}`$, we note that the bound state and band edge wave functions of one dimensional potentials are non-degenerate and have definite parity. Parity operation requires that $`\chi (y)=\chi (y)`$, which in turn gives
$$b_1=b_1^{}.$$
(28)
With the above constraint, (18) becomes
$$2b_1+n=d_1.$$
(29)
Finiteness of the wave function as $`x\mathrm{}`$, gives the values of $`d_1`$ in the two ranges as
$$d_1=\{\begin{array}{cc}\frac{1\pm 2s}{2}\text{for}0<s<1/2,& \\ & \\ \frac{12s}{2}\text{for}s>1/2.& \end{array}$$
(30)
From the parity constraint, one obtains the restriction on the values of the residues at the fixed poles as $`b_1=b_1^{}`$. Using these results, we proceed proceed to calculate the solutions for the two ranges.
### 3.3 Case 1 : Band spectrum
In the range $`0<s<1/2`$, we have seen that $`d_1`$ can take both the values of the residues. Taking all the possible combinations of the residues, with $`b_1=b_1^{}`$ and substituting them in (18), we evaluate $`n`$, the degree of the polynomial $`P_n(y)`$. There are four combinations forming four different sets, as given in the fifth column of table I.
Since $`n`$ needs to be positive, from table I, we pick only those sets which give a positive integral value for $`n`$. As seen earlier, $`\lambda `$ is a positive real constant. Thus, only the sets 1 and 2 will yield positive values for $`n`$ and hence; the other two sets are ruled out. Taking the values of $`b_1`$ and $`d_1`$ from the sets 1 and 2, substituting them in (29) and using the definition of $`\lambda `$ and $`s`$, we obtain the expressions for the energy eigenvalues corresponding to the two band edges of the $`n^{th}`$ band as,
$$E_n^\pm =\frac{\pi ^2}{2ma^2}\left(n+\frac{1}{2}\pm s\right)^2.$$
(31)
Here, $`E_n^\pm `$ correspond to the upper and lower band energies of the $`n^{th}`$ band. These results match with the solutions given in and .
The corresponding wave functions follow from (14) and (24):
$$\psi (y)=(y^2+1)^{b_1\frac{1}{2}}P_n(y).$$
(32)
To obtain the expression for the polynomial, we substitute $`\chi `$ from (14) in the QHJ equation, which gives a second order differential equation:
$`P_n^{\prime \prime }(y)+\left({\displaystyle \frac{4b_1y}{y^2+1}}\right)P_n^{}(y)+`$
$`\left({\displaystyle \frac{1/4s^2}{y^2+1}}+{\displaystyle \frac{4b_1^2y^2+2b_1(1y^2)+\lambda ^21}{(y^2+1)^2}}\right)P_n(y)=0.`$ (33)
The sets 1 and 2 have $`b_1=(1\lambda )/2`$, which yields
$$P_n^{\prime \prime }(y)+\left(\frac{2(1\lambda )y}{y^2+1}\right)P_n^{}(y)+\frac{\frac{1}{4}s^2+\lambda ^2\lambda }{(y^2+1)}P_n(y)=0.$$
(34)
From (31), one can see that $`\lambda `$ has two values $`\lambda =n\pm s+\frac{1}{2}`$. Substituting these in the above equation, one obtains two differential equations corresponding to the two energy eigenvalues $`E_n^\pm `$ :
$$(y^2+1)P_n^{\prime \prime }(y)+(12n2s)yP_n^{}(y)+n(n\pm 2s)P_n(y)=0.$$
(35)
Defining $`y=it`$, the above equation takes the form of the well known Jacobi differential equation
$$(1t^2)P_n^{\prime \prime }(t)+(\nu _1\nu _2t(\nu _1+\nu _2+2))P_n^{}(t)+n(n+\nu _1+\nu _2+1)P_n(t)=0,$$
(36)
with $`\nu _1=\nu _2=ns1/2`$, for the corresponding two $`\lambda `$ values. The expression for the two band edge wave functions for the $`n^{th}`$ band are given by,
$$\psi (y)=(y^2+1)^{\frac{\lambda }{2}}P_n^{\nu _1,\nu _2}(iy),$$
(37)
with their respective $`\nu _1,\nu _2`$ values corresponding to $`\lambda =n\pm s+1/2`$.
### 3.4 Case 2 : Bound state spectrum
We proceed in the same way as in case 1 i.e., take all possible combinations of $`b_1,b_1^{}`$ and $`d_1`$, keeping $`b_1=b_1^{}`$ in (29). Since $`d_1`$ can take only one value $`1/2s`$, only two sets are possible here. Out of these, the set corresponding to $`b_1=b_1^{}=(1\lambda )/2`$ alone, will give a positive value for $`n`$. Thus, substituting these values of residues in (29), one obtains the expression for the energy eigenvalue as
$$E_n=\frac{\pi ^2}{2ma^2}\left(\frac{1}{2}+n+\sqrt{\frac{1}{4}\frac{2mV_0a^2}{\pi ^2\mathrm{}^2}}\right)^2,$$
(38)
where $`n`$ can take positive integral values. Proceeding as above one obtains the Jacobi differential equation in terms of $`t`$ for the polynomial part :
$$(1t^2)P_n^{\prime \prime }(t)2t(ns+\frac{1}{2})P_n^{}(t)n(n+2s)P_n(t)=0.$$
(39)
The expression for the wave function is then given by
$$\psi (y)=(y^2+1)^{\frac{\lambda }{2}}P_n^{s_1,s_2}(iy),$$
(40)
where $`s_1=s_2=ns1/2`$.
Hence, as pointed out in the beginning, the two different sectors of the Hamiltonian, as a function of the coupling parameter and the eigenvalues emerge from general principles of QHJ formalism, relying on the singularity structure of the QMF function. The wave functions corresponding to the definite eigenvalues are obtained at the end, which match with the known results .
Conclusions
We have mapped the entire Scarf problem, containing the bound state and energy bands, to the zero energy sector of a different Hamiltonian, which is quasi-exactly solvable. This was achieved through point canonical transformations which led to the redistribution of singularities in the complex domain. The singularity structure of this new Hamiltonian is transparent enough to clearly isolate two different regimes, as a function of the coupling constant. When related to the original problem, they turn out to represent discrete levels and the band edges. It will be interesting to carefully analyze the equilibrium structure of the classical electrostatics problem, associated with the QES system, which leads to both bound states and band structure in the quantum domain. In light of the current interest in periodic potentials in BEC and photonic crystals, we hope the quantum Hamilton-Jacobi based treatment presented here is not only illuminating, but may also lead to development of new perturbative treatments for non-exactly solvable problems.
Acknowledgements
We are thankful to Dr.J. Banerji for a careful reading of the manuscript and R. Atre for his help during the course of this work.
|
warning/0506/cond-mat0506597.html
|
ar5iv
|
text
|
# Random quantum Ising chains with competing interactions
(22 Juin 2005)
## Abstract
In this paper we discuss the criticality of a quantum Ising spin chain with competing random ferromagnetic and antiferromagnetic couplings. Quantum fluctuations are introduced via random local transverse fields. First we consider the chain with couplings between first and second neighbors only and then generalize the study to a quantum analog of the Viana-Bray model, defined on a small world random lattice. We use the Dasgupta-Ma decimation technique, both analytically and numerically, and focus on the scaling of the lattice topology, whose determination is necessary to define any infinite disorder transition beyond the chain. In the first case, at the transition the model renormalizes towards the chain, with the infinite disorder fixed point described by Fisher. This corresponds to the irrelevance of the competition induced by the second neighbors couplings. As opposed to this case, this infinite disorder transition is found to be unstable towards the introduction of an arbitrary small density of long range couplings in the small world models.
Quantum fluctuations play a crucial role in the spin glass phases of the Sr-doped cuprate La<sub>2</sub>CuO<sub>4</sub> Kastner1998 , or the dipolar glass LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> in a transverse field Aeppli1987 . Randomly coupled quantum two level systems also appear in the understanding of the dielectric response of low temperature amorphous solids Ludwig2003 , and as the main low frequency source of decoherence of solid state quantum bits Galperin2003 . In all cases, the quantum fluctuations compete with the random couplings between the spin, and tend to disorder the corresponding random ordered phases.
One of the simplest random quantum model to study this competition is probably the random Ising spin model in a transverse magnetic field.
$$H=\underset{i,j}{}J_{ij}\sigma _i^z\sigma _{i+1}^z\underset{i}{}h_i\sigma _i^x,$$
(1)
where $`\sigma _i^x,\sigma _i^z`$ are the usual Pauli matrices, and the transverse fields $`h_i`$ are responsible for the quantum tunneling fluctuations between the up and down states of the Ising spins. In a pioneering work, Fisher has given asymptotically exact results for this random quantum Ising model with first neighbors random ferromagnetic bonds in one dimension Fisher1995 . By using a decimation technique developed by Dasgupta and Ma Dasgupta1980 , he described the infinite disorder quantum phase transition of this model. Some of the main features of this peculiar transition were a diverging dynamical exponent $`z`$, and very strong inhomogeneities manifesting through drastically different behavior between average and typical correlation functions.
Natural extensions of these results to higher dimensions have proved to be difficult. In particular, an analytical implementation of the Ma-Dasgupta decimation beyond the simple chain is extremely cumbersome. The reason is that any initial lattice except the chain is quickly randomized by the decimation. Thus one has to resort to a numerical implementation of this decimation Motrunich2000 ; Rieger2000 . For two-dimensional regular lattices, the results for the random ferromagnetic Ising model are consistent with the survival of an infinite disorder quantum phase transition, albeit with exponents different from the one-dimensional case Motrunich2000 . On the other hand, the quantum Ising spin glass, corresponding to the model (1) with both ferromagnetic (positive) and antiferromagnetic (negative) couplings, was also studied in two and three dimensions via Monte-Carlo simulations Rieger1996 . The numerical works found no sign of an infinite disorder quantum critical point. Since the results for the random ferromagnetic model are expected to extends to the quantum Ising spin glass, this discrepancy certainly deserves further work.
In this perspective, we investigate in this paper the stability of infinite disorder fixed point of the quantum Ising spin glass chain with respect to competing further neighbors couplings in two extreme cases. In a first step, we focus in the case where second neighbors couplings are present in model (1) besides the first neighbors couplings. The couplings are taken as either ferromagnetic or antiferromagnetic. By combining analytical (for small second neighbors couplings) and numerical decimation techniques we investigate the relevance of the presence of higher range couplings and their interplay with random signs in the couplings. We pay a special attention to the topology of the renormalized lattice, which appears crucial in the precise characterization of infinite disorder transitions.
A natural complement to this first case consists in considering this quantum Ising spin glass on a random network, obtained by adding a finite density of long range couplings between the chain’s sites. Indeed, our model can be considered as a quantum analog of the classical spin models of Viana and BrayViana1985 , although we keep a local regular topology besides the random long range couplings in our small world latticeAlbert2002 . In the case of classical spin glasses, these random lattice models are a natural extrapolation between the short-range model and its mean-field version. They undergo a finite temperature transition of the mean-field type, albeit with peculiarities induced by the finite connectivityNikoletopoulos . Wether the infinite disorder physics survives to this tendency towards mean-field like physics is the natural question we will consider.
In the first part of this letter, we consider the model (1) on a chain, with only first and second neighbors couplings (Zig-Zag ladder) footnote2 . Both the first and second neighbors couplings $`J_{i,i+1}^{(1)},J_{i,i+2}^{(2)}`$ can be antiferromagnetic ($`<0`$) with probability $`p`$, and ferromagnetic with probability $`1p`$. The $`|J_{i,i+1}^{(1)}|`$ are uniformly distributed between $`0`$ and $`1`$, the $`|J_{i,i+2}^{(2)}|`$ between $`0`$ and $`J_{max}^{(2)}`$, and the transverse fields $`h_i`$ between $`0`$ and $`h_{max}`$. Note that via an appropriate unitary transformation we can map this system onto one where only the second neighbors couplings can have both signs, but at the cost of a modification of the magnetic properties of the system. Hence for clarity, we prefer to consider only the more natural choice defined above.
We will analyze the low temperatures behavior of this system by means of the Dasgupta-Ma decimation technique Dasgupta1980 which was exploited by Fisher Fisher1995 in the case, among others, of the random ferromagnetic quantum Ising chain. Its extension to the present case of mixed coupling (anti-ferromagnetic and ferromagnetic) contains one supplementary rule as detailed below. The running energy scale $`\mathrm{\Omega }`$ is defined as the maximum of the amplitudes of bonds $`|J_{ij}|`$ and fields $`h_i`$. At each decimation step, if this maximum corresponds to a field $`h_i`$, the corresponding spin is frozen in the $`x`$ direction, generating new couplings $`\stackrel{~}{J}_{jk}=J_{jk}+(J_{ij}J_{ik}/\mathrm{\Omega })`$ between all pairs $`(j,k)`$ previously connected with the spin $`i`$. On the other hand, if the maximum is a ferromagnetic coupling $`J_{ij}`$, the two spins $`i`$ and $`j`$ are paired to form a new cluster $`[ij]`$ of magnetization $`\mu _{[ij]}=\mu _i+\mu _j`$ (where $`\mu _i`$ corresponds to the magnetization of cluster $`i`$), and coupling with site $`k`$ $`J_{[ij]k}=J_{ik}+J_{jk}`$ Fisher1995 . The new rule occurs when this maximum coupling is anti-ferromagnetic. In this case, if e.g the magnetization $`\mu _i`$ is larger than $`\mu _j`$, then the new cluster’s magnetization reads $`\mu _{[ij]}=\mu _i\mu _j`$, and the interaction with a third spin $`k`$ is $`J_{[ij]k}=J_{ik}J_{jk}`$. In both cases, the effective transverse field acting on the new cluster is $`h_{[ij]}=h_ih_j/\mathrm{\Omega }`$.
An analytical study of the scaling behavior of the model (1) under the above decimation rules is difficult even on the Zig-Zag ladder we consider. As mentioned in the introduction, couplings $`J_{ij}`$ are quickly generated on many length scales $`|ij|`$, and the initial lattice is quickly randomized (see below). To fix the notation and clarify the procedure, it is useful to start by considering the evolution under the RG of the first neighbors chain, extending the result of Ref. Fisher1995, to the presence of anti-ferromagnetic couplings. We introduce the convenient logarithmic variables $`\beta _i=\mathrm{ln}(\mathrm{\Omega }/h_i)`$, $`\zeta _{i,i+1}=\mathrm{ln}(\mathrm{\Omega }/|J_{i,i+1}|)`$ and scaling parameter $`\mathrm{\Gamma }:=\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })`$ where $`\mathrm{\Omega }_0`$ is the initial value of $`\mathrm{\Omega }`$. Their “distributions” densities are defined as $`(\beta ,\mathrm{\Gamma })`$ for the fields, $`𝒫^{(1+)}(\zeta ,\mathrm{\Gamma })`$ for the ferromagnetic bonds, and $`𝒫^{(1)}(\zeta ,\mathrm{\Gamma })`$ for the anti-ferromagnetic bonds. Note that while $`(\beta ,\mathrm{\Gamma })`$ is normalized, for the bonds only the sum $`𝒫^{(1)}(\zeta ,\mathrm{\Gamma })=𝒫^{(1+)}(\zeta ,\mathrm{\Gamma })+𝒫^{(1)}(\zeta ,\mathrm{\Gamma })`$ has a norm one. As can be deduced by a gauge transformation of (1), $`(\beta ,\mathrm{\Gamma })`$ and $`𝒫^{(1)}(\zeta ,\mathrm{\Gamma })`$ satisfy the same differential scaling equations than in the ferromagnetic case Fisher1995 provided we use the maximum instead of the sum in the above decimation rules, which is valid for broad enough distributions. Finally, the function $`𝒟(\beta ,\mathrm{\Gamma })=𝒫^{(1+)}(\zeta ,\mathrm{\Gamma })𝒫^{(1)}(\zeta ,\mathrm{\Gamma })`$ is found to satisfy the same scaling equation than $`𝒫^{(1)}`$. The fixed point $`=𝒫^{(1)}=𝒫^{}(x,\mathrm{\Gamma })=e^{x/\mathrm{\Gamma }}/\mathrm{\Gamma }`$ of the ferromagnetic chainFisher1995 is easily extended to the two following case : the above ferromagnetic point now corresponds to the solution $`𝒟=𝒫^{}`$, or $`𝒫^{(1+)}=𝒫^{(1)}=𝒫^{}`$, $`𝒫^{(1)}=0`$. As expected, it can be explicitely checked in the RG equations that this fixed point is unstable towards the proliferation of anti-ferromagnetic bonds. The new transition point corresponds to the solution $`𝒟=0`$, or $`𝒫^{(1+)}=𝒫^{(1)}=𝒫^{}/2`$, corresponding to an equal density of random positive and negative couplings. Hence, we will loosely call it the spin glass fixed point by analogy with the physics of the classical model in higher dimensions. The characteristics of the transtion from the the ferromagnetic to the disordered phase obtained by Fisher Fisher1995 translate to the present spin-glass fixed point into an average linear susceptibility (under the application of a small $`z`$ field $`\stackrel{~}{h}`$) which diverges as $`\chi (T)|\mathrm{ln}T|^{\varphi 2}/T`$ where $`\varphi =(1+\sqrt{5})/2`$ . Similarly, we extract the scaling behaviour of the average non linear susceptibility $`\chi _{nl}(T)=\left[\frac{^3}{\stackrel{~}{h}^3}|_{\stackrel{~}{h}=0}<M>(\stackrel{~}{h})\right]`$, where $`[\mathrm{}]`$ denotes an ensemble average, as $`\chi _{nl}(T)|\mathrm{ln}T|^{2\varphi 2}/T^3`$.
Having clearly defined the notation and fixed points for the chain, we can now study perturbatively their stability with respect to small second neighbors competing interactions. To first order, such an analysis can be conducted by considering the presence of $`J^{(2)}`$ negligible compared to the $`J^{(1)}`$, and checking whether this condition is self-consistently preserved under the re-scaling. More precisely, we will assume that (i) a $`J_{i,i+2}^{(2)}`$ will never constitute the highest energy in the system and therefore never be decimated (ii) in sums, the $`J_{i,i+2}^{(2)}`$ are negligible with respect to $`J_{i,i+1}^{(1)}`$ (iii) creation of third neighbour couplings out of second neighbour couplings can be neglected. As above, we define the distribution $`𝒫^{(2)}(\zeta ,\mathrm{\Gamma }):=𝒫^{(2+)}(\zeta ,\mathrm{\Gamma })+𝒫^{(2)}(\zeta ,\mathrm{\Gamma })`$ as the sum of “distributions” of positive and negative next nearest neighbour couplings. With the above hypothesis, its scaling behavior is found to be described by
$`{\displaystyle \frac{𝒫^{(2)}(\zeta )}{\mathrm{\Gamma }}}={\displaystyle \frac{𝒫^{(2)}(\zeta )}{\zeta }}𝒫^{(2)}(\zeta )\left(2(0)+𝒫^{(1)}(0)\right)`$
$`+2(0){\displaystyle _0^{\mathrm{}}}d\zeta _1d\zeta _2𝒫^{(1)}(\zeta _1)𝒫^{(2)}(\zeta _2)\delta (\zeta \zeta _1\zeta _2)`$
$`+𝒫^{(1)}(0)\delta (\zeta \mathrm{\Lambda })`$ (2)
where $`\mathrm{\Lambda }`$ is an arbitrary large constant which stands for the negligible $`J`$ in log. coordinates, and is taken to $`\mathrm{}`$ at the end of calculations. The $`\mathrm{\Gamma }`$ dependance of the distribution has been omitted for clarity. With the above hypothesis, the probability distributions for fields and nearest neighbour couplings still follow the equations for the chain. Hence, at the “Spin Glass” critical point, we can insert the scaling form $`=𝒫^{(1)}=𝒫^{}`$ in (2). It is usefull to split $`𝒫^{(2)}(\zeta ,\mathrm{\Gamma })`$ into a $`\mathrm{\Lambda }`$ independent part $`𝒫_i^{(2)}(z,\mathrm{\Gamma })`$ and $`𝒫_\mathrm{\Lambda }^{(2)}(z,\mathrm{\Gamma })`$. By denoting $`p(z,\mathrm{\Gamma })`$ the Laplace transform in $`\zeta `$ of $`𝒫^{(2)}(\zeta ,\mathrm{\Gamma })`$, we find that for $`z`$ and $`\mathrm{\Gamma }`$ finite and fixed, $`p_\mathrm{\Lambda }(z,\mathrm{\Gamma })0`$ when $`\mathrm{\Lambda }\mathrm{}`$. Then we show that the norm of the two parts of the solution satisfy : $`𝒫_i^{(2)}_\zeta =1𝒫_\mathrm{\Lambda }^{(2)}_\zeta =lim_{z0}p_i(z,\mathrm{\Gamma })=\mathrm{\Gamma }_0/\mathrm{\Gamma },`$ corresponding to a constant “decrease” of the couplings $`J^{(2)}`$. In this regime, the system “forgets” its initial conditions and flows to a general state gouverned by $`𝒫_\mathrm{\Lambda }^{(2)}`$. Consistency of conditions (i) to (iii) can also easily being checked from the properties of the Laplace transform. As a consequence, such small next nearest neighbor couplings correspond to an “irrelevant perturbation” at this infinite disorder fixed point.
To go beyond this perturbative analysis, we have studied the scaling behavior of the zig-zag ladder by implementing numerically the above renormalization rules. We start by choosing a random configuration of fields $`h_i`$ and couplings $`J_{i,i+1}^{(1)},J_{i,i+2}^{(2)}`$ according to the previous initial distributions probabilities. Then at each step, the energy scale is lowered and the number $`N`$ of spins is reduced by $`1`$ according to the decimations rules specified above. This process is continued up to the last remaining spin, and repeated for a number $`R=10^3`$ configurations. No assumption is made on the topology of the renormalized lattice, and we keep a priori all generated couplings. However, for practical reasons it appears necessary to restrict ourselves to energies larger than a lower cut-off $`\mathrm{\Omega }_{min}`$. With this procedure, the distributions $`𝒫(\zeta ,\mathrm{\Gamma }),(\beta ,\mathrm{\Gamma })`$ are correctly sampled below $`\mathrm{\Gamma }_{max}=\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega }_{min})`$ Motrunich2000 . For most of our results, this cut-off $`\mathrm{\Omega }_{min}`$ was maintained to negligible values, without any noticeable incidence on the results. For fixed $`J_{max}^{(2)}`$, the transition is reached by varying the maximum amplitude $`h_{max}`$ of the fields. We locate a putative infinite disorder phase transition by using the analogy with percolation Monthus1997 . At each decimation step $`i`$, corresponding to a system size $`N_0i`$, we consider the number of realizations $`n_h(i)`$ where a field was decimated at step $`i`$, the similarly for the bonds $`n_J(i)`$. At the transition, the ratio $`n_h(i)/n_J(i)`$ should become scale invariant, whereas it should diverge or decrease to zero respectively in the disordered or ordered random phases. Moreover, the scaling behavior of this ratio is an excellent way to check for possible finite size effects respective to the topology of the initial lattice. The inset of the figure 1 shows this scaling of the decimation ratio for two values around the candidate critical value of $`h_{max}`$. Once such candidates for the transition are determined, we have studied the scaling behavior of the distributions functions $`𝒫(\zeta ,\mathrm{\Gamma }),(\beta ,\mathrm{\Gamma })`$, of the distribution of magnetization $`\mu (\mathrm{\Gamma })`$, and number of active spins $`n(\mathrm{\Gamma })`$ in the clusters. This allows to characterize the criticality of the infinite disorder fixed point. Moreover, to fully characterize an infinite disorder fixed point beyond the simple chain, one should also be able to determine the renormalized topology of the critical lattice, and the associated correlations with the couplings. In a first attempt to study the scaling of this topology, we have followed the distribution of the connectivity of the lattice as the decimation goes on. The results, depicted on fig. 1, shows that while initially all sites have only $`4`$ neighbors, the distribution $`P(c)`$ flows towards an intermediate algebraic distribution at intermediates sizes. While highly connected sites appear, we find by varying our lower cut-off $`\mathrm{\Gamma }_{max}`$ that rather strong correlations exist between the bonds connecting these sites. And while the decimation is pursued, the distribution narrow back towards a delta function peaked on $`c=2`$, i.e the lattice is ultimately renormalized towards a chain. We thus find that for the $`J_1J_2`$ model, the infinite disorder fixed point is always given by the fixed point of the chain (see above and Fisher1995 ), in agreement with previous results on the similar ferromagnetic two-leg ladder Rieger2000
The previous results motivated the study of the opposite limit of long-range couplings competing with the initial couplings of the chain. Thus we naturally consider the hamiltonian (1) on a small world latticeAlbert2002 , where beyond the previous nearest neighbors couplings $`J_{i,i+1}^{(1)}`$, we add random infinite range couplings $`J_{i,j}^{LR}`$ between any two non-neighbor sites $`i`$ and $`j`$, with density $`p/N`$. In this paper, the existing couplings $`J_{i,j}^{LR}`$ and $`J_{i,i+1}^{(1)}`$ are distributed with the same uniform distribution between $`0`$ ad $`1`$. With these conventions, the average initial connectivity of this lattice is $`2+q`$. Results of the same numerical decimation procedure as above indicate a phase transition different from the previous one (Zig-Zag ladder). In particular, contrarily to the previous case, the distribution of connectivity of the renormalized lattice broadens without limit up to some finite size effects. This cross-over happens when the numerical upper bound of the renormalized distribution $`P(c)`$ becomes of the order of the system size. Once this happens, highly connected sites proliferate, leading to a mean-field like behavior.
In this paper we have shown how the presence of random signs and further neighbor couplings affect the critical behavior of the random quantum Ising chain. We have particularly focused on the topological properties of the renormalized lattice, and we have explicitly shown how the presence of second neighbors couplings (Zig-Zag ladder) leads to an asymptotic lattice equivalent to a simple chain, proving the irrelevance of the second neighbor couplings perturbation at the infinite disorder fixed point of the chain. On the other hand, the results of our numerical renormalization approach show that the inclusion of an arbitrary density of long range couplings in the chain modifies the scaling behavior of the lattice’s topology, and thus the associated critical behavior. These results stress the importance of determining the renormalized topological properties at any possible infinite disorder transition beyond the one-dimensional examples. In particular, the intermediate regime we have identified in our study of the Zig-Zag ladder opens the possibility of new infinite disorder scenarii for models with correlated long-range couplings. A natural extension of the present work would certainly focus on random algebraic interactions and the effect of the dimension, possibly relevant to the understanding of the dipolar glass LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> in a transverse field Aeppli1987 .
D. Carpentier and P. Pujol would like to acknowledge F. Leonforte for collaboration in a preliminary investigation related to the present work.
|
warning/0506/nucl-th0506088.html
|
ar5iv
|
text
|
# Multiplicity distribution and source deformation in full-overlap U+U collisions
## Abstract
We present a full Monte Carlo simulation of the multiplicity and eccentricity distributions in U+U collisions at $`\sqrt{s}=200A`$ GeV. While unavoidable trigger inefficiencies in selecting full-overlap U+U collisions cause significant modifications of the multiplicity distribution shown in Heinz:2005 , a selection of source eccentricities by cutting the multiplicity distribution is still possible.
In Heinz:2005 we advocated the use of full-overlap collisions between deformed uranium nuclei to probe open questions at the Relativistic Heavy Ion Collider (RHIC). Specifically, such collisions can be used to test the hydrodynamic behavior of elliptic flow to much higher energy densities than currently possible with non-central Au+Au collisions, and the large and strongly deformed reaction zones produced in such collisions will allow for a detailed examination of the path length dependence of the energy lost by a fast parton as it travels through the plasma created in the collision.
The calculations presented in Heinz:2005 were based on the assumption that full-overlap U+U collisions can be efficiently triggered on by using the zero degree calorimeters (ZDCs) of the RHIC experiments to discriminate against collisions with spectator nucleons flying down the beam pipe. By cutting the multiplicity distribution for the thus selected full-overlap collisions one can further select subevent classes with different spatial deformations of the created fireballs. In this short note we explore these assumptions in more quantitative detail, using a Monte Carlo simulation of the distributions of spectator nucleons and charged particle multiplicity for U+U collisions for arbitrary impact parameter and relative orientation between the two deformed uranium nuclei.
Our calculations are based on a Glauber model parametrization of the initial entropy production in these collisions, with standard Woods-Saxon form for the density distributions of the colliding nuclei (see KSH00 for details). \[Possible modifications arising from Color Glass Condensate (CGC) CGC initial conditions Hirano:2004rs will be shortly discussed at the end fn2 .\] The initial entropy density in the transverse plane at $`z=\mathrm{\hspace{0.17em}0}`$ (with $`z`$ denoting the beam direction) is determined by a combination of terms proportional to the wounded nucleon ($`n_{\mathrm{wn}}`$) and binary collision ($`n_{\mathrm{bc}}`$) densities:
$$s(𝒓_{};\mathrm{\Phi })=\kappa _s\left[\alpha n_{\mathrm{wn}}(𝒓_{};\mathrm{\Phi })+(1\alpha )n_{\mathrm{bc}}(𝒓_{};\mathrm{\Phi })\right].$$
(1)
Here $`\mathrm{\Phi }`$ is the angle between the beam direction and symmetry axis of one of the two U nuclei; in full-overlap collisions the symmetry axis of the other U nucleus lies in the same plane and forms an angle of $`\mathrm{\Phi }`$ or $`\pi \mathrm{\Phi }`$ with the beam axis. The normalization constant $`\kappa _s`$ in (1) is adjusted to reproduce the charged particle multiplicity density $`dN_{\mathrm{ch}}/dy`$ measured at midrapidity in central 200 $`A`$ GeV Au+Au collisions at RHIC PHOBOSv2 , assuming proportionality of $`dN_{\mathrm{ch}}/dy`$ with
the total entropy produced in the transverse plane. The wounded nucleon scaling fraction is tuned to $`\alpha =0.75`$ QGP3 to reproduce the centrality dependence of $`dN_{\mathrm{ch}}/dy`$ (see Figure 1). After fitting $`\kappa _s`$ and $`\alpha `$ to the Au+Au data, we use the same parameters to predict the multiplicities for U+U collisions at the same $`\sqrt{s}`$. The results for full-overlap U+U collisions are shown in Figure 1 by the dashed line; low (high) multiplicities correspond to side-on-side (tip-on-tip) collisions as indicated, due to their smaller (larger) binary collision contribution.
To determine the multiplicity distribution, we introduce Gaussian event-by-event fluctuations of the multiplicity $`ndN_{\mathrm{ch}}/dy`$ via Kharzeev:2000ph
$$\frac{dP}{dnd\mathrm{\Phi }}=A\mathrm{exp}\left\{\frac{(n\overline{n}(\mathrm{\Phi }))^2}{2a\overline{n}(\mathrm{\Phi })}\right\},$$
(2)
where $`\overline{n}(\mathrm{\Phi })`$ is the average charged particle multiplicity computed from Eq. (1) in a U+U collision with orientation angle $`\mathrm{\Phi }`$, and a width of $`a=0.6`$ has been shown to yield good agreement with PHOBOS data Kharzeev:2000ph . The multiplicity distribution is then obtained by integrating (2) over $`\mathrm{\Phi }`$. The resulting distribution Heinz:2005 is shown by the gray line in Figure 3 below. Its double hump structure results from the Jacobian $`d\overline{n}/d\mathrm{\Phi }`$, and its asymmetry is a consequence of the fluctuation width being proportional to the mean multiplicity $`\overline{n}`$. Note that the non-linear dependence of the charged multiplicity on the number of participant (wounded) nucleons, arising from the binary collision component in our parametrization (1), leads to a $`15\%`$ variation of the charged particle multiplicity among full-overlap U+U collisions as the relative orientation of their symmetry axes is varied over the accessible range.
The calculations of multiplicity and eccentricity distributions presented in Heinz:2005 rely on the assumption that, by monitoring spectator neutrons in the backward and forward zero degree calorimeters (ZDCs) of the RHIC experiments, full-overlap collisions can be perfectly distinguished from those collisions where the two nuclei are slightly misaligned. This is impossible in practice since even fully aligned collisions in general have a small number of spectator nucleons, arising from the dilute nuclear surface, and this number is larger for side-on-side than for tip-on-tip collisions (see Figure 1). Therefore, slightly misaligned tip-on-tip and fully aligned side-on-side collisions can have the same $`N_{\mathrm{part}}`$ and the same ZDC signal. To assess the contamination from collisions with imperfect overlap on the multiplicity and eccentricity distributions requires a more comprehensive study which includes non-central U+U collisions. This is the point of this short note.
A general U+U collision is parametrized by 5 parameters, the impact parameter $`b`$ and two Euler angles $`\mathrm{\Omega }=(\mathrm{\Phi },\beta )`$ for each nucleus describing the orientation of its symmetry axis relative to the beam axis and impact parameter direction. Equation (1) for the initial entropy density must therefore be generalized to
$$\begin{array}{c}s(𝒓_{};b,\mathrm{\Omega }_1,\mathrm{\Omega }_2)=\kappa _s[\alpha n_{\mathrm{wn}}(𝒓_{};b,\mathrm{\Omega }_1,\mathrm{\Omega }_2)\hfill \\ \hfill +(1\alpha )n_{\mathrm{bc}}(𝒓_{};b,\mathrm{\Omega }_1,\mathrm{\Omega }_2)].\end{array}$$
(3)
The multiplicity distribution is then calculated from
$$\frac{dP}{dn}=A^{}b𝑑bd^2\mathrm{\Omega }_1d^2\mathrm{\Omega }_2\mathrm{exp}\left\{\frac{\left(n\overline{n}(b,\mathrm{\Omega }_1,\mathrm{\Omega }_2)\right)^2}{2a\overline{n}(b,\mathrm{\Omega }_1,\mathrm{\Omega }_2)}\right\}.$$
(4)
Evaluating this 5-dimensional integral by Monte Carlo integration, we obtain the multiplicity distribution shown in Figure 2. Its right-most part contains the full-overlap collisions.
We can now try to select the latter from the overall event population by placing stringent cuts on the number of spectators ($`=2\times 238N_{\mathrm{part}}`$). The distribution of the number of spectators is shown in Figure 5 below. In Figure 3 we show the multiplicity distributions associated with the 5% and 0.5%, respectively, of events with the lowest spectator counts fn1 . It is immediately obvious that contamination from slightly misaligned collisions is sufficient to completely destroy the double-hump structure of the ideal full-overlap case, replacing it with a single peak.
By selecting low-spectator events, we bias the sample towards events with $`b0`$, $`\mathrm{\Phi }_{1,2}0`$, and the symmetry axes of the two nuclei approximately parallel. This suppresses the contribution from side-on-side configurations under the left peak of the idealized double-hump structure. At the same time, slightly misaligned tip-on-tip collisions fill in the dip between the two humps from the idealized case fn2 . The result is a single-peaked multiplicity distribution whose center moves left (towards lower multiplicities) as the cut on the number of spectator nucleons is loosened.
Nevertheless, for sufficiently tight spectator cuts, we still expect the collision events corresponding to the left edge of the multiplicity distributions shown in Figure 3 to have a larger contribution from strongly deformed side-on-side collisions than the events from the right edge (which will be mostly tip-on-tip collisions with small or zero source eccentricity). Following
our previous suggestion Heinz:2005 to select source eccentricities by cutting the multiplicity distribution of “zero spectator” collisions, we perform such cuts on the more realistic distributions shown in Figure 3. Figure 4 shows that it still possible in this way to select event classes with a given average source eccentricity: By taking the 0.5% of events with the lowest spectator count from Figure 3 (solid histogram) and cutting once more on the 5% of events with the lowest multiplicity, we obtain the eccentricity distribution shown by the black histogram in the bottom panel of Figure 4. This event class has an average source deformation $`ϵ_x`$ of about 18%, corresponding to Au+Au collisions with impact parameters around 5.5 fm. On the other hand, taking the same 0.5% spectator cut and selecting the 5% events with the largest multiplicities we obtain for the eccentricity distribution the gray histogram in the bottom Figure 4; this distribution peaks at $`ϵ_x=0`$ and has a very small average spatial deformation.
If one loosens the spectator cut to 5% instead of 0.5% (dotted histogram in Figure 3) and performs the same multiplicity selections (5% lowest or largest multiplicities, respectively), one obtains the eccentricity distributions shown in the top panel of Figure 4. Clearly, these distributions are much broader than with the tighter spectator cut, and the average eccentricities shift down from 17.7% to 14.2% for the low-multiplicity selection and up from 2.2% to 4.3% for the high-multiplicity selection. Note that, since the looser spectator cut allows for an increased contribution from non-zero impact parameters, the eccentricity of the nuclear overlap region can actually exceed the $`\mathrm{\hspace{0.17em}25}\%`$ ground state deformation of the single-uranium density distribution projected on the transverse plane. This gives rise to the right tail of the black histogram in the top panel of Figure 4. A typical event from this tail is shown in Figure 5. One sees that the 5% spectator cut allows for sizeable nonzero impact parameters and numbers of spectator nucleons, and that very tight ZDC cuts are required to ensure almost full overlap of the two uranium nuclei.
The detailed shapes of the eccentricity distributions shown in Figure 4 are expected to depend somewhat on our parametrization (1,3) of the initial transverse density distribution of the produced matter. It was shown in Ref. Hirano:2004rs that initial conditions motivated by the Color Glass Condensate picture of low-$`x`$ gluon saturation in large nuclei at high energies CGC produce transverse density distributions which fall off more steeply near the edge than the more Gaussian-like distributions Kolb:2001qz resulting from our Eqs. (1,3). This might result in somewhat larger eccentricities and narrower eccentricity distributions (i.e. better defined average eccentricities) than those shown in Figure 4. However, significantly higher statistics would likely be needed to clearly see such differences.
These results show that the suggestions made in Heinz:2005 for using U+U collisions to explore in more detail the ideal fluid dynamic nature of elliptic flow and the path length dependence of radiative parton energy loss are reasonably robust against trigger inefficiencies, and that a meaningful U+U collision program at RHIC is, in fact, feasible fn3 . Simultaneous strict cuts on small numbers of spectator nucleons and on charged particle multiplicity are necessary to select collisions which produce sources with well-defined and large spatial deformation; the histograms shown in the bottom panel of Figure 4 correspond to only 0.025% of all U+U collisions taking place in the accelerator. The top panel in Figure 4 shows that it is possible to loosen these tight cuts somewhat, at the expense of reducing the average spatial source deformation and introducing larger event-by-event fluctuations as well as an increased sensitivity to details of the Glauber model used for relating the relative nuclear orientation to the observed spectator nucleon and charged hadron multiplicities. We leave a further discussion of such model uncertainties (see also fn2 ) for later when U+U collisions become (hopefully) available.
We thank M. Gyulassy for stimulating discussions and G. Fai for an enlightening question which led to the discovery of an error in Eq. (4) in the originally submitted manuscript. This work was supported by the U.S. Department of Energy under contract DE-FG02-01ER41190.
|
warning/0506/gr-qc0506099.html
|
ar5iv
|
text
|
# N=1 de Sitter Supersymmetry Algebra
## 1 Introduction
Recent astrophysical data received from type Ia supernovas indicate that our universe might currently be in a de Sitter (dS) phase. Therefore it is important to find a formulation of de Sitter quantum field theory with the same level of completeness and rigor as for its Minkowskian counterpart. Some questions, however, are usually put forth for the non-existence of supersymmetry models with a positive cosmological constant, i.e. supersymmetry in de Sitter space. Such arguments are often based on the non-existence of Majorana spinors for $`O(4,1)`$ . Pilch et al have shown that if for every spinor, its independent charge-conjugate could be defined, de Sitter supergravity can be established with even $`N`$ .
Bros et al. presented a QFT of scalar free field in de Sitter space that closely mimics QFT in Minkowski space. We have generalized the Bros’s quantization of field, to quantization of fields with various spins in de Sitter space . In continuation of previous works where the charge-conjugate spinor had been defined , the supersymmetry in the ambient space notation has been studied in the present paper. Section two has been devoted to the discussion of the de Sitter group $`SO(1,4)`$, i.e. space-time symmetry of de Sitter space, and its universal covering group $`Sp(2,2)`$. Recalling the transformation properties of the spinor fields $`\psi (x),\overline{\psi }(x)`$, the charge conjugation symmetry of the de Sitter spinor field in ambient space notation is discussed in section $`3`$. The general de Sitter superalgebra is presented in section $`4`$. It is shown that a novel dS-superalgebra can be attained by the use of the spinor fields and charge conjugation in the ambient space notation. Finally, a brief conclusion and an outlook have been given in section $`5`$. To illustrate the novel algebra, the generalized Jacobi identities are calculated in appendix A.
## 2 de Sitter group
The de Sitter space is an elementary solution of the positive cosmological Einstein equation in the vacuum. It is conveniently seen as a hyperboloid embedded in a five-dimensional Minkowski space
$$X_H=\{x\mathrm{IR}^5;x^2=\eta _{\alpha \beta }x^\alpha x^\beta =H^2\},\alpha ,\beta =0,1,2,3,4,$$
(1)
where $`\eta _{\alpha \beta }=`$diag$`(1,1,1,1,1)`$. The de Sitter metrics reads
$$ds^2=\eta _{\alpha \beta }dx^\alpha dx^\beta _{x^2=H^2}=g_{\mu \nu }^{dS}dX^\mu dX^\nu ,\mu =0,1,2,3,$$
where the $`X^\mu `$’s are the $`4`$ space-time intrinsic coordinates on dS hyperboloid. Different coordinate systems can be chosen for $`X^\mu `$. A $`10`$-parameter group $`SO_0(1,4)`$ is the kinematical group of the de Sitter universe. In the limit $`H=0`$, this reduces to the Poincaré group. There are two Casimir operators
$$Q^{(1)}=\frac{1}{2}L_{\alpha \beta }L^{\alpha \beta },$$
$$Q^{(2)}=W_\alpha W^\alpha ,W_\alpha =\frac{1}{8}ϵ_{\alpha \beta \gamma \delta \eta }L^{\beta \gamma }L^{\delta \eta },$$
(2)
where $`ϵ_{\alpha \beta \gamma \delta \eta }`$ is the usual antisymmetrical tensor and the $`L_{\alpha \beta }`$’s are the infinitesimal generators, which obey the commutation relations
$$[L_{\alpha \beta },L_{\gamma \delta }]=i(\eta _{\alpha \gamma }L_{\beta \delta }+\eta _{\beta \delta }L_{\alpha \gamma }\eta _{\alpha \delta }L_{\beta \gamma }\eta _{\beta \gamma }L_{\alpha \delta }).$$
(3)
$`L_{\alpha \beta }`$ can be represented as $`L_{\alpha \beta }=M_{\alpha \beta }+S_{\alpha \beta }`$, where $`M_{\alpha \beta }=i(x_\alpha _\beta x_\beta _\alpha )`$ is the “orbital” part and $`S_{\alpha \beta }`$ is the “spinorial” part. The form of the $`S_{\alpha \beta }`$ depends on the spin of the field. For spin $`\frac{1}{2}`$ field, it can be defined as
$$S_{\alpha \beta }=\frac{i}{4}[\gamma _\alpha ,\gamma _\beta ],$$
(4)
where the five $`4\times 4`$ matrices $`\gamma ^\alpha `$ are the generators of the Clifford algebra based on the metric $`\eta _{\alpha \beta }`$:
$$\gamma ^\alpha \gamma ^\beta +\gamma ^\beta \gamma ^\alpha =2\eta ^{\alpha \beta }\text{I1},\gamma _{}^{\alpha }{}_{}{}^{}=\gamma ^0\gamma ^\alpha \gamma ^0.$$
(5)
An explicit and convenient representation is provided by
$$\gamma ^0=\left(\begin{array}{cccc}\text{I1}& 0\hfill & & \\ 0& \text{I1}\hfill & & \end{array}\right),\gamma ^4=\left(\begin{array}{cccc}0& \text{I1}\hfill & & \\ \text{I1}& 0\hfill & & \end{array}\right),$$
$$\gamma ^1=\left(\begin{array}{cccc}0& i\sigma ^1\hfill & & \\ i\sigma ^1& 0\hfill & & \end{array}\right),\gamma ^2=\left(\begin{array}{cccc}0& i\sigma ^2\hfill & & \\ i\sigma ^2& 0\hfill & & \end{array}\right),\gamma ^3=\left(\begin{array}{cccc}0& i\sigma ^3\hfill & & \\ i\sigma ^3& 0\hfill & & \end{array}\right),$$
(6)
where I1 is the unit $`2\times 2`$ matrix and $`\sigma ^i`$ are the Pauli matrices. This representation had been proved to be useful in analysis of the physical relevance of the group representation . In this representation,
$$\gamma _{}^{\alpha }{}_{}{}^{T}=\gamma ^4\gamma ^2\gamma ^\alpha \gamma ^2\gamma ^4.$$
The spinor wave equation in de Sitter space-time has been originally deduced by Dirac in 1935 , and can be obtained from the eigenvalue equation of the second order Casimir operator
$$(i\mathit{}\gamma .\overline{}+2i+\nu )\psi (x)=0,$$
(7)
where $`\mathit{}=\eta _{\alpha \beta }\gamma ^\alpha x^\beta `$ and $`\overline{_\alpha }=_\alpha +H^2x_\alpha x`$. Due to covariance of the de Sitter group, the adjoint spinor is defined as follows :
$$\overline{\psi }(x)\psi ^{}(x)\gamma ^0\gamma ^4.$$
(8)
Let us now recall the transformation properties of the spinor fields $`\psi (x)`$ and $`\overline{\psi }(x)`$. The two-fold, universal covering group of $`SO_0(1,4)`$, is the (pseudo-)symplectic group $`Sp(2,2)`$,
$$Sp(2,2)=\{g\mathrm{Mat}(2;\text{IHI}):detg=1,g^{}\gamma ^0g=\gamma ^0\},$$
(9)
where $`g^{}=^T\stackrel{~}{g}`$, $`{}_{}{}^{T}g`$ is the transposed of $`g`$ and $`\stackrel{~}{g}`$ the quaternionic conjugate of $`g`$. It should be noted that quaternionic $`𝒫`$ is
$$𝒫=(x^4,\stackrel{}{x})=x^4\text{I1}+ix^1\sigma ^1ix^2\sigma ^2+ix^3\sigma ^3,$$
(10)
where $`\sigma ^i`$ are the Pauli matrices and $`\stackrel{~}{𝒫}=(x^4,\stackrel{}{x})`$ is its conjugate. For obtaining the isomorphic relation between the two groups we define the matrices $`X`$ associated with $`xX_H`$ by:
$$X=\left(\begin{array}{cccc}x^0& 𝒫\hfill & & \\ \stackrel{~}{𝒫}& x^0\hfill & & \end{array}\right).$$
(11)
Through representation $`(6)`$ of the $`\gamma `$ matrices, $`X`$ can be written in the following form:
$$\mathit{}=x.\gamma =X\gamma ^0=\left(\begin{array}{cccc}x^0& 𝒫\hfill & & \\ \stackrel{~}{𝒫}& x^0\hfill & & \end{array}\right).$$
(12)
The transformation of $`X`$, under the action of the group $`Sp(2,2)`$ is
$$X^{}=gX\stackrel{~}{g}^t,\mathit{}^{}=g\mathit{}g^1.$$
(13)
For $`\mathrm{\Lambda }SO_0(1,4)`$ and $`gSp(2,2)`$ we have
$$x^\alpha =\eta ^{\alpha \beta }x_\beta ^{}=\frac{1}{4}tr(\gamma ^\alpha \gamma ^\beta )x_\beta ^{}=\frac{1}{4}tr(\gamma ^\alpha g\mathit{}g^1)$$
$$=\frac{1}{4}tr(\gamma ^\alpha g\gamma ^\beta g^1)x_\beta =\mathrm{\Lambda }^{\alpha \beta }(g)x_\beta .$$
(14)
For every $`gSp(2,2)`$, there exists a transformation $`\mathrm{\Lambda }SO_0(1,4)`$, which satisfies the following relations
$$\mathrm{\Lambda }_\beta ^\alpha (g)=\frac{1}{4}tr(\gamma ^\alpha g\gamma _\beta g^1),\mathrm{\Lambda }_\beta ^\alpha \gamma ^\beta =g\gamma ^\alpha g^1.$$
(15)
The isomorphic relation between the two groups is
$$SO_0(1,4)Sp(2,2)/\text{ZZ}_2.$$
(16)
The transformation laws for the $`\psi (x)`$ and its adjoint $`\overline{\psi }(x)`$, under which the de Sitter-Dirac equation is covariant, are :
$`\psi (x)\psi ^{}(x)=g^1\psi (\mathrm{\Lambda }(g)x),`$ (17)
$`\overline{\psi }(x)\overline{\psi }^{}(x)=\overline{\psi }(\mathrm{\Lambda }(g)x)i(g),`$ (18)
where $`i(g)\gamma ^4g\gamma ^4`$ . Similar to the Minkowskian space, it is useful to define $`g`$ by
$$g=\mathrm{exp}[\frac{i}{2}\omega ^{\alpha \beta }S_{\alpha \beta }],\omega ^{\alpha \beta }=\omega ^{\beta \alpha },$$
(19)
where $`\gamma ^0g^{}\gamma ^0=g^1`$, i.e. $`gSp(2,2).`$
## 3 Charge conjugation
Previously, the charge conjugation spinor $`\psi ^c`$ was calculated in the ambient space notation
$$\psi ^c=\eta _cC(\gamma ^4)^T(\overline{\psi })^T,$$
(20)
where $`\eta _c`$ is an arbitrary unobservable phase value, generally set to unity. In the present framework charge conjugation is an antilinear transformation. In the $`\gamma `$ representation $`(6)`$ we had obtained :
$$C\gamma ^0C^1=\gamma ^0,C\gamma ^4C^1=\gamma ^4$$
$$C\gamma ^1C^1=\gamma ^1,C\gamma ^3C^1=\gamma ^3,C\gamma ^2C^1=\gamma ^2.$$
(21)
In this representation, $`C`$ commutes with $`\gamma ^2`$, and anticommutes with other $`\gamma `$-matrices. Therefore the simple choice may be taken to be $`C=\gamma ^2`$, where the following relation is satisfied
$$C=C^1=C^T=C^{}.$$
(22)
This clearly illustrates the non-singularity of $`C`$.
The adjoint spinor, defined by $`\overline{\psi }(x)\psi ^{}(x)\gamma ^0\gamma ^4`$, transforms in a different way from $`\psi `$, under de Sitter transformation. It is easily shown that $`\psi ^c`$ transforms in the same way as $`\psi `$
$$\psi ^c(x^{})=g(\mathrm{\Lambda })\psi ^c(x).$$
The wave equation of $`\psi ^c`$ is different from the wave equation of $`\psi `$ by the sign of the ”charge” $`q`$ and $`\nu `$ . Thus it follows that if $`\psi `$ describes the motion of a dS-Dirac ”particle” with the charge $`q`$, $`\psi ^c`$ represents the motion of a dS-Dirac ”anti-particle” with the charge $`(q)`$. In other words $`\psi `$ and $`\psi ^c`$ can describe ”particle” and ”antiparticle” wave functions. $`\psi `$ and $`\psi ^c`$ are charge conjugation of each other
$$(\psi ^c)^c=C\gamma ^0\psi _{}^{c}{}_{}{}^{}=C\gamma ^0(C\gamma ^0\psi )=\psi .$$
(23)
## 4 N=1 Supersymmetry Algebra
Supersymmetry in de Sitter space has been studied by Pilch et al . Recently, supersymmetry has been investigated in constant curvature space as well . In this section we have presented the supersymmetry algebra in ambient space notation. It is shown that if the spinor field and the charge conjugation operators are defined in the ambient space notation, a novel de Sitter superalgebra can be attained.
In order to extend the de Sitter group, the generators of supersymmetry transformation $`Q_i^n`$ are introduced. Here $`i`$ is the spinor index ($`i=1,2,3,4`$) and $`n`$ is the supersymmetry index $`n=1,..,N`$. $`Q_i^n`$’s are superalgebra spinor generators which transform as de Sitter group spinors. The generators $`\stackrel{~}{Q}_i^n`$ are defined by
$$\stackrel{~}{Q}_i=\left(Q^T\gamma ^4C\right)_i=\overline{Q^c}_i,$$
(24)
where $`Q^T`$ is the transpose of $`Q`$. It can be shown that $`\stackrel{~}{Q}\gamma ^4Q`$ is a scalar field under the de Sitter transformation .
For $`N1`$, closure of algebra requires extra bosonic generators. These do not necessarily commute with other generators and consequently are not central charges. They are internal symmetry generators . These generators, shown by $`T_{mn}`$, commute with de Sitter generators.
Therefore the de Sitter superalgebra in four-dimensional space-time has the following generators:
* $`M_{\alpha \beta }`$, the generators of de Sitter group,
* the internal group generators $`T_{nm}`$, that are defined by the additional condition
$$T_{nm}=T_{mn};n,m=1,2\mathrm{}.N,$$
* the $`4`$-component dS-Dirac spinor generators,
$$Q_i^n,i=1,2,3,4,n=1,2,\mathrm{},N.$$
To every generator $`A`$, a grade $`p_a`$ is assigned. For the fermionic generator $`p_a=1`$, and for the bosonic generator $`p_a=0`$. The bilinear product operator is defined by
$$[A,B]=(1)^{p_a.p_b}[B,A].$$
(25)
The generalized Jacobi identities is
$$(1)^{p_a.p_c}[A,[B,C]]+(1)^{p_c.p_b}[C,[A,B]]+(1)^{p_b.p_a}[B,[C,A]]=0.$$
(26)
Using different generalized Jacobi identities, similar to the method presented by Pilch el al. , the full dS-superalgebra can be written in the following form:
$$[M_{\alpha \beta },M_{\gamma \delta }]=i(\eta _{\alpha \gamma }M_{\beta \delta }+\eta _{\beta \delta }M_{\alpha \gamma }\eta _{\alpha \delta }M_{\beta \gamma }\eta _{\beta \gamma }M_{\alpha \delta }),$$
$$[T_{rl},T_{pm}]=i(\omega _{rp}T_{lm}+\omega _{lm}T_{rp}\omega _{rm}T_{lp}\omega _{lp}T_{rm}),$$
$$[M_{\alpha \beta },T_{rl}]=0,$$
$$[Q_i^r,M_{\alpha \beta }]=(S_{\alpha \beta }Q^r)_i,[\stackrel{~}{Q}_i^r,M_{\alpha \beta }]=(\stackrel{~}{Q}^rS_{\alpha \beta })_i$$
$$[Q_i^r,T^{lp}]=(\omega ^{rl}Q_i^p\omega ^{rp}Q_i^l)$$
$$\{Q_i^r,Q_j^l\}=\omega ^{rl}\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right)_{ij}M_{\alpha \beta }+\left(\gamma ^4\gamma ^2\right)_{ij}T^{rl},$$
where $`S_{\alpha \beta }`$ is defined by $`(4)`$. The following relations are used to determine the above structure
$$\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right)^T=\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right),\left(\gamma ^\alpha \gamma ^4\gamma ^2\right)^T=\gamma ^\alpha \gamma ^4\gamma ^2,\left(\gamma ^4\gamma ^2\right)^T=\gamma ^4\gamma ^2.$$
(27)
It is necessary to obtain matrix $`\omega `$, which determines the structure of the internal group. For even $`N`$ de Sitter supersymmetry algebra, studied by Pilch et al , the matrix $`\omega `$ has been obtained explicitly. A new dS-superalgebra, defined in the ambient space notation, is introduced in this stage for $`N=1`$ case. In this case, $`T_{11}=0,\omega =1`$ and the simple de Sitter supersymmetry algebra is defined by the following relations:
$$[M_{\alpha \beta },M_{\gamma \delta }]=i(\eta _{\alpha \gamma }M_{\beta \delta }+\eta _{\beta \delta }M_{\alpha \gamma }\eta _{\alpha \delta }M_{\beta \gamma }\eta _{\beta \gamma }M_{\alpha \delta }),$$
(28)
$$\{Q_i,Q_j\}=\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right)_{ij}M_{\alpha \beta },$$
(29)
$$[Q_i,M_{\alpha \beta }]=(S_{\alpha \beta }Q)_i,[\stackrel{~}{Q}_i,M_{\alpha \beta }]=(\stackrel{~}{Q}S_{\alpha \beta })_i.$$
(30)
This can be proved by the use of generalized Jacobi identities (appendix).
Finally we present an explicit form of the supersymmetric generators which satisfy the above relations. We consider a superspace with bosonic coordinates $`x^\alpha `$ and fermionic coordinates $`\theta _i`$ where $`\theta _i`$ is a four component de Sitter-Dirac Grassmann spinor in the ambient space notation. The suitable representation of these superalgebra generators in superspace are provided by
$$\{\begin{array}{cc}M_{\alpha \beta }=i(x_\alpha \overline{}_\beta x_\beta \overline{}_\alpha )+\frac{}{\theta }S_{\alpha \beta }\theta ,\hfill & \\ Q=\gamma .\overline{}\theta +\mathit{}\frac{}{\stackrel{~}{\theta }},\hfill & \end{array}$$
(31)
where $`\frac{}{\stackrel{~}{\theta }_i}=\left(\gamma ^2\gamma ^4\right)_{ik}\frac{}{\theta _k}`$. Using the equation $`\{\theta _i,\frac{}{\theta _j}\}=\delta _{ij}`$ and the following identities ,
$$(S^{\alpha \beta }\gamma ^4\gamma ^2)_{ij}(S_{\alpha \beta })_{kl}+\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right)_{jk}(S_{\alpha \beta })_{il}+\left(S^{\alpha \beta }\gamma ^4\gamma ^2\right)_{ki}(S_{\alpha \beta })_{jl}=0,$$
$$(S^{\alpha \beta })_{ij}(S_{\alpha \beta })_{kl}+(S^{\alpha \beta })_{il}(S_{\alpha \beta })_{kj}=(\gamma ^\alpha )_{ij}(\gamma _\alpha )_{kl}+(\gamma ^\alpha )_{il}(\gamma _\alpha )_{kj},$$
it is straightforward to prove the above de Sitter supersymmetry algebra.
## 5 Conclusions
The formalism of the quantum field in de Sitter universe, in ambient space notation, is very similar to the quantum field formalism in Minkowski space. In this paper we present the de Sitter supersymmetry algebra in this notation, which is independent of the choice of the coordinate system. In addition, a novel superalgebra $`(2830)`$ has been obtained, which do not fall into the categories considered in previous works . The importance of this formalism may be shown further by the consideration of the linear quantum gravity and supergravity in de Sitter space, which lays a firm ground for further study of the evolution of the universe.
Acknowledgements: The authors would like to extend their gratitude to R. Kallosh for the useful discussions and S. Moradi for interest in this work, and the referee for his constructive suggestions.
## Appendix A $`N=1`$ de Sitter superalgebra
The generalized Jacobi identities $`(M,M,Q)`$ and $`(M,Q,Q)`$ can be easily utilized to prove the de Sitter supersymmetry algebra $`(2830)`$. The generalized Jacobi identity $`(M,M,Q)`$ is
$$[[M_{\alpha \beta },M_{\gamma \delta }],Q_i]+[[Q_i,M_{\alpha \beta }],M_{\gamma \delta }]+[[M_{\gamma \delta },Q_i],M_{\alpha \beta }]=0.$$
(32)
Using the equations $`(28)`$ and $`(30)`$, the following relations can be obtained:
$$[[M_{\alpha \beta },M_{\gamma \delta }],Q_i]=i\eta _{\alpha \gamma }[M_{\beta \delta },Q_i]i\eta _{\beta \delta }[M_{\alpha \gamma },Q_i]+i\eta _{\alpha \delta }[M_{\beta \gamma },Q_i]+i\eta _{\beta \gamma }[M_{\alpha \delta },Q_i],$$
$$[[Q_i,M_{\alpha \beta }],M_{\gamma \delta }]=[(S_{\alpha \beta })_i^jQ_j,M_{\gamma \delta }]=(S_{\alpha \beta })_i^j(S_{\gamma \delta })_j^kQ_k,$$
$$[[M_{\gamma \delta },Q_i],M_{\alpha \beta }]=[(S_{\gamma \delta })_i^jQ_j,M_{\alpha \beta }]=(S_{\gamma \delta })_i^j(S_{\alpha \beta })_j^kQ_k.$$
Implementing the above equations in the eq $`(32)`$,
$$i\eta _{\alpha \gamma }[M_{\beta \delta },Q_i]i\eta _{\beta \delta }[M_{\alpha \gamma },Q_i]+i\eta _{\alpha \delta }[M_{\beta \gamma },Q_i]+i\eta _{\beta \gamma }[M_{\alpha \delta },Q_i]+(S_{\alpha \beta })_i^j(S_{\gamma \delta })_j^kQ_k(S_{\gamma \delta })_i^j(S_{\alpha \beta })_j^kQ_k=0,$$
$$i\eta _{\alpha \gamma }(S_{\beta \delta })_i^jQ_j+i\eta _{\beta \delta }(S_{\alpha \gamma })_i^jQ_ji\eta _{\alpha \delta }(S_{\beta \gamma })_i^jQ_ji\eta _{\beta \gamma }(S_{\alpha \delta })_i^jQ_j+[(S_{\alpha \beta })_i^j(S_{\gamma \delta })_j^k(S_{\gamma \delta })_i^j(S_{\alpha \beta })_j^k]Q_k=0,$$
$$i[\eta _{\alpha \gamma }(S_{\beta \delta }Q)_i+\eta _{\beta \delta }(S_{\alpha \gamma }Q)_i\eta _{\alpha \delta }(S_{\beta \gamma }Q)_i\eta _{\beta \gamma }(S_{\alpha \delta }Q)_i]+[(S_{\alpha \beta }S_{\gamma \delta }S_{\gamma \delta }S_{\alpha \beta })Q]_i=0,$$
$$i[(\eta _{\alpha \gamma }S_{\beta \delta }+\eta _{\beta \delta }S_{\alpha \gamma }\eta _{\alpha \delta }S_{\beta \gamma }\eta _{\beta \gamma }S_{\alpha \delta })Q]_i+[(S_{\alpha \beta }S_{\gamma \delta }S_{\gamma \delta }S_{\alpha \beta })Q]_i=0,$$
it is shown that $`S_{\alpha \beta }`$’s satisfy the following commutation relation
$$[S_{\alpha \beta },S_{\gamma \delta }]=i(\eta _{\alpha \gamma }S_{\beta \delta }+\eta _{\beta \delta }S_{\alpha \gamma }\eta _{\alpha \delta }S_{\beta \gamma }\eta _{\beta \gamma }S_{\alpha \delta }).$$
This is none other than equation $`(3`$). This generalized Jacobi identity verifies the relation $`(30)`$, i.e. commutation relation of $`Q`$ and $`M`$.
The generalized Jacobi identity $`(M,Q,Q)`$ is
$$[\{Q_i,Q_j\},M_{\alpha \beta }]+\{[M_{\alpha \beta },Q_i],Q_j\}\{[Q_j,M_{\alpha \beta }],Q_i\}=0.$$
(33)
By the use of equations $`(29)`$ and $`(30)`$, we obtain
$$\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}[M_{\gamma \delta },M_{\alpha \beta }]\left(S_{\alpha \beta }\right)_{ik}\{Q_k,Q_j\}\left(S_{\alpha \beta }\right)_{ik}\{Q_k,Q_i\}=0,$$
$$\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}[M_{\gamma \delta },M_{\alpha \beta }]\left(S_{\alpha \beta }\right)_{ik}\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{kj}M_{\gamma \delta }\left(S_{\alpha \beta }\right)_{jk}\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ki}M_{\gamma \delta }=0,$$
$$\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}[M_{\gamma \delta },M_{\alpha \beta }]\left(S_{\alpha \beta }S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}M_{\gamma \delta }\left(S_{\alpha \beta }S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ji}M_{\gamma \delta }=0.$$
Utilizing equation $`(27)`$, we obtain
$$\left(S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}[M_{\gamma \delta },M_{\alpha \beta }]\left(S_{\alpha \beta }S^{\gamma \delta }\gamma ^4\gamma ^2\right)_{ij}M_{\gamma \delta }\left(S^{\gamma \delta }S_{\alpha \beta }\gamma ^2\gamma ^4\right)_{ij}M_{\gamma \delta }=0.,$$
$$S^{\gamma \delta }[M_{\gamma \delta },M_{\alpha \beta }][S_{\alpha \beta },S^{\gamma \delta }]M_{\gamma \delta }=0.$$
This generalized Jacobi identity once again, verifies the anti-commutation relation of $`Q_i`$ and $`Q_j`$.
|
warning/0506/physics0506020.html
|
ar5iv
|
text
|
# Permutation sampling in Path Integral Monte Carlo
## I Introduction
The Path Integral Monte Carlo method is arguably the most powerful numerical technique to calculate thermodynamic properties of continuum (i.e., non-lattice) quantum many-body systems at finite temperature.ceperley95 For Bose systems, it is the only known, generally applicable theoretical method essentially free from approximations. Numerical estimates yielded by PIMC are affected by a statistical error, as well as by systematic errors, due to the finite size of the simulated system and to imaginary time discretization. In most cases, however, the computational resources typically available nowadays allow one to render the size of all of these errors insignificant in practice.
The most notable application of PIMC to date, is the study of the superfluid (SF) transition in liquid <sup>4</sup>He by Ceperley and Pollock,ceperley86 whose results have become the standard reference for all theoretical calculations on SF helium; but numerous applications to other Bose systems have been reported in the literature, over the past two decades. No general formulation exists as yet of PIMC (nor of any other Quantum Monte Carlo method), capable of overcoming the sign problem, that has so far made it impossible to obtain equally high quality results for Fermi systems. Even for fermions, however, PIMC proves a valid option, allowing one to obtain approximate estimates, of accuracy at least comparable to that afforded by other methods.ceperley96
Physical effects of interest in quantum many-body systems are almost invariably associated with quantum statistics; for example, superfluidity in <sup>4</sup>He is intimately connected to long exchange cycles of helium atoms. Because a direct summation of all $`N!`$ permutations of $`N`$ indistinguishable particles is unfeasible, except for very small values of $`N`$, within PIMC quantum statistics is included by performing a statistical sampling of permutations. Thus, an all-important ingredient of any practical implementation of PIMC is an efficient procedure to carry out such a sampling.
Since the pioneering study of Ref. ceperley86, , there has been relatively little experimentation with implementations of PIMC differing, in some of the more important aspects, from the one described in Ref. ceperley95, , henceforth referred to as CP. The CP implementation has come to be regarded as “canonical”, especially when studying quantum many-body systems in the highly degenerate regime (i.e. at low temperature). It is based on an accurate (“pair-product”) high-temperature density matrix, allowing one to observe convergence of the physical estimates with a relatively low number of imaginary time “slices” (of the order of 40 for superfluid <sup>4</sup>He at a temperature $`T`$=1 K). Two slightly different procedures have been proposed and utilized, in the context of CP, to perform the sampling of permutations, both of which are thoroughly described in Ref. ceperley95, . To our knowledge, no systematic, quantitative assessment of the relative merits and advantages of these two sampling strategies has yet been offered; it is also unclear to what extent their effectiveness and applicability are problem-dependent, and/or hinge on the use of the above-mentioned high-temperature density matrix.
In this work, we illustrate a new method to sample permutations of indistinguishable particles in PIMC simulations. It bears some similarities with one of the two strategies described in Ref. ceperley95, , but differs from it in some important technical aspects. We also deem it easier to implement, and may be potentially more efficient, even though, naturally, this speculation will need to be supported by systematic comparisons with the other existing options.
As an illustrative application of our sampling method, we have carried out a PIMC simulation of liquid <sup>4</sup>He in the SF regime, i.e., we have repeated the original calculation of Ref. ceperley86, . SF helium is the accepted test bench for quantum Monte Carlo calculations, since it is the most extensively studied quantum fluid, for which effects of quantum statistics manifest themselves at the macroscopic level. In order to make the test more significant, and help expose any deficiency or merit of the permutation sampling procedure, we have not utilized the pair-product high-temperature density matrix; rather we opted for a much simpler form, which requires a substantially larger number of imaginary time slices, in order to observe convergence of the estimates. Besides providing results for energetic properties, known to be affected quantitatively by Bose statistics, and which we compare to those of Ref. ceperley86, , we have also attempted to furnish here some quantitative information, which should help assess the efficiency of our method in sampling the space of all possible entangled many-particle paths (i.e., including permutations). It is our hope that this will provide a baseline for future, more extensive comparisons of different approaches. Somewhat interestingly, our results show that the PIMC simulation of SF <sup>4</sup>He is feasible, albeit at a higher computational cost, even with a relatively simple PIMC implementation. Doubtless, this is also in part due to advances in computing hardware, which enable what may have been prohibitive two decades ago, when the first such simulation was carried out.
The remainder of this manuscript is organized as follows: in the next section, we provide a detailed description of our computational methodology. In the following sections, we illustrate our results, and outline our conclusions and outlook.
## II METHODOLOGY
The PIMC methodology is fairly mature, and extensively described in Ref. ceperley95, , to which we refer readers interested in a thorough, comprehensive illustration. Our specific implementation is largely based on the ideas and methods presented therein. Nevertheless, at the risk of some redundancy, we provide a somewhat detailed description of our implementation here. This will hopefully facilitate the task of others who may wish to repeat our study and/or experiment with our algorithm to sample permutations.
Consider a quantum many-body system of $`N`$ identical, point-like particles of mass $`m`$, described by the following Hamiltonian:
$$\widehat{H}=\lambda \underset{i=1}{\overset{N}{}}_i^2+\underset{i<j}{}V(|𝐫_i𝐫_j|)$$
(1)
where $`\lambda =\mathrm{}^2/2m`$. Implicit in the above model is the assumption that the interactions among particles can be accurately represented by a pairwise, central potential (the $`V`$ term in (1)). Although this is obviously an approximation, it is commonly made in theoretical studies of most quantum fluids. In any case, it is not a requirement for the applicability of PIMC, nor of our specific implementation. In the following, it is assumed that particles in the system obey Bose statistics.noteq The system is assumed to be enclosed in a vessel, shaped as a parallelepiped, with periodic boundary conditions in all directions, and to be held in thermal equilibrium at a temperature $`T`$.
The thermodynamic average of a physical quantity formally represented by an operator $`\widehat{O}`$ (for simplicity assumed diagonal in the coordinate representation) is expressed as follows:
$`O={\displaystyle \frac{1}{Z}}{\displaystyle 𝑑RO(R)\rho (R,R,\beta )}`$ (2)
where $`\beta `$=1/$`T`$ (we work with units where the Boltzmann constant $`k_B`$=1), and $`R\{𝐫_1,𝐫_2,\mathrm{}𝐫_N\}`$, is a configuration of the system, specified by the positions of all the $`N`$ particles. In Eq. (2), $`\rho (R,R^{},\beta )R|e^{\beta \widehat{H}}|R^{}`$ is the many-body density matrix, and $`Z=𝑑R\rho (R,R,\beta )`$ is the partition function.
A Monte Carlo evaluation of (2) consists of generating a large set of random many-particle configurations $`\{R_p\}`$, $`p=1,\mathrm{},M`$, statistically sampled from a probability density proportional to $`\rho (R,R,\beta )`$; the thermal average (2) can thus be estimated as a statistical average over the set of values $`\{O(R_p)\}`$.
An explicit expression for the density matrix $`\rho (R,R,\beta )`$ is unavailable for any non-trivial many-body system; however, one can still generate the set $`\{R_p\}`$, by sampling discrete many-particle paths $`X_p`$ through configuration space, i.e.,
$$X_p\{R_{0p},R_{1p},R_{2p},\mathrm{}R_{Lp}\}$$
(3)
Paths are formally defined in the imaginary time interval $`0\tau \beta `$, i.e., $`R_jR(j\delta \tau )`$, with $`L\delta \tau =\beta `$, and are randomly drawn from a probability distribution $`\overline{\rho }(X)`$ given by
$$\overline{\rho }(X)\overline{\rho }(R_0,R_1,R_2,\mathrm{}R_L)=\underset{j=0}{\overset{L1}{}}\rho _{}(R_j,R_{j+1},\delta \tau )$$
(4)
where $`\rho _{}`$ is an (analytically known) approximation to the true many-body density matrix, constructed to be asymptotically exact in the “high temperature” $`\delta \tau 0`$ limit. It can be shown that in that limit ($`L\mathrm{}`$), each configuration $`R_p`$ visited by paths is statistically sampled from a distribution proportional to $`\rho (R,R,\beta )`$.
The configuration $`R(\beta )`$, i.e., that corresponding to the end of the path in imaginary time, must coincide with $`R(0)`$, except for a permutation $`P`$ of the particle labels ($`1`$ through $`N`$). The possibility of permutations of particles must be allowed, in order to incorporate in the calculation the effects of particle indistinguishability and Bose statistics. Consequently, although many-particle paths are periodic in imaginary time, i.e., the configuration $`R(j\delta \tau +q\beta )`$ ($`q`$ being an arbitrary integer) is identical with $`R(j\delta \tau )`$ (in that particles occupy identical positions), individual particle labels can be different. Stated differently, single particle paths $`𝐫_1(j\delta \tau )\mathrm{}𝐫_N(j\delta \tau )`$ can become “entangled”, as a result of permutations.
Permutations normally become important at sufficiently low temperature; at high temperature, only the identity permutation contributes significantly to thermal averages. At low temperature, however, permutations underlie phenomena such as superfluidity and Bose Condensation in liquid helium and, presumably, in all other superfluids.
In an actual calculation implementing the above computational scheme, one must necessarily work with a finite value of $`L`$; in principle, one ought to carry out the $`L\mathrm{}`$ limit by extrapolating numerical results obtained with different values of $`L`$. In practice, this procedure proves quite cumbersome, especially when one is interested in many thermodynamic points. Thus, one typically performs all calculations (at a given temperature) with a single value of $`L`$, chosen sufficiently large so that estimates may be expected to coincide with the extrapolated values, within some small tolerance. For reasons of efficiency, it is desirable that such “optimal” value of $`L`$ not be too large (a few hundred slices at the most); thus, it is advantageous to work with as accurate as possible a “high-temperature density matrix” $`\rho _{}`$, which will allow one to achieve convergence of the numerical estimates without resorting to impractically large values of $`L`$.
The importance of this issue was demonstrated by Ceperley and Pollock,ceperley86 who proposed the following form for $`\rho _{}`$:
$$\rho _{}(R,R^{},\delta \tau )=A_F(R,R^{},\delta \tau )\left\{\underset{i<j}{}\mathrm{exp}\left[u(𝐫_{ij},𝐫_{ij}^{},\delta \tau )\right]\right\}$$
(5)
where $`𝐫_{ij}=𝐫_j𝐫_i`$, $`A_F(R,R^{},\delta \tau )=_{i=1}^N\rho _F(𝐫_i,𝐫_i^{},\delta \tau )`$, is the exact density matrix of a system of $`N`$ distinguishable, non-interacting particles, with
$$\rho _F(𝐫_i,𝐫_i^{},\delta \tau )=\left(\sqrt{4\pi \lambda \delta \tau }\right)^{3/2}\mathrm{exp}\left[\frac{(𝐫_i𝐫_i^{})^2}{4\lambda \delta \tau }\right]$$
(6)
and where $`u`$ is obtained by imposing that $`\rho _{}`$ be the exact density matrix for a system of two interacting particles. For PIMC calculations of highly quantal, hard-sphere-like systems such as condensed Helium, the form (5) for $`\rho _{}`$ affords a tremendous increase in efficiency, with respect to other, simpler forms for $`\rho _{}`$ (such as the so-called primitive approximation; for details, see Ref. ceperley95, ).
In this work, we have not made use of the high-temperature density matrix (5), choosing instead the following form:
$$\rho _{}(R_j,R_{j+1},\delta \tau )=A_F(R_j,R_{j+1},\delta \tau )\mathrm{exp}\left[\delta \tau U(R_j)\right]$$
(7)
where
$$U(R_j)=\frac{2V(R_j)}{3}+\stackrel{~}{V}(R_j)$$
(8)
$`V(R)_{i<j}V(|𝐫_i𝐫_j|)`$ being the total potential energy of the system in the configuration $`R_j`$, and
$`\stackrel{~}{V}(R_j)`$ $`=`$ $`{\displaystyle \frac{2V(R_j)}{3}}+{\displaystyle \frac{2\lambda (\delta \tau )^2}{9}}{\displaystyle \underset{i=1}{\overset{N}{}}}(_iV(R_j))^2`$ (9)
if $`j`$ is odd, and zero if $`j`$ is even. Here, $`_iV(R)`$ is the gradient of the total potential energy for the configuration $`R`$, with respect to the coordinates of the $`i`$th particle. This is a particular case of a more general expression, which can be shownmoron ; voth to be accurate up to terms of order $`\tau ^4`$ in the expansion of the exact density matrix $`\rho (R,R^{},\tau )`$ in powers of $`\tau `$.
Using the form (8) instead of the (superior) pair-product approximation (5), results in a substantially larger value of $`L`$ required to achieve convergence. For the specific physical system that we have chosen to test our algorithm, namely superfluid <sup>4</sup>He, the number $`L`$ of imaginary time slices needed is as much as 16 times greater than if (5) had been used. The reason for our choice is that our interest in primarily methodological. Specifically, we wish to separate the relative contributions to the effectiveness of a PIMC implementation, of the permutation sampling procedure and of the high-temperature density matrix utilized. A more stringent test is provided of our sampling scheme, if it can be shown to work satisfactorily with a fairly simple approximation for $`\rho _{}`$.
### II.1 Path Sampling
The generation of the set of many-particle paths $`\{X_p\}`$, with $`p=1,2,\mathrm{}M`$, can be conveniently carried out using the Metropolis algorithm. According to the standard procedure,allen one performs a random walk through the space of $`N`$-particle paths $`X`$, defined above, starting from an initial point $`X_{}`$. The $`X_p`$’s are then the points sequentially visited by the random walk.
Let $`X_l`$ be the $`l`$th element of the set $`\{X_p\}`$; in order to generate $`X_{l+1}`$, one samples a modification of the path $`X_l`$, involving new positions of one or more particles at several points (i.e., configurations) along $`X_l`$. Let $`X_l^{}`$ be the path arising from such modification of $`X_l`$, and let $`T(X_lX_l^{})`$ be the probability with which $`X_l^{}`$ is sampled from $`X_l`$. The proposed new path is accepted, thereby becoming the next point of the random walk (as well as the next element $`X_{l+1}`$ of the set $`\{X_p\}`$), with probability
$$W(X_lX_l^{})\frac{\overline{\rho }(X_l^{})}{\overline{\rho }(X_l)}\frac{T(X_l^{}X_l)}{T(X_lX_l^{})}$$
(10)
This is simply done by drawing a random number $`\chi `$ between zero and one; if $`W(X_lX_l^{})>\chi `$, then $`X_{l+1}X_l^{}`$, otherwise $`X_{l+1}X_l`$.
Of fundamental importance to the efficiency, unbiasedness and correctness of the algorithm, are the elementary moves whereby one generates the “trial” path $`X_l^{}`$ starting from $`X_l`$. In our PIMC implementation, two different types of moves are performed. A detailed descriptions of these moves is offered in the next two subsections.
#### II.1.1 “Wiggle” type moves
These moves modify the current path $`X_l`$ by just altering the path of one particle, randomly chosen. Random displacements are applied to a number $`s1`$ of consecutive positions of that particle along its path. This can be thought of as “chopping off” a portion of path, and replacing it with a different segment. The maximum number of positions modified by the update is $`L1`$, as the two ends of the portion are left unchanged. Because paths are periodic, it is possible to update a portion of path of a single particle that will include the zeroth or $`L`$th positions.noteb
In order to illustrate this type of move in detail, let us assume that a particle has been selected, and let the portion of path to be updated include the positions $`𝐫_{k+1}\mathrm{}𝐫_{k+s1}`$, where $`0k<L1`$ is an integer number, and $`s=2^m`$ is chosen so that $`2sL`$. Let $`𝐫_{k+1}^{}\mathrm{}𝐫_{k+s1}^{}`$ be the tentative new positions of the particle, selected according to some (yet unspecified) probabilistic criterion, expressed by a sampling function $`T`$. For notation purposes, we also define $`𝐫_k^{}=𝐫_k`$ and $`𝐫_{k+s}^{}=𝐫_{k+s}`$. Based on (8) and (10), the acceptance probability of the move will be
$$W=\frac{\left\{_{j=0}^{s1}\rho _F(𝐫_{k+j}^{},𝐫_{k+j+1}^{},\delta \tau )\right\}\mathrm{exp}\left[\delta \tau _{j=1}^{s1}U(R_{k+j}^{})\right]}{\left\{_{j=0}^{s1}\rho _F(𝐫_{k+j},𝐫_{k+j+1},\delta \tau )\right\}\mathrm{exp}\left[\delta \tau _{j=1}^{s1}U(R_{k+j})\right]}\frac{T(X^{}X)}{T(XX^{})}$$
(11)
having defined $`R^{}`$ as the configuration that differs from $`R`$ only by the displacement of the chosen particle from $`𝐫`$ to $`𝐫^{}`$, $`XR_0,R_1,\mathrm{}R_L`$ is the current path, whereas $`X^{}R_0,R_1,\mathrm{}R_k,R_{k+1}^{},\mathrm{}R_{k+s1}^{},R_{k+s},\mathrm{}R_L`$ is the proposed new path. There is considerable freedom in choosing the sampling probability $`T`$,allen but it is clearly advantageous to do so in a way that will simplify the expression (11). An obvious choice is
$`T(XX^{})={\displaystyle \underset{j=0}{\overset{s1}{}}}\rho _F(𝐫_{k+j}^{},𝐫_{k+j+1}^{},\delta \tau )`$ (12)
which reduces the acceptance probability (11) to
$$W=\mathrm{exp}\left[\delta \tau \underset{j=1}{\overset{s1}{}}\left(U(R_{k+j}^{})U(R_{k+j})\right)\right]$$
(13)
The probability $`T`$ so defined can be conveniently sampled by means of the “staging” algorithm.pollock84 The idea is as follows: $`T`$ is given by a product of $`s=2^m`$ Gaussian terms, namely
$`T(XX^{})\mathrm{exp}\left[{\displaystyle \frac{(𝐫_k^{}𝐫_{k+1}^{})^2}{4\lambda \delta \tau }}\right]\times \mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+1}^{}𝐫_{k+2}^{})^2}{4\lambda \delta \tau }}\right]\mathrm{}`$
$`\mathrm{}\times \mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+s1}^{}𝐫_{k+s}^{})^2}{4\lambda \delta \tau }}\right]`$ (14)
Using some simple algebranote2 it is possible to re-organize this product in the following, “hyerarchical” form:
$`T\mathrm{exp}\left[{\displaystyle \frac{(𝐫_k^{}𝐫_{k+s}^{})^2}{4s\lambda \delta \tau }}\right]\times \mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+s/2}^{}\overline{𝐫}_{k,k+s}^{})^2}{s\lambda \delta \tau }}\right]`$
$`\times \left\{\mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+s/4}^{}\overline{𝐫}_{k,k+s/2}^{})^2}{s\lambda \delta \tau /2}}\right]\mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+3s/4}^{}\overline{𝐫}_{k+s/2,k+s}^{})^2}{s\lambda \delta \tau /2}}\right]\right\}`$
$`\times \left\{\mathrm{exp}\right[{\displaystyle \frac{(𝐫_{k+s/8}^{}\overline{𝐫}_{k,k+s/4}^{})^2}{s\lambda \delta \tau /4}}\left]\mathrm{exp}\right[{\displaystyle \frac{(𝐫_{k+3s/8}^{}\overline{𝐫}_{k+s/4,k+s/2}^{})^2}{s\lambda \delta \tau /4}}]`$
$`\mathrm{exp}[{\displaystyle \frac{(𝐫_{k+5s/8}^{}\overline{𝐫}_{k+s/2,k+3s/4}^{})^2}{s\lambda \delta \tau /4}}]\mathrm{exp}[{\displaystyle \frac{(𝐫_{k+7s/8}^{}\overline{𝐫}_{k+3s/4,k+s}^{})^2}{s\lambda \delta \tau /4}}]\}`$
$`\times \left\{\mathrm{exp}\left[{\displaystyle \frac{(𝐫_{k+s/16}^{}\overline{𝐫}_{k,k+s/8}^{})^2}{s\lambda \delta \tau /8}}\right]\mathrm{}\mathrm{etc}\mathrm{}\right\}`$ (15)
where we have defined $`\overline{𝐫}_{\alpha ,\beta }^{}(𝐫_\alpha ^{}+𝐫_\beta ^{})/2`$. All distances are assumed to be computed with periodic boundary conditions, using the minimum image convention. Expression (II.1.1) immediately suggests a sequential, multi-level procedure to generate trial random positions of the particle being displaced. Since $`𝐫_k^{}=𝐫_k`$ and $`𝐫_{k+s}^{}=𝐫_{k+s}`$, the first factor in (II.1.1) does not enter the sampling in this type of move. Thus, one starts by generating a new “midpoint” position $`𝐫_{k+s/2}^{}`$, by sampling a three-dimensional Gaussian distribution function of semi-width $`\sigma _0=\sqrt{s\lambda \delta \tau /2}`$, centered at $`(𝐫_k^{}+𝐫_{k+s}^{})/2`$. It is customary to refer to the generation of the new midpoint as the zeroth level ($`l=0`$).
One then proceeds to the first level ($`l=1`$), where the two random positions $`𝐫_{k+s/4}^{}`$ and $`𝐫_{k+3s/4}^{}`$ are sampled from Gaussian distribution functions of semi-width $`\sigma _1=\sqrt{s\lambda \delta \tau /4}`$, centered at positions $`(𝐫_k^{}+𝐫_{k+s/2}^{})/2`$ and $`(𝐫_{k+s/2}^{}+𝐫_{k+s}^{})/2`$. At the next level ($`l=2`$), one generates four new positions and so on. The $`l`$th level involves the generation of $`2^l`$ new positions, sampled from Gaussian distribution functions of semi-width $`\sigma _l=\sqrt{2^{l1}s\lambda \delta \tau }`$. This “bisection” procedure ends when new positions of the particle have been generated at $`k+1,k+2,\mathrm{},k+s1`$. The last level is obviously $`l=m1`$.
The proposed new path $`X^{}`$ may be either accepted or rejected following an acceptance test based on (13). It proves much more efficient, however, to break down this global, final acceptance test, into $`m1`$ intermediate acceptance tests, each following every level of update. Specifically, after completing the $`l`$th level one proceeds to the next level with probability
$$W(ll+1)=\mathrm{exp}\left[\delta \tau \underset{jϵlevell}{}\left(U(R_{k+j}^{})U(R_{k+j})\right)\right]$$
(16)
aborting the the process (i.e., rejecting the proposed new path in its entirety) on the first negative outcome of an acceptance test. It is simple to see that the overall acceptance probability for the new path remains the same, given by (13), on breaking down the acceptance test by levels as explained above. The improvement in efficiency comes from the fact that the final acceptance is mostly influenced by the largest displacements, e.g., that of the midpoint. Thus, one can reject early (i.e., after the first level), and with relatively little computational effort, moves that most likely will eventually be rejected anyway.
The value of $`s`$ (namely the length of the portion of path that is updated) is set to ensure an optimally efficient sampling. The minimum possible value is $`s=2`$, which gives the highest acceptance rate, but at the same time also produces a modest path update (a single point of the path is modified). This becomes inefficient at low temperature, as paths can be fairly long (e.g., several hundred slices) and such “single-slice” updating can result in a very slow diffusion through configuration space, and consequently in undesirably long equilibration and auto-correlation times. It is therefore advantageous to take $`s`$ as large as possible, while still ensuring a reasonably high acceptance rate for multi-level moves (acceptance rate is a rapidly decreasing function of $`s`$). In our algorithm, we typically adjust $`s`$ so that the acceptance rate remains roughly between 20% and 50%.
#### II.1.2 “Permute” type moves
These moves involve a group of $`1<nN`$ particles. They are similar to the wiggle type moves, in that corresponding portions of the paths of the $`n`$ particles are modified, at $`s1`$ consecutive points. An additional aspect, however, is that the modified portion of the path of a particle in the group will connect, at $`k+s`$, to the path of a different particle, among the $`n`$ selected. This elementary move clearly allows one to sample permutations of the $`N`$ particles in the system, over the imaginary time interval $`[0,\beta ]`$.
The basic scheme of the move is as follows: first, a permutation of particle labels at $`k+s`$ is sampled; the number $`n`$ of particles involved in this permutation (henceforth referred to as the cycle) is not chosen a priori, but can vary from 2 to $`N`$. Once the permutation is selected, new single-particle paths are constructed, in a way completely analogous to that used in the “wiggle” moves. Finally, the new many-particle path $`X^{}`$ so obtained is accepted with probability given by (11). Again, it proves convenient to choose a sampling probability for the permutation that will in turn simplify the acceptance probability. Going back to Eq. (II.1.1), the first term of the product is now used to sample permutations, whereas the remaining terms are used to construct paths consistent with the permutations that have been sampled.
The sampling of a permutation is a recursive process in which particles are successively added to the cycle. The addition of a single particle includes an acceptance test, and the sampling of the particle from a table. One begins by selecting a random particle, say the $`\nu `$th for definiteness. Based on (II.1.1), a table is constructed, $`K_{\nu \omega }^{(1)}`$, with entries as follows:
$$K_{\nu \omega }^{(1)}=\rho _F(𝐫_{\nu k},𝐫_{\omega k+s},s\delta \tau )(1\delta _{\nu \omega })$$
(17)
where $`𝐫_{\nu k}`$ is the position of the $`\nu `$th particle particle at point $`k`$, whereas $`𝐫_{\omega k+s}`$ is that of the $`\omega `$th particle at point $`k+s.`$ note3 At this point, a first acceptance test is performed, namely the process will continue on to the next stage (i.e., selection of the permutation partner for particle $`\nu `$) with probability
$$C^{(1)}=\frac{Q_1}{Q_1+\rho _F(𝐫_{\nu k},𝐫_{\nu k+s},s\delta \tau )}$$
(18)
where $`Q_1=_\omega K_{\nu \omega }^{(1)}`$. If the acceptance test fails, then the process is aborted, i.e., the permutation move is rejected. Suppose, instead, that a positive outcome is obtained; an entry $`\alpha `$ is then sampled from the table $`K^{(1)}`$, with probability $`\mathrm{\Pi }_\alpha =K_{\nu \alpha }^{(1)}/Q_1`$. We see from (17) that particle $`\nu `$ itself is sampled with probability zero, i.e., the sampling of a “non-identical” permutation is forced here. The particle labeled $`\alpha `$ is selected as the second member of the permutation cycle being constructed. That means that, in the trial path $`X^{}`$, the path of particle $`\nu `$ will go through $`𝐫_{\nu k}`$ at the $`k`$th point and through $`𝐫_{\alpha k+s}`$ at point $`k+s`$. At this point, one has to sample a new position of particle $`\alpha `$ at $`k+s`$. Just as for particle $`\nu `$, one constructs a table
$$K_{\alpha \omega }^{(2)}=\rho _F(𝐫_{\alpha k},𝐫_{\omega k+s},s\delta \tau )(1\delta _{\alpha \omega })$$
(19)
and another acceptance test analogous to (18) is carried out, based on the probability
$$C^{(2)}=\frac{Q_2}{Q_2+\rho _F(𝐫_{\alpha k},𝐫_{\alpha k+s},s\delta \tau )}$$
(20)
with $`Q_2=_\omega K_{\alpha \omega }^{(2)}`$ (Again, the process is aborted if this acceptance test fails). An entry $`\mu `$ is sampled with probability $`\mathrm{\Pi }_\mu =K_{\alpha \mu }^{(2)}/Q_2`$; in this case, particle $`\nu `$ is sampled with finite probability, as the cycle can close on the initial particle, whereas particle $`\alpha `$ is now excluded from the sampling. At this point, if $`\mu =\nu `$, then the permutation cycle is closed, and it includes two particles, namely $`\nu `$ and $`\alpha `$. If, on the other hand, $`\mu \nu `$, then one must find another particle $`\gamma `$, which will become a member of the cycle, such that $`𝐫_{\mu k}=𝐫_{\gamma k+s}`$. Again, one constructs a table $`K^{(3)}`$ as above, the only difference being that now both $`\alpha `$ and $`\mu `$ are excluded from consideration, as both $`𝐫_{\alpha k+s}`$ and $`𝐫_{\mu k+s}`$ are already taken:
$$K_{\mu \omega }^{(3)}=\rho _F(𝐫_{\mu k},𝐫_{\omega k+s},s\delta \tau )(1\delta _{\alpha \omega }\delta _{\mu \omega })$$
(21)
A new acceptance/rejection test is performed, based on a probability $`C^{(3)}`$ defined analogously to (18)-(20):
$$C^{(3)}=\frac{Q_3}{Q_3+\rho _F(𝐫_{\mu k},𝐫_{\mu k+s},s\delta \tau )+\rho _F(𝐫_{\mu k},𝐫_{\alpha k+s},s\delta \tau )}$$
(22)
with $`Q_3=_\omega K_{\mu \omega }^{(3)}`$, and in case of success, one proceeds to sample entry $`\gamma `$ from table $`K^{(3)}`$, with probability $`\mathrm{\Pi }_\gamma =K_{\mu \gamma }^{(3)}/Q_3`$.
The basic idea should now be clear: This procedure is iterated until the cycle is finally closed, namely until particle $`\nu `$ is obtained from the sampling of the table $`K^{(n1)}`$. Two fundamental aspects of the above scheme to sample permutation cycles, are the exclusion from the tables $`K^{(n)}`$ of entries corresponding to particles already in the cycle (the $`\nu `$th is only excluded from $`K^{(1)}`$), and the acceptance tests based on $`C^{(n)}`$, preceding each new particle selection. The sum at the denominator of $`C^{(n)}`$ includes, besides $`Q_n`$, free-particle density matrices associated to all the entries excluded in the sampling table $`K^{(n)}`$.
Once a complete cycle has been obtained, one must construct trial paths for all particles in the cycle, consistent with the selected new positions at slice $`k+s`$. This second part is done in exactly the same way as for the “wiggle” moves, using the same sampling probabilities. Specifically, new midpoint positions are first sampled for all particles in the cycles; then, new positions at $`k+s/2`$ and $`k+3s/4`$ are sampled, and so on, with acceptance tests as in (16) after each level of path update. Note that the values of $`U`$ at points $`k`$ and $`k+s`$ remain unchanged, as no particle is displaced at these points; only particle labels are altered at point $`k+s`$. It is a simple matter to show that the path sampling probability arising from the above scheme is indeed consistent with (10).
The above scheme to sample permutations is similar to one described in Ref. ceperley95, ; the main difference, possibly significant, is that, in our procedure, one need not include, in any of the acceptance tests, a sum of terms representing all $`n`$ starting points of the cyclic permutation. Moreover, in our method $`n`$ distinct particles are sampled by construction, and the “identity” permutation, namely the one in which particle labels are left unchanged from $`k`$ to $`k+s`$, is excluded from the sampling.
Just as in the “wiggle” moves, $`s`$ must be chosen appropriately, namely, long enough that non-trivial permutation cycles can be sampled with appreciable probability. If $`s`$ is small, particularly if one is working with a small value of $`\delta \tau `$, the functions $`\rho _F`$ are negligible small for distances of the order of the average distance between particles, rendering it exceedingly unlikely to go beyond the first acceptance test (Eq. 18). On the other hand, taking $`s`$ too long, while allowing for large permutation cycles being sampled, results in very low overall acceptance for these cycles, much for the same reasons why acceptance falls for the “wiggle” moves as well if $`s`$ is too large. In the calculations whose results are illustrated in the next section, we have generally found that the optimal choice of $`s`$ is generally the same as for the “wiggle” moves.
Even when $`s`$ is optimally chosen, typical values of acceptance for permutations are low. One tries to keep the efficiency reasonable by attempting a large number permutational moves, which can be done fairly rapidly within the relatively simple scheme outlined above. We typically attempt several tens of thousands of permutations between two consecutive sets of “wiggle” type moves, in which full updates of the paths of all particles are attempted.
## III Results
We now describe the results of our PIMC simulations of condensed bulk <sup>4</sup>He at low temperature (1 K $`T`$ 4K), obtained with the algorithm described in the previous section. The model Hamiltonian for the system of interest is given by (1). For the purpose of comparing our results with existing calculations,ceperley95 ; ceperley86 we used an early version of the Aziz potentialaziz to describe the interaction between a pair of <sup>4</sup>He atoms ($`\lambda `$=6.0596 KÅ<sup>2</sup>).
We have computed several energetic and structural properties of the system. We have observed convergence of the energy estimates with a value of the “imaginary time step” $`\delta \tau `$=1/640 K<sup>-1</sup>. All of the results presented in this section are obtained with this value of $`\delta \tau `$. We estimate any residual, systematic error on the energy arising from our path discretization, to be worth no more than 0.1 K per <sup>4</sup>He atom. For structural properties, we have observed that estimates obtained with (up to four times) larger values of $`\delta \tau `$ are indistinguishable, within statistical uncertainties, from those obtained with the above-mentioned value of $`\delta \tau `$.
We found our optimal value of $`s`$, for both “wiggle” and “permute” type moves, to be $`s=6`$.notec Accordingly, the length of the portion of path that is updated on each move is $`2^s=64`$ imaginary time slices, which corresponds to an imaginary time interval of 0.1 K<sup>-1</sup>, with our choice of time step. Unless otherwise stated, the number of particles in our simulated system is $`N`$=64, as in Ref. ceperley86, , but results for $`N`$=216 were obtained as well..
### III.1 Energetics
The energy estimators utilized in this work are described, for instance, in Ref. voth, . Specifically, the average kinetic energy per particle $`K`$ is obtained as
$`K{\displaystyle \frac{1}{2\delta \tau }}{\displaystyle \frac{1}{4\lambda \delta \tau ^2}}(𝐫_k𝐫_{k+1})^2+{\displaystyle \frac{\lambda \delta \tau ^2}{9}}(V(R_{2k})^2`$ (23)
where $`\mathrm{}`$ stands for statistical average, $`(𝐫_k𝐫_{k+1})^2`$ is the square distance between the positions of a particle at adjacent points along the path, whereas the gradient of the potential energy in the third term is taken with respect to the coordinate of one of the particles, at an “even” slice. The potential energy per particle $`v`$ is instead obtained as
$$v\frac{1}{N}V(R_{2k1})$$
(24)
Both relations (23) and (24) are approximate, approaching the exact results only in the limit $`L\mathrm{}`$, $`\delta \tau 0`$. The above kinetic energy estimator is not the most efficient; it is known that the so-called “virial” estimator yields more accurate results (i.e., smaller statistical errors), given the same amount of computing time.cha99 However, the estimator (23) has been most commonly adopted in previous calculations of this type. In all of our calculations, we estimated the contribution to the potential energy attributed to particles outside the main simulation cell by assuming that the pair correlation function $`g(r)`$ equals one outside the cell.
Table 1 summarizes our results for the energetics of bulk <sup>4</sup>He, at different temperatures. Shown in parentheses are the corresponding results from Ref. ceperley86, . The estimates are in quantitative agreement, taking into account the statistical uncertainties of the two calculations. Amusingly, our statistical errors on the kinetic energy are not much smaller than those of Ref. ceperley86, , in spite of the fact that our calculation benefits of two more decades of advances in computing hardware. This is not too surprising, however, as the calculation of the kinetic energy, especially based on the estimator (23), is knowncha99 to be the place where the limitations are most evident of using less than optimal a high-temperature density matrix $`\rho _{}`$, such as the one used in this work. Still, the results of our calculation seem altogether satisfactory, giving us confidence that our PIMC implementation, including the new permutation sampling engine, performs correctly.
### III.2 Single-particle diffusion in imaginary time
We observe excellent agreement between our results and those of Ref. ceperley86, for structural properties, such as the pair correlation function (an example is given in Fig. 1); however, effects of quantum statistics on these quantities are small,ceperley95 and therefore their computation does not provide a particularly significant test of an algorithm to simulate indistinguishable quantum particles.
More telling are measures of the diffusion of particles in imaginary time. Fig. 2 shows results for the quantity $`D(\tau )`$, defined as
$$D(\tau )=\frac{\left(𝐫(\tau )𝐫(0)\right)^2}{6\lambda \tau }$$
(25)
where $`\mathrm{}`$ stands for statistical average. The two curves shown in the figures represent values of $`D`$ computed by PIMC for bulk <sup>4</sup>He (solid line), as well as for a system of distinguishable <sup>4</sup>He atoms, both at a temperature $`T`$=2 K. While in the first case <sup>4</sup>He atoms are treated as bosons, and therefore permutations are included, in the second case no permutations of particles are allowed.
Obviously, because in the latter case one must have $`𝐫(\beta )=𝐫(0)`$, i.e., single-particle paths must close onto themselves, it must be $`D(\beta )=0`$. On the other hand, if permutations are allowed, then single-particle paths can become entangled, and $`D(\beta )`$ may take on a finite value. Moreover, the value of $`D`$ is greater, at all imaginary times, in the case of Bose statistics; this is fairly intuitive, as the fact that particles are indistinguishable enhances the degree of delocalization of each individual particle.
### III.3 Superfluid density
We have also computed the <sup>4</sup>He superfluid density $`\rho _S`$, using the well-known “winding number” estimator.pollock87 At the lowest temperature considered in this work, namely $`T`$=1.1765 K, our result is $`1.02\pm 0.10`$, which is in agreement with experiment and with the PIMC result of Ref. pollock87, . We obtained this result with a number of slices $`L=544`$. It should also be mentioned that it appears possible to obtain a reasonably accurate estimate of $`\rho _S`$ using considerably fewer imaginary time slices (the result obtained with $`L=136`$ is indistinguishable, within statistical uncertainties, from the one quoted above), and that reducing $`L`$ also causes a significant reduction of the statistical error on $`\rho _S`$. In general, however, if $`L`$ becomes relatively large, namely of the order of a few hundred, lengthy simulations are required in order to reduce statistical error to an acceptable size (e.g., 0.05 or less). This problem seems common to other PIMC implementations as well, and it is not clear to us to what extent it may signal an inefficiency of our sampling method.
### III.4 Statistics of Permutations and Permutation Cycles
In order to characterize the performance of the permutation sampling algorithm, one may also look at quantities easily accessible in a simulation, which may not directly relate to anything measurable but provide a possible baseline for comparison of different algorithms. Table 2 provides statistics of permutation acceptance for a PIMC simulation of 64 <sup>4</sup>He atoms at $`T`$=1.1765 K (the lowest temperature considered here). The total number of permutations attempted in this run is $`4.5\times 10^5`$, and the fraction of accepted permutations (of any cycle length) is approximately 0.4%. Permutations were sampled over an imaginary time interval of length 0.1 K<sup>-1</sup>.
As one can see from the second column, 2-particle permutations are sampled overwhelmingly more than others; however, the rate of acceptance of attempted permutations is essentially constant, independent of $`n`$. This is found to be the case at all temperatures considered in this study. One may think that it would be advantageous to increase the rate at which permutations of more than two particles are sampled, since they presumably enhance the diffusion of the random walk through path space. Indeed, it is straightforward to generalize our sampling algorithm, so that permutations including more than two particles will be sampled more often. In practice, however, we found that including pair permutations is beneficial, in that it leads to a greater overall rate of acceptance of attempted permutations. How much of this is problem- or algorithm-dependent is difficult to say. Both quantities shown in Table 2 are rapidly decreasing functions of the temperature, as expected.
Although they are not sampled directly, permutation cycles involving large numbers of particles can and do occur, as a result of sampling many permutations involving few particles. Fig. 3 shows a histogram of probability for a particle to be part of a permutation cycle of length $`n`$ (i.e., involving $`n`$ particles) in a PIMC simulation carried out with the methodology illustrated above, at three different low temperatures below the $`\lambda `$-transition. Although the data are somewhat noisy, these results are in quantitative agreement with those of Ceperley ceperley95 for the same system, using the CP methodology. As the temperature is lowered, the probability that a particle will belong to a cycle of length $`n`$ becomes independent of $`n`$.
## IV Conclusions
A new algorithm to perform the sampling of permutations of indistinguishable particle in Path Integral Monte Carlo simulations was introduced. This procedure is similar, in spirit, to existing methods, but differs in some important aspects, and may have some advantages. We have tested it by performing a PIMC simulation of liquid <sup>4</sup>He at low temperature, in the superfuid regime. Aside from the permutation sampling scheme, the rest of the PIMC methodology utilized here is not optimized for <sup>4</sup>He calculations. In particular, it is worth repeating that much better options exist for the high-temperature density matrix, which can drastically reduce the number of time slices needed for convergence. Still, the calculation proves quite feasible with currently available, moderately powerful workstations. It should also be noted that, while the use of a more accurate high-temperature density matrix (specifically, the pair-product approximation (5)) greatly enhances the efficiency of calculations for a highly quantal, hard-sphere-like system such as helium, for other condensed systems such as molecular hydrogen, which feature a lesser degree of zero-point motion, or Coulomb system, for which the interaction potential is considerably less “stiff”, the high-temperature density matrix utilized here and Eq. (5) may be of comparable efficiency.
For comparison with existing calculations we limited the size of the system studied to $`N`$=64 particles, but it should be mentioned that simulations of systems with as many as four times more particles are also possible, with a reasonable amount of computer time (of the order of a month per thermodynamic point).
We have attempted to furnish as much quantitative information as possible, that may help assess the relative efficiency of the permutation scheme proposed here against existing ones. Obviously, a direct comparison of results provided by implementations only differing by the permutation scheme adopted, is also desirable. It is our hope that such a comparison will be soon carried out.
## Acknowledgments
This work was supported in part by the Petroleum Research Fund of the American Chemical Society under research grant 36658-AC5, by the Natural Sciences and Engineering Research council of Canada (NSERC) under research grant G121210893. The author wishes to acknowledge useful discussions with Saverio Moroni.
## TABLES
## FIGURE CAPTIONS
|
warning/0506/astro-ph0506518.html
|
ar5iv
|
text
|
# Measuring Galaxy Environments with Deep Redshift Surveys
## 1. Introduction
The observed properties of galaxies have long been known to depend upon the environment in which they are located. For instance, red, non-starforming galaxies (e.g. local ellipticals and lenticulars) are systematically over-represented in highly over-dense environments such as clusters (e.g. Davis & Geller, 1976; Dressler, 1980; Postman & Geller, 1984; Balogh et al., 1998). Recent studies have shown that the observed correlations between galaxy properties and environment are not limited to the cores of rich clusters, but extend to less dense domains including the outer regions of clusters, galaxy groups, and the field (e.g. Balogh et al., 1999; Carlberg et al., 2001; Blanton et al., 2003; Balogh et al., 2004; Croton et al., 2005).
There are a variety of physical processes that can readily explain these observational trends, including the action of dynamical friction, tidal stripping, or gas pressure in dense environments. These mechanisms, in combination with the hierarchical model of galaxy formation (Kauffmann et al., 1993; Somerville & Primack, 1999; Cole et al., 2000), in which galaxies form in less dense environments and are then accreted into larger groups and clusters, are generally consistent with the current observations. From the empirical evidence, however, it remains uncertain in what environment(s), by what mechanisms, and on what time-scales galaxies evolve from a field-like population to a cluster-like population. Still, the strong correlation of local galaxy density with galaxy properties over a broad range of environment does indicate that it plays an important role in galaxy formation and evolution.
To study galaxy properties spanning a broad and continuous range of local environment requires a thorough census of the 3-dimensional galaxy distribution over a large volume. Such data sets are only collected via large, systematic redshift surveys. For the nearby galaxy population, wide-field spectroscopic and photometric surveys (e.g. 2dFGRS, Colless et al. 2001 and SDSS, York et al. 2000) have provided excellent data sets for studying galaxy environments ranging from voids to rich clusters. Recent results from these large surveys have found that galaxy environments correlate strongly with the colors, luminosities, and morphologies of local galaxies (e.g. Balogh et al., 2004; Hogg et al., 2003; Blanton et al., 2003; Hogg et al., 2004).
With the advent of new high-redshift surveys (e.g. the VLT-VIMOS Deep Survey (VVDS), Le Fevre et al. 2003, 2004 and the DEEP2 Galaxy Redshift Survey (DEEP2), Davis et al. 2003; Faber et al. 2005), studies of galaxy environment will be able to extend beyond the local universe. Such deep, high-redshift surveys will provide a representative snap-shot of the galaxy population and corresponding local densities when the universe was half its present age, thereby permitting an investigation into the influence of environment upon galaxy evolution and formation. That is, extending environment studies to higher redshift will enable a determination of whether the correlations among galaxy properties observed in the local universe are the result of physical processes acting over the entire lifetime of the galaxy or whether the correlations were established during the early formation of the galaxy.
While redshift surveys have grown in scale and studies of galaxy environment have increased in prevalence, few published tests have detailed the degree to which environment measures are affected by survey limitations. For instance, the confined sky coverage of surveys introduces geometric distortions (or edge effects) which bias local density measures near boundaries or holes in the survey field. Environment statistics can also be impacted by the redshift precision and target selection requirements of a given survey. Mock galaxy catalogs provide an excellent means for testing the biases introduced to a given density measure by effects such as these.
In this paper, we test the applicability of several popular density estimators within deep redshift surveys utilizing the mock galaxy catalogs of Yan et al. (2004). In particular, we investigate the effects of redshift precision, survey field edges, redshift-space distortions, and target selection. We also devote specific attention to the DEEP2 survey with the goal of identifying the optimal density measure for use within the survey. Within this study, we do not consider global measures of environment trends (such as correlation functions), but instead focus on measures which can estimate environmental properties of individual objects. The outline of the paper is as follows. In the next section, we present the mock galaxy catalogs used to test environment measures at high redshift. Subsequently (§3), we describe the environment measures to be tested. In §4, we examine the significance of redshift precision in determining local galaxy densities. In §5, we conduct a detailed analysis of edge effects with respect to each environment parameter. In §6 and §7, we then address the influence of redshift-space distortions and target selection on the various environment estimators. In §8, the roles of completeness and the survey selection function are discussed. Finally in §9 and §10, we conclude with a summary of the applicability of each environment measure at high redshift and a discussion of the suitability of current deep surveys to measuring local galaxy densities.
## 2. Simulating Deep Redshift Surveys
Beginning with the Center for Astrophysics Redshift Survey (Davis et al., 1982), large redshift surveys have played a major role in studying galaxy properties, measuring cosmological parameters, and studying the large-scale structure of the universe. With improvements to astronomical instruments, local redshift surveys have ballooned in size and surveys at high redshift $`(z1)`$ have become possible. At present, high-$`z`$ surveys take two forms: (a) obtaining precise $`(\sigma _z0.001)`$ redshifts using spectroscopic observations of galaxies (e.g. DEEP2 and VVDS) and (b) using deep photometry in many passbands to make less precise $`(\sigma _z0.05)`$ photometric redshift estimates (e.g. COMBO-17, Wolf et al., 2003). Each of these has its advantages and disadvantages.
Even utilizing highly-multiplexed, multi-object spectrographs (e.g. DEIMOS, Faber et al. 2003) on large-aperture telescopes, a deep spectroscopic redshift survey requires a vast amount of telescope time and is invariably limited in the number of galaxies for which it can measure redshifts. Slit or fiber collisions constrain the number of objects able to be targeted during a given exposure while the forest of OH sky lines in the optical and infrared plus instrument defects and limited signal-to-noise cause redshifts to be missed for some percentage of targeted objects. Spectroscopic surveys benefit from a higher level of redshift precision which permits studies of kinematics within galaxies and galaxy groups, while also measuring spectral properties such as emission-line equivalent widths. On the other hand, using an imager with a large field-of-view and observing in many passbands, less precise photometric redshifts can be obtained for nearly all galaxies above a given magnitude limit in the targeted field. For this reason, photometric surveys are often able to build larger samples and are ideal for measuring the galaxy luminosity function or galaxy-galaxy lensing for which a high level of velocity accuracy is not necessary.
In this paper, we employ the simulated galaxy catalogs of Yan et al. (2004) to model both photometric and spectroscopic surveys at $`z1`$. The simulations and all work in this paper employ a $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`h=1`$, and $`\sigma _8=0.9`$. The mock catalogs are derived from N-body simulations by populating dark matter halos with galaxies according to a halo occupation distribution (HOD) function (Peacock & Smith, 2000; Seljak, 2000) which describes the probability distribution of the number of galaxies in a halo as a function of the host halo mass. The luminosities of galaxies are then assigned according to the conditional luminosity function (CLF) introduced by Yang et al. (2003), which describes the luminosity function in halos of mass $`M`$. Models for the HOD and the CLF are constrained from the 2dFGRS luminosity function (Madgwick et al., 2002) and two-point correlation function (Madgwick et al., 2003). By assuming that the manner in which dark matter halos are populated with galaxies does not evolve from $`z1`$ to $`z0`$ (Yan et al., 2003), the mock catalogs are built using simulation outputs at corresponding redshifts. The simulated galaxy catalogs show excellent agreement with the lower-z $`(0.7<z<0.9)`$ DEEP2 correlation function (Coil et al., 2004) and the COMBO-17 luminosity function (Wolf et al., 2003); they therefore should provide a realistic data set for studying measures of the environment of galaxies at $`z1`$. For further details regarding the construction of the mock catalogs, refer to Yan et al. (2003).
From the volume-limited mock catalogs, we are able to mimic a typical photometric redshift survey by selecting all galaxies above a given magnitude limit and applying to each galaxy redshift a random offset drawn from a Gaussian distribution with standard deviation $`\sigma _z`$. We utilize the DEEP2 survey as a model high-redshift spectroscopic survey. The volume-limited mock catalogs are selected according to the DEEP2 magnitude limit of $`R_{\mathrm{AB}}24.1`$ and passed through the DEEP2 target-selection and slitmask-making code, which is able to place approximately 60% of available targets on slitmasks for spectroscopy (Davis et al., 2004). Finally, 30% of objects are randomly rejected to reflect a conservative redshift success rate of $`70\%`$. The 12 mock catalogs cover fields of $`120^{}\times 30^{}`$ in area with a total of $`120`$ DEEP2 slitmasks tiling the 1 square degree. To simulate larger survey fields, we tiled multiple mock catalogs without overlap or discontinuity. Such large-field mocks were essential for studying edge-effects (§5) and for building large sample sizes.
In each mock catalog, we have a complete tally of the total galaxy distribution down to a luminosity of $`0.1\mathrm{L}_{}`$, along with subsets of objects which pass the DEEP2 target-selection criteria, were placed on a slitmask for observation, and yielded a successful redshift. Such a census enables detailed study of the survey selection function and the manner in which slitmask-making and target-selection affect environment statistics. Throughout this paper, we utilize several subsets drawn from the mock catalogs as described in Table 1. Note that for each mock galaxy, the simulations provide accurate positions in both real-space and redshift-space.
## 3. Environment Measures
The environment of a galaxy is typically defined in terms of the density of galaxies located in its immediate vicinity. However, a variety of density measures are often employed in estimating environment. For example, many previous analyses have focused on the identification of predefined groups or clusters of galaxies, which can be contrasted to those galaxies not inhabiting these over-dense regions – that is, the field population (e.g. Kuntschner et al., 2002; van den Bergh, 2002; Lewis et al., 2002; Christlein, 2000). Another approach is to instead derive a continuous measure of the galaxy density distribution, such as by measuring the distance to the $`n^{\mathrm{th}}`$-nearest neighbor (e.g. Gómez et al., 2003; Mateus & Sodré, 2004) or by directly smoothing the observed galaxy distribution on a fixed scale (e.g. Hogg et al., 2003; Beuing et al., 2002; Kauffmann et al., 2004). The underlying theme in each of these methods is that one requires a measure of the local number density of galaxies at the position of each galaxy in the sample.
In our analysis, we focus on density estimators that do not rely on identifying galaxy groups or clusters in any way. Lumping galaxies into predefined classifications provides a poorly sampled range of galaxy environments especially when compared to continuous measures of environment. At high redshift, where dense regions are commonly under-sampled and clusters and groups are less numerous, a more continuous definition of environment is all the more desirable. Still, identifying galaxies in groups at high $`z`$ is possible and has been tested in a separate paper (Gerke et al., 2005). In this analysis, we compare three popular density estimates: $`n^{\mathrm{th}}`$-nearest-neighbor distance, counts in an aperture, and the Voronoi volume. This set of measures is in no way presumed to be complete. Other promising methods for measuring local galaxy density, including using a Gaussian kernel to smooth the galaxy distribution over a given scale (e.g. Hogg et al., 2003; Balogh et al., 2004), are not discussed in this work.
### 3.1. $`n^{th}`$-Nearest-Neighbor Distance, $`D_n`$ and $`D_{p,n}`$
As first employed by Dressler (1980), the local galaxy density can be estimated using the distance to the $`n^{\mathrm{th}}`$-nearest, spectroscopically-observed neighbor of a given galaxy. Often, redshift information is simply employed to exclude foreground and background sources – by restricting neighbors to a given velocity interval – and the nearest-neighbor distance is measured in projection. Commonly, the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance, $`D_{p,n}`$, is expressed as a surface density, $`\mathrm{\Sigma }_n=n/(\pi D_{p,n}^2)`$. Measuring nearest-neighbor distances in projection is particularly useful when studying the density of galaxies in groups and clusters (e.g. Dressler, 1980; Lewis et al., 2002), where the appropriate velocity interval by which to exclude background and foreground galaxies may be selected according to the velocity dispersion of the group or cluster. In this manner, one may confidently exclude galaxies not associated with the group or cluster; furthermore, as shown in §6, measuring in projection minimizes the impact of redshift-space distortions.
For less dense environments or poorly sampled groups, there may be few neighbors within the selected velocity interval, causing $`D_{p,n}`$ to reflect the distance to other structures rather than the local density. In environments where working in projection is problematic, an alternative is to compute the $`n^{\mathrm{th}}`$-nearest-neighbor distance in 3-dimensions by searching in spherical apertures for the $`n^{\mathrm{th}}`$-nearest, spectroscopically-observed neighbor. Similar to its projected counterpart, the 3-dimensional (3-d) $`n^{\mathrm{th}}`$-nearest-neighbor distance, $`D_n`$, is often expressed as a number density, $`\rho _n=(3n)/(4\pi D_n^3)`$. Throughout this paper, all $`n^{\mathrm{th}}`$-nearest-neighbor distances are quoted in units of comoving $`h^1`$ Mpc and the symbols $`D_n`$ and $`D_{p,n}`$ are employed to denote the 3-dimensional and projected $`n^{\mathrm{th}}`$-nearest-neighbor distances, respectively.
To study the effectiveness of the 3-d and projected $`n^{\mathrm{th}}`$-nearest-neighbor distance measures at tracing the local density of galaxies in different environments, we compute both $`D_n`$ and $`D_{p,n}`$ for a DEEP2-selected sample consisting of 12,636 galaxies as drawn from a $`120^{}\times 60^{}`$ mock catalog. In Figure 1, we compare the values of $`D_n`$ and $`D_{p,n}`$ for each galaxy in the sample as measured using the redshift-space galaxy positions to the measured value of $`D_n`$ as computed using the real-space positions for each galaxy, which should reflect the true local environment. We find that at high densities, where redshift-space distortions are greater, the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance is superior at tracing the real-space density of galaxies but still suffers greatly from peculiar velocities. On scales corresponding to intermediate- and low-density environments, the 3-d measure of $`D_n`$ is a slightly more accurate estimate of the true galaxy distribution. For a DEEP2-selected mock catalog, $`15\%`$ of the observed sample resides in the regime $`(\mathrm{log}_{10}(D_5)0.5)`$ where the 3-d $`n^{\mathrm{th}}`$-nearest-neighbor distance saturates and loses sensitivity.
In addition to a lack of sensitivity on given scales, the behavior of the projected and the 3-d $`n^{\mathrm{th}}`$-nearest-neighbor distances depends upon the choice of $`n`$. A measure of the $`n^{\mathrm{th}}`$-nearest-neighbor distance effectively smoothes the galaxy distribution in a non-linear fashion according to the adopted value of $`n`$. If $`n`$ is chosen to be much larger than the richness of typical groups in the sample, then the $`n^{\mathrm{th}}`$-nearest-neighbor distances for galaxies in these groups will be pushed to erroneously high values, as it will reflect the distance to the next-nearest structure. In this work, we study both the projected and the 3-d methods for computing the $`n^{\mathrm{th}}`$-nearest-neighbor distance employing a variety of values for $`n`$. We limit most discussion to values of $`n=2,3,5`$ which correspond to the sizes of small groups detected in the DEEP2 survey (Gerke et al., 2005) and to the typical sizes of groups in the mock catalogs (see Fig. 2). For these values of $`n`$, the sensitivity of the $`n^{\mathrm{th}}`$-nearest-neighbor measure is rather independent of $`n`$. In under-dense environments, the dependence on $`n`$ is very weak, while in dense environments the strong clustering of groups (Padilla et al., 2004; Coil et al., 2005) arranges to curb $`D_n`$ for $`nn_{\mathrm{group}}`$. Also, in galaxy groups, redshift-space distortions are a much greater source of error in $`D_n`$ than small variations in the choice of $`n`$.
For the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance measure, we test the sensitivity of $`D_{p,n}`$ using line-of-sight velocity intervals ranging from $`\pm 750\mathrm{km}/\mathrm{s}`$ to $`\pm 2000\mathrm{km}/\mathrm{s}`$. As shown in Table 2, using a larger velocity interval by which to exclude foreground and background sources increases the accuracy of $`D_{p,n}`$ in dense environments but also sacrifices sensitivity at low densities. We find that for a DEEP2-selected sample, a velocity interval of $`\pm 1000\mathrm{km}/\mathrm{s}\pm 1500\mathrm{km}/\mathrm{s}`$ is best suited for a broad range of environments. Compared to photometric redshift errors (in the best datasets, $`\sigma _z6000\mathrm{km}/\mathrm{s}`$), the sizes of the tested line-of-sight velocity windows are small. However, larger velocity intervals sacrifice sensitivity on small scales and provide poorer measures of the local density about each galaxy; a window large enough not to be dominated by photometric redshift errors is also large compared to the typical length-scales of large-scale structure (e.g. the correlation length and typical void sizes), and thus provides a poor measure of environment.
### 3.2. Counts in an Aperture, $`C`$
Another method for estimating the local galaxy density is to count galaxies within a fixed metric aperture. For example, Hogg et al. (2003) count spectroscopically-observed galaxies in the SDSS within spheres of radii $`8h^1`$ comoving Mpc centered on each spectroscopically-observed galaxy. In high-redshift surveys where the survey field may cover $`1`$ square degree or less, such a large spherical aperture will be dramatically affected by the survey edges (see §5). For instance, within a $`30^{\prime \prime }\times 30^{\prime \prime }`$ field (i.e. $`20h^1`$ comoving Mpc on a side at $`z1`$), 81% of spherical apertures with a radius of $`1h^1`$ comoving Mpc will fit within the surveyed field, while for apertures of radius $`3h^1`$ and $`5h^1`$ comoving Mpc only 49% and 25% of the field, respectively, will be unaffected by edges. Furthermore, local studies indicate that larger apertures do not provide any advantage or additional information worth this high price. Both observational and theoretical studies, suggest that galaxy properties are more closely related to dark matter halo mass and small-scale environment than the large-scale environment of the galaxy (e.g. Lemson & Kauffmann, 1999; Blanton et al., 2004).
While choosing smaller spherical apertures would reduce the amount of survey volume affected by edges, apertures smaller than approximately $`\pm 1000\mathrm{km}/\mathrm{s}`$ along the line-of-sight will not be sensitive to galaxy groups or clusters. The counts in an aperture measure effectively smooths the data on some adopted scale, thereby losing sensitivity on smaller and larger scales. In our analysis, we employ a series of cylindrical apertures measuring 1–2 $`h^1`$ comoving Mpc transverse (radius) and $`\pm 500\mathrm{km}/\mathrm{s}\pm 2000\mathrm{km}/\mathrm{s}`$ along the line-of-sight. The dimensions of our cylindrical apertures are chosen to match the typical sizes of halos in the simulations (Yan et al., 2004).
### 3.3. The Voronoi Volume, $`V`$
The Voronoi volume is a geometric measure that has seen use from engineering and biology to astronomy (Ramella et al., 2001; Marinoni et al., 2002). Unlike counts in an aperture, the Voronoi volume does not smooth the galaxy distribution in any way. It provides a continuous, adaptive measure of galaxy density on all scales by measuring a unique volume about each spectroscopically-observed galaxy.
As illustrated in Figure 3, the Voronoi partition of space is the three-dimensional analogue of the two-dimensional Dirichlet tessellation, in which a plane containing a set of data points is divided into a set of polygons, each containing one of the points. A Voronoi polyhedron is the unique three-dimensional convex region of space surrounding a data point (the seed), such that within the polyhedron every point is closer to the seed than to any other data point. The faces of the Voronoi polyhedron are defined by the perpendicular bisecting planes of the vectors connecting the seed to its neighbors, where a seed’s neighbors are those points connected to it by the Delaunay complex – the set of tetrahedra whose vertices are at the data points and whose unique, circumscribing spheres contain no other data points. The Voronoi partition and Delaunay complex are thus geometrical duals of one another.
Computing the Voronoi partition for a galaxy redshift survey provides a natural way to measure the local density of galaxies, since the volume of a galaxy’s Voronoi polyhedron will vary inversely with the distance to its closest neighbors. For this reason, the Voronoi volume associated with each galaxy serves as a natural parameterization of that galaxy’s environment. Galaxies in dense regions will have small Voronoi volumes, while isolated galaxies will have larger volumes. Our methods for computing the Voronoi partition are identical to those of Marinoni et al. (2002), and we refer the reader to that work for computational details and for further discussion of the usefulness and historical context of the Voronoi partition and Delaunay complex. In this paper, we will employ the symbol $`V`$ to denote the Voronoi volume of a given galaxy and all Voronoi volumes are measured in units of comoving $`(h^1\mathrm{Mpc})^3`$.
## 4. Redshift Precision and Target Selection Rate: Photometric versus Spectroscopic Surveys
As discussed in §2, photometric and spectroscopic redshift surveys differ in the precision with which they are able to measure galaxy redshifts. To test the significance of redshift precision in measuring local galaxy environment, we have produced a variety of mock surveys with differing characteristics. First, we simulate two photometric redshift surveys which mimic the varying precisions of the COMBO-17 photometric redshift survey quoted in the literature. Our first simulated photometric redshift survey adopts a magnitude limit of $`R_{\mathrm{AB}}24.1`$ and redshift uncertainty of $`\sigma _z0.02`$ which reaches equally deep and is more precise than the COMBO-17 specifications of $`R_{vega}24`$, $`\sigma _z0.03`$ as given by Wolf et al. (2003). In addition, we simulate a photometric redshift survey with the same magnitude limit of $`R_{\mathrm{AB}}24.1`$ and a redshift precision of $`\sigma _z0.05`$ as specified for COMBO-17 by Taylor et al. (2004). Both magnitude-limited samples are drawn from the same volume-limited mock catalog and include 22,961 galaxies covering a $`120^{}\times 30^{}`$ field. Note that our assumed redshift uncertainties are lower limits to the redshift precision for the COMBO-17 survey. As discussed by Bell et al. (2004), the photometric redshift precision for COMBO-17 depends strongly on the galaxy type and redshift; galaxies such as starbursts which lack a strong 4000Å-break yield redshifts with much greater uncertainties $`(\sigma _z0.1)`$, while at higher redshifts K-correction uncertainties introduce systematic redshift errors.
Running the same volume-limited galaxy catalog through the DEEP2 target-selection and slitmask-making software and assuming a conservative redshift success rate (see §2), we also produce a mock spectroscopic sample with redshift precision of $`\sigma _z0.0001`$, mimicing the DEEP2 redshift survey. This DEEP2-selected spectroscopic sample includes 9,302 galaxies covering the same 1 square degree field (sampling $`50\%`$ of galaxies to the magnitude limit). To simulate the VVDS “deep survey” in the CDF-S (Le Fevre et al., 2004), we randomly select 25% of objects to the same $`R_{\mathrm{AB}}24.1`$ magnitude limit. This “VVDS-selected” sample is an optimistic simulation of VVDS, assuming a 100% redshift success rate (Vanzella et al., 2004) and ignoring differences in the bandpass used. Lastly, as a comparison sample, we select the full magnitude-limited mock (22,961 galaxies at $`R_{\mathrm{AB}}24.1`$) assigning redshifts according to the real-space positions of the galaxies as defined in the mock simulations. Each environment estimator ($`n^{\mathrm{th}}`$-nearest-neighbor distance, Voronoi volume, and counts in an aperture) is then computed on the photometric, spectroscopic, and real-space galaxy samples.
For this comparison, we restrict our analysis to the redshift range $`0.7<z<1.4`$ and to only those galaxies at transverse distances of greater than $`4h^1`$ comoving Mpc from the nearest edge in the survey volume. These restrictions make edge effects in both the redshift and transverse directions negligible, but do not introduce any selection biases (see §5). Note that the mock catalogs are not subject to interior edges; that is, the simulations cover a contiguous 1 square degree of sky with no holes.
We find that the precision of even the best photometric redshifts is not sufficient to measure local galaxy environments. Figure 4 shows the comparison between Voronoi volumes, $`V`$, as measured using the real-space galaxy positions compared to those calculated using the observed redshifts for two representative surveys. Even assuming redshift errors as small as $`\sigma _z0.02`$, the environment measured in a photometric redshift survey is insensitive for all but the very lowest density environments; the Spearman ranked correlation coefficient between the real-space and photometric measures of Voronoi volumes is $`\rho =0.4`$. For the spectroscopic survey, redshift-space distortions introduce some scatter at high densities, but the overall distribution of environments is well measured. In all, the Voronoi volumes measured from the observed spectroscopic redshift distribution trace the real-space Voronoi volumes with much greater precision, yielding a correlation coefficient of $`\rho =0.73`$. Very similar results are observed for the $`D_n`$, $`D_{p,n}`$, and $`C`$ environment estimators.
In Table 3, we expand our analysis to a better sampled range of redshift uncertainties. Even if the precision of photometric redshifts is greatly improved – by a factor of 2 or 4 – we find that low-resolution spectroscopic surveys with galaxy sampling comparable to DEEP2 provide a significantly better trace of the 3-dimensional galaxy environment. It is only at very high redshift precisions $`(\sigma _z0.005)`$ and when measuring densities in projection that photometric redshift surveys are able to rival their spectroscopic counterparts as probes of galaxy environment.
Among the spectroscopic redshift surveys simulated, the greater sampling and improved redshift precision of the DEEP2 survey prove significantly superior to the VVDS in tracing the real-space density of galaxies. On the other hand, at precisions better than the $`30\mathrm{km}/\mathrm{s}`$ uncertainty in DEEP2 redshifts, redshift-space distortions dominate the ability to measure local densities and thereby limit any advantage of improved redshift measurements (see Table 3). Note that the galaxy samples in Table 3 transition from a magnitude-limited $`(R_{\mathrm{AB}}24.1)`$ sample at low redshift precision to mimic photometric or grism spectroscopic redshift surveys to a sample selected using the DEEP2 target-selection and slitmask-making procedures or a VVDS-like selection to simulate higher-resolution spectroscopic surveys, which have superior redshift precision but lower sampling rates.
## 5. Edge Effects
When measuring galaxy densities within any survey, one must always be careful of edge effects introduced by the limited area of sky covered in the survey. Even using the largest optical telescopes and an instrument with a generous field-of-view, a deep redshift survey is limited in its ability to cover large regions. Furthermore, to minimize the effects of cosmic variance on the data set, a high-redshift survey is likely to spread the sky coverage over several fields. This limits the amount of contiguous sky coverage and increases the proportion of the survey area that is near an edge. In addition to the edges created by the chosen geometry of the survey field(s), edges and holes can be created in the data set by effects such as bright stars in a field or problematic regions in photometric detection which prohibit any galaxies from being targeted there. To start, we will restrict our discussion to survey edges in the plane of the sky, but later discussion will address edges in the line-of-sight direction.
The general effect of edges on each density estimator is to push the measurement towards lower density. To quantify the degree to which each environment measure is affected by edges we compute each measure on a large DEEP2-selected spectroscopic mock galaxy catalog covering a wide field, $`120^{}\times 90^{}`$. From the center of this larger simulation we extract a smaller rectangular survey field, covering $`40^{}\times 30^{}`$, and rerun each environment measure on this data set. For every galaxy in the smaller survey field, we then have measurements of environment unaffected by edges (when measured on the larger sample) and measurements in which survey edges play a greater role (when measured on the smaller field). Trimming to the redshift range $`0.7<z<1.4`$, the smaller sample consists of 2,803 galaxies. In the following subsections, we discuss how each environment measure is affected by the edges of the survey region on the plane of the sky. Note that we also ran tests incorporating holes and irregular survey edges with very little change in the relative results for the tested environment measures.
For each galaxy in the studied sample, we express the difference in a given environment measure due to survey edges as a fraction of the width of the distribution for that measure. Specifically, we define the percent change in environment measure $`X`$ for galaxy $`q`$ by
$$\mathrm{\Delta }_e(X)=\frac{\mathrm{log}_{10}\left(X_{2,q}\right)\mathrm{log}_{10}\left(X_{1,q}\right)}{\sigma _1}100\%$$
(1)
where $`X_{2,q}`$ is the measure of $`X`$ for galaxy $`q`$ computed on the smaller mock, $`X_{1,q}`$ is similarly computed on the wide-field mock, and $`\sigma _1`$ is a measure of the Gaussian width of the logarithmic distribution of environment measure $`X`$ as calculated in the larger simulation. Quantifying the change in each environment measure in this fashion enables the role of edge effects to be compared between different environment estimators in a uniform manner.
### 5.1. Survey Edges and $`n^{th}`$-Nearest-Neighbor Distance
The $`n^{\mathrm{th}}`$-nearest-neighbor distance environment measure – in both projection and 3-dimensions – is affected by edges in a predictable manner. Any galaxy with an edge located closer than the measured $`n^{\mathrm{th}}`$-nearest-neighbor distance must be affected by the survey edges. However, to remove all such galaxies based on this criterion ($`D_n>D_{\mathrm{edge}}`$ or $`D_{p,n}>D_{\mathrm{edge}}`$) biases the sample towards over-dense environments by excluding less-dense regions over a greater volume than more-dense regions. To avoid such biasing of the sample, a simple cut in edge distance can be made – excluding all galaxies within some distance of the nearest edge. This cut introduces no environment-dependent bias, but does allow some contamination to the sample at under-dense environments depending on the severity of the cut. In our simulations, we find that removing all galaxies within $`2h^1`$ comoving Mpc of a survey edge creates a catalog with minimal contamination (roughly 5% of the sample has $`\mathrm{\Delta }_e(D_3)>10\%`$) and still retains 65% of the data set. Relaxing the constraint to $`D_{\mathrm{edge}}>1h^1`$ comoving Mpc, the level of contamination in the sample doubles to $`10\%`$ with $`\mathrm{\Delta }_e(D_3)>10\%`$ while the percentage of the sample retained increases to $`85\%`$.
As illustrated in Figure 5, we find that edge-effects show a clear dependence on $`n`$; the level of contamination due to survey edges in the plane of the sky increases by roughly a factor of two for $`D_5`$ relative to $`D_3`$. Also, for a fixed value of $`n`$, the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance is slightly more robust to edge-effects in the regime where sample size is maximized ($`D<2h^1\mathrm{Mpc}`$ in Figure 5).
### 5.2. Survey Edges and Counts in an Aperture
For an aperture of fixed comoving size, the edge effects upon the counts in an aperture density measurement are easily understood and cleaned from the sample. Only galaxies located within $`r_\mathrm{t}`$ of an edge are affected, where $`2r_\mathrm{t}`$ is the transverse diameter of the chosen aperture. Thus, by removing any galaxies within $`r_\mathrm{t}`$ of an edge, the sample is entirely devoid of edge-affected galaxies. Such a trimming of the data set does not introduce a selection effect, that is, there is no bias towards environments of a given sort.
In our simulated spectroscopic data set of 2,803 galaxies, $`15\%`$ of the sample are positioned within $`r_\mathrm{t}`$ of a survey edge using a cyclindrical aperture of $`r_\mathrm{t}=1h^1`$ comoving Mpc. However, only for a scant $`3\%`$ of the sample did we find $`\mathrm{\Delta }C0`$ when comparing measurements of $`C`$ made on the smaller simulation to those made on the wide-field sample. One possible means for salvaging some edge-affected galaxies would be to scale the measured counts in each aperture by the amount of the aperture contained within the survey area. Due to the low rate at which $`C`$ is perturbed by survey edges, however, such a correction actually causes an overestimate of $`C`$ for the majority of galaxies near the edge $`(D_{\mathrm{edge}}<r_\mathrm{t})`$ of the survey field.
### 5.3. Survey Edges and Voronoi Volumes
Due to the geometrical complexity of the Voronoi tesselation, understanding the effects of edges on the calculated Voronoi volumes is less straightforward than for the previously discussed density measures. For galaxies very close to exterior edges in the survey field, Voronoi volumes can be unbounded and such galaxies should be consequently discarded from the sample. On a more subtle level, edge effects – including interior edges – will also cause volumes to be increased in size while the volumes still remain bounded. Some of these edge-affected volumes can be detected as having Voronoi vertices outside of the survey field. However, many other edge-affected Voronoi volumes are not detectable in such a manner.
As illustrated in Figure 6, even excluding galaxies located near a survey edge (e.g. within $`2h^1`$ comoving Mpc), the distribution of Voronoi volumes is still greatly affected by edges with a bias towards large volumes being pushed to even larger values (see Fig. 6). It is possible to minimize this effect by retaining only galaxies with $`V`$ below some limit. In our simulations, by truncating at $`\mathrm{log}_{10}(V)=3.1`$, the amount of contamination due to edge effects can be reduced to approximately 20% of the sample with $`\mathrm{\Delta }_e(V)>10\%`$. While making such cuts according to distance to the nearest edge $`(D_{\mathrm{edge}}>2h^1\mathrm{comoving}\mathrm{Mpc})`$ and Voronoi volume $`(\mathrm{log}_{10}(V)<3.1)`$ effectively reduces the number of edge-affected galaxies in the sample, it also restricts the dynamic range of the Voronoi measure and considerably reduces the size of the sample. For our simulated spectroscopic sample of 2,803 galaxies, the Voronoi volume density measure was the most dramatically affected by edges with $`45\%`$ of the sample being corrupted, that is, having $`\mathrm{\Delta }_e(V)>10\%`$.
### 5.4. Effects of Finite Redshift Range
As a secondary effect, the finite redshift range probed by any survey imposes edges in the line-of-sight direction. The role of these edges is more easily handled by restraining all scientific analysis to a limited, well-sampled range of redshifts. In the DEEP2 survey, the ability to measure redshifts at $`z>1.4`$ or $`z<0.7`$ decreases significantly as the \[OII\] emission-line doublet leaves the observed optical window. In a DEEP2-selected mock catalog, we find that by restricting our sample to those galaxies at $`0.7<z<1.4`$, less than 1% of the sample has a $`5^{\mathrm{th}}`$-neareast-neighbor distance, $`D_5`$, greater than the distance to the $`z=0.7`$ or $`z=1.4`$ edge. Similar contamination rates are found for the other environment estimators.
A second concern for spectroscopic redshift surveys is the possibility of missing redshifts over specific wavelength intervals, especially in the far-optical and near-infrared where OH sky lines can dominate the spectrum. At $`z1`$ where optical surveys often rely upon a singular spectral feature (e.g. the \[OII\] doublet at $`\lambda _{\mathrm{rest}}3727\mathrm{\AA }`$) for redshift measurements, a hole in wavelength coverage translates directly into a hole in the survey’s redshift sampling. For the DEEP2 survey, the high-resolution $`(R5000)`$ of the DEIMOS data minimizes this effect, as the sky lines are then narrower than the components and spacing of the \[OII\] doublet. In truth, the DEEP2 redshift distribution exhibits no significant cross-correlation with a sky spectrum mapped to redshift according to either the central wavelength of the \[OII\] doublet or the wavelengths of either subcomponent (Newman et al., 2005). However, for lower resolution surveys such as the VVDS, windows of redshift insensitivity may be a concern that must be addressed in measuring galaxy densities.
## 6. Redshift-Space Distortions
While spectroscopic redshift surveys are able to measure galaxy redshifts with great precision, redshift measurements by nature are measurements of velocity and not distance. Accordingly, converting differences in redshift to relative line-of-sight distances is subject to the peculiar velocities of the galaxies. Such peculiar motions are greatest in dense regions such as groups or clusters where the velocity dispersion of the group causes the inter-member spacing to be larger in redshift space than in real space. Due to this environmental dependency of redshift-space distortions, it is essential to understand the manner in which they affect a given galaxy density measure.
Within a mock DEEP2-selected spectroscopic galaxy catalog covering $`120^{}\times 90^{}`$, we compute each environment measure using the both real-space positions of the galaxies and the observed redshift-derived positions. Restraining our analysis to galaxies at edge distances greater than $`4h^1`$ comoving Mpc and within the redshift range $`0.7<z<1.4`$, we quantify the effect of redshift-space distortions on each environment estimator by calculating the change in each measure as computed on the real-space mock relative to the corresponding measure derived from the observed spectroscopic mock. As in §5, we express the difference in a given environment measure as a fraction of the width of the real-space distribution for that environment measure. For example, the percent change in environment measure $`X`$ for galaxy $`q`$ is given by
$$\mathrm{\Delta }_z(X)=\frac{\mathrm{log}_{10}\left(X_{z,q}\right)\mathrm{log}_{10}\left(X_{R,q}\right)}{\sigma _R}100\%$$
(2)
where $`X_{z,q}`$ is the measure of $`X`$ for galaxy $`q`$ computed from the redshift-derived position and $`X_{R,q}`$ is similarly computed from the real-space position. The width, $`\sigma _R`$, is determined via a Gaussian fit to the logarithmic distribution of environment measure $`X`$ for all galaxies in the real-space simulation. Here, $`\sigma _R`$ can be measured on the real-space distribution of $`\mathrm{log}_{10}(X)`$ or the redshift-space distribution with negligible difference for the DEEP2-selected sample.
As illustrated in Figure 7, the Voronoi volume and the 3-dimensional $`3^{\mathrm{rd}}`$\- or $`5^{\mathrm{th}}`$-nearest-neighbor distances are similarly affected by redshift-space distortions. For each measure, the effects of redshift-space distortions are non-negligible and as shown in Figure 8 are greatest in over-dense environments. In comparison to the $`V`$ and 3-d $`D_n`$ measures, the counts in an aperture density estimator, $`C`$, and projected $`n^{\mathrm{th}}`$-nearest-neighbor measure, $`D_{p,n}`$, are less affected by the “fingers-of-god” due to their effective smoothing in the redshift direction; by definition, the projected estimators, $`C`$ and $`D_{p,n}`$ forfeit sensitivity in the redshift direction, which reduces their susceptability to redshift-space distortions. For nearly 80% of the sample, $`C`$ is unaffected $`(\mathrm{\Delta }C=0)`$ by peculiar motions when using a cylindrical aperture with a length of $`\pm 1000\mathrm{km}/\mathrm{s}`$ and diameter of $`1h^1`$ comoving Mpc. The sensitivitity of $`C`$ to redshift-space distortions is somewhat dependent upon the choice of the aperture size in the line-of-sight direction such that a smaller aperture is more adversely affected. For more than 90% of galaxies in our sample, we find $`|\mathrm{\Delta }C|1`$, again using an aperture with length of $`\pm 1000\mathrm{km}/\mathrm{s}`$. Similar results are found for the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance measure; more than 80% of the sample meets the criterion $`\mathrm{\Delta }_z(D_{p,n})<5\%`$ for $`n=3,5`$ and using a velocity interval of $`\pm 1000\mathrm{km}/\mathrm{s}`$ to exclude foreground and background sources.
## 7. Target Selection and Observation
There are inevitable trade-offs between the number density of sources targeted for observation and the area of sky covered in a redshift survey. Clustering of high-redshift galaxies and fiber or slit collisions on multi-object spectrographs conspire to limit the fraction of target objects which a survey can observe at one time. Furthermore, not every object targeted will successfully yield a redshift, generally due to finite signal-to-noise and instrumental effects.
The DEEP2 redshift survey will target $`50,000`$ galaxies covering $`3.5`$ square degrees of sky, over 80 nights on the Keck II telescope (Davis et al., 2004; Faber et al., 2005). This impressive survey, however, will only target approximately 60% of available high-redshift $`(0.7<z<1.4)`$ galaxies in its four fields and successfully measure redshifts for about 75% of targeted galaxies. The DEEP2 survey is designed with the goals of studying large-scale structure and galaxy properties at high redshift. Thus, the survey design attempts to maximize the number of redshifts obtained, to sample the galaxy distribution over a broad range of length scales, and to minimize the effects of cosmic variance. Due to slit collisions on DEIMOS slitmasks, the DEEP2 survey systematically under-samples regions of sky with a high surface density of galaxies (see Figure 9). It is critical for a study of galaxy environments to understand how this bias may affect the environment measured used.
While the detailed effects of target selection and redshift incompleteness are clearly specific to each survey, the goal of this section is to understand how a given environment measure is affected by the limited sampling common to all deep redshift surveys. In this work, we adopt the DEEP2 survey as a representative high-redshift, spectroscopic survey. As discussed in §2, the DEEP2 survey targets all galaxies at $`R_{\mathrm{AB}}24.1`$ according to a probabilistic algorithm which preferentially selects high-$`z`$ galaxies. Applying the DEEP2 target-selection and slitmask-making algorithms to a magnitude-limited $`(R_{\mathrm{AB}}24.1)`$ mock catalog covering $`40^{}\times 30^{}`$, the simulated survey targets and successfully measures redshifts for 2,839 of the 5,866 galaxies in the field and with $`0.7<z<1.4`$ (assuming a redshift success rate of $`70\%`$). To study the combined effects of target selection, slitmask making, and redshift success, we compute each environment measure on the DEEP2-selected sample and on the full magnitude-limited sample.
As illustrated in Figure 10, the counts in an aperture measure, $`C`$, shows no indication of an environment-dependent bias in the DEEP2 target selection. If DEEP2 severely under-samples dense environments, then we would expect to see a saturation in the observed value of $`C`$ at high densities relative to the estimation of $`C`$ computed on the magnitude-limited sample. Instead, we find a linear relation extending to dense environments which follows the $`50\%`$ overall completeness of the DEEP2-selected mock catalog.
Due to the fixed comoving aperture size of the counts in an aperture environment measure, it probes the same physical scale independent of how the targeted galaxies are selected. The $`n^{\mathrm{th}}`$-nearest-neighbor distance measure, on the other hand, can sample systematically different effective scales depending on the galaxy sampling. As illustrated in Figure 11, within the observed spectroscopic sample, the $`5^{\mathrm{th}}`$-nearest-neighbor distance is roughly tracing the $`10^{\mathrm{th}}`$-nearest-neighbor distance in the magnitude-limited mock; this is sensible, as the DEEP2-selected mock samples $`50\%`$ of galaxies and what was the $`10^{\mathrm{th}}`$-nearest neighbor in the magnitude-limited sample will typically be the $`5^{\mathrm{th}}`$ in the DEEP2-selected sample. Similarly, the $`2^{\mathrm{nd}}`$\- and $`3^{\mathrm{rd}}`$-nearest-neighbor distances are effectively tracing the $`4^{\mathrm{th}}`$\- and $`6^{\mathrm{th}}`$-nearest-neighbor distances, respectively, in the magnitude-limited mock (see Figure 11). While target selection and slitmask making affect the scale on which the $`n^{\mathrm{th}}`$-nearest-neighbor distance samples the galaxy distribution, they do so in a manner which does not depend upon environment. Thus, while DEEP2 under-samples dense *regions of sky*, the survey does not under-sample dense *environments* (see Figure 12).
The limited galaxy sampling of the DEEP2 survey causes the calculated Voronoi volumes to be systematically larger than if computed on the full magnitude-limited sample. We find that the level to which a given Voronoi volume is affected by the DEEP2 sample selection does not depend on $`V`$ or redshift; the limited sampling of the DEEP2 survey simply introduces a random scatter towards larger volumes. Similar to $`D_n`$ and $`C`$, we see no evidence for an environment-dependent bias due to the DEEP2 target-selection procedures.
## 8. Correcting for the Survey Selection Function
For any magnitude-limited survey, the fraction of total galaxies in a volume-limited sample within the magnitude limits varies – commonly decreasing – as a function of redshift. This variable sampling of the galaxy distribution as a function of redshift causes measurements of galaxy densities to depend strongly on $`z`$. For instance, if a survey under-samples at higher redshift, then estimates of $`D_n`$ and $`V`$ at high $`z`$ will be artificially inflated relative to estimates at low $`z`$; similarly $`C`$ will be under-estimated at higher redshift. Often, magnitude-limited redshift surveys are trimmed to a volume-limited sub-sample to avoid these issues. At high redshift, however, this can dramatically reduce the sample size; for example, selecting a vol-limited sub-sample within a DEEP2-selected mock catalog excludes as much as 40% of the observed galaxies. Furthermore, over regimes where luminosity evolution is significant $`(\mathrm{\Delta }_z0.1)`$, even defining a volume-limited sample can be problematic.
To utilize the entire survey sample or for surveys that do not follow a simple magnitude-limited target selection, the variations in the galaxy sampling rate with redshift may be quantified in terms of a survey selection function, $`s(z)`$, with which density estimates (number of galaxies per comoving volume or number of galaxies per projected comoving area) may be corrected as follows:
$$X_0(\alpha ,\delta ,z)=\frac{X_z(\alpha ,\delta ,z)}{s(z)w(\alpha ,\delta )},$$
(3)
where $`X_z`$ is the density estimate computed from the observed redshift distribution, $`w`$ is a 2-dimensional survey completeness map which accounts for variation in redshift completeness from field to field within the survey, and $`X_0`$ is the corrected density estimate.
There are several ways to determine the selection function of a survey. The most common approach is to first estimate the galaxy luminosity function (LF) for all of the galaxies in the redshift survey, and then use it to predict the redshift distribution of the sample (e.g. Madgwick et al., 2003). However, unless evolution is correctly incorporated, the LF will not alone be able to correctly predict the redshift dependence of the observed number density of galaxies in a deep survey. Furthermore, working at high redshift, it becomes increasingly difficult to constrain the galaxy LF since observations are limited to the brightest sources, thereby making estimations of the characteristic luminosity, $`M_{}(z)`$, and faint-end slope, $`\alpha (z)`$, less secure (e.g. Willmer et al., 2005; Wolf et al., 2003; Bell et al., 2004). For this reason, the selection function for a survey is often estimated by smoothing the observed number density of galaxies, $`n(z)`$, as a function of redshift and then normalizing according to an assumed dependence of the comoving number density of galaxies on $`z`$ (e.g. Coil et al., 2004). This has the disadvantage that density inhomogeneities in the survey will somewhat affect the derived redshift distribution, even with large smoothing kernels, due to the strength of cosmic variance; also, kernels large enough to minimize this will distort real features in $`n(z)`$, especially where there are large gradients. On the other hand, any evolution in the observed number density of galaxies with $`z`$ will be automatically incorporated into the estimation of $`s(z)`$.
Yet another approach to estimating $`s(z)`$ is to compute an analytical fit to the observed data from which the selection function is then derived (e.g. Cooper et al., 2005; Faber et al., 2005). Similar to a selection function estimated from smoothing the observed $`n(z)`$ distribution, an analytical fit to the data – or “fitting” method for estimating $`s(z)`$ – is subject to cosmic variance, but to a much smaller degree than the “smoothing” method, as small-scale variations in $`n(z)`$ which do not match the model are not allowed.
In this work, we estimate the survey selection function for the mock DEEP2 survey according to four different prescriptions: (a) estimating $`s(z)`$ by smoothing the observed $`n(z)`$ distribution in a DEEP2-selected mock catalog $`(120^{}\times 30^{})`$ assuming no evolution in the comoving number density of galaxies with redshift, using a similar algorithm as Coil et al. (2004) (“smoothing” method), (b) fitting for the selection function assuming a functional form for the redshift dependence of the successfully-observed $`dN/dz`$ and again assuming no evolution in the comoving number density of galaxies (“fitting” method), (c) determining the true selection function by computing the number density of available targets over many DEEP2 mock pointings relative to the volume-limited number density of galaxies in the mocks, and (d) deriving $`s(z)`$ from the conditional LF assumed in constructing the mock catalogs. This last approach is identical to the commonly-used method of estimating the selection function using the measured LF and predicting the redshift distribution of the underlying galaxy population. The first two methods, (a) and (b), are analagous to methods one might use to derive $`s(z)`$ solely from the observational data in a deep redshift survey and are subject to cosmic variance and uncertainties in the assumed normalization and redshift dependence of the comoving number density of galaxies. The latter two methods are effectively identical, and test that the mock catalogs are working as advertised.
In Figure 13, we present the mock selection functions derived using each of the methods described above. In general, the agreement between the different approaches for determining $`s(z)`$ is quite good. At the highest redshifts $`(z>1.0)`$, the different estimations of the selection function differ due to differences in the assumed comoving number density of galaxies at high redshift. The conditional LF adopted in constructing the mock catalogs yields evolution that produces a decrease in the comoving number density, $`\nu (z)`$, of galaxies at $`z>1.0`$, while in estimating $`s(z)`$ from the observed mock DEEP2-selected redshift distribution we assume a constant form for $`\nu (z)`$. Both estimations of the comoving number density at $`z>1.0`$ are consistent with existing observational evidence (Wolf et al., 2003; Willmer et al., 2005).
The footprint of large-scale structure on a selection function derived from smoothing the observed DEEP2-selected mock $`n(z)`$ distribution is reduced if a large smoothing kernel is used; here, we apply two successive smoothing windows of width $`\mathrm{\Delta }z=0.15`$. If the smoothing kernel is too small, the presence or absence of structures such as filaments or walls (i.e. cosmic variance) will cause us to overestimate or underestimate the fraction of galaxies sampled at a given redshift. Accordingly, over- or under-densities of galaxies will be inappropriately reduced in amplitude when corrected by the survey selection function; e.g. the presence of a filament will push the measured $`s(z)`$ up at its redshift, reducing the corrected density measured, $`X_0`$, artificially. Any smoothing large enough to erase the effects of cosmic variance in a survey covering a few square degrees will, unfortunately, cause flattening in the shape of $`s(z)`$, especially near the limits of the redshift range probed. Due to the drawbacks of smoothing, fitting for the selection function as detailed above is often a superior method for estimating $`s(z)`$ from the observed data, but it does require assumptions about the form of $`dN/dz`$, which smoothing does not.
To study the effectiveness of correcting the measured galaxy densities by the factor of $`1/s(z)`$ (see eq. 3), we have computed the projected $`7^{\mathrm{th}}`$-nearest-neighbor surface density, $`\mathrm{\Sigma }_7`$, within a volume-limited mock catalog covering $`120^{}\times 30^{}`$ of sky. We then compare this to the projected $`3^{\mathrm{rd}}`$-nearest-neighbor surface density, $`\mathrm{\Sigma }_3`$, for those galaxies successfully observed in the DEEP2-selected sample. We then correct these “observed” $`\mathrm{\Sigma }_3`$ values using each of the $`s(z)`$ shown in Figure 13, and also attempt an empirical correction. This correction is given by dividing each observed density value by the median $`\mathrm{\Sigma }_3`$ for galaxies at that redshift where the median is computed in a bin of $`\mathrm{\Delta }z=0.04`$. Correcting the measured densities in this manner converts the $`\mathrm{\Sigma }_3`$ values into measures of over-density relative to the median density and is similar to the methods employed by Hogg et al. (2003) and Blanton et al. (2003).
Figure 14 illustrates the effectiveness of each selection function at reproducing the redshift dependence of the galaxy density distribution as measured in the volume-limited sample. Within redshift bins of $`\mathrm{\Delta }z=0.02`$, we compute the difference in median density between the corrected $`\mathrm{\Sigma }_3`$ values and the median density, $`\mathrm{\Sigma }_7`$, of objects in the volume-limited sample. While each of the methods for estimating the survey selection function is an improvement over the uncorrected density distribution, an empirical correction (as described in the previous paragraph) which removes all $`z`$-dependence in the observed density distribution is at least as effective as correcting using a selection function (see Table 4).
## 9. Discussion
Every environment measure that we have considered has its advantages and disadvantages. The counts in an aperture measure, $`C`$, lacks sensitivity in low-density environments and while not lacking in dynamic range, it provides a non-continuous (or quantized) measure of environment, a particular disadvantage if the typical value of $`C`$ is small. It is best suited for working in dense environments where $`C`$ is more robust to redshift-space distortions than other measures and for analyses in which one wishes to classify a sample into coarse density bins or classes. The counts in an aperture statistic is unique among the environment estimators tested in that it measures the galaxy density on a clearly defined, fixed length scale. In contrast, the projected and 3-dimensional $`n^{\mathrm{th}}`$-nearest-neighbor distance measures probe the radius enclosing some total number of galaxies and are not direct density measures. The $`C`$ parameter also provides a great advantage via its robustness to survey edge effects.
Similar to the counts in an aperture statistic, the projected $`n^{\mathrm{th}}`$-nearest neighbor distance measure is well suited for measuring density in groups and clusters. However, unlike $`C`$, the projected $`D_{p,n}`$ parameter provides a continuous estimate of galaxy density extending to under-dense environments where it still provides a reasonably accurate measure. While slightly more robust to edges than its 3-dimensional counterpart, $`D_n`$, the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance is more prone to survey edge contamination than the counts in an aperture statistic. Figure 15 shows the correlation between $`D_{p,3}`$ and $`C`$ as computed in a $`40^{}\times 30^{}`$ simualted DEEP2 pointing. The saturation of $`C`$ in less-dense regions is striking and proves to be a significant drawback for a density estimator which is otherwise extremely robust to survey edges and redshift-space distortions.
The 3-dimensional $`n^{\mathrm{th}}`$-nearest neighbor distance and Voronoi volume statistics are the best suited for measuring under-dense environments. In groups and clusters, however, these density estimators are significantly affected by redshift-space distortions, causing each measure to become saturated. As illustarted in Figure 16, far removed from survey edges, the Voronoi volume and $`D_5`$ measures agree very well over all environments observed in the DEEP2-selected mock catalogs. However, for the simulated DEEP2 survey data, the $`n^{\mathrm{th}}`$-nearest-neighbor distance is much more robust to edge effects and is less expensive to calculate.
For studies of environment at high redshift, including analysis in the DEEP2 survey, we conclude that among the environment measures tested the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance provides the most accurate estimate of local galaxy density over the broadest range of scales. For work in dense environments, the $`D_{p,n}`$ offers great robustness to redshift-space distortions and maintains a reasonably high level of accuracy in under-dense environments. While $`D_{p,n}`$ can be affected by survey edges, contamination from geometric distortions is easily understood and effectively minimized without dramatically reducing the galaxy sample.
## 10. Conclusions
We have studied the applicability of several galaxy-density estimators within deep redshift surveys at $`z1`$ utilizing the mock galaxy catalogs of Yan et al. (2004). We conclude as follows:
1. Photometric redshifts derived from multi-band photometry $`(\sigma _z0.02)`$ are not suitable for measuring galaxy densities. Current photometric redshift surveys such as COMBO-17 do not have the redshift precision needed to study environment at high redshift. While more costly to obtain, spectroscopic redshifts are requisite to accurately probe the local galaxy environment in a large survey.
2. With the exception of the counts in an aperture estimator, $`C`$, survey field edges are a major source of contamination for each environment measure tested. To reduce these edge effects without biasing the sample, all galaxies within some comoving distance ($`12h^1`$ comoving Mpc for DEEP2) of a transverse survey edge should be rejected. At $`z1`$, excluding all galaxies within $`1h^1`$ comoving Mpc of an edge removes roughly 0.05 degrees along each dimension of the survey field. For smaller high-redshift surveys, such as CFRS (Lilly et al., 1995) or CNOC2 (Yee et al., 2000), edge effects introduce contamination to a considerable portion of the survey data set, thereby limiting the statistical power of the samples. Likewise, for a survey of the GOODS-North field (Giavalisco et al., 2004), edge effects would bias density measurements over $`75\%`$ of the field. Testing each environment measure on a simulated DEEP2-selected mock sample $`(40^{}\times 30^{})`$, the Voronoi volume is most severely affected by edges, with more than 2 times as much contamination from edge effects than $`D_n`$ or $`D_{p,n}`$. The counts in an aperture measure displays the best behavior near edges of a survey field, with a nearly negligible portion of the sample contaminated in our simulations.
3. Redshift-space distortions are a significant and fundamental roadblock to measuring accurate galaxy densities in over-dense environments. The $`n^{\mathrm{th}}`$-nearest-neighbor distance measured in 3-dimensions and the Voronoi volume are most greatly affected, while estimators such as projected $`n^{\mathrm{th}}`$-nearest-neighbor distance and counts in an aperture – which smooth the galaxy distribution along the line-of-sight – are less affected by the “fingers-of-god”. Still, it should be noted that less than 15% of a simulated $`R_{\mathrm{AB}}24.1`$ galaxy sample occupies environments at which the $`V`$ and $`D_n`$ statistics saturate due to redshift-space distortions.
4. The target selection algorithm employed by a survey could lead to environment-dependent biases in the observed galaxy sample. The DEEP2 survey, which slightly under-samples dense regions of sky, is equally sensitive at high and low densities. That is, we find that the DEEP2 survey equally samples all environments at $`z1`$ (see Figure 12). Also, we find that while the DEEP2 survey samples only $`50\%`$ of galaxies at $`z1`$, this uniform incompleteness simply introduces a random scatter in the measured environments and does not introduce an environment-dependence bias.
5. In examining the evolution of galaxy environments as a function of redshift, the estimation of the survey selection function plays a critical role. Uncertainties in the comoving number density of galaxies at high $`z`$ make comparisons over large redshift intervals $`(\mathrm{\Delta }z0.5)`$ problematic. Apart from such ambiguities, simple empirical corrections for densities as a function of redshift are highly effective.
6. For the DEEP2 Galaxy Redshift Survey, the projected $`n^{\mathrm{th}}`$-nearest-neighbor distance provides the most accuracte estimate of local galaxy density over a continuous and broad range of scales. The $`D_{p,n}`$ statistic is reasonably robust to redshift space distortions and still effective at tracing galaxy environments in under-dense regions.
7. Among current data sets at high redshift, we find the DEEP2 Galaxy Redshift Survey provides the best opportunity for measuring accurate galaxy environments over a broad and continuous range of scales. The high sampling rate and excellent redshift-precision of DEEP2 enable environments to be measured in even the most over-dense regions and yield improved accuracy over other deep surveys. Furthermore, DEEP2’s high-precision redshifts and large survey area (3.5 square degrees) minimize the effects of edges in both the transverse and line-of-sight directions.
We wish to thank Chris Marinoni for providing his Voronoi-Delaunay method group-finding code. This work was supported in part by NSF grant AST00-71048. JAN and DSM acknowledge support by NASA through Hubble Fellowship grants HST-HF-01165.01-A and HST-HF-01163.01-A, respectively, awarded by the Space Telescope Science Institute, which is operated by AURA Inc. under NASA contract NAS 5-26555. BFG acknowledges support from a NSF Fellowship. MCC thanks Mike Blanton for useful discussions about this work. MCC also thanks Josh Simon and Alison Coil for careful reading of this manuscript and many insightful suggestions which have improved this work. Lastly, the authors are humbly indebted to Steve Dawson for his invaluable assistance with various technical aspects of this effort.
|
warning/0506/quant-ph0506175.html
|
ar5iv
|
text
|
# Noise-free high-efficiency photon-number-resolving detectors
(Received 23 March 2005; revised 15 April 2005; published 17 June 2005)
## Abstract
High-efficiency optical detectors that can determine the number of photons in a pulse of monochromatic light have applications in a variety of physics studies, including post-selection-based entanglement protocols for linear optics quantum computing and experiments that simultaneously close the detection and communication loopholes of Bell’s inequalities. Here we report on our demonstration of fiber-coupled, noise-free, photon-number-resolving transition-edge sensors with $`88\%`$ efficiency at $`1550`$ nm. The efficiency of these sensors could be made even higher at any wavelength in the visible and near-infrared spectrum without resulting in a higher dark-count rate or degraded photon-number resolution.
High-efficiency, photon-number-resolving detectors can transform the field of quantum optics. One of many experiments they can enable is linear optics quantum computing, which requires postselection based on photon number Knill et al. (2001). So far, researchers have succeeded in implementing simple two-qubit gates by using conventional detectors that cannot distinguish between one and two photons Gasparoni et al. (2004); O’Brien et al. (2003), but moving beyond two qubits requires high-efficiency, low dark-count rate, photon-number-resolving detectors. Such detectors can be used to herald multiphoton path-entangled states from the output of a parametric down-conversion crystal Eisenberg et al. (2005), and these entangled states could be used for applications ranging from quantum cryptography Ekert91 ; BBM92 to lithography beyond the diffraction limit Boto et al. (2000)<sup>1</sup><sup>1</sup>1In order to be practical, quantum lithography requires the development of a bright source of entangled photons.. Photon-number-resolving detectors can also be used to verify the quality of single-photon sources, a necessity for secure information transfer in some quantum key distribution protocols. Furthermore, low-noise, high-efficiency detectors that operate at telecommunication wavelengths can significantly extend the length of a secure link in fiber quantum key distribution where light must be transmitted over large distances Lutkenhaus (1999).
Conventional detectors that operate in the visible and near-infrared, such as avalanche photodiodes and photomultiplier tubes, may be single-photon sensitive, but they cannot reliably determine the number of photons in a pulse of lightOldham et al. (1972); Engstrom (1980); Levine and Bethea (1984). In principle, beam-splitters and single-photon sensitive detectors can simulate a photon-number resolving detector Achilles et al. (2004), but the probability of correctly identifying an $`N`$-photon event drops exponentially with $`N`$, even for detectors with $`100\%`$ efficiency, because it is impossible to control the path each photon takes at a beamsplitter. Novel technologies such as the visible light photon counter Waks et al. (2003) have some photon-number resolution ability, but operating them at maximum detection efficiency introduces dark counts at rates greater than 10 kHz, and these detectors are sensitive in the visible spectrum only.
Superconducting transition-edge sensors (TESs) have photon-number resolution with negligibly-low dark counts. The TESs discussed in this paper are quantum calorimeters optimized for detection of near-infrared and visible photons Irwin (1995); Cabrera et al. (1998); Miller et al. (2003). The main components of a calorimeter are the absorber, a thermometer, and a weak link to a thermal heat sink, as shown in Fig. 1. When energy impinges on the absorber, it heats up quickly and slowly cools through the weak thermal link, and the temperature change is measured by the thermometer. Detection of visible and near-infrared light at the single-photon level places stringent requirements on the heat capacity and thermometry. For the TESs described here, the electron subsystem in a thin film of tungsten plays the parts of both the absorber and thermometer. The detector is cooled below its superconducting transition temperature and a voltage bias is applied to increase the electron temperature above that of the substrate. At low temperatures, the electrons in tungsten have anomalously low thermal coupling to the phonons, providing the weak thermal link, and the rapid change in resistance near the superconducting critical temperature results in a very sensitive measure of temperature. The temperature change due to energy deposition by a photon results in a change in resistance, and the current change in the voltage-biased detector is measured with a superconducting quantum-interference device (SQUID) array Huber et al. (2001). The change in temperature (and thus current) is proportional to the photon energy, so the sensor can resolve the number of photons in a pulse of monochromatic light.
The detection efficiency of a bare thin film tungsten sensor $`20\mathrm{nm}`$ thick is $`15`$ to $`20\%`$ at visible and near-infrared wavelengths and is limited by reflection from the front surface and transmission through the film. However, every photon that is absorbed by the tungsten leads to a change in temperature of the electrons, so the detection probability can be increased by embedding the tungsten detector in a stack of optical elements that enhance the absorption of the light in the tungsten. The TESs discussed in this letter measure $`25\mu \mathrm{m}`$ by $`25\mu \mathrm{m}`$ and are approximately $`20\mathrm{nm}`$ thick with superconducting critical temperatures of $`110\pm 5\mathrm{mK}`$. They are embedded in structures that are designed to maximize absorption at $`1550\mathrm{nm}`$, a wavelength of particular interest for telecommunications. Existing semiconductor-based detectors have low ($`1020\%`$) efficiency and high ($`1020\mathrm{kHz}`$) dark-count rates at this wavelength. A detailed description of the structures is presented elsewhere Rosenberg et al. (June 2005). These sensors have thermal decay times as short as $`5\mu s`$ and provide excellent discrimination between multi-photon events. Figure 2 shows data from a TES embedded in an optical structure designed to enhance the absorption of light at a wavelength of $`1550`$ nm. The sensor was illuminated with a pulsed source of $`1550\mathrm{nm}`$ photons. The histogram displays the distribution of pulse heights after the data were corrected for the non-linearity in the temperature dependence of the resistance in the superconducting transition. The full width at half maximum (FWHM) of the zero-, one-, two-, three- and four-photon peaks are $`0.13`$ eV, $`0.20`$ eV, $`0.25`$ eV, $`0.34`$ eV, and $`0.45`$ eV, respectively. The increase in the FWHM of the peaks with increasing energy is due to the device non-linearity mentioned above. Measurements and simulations of the optical properties of the various layers indicate that the total expected efficiency of the sensor, neglecting system losses, is $`92\%`$.
Photons were coupled to the detector through $`9\mu \mathrm{m}`$ core single-mode fiber with an anti-reflective coating for $`1550\mathrm{nm}`$. The fiber was held $`5075\mu \mathrm{m}`$ above the detector and aligned at room temperature by backside through-chip imaging. Focusing the light from the fiber was not necessary because the spot size was small enough at this distance that greater than $`99\%`$ of the light was incident on the detector. The housing holding the fiber was clamped in place and the detector was then cooled to less than $`100\mathrm{mK}`$ in an adiabatic demagnetization refrigerator. Because the applied voltage bias keeps the electrons in the superconducting-to-normal transition, the detector is not sensitive to slight fluctuations in the cryostat temperature as long as the temperature is well below the superconducting transition temperature of $`110\pm 5\mathrm{mK}`$.
Coupling and alignment losses reduce the measured efficiency of the detector from the expected $`92\%`$. To minimize connection losses, the fiber from the detector and the fiber going to room temperature were fused together in the cold space of the cryostat. The typical loss for a fiber fuse is approximately $`0.5\%`$. We measured the room temperature loss from outside the cryostat to the sample space to be $`2.3\%`$. Tests to determine the loss in a loop of fiber that passes through the cold space of the refrigerator indicated that the loss does not change when the fiber is cooled. Thermal cycling of the fiber-coupled detector did not change its efficiency, and we measured greater than $`80\%`$ efficiency for several different detectors, indicating that our alignment method is robust and that the fiber-to-detector alignment does not degrade when cooled.
Measuring the efficiency of the detectors is nontrivial due to the low power levels involved and the introduction of loss through fiber connectors. The relatively slow pulse decay (several microseconds) and the desire to avoid pulse pile-up requires the use of subfemtowatt average optical power levels. At present, commercial power meters do not have the sensitivity required to measure such low levels. To circumvent this problem, we calibrated a series of programmable optical attenuators using a calibrated power meter well within its linear regime, as shown in Fig. 3. The attenuator calibration and efficiency measurements were performed using a laser with a center wavelength of $`1550\mathrm{nm}`$ and FWHM of $`0.05\mathrm{nm}`$.
The efficiency measurements presented here were performed in continuous-wave operation at several different power levels to ensure that there was no dependence of measured efficiency on power level, as shown in Fig. 4. Power levels were adjusted by means of the calibrated programmable attenuators. The number of pulses with pulse heights within $`\pm 3\sigma `$ ($`\sigma 0.07\mathrm{eV})`$ of the one-photon peak was recorded at each power level for $`100\mathrm{s}`$. The background rate in the same energy range, which was approximately $`400\mathrm{Hz}`$ and was due to blackbody radiation from room temperature surfaces A. J. Miller et al., unpublished <sup>1</sup><sup>1</sup>1Note that the background counts were not from detector noise; rather, they resulted from photons from the low-energy tail of the blackbody distribution propagating through the optical fiber., was measured periodically between data sets and subtracted from the raw count rate. The procedure outlined above provides a measure of the number of single-photon events from the laser. At the higher power levels, counts were present at energies greater than $`3\sigma `$ above the mean energy of the one-photon peak. The spectral weight at these energies results from pulse pile-up, the arrival of a second photon before the system has recovered from a previous event. To correct for this effect, we included the high-energy counts as two-photon events. This correction is smallest ($`1\%`$ of the total counts) at the lowest power levels and remains below $`2\%`$ at even the highest power level.
Fig. 4 shows the detection efficiency as a function of power level, with error bars given by the uncertainties due to Poisson statistics. Not shown are the uncertainties resulting from fiber bends at room temperature. We have observed that small bends in the fiber can easily lower the measured efficiency by up to three per cent, and the scatter in efficiency measurements is most likely due to slight shifts in the fiber position. All the efficiency values presented are relative to the amount of light in the fiber at point A in Fig. 3.
The measured system efficiency of $`88.6\pm 0.4\%`$ is consistent with measurements and simulations of the optical elements and the system losses. This detection efficiency exceeds the threshold of $`83\%`$ required to close the detection loophole in an experiment testing Bell’s inequalities. This enables an experiment that would simultaneously close both the detection and communication loopholes, decisively refuting a local realism interpretation of quantum mechanics.
Increasing the detection efficiency beyond $`88.6\%`$ at $`1550\mathrm{nm}`$ is in principle simple and involves fabricating an optical structure with more layers and finer control over layer thickness. Similarly, it should be possible to produce near unity-efficiency detectors at any wavelength in the ultraviolet to near-infrared frequency range with this technique. Simulations indicate the possibility of increasing the efficiency well above $`99\%`$ at any given wavelength in this spectrum, making these detectors an extremely valuable tool for quantum optics and quantum information processing.
###### Acknowledgements.
The authors thank ARDA for financial support, Alan Migdall, Richard Mirin, John Martinis, Alexander Sergienko and Erich Grossman for valuable technical discussions, and Marty Gould of Zen Machining. D. R. is supported by the DCI postdoctoral program.
|
warning/0506/quant-ph0506023.html
|
ar5iv
|
text
|
# Operator Quantum Error Correcting Subsystems for Self-Correcting Quantum Memories
## I Introduction
In the early days of quantum computation, the analog nature of quantum information and quantum transforms, as well as the effect of noise processes on quantum systems, were thought to pose severe obstaclesUnruh (1995); Landauer (1996) towards the experimental realization of the exponential speedups promised by quantum computers over classical computersBernstein and Vazirani (1993); Simon (1994); Shor (1994); Raz (1999). Soon, however, a remarkable theory of fault-tolerant quantum computationShor (1995); Steane (1996); Knill and Laflamme (1997); Bennett et al. (1996); Knill et al. (1998a, b); Shor (1996); Aharonov and Ben-Or (1997); Gottesman (1998); Kitaev (2003); Preskill (1998) emerged which dealt with these problems and showed that quantum computers are indeed more similar to probabilistic classical computers than to analog devices. Analog computers have a computational power which is dependent on a lack of noise and on exponential precision, whereas probabilistic classical computers can be error corrected and made effectively digital even in the presence of noise and non-exponential precision. The theory of fault-tolerant quantum computation establishes that quantum computers are truly digital devices deserving of the moniker computer. An essential idea in the development of the theory of fault-tolerant quantum computation was the notion that quantum information could be encoded into subspacesShor (1995); Steane (1996); Knill and Laflamme (1997); Bennett et al. (1996) (quantum error correcting codes) and thereafter protected from degradation via active procedures of detection and correction of errors. Encoding quantum information into subspaces, however, is not the most general method of encoding quantum information into a quantum system. The most general notion for encoding quantum information is to encode the information into a subsystem of the quantum systemKnill and Laflamme (1997); Knill et al. (2000); Viola et al. (2001). This has been perhaps best exploited in the theory of noiseless subsystemsKnill et al. (2000); Kempe et al. (2001); Bacon (2001); Zanardi (2001); Zanardi and Lloyd (2003); Shabani and Lidar (2005) and dynamic recoupling schemesViola et al. (2000); Zanardi (2001); Wu and Lidar (2002). Recently a very general notion of quantum error correction has appeared under the moniker of “operator quantum error correction.”Kribs et al. (2005a, b) In this work the possibility of encoding into subsystems for active error correction is explicitly examined. While it was foundKribs et al. (2005a, b) that the notion of a encoding into a subsystem does not lead to new codes (all subsystem codes could be thought of as arising from subspace codes), encoding into a subsystem does lead to different recovery procedures for quantum information which has been encoded into a subsystem. Hence operator quantum error correcting codes, while not offering the hope of more general codes, do offer the possibility of new quantum error correcting routines, and in particular to the possibility of codes which might help improve the threshold for fault-tolerant quantum computation due to the lessened complexity of the error correcting routine.
In this paper we present two examples of operator quantum error correcting codes which use subsystem encodings. The codes we present have the interesting property that the recovery routine does not restore information encoded into a subspace, but recovers the information encoded into a subsystem. Using the $`[n,d,k]`$ labelling a quantum error correcting codes, where $`n`$ is the number of qubits used in the code, $`d`$ is the distance of the code, and $`k`$ is the number of encoded qubits for the code, our codes are $`[n^2,n,1]`$ and $`[n^3,n,1]`$ quantum error correcting codes. The subsystem structure of our codes is explicitly exploited in the recovery routine for the code, and because of this they are substantially simpler than any subspace code derived from these codes.
While the two codes we present are interesting in there own right, there is a further motivation for these codes above and beyond their exploitation of the subsystem structure in the recovery routine. The two operator quantum error correcting subsystems we present are motivated by two interesting Hamiltonians defined on two and three-dimensional square and cubic lattices of qubits with certain anisotropic spin-spin interactionsBacon (2001). The three-dimensional version of this system is particularly intriguing since it offers the possibility of being a self-correcting quantum memory. In a self-correcting quantum memory, quantum error correction is enacted not by the external control of a complicated quantum error correction scheme, but instead by the natural physics of the device. Such a quantum memory offers the possibility for removing the need for a quantum microarchitecture to perform quantum error correction and could therefore profoundly speed up the process of building a quantum computer. In this paper we present evidence, in the form of a simple mean field argument, that the three-dimensional system we consider is a self-correcting quantum memory. We also show that the operator error correcting subsystem structure of this code is an important component to not only the self-correcting properties of this system, but also to encoding and decoding information in this system.
The organization of the paper is as follows. In Section II we review the notion of encoding information into a subsystem and discuss the various ways in which this has been applied to noiseless subsystems and dynamic recoupling methods for protecting quantum information. Next, in Section III, we discuss how operator error correcting subsystems work and how they differ from standard quantum error correcting codes. Our first example of a operator quantum error correcting subsystem is presented in Section IV where we introduce an example on a square lattice. The second, and more interesting, example of a operator quantum error correcting subsystem is given in Section V where we discuss an example on a cubic lattice. In Section VI we introduction the notion of a self-correcting quantum memory and present arguments that a particular Hamiltonian on a cubic lattice related to our cubic lattice subsystem is self-correcting. We conclude in Section VII with a discussion of open problems and the prospects for operator quantum error correcting subsystems and self-correcting quantum memories.
## II Subsystem Encoding
Consider two qubits. The Hilbert space of these qubits is given by $`^2^2`$. Pick some fiducial basis for each qubit labelled by $`|0`$ and $`|1`$. One way to encode a single qubit of information into these two qubits is to encode the information into a subspace of the joint system. For example, we can define the logical basis states $`|0_L=\frac{1}{\sqrt{2}}(|01|10)`$ and $`|1_L=|11`$ such that a single qubit of information can be encoded as $`\alpha |0_L+\beta |1_L=\frac{\alpha }{\sqrt{2}}(|01|10)+\beta |11`$. This is an example of idea of encoding quantum information into a subspace, in this case the subspace spanned by $`|0_L`$ and $`|1_L`$. But another way to encode a single qubit of information is to encode this information into one of the two qubits. In particular if we prepare the state $`|\psi (\alpha |0+\beta |1)`$ for an arbitrary single qubit state $`|\psi `$, then we have also encoded a single qubit of information in our system. This time, however, we have encoded in the information into a subsystem of the system. It is important to note that the subsystem encoding works for an arbitrary state $`|\psi `$. If we fix $`|\psi `$ to some known state, then we are again encoding into a subspace. We reserve the nomenclature of “encoding into a subsystem” to times in which $`|\psi `$ is arbitrary.
More generally, if we have some Hilbert space $``$, then a subsystem $`𝒞`$ is a Hilbert space arising from $``$ asKnill and Laflamme (1997); Knill et al. (2000)
$$=(𝒞𝒟).$$
(1)
Here we have taken our Hilbert space and partitioned it into two subspaces, $``$ and a subspace perpendicular to $``$. On this perpendicular subspaces, we have introduced a tensor product structure, $`𝒞𝒟`$. We can then encode information into the first subsystem $`𝒞`$. This can be achieved by preparing the quantum information we wish to encode $`\rho _C`$ into the first subsystem, $`𝒞`$ along with any arbitrary state $`\rho _D`$ into the second subsystem $`𝒟`$:
$$\rho =(\rho _C\rho _D)0.$$
(2)
The fact that quantum information can most generally be encoded into a subsystem was an essential insight used in the construction of noiseless (decoherence-free) subsystemsKnill et al. (2000); Kempe et al. (2001); Bacon (2001); Zanardi (2001); Zanardi and Lloyd (2003); Shabani and Lidar (2005). Suppose we have a system with Hilbert space $`_S`$ and an environment with Hilbert space $`_E`$. The coupling between these two systems will be described by an interaction Hamiltonian $`H_{int}`$ which acts on the tensor product of these two spaces $`_S_B`$. The idea of a noiseless subsystem is that it is often the case that there is a symmetry of the system-environment interaction such that the action of the interaction Hamiltonian factors with respect to some subsystem structure on the system’s Hilbert space,
$$H_{int}=\underset{\alpha }{}\left[\left(I_dD_\alpha \right)E_\alpha \right]B_\alpha ,$$
(3)
where $`I_d`$ is the d-dimensional identity operator acting on the subsystem code space $`𝒞`$, $`D_\alpha `$ acts on the subsystem $`𝒟`$, $`E_\alpha `$ acts on the orthogonal subspace $``$, and $`B_\alpha `$ operates on the environment Hilbert space $`_E`$. When our interaction Hamiltonian possesses a symmetry leading to such a structure, then, if we encode quantum information into $`𝒞`$, this information will not be effected by the system-environment coupling. Thus information encoded in such a subsystem will be protected from the effect of decoherence and hence exists in a noiseless subsystem. Noiseless subsystems were a generalization of decoherence-free subspacesZanardi and Rasetti (1997); Lidar et al. (1998), this latter idea occurring when the subsystem structure is not exploited, $`𝒟=`$, and then encoding quantum information is simply encoding quantum information into a subspace. Subsystems have also been used in dynamic recoupling techniquesViola et al. (2000); Zanardi (2001); Wu and Lidar (2002) where symmetries are produced by an active symmetrization of the system’s component of the system-environment evolution.
## III Operator Quantum Error Correcting Subsystems
Here we examine the implications of encoding information into a subsystem for quantum error correcting protocolsKribs et al. (2005a, b). Suppose that we encode quantum information into a subsystem $`𝒞`$ of some quantum system with full Hilbert space $`=(𝒞𝒟)`$. Now suppose some quantum operation (corresponding to an error) occurs on our system. Following the standard quantum error correcting paradigm, we then apply a recovery procedure to the system. When $`𝒟=`$, i.e. when we are encoding into a quantum error correcting subspace, then the quantum error correcting criteria is simply that the effect of the error process followed by the recovery operation should act as identity on this subspace. If we encode information into a subsystem, however, this criteria is changed to only requiring that the recovery operation should act as identity on the subsystem $`𝒞`$. In particular we do not care if the effect of an error followed by our recovery procedure enacts some nontrivial procedure on the $`𝒟`$ subsystem. In fact our error correcting procedure may induce some nontrivial action on the $`𝒟`$ subsystem in the process of restoring information encoded in the $`𝒞`$ subsystem.
How does the above observation modify the basic theory of quantum error correcting codes? In standard quantum error correction, we encode into some error correcting subspace with basis $`|i`$. The necessary and sufficient condition for there to be a procedure under which quantum information can be restored under a given set of errors $`E_a`$ is given by
$$i|E_a^{}E_b|j=\delta _{i,j}c_{a,b}.$$
(4)
For the case of encoding into a subsystem this necessary and sufficient condition is modified as follows. Let $`|i|k`$ denote a basis for the subspace $`𝒞𝒟`$. Then Kribs et al.Kribs et al. (2005a, b) showed a necessary conditionKnill and Laflamme (1997); Bennett et al. (1996) for the quantum error correcting is given by
$$\left(i|k|\right)E_a^{}E_b\left(|j|l\right)=\delta _{i,j}m_{a,b,k,l}.$$
(5)
That this condition is also sufficient has recently also been shown by Nielsen and PoulinNielsen and Poulin (2005). As noted in Kribs et al. (2005a, b), a code constructed from the subsystem operator quantum error correcting criteria can always be used to construct a subspace code which satisfies the subspace criteria Eq. (4). We note, however, that while this implies that the notion of using subsystems for quantum error correction does not lead to new quantum error correcting codes above and beyond subspace encodings, the codes constructed which exploit the subsystem structure have error recovery routines which are distinct from those which arise when encoding into a subspace. In particular, when one encodes into a subsystem, the recovery routine does not need to fix errors which occur on other subsystems. Below we will present examples of subsystem encodings in which the subsystem structure of the encoding is essential not for the existence of the quantum error correcting properties, but it essential for the simple recovery routine we present.
## IV Two-Dimensional Operator Quantum Error Correcting Subsystem
Here we construct an operator quantum error correcting subsystem for a code which lives on a two-dimensional square lattice. This code makes explicit use of the subsystem structure in its error recovery procedure. A familiarity with the stabilizer formalism for quantum error correcting codes is assumed (see Gottesman (1997); Nielsen and Chuang (2000) for overviews.)
### IV.1 Preliminary Definitions
Consider a square lattice of size $`n\times n`$ with qubits located at the vertices of this lattice. Let $`O_{i,j}`$ denote the operator $`O`$ acting on the qubit located at the $`i`$th row and $`j`$th column of this lattice tensored with identity on all other qubits. Recall that the Pauli operators on a single qubit are
$`X=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],iY=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],\mathrm{and}Z=\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right].`$ (12)
It is convenient to use a compact notation to denote Pauli operator on our $`n^2`$ qubits by using two $`n^2`$ bit strings,
$$P(a,b)=\underset{i,j=1}{\overset{n}{}}X_{i,j}^{a_{i,j}}Z_{i,j}^{b_{i,j}}=\underset{i,j=1}{\overset{n}{}}\{\begin{array}{cc}X_{i,j}\hfill & \mathrm{if}a_{i,j}=1\mathrm{and}b_{i,j}=0\hfill \\ Z_{i,j}\hfill & \mathrm{if}a_{i,j}=0\mathrm{and}b_{i,j}=1\hfill \\ iY_{i,j}\hfill & \mathrm{if}a_{i,j}=1\mathrm{and}b_{i,j}=1\hfill \end{array},$$
(13)
where $`a,b_2^{n^2}`$ are $`n`$ by $`n`$ matrices of bits. Together with a phase, $`i^\varphi `$, $`\varphi _4`$, a generic element of the Pauli group on our $`n^2`$ qubits is given by $`i^\varphi P(a,b)`$. We will often refer to a Pauli operator as being made up of $`X`$ and $`Z`$ operators, noting that when both appear, the actual Pauli operator is the $`iY`$ operator.
We begin by defining three sets of operators which are essential to understanding the subsystem structure of our qubits. Each of these sets will be made up of Pauli operators.
The first set of Pauli operators which will concern us, $`𝒯`$, is made up of Pauli operators which have an even number of $`X_{i,j}`$ operators in each column and an even number of $`Z_{i,j}`$ operators in each row:
$`𝒯=\left\{(1)^\varphi P(a,b)\right|\varphi _2,{\displaystyle \underset{i=1}{\overset{n}{}}}a_{i,j}=0,\mathrm{and}{\displaystyle \underset{j=1}{\overset{n}{}}}b_{i,j}=0\},`$ (14)
where $``$ denotes the binary exclusive-or operation (we use it interchangeably with the direct sum operation with context distinguishing the two uses.) Note that these operators form a group under multiplication. This group can be generated by nearest neighbor operators on our cubic lattice
$$𝒯=X_{i,j}X_{i+1,j},Z_{j,i}Z_{j,i+1},i_{n1},j_n.$$
(15)
Examples of elements of the group $`𝒯`$ are diagrammed in Fig. 1.
The second set of Pauli operators we will be interested in is a subset of $`𝒯`$, which we denote $`𝒮`$. $`𝒮`$ consists of Pauli operators which are made up of an even number of rows consisting entirely of $`X`$ operators and an even number of columns consisting entirely of $`Z`$ operators,
$$𝒮=\left\{P(a,b)\right|\underset{i=1}{\overset{n}{}}\left(\underset{j=1}{\overset{n}{}}a_{i,j}\right)=0,\underset{j=1}{\overset{n}{}}\left(\underset{i=1}{\overset{n}{}}b_{i,j}\right)=0\}$$
(16)
where $``$ is the binary and operation. $`𝒮`$ is also a group. In fact it is an Abelian subgroup of $`𝒯`$. Further all of the elements of $`𝒮`$ commute not just with each other, but with all of the elements of $`𝒯`$. It can be generated by nearest row and column operators,
$$𝒮=\underset{i=1}{\overset{n}{}}X_{j,i}X_{j+1,i},\underset{i=1}{\overset{n}{}}Z_{i,j}Z_{i,j+1},j_{n1}.$$
(17)
These generators will be particularly important for us, so we will denote them by
$`S_i^X={\displaystyle \underset{j=1}{\overset{n}{}}}X_{i,j}X_{i+1,j},\mathrm{and}S_j^Z={\displaystyle \underset{i=1}{\overset{n}{}}}Z_{i,j}Z_{i,j+1}.`$ (18)
$`𝒮`$ is a stabilizer group familiar from the standard theory of quantum error correcting codes. An example of an element in $`𝒮`$ is given in Fig. 2
The final set of operators which we will consider, $``$, is similar to $`𝒮`$ except that the evenness condition becomes an oddness condition,
$`=\left\{(1)^\varphi P(a,b)\right|\varphi _2,{\displaystyle \underset{i=1}{\overset{n}{}}}\left({\displaystyle \underset{j=1}{\overset{n}{}}}a_{i,j}\right)=1,{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \underset{i=1}{\overset{n}{}}}b_{i,j}\right)=1\}.`$ (19)
This set does not by itself form a group, but together with $`𝒮`$ it does form a group. This combined group is not Abelian. $``$ has the property that all of its elements commute with those of $`𝒯`$ and $`𝒮`$. A nontrivial element of $``$ is given in Fig. 3.
### IV.2 Subsystem Structure
We will now elucidate how $`𝒯`$, $`𝒮`$, and $``$ are related to a subsystem structure on our $`n^2`$ qubits. Let $`=(^2)^{n^2}`$ denote the Hilbert space of our $`n^2`$ qubits.
We first note that since $`𝒮`$ consists of a set of mutually commuting observables, we can use these observables to label subspaces of $``$. In particular we can label these subspaces by the $`2(n1)`$ $`\pm 1`$-valued eigenvalues of the $`S_i^X`$ and $`S_i^Z`$ operators. Let us denote these eigenvalues by $`s_i^X`$ and $`s_i^Z`$ respectively and the length $`n1`$ string of these $`\pm 1`$ eigenvalues by $`s^X`$ and $`s^Z`$. We can thus decompose $``$ into subspaces as
$$=\underset{s_1^X,\mathrm{},s_{n1}^X,s_1^Z,\mathrm{},s_{n1}^Z=\pm 1}{}_{s^X,s^Z}=\underset{s^X,s^Z}{}_{s^X,s^Z}.$$
(20)
By standard arguments in the stabilizer formalism, each of the $`_{s^X,s_Z}`$ subspaces is of dimension $`d=2^{n^22(n1)}`$. Just for completeness, we note that the operators $`S_i^X`$ and $`S_i^Z`$ act under this decomposition as
$`S_i^X`$ $`=`$ $`{\displaystyle \underset{s^X,s^Z}{}}s_i^XI_{2^{n^22(n1)}},`$
$`S_i^Z`$ $`=`$ $`{\displaystyle \underset{s^X,s^Z}{}}s_i^ZI_{2^{n^22(n1)}}.`$ (21)
Now examine the two groups $`𝒯`$ and the group generated by elements of $``$ and $`𝒮`$. Both of these groups are non-Abelian. All of the elements of $`𝒯`$ and $``$ commute with elements of $`𝒮`$. Further, all of the elements of $`𝒯`$ and $``$ commute with each other. This implies, via Schur’s lemmaCornwell (1997); James and Liebeck (2001), that $``$ and $`𝒯`$ must be represented on $`_{s^X,s^Z}`$ by a subsystem action. In particular the full Hilbert space splits as
$$=\underset{s^X,s^Z}{}_{s^X,s^Z}^𝒯_{s^X,s^Z}^{},$$
(22)
such that operators from $`T𝒯`$ act on the first tensor product
$$T=\underset{s^X,s^Z}{}T_{s^X,s^Z}I_2,T𝒯,$$
(23)
and the operators from $`L`$ act on the second tensor product
$$L=\underset{s^X,s^Z}{}I_{2^{(n1)^2}}L_{s^X,s^Z},L.$$
(24)
Here we have assigned dimensions $`2^{(n1)^2}`$ and $`2`$ to these tensor product spaces. To see why these dimensionalities arise we appeal to the stabilizer formalism. We note that modulo the stabilizer structure of $`𝒮`$, $``$ is a single encoded qubit. Similarly if one examines the following set of $`(n1)^2`$ operators from $`𝒯`$,
$`\overline{Z}_{i,j}=Z_{i,j}Z_{i,j+1},\overline{X}_{i,j}={\displaystyle \underset{k=1}{\overset{j}{}}}X_{i,k}X_{n1,k},`$ (25)
where $`i_{n1},j_{n1}`$, one finds that modulo the stabilizer they are equivalent to $`(n1)^2`$ encoded Pauli operators.
The subsystem code we now propose encodes a single qubit into the Hilbert space $`_{s^X,s^Z}^{}`$ with $`s_i^X=s_i^X=+1,i_{n1}`$ (choices with other $`\pm 1`$ choices form an equivalent code in the same way that stabilizer codes can be chosen for different stabilizer generator eigenvalues.) The code we propose thus encodes one qubit of quantum information into a subsystem of the $`n^2`$ “bare” qubits. We stress that the encoding we perform is truly a subsystem encoding: we do not care what the state of the $`_{s^X,s^Z}^𝒯`$ subsystem is. For simplicity it may be possible to begin by encoding into a subspace which includes our particular subsystem (i.e. by fixing the state on $`_{s^X,s^Z}^𝒯`$) but this encoding is not necessary and indeed, after our recover routine for the information encoded into $`_{s^X,s^Z}^{}`$, we will not know the state of the $`_{s^X,s^Z}^𝒯`$ subsystem. We will denote the Hilbert space $`_{s^X,s^Z}^{}`$ with $`s_i^X=s_i^X=+1,i_{n1}`$ by $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$.
### IV.3 Subsystem Error Correcting Procedure
If we encode quantum information into the subsystem $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$, then what sort of error correcting properties does this encoding result in? We will see that the $`𝒮`$ operators can be used to perform an error correcting procedure which restores the information on $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$, but which often acts nontrivially on the subsystem $`_{s^X=s^Z=\{+1\}^{n1}}^𝒯`$. This exploitation of the subsystem structure in the correction procedure is what distinguishes our subsystem operator quantum error correcting code from standard subspace quantum error correcting codes.
Suppose that a Pauli error $`P(a,b)`$ occurs on our system. For a Pauli operator $`P(a,b)`$, define the following error strings:
$$e_j(a)=\underset{i=1}{\overset{n}{}}a_{i,j},f_i(b)=\underset{j=1}{\overset{n}{}}b_{i,j}.$$
(26)
Notice that if $`e_j=f_i=0,i,j`$, then this implies that $`P(a,b)`$ is in the set $`𝒯`$. Further note that in this case, the effect of $`P(a,b)`$ is to only act on the $`_{s^X,s^Z}^𝒯`$ subsystems, i.e. $`P(a,b)`$ is block diagonal under our subsystem decomposition, Eq. (22), acting as
$$P(a,b)=\underset{s^X,s^Z}{}E_{s^X,s^Z}(a,b)I_2.$$
(27)
where $`E_{s^X,s^Z}(a,b)`$ is a nontrivial operator depending on the subspace labels $`s^X,s^Z`$ and the type of Pauli error $`(a,b)`$. Therefore errors of this form ($`e_j=f_i=0`$) do not cause errors on our information encoded in $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$. With respect to the errors of this form, the information is encoded into a noiseless subsystemKnill et al. (2000).
Returning now to the case of a general $`P(a,b)`$, from the above argument we see that if we can apply a Pauli operator $`Q(c,d)`$ such that $`Q(c,d)P(a,b)`$ is a new error, call it $`R(a^{},b^{})`$, which has error strings $`e_j^{}(a^{})=f_i^{}(b^{})=0,i,j`$, then we will have a procedure for fixing the error $`P(a,b)`$, modulo the subsystem structure of our encoded quantum information. In other words, our error correcting procedure need not result in producing the identity action on the subspace labelled by $`s_i^X=s_i^Z=+1,i_{n1}`$, but need only produce the identity action on the subsystem $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$. We can perform just such a procedure by using the elements of $`𝒮`$ as a syndrome for which errors of small enough size can be corrected.
To see how this works, suppose $`P(a,b)`$ occurs on our system. Then note that measuring $`S_i^X`$ is equivalent to determining
$$\underset{j=1}{\overset{n}{}}(b_{i,j}b_{i+1,j})=f_i(b)f_{i+1}(b),$$
(28)
and similarly measuring $`S_j^Z`$ is equivalent to determining
$$\underset{i=1}{\overset{n}{}}(a_{i,j}a_{i,j+1})=e_j(a)e_{j+1}(a).$$
(29)
Note that all $`2(n1)`$ of these measurements can be performed simultaneously since the elements of $`𝒮`$ all commute with each other. We wish to use these measurement outcomes to restore the system to $`e_j(a)=f_i(b)=0`$ (if possible.)
To see how to do this, treat the $`f_i(b)`$ as a $`n`$ bit codeword for a simple redundancy code (i.e. the two codewords are $`f_i(b)=0,i`$ and $`f_i(b)=1,i`$). A similar procedure will hold for the $`e_i(b)`$. Measuring the $`n1`$ operators $`S_i^X`$ is equivalent to measuring the syndrome of our redundancy code. In particular we can use this syndrome to apply an error correcting procedure for the $`f_i(b)`$ bit strings. The result of this correction procedure is to restore the system to either the codeword $`f_i(b)=0,i`$ or the codeword $`f_i(b)=1,i`$. The former corresponds to an error correction procedure which can succeed (given that an equivalent procedure for the $`e_j(a)`$ bit strings also succeeds), whereas the latter procedure is one where the error correction procedure will fail. Notice that our error correcting procedure, when it succeeds, is only guaranteed to restore the system to $`f_i(b)=0`$ and $`e_j(b)=0`$, and thus the full effect of the procedure may be to apply some nontrivial operator to the $`_{s^X,s^Z}^𝒯`$ subsystems.
Let us be more detailed in describing the error correcting procedure for the $`f_i(b)`$ code words. Let $`s_i^X`$ be the result of our measurements of the $`S_i^X`$ operators. Give the $`s_i^X`$ we can construct two possible bit strings $`f_i^{}`$ and $`\neg f_i^{}`$ ($`\neg `$ denotes the negation operation) consistent with these measurements. Let $`H(f^{})`$ and $`H(\neg f^{})`$ denote the Hamming weight of these bit strings (i.e. $`H(f^{})`$ is the number of $`1`$s in the $`n`$ bits $`f_i^{}`$) and define $`f^{\prime \prime }`$ to be the bit string $`f^{}`$ or $`\neg f^{}`$ with the smallest Hamming weight. We now apply an operation consisting only of $`Z_{i,j}`$ operators. In particular we apply the operator
$$Q_1(f^{\prime \prime })=\underset{i=1}{\overset{n}{}}Z_{i,j_0}^{\delta _{f_i^{\prime \prime },1}},$$
(30)
for any fixed column index $`j_0`$. The operator $`Q_1(f^{\prime \prime })P(a,b)`$ is then seen to be of one of two forms: either this new operator has the $`Z`$ error string equal to all zeros or all ones. In the first case we have successfully restored the system to the all $`f_i(b)=0`$ codeword, whereas for the second case, we have failed. How many $`Z`$ errors can be corrected in this fashion? If $`P(a,b)`$ consisted of $`Z`$ errors $`b`$ with an error string $`f_i(b)`$ with a Hamming weight of this string $`H(f)`$ which is less than or equal to $`\frac{n1}{2}`$, then the correction procedure will succeed. Thus the code we have constructed is a $`[n^2,n,1]`$ code: it encodes a single qubit into $`n^2`$ qubits and has a distance $`n`$.
Above we have focused on the case of Pauli $`Z`$ errors. Clearly an analogous argument holds for Pauli $`X`$ errors (with the role of the rows and columns reversed.) Further, Pauli $`Y`$ errors are taken care of by the combined action of these two procedures. By the standard arguments of digitizing errors in quantum error correcting codes, we have thus shown how our operator quantum error correcting subsystem code can correct up to $`\frac{n1}{2}`$ arbitrary single qubit errors.
### IV.4 Logical Operators
We comment here on the logical operators (operators which act on the encoded subsystem) for this code. From our analysis of the subsystem structure, it is clear that elements of $``$ act on the subsystem. Thus, for instance, the effect of a row of Pauli $`X`$ operators is to enact an encoded Pauli $`X`$ operation on the coded subsystem. We can choose a labelling of the subsystem such that
$$\overline{X}=\underset{j=1}{\overset{n}{}}X_{1,j}=\underset{s^X,s^Z}{}I_{2^{(n1)^2}}X,$$
(31)
while in the same basis the effect of a column of Pauli $`Z`$ operators is to enact and an encoded Pauli $`Z`$ on the coded subsystem,
$$\overline{Z}=\underset{i=1}{\overset{n}{}}Z_{i,1}=\underset{s^X,s^Z}{}I_{2^{(n1)^2}}Z.$$
(32)
These two operators can then be used to enact any Pauli operator on the encoded quantum information. Notice that other elements of $``$ also act as encoded Pauli operators on $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$ (but act with differing signs on the other $`s^X,s^Z`$ labelled subspaces.) An important property of these logical operators is that they can be enacted by performing single qubit operators and no coupling between the different qubits is needed. This is important because it will allow us to assume an independent error model for error which occur when we imprecisely implement these gates on our encoded quantum information. Not only can the above construction be used to implement the Pauli operators on our subsystem code, it can also be used to measure the Pauli operators on our subsystem code.
Another easily implementable operation on our code is a logical controlled-not. Suppose we take two identically sized two-dimensional codes and stack them on top of each other. Then the application of a transverse controlled-not operator between all $`n^2`$ of these two systems will enact a logical controlled-not between the two encoded qubits. To see this note that if we treat the elements of the set $`𝒮`$ as a stabilizer code, then these transverse operators preserve the combined stabilizer $`𝒮\times 𝒮`$ and that the action of the $`n^2`$ controlled-not gates do not mix the $`\times `$ and $`𝒯\times 𝒯`$ operators.
GottesmanGottesman (1998, 1997) has shown that given the ability to measure and apply the encoded Pauli operators along with the ability to perform a controlled NOT on a stabilizer code, one can perform any encoded operation which is in the normalizer of the Pauli group (i.e. the gate set relevant to the Gottesman-Knill theoremGottesman (1997)).
We have seen how to implement encoded Pauli operators and the controlled-NOT on the information encoded into our subsystem. These operations do not allow for universal quantum computation, so an important open question for our subsystem code is to find an easily implementable method for completing this gateset to a universal set of gates.
### IV.5 Hamiltonian Model of the two-dimensional Subsystem Code
An interesting offshoot the above two-dimensional operator quantum error correcting subsystem is the analysis of a particularly simple Hamiltonian whose ground state has a degeneracy which corresponds to the subsystem code. We introduce this Hamiltonian here in order to make our analysis of a similar Hamiltonian for our three-dimensional subsystem code more transparent. The Hamiltonian is given by nearest neighbor interactions constructed entirely from operators in the set $`𝒯`$,
$$H=\lambda \underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n1}{}}\left(Z_{i,j}Z_{i,j+1}+X_{j,i}X_{j+1,i}\right).$$
(33)
Since $`H`$ is constructed entirely from elements of $`𝒯`$, this Hamiltonian can be decomposed as
$$H=\underset{s^X,s^Z}{}H_{s^X,s^Z}I_2.$$
(34)
To understand the exact nature of this Hamiltonian, we would need to diagonalize each of the $`H_{s^X,s^Z}`$. What can be said, however, is that the ground state of the system will arise as the ground state of one or more of the $`H_{s^X,s^Z}`$ (numerical diagonalization of systems with a few qubits show that the ground state comes from only the $`s_i^X=s_i^Z=+1`$ subsystem and we conjecture that this subspace always contains the ground state.) If $`k`$ of the $`H_{s^X,s^Z}`$ contribute to the ground state, the degeneracy of the ground state will be $`2k`$ due to the subsystem corresponding to the $``$ operators. Thus we see that we can encode quantum information into the subsystem degeneracy of this Hamiltonian, in the similar manner that information is encoded into the ground state of a Hamiltonian related to the toric codeKitaev (2003); Dennis et al. (2002); Ogburn and Preskill (1999). However, we do not know whether this system exhibits a gap in its excitation spectrum similar to that which exists in the toric codes. In Section VI we return will introduce a similar Hamiltonian for our three-dimensional operator quantum error correcting subsystem.
## V Three-dimensional Operator Quantum Error Correcting Subsystem
We now turn to a three-dimensional operator quantum error correcting subsystem which is a generalization of the two-dimensional subsystem code we presented above. In particular, whereas the construction for the two-dimensional model relied on the structure of $`𝒯`$ containing Pauli operators with even number of Pauli $`Z`$’s in a row and even number of Pauli $`X`$’s in a column, in the three-dimensional case we rely on a new set of operators with an even number of Pauli $`Z`$’s in the $`yz`$ plane and an even number of Pauli $`X`$’s in the $`xy`$ plane.
Consider a cubic lattice of size $`n\times n\times n`$ with qubits located at the vertices of this lattice and let $`n`$ be odd. Let $`O_{i,j,k}`$ denote the operator $`O`$ acting on the qubit located at the $`(i,j,k)`$th lattice site tensored with identity on all other qubits. We again use a compact notation to denote Pauli operator on our $`n^3`$ qubits by using two $`n^3`$ bit strings,
$$P(a,b)=\underset{i,j,k=1}{\overset{n}{}}X_{i,j,k}^{a_{i,j,k}}Z_{i,j,,k}^{b_{i,j,k}}=\underset{i,j,k=1}{\overset{n}{}}\{\begin{array}{cc}X_{i,j,k}\hfill & \mathrm{if}a_{i,j,k}=1\mathrm{and}b_{i,j,k}=0\hfill \\ Z_{i,j,k}\hfill & \mathrm{if}a_{i,j,k}=0\mathrm{and}b_{i,j,k}=1\hfill \\ iY_{i,j,k}\hfill & \mathrm{if}a_{i,j,k}=1\mathrm{and}b_{i,j,k}=1\hfill \end{array},$$
(35)
where $`a,b_2^{n^3}`$ are $`n`$ by $`n`$ by $`n`$ arrays of bits.
As in the two-dimensional case, we will define three sets of operators, $`𝒯_3`$, $`𝒮_3`$, and $`_3`$ which are essential to understanding the subsystem structure of our qubits.
The first set of Pauli operators which will concern us, $`𝒯_3`$, is made up of Pauli operators which have an even number of $`X_{i,j,k}`$ operators in each xy-plane and an even number of $`Z_{i,j,k}`$ operators in each yz-plane:
$$𝒯_3=\left\{(1)^\varphi P(a,b)\right|\varphi _2,\underset{i,j=1}{\overset{n}{}}a_{i,j,k}=0,\mathrm{and}\underset{j,k=1}{\overset{n}{}}b_{i,j,k}=0\}.$$
(36)
These operators, like the analogous two-dimensional $`𝒯`$ form a group under multiplication. This group can be generated by nearest neighbor operators on our cubic lattice
$$𝒯_3=X_{k,i,j}X_{k+1,i,j},X_{i,k,j}X_{i,k+1,j},Z_{i,j,k}Z_{i,j,k+1},Z_{i,k,j}Z_{i,k+1,j},i,j_n,k_{n1}$$
(37)
The second set of Pauli operators we will be interested in is a subset of $`𝒯_3`$, which we denote $`𝒮_3`$. $`𝒮_3`$ consists of Pauli operators which are made up of an even number of xy-planes made entirely of Pauli $`Z`$ operators and an even number of yz-planes made entirely of Pauli $`X`$ operators:
$$𝒮_3=\left\{P(a,b)\right|\underset{i=1}{\overset{n}{}}\left(\underset{j,k=1}{\overset{n}{}}a_{i,j,k}\right)=0,\underset{k=1}{\overset{n}{}}\left(\underset{i,j=1}{\overset{n}{}}b_{i,j,k}\right)=0\}.$$
(38)
$`𝒮_3`$ is an Abelian subgroup of $`𝒯_3`$ and all of the elements of $`𝒮_3`$ commute with all of the elements of $`𝒯_3`$. It can be generated by nearest xy-plane and yz-plane operators:
$$𝒮_3=\underset{i,j=1}{\overset{n}{}}X_{i,j,k}X_{i,i,k+1},\underset{i,j=1}{\overset{n}{}}Z_{k,i,j}Z_{k+1,i,j},k_{n1}.$$
(39)
We label these generators, as before:
$`S_k^X={\displaystyle \underset{i,j=1}{\overset{n}{}}}X_{i,j,k}X_{i,j,k},\mathrm{and}S_i^Z={\displaystyle \underset{j,k=1}{\overset{n}{}}}Z_{i,j,k}Z_{i+1,j,k}.`$ (40)
$`𝒮_3`$ is again a stabilizer group.
The final set of operators which we will consider, $`_3`$, is similar to $`𝒮_3`$ except that the evenness condition becomes and oddness condition
$$_3=\left\{(1)^\varphi P(a,b)\right|\varphi _2,\underset{i=1}{\overset{n}{}}\left(\underset{j,k=1}{\overset{n}{}}a_{i,j,k}\right)=1,\underset{k=1}{\overset{n}{}}\left(\underset{i,j=1}{\overset{n}{}}b_{i,j,k}\right)=1\}.$$
(41)
$`_3`$ together with $`𝒮_3`$ forms a group and all of the elements of $`_3`$ commute with those of $`𝒯_3`$.
### V.1 Subsystem Structure
All three of the sets, $`𝒯_3`$, $`𝒮_3`$, and $`_3`$ will play a directly analogous role to the sets $`𝒯`$, $`𝒮`$, and $``$ in our two-dimensional model. In particular if we let $``$ denote the Hilbert space of our $`n^3`$ qubits, then we can partition this space into subspaces labelled by the $`2(n1)`$ different $`\pm 1`$ eigenvalues of the operators $`S_k^X`$ and $`S_i^Z`$ of Eq. 40. Again we will label these eigenvalues by $`s_k^X`$ and $`s_i^Z`$, with $`s^X`$ and $`s^Z`$ labelling these strings. The Hilbert space of the system then decomposes as
$$=\underset{s_1^X,\mathrm{},s_{n1}^X,s_1^Z,\mathrm{},s_{n1}^Z=\pm 1}{}_{s^X,s^Z}=\underset{s^X,s^Z}{}_{s^X,s^Z}.$$
(42)
Again, the $`_{s^X,s^Z}`$ subspaces have a tensor product structure, $`_{s^X,s^Z}=_{s^X,s^Z}^𝒯_{s^X,s^Z}^𝒯`$, such that elements of $`_3`$ act as
$$L=\underset{s^X,s^Z}{}I_{2^{n^32n+1}}L_{s^X,s^Z},L_3,$$
(43)
and those of $`𝒯_3`$ act as
$$T=\underset{s^X,s^Z}{}T_{s^X,s^Z}I_2,T𝒯_3.$$
(44)
### V.2 Subsystem Quantum Error Correcting Procedure
The subsystem error correcting procedure for the three-dimensional code nearly directly mimics that of the two-dimensional code. Here we discuss how the subsystem error correcting procedure works without going into the details as we did in the two-dimensional case. The three-dimensional procedure is nearly identical to that of the two-dimensional procedure with the sets $`𝒯`$, $`𝒮`$, and $``$ interchanged with the sets $`𝒯_3`$, $`𝒮_3`$, and $`_3`$ respectively.
We will again encode our quantum information into the $`_{s^X=s^Z=\{+1\}^{n1}}^{}`$ subsystem. Whereas for the two-dimensional code we defined error strings for the rows and column conditions of the set $`𝒯`$, now we define error strings for the $`xy`$ and $`yz`$ plane conditions of the set $`𝒯_3`$. If the Pauli error $`P(a,b)`$ occurs on our system, then we can define the two error strings
$$e_k(a)=\underset{i,j=1}{\overset{n}{}}a_{i,j,k}\mathrm{and}f_i(b)=\underset{j,k=1}{\overset{n}{}}b_{i,j,k}.$$
(45)
Pauli errors with $`e_k(a)=f_i(b)=0`$ are errors from $`𝒯_3`$ and act trivially on the information encoded into the $`_{s_k^X=s_i^Z=+1}^{}`$ subsystem.
The quantum error correction procedure is then directly analogous to the one for the two-dimensional code. We measure the $`s_k^X`$ and $`s_i^Z`$ operators and treat these as nearest neighbor parity checks for a redundancy code on the $`f_i(b)`$ and $`e_k(a)`$ respectively. Then in direct analogy with the two-dimensional code, we can apply a subsystem error correcting procedure which restores the system modulo the subsystem structure. The three-dimensional code is a $`[n^3,n,1]`$ code. We have thus gained nothing in terms of the distance of the code, but, as we will argue in the next Section, the three-dimensional code when converted to a Hamiltonian whose ground state is the subsystem code has intriguing features not found in the two-dimensional code.
### V.3 Logical Operators
The logical operators for the three-dimensional code are directly analogous to those in the two-dimensional code. As in the two-dimensional code, operators from $`_3`$ can be used to enact Pauli operators on the information encoded into the $`_{s_k^X=s_i^Z=+1}^{}`$ subsystem. Similarly, we can enact a controlled-not between two encoded qubits by performing $`n^3`$ controlled-not gates between two identical copies of the code. This allows us to again perform any operation in the normalizer of the Pauli group on our encoded qubits.
## VI Self-Correction in the Three-dimensional Example
In the 1930s, when Alan Turing wrote his now classic papersTuring (1936, 982) laying out the foundations of computer science, there was absolutely no reason to believe that any computing device such as the one described by Turing could actually be built. One of the foremost problems, immediately apparent to the engineers of the day, was the lack of reliable components out of which a computer could be built. Von Neumann solved this problem, in theory, by showing that robust encoding of the classical information could be used to overcome errors and faulty components in a computervon Neumann (1956). Despite Von Neumann’s theoretical ideas, it took the invention of the transistor and the integrated circuit, to mention only the broadest innovations, in order to bring forth the technological movement now known as the computer revolution. The overarching result of the technological innovations responsible for the computer revolution was the development of techniques which exhibited Von Neumann’s theoretical ideas in a natural setting. Modern computers naturally correct errors in both the storage and manipulation of classical information. The task of robust storage and manipulation of the data is essentially guaranteed by the physics of these devices. There are distinct physical reasons why robust storage and manipulation of classical information is possible.
If there are distinct physical reasons why robust storage and manipulation of classical information is possible, an obvious question to ask in the quantum information sciences is whether we can mimic these effects in the quantum domain. Do there exist, or can we engineer, physical systems whose physics ensures that the robust storage and manipulation of quantum information is possible? In this section, we will present evidence, in the form of a mean field argument, that a Hamiltonian related to the three-dimensional subsystem code might be exactly this type of system.
### VI.1 Self-Correcting Quantum Memories
The traditional approach to building a robust fault-tolerant quantum computer imagines building the computer using a complex microarchitecture of quantum error correcting fault-tolerant procedures. This poses a severe technological overhead of controlling thousands of qubits in a complex manner, simply to get a single robust qubit. KitaevKitaev (1995, 2003) was the first to suggest that an alternative, less complex, method to constructing a fault-tolerant quantum computer might be possible. Kitaev showed that there exists a quantum error correcting code, the toric code, which is the degenerate ground state of a certain four body spatially local Hamiltonian on a two-dimensional lattice of qubits. Kitaev imagined encoding quantum information into the ground state of this system and then, because there is an energy gap between the ground state of this system and the first excited state and because the errors which will destroy quantum information consist of error which scale like the size of the lattice, this quantum information would be protected from decoherence due to the environment as long as the temperature of the environment was sufficiently low.
It is important to note that the Hamiltonian implementation of Kitaev’s toric code (by which we mean encoding information into a physical system governed by the four-body Hamiltonian associated with the toric code), while providing a mechanism for the robust storage of quantum information, does not provide a full fault-tolerant method for quantum computation. The reason for this is that during the implementation of the physical processes which manipulate the information encoded into the ground state of the system, real excitations will be created which will disorder the system. This distinction has been confused in a manner because the toric codes can be used to construct a fault-tolerant quantum computer, but only with the aid of external quantum control which serves to identify and correct errors which occur during the manipulation of the information encoded into the system (see for example Dennis et al. (2002)).
The idea of a self-correcting quantum memory is to overcome the limitations of Kitaev’s original model by constructing a physical system whose energy levels not only correspond to a quantum error correcting code (in our case a subsystem code), but which also uses the energetics of this system to actively correct real errors created when the quantum information is being manipulated. Thus, while in the toric codes, a single real error on the system can create excitations which can disorder the system, in a self-correcting system, a single real error on the system cannot disorder the system.
In order to explain the distinction of a self-correcting memory from the original toric code we will compare the situation to that of the one dimensional and two-dimensional classical ferromagnetic Ising model. These models will be analogous to the toric code Hamiltonian model and a self-correcting Hamiltonian model, respectively.
Recall that in a ferromagnetic Ising model one takes a lattice of classical spins and these are coupled by Ising interactions between the neighboring spins via a Hamiltonian
$$H=\frac{J}{2}\underset{i,j}{}s_is_j,$$
(46)
where $`s_i\{\pm 1\}`$ are the spin variables, the sum $`i,j`$ is over neighbors in the lattice, and $`J>0`$. Notice that the ground state of this Hamiltonian corresponds to the uniform states $`s_i=+1,i`$ or $`s_i=1,i`$. These, of course, are also the two codewords for a classical redundancy code. Thus we can imagine that we encode classical information into the ground state of this Hamiltonian in direct analogy to the way in which information (but quantum this time) is encoded into ground state of the toric code. Errors on the Ising codes are just bit flips. From hereon out when we refer to the Ising model we will implicitly be discussing the ferromagnetic Ising model.
Recall some basic properties of the one and two-dimensional Ising models (see, for example Plischke and Bergersen (1994)). We begin by discussing the thermal equilibrium values of the total magnetization,
$$M=\underset{i}{}s_i,$$
(47)
of these models. In one dimension, for any $`T>0`$, the total magnetization of the Ising model vanishes in thermal equilibrium, whereas for the two-dimensional Ising model, the magnetization is zero above some critical temperature $`T_c`$ and below this temperature, two magnetizations of equal magnitude and opposite sign are maintained. Since the magnetization is a measure of the information recorded in the redundancy code, we see that if we encode information into the ground state of the one dimensional Ising model and this system is allowed to reach thermal equilibrium, then this information will be destroyed. On the other hand, for the two-dimensional Ising model, if we encode information into the ground state and the system is below the critical temperature $`T_c`$, then this information will be maintained. Above $`T_c`$, like the one-dimensional Ising model, the information will be destroyed. From the point of view of storing the information in the thermal states of these models, the two-dimensional Ising model is a robust medium, but the one-dimensional Ising model is not.
But what about the properties of the Ising models on the way to reaching equilibrium (i.e. during the time evolution with the environment)? In the one dimensional case we find that the system will generically (depending on the exact method of relaxation) take a time which is suppressed like a Boltzman factor $`e^{J/T}`$. Thus at low enough temperature, we can encode information into the ground state of the one dimensional Ising model and it will be protected for a long time. While the scaling of this decay rate is favorable in the temperature T, this type of approach is different from what is done in standard error correction where larger redundancy can be used to overcome errors without changing the error rate (as long as that error rate is below the threshold for the error correcting code.)
What happens for the time evolution of a two-dimensional Ising model? If we start the system in one of the redundancy code states, then far below $`T_c`$ the system will relax quickly to the closest of the two equilibrium states with a large total magnetization. As we raise the temperature closer to $`T_c`$, this relaxation will slow down. Above $`T_c`$ the picture is similar to that of the one dimensional Ising model that if we are close to $`T_c`$, then the relaxation to vanishing magnetization is suppressed like $`e^{J/(TT_c)}`$.
What are the main reasons for the differences in the ability to store information in the one and two-dimensional Ising models? A rough heuristic of what is happening is that in the two-dimensional Ising model, the errors self-correctBarnes and Warren (2000). Consider starting the one dimensional Ising model in the all $`s_i=+1`$ state. Now flip one of the spins at the end of the chain. This will cost an energy $`J`$. Flipping the neighbor of this spin will then cost no energy. Proceeding along the chain in this manner one sees that one can expend energy $`J`$ to turn the system from the codeword all $`s_i=+1`$ state to the all $`s_i=1`$ state. Thus the environment need only supply this energy to disorder the system. However, in the two-dimensional Ising model, something different happens. Suppose again that we start in the all $`s_i=+1`$ state. Here if we flip a single spin (say on the boundary of the lattice) then the energy required to flip this spin is $`J`$ times the number of bonds this spin has with its neighbors. Now flipping a neighbor will cost energy: in the two-dimensional model the energy cost of flipping a connected domain of spins is proportional to the perimeter of this domain. Since to get from the all $`s_i=+1`$ codeword to the all $`s_i=1`$ codeword we need to build a domain of size the entire lattice, we see that we will require at least an energy times the size of the lattice to disorder the system. Suppose, now that errors are happening at some rate to all of the spins in the Ising models. In the one dimensional model, once one creates a single error, there is no energy barrier to disordering the system. In the two-dimensional model, however, there is now an energy barrier. In particular, the system coupled to its environment will not only cause errors, but will also cause errors to be corrected by shrinking the domains of flipped errors. As long as the error rate is not too strong (which corresponds loosely to being below the critical temperature $`T_c`$) the pathways that fix the error will dominate the actual creation of errors.
Thus we see that a two-dimensional Ising model operating below the critical temperature is performing classical error correction on information stored in a redundancy code. In the one-dimensional Ising model and in the two-dimensional Ising model above the critical temperature, there is suppression due to a Boltzman factor, but there is no self-correction (or the self-correction is not fast enough) and the information stored in the redundancy code is destroyed.
The two-dimensional anyon models of topological quantum computing and variationsKitaev (1995); Ogburn and Preskill (1999); Dennis et al. (2002); Kitaev (2003); Wang et al. (2003); Freedman et al. (2001); Nayak and Shtengel (2001); Brink and Wang (2003); Freedman (2003); Freedman et al. (2003); Hamma et al. (2004); Freedman et al. (2004a, b); Freedman et al. (2005); Kitaev (2005), including Kitaev’s toric model, all share the property with the one dimensional Ising model that the system can disordered using only an energy proportional to the gap in the Hamiltonian. (The models in Hamma et al. (2004) contain errors similar to those in the two dimensional Ising model for certain types of quantum errors, but not for both phase and bit flip errors.) This can provide protection via an exponential suppression due to a Boltzman factor, but this does not provide indefinite correction. The idea of a self-correcting quantum memory, however is to mimic the two-dimensional Ising model. In particular below a critical temperature, quantum information stored in the system should persist even when the system is in thermal equilibrium and further the system should have a mechanism whereby real errors are corrected by the energetics of the system faster than the real errors occur when operating below the critical temperature.
Finally, we note that there is one system which is widely suspected to be a self-correcting quantum memory. This is a version of the toric code on a four-dimensional latticeDennis et al. (2002). The problem with this model is that it exists in four dimensions and that it requires greater than two-qubit interactions in order to implement and is therefore not realistic for practical implementations. Our original motivation for considering the subsystem codes presented in this paper was to obtain a self-correcting system in the realistic setting of three or lower dimensions and using two-qubit interactions.
### VI.2 The Three-dimensional Hamiltonian
Next we turn our attention to a system which may be a self-correcting quantum memory. We provide evidence for this by showing that in a mean field approximation this Hamiltonian has properties for the expectation value of its energy which is similar to the energetics of the two-dimensional Ising model.
Consider the following Hamiltonian on a cubic lattice of qubits constructed exclusively from elements of $`𝒯_3`$,
$$H=\lambda \underset{i,j=1}{\overset{n}{}}\underset{k=1}{\overset{n1}{}}\left(X_{k,i,j}X_{k+1,i,j}+X_{i,k,j}X_{i,k+1,j}+Z_{i,k,j}Z_{i,k+1,j}+Z_{i,j,k}Z_{i,j,k+1}\right),$$
(48)
with $`\lambda >0`$. This Hamiltonian consists of Ising couplings along the direction $`x`$ in the $`xy`$-plane and along the direction $`z`$ in the $`yz`$-plane. As in the two-dimensional Hamiltonian, Eq. (33), we can use the fact that $`H`$ is a sum of operators from $`𝒯_3`$ to block diagonalize $`H`$ with respect to the subsystem structure of our three-dimensional operator quantum error correcting subsystem:
$$H=\underset{s^X,s^Z}{}H_{s^X,s^Z}I_2.$$
(49)
Information can then be encoded into the $`I_2`$ subsystem. We would like to show that if we perform such an encoding, then the information stored in this subsystem will be protected from the effect of general quantum errors in a manner similar to that of the two-dimensional Ising model.
### VI.3 The Mean Field Argument
Ignore, for the moment the boundary conditions for the Hamiltonian and suppose that the ground state of Eq. (48) has the following properties for the expectation values of the Ising bonds in the system,
$`X_{k,i,j}X_{k+1,i,j}_G=c_{xx}`$
$`X_{i,k,j}X_{i,k+1,j}_G=c_{xy}`$
$`Z_{i,k,j}Z_{i,k+1,j}_G=c_{zy}`$
$`Z_{i,j,k}Z_{i,j,k+1}_G=c_{zz},`$
where $`c_{\alpha \beta }>0`$. In such a phase, the ground state energy is given by
$$E_G=H_G=\lambda n^2(n1)(c_{xx}+c_{xy}+c_{zy}+c_{zz}).$$
(51)
Now consider the effect of a single Pauli error on a single qubit in the lattice (assume this is away from the boundary). For example, consider a Pauli $`X`$ error. Using the fact that $`X`$ commutes with Ising bonds oriented along the $`x`$ direction but anticommutes with Ising bonds oriented along the $`z`$ direction, it then follows that expectation value of the energy of the system changes to
$$E_1=H_1=E_G+2\lambda (c_{zy}+c_{zy}+c_{zz}+c_{zz}),$$
(52)
where we have separated out each term arising from each of the four $`z`$ direction Ising bonds which connect to the lattice site where the $`X`$ error has occurred. Thus we see that a single bit flip error on the ground state will cause the expectation value of the energy to increase. Now consider the effect of applying a second Pauli $`X`$ error which is a nearest neighbor to the original spin in the same $`yz`$ plane. Now the Ising bond between these two errors does not contribute to the change in the expectation value of the ground state, but all of the $`z`$ direction Ising bonds around the perimeter of the two flipped spins do contribute. Thus, for example, if the spins are neighbors along the $`y`$ direction, the expectation value of the energy of this new state is
$$E_2=H_2=E_G+2\lambda (2c_{zy}+4c_{zz}).$$
(53)
Generalizing the above argument we see that a connected domain of Pauli $`X`$ errors in an $`yz`$ plane will result in an energy increase proportional to the perimeter of the domain. A similar argument will hold for Pauli $`Z`$ errors in an $`xy`$ plane, but now the Ising bonds along the $`x`$ direction will contribute to the change in expectation value, and those along the $`z`$ will not contribute.
Suppose that a general error $`P(a,b)`$ occurs on the ground state of $`H`$. Then in each plane $`yz`$ plane the part of the error coming from $`a`$ will produce excitations whose expectation of the energy scales like the perimeter of domains of errors in $`a`$ and similarly for each $`xy`$ plane, but now for the $`b`$ component of the errors. From this argument we see that, at least for the expectation value of the energy, the system looks very similar to the two-dimensional Ising model. But now instead of only bit flip errors, more general quantum errors produce changes in energy of the system which are proportional to the perimeter of the erred domain. Thus we argue that this provides evidence that the model we have presented will be self-correcting. For the same reasons that the two-dimensional Ising model will not disorder the classical information stored in a redundancy code (i.e. since our errors require (expected) energy proportional to the perimeter of the erred domain) we expect the quantum information stored in our system will not disorder up to some critical temperature.
Of course, the evidence we have provided is based on numerous assumptions arising from our mean field model. First of all we have ignored boundary conditions. It is possible that certain edge states could disorder the system. Secondly we have only made arguments about the expectation value of the energy after error have occurred to the ground state. This doesn’t give us concrete information about the energy level structure of our Hamiltonian. It could be that while the expectation value scales like the perimeter, there are actually error pathways whose energetics are much less favorable. Third we have assumed the existence of a phase with the desire expectation values of the bond energies. This phase may not exist, i.e. it may be that in the thermodynamic limit, the expectations values all vanish. This would totally invalidate the mean-field argument we have given above. Given these caveats our mean field argument only suggests that the system will be self correcting. Clearly rigorously establishing whether our memory is self-correcting is a challenging open problem.
### VI.4 The Quantum Error Correcting Order Parameter
In the Ising model, an indication that the information stored in a redundancy code is still there after thermalization is the persistence of the total magnetization of the system. In particular, for the two-dimensional Ising model at a temperature between zero and the critical temperature, the magnetization in equilibrium is never exactly equal to its maximal value, $`\pm n^2`$. This is because there are always small domains of flipped spins due to the interaction of the system with its environment. However, a magnetization which is different from zero may be interpreted as a measurement of the majority vote for the redundancy code in this system. If we are going to establish that our quantum system is actually self-correcting, it is important to identify an order parameter for our system which can be used to reveal the presence of quantum information in our system. This can be done using the operator quantum error correcting subsystem properties of the three-dimensional code.
Suppose we encode quantum information into a quantum error correcting code and apply a number of quantum errors less than the number which the code has been designed to correct. We know from the theory of quantum error correction that the encoded information in this system can be recovered by the measurement of an appropriate error syndrome and the application of the appropriate recovery procedure. We can use this to construct an order parameter for any quantum error correcting code.
A note about the nature of order parameters for quantum information before we describe this parameter for our three-dimensional model. In the classical Ising model, we found that below the critical temperature there was a bifurcation of the system into two magnetizations of equal and opposite magnitude. Since classical information is based upon a bit, we are not surprised to find that such a bifurcation into two states happens. Actually, depending on the initial state of the system before it is thermalized, the thermal state of the system can be any value between these two extremes. But if we start our system by encoding into one of the two codestates (the all $`\pm 1`$ state) then only the two bifurcated values will result after thermalization. What is the analogous situation for quantum information? For quantum information, we must show not that the bifurcation happens for a bit of encoded information but instead for a qubit of encoded information. Since a qubit is parameterized by the Bloch-sphere, one might expect that one needs an order parameter with similar properties. Such an order parameter can be constructed, but we can get away with examining fewer parameters in order to show the robustness of the quantum information. In particular if we make measurements along the $`x`$, $`y`$, and $`z`$ directions of the quantum information, then because of the linearity of the density operator we can use this to show that the information has been preserved. In particular, we can imagine encoding into one of the eigenstates of Pauli operators along these directions and looking at this system after it has thermalized. Notice now that instead of a single order parameter, we will have three order parameters. In order to demonstrate the self-correcting nature of our three-dimensional Hamiltonian, we will need to show that the expectation value of these three order parameters each bifurcate below some critical temperature.
Consider, now, the active recovery procedure for our three-dimensional subsystem code. We begin by measuring the $`S_k^X`$ and $`S_i^Z`$ operators. Given these measurements, we can, as in the active recovery procedure, deduce an appropriate recovery operator to restore the information originally encoded into a subsystem, modulo the subsystem structure of the system. If we were to apply this syndrome and measure the encoded Pauli logical operators for the code (which were given in Sec. V.3) then this would serve as an order parameter for our system. We note, however that if we are simply interested in measuring the Pauli logical operators and not in fully restoring the information into the original subsystem encoding, we do not need to actually apply the syndrome. This is because the syndrome we diagnose will either commute or anti-commute with the encoded Pauli operator we wish to measure. Thus given the syndrome we can, instead of applying the appropriate recovery operation, simply flip or not flip the answer we get from measuring the encoded Pauli operator depending on the exact syndrome measured.
Now we can explicitly describe our order parameters. Suppose we measure the $`S_k^X`$ and $`S_i^Z`$ operators and obtain the values $`s_k^X`$ and $`s_i^Z`$ and we measure one of the encoded Pauli operators. Suppose, for example, that this encoded Pauli operator is an encoded $`X`$ operator which we measure by measuring all of the Pauli $`X`$ operators in a fixed $`xy`$ plane. Given the $`s_k^X`$ we can deduce whether an even or an odd number of Pauli $`Z`$ errors will need to apply to the fixed $`xy`$ plane to restore the quantum information in the subsystem (if possible). If this number is odd, then we simply flip the value of we measure for the encoded Pauli $`X`$ operator and if this number is even, we do not flip the value for the encoded $`X`$ operator. Thus we see that we measure the $`x`$ directional order parameter from our system by measuring the $`S_k^X`$ and the encoded $`X`$ operator and, as a function of these values, produce a single number representing the $`x`$ directional order parameter. Similar comments hold for the order parameters along the other cardinal directions.
## VII Conclusion
In this paper we have constructed a new class of quantum error correcting procedures based upon the notion of encoding quantum information into a subsystem. By encoding into a subsystem, we were able to demonstrate a recovery routine which explicitly used the subsystem structure. The three-dimensional code we constructed was shown to be related to a three-dimensional spin lattice system which we gave evidence for being a self-correcting quantum memory. We close by remarking on some open problems for this three-dimensional system and some thoughts about future directions for constructing self-correcting quantum systems.
The first open question concerns the implementation of our model in a physical system. A particularly promising system for such an implementation is with ultracold atoms trapped in an optical latticeDuetsch et al. (2000). Duan, Demler, and LukinDuan et al. (2003) showed how to simulate a large class of spin-spin interactions for these systems. An open question is whether their techniques allow one to implement our three-dimensional anisotropic spin-spin Hamiltonian. Of particular concern is the magnitude of the spin-spin coupling which one can achieve in these models. This will directly effect the critical temperature for any self-correction that occurs in the system. A further concern for the physical implementation in an optical lattice is the ability to measure the syndrome operators $`S_k^X`$ and $`S_i^Z`$ along with appropriate logical Pauli operators. Finally one would like to understand how to produce an effective controlled-not coupling between two such encoded lattices. Solutions to all of these problems would allow one to propose an experiment in which a self-correcting quantum memory could be demonstrated.
A second open question is, of course, whether our three-dimensional system is indeed self-correcting. Noting that the Hamiltonian for this system does not possess a sign-problem, one approach to verifying this question would proceed by performing quantum Monte Carlo simulations of this system. The order parameters we have described could then be simulated at finite temperature and evidence for self-correction could then be examined. Another promising approach is to use recent new ideas in the density matrix renormalization groupVidal (2003); Verstraete and Cirac (2004) to simulate the thermal properties of this system.
Another important question is whether one can design a self-correcting quantum systems in two dimensions. This would be particularly desirable if one wishes to physically implement the self-correction in a solid-state system.
Finally a large open question is what role operator quantum error correcting subsystem codes can play in quantum information science. What do other such subsystem codes look like? Since these codes have important properties due to the fact that they exploit degenerate quantum codes, can subsystem codes be used to beat the quantum Hamming boundGottesman (1997)? Further, results which relied on showing that there were no subspace codes with certain properties (for example, as in Reimpell and Werner (2005)) need to be reexamined in light of the existence of operator quantum error correcting subsystem codes.
###### Acknowledgements.
This work was supported in part by the DARPA QuIST program under grant AFRL-F30602-01-2-0521. D.B. acknowledges the support of a postdoctoral fellowship at the Santa Fe Institute, during which some of initial work on this paper was carried out. D.B. also thanks Kaveh Khodjasteh, Daniel Lidar, Michael Nielsen, and David Poulin for useful discussions.
|
warning/0506/gr-qc0506051.html
|
ar5iv
|
text
|
# Decoherence during Inflation: the generation of classical inhomogeneities
## I Introduction
The emergence of classical physics from quantum behaviour is important for several physical phenomena in the early Universe. This is beyond the fundamental requirement that only after the Planck time can the metric of the Universe be assumed to be classical. For example, the inflationary era is assumed to have been induced by scalar inflaton fields, with simple potentials linde ; Boya2004 . Such fields are typically assumed to have classical behaviour, although in principle a full quantum description should be used. In fact, the origin of large scale structure in the Universe can be traced back to quantum fluctuations that, after crossing the Hubble radius, were frozen and became classical, stochastic, inhomogeneities todosLSS .
It is generally assumed that several phase transitions have occurred during the expansion of the Universe old . As in the case for the inflaton fields, the (scalar) order parameter fields that describe these transitions are described classically. However, the description of early universe phase transitions from first principles is intrinsically quantum mechanical cormier . As a specific application Kibble of the previous point, the very notion of topological defects (e.g. strings and monopoles) that characterize the domain structure after a finite-time transition, and whose presence has consequences for the early universe, is based on this assumption of classical behaviour for the order parameter vilen , as it distributes itself between the several degenerate ground states of the ordered system.
In previous publications, one of us has analyzed the emergence of a classical order parameter during a second order phase transition and the role of decoherence in the process of topological defect formation lomplb ; deconpb ; diana ; lomplb2 .
In the present paper our concern is directly related with the first point above, the quantum to classical transition of the inflaton. Any approach must take into account both the quantum nature of the scalar field and the non-equilibrium aspects of the process calzettahu95 . The problem of the quantum to classical transition in the context of inflationary models was first addressed by Guth and Pi guthpi . In that work, the authors used an inverted harmonic oscillator as a toy model to describe the early time evolution of the inflaton, starting from a Gaussian quantum state centered on the maximum of the potential. They subsequently showed that, according to Schrödinger’s equation, the initial wave packet maintains its Gaussian shape (due to the linearity of the model). Since the wave function is Gaussian, the Wigner function is positive for all times. Moreover, it peaks on the classical trajectories in phase space as the wave function spreads. The Wigner function can then be interpreted as a classical probability distribution for coordinates and momenta, showing sharp classical correlations at long times. In other words, the initial Gaussian state becomes highly squeezed and indistinguishable from a classical stochastic process. In this sense, one recovers a classical evolution of the inflaton rolling down the hill.
A similar approach has been used by many authors to describe the appearance of classical inhomogeneities from quantum fluctuations in the inflationary era staro ; Mopntemayor . Indeed, a massless free field $`\varphi `$ in an expanding universe can be written as $`\varphi =a^1\psi `$ where $`a`$ is the scale factor and the Fourier modes of the field $`\psi `$ satisfy the linear equation
$$\psi _k^{\prime \prime }+(k^2\frac{a^{\prime \prime }}{a})\psi _k=0.$$
(1)
For sufficiently long-wavelengths ($`k^2a^{\prime \prime }/a`$), this equation describes an unstable oscillator. If one considers an initial Gaussian wave function, it will remain Gaussian for all times, and it will spread with time. As with the toy model of Guth and Pi, one can show that classical correlations do appear, and that the Wigner function can again be interpreted as a classical probability distribution in phase space. (It is interesting to note that a similar mechanism can be invoked to explain the origin of a classical, cosmological magnetic field from amplification of quantum fluctuations).
However, classical correlations are only one aspect of classical behaviour. It was subsequently recognized that, in order to have a complete classical limit, the role of the environment is crucial, since its interaction with the system distinguishes the field basis as the pointer basis kiefer . \[We are reminded that, even for the fundamental problem of the space-time metric becoming classical, simple arguments based on minisuperspace models suggest that the classical treatment is only correct because of the interaction of the metric with other quantum degrees of freedom halli .\]
While these linear instabilities cited above characterize free fields, the approach fails when interactions are taken into account. Indeed, as shown again in simple quantum mechanical models (e.g. the anharmonic inverted oscillator), an initially Gaussian wave function becomes non-Gaussian when evolved numerically with the Schrödinger equation. The Wigner function now develops negative parts, and its interpretation as a classical probability breaks down diana . One can always force the Gaussianity of the wave function by using a Gaussian variational wave function as an approximate solution of the Schrödinger equation, but this approximation deviates significantly from the exact solution as the wave function probes the non-linearities of the potential diana ; stancioff .
When interactions are taken into account, classical behaviour is recovered only for “open systems”, in which the observable degrees of freedom interact with their environment. When this interaction produces both a diagonalization of the reduced density matrix and a positive Wigner function, the quantum to classical transition is completed giulinibook .
In Ref. diana has been studied an anharmonic inverted oscillator coupled to a high temperature environment. Under some considerations, it was shown that the system becomes classical very quickly, even before the wave function probes the non-linearities of the potential. Being an early time event, the quantum to classical transition can now be studied perturbatively. In general, recoherence effects are not expected nuno . Taking these facts into account, we have extended the approach to field theory models lomplb ; deconpb . In field theory, one is usually interested in the long-wavelengths of the order parameter. Even the early universe is replete with fields of all sorts which comprise a rich environment, in the inflationary example, we considered a model in which the system-field interacts with the environment-field, including only its own short-wavelengths. This is enough during inflation. Assuming weak self-coupling constant (a flat inflaton potential) we have shown that decoherence is an event shorter than the time $`t_{\mathrm{end}}`$, which is a typical time-scale for the duration of inflation. As a result, perturbative calculations are justified deconpb . Subsequent dynamics can be described by a stochastic Langevin equation, the details of which are only known for early times matacz .
In our approach, the quantum to classical transition is defined by the diagonalization of the reduced density matrix. In phase transitions the separation between long and short-wavelengths is determined by their stability, which depends on the parameters of the potential. During Inflation, this separation is set by the existence of the Hubble radius. Modes cross the apparent horizon during their evolution, and they are usually treated as classical. The main motivation of this paper is to present a formal way to understand this statement within the open quantum system approach. In the last sense, decoherence is the critical ingredient if we are to dynamically demonstrate the quantum-to-classical transition of the open system.
The splitting between short and long-wavelength modes may be done using a time-dependent or time-independent comoving cut-off as well as a smoother window function matacz ; Riotto . Since any time-dependent splitting produce an effective and arbitrary (split-dependent) interaction between the system and environment degrees of freedom, which can only be discarded by making some additional assumptions, we consider convenient to use a time-independent comoving cut-off, as it has been done in Refs. HUPAZYHANG ; lombmazz
The paper is organized as follows. In Section II we introduce our model. This is a theory containing a real system-field, massless and minimally coupled to a fixed de Sitter background. We compute the influence functional by integrating out the environmental sector of the field, composed by the short-wavelength modes. Section III is dedicated to reviewing the evaluation of the master equation and the diffusion coefficients which are relevant in order to study decoherence. In Section IV we analyze the diffusion coefficients and evaluate upper bounds on the decoherence times. As we will see, decoherence takes place before the end of the inflationary period. Section V is concerned with the effective stochastic evolution of the system. We present the renormalized stochastic Langevin equation for an homogeneous system-field and then we analyze the influence of the environment on the power spectrum for some modes in the system. Section VI contains our final remarks. Two short appendices fill in some of the detail. Throughout the paper we use units such that $`\mathrm{}=c=1`$.
## II The Influence functional and the density matrix
Let us consider a massless quantum scalar field, minimally coupled to a de Sitter spacetime $`ds^2=a(\eta )[d\eta ^2d\stackrel{}{x}^2]`$ (where $`\eta `$ is the conformal time, $`d\eta =dt/a(t)`$ with $`t`$ the cosmic time), with a quartic self-interaction. The classical action is given by
$$S[\varphi ]=d^4xa^4(\eta )\left[\frac{\varphi _{}^{}{}_{}{}^{2}}{2a^2(\eta )}\frac{\varphi ^2}{2a^2(\eta )}\lambda \varphi ^4\right],$$
(2)
where $`a(\eta )=1/(H\eta )`$ and $`\varphi ^{}=d\varphi /d\eta `$ ($`a(\eta _i)=1`$ \[$`\eta _i=H^1`$\], and $`H^1`$ is the Hubble radius). Let us make a system-environment field splitting lombmazz
$$\varphi =\varphi _<+\varphi _>,$$
(3)
where the system-field $`\varphi _<`$ contains the modes with wave vectors shorter than a critical value $`\mathrm{\Lambda }2\pi /\lambda _c`$, while the environment-field $`\varphi _>`$ contains wave vectors longer than $`\mathrm{\Lambda }`$. As we set $`a(\eta _i)=1`$, a physical length $`\lambda _{\mathrm{phys}}=a(\eta )\lambda `$ coincides with the corresponding comoving length $`\lambda `$ at the initial time. Therefore, the splitting between system and environment gives a system sector constituted by all the modes with physical wavelengths shorter than the critical length $`\lambda _c`$ at the initial time $`\eta _i`$.
After splitting, the total action (2) can be written as
$$S[\varphi ]=S_0^<[\varphi _<]+S_0^>[\varphi _>]+S_{\mathrm{int}}[\varphi _<,\varphi _>],$$
(4)
where $`S_0`$ denotes the free field action and the interaction terms are given by
$`S_{\mathrm{int}}[\varphi _<,\varphi _>]=\lambda {\displaystyle }d^4xa^4(\eta )\{\varphi _<^4(x)+\varphi _>^4(x)`$ (5)
$`+`$ $`6\varphi _<^2(x)\varphi _>^2(x)+4\varphi _<^3(x)\varphi _>(x)+4\varphi _<(x)\varphi _>^3(x)\}.`$
The total density matrix elements (for the system and environment fields) are defined as
$$\rho [\varphi _<^+,\varphi _>^+|\varphi _<^{},\varphi _>^{};\eta ]=\varphi _<^+\varphi _>^+|\widehat{\rho }[\eta ]|\varphi _<^{}\varphi _>^{},$$
(6)
where $`|\varphi _<^\pm `$ and $`|\varphi _>^\pm `$ are the eigenstates of the field operators $`\widehat{\varphi }_<`$ and $`\widehat{\varphi }_>`$, respectively. For simplicity, we will assume that the interaction is turned on at the initial time $`\eta _i`$ and that, at this time, the system and the environment are not correlated (we ignore, for the moment, the physical consequences of such a choice, it has been discussed in deconpb ). Therefore, the total density operator can be written as the product of the density operator for the system and for the environment
$$\widehat{\rho }[\eta _i]=\widehat{\rho }_>[\eta _i]\widehat{\rho }_<[\eta _i].$$
(7)
We will further assume that the initial state of the environment is the Bunch-Davies vacuum BirreLLDavies .
We are interested in the influence of the environment on the evolution of the system. Therefore the reduced density matrix is the object of relevance. It is defined by
$$\rho _\mathrm{r}[\varphi _<^+|\varphi _<^{};\eta ]=𝒟\varphi _>\rho [\varphi _<^+,\varphi _>|\varphi _<^{},\varphi _>;\eta ].$$
(8)
The reduced density matrix evolves in time by means of
$`\rho _\mathrm{r}[\varphi _{<f}^+|\varphi _{<f}^{};\eta ]`$ $`=`$ $`{\displaystyle 𝑑\varphi _{<i}^+𝑑\varphi _{<i}^{}\rho _r[\varphi _{<i}^+|\varphi _{<i}^{};\eta _i]}`$
$`\times `$ $`J_\mathrm{r}[\varphi _{<f}^+,\varphi _{<f}^{};\eta |\varphi _{<i}^+,\varphi _{<i}^{};\eta _i],`$
where $`J_\mathrm{r}`$ is the reduced evolution operator
$`J_\mathrm{r}[\eta |\eta _i]`$ $`=`$ $`{\displaystyle _{\varphi _{<i}^+}^{\varphi _{<f}^+}}𝒟\varphi _<^+{\displaystyle _{\varphi _{<i}^{}}^{\varphi _{<f}^{}}}𝒟\varphi _<^{}`$ (10)
$`\times `$ $`\mathrm{exp}\{i(S^<[\varphi _<^+]S^<[\varphi _<^{}])\}F[\varphi _<^+,\varphi _<^{}].`$
The influence functional (or Feynman-Vernon functional) $`F[\varphi _<^+,\varphi _<^{}]`$ is defined as
$`F[\varphi _<^+,\varphi _<^{}]={\displaystyle 𝑑\varphi _{>i}^+𝑑\varphi _{>i}^{}\rho _>[\varphi _{>i}^+,\varphi _{>i}^{},\eta _i]𝑑\varphi _{>f}}`$
$`\times {\displaystyle _{\varphi _{>i}^+}^{\varphi _{>f}}}𝒟\varphi _>^+{\displaystyle _{\varphi _{>i}^{}}^{\varphi _{>f}}}𝒟\varphi _>^{}\mathrm{exp}\left\{i(S[\varphi _>^+]+S_{\mathrm{int}}[\varphi _<^+,\varphi _>^+])\right\}`$
$`\times \mathrm{exp}\left\{i(S[\varphi _>^{}]+S_{\mathrm{int}}[\varphi _<^{},\varphi _>^{}])\right\}.`$ (11)
This functional takes into account the effect of the environment on the system. The influence functional describes the averaged effect of the environmental degrees of freedom on the system degrees of freedom to which they are coupled. With this functional, one can identify a noise and dissipation kernel related by some kind of fluctuation-dissipation relation. This relation is important when one is interested in possible stationary states where a balance is eventually reached. During inflation we have the field (inflaton) on a very flat potential, away from its minimum, and we are, in general, only interested in the dynamics over some relatively small time. For example, we would neglect dissipation during the slow-roll period; but it is not correct during the eventual reheating phase.
We define the influence action $`\delta A[\varphi _<^+,\varphi _<^{}]`$ and the coarse grained effective action (CGEA) $`A[\varphi _<^+,\varphi _<^{}]`$ as
$$F[\varphi _<^+,\varphi _<^{}]=\mathrm{exp}\{i\delta A[\varphi _<^+,\varphi _<^{}]\},$$
(12)
$$A[\varphi _<^+,\varphi _<^{}]=S[\varphi _<^+]S[\varphi _<^{}]+\delta A[\varphi _<^+,\varphi _<^{}].$$
(13)
We will calculate the influence action perturbatively in $`\lambda `$ and we will consider only terms up to order $`\lambda ^2`$ and one loop in the $`\mathrm{}`$ expansion. The influence action has the following form lombmazz :
$`\delta A[\varphi _<^+,\varphi _<^{}]`$ $`=`$ $`S_{\mathrm{int}}[\varphi _>^+,\varphi _<^+]_0S_{\mathrm{int}}[\varphi _>^{},\varphi _<^{}]_0`$
$``$ $`iS_{\mathrm{int}}[\varphi _>^+,\varphi _<^+]S_{\mathrm{int}}[\varphi _>^{},\varphi _<^{}]_0`$
$`+`$ $`iS_{\mathrm{int}}[\varphi _>^+,\varphi _<^+]_0S_{\mathrm{int}}[\varphi _>^{},\varphi _<^{}]_0`$
$`+`$ $`{\displaystyle \frac{i}{2}}\left\{S_{\mathrm{int}}[\varphi _>^+,\varphi _<^+]_0^2S_{\mathrm{int}}[\varphi _>^+,\varphi _<^+]^2_0\right\}`$
$`+`$ $`{\displaystyle \frac{i}{2}}\left\{S_{\mathrm{int}}[\varphi _>^{},\varphi _<^{}]_0^2S_{\mathrm{int}}[\varphi _>^{},\varphi _<^{}]^2_0\right\}`$
where $`_0`$ is the quantum average with respect to the free field action of the environment, defined as
$`B[\varphi _>^+,\varphi _>^{}]_0={\displaystyle 𝑑\varphi _{>i}^+𝑑\varphi _{>i}^{}𝑑\varphi _{>f}}`$ (15)
$`\times {\displaystyle _{\varphi _{>i}^+}^{\varphi _{>f}}}D\varphi _>^+{\displaystyle _{\varphi _{>i}^{}}^{\varphi _{>f}}}D\varphi _>^{}B[\varphi _>^+,\varphi _>^{}]`$
$`\times \mathrm{exp}\left\{i\left(S^>[\varphi _>^+]S^>[\varphi _>^{}]\right)\right\}\varphi _{>i}^+|\widehat{\rho }_>[\eta _i]|\varphi _{>i}^{}.`$
Here $`\widehat{\rho }_>`$ is the Bunch-Davies vacuum state assumed for the environment.
The influence functional can be computed, and the result is
$`\mathrm{Re}\delta A`$ $`=`$ $`\lambda {\displaystyle }d^4x_1a^4(\eta )\{2\mathrm{\Delta }_4(x_1)`$ (16)
$``$ $`12\mathrm{\Delta }_2(x_1)iG_{++}^\mathrm{\Lambda }(x_1,x_1)\}`$
$`+`$ $`\lambda ^2{\displaystyle d^4x_1d^4x_2a^4(\eta _1)a^4(\eta _2)\mathrm{\Theta }(\eta _1\eta _2)}`$
$`\times `$ $`\{64\mathrm{\Delta }_3(x_1)\mathrm{Re}G_{++}^\mathrm{\Lambda }(x_1,x_2)\mathrm{\Sigma }_3(x_2)`$
$`+`$ $`288\mathrm{\Delta }_2(x_1)\mathrm{Im}G_{++}^{\mathrm{\Lambda }2}(x_1,x_2)\mathrm{\Sigma }_2(x_2)\},`$
$`\mathrm{Im}\delta A`$ $`=`$ $`\lambda ^2{\displaystyle d^4x_1d^4x_2a^4(\eta _1)a^4(\eta _2)}`$ (17)
$`\times `$ $`\{32\mathrm{\Delta }_3(x_1)\mathrm{Im}G_{++}^\mathrm{\Lambda }(x_1,x_2)\mathrm{\Delta }_3(x_2)`$
$``$ $`144\mathrm{\Delta }_2(x_1)\mathrm{Re}G_{++}^{\mathrm{\Lambda }2}(x_1,x_2)\mathrm{\Delta }_2(x_2)\},`$
where $`x_j`$ denotes ($`\eta _j`$,$`\stackrel{}{x}_j`$), $`\mathrm{\Theta }(x)`$ is the Heaviside step function, and the integrations in time run from $`\eta _i`$ to $`\eta `$. $`G_{++}^\mathrm{\Lambda }(x_1,x_2)i\varphi _>^+(x_1)\varphi _>^+(x_2)_0`$ is the relevant short-wavelength closed time-path correlator (it is proportional to the Feynmann propagator of the environment field, where the integration over momenta is restricted by the presence of the infrared cut-off $`\mathrm{\Lambda }`$), and we have defined
$$\mathrm{\Delta }_n=\frac{1}{2}(\varphi _<^{+n}\varphi _<^n),\mathrm{\Sigma }_n=\frac{1}{2}(\varphi _<^{+n}+\varphi _<^n),$$
(18)
with $`n=1,2,3`$.
## III Master equation and diffusion coefficients
In this Section we obtain the evolution equation for the reduced density matrix (master equation), paying particular attention to the diffusion terms, which are responsible for decoherence. To do so, we closely follow the quantum Brownian motion (QBM) example unruh ; qbm , translated into quantum field theory lomplb ; lombmazz .
The first step in the evaluation of the master equation is the calculation of the density matrix propagator $`J_\mathrm{r}`$ from Eq.(10). In order to solve the functional integration which defines the reduced propagator, we perform a saddle point approximation, assuming the classical field configuration dominates functional integrals,
$$J_\mathrm{r}[\varphi _{<f}^+,\varphi _{<f}^{},\eta |\varphi _{<i}^+,\varphi _{<i}^{},\eta _i]\mathrm{exp}iA[\varphi _{<\mathrm{cl}}^+,\varphi _{<\mathrm{cl}}^{}],$$
(19)
where $`\varphi _{<\mathrm{cl}}^\pm `$ is the solution of the semiclassical equation of motion $`\delta ReA/\delta \varphi _<^+|_{\varphi _<^+=\varphi _<^{}}=0`$ with boundary conditions $`\varphi _{<\mathrm{cl}}^\pm (\eta _i)=\varphi _{<i}^\pm `$ and $`\varphi _{<\mathrm{cl}}^\pm (\eta )=\varphi _{<f}^\pm `$. Since we are working up to $`\lambda ^2`$ order, we can evaluate the influence functional using the solutions of the free field equations. This classical equation is $`\varphi _<^{^{\prime \prime }}+2\varphi _<^{^{}}^2\varphi _<=0,`$ ($`=a^{}(\eta )/a(\eta )`$). A Fourier mode $`\psi _\stackrel{}{k}`$ of the field $`\psi a(\eta )\varphi _<`$, satisfies
$$\psi _\stackrel{}{k}^{^{\prime \prime }}+\left(k^2\frac{2}{\eta ^2}\right)\psi _\stackrel{}{k}=0,$$
(20)
where we have used the fact that $`a^{\prime \prime }/a=2/\eta ^2`$. It is important to note that for long-wavelength modes, $`k2/\eta ^2`$, Eq. (20) describes an unstable (upside-down) harmonic oscillator guthpi .
The classical solution for a mode in the system can be written as
$$\varphi _\stackrel{}{k}^{\pm \mathrm{cl}}(\eta ^{})=\varphi _{<i}^\pm (\stackrel{}{k})u_1(\eta ^{},\eta )+\varphi _{<f}^\pm (\stackrel{}{k})u_2(\eta ^{},\eta ),$$
(21)
where
$`u_1=`$ $`{\displaystyle \frac{\mathrm{sin}[k(\eta \eta ^{})](\frac{1}{k}+k\eta \eta ^{})+\mathrm{cos}[k(\eta \eta ^{})](\eta ^{}\eta )}{\mathrm{sin}[k(\eta ^{}\eta _i)](\frac{1}{k}+k\eta _i\eta ^{})+\mathrm{cos}[k(\eta ^{}\eta _i)](\eta _i\eta ^{})}},`$
$`u_2=`$ $`{\displaystyle \frac{\mathrm{sin}[k(\eta _i\eta ^{})](\frac{1}{k}+k\eta ^{}\eta _i)+\mathrm{cos}[k(\eta ^{}\eta _i)](\eta ^{}\eta _i)}{\mathrm{sin}[k(\eta _i\eta )](\frac{1}{k}+k\eta \eta _i)+\mathrm{cos}[k(\eta \eta _i)](\eta \eta _i)}}.`$
We will assume that the system-field contains only one Fourier mode with $`\stackrel{}{k}=\stackrel{}{k}_0`$. This is a sort of “minisuperspace” approximation for the system-field that will greatly simplify the calculations, therefore we assume
$$\varphi _{<\mathrm{cl}}^\pm (\stackrel{}{x},\eta ^{})=\varphi _{\stackrel{}{k}_0}^{\pm \mathrm{cl}}(\eta ^{})\mathrm{cos}(\stackrel{}{k}_0.\stackrel{}{x}),$$
(23)
where $`\varphi _{\stackrel{}{k}_0}^{\pm \mathrm{cl}}`$ is given by (21).
In order to obtain the master equation we must compute the final time derivative of the propagator $`J_\mathrm{r}`$. After that, all the dependence on the initial field configurations $`\varphi _{<i}^\pm `$ (coming from the classical solutions $`\varphi _<^{\pm \mathrm{cl}}`$) must be eliminated. Following the same procedure outlined in previous publications deconpb , we can prove that the free propagator satisfies
$`\varphi _{k_0}^{\pm \mathrm{cl}}(\eta ^{})J_0`$ $`=`$ $`\{u_2(\eta ^{},\eta )\varphi _{<f}^\pm `$
$``$ $`{\displaystyle \frac{2a^2(\eta )u_2^{}(\eta ,\eta )u_1(\eta ^{},\eta )\varphi _{<f}^\pm }{a^2(\eta )u_1^{}(\eta ,\eta )a^2(\eta _i)u_2^{}(\eta _i,\eta )}}`$
$``$ $`{\displaystyle \frac{4iu_1(\eta ^{},\eta )V^1}{[a^2(\eta )u_1^{}(\eta ,\eta )a^2(\eta _i)u_2^{}(\eta _i,\eta )]}}_{\varphi _{<\mathrm{f}}^\pm }\}J_0,`$
where a prime now stands for a derivative with respect to $`\eta ^{}`$ and the spatial volume $`V`$ appears because of we are considering only one Fourier mode for the system. These identities allow us to remove the initial field configurations $`\varphi _{<i}^\pm `$, by expressing them in terms of the final amplitudes $`\varphi _{<f}^\pm `$ and the derivatives $`_{\varphi _{<f}^\pm }`$, and obtain the master equation.
The full equation is very complicated and, as for quantum Brownian motion, it depends on the system-environment coupling. In what follows we will compute the diffusion coefficients for the different couplings described in the previous section. As we are solely interested in decoherence, it is sufficient to calculate the correction to the usual unitary evolution coming from the imaginary part of the influence action. The result reads
$`i`$ $`_\eta \rho _r[\varphi _{<f}^+|\varphi _{<f}^{};\eta ]=\varphi _{<f}^+|[\widehat{H}_{\mathrm{ren}},\widehat{\rho }_r]|\varphi _{<f}^{}`$ (25)
$``$ $`i\left[\mathrm{\Gamma }_1D_1(\stackrel{}{k_0},\eta ,\mathrm{\Lambda })+\mathrm{\Gamma }_2D_2(\stackrel{}{k_0},\eta ,\mathrm{\Lambda })\right]\rho _r[\varphi _{<f}^+|\varphi _{<f}^{};\eta ]`$
$`+`$ $`\mathrm{},`$
where we have defined $`\mathrm{\Gamma }_1=\frac{\lambda ^2V}{H^2}(\varphi _{<f}^{+}{}_{}{}^{3}\varphi _{<f}^{}{}_{}{}^{3})^2`$ and $`\mathrm{\Gamma }_2=\frac{\lambda ^2V}{4}(\varphi _{<f}^{+}{}_{}{}^{2}\varphi _{<f}^{}{}_{}{}^{2})^2`$. The ellipsis denotes other terms coming from the time derivative that not contribute to the diffusive effects. This equation contains time-dependent diffusion coefficients $`D_j`$. Up to one loop, only $`D_1`$ and $`D_2`$ survive. Coefficient $`D_1`$ is related to the interaction term $`\varphi _<^3\varphi _>`$, while $`D_2`$ to $`\varphi _<^2\varphi _>^2`$. These coefficients can be (formally) written as
$`D_1(\stackrel{}{k_0},\eta ,\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{H^2}{2}}{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta )a^4(\eta ^{})F_{\mathrm{cl}}^3(\eta ,\eta ^{},k_0)`$
$`\times `$ $`ImG_{++}^\mathrm{\Lambda }(\eta ,\eta ^{},3\stackrel{}{k_0})\mathrm{\Theta }(3k_0\mathrm{\Lambda }),`$
and
$`D_2(\stackrel{}{k_0},\eta ,\mathrm{\Lambda })`$ $`=`$ $`36{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta )a^4(\eta ^{})F_{\mathrm{cl}}^2(\eta ,\eta ^{},k_0)`$
$`\times `$ $`[ReG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},2\stackrel{}{k_0})+2ReG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},0)],`$
with the function $`F_{\mathrm{cl}}`$ defined by
$$F_{\mathrm{cl}}(\eta ,\eta _i,k_0)=\frac{\mathrm{sin}[k_0(\eta \eta _i)]}{k_0\eta }+\frac{\eta _i\mathrm{cos}[k_0(\eta \eta _i)]}{\eta }.$$
(28)
The explicit expressions of these coefficients are complicated functions of conformal time, the particular mode $`k_0`$, and the cut-off $`\mathrm{\Lambda }`$, and we show them in Appendix A.
It is important to note that here we are only studying the effect of normal diffusion terms, even it is known that anomalous diffusion terms can also be relevant at zero temperature. Analysis done in Ref. qbm suggests that anomalous diffusion for a supraohmnic environment is only relevant on a small transient and decoherence for unstable long-wavelength modes are driven by normal diffusion coefficients paula .
## IV Decoherence
Coherences are destroyed by diffusion terms. This process is evident after considering the following appoximate solution to the master equation
$`\rho _\mathrm{r}`$ $`[\varphi _<^+,\varphi _<^{};\eta ]\rho _\mathrm{r}^\mathrm{u}[\varphi _<^+,\varphi _<^{};\eta ]`$ (29)
$`\times \mathrm{exp}\left[{\displaystyle \underset{j}{}}\mathrm{\Gamma }_\mathrm{j}{\displaystyle _{\eta _i}^{\eta _f}}𝑑\eta D_\mathrm{j}(k_0,\mathrm{\Lambda },\eta )\right],`$
where $`\rho _\mathrm{r}^\mathrm{u}`$ is the solution of the unitary part of the master equation (i.e. without environment), and $`\mathrm{\Gamma }_j`$ includes the coefficients in front each diffusion term in Eq.(25). The system will decohere when the non-diagonal elements of the reduced density matrix are much smaller than the diagonal ones.
The decoherence time-scale sets the time after which we have a classical field configuration, and it can be defined as the solution to
$`1`$ $``$ $`{\displaystyle \underset{j}{}}\mathrm{\Gamma }_j{\displaystyle _{\eta _i}^{\eta _d}}𝑑\eta D_j(k_0,\mathrm{\Lambda },\eta )`$
$``$ $`\mathrm{\Gamma }_l{\displaystyle _{\eta _i}^{\eta _d}}𝑑\eta D_l(k_0,\mathrm{\Lambda },\eta ),`$
where the inequality is valid for any particular $`j=l`$. That is, the interactions with the environment have a cumulative effect on the onset of classical behaviour, i.e. the inclusion of a further interaction term reduces the decoherence time $`\eta _d`$. Therefore, in order to find upper bounds to $`\eta _d`$, we define the decoherence time $`\eta _{d_j}`$ coming from each diffusion coefficient by
$$1\mathrm{\Gamma }_j_{\eta _i}^{\eta _{d_j}}𝑑\eta D_j(k_0,\mathrm{\Lambda },\eta ),$$
(31)
with $`j=1`$,$`2`$.
### IV.1 Diffusion terms: Numerical results and analytic approximations
In this subsection we will analyze the behaviour of each diffusion coefficient as a function of the Fourier mode $`k_0`$ (considered for the system in Eq.(23)) and the cut-off $`\mathrm{\Lambda }`$. We will also analyze the temporal evolution of the coefficients and their integration in time for fixed values of $`k_0`$ and $`\mathrm{\Lambda }`$. We will present simple analytical approximations to the coefficients which can be used in Eq.(31) instead of the full expressions to estimate the decoherence time-scale.
We define the dimensionless quantity $`𝒩[k_0,\eta ]`$ which is the number of e-foldings between the time $`\eta _{k_0}`$ when the mode $`k_0`$ crosses the Hubble radius (i.e., $`|k_0\eta _{k_0}|=1`$) and any time $`\eta `$ during inflation,
$$𝒩[k_0,\eta ]\mathrm{ln}\left|\frac{\eta _{k_0}}{\eta }\right|=\mathrm{ln}|k_0\eta |.$$
(32)
This quantity has the special feature that its sign indicates whether the mode is inside ($`𝒩[k_0,\eta ]<0`$) or outside ($`𝒩[k_0,\eta ]>0`$) the Hubble radius.
Let us first consider the diffusion coefficient $`D_1`$, which comes from the interaction term $`\varphi _<^3\varphi _>`$. Because of we are considering only one Fourier mode for the system, with wave vector $`\stackrel{}{k_0}`$, and the environment-field contains only modes with $`k>\mathrm{\Lambda }`$, this coefficient is different from zero only if $`\mathrm{\Lambda }/3<k_0<\mathrm{\Lambda }`$ (i.e., $`\varphi _<`$ is only coupled with the $`\stackrel{}{k}=3\stackrel{}{k_0}`$ mode of the environment).
For this coefficient we can obtain an exact analytical expression from Eq.(73) (see Appendix A). In Fig. 1 we have plotted this expression as a function of $`k_0`$ for a particular value of the conformal time ($`H\eta =1/2`$). For later times, the graphs are qualitatively similar but $`D_1`$ oscillates more rapidly (since we have obtained our results by perturbative calculations, they are not valid at large times). The coefficient decreases with $`k_0`$ and takes its maximum value for $`k_0\mathrm{\Lambda }/3`$, implying that this kind of interaction produce more decoherence for small values of $`k_0`$ and, from the above discussion, for small values of $`\mathrm{\Lambda }`$.
For practical purposes, we consider the following simple approximation to $`D_1`$:
$$D_1^{\mathrm{approx}}(k_0,\eta ,\mathrm{\Lambda })=\frac{1}{100}\frac{(1+H\eta )}{H^4\eta ^7k_0^3}\mathrm{\Theta }(3k_0\mathrm{\Lambda }).$$
(33)
As we can see from Fig. 2, this approximation is less close to the exact coefficient for big values of $`k_0`$. It is important to note from Eq.(31) that if the approximation is a lower bound to the coefficient, then it will be useful to calculate an upper bound to $`\eta _{d_1}`$.
According to the definition of $`\eta _{d_1}`$, given $`k_0`$, $`\mathrm{\Lambda }`$ and $`\mathrm{\Gamma }_1`$, we can estimate this time by integrating $`D_1`$ over the conformal time $`\eta `$ and plotting this temporal integral as a function of $`𝒩[k_0,\eta _{d_1}]`$. Fig. 3 shows such plots for two particular values of $`k_0/H`$, where we have added the curves of the approximation (33). As we have illustrated with these examples, the approximation to $`D_1`$ is useful to estimate the order of magnitude of the corresponding decoherence time $`\eta _{d_1}`$.
Let us now examine the behaviour of the coefficient $`D_2`$ which is associated with the interaction term $`\varphi _<^2\varphi _>^2`$. Since the interaction is now quadratic in $`\varphi _>`$, there are no restrictions on the values of $`k_0`$ such that $`D_20`$. Hence this coefficient can affect the coherence of all modes in the system, therefore it is the most important in our model.
The dependence of $`D_2`$ with $`\mathrm{\Lambda }`$ is showed in Fig. 4 for three different values of $`k_0`$ and a fixed time $`\eta =10^3H^1`$. As we can see in this figure, for modes with $`k_0<<\mathrm{\Lambda }`$ the coefficient is weakly dependent of the value of the critical wave vector $`\mathrm{\Lambda }`$.
In Figs. 5 and 6 we have plotted the coefficient $`D_2`$ as a function of $`k_0/H`$ for fixed values of $`\eta `$ and $`\mathrm{\Lambda }`$. Fig. 5 shows that modes within the Hubble radius ($`|k_0\eta |>1`$) the diffusion coefficient is an oscillatory function and it has a maximum when $`k_0\mathrm{\Lambda }`$. The same behaviour was noted for conformally coupled fields lombmazz . Physically, this fact can be interpreted in terms of particle creation in the environment due to its interaction with the system.
For these modes, a very simple but good approximation to $`D_2`$ is given by
$$D_2^{\mathrm{}}(\eta )=\frac{27}{2\pi }\frac{1}{(H\eta )^4},$$
(34)
where the upper-script $`\mathrm{}`$ stands for “local”, since it corresponds to approximate the whole diffusion coefficient by the term in Eq.(III) containing the Dirac delta function (see details in Appendix A). On the other hand, Fig. 6 shows the dependence of $`D_2`$ with $`k_0`$ for modes outside the Hubble radius at $`\eta `$. For these modes we note that the diffusive effects are more important for the smallest values of $`k_0`$, which are most sensitive at the expansion of the universe.
Note that the classical equation (20) for each mode of the free field $`\psi `$ (defined as $`\psi =a(\eta )\varphi `$) describes an stable (unstable) oscillator if the mode $`k_0`$ satisfies $`|k_0\eta |>>1`$ ($`|k_0\eta |<<1`$). Therefore, it was expectable that the diffusion coefficients mirrors this fact.
For modes far outside ($`|k_0\eta |<<1`$) the Hubble radius we have found an asymptotic approximation of $`D_2`$, useful if $`\mathrm{\Lambda }k_0`$, which can be written as
$$D_2^a(\eta ,\mathrm{\Lambda })=\frac{1}{\pi ^2H^4\mathrm{\Lambda }^3\eta ^7}.$$
(35)
Fig. 7 shows a graph of the coefficient $`D_2`$ as a function of $`𝒩[k_0,\eta _{d_2}]`$ for a particular value of $`k_0`$ and $`\mathrm{\Lambda }`$, and the curves of the two approximations above. From this figure, we can distinguish two different regimes: one is well described by the local approximation ($`|k_0\eta |>1`$) and the other by the asymptotic one ($`|k_0\eta |<<1`$). This behaviour of $`D_2`$ indicates that the decoherence process becomes faster once the mode crosses the Hubble radius. In addition, we can see that a better approximation to $`D_2`$ is obtained by averaging the local and the asymptotic ones, that is
$$D_2^{approx}=\frac{D_2^a+D_2^{\mathrm{}}}{2},$$
(36)
which is also showed in the same figure.
The dependence of $`D_2`$ with $`k_0`$ for a fixed value of $`𝒩[k_0,\eta ]`$ is shown in Fig. 8, in which we see that the approximation is less close to the numerical curve for big values of $`k_0`$, but it bounds $`D_2`$ from below.
In Figs. 9 and 10 we have plotted the temporal integral of the coefficient $`D_2`$ as a function of $`𝒩[k_0,\eta _{d_2}]`$ for different values of $`k_0`$ and $`\mathrm{\Lambda }`$, where we have also plotted the curves obtained using the approximate expression given in Eq.(36). In view of our analysis, it is reasonable to use that approximation to estimate the decoherence time-scale $`\eta _{d_2}`$, at least for values of $`\mathrm{\Lambda }`$ smaller than $`10^3k_0`$.
### IV.2 The decoherence time
So far in this Section we have analyzed the behaviour of the diffusion coefficients, and we have shown the approximations considered are useful to estimate the decoherence time associated with each diffusion coefficient (or interaction term). Let us now work at the level of the order of magnitude and apply these results to quantify the decoherence time $`\eta _d`$.
As we have already mentioned, $`\eta _{d_1}`$ and $`\eta _{d_2}`$ are upper bounds to $`\eta _d`$. The time-scale given by $`\eta _d(k_0)`$ sets the time after which we are able to distinguish between two different amplitudes of the Fourier mode $`k_0`$ within the volume $`V`$. Thus, the maximum value of $`\eta _d(k_0)`$ with $`k_0<\mathrm{\Lambda }`$ corresponds to an upper bound to the decoherence time for the system-field (since, in principle, $`\varphi _<`$ contains all these modes with amplitudes different from zero).
In order to quantify decoherence times $`\eta _{d_{1,2}}`$ we have to fix the values of $`\mathrm{\Gamma }_{1,2}`$ (i.e., we have to assume values to $`\lambda `$, $`V`$, $`\varphi _{<f}^+`$, and $`\varphi _{<f}^{}`$). For this, as a first approximation and since we are considering a fixed de Sitter background, we will assume that the slow-roll conditions are satisfied at least up to times of the order of the decoherence time-scale. We will choose typical values for the parameters of the model ($`\lambda ,V`$) and for the elements of the reduced density matrix ($`\varphi _{<f}^+`$,$`\varphi _{<f}^{}`$).
The slow-roll conditions are usually written as Peacock ; langlois :
$$\epsilon _U=\frac{m_{\mathrm{pl}}^2}{16\pi }\left(\frac{U^{}}{U}\right)^2<<1,\eta _U=\frac{m_{\mathrm{pl}}^2}{8\pi }\left(\frac{U^{\prime \prime }}{U}\right)<<1,$$
(37)
where $`m_{\mathrm{pl}}G^{1/2}`$ is the Plank mass, $`U=\lambda \varphi _0^4`$, $`U^{}=dU/d\varphi _0`$, and $`\varphi _0`$ is the classical inflaton field. If we also assume that $`\varphi _0`$ is homogeneous on physical scales $`\lambda _{\mathrm{ph}}>>(\sqrt{\lambda }\varphi _0)^1`$, these conditions imply that the classical configuration field $`\varphi _0`$ and the Hubble rate satisfy
$$H^2\frac{8\pi U}{3m_{\mathrm{pl}}^2};\frac{d\varphi _0}{dt}\frac{U^{}}{3H}.$$
(38)
Defining the end of the inflationary period setting $`ϵ_U1`$, one can set $`\varphi _0(N_\eta )\sqrt{(N_\eta +1)/\pi }m_{\mathrm{pl}}`$, where $`N_\eta =\mathrm{ln}a(\eta _f)/a(\eta )`$; typically the e-fold number $`N_{\eta _i}N60`$ langlois . Thus, we assume the mean value of the system-field at time of decoherence is $`\varphi _0(N_{\eta _d})`$ ($`\mathrm{\Sigma }_{\varphi _{<f}}(\varphi _{<f}^++\varphi _{<f}^{})/2\varphi _0(N_{\eta _d})`$). Using this mean value and Eq.(38) we can write $`H^28\lambda N_{\eta _d}^2m_{\mathrm{pl}}^2/(3\pi )`$. In order to choose a typical order of magnitude for the fluctuations of the system-field, let us use the fact that the amplitude of the so-called primordial density perturbations $`\delta `$ can be inferred to be of order $`10^5`$, and that $`\delta H/\dot{\varphi _0}\delta \varphi `$, where $`\delta \varphi `$ is the amplitude of the inflaton field fluctuations Peacock ; langlois . Thus, with the use of these constraints and Eq.(38) we set $`\mathrm{\Delta }_{\varphi _{<f}}(\varphi _{<f}^+\varphi _{<f}^{})10^5\varphi _0(N_{\eta _d})/N_{\eta _d}`$. Since $`V`$ is the spatial volume inside which there are no coherent superpositions of macroscopically distinguishable states for the system-field, it is reasonable to choose $`V=vH^3`$, with $`v1`$ (which corresponds, if $`N60`$, to the comoving volume $`V(a_0H_0)^3`$, where $`H_0`$ and $`a_0`$ are the present values of the Hubble rate and the scalar factor, respectively) .
From previous considerations and according to the condition (31), the times $`\eta _{d_{1,2}}`$ can be estimated as the solution to
$`{\displaystyle _{\eta _i}^{\eta _{d_1}}}H𝑑\eta D_1(\eta )`$ $`{\displaystyle \frac{H}{\mathrm{\Gamma }_1}}2\times 10^{15}{\displaystyle \frac{(\lambda 10^5)}{v}}\left({\displaystyle \frac{N}{60}}\right)^5`$ (39a)
$`{\displaystyle _{\eta _i}^{\eta _{d_2}}}H𝑑\eta D_2(\eta )`$ $`{\displaystyle \frac{H}{\mathrm{\Gamma }_2}}9\times 10^{17}{\displaystyle \frac{1}{v}}\left({\displaystyle \frac{N}{60}}\right)^4,`$ (39b)
where we have used the fact that $`N_{\eta _d}N`$. The last term in (39a) follows after taking $`(\varphi _{<f}^{+3}\varphi _{<f}^3)\mathrm{\Delta }_{\varphi _{<f}}\mathrm{\Sigma }_{\varphi _{<f}}^2`$. For example, from Fig. 10 we can see that the temporal integral of $`D_2`$ takes a value of order $`10^{17}`$ for $`𝒩[k_0=H,\eta _{d_2}]`$ between $`7`$ and $`8`$. As $`D_2`$ is weakly dependent with $`\mathrm{\Lambda }`$ for $`k_0<H`$, we can say that it is valid for $`\mathrm{\Lambda }H`$. Thus, since $`D_2`$ decrease with $`k_0`$ for $`|k_0\eta |<1`$, we can conclude that the decoherence time for those modes with $`k_0<H`$ are smaller.
Substituting the approximation given in Eq.(33) into the left-hand side of (39a) and assuming that $`|H\eta _{d_1}|<<1`$, we get
$`t_{d_1}`$ $``$ $`{\displaystyle \frac{1}{6H}}\mathrm{ln}\left(600\alpha ^3{\displaystyle \frac{H}{\mathrm{\Gamma }_1}}\right)`$
$``$ $`{\displaystyle \frac{7}{H}}+{\displaystyle \frac{1}{6H}}\mathrm{ln}\left(\lambda 10^5\right)+{\displaystyle \frac{1}{6H}}\mathrm{ln}\left({\displaystyle \frac{\alpha ^3}{v}}\left({\displaystyle \frac{N}{60}}\right)^5\right),`$
where $`t_{d_1}=H^1\mathrm{ln}a(\eta _{d_1})`$, and we have defined $`k_0=\alpha H`$.
With the use of the approximation in Eq.(36) we obtain a quadratic equation in $`xa^3(\eta _{d_2})=\mathrm{exp}(3Ht_{d_2})`$, which is simple to solve for $`t_{d_2}`$. The result is
$`x`$ $``$ $`{\displaystyle \frac{27}{4}}\pi \sigma ^3\left(\sqrt{1+{\displaystyle \frac{8}{27\sigma ^3}}\left[{\displaystyle \frac{8H}{9\mathrm{\Gamma }_2}}+{\displaystyle \frac{1}{\pi }}+{\displaystyle \frac{2}{27\pi ^2\sigma ^3}}\right]}1\right)`$ (41)
$`<`$ $`{\displaystyle \frac{27}{4}}\pi \sigma ^3\left(\sqrt{1+{\displaystyle \frac{8}{27\sigma ^3}}\left[{\displaystyle \frac{8H}{9\mathrm{\Gamma }_2}}+{\displaystyle \frac{1}{\pi }}+{\displaystyle \frac{2}{27\pi ^2\sigma ^3}}\right]}\right),`$
where $`\sigma \mathrm{\Lambda }/H`$. For the sake of simplicity let us set $`\sigma 1`$. Since $`H/\mathrm{\Gamma }_2>>1`$, it yields
$`t_{d_2}`$ $``$ $`{\displaystyle \frac{1}{6H}}\mathrm{ln}\left(12\pi ^2\sigma ^3{\displaystyle \frac{H}{\mathrm{\Gamma }_2}}\right)`$ (42)
$``$ $`{\displaystyle \frac{7.7}{H}}+{\displaystyle \frac{1}{6H}}\mathrm{ln}\left({\displaystyle \frac{\sigma ^3}{v}}\left({\displaystyle \frac{N}{60}}\right)^4\right).`$
Assuming $`N=Ht_{\mathrm{end}}60`$ as an estimative scale to the end of inflationary period and values of $`\lambda 10^5`$, we obtain
$$\frac{t_{d_1}}{t_{\mathrm{end}}}\frac{7}{60}+\frac{1}{120}\mathrm{ln}\left(\frac{\alpha }{v^{1/3}}\right),$$
(43)
which makes sense only if $`k_0<\mathrm{\Lambda }<3k_0`$ ($`\alpha <\sigma <3\alpha `$), and
$$\frac{t_{d_2}}{t_{\mathrm{end}}}\frac{2}{15}+\frac{1}{120}\mathrm{ln}\left(\frac{\sigma }{v^{1/3}}\right).$$
(44)
From scales $`t_{d_1}`$ and $`t_{d_2}`$ we conclude that if one set $`\mathrm{\Lambda }H`$, the decoherence time-scale for the system-field is shorter than the minimal duration of inflation for all the wave-vectors within the system sector.
## V Effective dynamical evolution of the system-field in a fixed de Sitter background
In this Section we concern ourselves with the time evolution of the system-field. After reviewing the phenomenological way to describe the stochastic dynamical evolution of the system, we present the renormalized “semiclassical-Langevin” equation. This equation can be used to describe the dynamical evolution of the classical configurations of the system-field and hence is useful once all the modes in the system have lost coherence. As a first step to understand the generation of classical inhomogeneities from quantum fluctuations, we then consider a simple situation in which we analyze the influence of the environment on the power spectrum for modes inside the system sector.
### V.1 Effective dynamical evolution, noise and expectation values
The real part of the influence action contains divergent terms and should be renormalized. The imaginary part is finite and is associated with the decoherence process. It is well known that the terms of the imaginary part that come from a given interaction term in the original action can be viewed as arising from a noise source lomplb ; GMuler . In our case there are two such sources $`\xi _2`$ and $`\xi _3`$, which are associated with the interaction terms $`\varphi _<^2\varphi _>^2`$ and $`\varphi _<^3\varphi _>`$ respectively. That is, the imaginary part of the influence action can be rewritten as
$$Im\delta A=\mathrm{ln}\left(F[\mathrm{\Delta }_2]F[\mathrm{\Delta }_3]\right),$$
(45)
where $`F[\mathrm{\Delta }_n]`$ ($`n=2,3`$) is the characteristic functional of the noise $`\xi _n`$, which is related with Gaussian functional probability distribution $`P[\xi _n]`$ as
$`F[\mathrm{\Delta }_n]`$ $`=`$ $`{\displaystyle D\xi _nP[\xi _n]\mathrm{exp}\left\{id^4x\mathrm{\Delta }_n(x)\xi _n(x)\right\}},`$
$`P[\xi _n]`$ $`=`$ $`N_n\mathrm{exp}\{{\displaystyle \frac{1}{2}}{\displaystyle }d^4x_1{\displaystyle }d^4x_2\xi _n(x_1)`$ (46)
$`\times `$ $`\nu _n^1(x_1,x_2)\xi _n(x_2)\},`$
where $`N_n`$ is a normalization factor and $`\nu _n^1`$ is the functional inverse of the noise kernel $`\nu _n`$:
$`\nu _2(x_1,x_2)`$ $`=288\lambda ^2a^4(\eta _1)a^4(\eta _2)ReG_{++}^{\mathrm{\Lambda }2}(x_1,x_2),`$ (47a)
$`\nu _3(x_1,x_2)`$ $`=64\lambda ^2a^4(\eta _1)a^4(\eta _2)ImG_{++}^\mathrm{\Lambda }(x_1,x_2).`$ (47b)
The Gaussian noise field $`\xi _n(x)`$ is completely characterized by
$`\xi _n(x)_P`$ $`=0,`$ (48a)
$`\xi _n(x_1)\xi _n(x_2)_P`$ $`=\nu _n(x_1,x_2),`$ (48b)
where with $`_P`$ we are denoting average over all realizations of $`\xi _n(x)`$.
The functional variation
$$\frac{\delta S_{\mathrm{eff}}}{\delta \varphi _<^+}|_{\varphi _<^+=\varphi _<^{}}=0,$$
(49)
yields the “semiclassical-Langevin” equation for the system-field, which is only valid once the system has become classical.
With the identifications above, the reduced density matrix can be rewritten as BLHU :
$`\rho _\mathrm{r}[\varphi _{<f}^+|\varphi _{<f}^{};\eta ]`$ $`=`$ $`{\displaystyle D\xi _2P[\xi _2]D\xi _3P[\xi _3]}`$
$`\times `$ $`\rho _\mathrm{r}[\varphi _{<f}^+|\varphi _{<f}^{},\xi _2,\xi _3;\eta ],`$
with
$`\rho _\mathrm{r}[\varphi _{<f}^+|\varphi _{<f}^{},\xi _2,\xi _3;\eta ]{\displaystyle 𝑑\varphi _{<i}^+𝑑\varphi _{<i}^{}_{\varphi _{<i}^+}^{\varphi _{<f}^+}D\varphi _<^+}`$
$`\times {\displaystyle _{\varphi _{<i}^{}}^{\varphi _{<f}^{}}}D\varphi _<^{}\rho _\mathrm{r}[\varphi _{<i}^+|\varphi _{<i}^{};\eta _i]\mathrm{exp}\left\{iS_{\mathrm{eff}}[\varphi _<^+,\varphi _<^{},\xi _2,\xi _3]\right\},`$
where the effective action $`S_{\mathrm{eff}}`$ is given by
$`S_{\mathrm{eff}}[\varphi _<^+,\varphi _<^{},\xi _2,\xi _3]`$ $`=`$ $`Re\{A[\varphi _<^+,\varphi _<^{}]\}`$
$``$ $`{\displaystyle d^4x[\mathrm{\Delta }_2(x)\xi _2(x)+\mathrm{\Delta }_3(x)\xi _3(x)]}.`$
Here $`A`$ is the CGEA specified in Eq.(13). Thus, the full expectation value of any operator $`\widehat{Q}[\varphi _<]`$ can be written as
$`\widehat{Q}[\varphi _<]`$ $`=`$ $`{\displaystyle D\xi _2P[\xi _2]D\xi _3P[\xi _3]𝑑\varphi _<}`$ (53)
$`\times `$ $`\rho _\mathrm{r}[\varphi _<|\varphi _<,\xi _2,\xi _3;\eta ]Q[\varphi _<]`$
$``$ $`\widehat{Q}[\varphi _<]_q_P,`$
where $`_q`$ is an usual quantum average for a system-field subjected to external stochastic forces.
### V.2 Renormalized equation of motion for the system
Taking the functional variation as in Eq.(49) we obtain
$`\varphi _<^{\prime \prime }(\eta ,\stackrel{}{x})\mathrm{}\varphi _<(\eta ,\stackrel{}{x})+2\varphi _<^{}(\eta ,\stackrel{}{x})`$
$`+\mathrm{\hspace{0.33em}4}\lambda a^2(\eta )\varphi _<(\eta ,\stackrel{}{x})[\varphi _<^2(\eta ,\stackrel{}{x})3iG_{++}^\mathrm{\Lambda }(\eta ,\eta ,\stackrel{}{0})]`$
$`\mathrm{\hspace{0.33em}96}\lambda ^2a^2(\eta ){\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta ^{})\varphi _<^2(\eta ,\stackrel{}{x})`$
$`\times {\displaystyle }d^3y\varphi _<^3(\eta ^{},\stackrel{}{y})ReG_{++}^\mathrm{\Lambda }(\eta ,\eta ^{},\stackrel{}{x}\stackrel{}{y})`$
$`\mathrm{\hspace{0.33em}288}\lambda ^2a^2(\eta ){\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta ^{})\varphi _<(\eta ,\stackrel{}{x})`$
$`\times {\displaystyle }d^3y\varphi _<^2(\eta ^{},\stackrel{}{y})ImG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},\stackrel{}{x}\stackrel{}{y})`$
$`=\xi _2(\eta ,\stackrel{}{x}){\displaystyle \frac{\varphi _<(\eta ,\stackrel{}{x})}{a^2(\eta )}}{\displaystyle \frac{3}{2}}\xi _3(\eta ,\stackrel{}{x}){\displaystyle \frac{\varphi _<^2(\eta ,\stackrel{}{x})}{a^2(\eta )}}.`$ (54)
This equation contains divergences. In order to renormalize it we use the method of adiabatic substraction with dimensional regularization, which works at the level of the field equation and is particularly useful for solving the equation numerically RenoPazMazziCarmen . To simplify the task we assume the system-field to be homogeneous enough so that the spatial derivatives and the term in Eq.(V.2) coming from the interaction $`\varphi _>\varphi _<^3`$ are negligible. Note that for an homogeneous system-field there is no contribution from this interaction term up to one loop order, due to orthogonality of the Fourier modes (see, for instance, Eq.(17)). Details of the renormalization procedure are relegated to Appendix B. The renormalized equation is
$`\varphi _<^{\prime \prime }(\eta )+[\mathrm{\Delta }M^2(\eta )+\mathrm{\Delta }\mathrm{\Sigma }(\eta )R]a^2(\eta )\varphi _<(\eta )`$ (55)
$`+\mathrm{\hspace{0.33em}2}\varphi _<^{}(\eta )+4[\lambda +\mathrm{\Delta }\stackrel{~}{\lambda }(\eta )]a^2(\eta )\varphi _<^3(\eta )`$
$`+{\displaystyle \frac{36\lambda ^2}{\pi ^2}}\varphi _<(\eta )a^2(\eta )\{{\displaystyle _{\eta _i}^\eta }d\eta ^{}\varphi _<^2(\eta ^{})𝒥(\eta ,\eta ^{})`$
$`+{\displaystyle _{\eta _i}^\eta }d\eta ^{}{\displaystyle \frac{\varphi _<(\eta ^{})\varphi _{}^{}{}_{<}{}^{}(\eta ^{})}{Ha(\eta )}}(\eta ,\eta ^{})\}=\stackrel{~}{\xi }_2(\eta ){\displaystyle \frac{\varphi _<(\eta )}{a^2(\eta )}}.`$
where $`R=12H^2`$; we have redefined the noise source as
$$\stackrel{~}{\xi }_2(\eta )=\frac{1}{V}_Vd^3x\xi _2(\eta ,\stackrel{}{x}),$$
(56)
whose correlation function is given by
$`\stackrel{~}{\xi }_2(\eta _1)\stackrel{~}{\xi }_2(\eta _2)_P=a^2(\eta _1)a^2(\eta _2){\displaystyle \frac{36\lambda ^2}{\pi ^2V}}\{{\displaystyle \frac{\pi }{2}}\delta [(\eta _1\eta _2)]`$
$`+{\displaystyle \frac{\mathrm{cos}[2\mathrm{\Lambda }(\eta _1\eta _2)]}{\mathrm{\Lambda }}}\left[{\displaystyle \frac{2}{\eta _1\eta _2}}+{\displaystyle \frac{(\eta _1\eta _2)^2}{3\eta _1^2\eta _2^2}}+{\displaystyle \frac{1}{3\mathrm{\Lambda }^2\eta _1^2\eta _2^2}}\right]`$
$`(\pi 2𝒮i[2\mathrm{\Lambda }|\eta _1\eta _2|])\left[{\displaystyle \frac{|\eta _1\eta _2|^3}{3\eta _1^2\eta _2^2}}+{\displaystyle \frac{|\eta _1\eta _2|}{\eta _1\eta _2}}\right]`$
$`\mathrm{sin}[2\mathrm{\Lambda }(\eta _1\eta _2)][{\displaystyle \frac{1}{2(\eta _1\eta _2)}}{\displaystyle \frac{2(\eta _1\eta _2)}{3\mathrm{\Lambda }^2\eta _1^2\eta _2^2}}]\};`$
and we have also defined the following functions:
$`\mathrm{\Delta }\stackrel{~}{\lambda }(\eta )`$ $`=\mathrm{\Delta }\lambda +{\displaystyle \frac{9\lambda ^2}{2\pi ^2}}\left[{\displaystyle \frac{\gamma }{2}}\mathrm{ln}\left|{\displaystyle \frac{a(\eta )\mu }{2\mathrm{\Lambda }}}\right|\right],`$ (57a)
$`\mathrm{\Delta }\mathrm{\Sigma }(\eta )`$ $`=\mathrm{\Delta }\xi {\displaystyle \frac{\lambda }{4\pi ^2}}\left[{\displaystyle \frac{\gamma }{2}}\mathrm{ln}\left|{\displaystyle \frac{a(\eta )\mu }{2\mathrm{\Lambda }}}\right|\right],`$ (57b)
$`\mathrm{\Delta }M^2(\eta )`$ $`=\mathrm{\Delta }m^2{\displaystyle \frac{18\lambda ^2}{\pi ^2}}{\displaystyle \frac{a^2(\eta _i)}{a^2(\eta )}}\varphi _<^2(\eta _i)`$ (57c)
$`\times 𝒞i[2\mathrm{\Lambda }(\eta \eta _i)]{\displaystyle \frac{3\lambda \mathrm{\Lambda }^2}{2\pi ^2a^2(\eta )}},`$
$`(\eta ,\eta ^{})`$ $`={\displaystyle \frac{\eta }{\eta _{}^{}{}_{}{}^{2}}}𝒞i[2\mathrm{\Lambda }(\eta \eta ^{})],`$ (57d)
$`𝒥(\eta ,\eta ^{})`$ $`=𝒥_1(\eta ,\eta ^{})+𝒥_2(\eta ,\eta ^{})+𝒥_3(\eta ,\eta ^{}),`$ (57e)
$`𝒥_1(\eta ,\eta ^{})`$ $`=𝒞i[2\mathrm{\Lambda }(\eta \eta ^{})]\left[{\displaystyle \frac{2(\eta ^3\eta _{}^{}{}_{}{}^{3})}{3\eta _{}^{}{}_{}{}^{4}}}+{\displaystyle \frac{\eta ^2}{\eta _{}^{}{}_{}{}^{3}}}\right],`$ (57f)
$`𝒥_2(\eta ,\eta ^{})`$ $`={\displaystyle \frac{2(\eta \eta ^{})}{3\eta _{}^{}{}_{}{}^{4}\mathrm{\Lambda }^2}}[\mathrm{cos}[2\mathrm{\Lambda }(\eta \eta ^{})]`$ (57g)
$`{\displaystyle \frac{\mathrm{sin}[2\mathrm{\Lambda }(\eta \eta ^{})]}{2\mathrm{\Lambda }(\eta \eta ^{})}}],`$
$`𝒥_3(\eta ,\eta ^{})`$ $`={\displaystyle \frac{\mathrm{sin}[2\mathrm{\Lambda }(\eta \eta ^{})]}{\mathrm{\Lambda }}}\left[{\displaystyle \frac{(\eta \eta ^{})^2}{3\eta _{}^{}{}_{}{}^{4}}}+{\displaystyle \frac{2\eta }{\eta _{}^{}{}_{}{}^{3}}}\right],`$ (57h)
where $`𝒞i[x]`$ is the cosine integral function Abramob . Notice that these functions are logarithmically divergent in the limit $`\mathrm{\Lambda }0`$, which is the well-known infrared divergence RenoPazMazziCarmen .
The functions above contain useful information. Particularly, they allows us to examine the conditions under which the loop expansion breaks down. In order to estimate the time after which the one-loop terms become of the same order of magnitude as the classical ones we may compare their time-dependent parts. For example, from Eq.(57) we see that the time-dependent part of $`\mathrm{\Delta }\stackrel{~}{\lambda }(\eta )`$ ($`\mathrm{\Delta }\mathrm{\Sigma }(\eta )`$) is the order $`\lambda `$ (one) for $`Ht1/\lambda `$ and $`\mathrm{\Delta }M^2(\eta )`$ is important only at the initial time ($`\eta \eta _i`$). In the same way, we can use the term $`\lambda a^2\varphi _<^3`$ to compare it with the remaining ones. In Fig. 11 we have plotted the $`𝒥_i`$ functions, where we can see that $`𝒥_1`$ dominates all others. It allow us to make the following approximation:
$`{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}\varphi _<^2(\eta ^{})𝒥(\eta ,\eta ^{})`$ $``$ $`\varphi _<^2(\eta ){\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}𝒥_1(\eta ,\eta ^{})`$ (58)
$``$ $`\varphi _<^2(\eta ){\displaystyle \frac{1}{3}}(\mathrm{ln}|\mathrm{\Lambda }\eta |)^2,`$
where the last term is a simple long-time asymptotic expression. With the use of this approximation we obtain that the two terms in question are of the same order when $`|\mathrm{ln}|\mathrm{\Lambda }\eta ||1/\sqrt{\lambda }`$ and, for $`\mathrm{\Lambda }H`$, $`Ht1/\sqrt{\lambda }`$.
As it is shown in Fig. 12, the $``$ function peaks at $`\eta \eta ^{}`$ and hence we can approximate
$`{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}{\displaystyle \frac{\varphi _<(\eta ^{})\varphi _{}^{}{}_{<}{}^{}(\eta ^{})}{Ha(\eta )}}`$ $``$ $`{\displaystyle \frac{\varphi _<(\eta )\varphi _{}^{}{}_{<}{}^{}(\eta )}{Ha(\eta )}}{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}`$
$``$ $`{\displaystyle \frac{\varphi _<(\eta )\varphi _{}^{}{}_{<}{}^{}(\eta )}{Ha(\eta )}}`$
$`\times `$ $`(\gamma +\mathrm{ln}|2\mathrm{\Lambda }\eta |),`$ (60)
where the last expression corresponds to a long-time approximation. As for the $`𝒥`$ term, using the long-time approximation we get $`\mathrm{ln}|\mathrm{\Lambda }\eta |\varphi _<H/(\lambda \dot{\varphi }_<)`$. If in addition we assume that the slow-roll conditions are satisfied, this time can be estimated as $`|\mathrm{ln}|\mathrm{\Lambda }\eta ||1/(\lambda \epsilon _U)`$ and, for $`\mathrm{\Lambda }H`$, $`Ht1/(\lambda \epsilon _U)`$, where $`\epsilon _U`$ is the slow-roll parameter given in Eq.(37).
From previous discussion, we can see that the additional terms are only important at times longer than the typical time-scale associated with the decoherence process.
### V.3 Generation of inhomogeneities: role of the noise
Having reached this point, we ask ourself about the primordial inhomogeneities. Certainly, the model we are considering is too simplified to obtain any property of the power spectrum of such inhomogeneities. Nevertheless, we can discuss some relevant aspects about the way of connecting the initial quantum fluctuations with those which eventually become classical and statistical. With this purpose, let us compute the power spectrum for some of the Fourier modes in the system working in an approximate way and making certain assumptions, but taking into account the decoherence process.
According to the results of the previous Section, the decoherence process occurs at different rates for each mode in the system. Provided that decoherence is more effective for modes outside the Hubble radius, we consider the following situation. We suppose that the “relevant” modes are in the Bunch-Davies state (i.e., in equilibrium with the environment) at the initial time $`\eta _i`$. Consequently, the matrix elements of the initial density operator for each of these modes are Gaussian functions in a mode-amplitude basis. As a first approximation we neglect the non-linearities which might affect the Gaussian form of those matrix elements.
We are now interested only in computing the power spectrum of those “relevant” modes up to $`\mathrm{}`$ and $`\lambda ^2`$ order. To carry this out, we split the system-field as $`\varphi _<=\varphi _0(\eta )+\delta \varphi _<`$, where we identify $`\varphi _0(\eta )`$ as a classical background field which satisfies slow-roll conditions. The power spectrum of the quantum field fluctuations $`\delta \varphi _<`$ may be expressed as $`P_\varphi (k)=2\pi ^2k^3\mathrm{\Delta }_\varphi ^2(k)`$, with $`\mathrm{\Delta }_\varphi ^2(k)`$ defined by
$$\delta \varphi _<(\stackrel{}{x})\delta \varphi _<(\stackrel{}{x}+\stackrel{}{r})=d^3k\frac{\mathrm{\Delta }_\varphi ^2(k)}{4\pi k^3}\mathrm{exp}(i\stackrel{}{k}\stackrel{}{r}),$$
(61)
where $``$ is the full average in Eq.(53).
Consistently with the assumption that the “relevant” matrix elements remain Gaussian, we expand the semiclassical equation (V.2) up to linear order in the mode amplitude of interest $`\delta \varphi _<(\stackrel{}{k})`$. Through this procedure we obtain
$`\varphi _0^{\prime \prime }(\eta )+2\varphi _0^{}(\eta )+4\lambda a^2\varphi _0^3(\eta )=0,`$ (62a)
$`\delta \varphi _<^{\prime \prime }(\stackrel{}{k},\eta )+[k^2+12\lambda a^2\varphi _0^2(\eta )]\delta \varphi _<(\stackrel{}{k},\eta )`$
$`+2\delta \varphi _<^{}(\stackrel{}{k},\eta )={\displaystyle \frac{\xi _2(\stackrel{}{k},\eta )}{a^2}}\varphi _0(\eta ),`$ (62b)
where we have discarded the terms which do not contribute to the power spectrum up to $`\mathrm{}`$ order. The term with the $`\xi _3`$ noise source gives a zero contribution due to our approximations and the orthogonality of the Fourier modes. It is important to note the presence of the $`\xi _2`$ noise source, which ensure the decoherence process occurs as the matrix elements are evolved in time.
In order to obtain the power spectrum let us split $`\delta \varphi _<(\stackrel{}{k},\eta )=\delta \varphi _<^\xi (\stackrel{}{k},\eta )+\delta \varphi _<^q(\stackrel{}{k},\eta )`$, with $`\delta \varphi _<^\xi (\stackrel{}{k},\eta )\delta \varphi _<(\stackrel{}{k},\eta )_q`$, where $`_q`$ is the quantum average defined in Eq.(53). Because of the assumption of linearity, this quantum average satisfies the semiclassical equation (62), whose solution can be written as the sum of the homogeneous solution $`\delta \varphi _<^h(k,\eta )`$ and a particular solution $`\delta \varphi _<^p(\stackrel{}{k},\eta )`$. The former is given by
$$\delta \varphi _<^h(k,\eta )=a^1(\eta )[\alpha _k\sqrt{|\eta |}J_\nu +\beta _k\sqrt{|\eta |}J_\nu ],$$
(63)
where $`\alpha _k`$ and $`\beta _k`$ are constants of integration, and $`\nu =\sqrt{\frac{9}{4}ϵ}`$, with $`ϵ=6\lambda \varphi _0^2/H^2`$. Setting $`\nu \frac{3}{2}`$, a particular solution is
$$\delta \varphi _<^p(\stackrel{}{k},\eta )=_{\eta _i}^\eta 𝑑\eta ^{}g(k,\eta ,\eta ^{})\xi _2(\stackrel{}{k},\eta ^{})\varphi _0(\eta ^{}),$$
(64)
where
$`g(k,\eta ,\eta ^{})`$ $`=`$ $`{\displaystyle \frac{1}{a(\eta )a(\eta ^{})}}[{\displaystyle \frac{\mathrm{sin}[k(\eta \eta ^{})]}{k}}(1+{\displaystyle \frac{1}{k^2\eta \eta ^{}}})`$ (65)
$``$ $`{\displaystyle \frac{\mathrm{cos}[k(\eta \eta ^{})]}{k^2\eta \eta ^{}}}(\eta \eta ^{})].`$
With the use of the initial conditions $`\delta \varphi _<(\stackrel{}{k},\eta _i)_q=\dot{\delta \varphi _<}(\stackrel{}{k},\eta _i)_q=0`$, we obtain that $`\delta \varphi _<^\xi (\stackrel{}{k},\eta )_P=0`$, and thus $`\delta \varphi _<^\xi (\stackrel{}{k},\eta )=\delta \varphi _<^p(\stackrel{}{k},\eta )`$. Within these approximations, the result is analogous to that for the linear quantum Brownian motion (QBM) qbm ; Morika . Therefore, the quantity $`\mathrm{\Delta }_\varphi ^2(k)`$ can be written so that it receives two contributions:
$$\mathrm{\Delta }_\varphi ^2(k)=\mathrm{\Delta }_{\varphi ^q}^2(k)+\mathrm{\Delta }_{\varphi ^\xi }^2(k).$$
(66)
The first one comes from the unitary evolution of the initial density matrix, i.e., it is the usual quantum result for the case of the free field langlois : $`\mathrm{\Delta }_{\varphi ^q}^2(k)=\left(H/2\pi \right)^2(1+k^2\eta ^2)`$. The second one appears due to the $`\xi _2`$ noise source and can be computed through
$`\delta \varphi _<^\xi (\stackrel{}{k_1},\eta )\delta \varphi _<^\xi (\stackrel{}{k_2},\eta )_P`$ $`=`$ $`{\displaystyle _{\eta _i}^\eta }𝑑\eta _1{\displaystyle _{\eta _i}^\eta }𝑑\eta _2\varphi _0(\eta _1)\varphi _0(\eta _2)`$ (67)
$`\times `$ $`\xi _2(\stackrel{}{k_1},\eta _1)\xi _2^{}(\stackrel{}{k_2},\eta _2)_P`$
$`\times `$ $`g(k_1,\eta ,\eta _1)g(k_2,\eta ,\eta _2),`$
where
$`\xi _2(\stackrel{}{k_1},\eta _1)\xi _2^{}(\stackrel{}{k_2},\eta _2)_P=288\lambda ^2a^4(\eta _1)a^4(\eta _1)`$
$`\times (2\pi )^3\delta ^3(\stackrel{}{k_1}\stackrel{}{k_2})ReG_{++}^{\mathrm{\Lambda }2}(\eta _1,\eta _2,\stackrel{}{k_1}).`$ (68)
Thus $`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ can be expressed as
$`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ $`=`$ $`\lambda ^2{\displaystyle \frac{144}{\pi ^2}}k^3{\displaystyle _{\eta _i}^\eta }𝑑\eta _1{\displaystyle _{\eta _i}^\eta }𝑑\eta _2a^4(\eta _1)a^4(\eta _2)`$ (69)
$`\times `$ $`\varphi _0(\eta _1)\varphi _0(\eta _2)g(k,\eta ,\eta _1)g(k,\eta ,\eta _2)`$
$`\times `$ $`ReG_{++}^{\mathrm{\Lambda }2}(\eta _1,\eta _2,\stackrel{}{k}).`$
The Fourier transform $`ReG_{++}^{\mathrm{\Lambda }2}(\eta _1,\eta _2,\stackrel{}{k})`$ (where $`k<\mathrm{\Lambda }`$) can be obtained from Eq.(79) (see Appendix A) replacing $`2k_0`$ by $`k`$.
Since the additional contribution $`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ to the power spectrum is of order $`\lambda ^2`$, it is expected to be negligible. As we are assuming that $`\varphi _0`$ is a slowly varying field, we can compare the relative order of magnitude of both contribution setting $`\varphi _0(\eta _1)\varphi _0(\eta _2)\varphi _0`$ and $`H^2\lambda \varphi _0^4/m_{\mathrm{pl}}^2`$, so that $`\varphi _0`$ can be taken out from the integration in Eq.(69). Thereby, when $`|k\eta |<1`$, the contributions $`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ is negligible compared to the usual $`\mathrm{\Delta }_{\varphi ^q}^2(k)`$ if $`(\mathrm{\Delta }_{\varphi ^\xi }(k)/\lambda \varphi _0)^2<<\lambda ^1(\varphi _0/m_{\mathrm{pl}})^2`$. Since the last quotient is usually much bigger than one, this condition is typically satisfied if $`(\mathrm{\Delta }_{\varphi ^\xi }(k)/\lambda \varphi _0)^2<1`$. As in the example shown in Fig. 13, this is the case and hence the additional contribution can be neglected.
On the other hand, the usual contribution $`\mathrm{\Delta }_{\varphi ^q}^2(k)`$ is independent of $`k`$ for a fixed value of $`k\eta `$, corresponding to a nearly scale-invariant spectrum, whereas $`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ depends on $`k`$ and $`\mathrm{\Lambda }`$ (see Figs. 14 and 15). To see this more clearly, it is useful to rewrite $`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ as
$`\mathrm{\Delta }_{\varphi ^\xi }^2(k)`$ $`=`$ $`{\displaystyle \frac{144}{\pi ^2}}{\displaystyle _{k\eta _i}^{k\eta }}{\displaystyle \frac{dx_1}{x_1^4}}{\displaystyle _{k\eta _i}^{k\eta }}{\displaystyle \frac{dx_2}{x_2^4}}\varphi _0(\eta _1)\varphi _0(\eta _2)`$ (70)
$`\times `$ $`f(k\eta ,x_1)f(k\eta ,x_2)F(x_1,x_2,{\displaystyle \frac{\mathrm{\Lambda }}{k}}),`$
where $`fk^3Hg`$ and $`Fk^3H^1ReG_{++}^{\mathrm{\Lambda }2}`$, which depend on $`k`$ not only through the $`k\eta `$ combination, but also through the combinations $`k\eta _i`$ and $`\mathrm{\Lambda }/k`$. It is well-known that a finite duration of the inflation stage produces a departure of the power spectrum from the scale-invariant one Boya2004 . The breaking of the scale invariance by the presence of the infrared cut-off $`\mathrm{\Lambda }`$ is similar to the one found in Ref. CalzGoron .
It is important to note that we have assumed a Gaussian initial condition for the reduced density matrix elements and that the Gaussian form of these elements is not affected by the non-linearities of the interactions. However, as was pointed out in Ref. diana , it is expectable that when the non-linearities become important, the Gaussian approximation breaks down and therefore the associated Wigner functional becomes non-positive. Nevertheless, as well as for a given mode the Gaussian approximation remains valid up to times longer than the decoherence time-scale, one can use the classical description for it even in the non-linear regime.
## VI Final remarks
Let us summarize the results contained in this paper. After the integration of the high frequency modes in Section II, we obtained the CGEA for the modes whose wave vector is shorter than a critical value $`\mathrm{\Lambda }`$. From the imaginary part of the CGEA we obtained, in Section III, the diffusion coefficients of the master equation. System and environment are two sectors of a single scalar field, and the results depend on the “size” of these sectors, which is fixed by the cut-off $`\mathrm{\Lambda }`$.
To analyze the decoherence process, in Section IV we evaluated the diffusion coefficients and its integrations over the conformal time. It allows us to conclude that if we consider a cut-off $`\mathrm{\Lambda }H`$, those modes with wave vector $`k_0\mathrm{\Lambda }`$ are the more affected by diffusion through one of the coefficients. For these modes, we saw that the effect is weakly dependent of the critical wave vector $`\mathrm{\Lambda }`$. If one consider a cut-off $`\mathrm{\Lambda }H`$, and modes $`H<k_0<\mathrm{\Lambda }`$, diffusive effects are larger for those modes in the system whose wave vector is close to the critical $`\mathrm{\Lambda }`$ lombmazz .
We presented the complete expression of the diffusive terms (in Appendix A), and also some simple analytical approximations to the coefficients, which are useful to make an evaluation of the decoherence time-scale. We performed an extensive analysis of the evaluation of the time-scale for the decoherence process for a typical case of interest. In such a case, we obtained that for a given mode $`k_0<\mathrm{\Lambda }H`$, decoherence is effective by the time in which inflation is ending.
In Section V we analyzed the effective evolution of the system. Assuming an homogeneous system-field we first presented the explicit form of the renormalized stochastic Langevin equation and we analyzed the relative importance of the terms appearing in that equation due to the system-environment interaction. From such analysis we conclude that those terms are of the same order of magnitude than the classical ones for times longer than the typical decoherence time-scale.
We them considered inhomogeneity generation in a simple particular situation, in which we analyzed the influence of the environment on the power spectrum for some modes in the system. In that situation, splitting the spectrum as the sum of the contribution coming from the unitary evolution of the Bunch-Davies initial condition plus the one which appears because of the system-environment interaction, we found that the latter is negligible compared to the former. In spite of this result, we have remarked that the system-environment interaction is essential to have a complete quantum to classical transition which allows a late-time classical treatment of the system degrees of freedom.
As for future work, we consider that it is worth applying the same procedure used in this paper to inflationary models involving two or more interacting fields, such as some hybrid model Linde1994 .
## VII Acknowledgments
This work is supported by UBA, CONICET and Fundación Antorchas; Argentina. We specially thank M. Zaldarriaga for very useful comments and suggestions. We also thank comments from E. Calzetta, B.L. Hu, C. Kiefer, F. Mazzitelli, N. Mavromatos, and R. Rivers.
## Appendix A Diffusion Coefficients
In this appendix we describe some technical details about the computation of the diffusion coefficients $`D_1`$ and $`D_2`$, starting from Eq.(III) and (III) respectively.
In order to evaluate the coefficient $`D_1`$, we have to perform the Fourier transform of the imaginary part of the propagator for the environment $`G_{++}^\mathrm{\Lambda }`$ (see Eq.(III)). This propagator can be expressed in terms of the mode functions $`\varphi _k`$ given by
$`\varphi _\stackrel{}{k}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle \frac{\mathrm{exp}\{ik\eta \}}{\sqrt{2k}a(\eta )}}\left(1{\displaystyle \frac{i}{k\eta }}\right)\mathrm{exp}\{i\stackrel{}{k}\stackrel{}{x}\}`$ (71)
$``$ $`\stackrel{~}{\varphi }_k(\eta )\mathrm{exp}\{i\stackrel{}{k}\stackrel{}{x}\},`$
which corresponds to the Bunch-Davies vacuum assumed for the environment-field. The imaginary part of the propagator reads
$$ImG_{++}^\mathrm{\Lambda }(x_1,x_2)=_{k>\mathrm{\Lambda }}d^3kRe\{\varphi _\stackrel{}{k}(x_1)\varphi _\stackrel{}{k}^{}(x_2)\}.$$
(72)
Substituting this expression into Eq.(III) we obtain
$`D_1(\stackrel{}{k_0},\eta ,\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{k_0^4}}{2}}{\displaystyle _x^{\stackrel{~}{k_0}}}dy[{\displaystyle \frac{\mathrm{cos}[3\mathrm{\Delta }_{xy}]}{6x^3y^3}}(1+{\displaystyle \frac{1}{9xy}})`$ (73)
$`+`$ $`{\displaystyle \frac{\mathrm{sin}[3\mathrm{\Delta }_{xy}]\mathrm{\Delta }_{xy}}{18x^4y^4}}]F^3_{cl}(x,y)\mathrm{\Theta }(3\stackrel{~}{\mathrm{\Lambda }}),`$
with
$$F_{cl}(x,y)=\frac{\mathrm{sin}[\mathrm{\Delta }_{xy}]}{x}+\frac{y\mathrm{cos}[\mathrm{\Delta }_{xy}]}{x},$$
(74)
where to reduce the notation we have defined the following set of dimensionless variables:
$`\stackrel{~}{k_0}={\displaystyle \frac{k_0}{H}},\stackrel{~}{\mathrm{\Lambda }}={\displaystyle \frac{\mathrm{\Lambda }}{k_0}},y=|k_0\eta ^{}|,x=|k_0\eta |,\mathrm{\Delta }_{xy}=xy.`$
The integrations above can be performed exactly and the result can be numerically integrated over the conformal time $`\eta `$.
To obtain $`D_2`$ from Eq.(III) we need to compute the Fourier transform of the real part of $`G_{++}^{\mathrm{\Lambda }2}`$ in $`\stackrel{}{k}=2\stackrel{}{k_0}`$ and $`\stackrel{}{k}=\stackrel{}{0}`$. This real part is given by
$`ReG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},2\stackrel{}{k_0})=(2\pi )^3{\displaystyle _{k>\mathrm{\Lambda }}}d^3k{\displaystyle _{k^{}>\mathrm{\Lambda }}}d^3k^{}`$
$`\times \delta ^3\left(\stackrel{}{k}+\stackrel{}{k}^{}2\stackrel{}{k_0}\right)Re\{\stackrel{~}{\varphi }_k(\eta )\stackrel{~}{\varphi }_k^{}(\eta ^{})\stackrel{~}{\varphi }_k^{}(\eta )\stackrel{~}{\varphi }_k^{}^{}(\eta ^{})\},`$ (75)
where $`\stackrel{~}{\varphi }_k`$ is the time-dependent mode function defined in Eq.(71) and $`\stackrel{~}{\varphi }_k^{}`$ its complex conjugate. This Fourier transform can be expressed in terms of integrals of the form,
$`I_{n,m}^C`$ $`=`$ $`{\displaystyle \frac{k_0^{m+n3}}{\pi }}{\displaystyle _{k>\mathrm{\Lambda }}}d^3k{\displaystyle _{k^{}>\mathrm{\Lambda }}}d^3k^{}\delta ^3\left(\stackrel{}{k}+\stackrel{}{k}^{}2\stackrel{}{k_0}\right)`$ (76)
$`\times `$ $`{\displaystyle \frac{\mathrm{cos}[(k+k^{})(\eta \eta ^{})]}{k^nk_{}^{}{}_{}{}^{m}}}`$
$`=`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}+2}^+\mathrm{}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle _{u2}^{u+2}}{\displaystyle \frac{dz}{z^{m1}}}\mathrm{cos}[(u+z)\mathrm{\Delta }_{xy}]`$
$`+`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^{\stackrel{~}{\mathrm{\Lambda }}+2}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^{u+2}}{\displaystyle \frac{dz}{z^{m1}}}\mathrm{cos}[(u+z)\mathrm{\Delta }_{xy}],`$
where $`n`$ and $`m`$ are integer numbers (only $`m=3`$ result to be necessary). The second equality follows after the change of variables $`u=k_0^1k`$ and $`z=k_0^1|\stackrel{}{k}2\stackrel{}{k_0}|=k_0^1\sqrt{k^24kk_0\mathrm{cos}(\theta )+4k_0^2}`$, where $`\theta `$ is the angle between $`\stackrel{}{k}`$ and $`\stackrel{}{k_0}`$. We also define the integral $`I_{n,m}^S`$ as the one obtained from $`I_{n,m}^C`$ by replacing the cosine function by sine.
An integration of $`I_{n,m}^{C,S}`$ by parts yields
$`I_{n,m}^C+\mathrm{\Delta }_{xy}{\displaystyle \frac{I_{n,m1}^S}{2m}}`$ $`=`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^+\mathrm{}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{cos}[2(u+1)\mathrm{\Delta }_{xy}]}{(2m)(u+2)^{m2}}}`$ (77)
$``$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}+2}^+\mathrm{}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{cos}[2(u1)\mathrm{\Delta }_{xy}]}{(2m)(\stackrel{~}{\mathrm{\Lambda }}+u)^{m2}}}`$
$``$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^{\stackrel{~}{\mathrm{\Lambda }}+2}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{cos}[2(u+1)\mathrm{\Delta }_{xy}]}{(2m)\stackrel{~}{\mathrm{\Lambda }}^{m2}}},`$
$`I_{n,m}^S\mathrm{\Delta }_{xy}{\displaystyle \frac{I_{n,m1}^C}{2m}}`$ $`=`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^+\mathrm{}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{sin}[2(u+1)\mathrm{\Delta }_{xy}]}{(2m)(u+2)^{m2}}}`$ (78)
$`+`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}+2}^+\mathrm{}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{sin}[2(u1)\mathrm{\Delta }_{xy}]}{(2m)(\stackrel{~}{\mathrm{\Lambda }}+u)^{m2}}}`$
$`+`$ $`{\displaystyle _{\stackrel{~}{\mathrm{\Lambda }}}^{\stackrel{~}{\mathrm{\Lambda }}+2}}{\displaystyle \frac{du}{u^{n1}}}{\displaystyle \frac{\mathrm{sin}[2(u+1)\mathrm{\Delta }_{xy}]}{(2m)\stackrel{~}{\mathrm{\Lambda }}^{m2}}}.`$
With the use of these properties and definitions, the Fourier transform (A) becomes
$`ReG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},2\stackrel{}{k_0})={\displaystyle \frac{H^4x^2y^2}{8(2\pi )^2k_0^3}}\left\{I_A+{\displaystyle \frac{2I_B}{xy}}+{\displaystyle \frac{I_C}{x^2y^2}}\right\}`$ (79)
where we have defined:
$`I_A=`$ $`I_{11}^C`$ (80a)
$`I_B=`$ $`I_{31}^C\mathrm{\Delta }_{xy}I_{21}^S`$ (80b)
$`I_C=`$ $`I_{33}^C2\mathrm{\Delta }_{xy}I_{32}^S\mathrm{\Delta }_{xy}^2I_{22}^C`$ (80c)
These last integrals are easily computed, with the result
$`I_A`$ $`=`$ $`2\pi \delta \left(\mathrm{\Delta }_{xy}\right)+{\displaystyle \frac{\mathrm{cos}[2(\stackrel{~}{\mathrm{\Lambda }}+1)\mathrm{\Delta }_{xy}]\mathrm{cos}[2\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}]}{\mathrm{\Delta }_{xy}^2}},`$ (81)
$`I_B`$ $`=`$ $`\mathrm{cos}[2\mathrm{\Delta }_{xy}]\left\{𝒞i[2(\stackrel{~}{\mathrm{\Lambda }}+2)|\mathrm{\Delta }_{xy}|]𝒞i[2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{\Delta }_{xy}|]\right\}`$ (82)
$``$ $`\mathrm{sin}[2|\mathrm{\Delta }_{xy}|]\left\{\pi 𝒮i[2(\stackrel{~}{\mathrm{\Lambda }}+2)|\mathrm{\Delta }_{xy}|]𝒮i[2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{\Delta }_{xy}|]\right\}`$
$`+`$ $`{\displaystyle \frac{\mathrm{sin}[2(\stackrel{~}{\mathrm{\Lambda }}+1)\mathrm{\Delta }_{xy}]}{\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}}}{\displaystyle \frac{\mathrm{sin}[2\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}]}{\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}}},`$
$`I_C`$ $`=`$ $`\left\{\pi 𝒮i[2(\stackrel{~}{\mathrm{\Lambda }}+2)|\mathrm{\Delta }_{xy}|]𝒮i[2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{\Delta }_{xy}|]\right\}`$ (83)
$`\times `$ $`\left(\mathrm{cos}[2\mathrm{\Delta }_{xy}]|\mathrm{\Delta }_{xy}|{\displaystyle \frac{\mathrm{sin}[2|\mathrm{\Delta }_{xy}|]}{2}}\right)`$
$`+`$ $`\left\{𝒞i[2(\stackrel{~}{\mathrm{\Lambda }}+2)|\mathrm{\Delta }_{xy}|]𝒞i[2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{\Delta }_{xy}|]\right\}`$
$`\times `$ $`\left(\mathrm{sin}[2\mathrm{\Delta }_{xy}]\mathrm{\Delta }_{xy}+{\displaystyle \frac{\mathrm{cos}[2\mathrm{\Delta }_{xy}]}{2}}\right)`$
$`+`$ $`{\displaystyle \frac{\mathrm{cos}[2\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}]}{\stackrel{~}{\mathrm{\Lambda }}^2}}{\displaystyle \frac{\mathrm{cos}[2(\stackrel{~}{\mathrm{\Lambda }}+1)\mathrm{\Delta }_{xy}]}{\stackrel{~}{\mathrm{\Lambda }}}},`$
where $`𝒮i[x]`$ and $`𝒞i[x]`$ are the sine and the cosine integral, respectively Abramob .
The Fourier transform in $`\stackrel{}{k}=\stackrel{}{0}`$ is much simpler to compute, and the result is
$`ReG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},\stackrel{}{0})={\displaystyle \frac{H^4x^2y^2}{2(2\pi )^2k_0^3}}\{{\displaystyle \frac{\pi }{2}}\delta \left(\mathrm{\Delta }_{xy}\right)`$
$`\mathrm{sin}[2\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}]\left({\displaystyle \frac{1}{2\mathrm{\Delta }_{xy}}}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{\mathrm{\Delta }_{xy}}{\stackrel{~}{\mathrm{\Lambda }}^2x^2y^2}}\right)`$
$`+{\displaystyle \frac{\mathrm{cos}[2\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Delta }_{xy}]}{\stackrel{~}{\mathrm{\Lambda }}}}\left({\displaystyle \frac{2}{xy}}+{\displaystyle \frac{\mathrm{\Delta }_{xy}^2}{3x^2y^2}}+{\displaystyle \frac{1}{3\stackrel{~}{\mathrm{\Lambda }}^2x^2y^2}}\right)`$
$`(\pi 2𝒮i[2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{\Delta }_{xy}|]){\displaystyle \frac{|\mathrm{\Delta }_{xy}|}{xy}}({\displaystyle \frac{\mathrm{\Delta }_{xy}^2}{3xy}}+1)\}.`$ (84)
Substitution of Eqs.(79) and (A) into Eq.(III) yields an expression for $`D_2`$ that allows to evaluate it and its integration over the conformal time $`\eta `$ numerically.
## Appendix B Renormalization
In this Appendix we obtain the renormalized semiclassical equation for the system-field given in Eq.(55). To do so, we follow the same procedure as for the renormalization of the evolution equation for the mean value $`\widehat{\varphi }`$ RenoPazMazziCarmen .
Starting with the bare action for the field $`\varphi `$, the semiclassical equation for the homogeneous system-field $`\varphi _<`$ reads
$`\varphi _<^{\prime \prime }(\eta )+[m_0^2+\xi _0R]a^2(\eta )\varphi _<(\eta )+2\varphi _<^{}(\eta )`$
$`+4\lambda _0a^2(\eta )\varphi _<^3(\eta )12\lambda a^2(\eta )\varphi _<(\eta )iG_{++}^\mathrm{\Lambda }(\eta ,\eta ,\stackrel{}{0})`$
$`288\lambda ^2a^2(\eta )\varphi _<(\eta ){\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta ^{}){\displaystyle d^3y\varphi _<^2(\eta ^{})}`$
$`\times ImG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},\stackrel{}{x}\stackrel{}{y})=\xi _2(\eta ,\stackrel{}{x}){\displaystyle \frac{\varphi _<(\eta )}{a^2(\eta )}},`$ (85)
where $`m_0`$, $`\xi _0`$ and $`\lambda _0`$ are the bare constants, and $`R=12H^2`$. Here we have replaced $`\lambda _0`$ by the renormalized constant $`\lambda `$ in the one-loop terms.
Let us add and substrate the following terms to the left-hand side of Eq.(B):
$$12\lambda a^2(\eta )\varphi _<(\eta )\widehat{\varphi }_>^2_{ad2}$$
(86)
where $`\widehat{\varphi }_>^2_{ad2}`$ is the adiabatic expansion of the expectation value $`\widehat{\varphi }_>^2`$ up to the second adiabatic order RenoPazMazziCarmen .
In what follows we will show how the infinities in Eq.(B) are cancelled by the terms subtracted. We will compute the divergent terms added via dimensional regularization, which will thus be able to be absorbed in the bare parameters as usual.
The expansion of $`\widehat{\varphi }_>^2`$ up to the second adiabatic order yields:
$$\varphi _>^2_{ad2}\varphi _>^2_{ad2}^F+\varphi _>^2_{ad2}^I,$$
(87)
with
$`\varphi _>^2_{ad2}^F=`$ $`{\displaystyle \frac{a}{2}}{\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{d^3k}{(2\pi a)^3}}\{{\displaystyle \frac{1}{\omega _k}}{\displaystyle \frac{\left(\xi _0\frac{1}{6}\right)a^2R}{2\omega _k^3}}`$ (88a)
$`+`$ $`{\displaystyle \frac{m_0^2}{4\omega _k^5}}[a_{}^{}{}_{}{}^{2}+aa^{\prime \prime }]{\displaystyle \frac{5m_0^4}{8\omega _k^7}}(aa^{})^2\},`$
$`\varphi _>^2_{ad2}^I=`$ $`{\displaystyle \frac{a}{2}}{\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{d^3k}{(2\pi a)^3}}{\displaystyle \frac{6\lambda _0a^2}{\omega _k^3}}\varphi _<^2.`$ (88b)
where $`\omega _k^2=k^2+a^2m_0^2`$, and the integrations over the wave vector $`\stackrel{}{k}`$ are restricted by $`k>\mathrm{\Lambda }`$.
In order to use dimensional regularization, we perform the integrations above over all wave vectors with $`k0`$ and then we subtract the ones restricted by $`k<\mathrm{\Lambda }`$ with the result
$`\varphi _>^2_{ad2}^F`$ $`={\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \frac{R}{12}}\left[{\displaystyle \frac{1}{n4}}+{\displaystyle \frac{\gamma }{2}}\mathrm{ln}\left|{\displaystyle \frac{a(\eta )\mu }{2\mathrm{\Lambda }}}\right|\right]`$
$`{\displaystyle \frac{\mathrm{\Lambda }^2}{8\pi ^2a^2(\eta )}},`$ (89a)
$`\varphi _>^2_{ad2}^I`$ $`={\displaystyle \frac{3\lambda }{2\pi ^2}}\varphi _<^2(\eta )\left[{\displaystyle \frac{1}{n4}}+{\displaystyle \frac{\gamma }{2}}\mathrm{ln}\left|{\displaystyle \frac{a(\eta )\mu }{2\mathrm{\Lambda }}}\right|\right],`$ (89b)
where $`\gamma `$ is the Euler’s constant and we have replaced the bare parameters with their renormalized counterparts ($`m=0`$, $`\xi =0`$ and $`\lambda `$). It is important to notice that only after subtracting the bare parameters can be replaced by the renormalized ones, since one of the integrals over $`k<\mathrm{\Lambda }`$ has an infrared logarithmic divergence.
Writing the bare parameters in terms of the renormalized ones plus conterterms,
$$m_0^2=0+\delta m^2,\xi _0=0+\delta \xi ,\lambda _0=\lambda +\delta \lambda ,$$
we can see that the divergences appearing for $`n4`$ are cancelled with the use of the following counterterms:
$`\delta m^2=`$ $`\mathrm{\Delta }m^2,`$ (90a)
$`\delta \xi =`$ $`{\displaystyle \frac{\lambda }{4\pi ^2}}\left[{\displaystyle \frac{1}{n4}}\right]+\mathrm{\Delta }\xi ,`$ (90b)
$`\delta \lambda =`$ $`{\displaystyle \frac{9\lambda ^2}{2\pi ^2}}\left[{\displaystyle \frac{1}{n4}}\right]+\mathrm{\Delta }\lambda ,`$ (90c)
where $`\mathrm{\Delta }m^2`$, $`\mathrm{\Delta }\xi `$ and $`\mathrm{\Delta }\lambda `$ remain finite as $`n4`$.
Replacing the bare parameters by the renormalized ones in the integrals over $`k>\mathrm{\Lambda }`$ of Eq.(88) we obtain
$`\varphi _>^2_{ad2}^F`$ $`=iG_{++}^\mathrm{\Lambda }(\eta ,\eta ,\stackrel{}{0}),`$ (91a)
$`\varphi _>^2_{ad2}^I`$ $`={\displaystyle \frac{3\lambda }{2\pi ^2}}{\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{dk}{k}}\varphi _<^2(\eta ).`$ (91b)
From equation above it is simple to note that the order $`\lambda `$ contribution in Eq.(B) is completely cancelled by the one with $`\varphi _>^2_{ad2}^F`$ (see Eq.(86)).
In order to separate the divergent part from the order $`\lambda ^2`$ contribution of Eq.(B), we write the propagator as
$`{\displaystyle d^3yImG_{++}^{\mathrm{\Lambda }2}(\eta ,\eta ^{},\stackrel{}{x},\stackrel{}{y})}=(2\pi )^3{\displaystyle _{k>\mathrm{\Lambda }}}d^3kIm[\stackrel{~}{\varphi }_\stackrel{}{k}^2\stackrel{~}{\varphi }_\stackrel{}{k}^2]`$
$`={\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{dk}{8\pi ^2}}\{{\displaystyle \frac{\mathrm{cos}[2k(\eta \eta ^{})]}{a^2(\eta )a^2(\eta ^{})}}[{\displaystyle \frac{2(\eta \eta ^{})}{k\eta \eta ^{}}}+{\displaystyle \frac{2(\eta \eta ^{})}{k^3\eta ^2\eta _{}^{}{}_{}{}^{2}}}]`$
$`{\displaystyle \frac{\mathrm{sin}[2k(\eta \eta ^{})]}{a^2(\eta )a^2(\eta ^{})}}[1+{\displaystyle \frac{1}{k^4\eta ^2\eta _{}^{}{}_{}{}^{2}}}+{\displaystyle \frac{2}{k^2\eta \eta ^{}}}{\displaystyle \frac{(\eta \eta ^{})^2}{k^2\eta ^2\eta _{}^{}{}_{}{}^{2}}}]\}`$
$`I_D+I_{ND},`$
where $`\stackrel{~}{\varphi }_\stackrel{}{k}`$ and $`\stackrel{~}{\varphi }_\stackrel{}{k}^{}`$ are the mode function and its complex conjugate respectively, given in Eq.(71). The only divergent term is
$$I_D_{k>\mathrm{\Lambda }}\frac{dk}{8\pi ^2}\frac{\mathrm{sin}[2k(\eta \eta ^{})]}{a^2(\eta )a^2(\eta ^{})}.$$
(92)
With the use of the definition of $`I_D`$ and $`I_{ND}`$, we can write the order $`\lambda ^2`$ contribution as
$``$ $`288\lambda ^2a^2(\eta ){\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta ^{})\varphi _<(\eta ){\displaystyle d^3y\varphi _<^2(\eta ^{})ImG_{++}^{\mathrm{\Lambda }2}}`$ (93)
$``$ $`12\lambda a^2(\eta )\varphi _<(\eta )\left([\varphi _>^2]_D^I+[\varphi _>^2]_{ND}^I\right),`$
with
$`[\varphi _>^2]_D^I=24\lambda {\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}a^4(\eta ^{})\varphi _<^2(\eta ^{})I_D`$
$`={\displaystyle \frac{3\lambda }{\pi ^2}}{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}{\displaystyle \frac{a^2(\eta ^{})\varphi _<^2(\eta ^{})}{a^2(\eta )}}{\displaystyle _{k>\mathrm{\Lambda }}}𝑑k\mathrm{sin}[2k(\eta \eta ^{})]`$
$`={\displaystyle \frac{3\lambda }{2\pi ^2}}\varphi _<^2(\eta ){\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{dk}{k}}+{\displaystyle \frac{3\lambda }{2\pi ^2}}{\displaystyle \frac{a^2(\eta _i)\varphi _<^2(\eta _i)}{a^2(\eta )}}`$
$`\times {\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{dk}{k}}\mathrm{cos}[2k(\eta \eta _i)]`$
$`+{\displaystyle \frac{3\lambda }{2\pi ^2}}{\displaystyle _{\eta _i}^\eta }𝑑\eta ^{}{\displaystyle \frac{(a^2(\eta ^{})\varphi _<^2(\eta ^{}))^{}}{a^2(\eta )}}{\displaystyle _{k>\mathrm{\Lambda }}}{\displaystyle \frac{dk}{k}}\mathrm{cos}[2k(\eta \eta ^{})],`$
where the last equality follows after performing an integration by parts. In this expression we can see that the first term after the last equality is $`\varphi _>^2_{ad2}^I`$ and hence all infinities are cancelled. Finally, performing the $`I_{ND}`$ integral and reordering the terms we obtain the explicit form of the renormalized equation given in Eq.(55).
|
warning/0506/quant-ph0506021.html
|
ar5iv
|
text
|
# Condition and capability of quantum state separation
## I Introduction
A fundamental difference between quantum mechanics and the classical correspondence is that in the former, a system can be not only in a basis state but also in a state which is a linear combination, or ‘superposition’, of different basis states. Quantum computation and quantum information processing benefit extremely from superposition since performing a quantum operation on a superposition is equivalent to performing the same operation synchronously on all of the basis states constituting this superposition. One of the most famous examples is Shor’s quantum factoring algorithm SH94 . On the other hand, however, the existence of superposition in quantum mechanics also puts many constraints on the physically realizable information processing tasks, when we have only limited information about the original state of the system that we are concerned with. Take quantum cloning, perhaps the most fundamental task in quantum computation and quantum information processing, as an example. When the state to be cloned is thoroughly known, it can be perfectly cloned by using a state-dependent cloning machine (In fact, since the state is known, we can prepare as many copies of it as needed. The reason behind it is in fact that classical information can be cloned arbitrarily). Here and in the rest of this paper, by ‘perfectly’ we mean the information processing task is realized with certainty and without any approximation or error. Suppose further we want to build a universal cloning machine for different pure states, then only if these states are linearly independent that a desired exact cloning machine exists even in a probabilistic manner DG98 . The possibility to obliviously clone states from a linearly dependent set is forbidden by the linearity of quantum operations. Another result in Ref. DG98 which receives less attention than it deserves is the converse of the above statement. That is, when the possible states of the original system are linearly independent, then $`1N`$ probabilistic cloning is possible for any $`N1`$. In this paper, we generalize this result to show that the linear independency of the original states is enough to make any information processing tasks possible in a probabilistic manner.
Another fundamental task in quantum computation and quantum information processing is quantum discrimination. Given that the system of interest is prepared in one of some possible states, the purpose of discrimination is to tell which state the system is actually in. Rather surprisingly, these two seemingly, at least at first glance, very different tasks are closely related. A quantum system can be perfectly cloned WZ82 ; DI82 (resp. perfectly discriminated HE76 ) if and only if the possible states of the system are orthogonal; and it can be conclusively cloned DG98 (resp. unambiguously discriminated DG98 ; CH98 ) if and only if the possible states are linearly independent. Furthermore, Duan and Guo DG98 pointed out that exact $`1\mathrm{}`$ cloning and unambiguous discrimination can be simulated by each other; a more delicate and quantitative connection between these two tasks was investigated in Ref. CB98 .
Motivated by this connection, Chefles and Barnett proposed a generalized way, namely quantum separation, to deal with quantum exact cloning and quantum unambiguous discrimination uniformly CB98 . To be specific, suppose a quantum system is prepared in one of the two states $`|\psi _1`$ and $`|\psi _2`$ but we do not know exactly which one. A quantum separation performed on this system then leads, generally in a probabilistic but conclusive manner, the system into $`|\psi _i^{}`$ provided that originally it is in the state $`|\psi _i`$, for $`i=1,2`$. In their paper, Chefles and Barnett put a constraint that the desired states $`|\psi _1^{}`$ and $`|\psi _2^{}`$ should satisfy the condition that
$$|\psi _1^{}|\psi _2^{}||\psi _1|\psi _2|,$$
(1)
just as in the cases of exact cloning and unambiguous discrimination. That is also why they called this process ‘separation’ since decrease of the inner product means that these two states become more distinct or separable. In the present paper, we generalize this concept in two ways. First, we get rid of the constraint in Eq.(1) to consider more general physical processes, although we still use the term ‘separation’ for convenience. Second, we generalize separation to the case of multiple mixed states.
To be specific, we give the formal definition of quantum separation as follows. Suppose a quantum system is prepared in a state secretly chosen from $`\rho _1,\mathrm{},\rho _n`$. A quantum separation is a physically realizable process which, generally in a probabilistic but conclusive manner, leads $`\rho _i`$ to $`\rho _i^{}`$ for some quantum states $`\rho _1^{},\mathrm{},\rho _n^{}`$. Recall that any physically realizable process is completely positive and trace preserving, and so can be represented by Kraus operator-sum form KR83 . That is, there exist quantum operators $`A_{Sk},A_{Fk}`$ such that
$$A_{Sk}\rho _iA_{Sk}^{}=s_{ik}\rho _i^{},$$
(2)
$$A_{Fk}\rho _iA_{Fk}^{}=f_{ik}\sigma _{ik}$$
(3)
for some nonnegative real numbers $`s_{ik}`$ and $`f_{ik}`$, and mixed states $`\sigma _{ik}`$, where $`i=1,\mathrm{},n`$. Here the subscript $`S`$ and $`F`$ denote success and failure, respectively. Intuitively, Eq.(2) means that if the separation succeeds, the system evolves into $`\rho _i^{}`$ provided that it is originally in the state $`\rho _i`$. Notice that there may be more than one operator, indexed by $`k`$, corresponding to successfully separating $`\rho _i`$ or getting an inconclusive result. By appending the shorter group with zero operators, we can assume that the range of $`k`$ is taken the same for success and failure. Furthermore, these operators should satisfy the completeness relation
$$\underset{k}{}(A_{Sk}^{}A_{Sk}+A_{Fk}^{}A_{Fk})=I.$$
(4)
Here $`I`$ is the identity operator.
Since no constraints are put on the output states in the general framework, we can in fact represent any oblivious computation and information process by quantum state separation. To see the power of this framework more explicitly, let us examine some special cases. It is easy to check that exact $`1N`$ cloning is a special case of quantum separation by letting the desired state $`\rho _i^{}`$ be $`\rho _i^N`$ while unambiguous discrimination is the case when all $`\rho _i^{}`$ are orthogonal such that there exists a quantum measurement which can further discriminate them perfectly. Furthermore, suppose all $`\rho _i`$ lie in a Hilbert space $``$. The $`1N`$ mixed state broadcasting BC96 can be involved in the general framework of quantum state separation by requiring that each $`\rho _i^{}`$ lies in the Hilbert space $`^N`$ and the reduced density matrices of $`\rho _i^{}`$ obtained by tracing over any $`N1`$ subsystems equal to $`\rho _i`$. Note also that unambiguous filtering TB03 , unambiguous comparison BC03 , and unambiguous subset discrimination ZY02 are all special cases of unambiguous discrimination between mixed states, which has received much attention in recent years \[13-19\]. By considering quantum state separation, we can deal with all these information processing processes in a uniform and more general way.
The aim of this paper is to examine the conditions and the capability of quantum information processing in the framework of state separation. In Sec. II, we show that in order to physically realize a universal and conclusive information processing task on an unknown system, linearity is in fact the only constraint. In other words, when the possible states of the unknown system are linearly independent, then any separation with any output states is possible. In Sec. III, we derive a lower bound on the average failure probability of any physically realizable quantum separation, when the mixed state case is considered.
## II Conditions of state separation
In this section, we derive some necessary and sufficient conditions for quantum separation to be physically realizable. First, when the final states are specified, we have the following theorem for the pure state case.
###### Theorem 1
Given two sets of pure states $`|\psi _1,\mathrm{},|\psi _n`$ and $`|\psi _1^{},\mathrm{},|\psi _n^{}`$. There exists a quantum separation which can lead $`|\psi _i`$ to $`|\psi _i^{}`$ if and only if
$$X\sqrt{\mathrm{\Gamma }}X^{}\sqrt{\mathrm{\Gamma }}0$$
(5)
for some positive definite diagonal matrix $`\mathrm{\Gamma }=\mathrm{diag}(\gamma _1,\mathrm{},\gamma _n)`$, where $`n\times n`$ matrices $`X=[\psi _i|\psi _j]`$ and $`X^{}=[\psi _i^{}|\psi _j^{}]`$. Here by $`M0`$ we mean that the matrix $`M`$ is positive semidefinite, i.e., for any $`n`$-dimensional complex vector $`\alpha `$, $`\alpha M\alpha ^{}0`$.
To prove this theorem, we introduce first a lemma proven in Ref. DG98 :
###### Lemma 1
For any two sets of pure states $`|\psi _1,\mathrm{},|\psi _n`$ and $`|\psi _1^{},\mathrm{},|\psi _n^{}`$, if
$$\psi _i|\psi _j=\psi _i^{}|\psi _j^{}$$
(6)
for any $`i,j=1,\mathrm{},n`$, then there exits a unitary operator $`U`$ such that $`U|\psi _i=|\psi _i^{}`$.
We learn from this lemma that in pure state case, the only thing determining whether or not there exists a unitary evolution between two sets of states is the inner products of all pairs of states from the same set. This is a remarkable property of pure state evolution. When mixed states are considered, things become more complicated and many more facts other than fidelities between different states must be involved to determine the existence of such a unitary transformation. That is also why we consider only pure state case here.
Having the above lemma as a tool, we can prove Theorem 1 as follows:
Proof of Theorem 1. By definition, there exist quantum operators $`A_{Sk}`$ and $`A_{Fk}`$ satisfying Eq.(4) such that
$$A_{Sk}|\psi _i=\sqrt{s_{ik}}|\psi _i^{}$$
(7)
$$A_{Fk}|\psi _i=\sqrt{f_{ik}}|\varphi _{ik},$$
(8)
for some state $`|\varphi _{ik}`$, where $`0<s_{ik}1`$ and $`0f_{ik}<1`$. For any $`n`$-dimensional complex vector $`\alpha =(\alpha _1,\mathrm{},\alpha _n)`$, let $`|\mathrm{\Psi }=_{i=1}^n\alpha _i|\psi _i`$. Notice that $`A_{Fk}^{}A_{Fk}`$ is positive semidefinite for any $`k`$. It follows that
$$\begin{array}{cc}\hfill 0& \mathrm{\Psi }|\underset{k}{}A_{Fk}^{}A_{Fk}|\mathrm{\Psi }\hfill \\ & \\ & =\mathrm{\Psi }|I\underset{k}{}A_{Sk}^{}A_{Sk}|\mathrm{\Psi }\hfill \\ & \\ & =\mathrm{\Psi }|\mathrm{\Psi }\underset{k}{}\mathrm{\Psi }|A_{Sk}^{}A_{Sk}|\mathrm{\Psi }\hfill \\ & \\ & =\underset{i,j}{}\alpha _i^{}\alpha _j\psi _i|\psi _j\underset{k}{}\underset{i,j}{}\alpha _i^{}\alpha _j\sqrt{s_{ik}s_{jk}}\psi _i^{}|\psi _j^{}\hfill \\ & \\ & =\alpha X\alpha ^{}\alpha \underset{k}{}\sqrt{S_k}X^{}\sqrt{S_k}\alpha ^{}\hfill \\ & \\ & \alpha X\alpha ^{}\alpha \sqrt{S_1}X^{}\sqrt{S_1}\alpha ^{}.\hfill \end{array}$$
(9)
Here, $`S_k=\mathrm{diag}(s_{1k},\mathrm{},s_{nk})`$ are $`n\times n`$ diagonal matrices. The last line of Eq.(9) follows from the fact that for any $`k`$, $`\sqrt{S_k}X^{}\sqrt{S_k}`$ is positive semidefinite. From the arbitrariness of $`\alpha `$, we derive that
$$X\sqrt{S_1}X^{}\sqrt{S_1}0,$$
(10)
which completes the proof of the necessity part.
The proof of the sufficiency part is almost the same as the proof of that linear independency implies capability of exact cloning in Ref. DG98 . To be complete, we outline here the main steps.
To show the existence of a desired separation under the assumption of Eq.(5), we need only to prove that there exists a unitary transformation $`U`$ such that for any $`i=1,\mathrm{},n`$,
$$U|\psi _i_A|\mathrm{\Sigma }_B|P_P=\sqrt{\gamma _i}|\psi _i^{}_{AB}|P_0_P+\underset{k=1}{\overset{n}{}}c_{ik}|\mathrm{\Phi }_i_{AB}|P_k_P,$$
(11)
where $`|P_0,|P_1,\mathrm{},|P_n`$ are orthonormal states in the probe system $`P`$, and $`|\mathrm{\Phi }_i_{AB}`$ are normalized but not necessarily orthogonal states. Here the subscript $`B`$ denotes an ancillary system and $`|\mathrm{\Sigma }`$ is a standard ‘blank’ state (in some cases, say unambiguous discrimination, the ancillary system is unnecessary). After the unitary evolution described by Eq.(11), a projective measurement which consists of $`|P_0P_0|`$ and $`I|P_0P_0|`$ is performed on probe system $`P`$. If the outcome corresponding to $`I|P_0P_0|`$ occurs, the separation fails; otherwise this separation succeeds and the secretly chosen state $`|\psi _i`$ conclusively evolves into the desired state $`|\psi _i^{}`$.
In the following, we show the existence of the unitary transformation $`U`$ in Eq.(11). Taking the inter-inner products of the both sides of Eq.(11) for different $`i`$ and $`j`$, we have the matrix equation
$$X=\sqrt{\mathrm{\Gamma }}X^{}\sqrt{\mathrm{\Gamma }}+CC^{},$$
(12)
where $`n\times n`$ matrix $`C=[c_{ij}]`$. From Lemma 1, the only thing left is to show the existence of the matrix $`C`$. But from Eq.(5), the positive semidefinite matrix $`X\sqrt{\mathrm{\Gamma }}X^{}\sqrt{\mathrm{\Gamma }}`$ can be diagonalized by a unitary matrix $`V`$ as
$$V(X\sqrt{\mathrm{\Gamma }}X^{}\sqrt{\mathrm{\Gamma }})V^{}=\mathrm{diag}(c_1,\mathrm{},c_n)$$
(13)
for some nonnegative numbers $`c_1,\mathrm{},c_n`$. So we need only set $`C=V^{}\mathrm{diag}(\sqrt{c_1},\mathrm{},\sqrt{c_n})V^{}`$ and then the sufficiency part of the theorem is proven. $`\mathrm{}`$
Theorem 1 tells us when a $`given`$ separation can be physically realized in pure state case. The following theorem, however, gives a necessary and sufficient condition under which $`any`$ quantum separation is realizable on a given system in the general case of mixed states. To begin with, we introduce some notations. For a density matrix $`\rho `$, we denote by $`\mathrm{supp}(\rho )`$ the support space of $`\rho `$. That is, the space spanned by all eigenvectors with nonzero corresponding eigenvalues of $`\rho `$. Furthermore, by $`\mathrm{supp}(\rho _1,\mathrm{},\rho _n)`$ we denote the support space spanned by eigenvectors of $`\rho _1,\mathrm{},\rho _n`$ with nonzero corresponding eigenvalues.
###### Theorem 2
Suppose a quantum system is prepared in a state secretly chosen from $`\rho _1,\mathrm{},\rho _n`$. Let $`S=\{\rho _1,\mathrm{},\rho _n\}`$ and $`S_i=S\backslash \{\rho _i\}`$. Then
1) any state separation on this system is possible (that is, for any states $`\rho _1^{},\mathrm{},\rho _n^{}`$, there exists a separation which leads $`\rho _i`$ conclusively to $`\rho _i^{}`$) if and only if $`\mathrm{supp}(S)\mathrm{supp}(S_i)`$ for any $`i=1,\mathrm{},n`$.
2) Furthermore, if $`\mathrm{supp}(S)=\mathrm{supp}(S_i)`$ for some $`i`$ and there exists a separation which leads $`\rho _i`$ conclusively to $`\rho _i^{}`$ for some quantum states $`\rho _1^{},\mathrm{},\rho _n^{}`$, then $`\mathrm{supp}(S^{})=\mathrm{supp}(S_i^{})`$, where $`S^{}=\{\rho _1^{},\mathrm{},\rho _n^{}\}`$ and $`S_i^{}=S^{}\backslash \{\rho _i^{}\}`$.
Proof. The necessity part of 1) is obvious, since we can take special cases of quantum separation, say unambiguous discrimination, to show that $`\mathrm{supp}(S)\mathrm{supp}(S_i)`$ (for the condition under which unambiguous discrimination between mixed states is possible, we refer to Ref. FD04 ).
To prove the sufficiency part of 1), suppose that $`\mathrm{supp}(S)\mathrm{supp}(S_i)`$ for any $`i=1,\mathrm{},n`$. Then from Ref. FD04 , there exist $`n`$ positive real numbers $`\gamma _1,\mathrm{},\gamma _n`$ such that we can unambiguously discriminate $`\rho _i`$ with probability $`\gamma _i`$. Once the state $`\rho _i`$ is identified, we can prepare $`\rho _i^{}`$ with certainty by a physical realizable process (which may be dependent on $`\rho _i^{}`$). So by combining these two steps together, we construct a protocol which leads $`\rho _i`$ to $`\rho _i^{}`$ with positive probability $`\gamma _i`$.
Now we prove 2) by contradiction. Suppose $`\mathrm{supp}(S^{})\mathrm{supp}(S_i^{})`$. Then there exists a pure state $`|\varphi `$ which is in $`\mathrm{supp}(\rho _i)`$ but not in $`\mathrm{supp}(S_i^{})`$. So we can construct a positive-operator valued measurement comprising $`|\varphi \varphi |`$ and $`I|\varphi \varphi |`$ to unambiguously discriminate $`\rho _i`$ from the other $`n1`$ states with a positive probability. Notice that an unambiguous discrimination is also a quantum separation. Combining these two separation processes together we get a new one which can discriminate unambiguously the state $`\rho _i`$ from other states with a positive probability. That is a contradiction with the assumption that $`\mathrm{supp}(S)=\mathrm{supp}(S_i)`$. $`\mathrm{}`$
Notice that when $`\rho _1=|\psi _1\psi _1|,\mathrm{},\rho _n=|\psi _n\psi _n|`$ are all pure states, the condition that $`\mathrm{supp}(S)\mathrm{supp}(S_i)`$ for any $`i=1,\mathrm{},n`$ is equivalent to that $`|\psi _1,\mathrm{},|\psi _n`$ are linearly independent. So we have the following corollary which has more physical intuition.
###### Corollary 1
Suppose a quantum system is prepared secretly in one of the states $`|\psi _1,\mathrm{},|\psi _n`$. Then
1) any state separation on this system is possible if and only if $`|\psi _1,\mathrm{},|\psi _n`$ are linearly independent.
2) Furthermore, if $`|\psi _1,\mathrm{},|\psi _n`$ are linearly dependent and there exists a separation which leads $`|\psi _i`$ conclusively to $`|\psi _i^{}`$ for some quantum states $`|\psi _1^{},\mathrm{},|\psi _n^{}`$, then $`|\psi _1^{},\mathrm{},|\psi _n^{}`$ are also linearly dependent.
The two statements in Corollary 1 are complementary with each other. Statement 2) tells us the constraints on realizable information processing tasks when the system we are concerned with is in a state coming secretly from a linearly dependent set. On the other hand, statement 1) shows that linear dependency is actually the only case in which physically realizable information processing tasks will be constrained. That is, if the state of the original system is prepared secretly in one of linearly independent pure states, then any tasks, represented by our generalized separation with arbitrary outcome states, are probabilistically and conclusively realizable.
From Theorem 1, we get the following direct corollary:
###### Corollary 2
For any set $`S=\{\rho _1,\mathrm{},\rho _n\}`$ of quantum states, the following statements are equivalent:
1) The states secretly chosen from $`S`$ can be unambiguously discriminated.
2) The states secretly chosen from $`S`$ can be conclusively cloned.
3) The set $`S`$ can evolve, through appropriate separation processes, into any set $`S^{}=\{\rho _1^{},\mathrm{},\rho _n^{}\}`$ of quantum states, where $`\rho _i`$ becomes $`\rho _i^{}`$ for any $`i=1,\mathrm{},n`$.
Informally, from this corollary, exact cloning and unambiguous discrimination put the strongest constraints on the possible states the original system can be prepared in.
## III Lower bound on average failure probability
Theorem 1 gives a necessary and sufficient condition under which a given separation can be realized for a given original system, when the case of $`pure`$ state is considered. The general case where the state of the system we are concerned with comes from a mixed state set is, however, not investigated. Actually, it is unlikely that there exists a corresponding condition for mixed states due to lack of a result similar to Lemma 1. However, we can still derive a lower bound on the average failure probability of any separation once it is realizable.
###### Theorem 3
Suppose a quantum system is prepared in a state secretly chosen from $`\rho _1,\mathrm{},\rho _n`$ with respective $`a`$ $`priori`$ probabilities $`\eta _1,\mathrm{},\eta _n`$, and there exists a separation which leads $`\rho _i`$ to $`\rho _i^{}`$ for some quantum states $`\rho _1^{},\mathrm{},\rho _n^{}`$. Then the average failure probability $`P_f`$ of this separation satisfies
$$P_f\sqrt{\frac{n}{n1}\underset{(i,j)\mathrm{\Delta }}{}\eta _i\eta _j\left(\frac{F(\rho _i,\rho _j)F(\rho _i^{},\rho _j^{})}{1F(\rho _i^{},\rho _j^{})}\right)^2},$$
(14)
where the index set $`\mathrm{\Delta }=\{(i,j):ij\mathrm{and}F(\rho _i^{},\rho _j^{})F(\rho _i,\rho _j)\}`$.
Proof. From the assumption, there exist quantum operators $`A_{Sk}`$ and $`A_{Fk}`$ satisfying the completeness relation Eq.(4), such that Eqs.(2) and (3) hold. It is easy to check that
$$P_f=\underset{i,k}{}\eta _if_{ik}$$
(15)
and for any $`i=1,\mathrm{},n`$,
$$\underset{k}{}(s_{ik}+f_{ik})=1.$$
(16)
By Cauchy-Schwarz inequality,
$$\begin{array}{cc}\hfill P_f^2& \frac{n}{n1}\underset{ij}{}\eta _i\eta _j\left(\underset{k}{}f_{ik}\right)\left(\underset{k}{}f_{jk}\right)\hfill \\ & \\ & \frac{n}{n1}\underset{ij}{}\eta _i\eta _j\left(\underset{k}{}\sqrt{f_{ik}f_{jk}}\right)^2.\hfill \end{array}$$
(17)
From Eq.(2) and Polar decomposition theorem, we have
$$A_{Sk}\sqrt{\rho _i}=\sqrt{A_{Sk}\rho _iA_{Sk}^{}}U_{ik}=\sqrt{s_{ik}}\sqrt{\rho _i^{}}U_{ik}$$
(18)
for some unitary matrix $`U_{ik}`$. And similarly, Eq.(3) implies that
$$A_{Fk}\sqrt{\rho _i}=\sqrt{A_{Fk}\rho _iA_{Fk}^{}}V_{ik}=\sqrt{f_{ik}}\sqrt{\sigma _{ik}}V_{ik}$$
(19)
for some unitary matrix $`V_{ik}`$.
Recall that for any density matrices $`\rho `$ and $`\sigma `$, the fidelity $`F(\rho ,\sigma )=\mathrm{max}_U|\mathrm{Tr}(\sqrt{\rho }\sqrt{\sigma }U)|`$ , where the maximum is taken over all unitary matrix $`U`$. For any $`ij`$, let us take $`U_i^j`$ such that $`F(\rho _i,\rho _j)=|\mathrm{Tr}(\sqrt{\rho _i}\sqrt{\rho _j}U_i^j)|`$. Then
$$\mathrm{Tr}(\sqrt{\rho _i}A_{Sk}^{}A_{Sk}\sqrt{\rho _j}U_i^j)=\sqrt{s_{ik}s_{jk}}\mathrm{Tr}(U_{ik}^{}\sqrt{\rho _i^{}}\sqrt{\rho _j^{}}U_{jk}U_i^j)$$
(20)
$$\mathrm{Tr}(\sqrt{\rho _i}A_{Fk}^{}A_{Fk}\sqrt{\rho _j}U_i^j)=\sqrt{f_{ik}f_{jk}}\mathrm{Tr}(V_{ik}^{}\sqrt{\sigma _{ik}}\sqrt{\sigma _{jk}}V_{jk}U_i^j).$$
(21)
Summing up Eqs.(20) and (21) for all $`k`$ and noticing Eq.(4), we have
$$\begin{array}{cc}\hfill F(\rho _i,\rho _j)=& |\underset{k}{}(\sqrt{s_{ik}s_{jk}}\mathrm{Tr}(\sqrt{\rho _i^{}}\sqrt{\rho _j^{}}W_{ijk})\hfill \\ & \\ & +\sqrt{f_{ik}f_{jk}}\mathrm{Tr}(\sqrt{\sigma _{ik}}\sqrt{\sigma _{jk}}W_{ijk}^{}))|,\hfill \end{array}$$
(22)
where $`W_{ijk}=U_{jk}U_i^jU_{ik}^{}`$ and $`W_{ijk}^{}=V_{jk}U_i^jV_{ik}^{}`$ are unitary matrices. We further derive that
$$\begin{array}{ccc}\hfill F(\rho _i,\rho _j)& \hfill & \underset{k}{}\sqrt{s_{ik}s_{jk}}|\mathrm{Tr}(\sqrt{\rho _i^{}}\sqrt{\rho _j^{}}W_{ijk})|\hfill \\ & & \\ & & +\underset{k}{}\sqrt{f_{ik}f_{jk}}|\mathrm{Tr}(\sqrt{\sigma _{ik}}\sqrt{\sigma _{jk}}W_{ijk}^{})|\hfill \\ & & \\ & \hfill & \underset{k}{}\sqrt{s_{ik}s_{jk}}F(\rho _i^{},\rho _j^{})\hfill \\ & & \\ & & +\underset{k}{}\sqrt{f_{ik}f_{jk}}F(\sigma _{ik},\sigma _{jk})\hfill \\ & & \\ & \hfill & \underset{k}{}\sqrt{s_{ik}s_{jk}}F(\rho _i^{},\rho _j^{})+\underset{k}{}\sqrt{f_{ik}f_{jk}}.\hfill \end{array}$$
(23)
Notice that
$$\begin{array}{cc}\hfill \underset{k}{}\sqrt{s_{ik}s_{jk}}& \underset{k}{}\frac{s_{ik}+s_{jk}}{2}\hfill \\ & \\ & =1\underset{k}{}\frac{f_{ik}+f_{jk}}{2}\hfill \\ & \\ & 1\underset{k}{}\sqrt{f_{ik}f_{jk}}.\hfill \end{array}$$
(24)
Substituting Eq.(24) into Eq.(23), we have
$$\underset{k}{}\sqrt{f_{ik}f_{jk}}\frac{F(\rho _i,\rho _j)F(\rho _i^{},\rho _j^{})}{1F(\rho _i^{},\rho _j^{})}$$
(25)
Taking Eq.(25) for $`(i,j)\mathrm{\Delta }`$ back into Eq.(17) and noticing that $`_k\sqrt{f_{ik}f_{jk}}0`$ for $`(i,j)\mathrm{\Delta }`$, we arrive at the desired bound,
$$P_f\sqrt{\frac{n}{n1}\underset{(i,j)\mathrm{\Delta }}{}\eta _i\eta _j\left(\frac{F(\rho _i,\rho _j)F(\rho _i^{},\rho _j^{})}{1F(\rho _i^{},\rho _j^{})}\right)^2}.$$
(26)
That completes the proof. $`\mathrm{}`$
Following the argument behind Theorem 3 in Ref. FD04 , we can derive a series of lower bounds on the average failure probability. For the sake of completeness, we outline the derivation as follows. Define
$$M_t=\underset{i}{}\eta _i^{2t}(\underset{k}{}f_{ik})^{2t}$$
(27)
and
$$N_t=\underset{ij}{}\eta _i^t\eta _j^t(\underset{k}{}f_{ik})^t(\underset{k}{}f_{jk})^t.$$
(28)
Then $`M_t=\sqrt{N_{2t}+M_{2t}}`$ and by Cauchy inequality, $`M_tN_t/(n1)`$. So for any $`r0`$,
$$\begin{array}{cc}P_f^2\hfill & =N_1+M_1=N_1+\sqrt{N_2+M_2}=\mathrm{}\hfill \\ & \\ & =N_1+\sqrt{N_2+\sqrt{\mathrm{}+\sqrt{N_{2^r}+M_{2^r}}}}\hfill \\ & \\ & N_1+\sqrt{N_2+\sqrt{\mathrm{}+\sqrt{\frac{n}{n1}N_{2^r}}}}.\hfill \end{array}$$
(29)
If we further define
$$C_t=\underset{(i,j)\mathrm{\Delta }}{}\eta _i^t\eta _j^t\left(\frac{F(\rho _i,\rho _j)F(\rho _i^{},\rho _j^{})}{1F(\rho _i^{},\rho _j^{})}\right)^{2t},$$
(30)
then from Eq.(25) and the fact that $`_k\sqrt{f_{ik}f_{jk}}0`$ for $`(i,j)\mathrm{\Delta }`$, we have $`N_tC_t`$. Consequently, the promised lower bounds on the average failure probability $`P_f`$ can be derived as
$$P_fP_f^{(r)}\sqrt{C_1+\sqrt{\mathrm{}+\sqrt{\frac{n}{n1}C_{2^r}}}}.$$
(31)
The bound presented in Eq.(14) is just the special case of the above bounds when $`r=0`$. Note that $`P_f^{(0)}P_f^{(1)}\mathrm{}`$ by Cauchy-Schwarz inequality. When $`r`$ increases, the bound becomes better and better; and the limit when $`r`$ tends to infinity is the best bound we can derive using this method.
Now let us analyze the bound in Eq.(14) carefully. First, note that when pure state separation is considered, Qiu obtained in Ref. QI02 a lower bound on the average failure probability which reads
$$1\frac{1}{n1}\underset{i<j}{}\frac{\eta _i+\eta _j2\sqrt{\eta _i\eta _j}|\psi _i|\psi _j|}{1|\psi _i^{}|\psi _j^{}|}.$$
(32)
It is easy by using Cauchy-Schwarz inequality to check that our bound presented in Eq.(14) is better in general than the one in Eq.(32). On the other hand, in the case of $`MN`$ ($`MN`$) exact cloning, where the original state and the final state are, respectively, $`\rho _i^M`$ and $`\rho _i^N`$ for $`i=1,\mathrm{},n`$, and so $`F(\rho _i^{},\rho _j^{})F(\rho _i,\rho _j)`$ holds for any $`ij`$. So we have actually derived a lower bound on the average failure probability of exact $`MN`$ cloning as
$$P_f^{EC}\sqrt{\frac{n}{n1}\underset{ij}{}\eta _i\eta _j\left(\frac{F(\rho _i,\rho _j)^MF(\rho _i,\rho _j)^N}{1F(\rho _i,\rho _j)^N}\right)^2}.$$
(33)
When $`\rho _i=|\psi _i\psi _i|`$ are pure states and $`\eta _1=\mathrm{}=\eta _n=1/n`$, this bound can be shown better than
$$1\frac{2}{n(n1)}\underset{i<j}{}\frac{1|\psi _i|\psi _j|^M}{1|\psi _i|\psi _j|^N},$$
(34)
which was derived in Ref. CB98 . Finally, in the case of unambiguous discrimination, where the final states $`\rho _i^{}`$ are orthogonal to each other, the bound in Eq.(14) turns out to be
$$P_f^{UD}\sqrt{\frac{n}{n1}\underset{ij}{}\eta _i\eta _jF(\rho _i,\rho _j)^2},$$
(35)
coinciding with that obtained in Ref. FD04 . It is also worth noting that the bound can further degenerate to the Jaeger-Shimony bound $`12\sqrt{\eta _1\eta _2}|\psi _1|\psi _2|`$ for two pure states JS95 and the IDP bound $`1|\psi _1|\psi _2|`$ for two pure states with equal $`a`$ $`priori`$ probabilities IV87 ; DI88 ; PE88 .
## IV Conclusion
To conclude, by deriving a necessary and sufficient condition for any quantum separation to be physically realizable, we show that in probabilistic manner, linearity is in fact the only one that restricts the physically realizable tasks. That is, when a system is prepared in a state secretly chosen from a linearly independent pure state set, then any generalized state separation is physically realizable with a positive probability. A lower bound on the average failure probability of any quantum state separation is also derived and special cases of this bound are analyzed.
The authors thank the colleagues in the Quantum Computation and Quantum Information Research Group for useful discussion. This work was partly supported by the Natural Science Foundation of China (Grant Nos. 60273003, 60433050, and 60305005). R. Duan acknowledges the financial support of Tsinghua University (Grant No. 052420003).
|
warning/0506/cond-mat0506716.html
|
ar5iv
|
text
|
# Phase Diagrams and Crossover in Spatially Anisotropic d=3 Ising, XY Magnetic and Percolation Systems: Exact Renormalization-Group Solutions of Hierarchical Models
## I Introduction
Spatially anisotropic systems greatly enrich our experience of collective phenomena, as exemplified by high-$`T_c`$ superconducting materials, in which the couplings along one direction are much weaker than those in the perpendicular plane. Anisotropic systems are also intriguing from a conceptual point of view, since vastly different critical phenomena are known to happen in different spatial dimensions, whereas, between $`d`$-dimensional systems stacked along a new direction, even the weakest coupling, while not affecting the critical temperature, induces $`(d+1)`$-dimensional critical behavior. Calculational results that yield the global phase diagram of anisotropic systems and thus provide a unified connected picture of the various anisotropic and isotropic behaviors at different subdimensions and at the full dimension have been rare and mostly confined to $`d=2`$. In the present study, we obtain global phase diagrams for a variety of anisotropic $`d=3`$ systems: Ising magnetic, XY magnetic, and percolation systems. Anisotropy along one direction (uniaxial) and full anisotropy, in which the couplings along each direction are different, are studied, yielding global phase diagrams. We use hierarchical models, which yield exact renormalization-group solutions Berker ; Kaufman ; Kaufman2 . Thus, the construction of hierarchical lattices that incorporate correct dimensional reductions is an important step of the study. The exact solutions of hierarchical models can simultaneously be considered approximate position-space renormalization-group solutions of models on naturally occurring lattices Berker . The method developed in this study will be employed to extend, from isotropic to anisotropic systems, the renormalization-group solutions of the tJ and Hubbard models of electronic conduction Falicov ; Hinczewski ; Hinczewski2 .
## II Anisotropic Hierarchical Lattices
Hierarchical lattices are constructed by repeatedly self-imbedding a graph. These provide exactly solvable models, with which complex problems can be studied and understood. For example, frustrated McKay , spin-glass Migliorini , random-bond Andelman and random-field Falicov2 , Schrödinger equation Domany , lattice-vibration Langlois , dynamic scaling Stinchcombe , aperiodic magnet Haddad , complex phase diagram Le , and directed-path daSilveira systems, etc., have been solved on hierarchical lattices.
In this study, we construct anisotropic hierarchical lattices by the parallel, mutual imbedding of several graphs. In each imbedding step, $`b`$ and $`b^d`$ respectively are the length and volume rescaling factors. We illustrate the method by the simplest case of the anisotropic $`d=2`$ lattice, before moving on to the uniaxially or fully anisotropic $`d=3`$ lattices. The parallel, mutual imbeddings of the two graphs shown in Fig.1 provide an anisotropic $`d=2`$ hierarchical lattice. If either of the couplings ($`K_x,K_y`$) is set to zero, the remaining coupling constitutes a one-dimensional lattice. When the couplings are of equal strength, $`K_x=K_y`$, the two directions, represented by the two imbedding sequences, are equivalent and the lattice is isotropic $`d=2`$. This lattice will be referred to as $`A_2`$.
It is thus seen that generally our requirements in the construction of anisotropic hierarchical lattices are (1) the proper reduction to the lower dimension when one (or more, see below) of the couplings is set to zero and (2) the restitution of an isotropic lattice when the couplings are of equal strength.
The parallel, mutual imbeddings of the two graphs shown in Fig.2 provide a uniaxially anisotropic $`d=3`$ hierarchical lattice. If the coupling $`K_z`$ is set to zero, the coupling $`K_{xy}`$ constitutes an isotropic two-dimensional lattice. If the coupling $`K_{xy}`$ is set to zero, the coupling $`K_z`$ constitutes a one-dimensional lattice. When the couplings are of equal strength, $`K_{xy}=K_z`$, the $`z`$ direction, represented by the first imbedding sequence, and the $`x,y`$ directions, represented by the second imbedding sequence, are all equivalent and the lattice is isotropic $`d=3`$. This lattice will be referred to as $`U_3`$. In Fig.3, the parallel, mutual imbeddings of the top graph with either one of the following graphs also provide uniaxially anisotropic $`d=3`$ hierarchical lattices. If the last graph is used, isotropy is not restored when $`K_{xy}=K_z`$; this lattice is nevertheless included, for comparison, in our study. These lattices, differentiated by the choice of the second imbedding graph, will be respectively referred to as $`U_{3a},U_{3b},U_{3c}`$.
A fully anisotropic $`d=3`$ hierarchical lattice is provided in Fig.4 by each shown imbedding in parallel with the two imbeddings obtained by permuting $`K_x`$ (full line), $`K_y`$ (dashed), and $`K_z`$ (dotted). If any one the couplings $`K_u`$ is set to zero, the remaining two couplings constitute an anisotropic two-dimensional lattice. If any two of the couplings are set to zero, the remaining coupling constitutes a one-dimensional lattice. When the couplings are of equal strength, $`K_x=K_y=K_z`$, the three directions, represented by the three imbedding sequences, are equivalent and the lattice is isotropic $`d=3`$. These lattices will be referred to as $`A_{3a}`$ and $`A_{3b}`$.
The anisotropic systems that we study are located on the anisotropic lattices constructed above. These hierarchical models admit exact renormalization-group solutions, with recursion relations obtained by decimations in direction opposite to their construction direction. The exact solutions of hierarchical models can simultaneously be considered approximate position-space renormalization-group solutions of models on naturally occurring lattices. In fact, the recursion relations obtained for the models below correspond to Migdal-Kadanoff Migdal ; Kadanoff approximate recursion relations, which are hereby generalized to anisotropic systems.
## III Anisotropic Ising Magnets
The Ising model is defined by the Hamiltonian
$`\beta H=`$ $`{\displaystyle \underset{u}{}}K_u{\displaystyle \underset{ij_u}{}}s_is_j,`$ (1)
where, at each lattice site $`i`$, $`s_i=\pm 1`$, and $`<ij>_u`$ denotes summation over bonds of type $`u`$. The various decimations in the models are composed of two elementary steps, $`K=K_u+K_v`$ for bonds in parallel and $`K=\mathrm{tanh}^1(\mathrm{tanh}K_u+\mathrm{tanh}K_v)`$ for bonds in series, where $`K`$ is the effective coupling of the combined bonds.
The phase boundaries for the Ising model on the $`d=2`$ anisotropic hierarchical lattices $`A_2,A_{3a},`$ and $`A_{3b}`$ (setting $`K_z=0`$ in the latter two) are given in Fig.5, along with the exact result for the anisotropic square lattice Onsager . The renormalization-group flows are indicated on the phase boundary of the hierarchical models. The fixed point occurs at isotropy, $`K_x=K_y`$, to which the $`d=2`$ anisotropic critical points flow, thereby sharing the same critical exponents.
The phase boundaries for the Ising model on the $`d=3`$ uniaxially anisotropic hierarchical lattices $`U_3,U_{3a},U_{3b},U_{3c},A_{3a}`$, and $`A_{3b}`$ (setting $`K_x=K_y`$) are given in Fig.6. The exact phase transition points for the square Onsager and cubic Ferrenberg lattices are also shown. For each model, the phase transitions at $`d=1`$ (at infinite coupling) and $`d=2`$ cross over to $`d=3`$ criticality, which is thus universal for all $`d=3`$ anisotropic and the $`d=3`$ isotropic cases.
The phase boundary surface for the Ising model on the $`d=3`$ fully anisotropic hierarchical lattice $`A_{3b}`$ is given in Fig.7. The dashed lines on the planes are the exact $`d=2`$ solutions for the square lattice Onsager . Again, all points on the critical surface of the $`d=3`$ fully anisotropic model flow onto the fixed point located at isotropy, thereby sharing its critical exponents. The critical exponents found for this model are $`y_T=0.69,y_H=1.68`$ for $`d=2`$ (for the square lattice $`y_T=1,y_H=1.875`$ Onsager ) and $`y_T=0.92,y_H=2.20`$ for $`d=3`$ (for the cubic lattice $`y_T=1.59,y_H=2.50`$ Moore ; Gaunt ).
## IV Anisotropic XY Magnets
The XY model is defined by the Hamiltonian
$`\beta H=`$ $`{\displaystyle \underset{u}{}}J_u{\displaystyle \underset{ij_u}{}}𝐬_i𝐬_j={\displaystyle \underset{u}{}}J_u{\displaystyle \underset{ij_u}{}}\mathrm{cos}(\theta _i\theta _j),`$ (2)
where at each lattice site $`i`$, $`𝐬_i`$ is a unit vector confined to the $`xy`$ plane at angle $`\theta _i`$ to the $`x`$ axis and $`<ij>_u`$ denotes summation over bonds of type $`u`$. Under renormalization-group transformations, the coupling between nearest-neighbor sites takes the general form of a function $`V_u(\theta _i\theta _j)`$. The various decimations in the models are composed of two elementary steps,
$`V=`$ $`V_u+V_v\text{and}`$ (3)
$`V(\theta _i\theta _k)=`$ $`\mathrm{ln}{\displaystyle _0^{2\pi }}𝑑\theta _j\mathrm{exp}[V_u(\theta _i\theta _j)+V_v(\theta _j\theta _k)],`$
respectively for bonds in parallel and in series, where $`V`$ is the effective coupling of the combined bonds. In terms of Fourier components,
$`f_u(s)=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta }{2\pi }}e^{is\theta }\mathrm{exp}[V_u(\theta )V_u(0)],`$ (4)
$`\mathrm{exp}[V_u(\theta )`$ $`V_u(0)]={\displaystyle }_se^{is\theta }f_u(s),`$
Eqs.(3) respectively are
$`f(s)=`$ $`{\displaystyle \underset{p}{}}f_u(p)f_v(sp)\text{and}`$ (5)
$`f(s)=`$ $`f_u(s)f_v(s),`$
in a form that is more conveniently followed in our calculations. The phase boundaries for the XY model on the $`d=3`$ uniaxially anisotropic hierarchical lattices $`U_3,U_{3a},U_{3b},`$ and $`U_{3c}`$ are given in Fig.8. The exact phase transition points for the square Hasenbusch and cubic Ferer lattices are also shown.
In $`d=2`$, namely along the horizontal axis, above a critical interaction strength marked by the squares on the figure, the systems exhibit algebraic order Kosterlitz ; Jose ; Berker2 : The starting Hamiltonian \[Eq.(2)\] flows to a Villain potential Villain ,
$`f_V(s)=`$ $`A\mathrm{exp}(s^2/2J_V),`$ (6)
exhibiting a fixed-line behavior parametrized by $`J_V`$. This corresponds to a system without true long-range order, namely with zero magnetization, but infinite correlation length and algebraic order in which the correlations decay as a power law. In $`d=3`$, true long-range order occurs: points in the ferromagnetic phase renormalize to a delta function potential; points on the phase boundaries renormalize to single true fixed potential, shown in Fig.9, differing from the Villain potential as also seen on the figure. The behavior here for $`d=2`$ is not true fixed-line behavior. After tens of thousands of renormalization-group iterations (corresponding to a scale change factor of 210,000), the Villain potential decays Berker2 to a disordered sink with $`V(\theta )_{max}V(\theta )_{min}<10^4`$. The sharp change in the necessary number of iterations, as seen in Fig.10, indicates the onset of effective algebraic order.
As seen in Fig.8, for each XY model, the phase transitions at $`d=1`$ (at infinite coupling) and $`d=2`$ (onset of algebraic order) cross over to $`d=3`$ criticality, which is thus universal for all $`d=3`$ anisotropic and the $`d=3`$ isotropic cases.
## V Anisotropic Percolation
Anisotropic percolation is defined such that on each connection of direction $`u`$, a bond exist with probability $`p_u`$. The various decimations in the models are composed of two elementary steps, $`p=p_up_v+p_u(1p_v)+p_v(1p_u)`$ for connections in parallel and $`p=p_up_v`$ for connections in series, where $`p`$ is the effective connectedness probability of the combined connections. The phase diagram for percolation on the $`d=2`$ anisotropic hierarchical lattice ($`A_2`$) is given in Fig.11. The percolation fixed point occurs at isotropy, $`p_x=p_y`$, to which the $`d=2`$ anisotropic percolation onsets flow, thereby sharing the same critical exponents. The phase boundaries for percolation on the $`d=3`$ uniaxially anisotropic hierarchical lattices $`U_3,U_{3a},U_{3b}`$, and $`U_{3c}`$ are given in Fig.12. The percolation points for the cubic and square lattices are also shown Essam . For each model, percolation onset at $`d=1`$ (at $`p_z=1`$) and $`d=2`$ cross over to $`d=3`$ percolation onset, which is thus universal for all $`d=3`$ anisotropic and the $`d=3`$ isotropic cases.
This research was supported by the Scientific and Technical Research Council of Turkey (TÜBITAK) and by the Academy of Sciences of Turkey.
|
warning/0506/hep-lat0506021.html
|
ar5iv
|
text
|
# Chiral properties of two-flavor QCD in small volume and at large lattice spacing
## I Introduction
Chiral symmetry is a fundamental part of the theory of strong interactions and therefore should be respected when putting QCD on the lattice. It was realized more than two decades ago, that this can be done with Dirac operators $`D`$ which (at zero quark mass) obey the Ginsparg–Wilson equation Ginsparg:1981bj
$$D\gamma _5+\gamma _5D=\frac{a}{R_0}D\gamma _5D$$
(1)
where $`a`$ is the lattice spacing and $`R_0`$ the radius of the Ginsparg–Wilson circle. A decade ago, this constraint was realized by overlap fermions Neuberger:1997fp ; Neuberger:1998my , fixed-point fermions Hasenfratz:1993sp and domain wall fermions Kaplan:1992bt ; Furman:1994ky . Among those, only domain wall fermions have been used for some time in simulations on the lattice which include the fermionic determinant. Only recently the first steps toward dynamical simulations using overlap fermions have been taken Bode:1999dd ; Fodor:2003bh ; Fodor:2004wx ; DeGrand:2004nq ; Cundy:2005pi . Because of the high cost of applying the Dirac operator these are still limited to a small volume. In this paper we present results from simulations with two flavors of dynamical overlap fermions in a small box and at a large lattice spacing. We measure the topological susceptibility and, by comparing the low-lying eigenvalue distribution to random matrix theory, the chiral condensate. We also fix the lattice spacing and look for dynamical effects in meson two-point functions. The measurement of the susceptibility is greatly facilitated as compared to, e.g., simulations using highly improved staggered fermions Aubin:2004qz . With the overlap operator, we can use the same Dirac operator in the simulation and in the definition of the topological charge via the index theorem. The topological charge defined in this way thus directly influences the weight of a configuration.
A more technical consequence of that is the discontinuity of the fermionic determinant as a function of the gauge variables. The surfaces in gauge field space at which the fermionic action is discontinuous coincide with the change in topology defined by the index of the Dirac operator. Although this does not pose a problem in principle, in practice changing the topological sector turns out to be a significant problem. As recent simulations with highly improved staggered fermions show, this problem itself, however, is not restricted to chiral formulations of QCD on the lattice Bernard:2003gq .
In a recent paper DeGrand:2004nq , we described the setup of our simulations of two degenerate flavors of overlap fermions. We studied the impact of fat (stoutMorningstar:2003gk ) links and multiple pseudo-fermion fields. Some improvements to the algorithm will be reported in the following. The simulations presented in this paper used very limited computer resources, i.e. half a year on an array of 12 3.2 Ghz Pentium-IVE’s. With that we present some results on $`8^4`$ lattices and a relatively coarse lattice spacing of around $`a=0.16\mathrm{fm}`$ and a quark mass down to about $`35\mathrm{MeV}`$.
This paper is organized as follows: We first review the algorithm, describe the improvements over our previous publication and give the parameters of our simulation. Then in Sec. IV we attempt to set the lattice spacing by determining the Sommer parameter $`r_0`$ and proceed with the extraction of the topological susceptibility in Section V. In Section VI we look at meson two-point functions. By comparing to quenched results on matched lattices, we demonstrate the impact of dynamical fermions on the scalar correlator. Finally, we compare the low-lying spectrum of the overlap operator with the predictions from random matrix theory. A note in caution: for all these investigations, the volume which we are simulating is too small. Future simulations will improve on this. Here we want to demonstrate that simulations with dynamical overlap fermions are possible and give results which match our expectations of full QCD in a small box.
## II Definitions and Algorithm
Let us fix the conventions and describe the algorithm together with the improvements since our previous publication DeGrand:2004nq , to which we refer the reader for more details. The massive overlap operator is given by Neuberger:1997fp ; Neuberger:1998my
$$D_{ov}(m)=(R_0\frac{m}{2})\left[1+\gamma _5ϵ(h(R_0))\right]+m$$
(2)
with $`ϵ(h)=h/\sqrt{h^2}`$ the sign function of the Hermitian kernel operator $`h=\gamma _5d`$ which is taken at the negative mass $`R_0`$. $`R_0`$ is the radius of the Ginsparg–Wilson circle. We are using a planar kernel Dirac operator $`d`$ with nearest and next-to-nearest (“$`\sqrt{2}`$”) interactions. For details see Ref. DeGrand:2004nq . The sign function is computed using the Zolotarev approximation with an exact treatment of the low-lying eigenmodes $`|\lambda `$ of $`h(R_0)`$
$$ϵ(h(R_0))=\underset{i}{}\mathrm{sign}(\lambda _i)|\lambda _i\lambda _i|\underset{i=1}{\overset{n_{\mathrm{eig}}}{}}\mathrm{sign}(\lambda _i)|\lambda \lambda |+h\underset{j}{}\frac{b_j}{h^2+c_j}(1\underset{i=1}{\overset{n_{\mathrm{eig}}}{}}|\lambda _i\lambda _i|).$$
(3)
For our simulation we use the Hybrid Monte Carlo (HMC) algorithm Duane:1987de as modified for overlap fermions by Ref. Fodor:2003bh . The effective action is given by
$$S_{\mathrm{eff}}=\varphi _0^+H^2(m_0)\varphi _0+\underset{i=1}{\overset{n_{pf}1}{}}\varphi _i^+H^2(m_i)H^2(m_{i1})\varphi _i+S_g[U]$$
(4)
with $`H^2(m)=D_{ov}(m)^+D_{ov}(m)`$ the square of the hermitian overlap Dirac operator. $`S_g[U]`$ is the gluonic action. The $`\varphi _i`$ are the $`n_{pf}`$ pseudo-fermion fields used to include the contribution of the fermion determinant. The use of several of these fields has been suggested in Hasenbusch:2001ne ; Hasenbusch:2002ai . It improves the stochastic estimate of the determinant. We studied its effects extensively in Ref. DeGrand:2004nq .
With the HMC algorithm, an ensemble distributed according to this effective action is generated by updating the gauge fields using molecular dynamics trajectories with a final accept-reject step. At the beginning of the trajectory one chooses momenta $`\pi `$ conjugate to the gauge fields and also refreshes the pseudo-fermions. One then integrates the resulting equations of motion (treating $`S_{\mathrm{eff}}`$ as the potential) numerically in some fictitious simulation time $`\tau `$. They result from the requirement that the total ‘energy’ $`=\pi ^2/2+S_{\mathrm{eff}}`$ is conserved. At the end one applies an accept/reject step with acceptance probability $`P=\mathrm{min}[1,\mathrm{exp}(\mathrm{\Delta })]`$ which corrects for the errors in the numerical integration.
An important difference between conventional fermions and overlap fermions is that the effective action $`S_{\mathrm{eff}}`$ is discontinuous for the latter. The discontinuity has its origin in the sign function in the definition of the Dirac operator Eqs. (2) and (3). The fermionic action is discontinuous at the surfaces where the kernel operator $`h(R_0)`$ has a zero eigenvalue. According to the index theorem Hasenfratz:1998ri , these are also the surfaces where the topological charge (as seen by the fermions) changes. Ref. Fodor:2003bh gives the prescription for how to deal with this situation. The molecular dynamics evolution can be thought of as resembling that of a classical particle in the presence of a potential step. If the step in the action is too big for the particle to get across it, the particle is reflected, i.e. the momentum component normal to the zero eigenvalue surface is reversed. On the other hand, if the normal component is large enough to change topological sector, the normal component is reduced such that energy is conserved. Following Ref. Fodor:2003bh this is called a refraction. The momentum is then altered according to
$$\mathrm{\Delta }\pi =\{\begin{array}{cc}NN|\pi +N\mathrm{sign}N|\pi \sqrt{N|\pi ^22\mathrm{\Delta }S_f}\hfill & \text{if }N|\pi ^2>2\mathrm{\Delta }S_f\hfill \\ 2NN|\pi \hfill & \text{if }N|\pi ^22\mathrm{\Delta }S_f\hfill \end{array}$$
(5)
with $`N`$ the vector normal to the zero eigenvalue surface, $`\pi `$ the momentum and $`\mathrm{\Delta }S_f`$ the discontinuity in the fermionic action. To monitor whether an eigenvalue has changed sign, one thus has to compute some number $`n_{\mathrm{eig}}`$ of the lowest eigenmodes of $`h(R_0)`$ and see whether the eigenvalue of any of them changes sign. The matching is done by computing the scalar products of the low modes before and after a step. Since the eigenmodes are needed anyway to precondition the construction of the sign function, there is virtually no overhead associated with this test.
Note that this part of the algorithm potentially scales with the square of the volume: The cost of determining the height of the step is at least proportional to the volume. The number of times this procedure has to be executed can be assumed to be proportional to the density of eigenmodes of the kernel operator at the origin, which in turn might be proportional to the volume. It is therefore pivotal to keep the cost of this step as low as possible.
A major improvement in the algorithm is the way in which we compute the height of the step. In our previous publication, we ran two conjugate gradients to compute
$$\mathrm{\Delta }S=\mathrm{\Delta }\left[\varphi _0^+H^2(m_0)\varphi _0+\underset{i=1}{\overset{N1}{}}\varphi _i^+H^2(m_i)H^2(m_{i1})\varphi _i\right]$$
(6)
where the difference is taken between the fermionic actions for which only the sign of the lowest eigenmode (the one which becomes zero on the surface) is changed without changing any of the other modes or the gauge configuration. This was very expensive since we have to decide frequently whether to refract or reflect. Because the square of the hermitian Dirac operator in one chiral sector $`\sigma =\pm 1`$ is of the form
$$H_\sigma ^2(m)=2(R_0^2\frac{m^2}{4})P_\sigma \left[1+\sigma \underset{i}{}ϵ(\lambda _i)|\lambda _i\lambda _i|\right]P_\sigma +m^2$$
(7)
with $`|\lambda `$ the eigenvectors of $`h(R_0)`$ and $`P_\sigma =\frac{1}{2}(1+\sigma \gamma _5)`$ the projector on the chiral sector, the sign change of the crossing mode amounts to the change
$$H_\sigma ^2(m)H_\sigma ^2(m)\pm (4R_0^2m^2)P_\sigma |\lambda _0\lambda _0|P_\sigma $$
(8)
with $`|\lambda _0`$ the zero mode and the sign being minus the product of the sign associated with the chiral sector and the sign of the eigenmode before the step. Thus, we can use the Sherman-Morrison formula Golub to compute the height of the step. The result is
$$\mathrm{\Delta }\left[\varphi |P_\sigma \frac{1}{H_\sigma (m)^2}P_\sigma |\varphi \right]=\frac{(4R_0^2m^2)}{1\pm (4R_0^2m^2)\lambda _0|P_\sigma H_\sigma ^2(m)P_\sigma |\lambda _0}|\varphi |P_\sigma \frac{1}{H_\sigma (m)^2}P_\sigma |\lambda _0|^2.$$
(9)
This has the additional advantage that instead of running two conjugate gradients for each pseudo-fermion field, we only have to invert once, using the eigenmode which changes sign as a source. Because the height of the step is directly computed, one also has better control over the accuracy — compared to taking the difference of two approximate quantities.
One can further exploit Eq. 9 by realizing that the cost the of the inversion of $`H^2(m)`$ depends on the chirality of the source and the topological sector in which one is inverting. One can compute the height of the step from either side of the surface. In the chiral sector with zero-modes, it is cheaper to invert in the topological sector of lower charge. However, since the zero-modes push the spectrum of the other modes up, the conditioning number of $`H_\sigma ^2(m)`$ in the sector without zero-modes is lower on the side of the surface with higher topology. A second advantage of the use of this formula is that we can monitor the step-height during the CG iterations and terminate the iteration when it has become clear that we are going to reflect.
Finally, let us report a small improvement in the computation of the fermion force $`\delta S_f[U]`$. The derivative of the Zolotarev part of the approximation to the sign function Eq. (3) has been given in many places. For each pseudo-fermion field and each order in the rational approximation the formula has a term $`1/(h^2c)`$, the derivative of which is
$$\frac{1}{h^2c}(h\delta h+\delta hh)\frac{1}{h^2c}.$$
(10)
One thus has to invert the kernel action, which can be done simultaneously for all shifts using a multi-mass algorithm. The computation of $`\delta h`$ follows via standard methods from Ref. Gottlieb:1987mq . However, due to the many shifts in the Zolotarev approximation and several pseudo-fermion fields, this part is a non-negligible contribution to the total cost of the simulation.
More difficult is the derivative of the projector term $`P_\lambda =|\lambda \lambda |`$. This derivative is basically given by first-order perturbation theory (See Ref. Narayanan:2000qx )
$$\delta P_\lambda =\frac{1}{\lambda h}(1P_\lambda )\delta hP_\lambda +P_\lambda \delta h^+\frac{1}{\lambda h}(1P_\lambda ).$$
(11)
Because $`(\lambda h)`$ is singular, its inversion is problematic even though the contribution of the eigenmode with eigenvalue $`\lambda `$ is projected out of the source: since both $`\lambda `$ and $`|\lambda `$ are only approximately known, one faces a “zero divided by zero” problem. In our previous publication, we therefore shifted the pole position, performed the inversion, and interpolated our result. This turned out to be insufficiently stable. Now we use a Chebychev approximation to the inverse of $`h^2\lambda ^2`$ in the range such that the inverse is precise outside the known eigenvalues of $`h(R_0)`$. The advantage is that this approximation is finite at $`h^2\lambda ^2=0`$. The problem is thus reduced to a contribution from the eigenvalue $`\lambda `$ mode which is zero times some finite number given by the required accuracy and the range in which one computes the eigenvalues explicitly.
Eq. 7 provides us with a tantalizing result, which we do not know how to apply in the context of HMC: an exact formula for the ratio of the fermion determinants on either side of the topology-changing boundary:
$$\frac{det\stackrel{~}{H}_\sigma ^2(m)}{detH_\sigma ^2(m)}=1\pm (4R_0^2m^2)\lambda _0|P_\sigma \frac{1}{H_\sigma ^2(m)}P_\sigma |\lambda _0.$$
(12)
HMC does not ever use the exact determinant as part of the simulation. Instead, it generates configurations whose statistical weight is controlled by the effective action Eq. 4. In an algorithm which does approximate the determinant directly, like the $`R_0`$ algorithm, it seems straightforward; one would just use the logarithm of Eq. 12 as the step. However, we are unwilling to abandon HMC for two reasons: First, we feel that we benefit substantially from the fact that HMC is exact and therefore do not want to run an $`R`$ type algorithm which cannot be made exact in an obvious way. Second, we gain considerable speed using HMC over an R algorithm because in HMC we can use a previous solution to an inverse of $`H(m)^2`$ to begin the computation of the new force. This point deserves further investigation.
## III Simulation Parameters and Performance of the Algorithm
We simulate on $`8^4`$ lattices at one value of the gauge coupling $`\beta =7.2`$ which we chose to be roughly at the $`N_t=6`$ phase transition (for and overview over the lattice spacings see Table 3). We use the Lüscher–Weisz gauge action Luscher:1984xn with the tadpole improved coefficients of Ref. Alford:1995hw . Instead of determining the fourth root of the plaquette expectation value $`u_0=(U_{pl}/3)^{1/4}`$ self-consistently, we set it to 0.86 for all our runs as we did in our previous publication.
Our kernel operator $`d`$ is constructed from gauge links to which two levels of isotropic stout blocking Morningstar:2003gk have been applied. The blocking parameter $`\rho `$ is set to 0.15.
We report on simulations at three values of the bare sea quark mass $`am_q=0.03`$, $`0.05`$ and $`0.1`$. Based on measured lattice spacings from the Sommer parameter and the perturbative calculation of matching factors reported in the Appendix, we believe that these values correspond to $`\overline{MS}`$ quark masses of about 35, 55 and 100 MeV. An overview of our collected statistics is given in Table 1. The trajectories all have length one, divided in 20 elementary time steps. We use a Sexton-Weingarten Sexton:1992nu integration scheme in order apply a smaller time step for the gauge field integration. We perform 12 applications of the gauge force and gauge field update per elementary time step. Again, see Ref. DeGrand:2004nq for details.
In order to monitor whether an eigenvalue has changed sign, we compute the lowest $`8`$ eigenmodes of $`h(R_0)`$ in each step. These are also used to precondition the construction of the sign function.
For our analysis, we typically discard the first 100 configurations in a stream and separate two consecutive measurements by 5 trajectories. The separation of the configurations is based on our measurement of the auto-correlation time of the plaquette. It is around $`5(1)`$ varying little with the quark mass. We observed no significant difference in acceptance rate between the quark masses. However, the auto-correlation time of the topological charge differs enormously. This is due to the fact that the estimation of the step height of the fermionic action intrinsic to HMC is much more subject to fluctuation for smaller masses than for larger ones. As argued in our previous publication, a poor estimator results in reflections. We partially address this problem by the use of multiple pseudo-fermion fields, but the problem remains. In Fig. 1 we show a scatter plot of the real change in the determinant at the step, from Eq. 12, as compared to the stochastic estimate with three pseudo-fermion fields. We subtracted the normal component $`N,\pi ^2`$ of the momentum so that, according to Eq. (5), negative values allow refractions. However, $`N,\pi ^2`$ is almost always smaller than 10. The stochastic estimate of the step height has a wide spread but is typically large. Since $`\mathrm{exp}(\mathrm{\Delta }S)`$ will average to the ratio of the fermionic determinants on both sides of the step, this is a consequence of the large fluctuation in the estimator. (A few small values of $`\mathrm{\Delta }S`$ have to be compensated by a large number of large ones, for which $`\mathrm{exp}(\mathrm{\Delta }S)`$ is approximately zero.)
The low correlation between the estimator and the physical step height Eq. (12) shows up in the large auto-correlation time of the topological charge, whose time history is shown in Fig. 2. Even though part of it is physics — lighter quarks make it harder to get from, e.g., $`\nu =0`$ to $`\nu =\pm 1`$ — the height of the step grows with $`1/m^2`$ instead of the expected determinant ratio, $`\mathrm{log}m`$. Since the normal component of the momentum is roughly independent of the quark mass, it becomes more and more difficult to change topology (also see discussion in Sec. V). The large auto-correlation time for the topology is a phenomenon that is also known with other fermions, e.g. improved staggered quarks. To the extent that these formulations know about topology, the step in the fermion action for the overlap might be replaced for them by a steep region which approximates the step. The result is the same: if the approximation of the determinant is bad, the step is overestimated most of the time and one does not change topology.
Let us finally take a look at the relative cost of the various ingredients of the algorithm. In Table 2 we list the cost per call and the fraction of the total cost of each of the major parts of the program. The conjugate gradient inversion of $`H_{\mathrm{ov}}^2`$ needed for the computation of the force, starting action and the reflections/refractions takes by far the largest fraction. Even though one inverts only on one source (the zero-mode of $`h(R_0)`$) in the reflection/refraction routine, these inversions are very expensive since there is no good starting vector. Therefore, they alone take a fifth to a quarter of the total cost, depending on the quark mass. The inversions can be cheaper for the lighter quark mass because there is less topology; the conditioning number of $`H_{\mathrm{ov}}`$ for $`\nu =0`$ is lower than for $`|\nu |>0`$.
The cost of computing the $`8`$ eigenvectors of the kernel operator is small, about $`10\%`$. It is cheap because the eigenmodes of the kernel do not change much during the evolution; one thus has good starting vectors. We need them to precondition the construction of the overlap operator and to monitor whether an eigenvalue has changed sign. Finally, the computation of the fermion force (outside the inversion of $`H_{\mathrm{ov}}^2`$) takes about a sixth of the total time. This is due to the inversion of the kernel operator, the computation of $`\psi _i^+\delta h(R_0)\psi _i`$ for each of the poles in the Zolotarev approximation of the sign function and the inversion of $`h(R_0)`$ for the projector term in the sign function as discussed at the end of Section II.
## IV Setting the scale
In order to get an idea at which lattice spacing we are simulating, we measured the heavy quark potential and extracted the Sommer parameter $`r_0`$ Sommer:1993ce . (It is defined as the distance at which the force between two heavy quarks $`F(r)`$ with distance $`r`$ satisfies $`r_0^2F(r_0)=1.65`$.) This is necessary because we expect a substantial shift due to the dynamical fermions Hasenfratz:1993az with respect to the quenched lattice spacing Gattringer:2001jf . Unfortunately, our lattice is relatively small. It is thus not possible to estimate the uncertainty in the extraction of the potential from the data alone. To illustrate this we show in Fig. 3 the potential extracted from Wilson loops $`W[r,t]`$ which have been constructed from HYP smeared linksref:HYP . The two sets correspond to single exponential fits to the $`W[r,t]`$ with $`t[2,4]`$ and $`t[3,4]`$, respectively. The latter set gives a significantly smaller potential. The curves represent fits of the form
$$V(r)=\frac{A}{r}+Br+C+Df(r)$$
(13)
where $`f(r)`$ is a perturbatively determined correction Hasenfratz:2001tw to account for the effect of the HYP links on short distances. The results are given in Tab. 3.
To estimate the systematic error from the small volume, we have generated a quenched set of $`8^4`$ and $`12^4`$ lattices at $`\beta =7.77`$, $`u_0=0.887`$ with a similar lattice spacing ($`a=0.163(1)`$). Analyzing these lattices, we find that the $`r_0`$ extracted from the $`8^4`$ lattices with fit range $`t[3,4]`$ is the same within error-bars as the one extracted from the $`12^4`$ and fit ranges starting at $`3`$ or $`4`$ ranging to $`5`$ or $`6`$. In the following, we will therefore work with the lattice spacing extracted form the $`t[3,4]`$ fit range.
The dimensionless quantity $`r_0\sqrt{\sigma }`$ has been used in the past to quantify the impact of dynamical quarks on the shape of the potential Bernard:2001av . We find $`r_0\sqrt{\sigma }=1.10(1)`$ almost independent of the sea quark mass and $`1.18(1)`$ for our quenched ensemble. On larger lattices and a finer lattice spacing, Ref. Bernard:2001av found a quenched value of about $`1.16`$ and 1.128 for two flavors of dynamical staggered quarks. From Ref. AliKhan:2001tx we get a value of about 1.14 on at a similar lattice spacing with two flavors of dynamical clover Wilson fermions. Given the systematic and statistical errors these values agree well with our findings.
## V Topological susceptibility
The topological susceptibility illustrates the strengths and weaknesses of our simulation. In contrast to all simulations with non-chiral actions, the measurement of the topological charge in an overlap simulation is trivial. It can even be done during the simulation by monitoring zero crossings of the smallest eigenmode of the kernel operator (if the topological charge has been determined once at the beginning of the simulation). However, as we have already remarked, the autocorrelation time of the topological charge during the simulation is annoyingly long.
We begin by showing (in Fig. 2) time histories of the topological charge for the different simulations performed at our three values of dynamical quark mass. Histograms of the topological charge (including the thermalization runs) are shown in Fig. 5. The topology is recorded at the end of each trajectory (rather than at the end of each time-step). No sophisticated analysis is needed to see that the autocorrelation time grows as the quark mass falls, and Fig. 4 shows that the mean time between topological changes varies inversely with the square of the quark mass.
With such long times between tunnelings, we are concerned about thermalization effects in our data. At $`am_q=0.1`$ $`\nu `$ is zero within statistical uncertainty (it is -0.05(13) throughout the run) and $`\nu ^2`$ seems to have stabilized after 50 trajectories are discarded, so we cut the data there. At $`am_q=0.03`$ $`\nu `$ is -0.15(11) when 100 trajectories are discarded from each run and 0.04(12) when 150 are discarded; $`\nu ^2`$ shows little variation with cuts of more than 50 initial trajectories per stream, and so we dropped the first 100 trajectories. At $`am_q=0.05`$ the situation is similar; we see little variation cutting more than 90-100 configurations and again dropped 100 before averaging.
We find topological susceptibilities of $`\chi a^4=2.17(29)\times 10^4`$ at $`am_q=0.1`$, $`\chi a^4=1.37(39)\times 10^4`$ at $`am_q=0.05`$, and $`\chi a^4=1.02(24)\times 10^4`$ at $`am_q=0.03`$.
For the quenched data we are able to space configurations used in the analysis far enough apart in simulation time as to be essentially uncorrelated. We measured a lattice topological susceptibility of $`\chi a^4=6.13(76)\times 10^4`$. With the quenched $`r_0/a=3.08`$, this is $`\chi (191`$ MeV$`)^4`$, which is quite consistent with typical quenched results, e.g. Gattringer:2002mr .
We attempt to translate this data into dimensionless units in order to facilitate comparisons with other measurements. We take our measurements of $`r_0/a`$ from the previous section to compute $`\chi r_0^4`$ and do the same with our quark masses, using the Z-factor as described in the Appendix to convert them to $`\mu =2`$ GeV $`\overline{MS}`$ values. We present our results in Fig. 6. DürrDurr:2001ty has presented a phenomenological interpolating formula for the mass dependence of the topological susceptibility, in terms of the condensate $`\mathrm{\Sigma }`$ and quenched topological susceptibility $`\chi _q`$,
$$\frac{1}{\chi }=\frac{N_f}{m_q\mathrm{\Sigma }}+\frac{1}{\chi _q}.$$
(14)
Taking $`\mathrm{\Sigma }`$ from our RMT analysis in the next section ( $`r_{0}^{}{}_{}{}^{3}\mathrm{\Sigma }=0.43`$) produces the curve shown in the figure.
Lattice results presented elsewhere typically use the pseudo-scalar mass as the ordinate. As we will describe in the next section, we do not have reliable pseudo-scalar masses because our lattices are too small. However, we can use the Dürr formula as a benchmark to compare with other simulations. Our results are in rough agreement in magnitude with those of another dynamical overlap simulation, at lattice spacing $`a0.25`$ fm, by Fodor, et. al.Fodor:2004wx . The two simulations both lie below the Dürr curve.
Data from simulations with two flavors of ordinary staggered simulations and was analyzed and published by (among others) A. HasenfratzHasenfratz:2001wd . At finite lattice spacing all this data lies far above the Dürr curve, and is not too different from the quenched result. Other dynamical fermion data with non-chiral fermionsAllton:2004qq is also high with respect to the Dürr formula and to overlap data. It is easy to imagine that simulations with non-chiral actions would overestimate the topological susceptibility since their massless Dirac operators do not have exact zero modes.
We do not think that the small volumes of our simulation have suppressed the susceptibility since our quenched simulations are not anomalously low. Little is known about the scaling properties of overlap actions, and our results and those of Ref. Fodor:2004wx have large lattice spacings. We are of course not satisfied with the quality of our data from the point of view of autocorrelations, lattice spacing, extraction of hadron masses, and simulation volume.
## VI Meson two-point functions
Let us now turn to meson two-point functions. The purpose of doing so is two-fold. First, we want to get an idea of the pion masses at which we are simulating. This will not work very well since the three volume is small and the time extent is far too small for the excited states to decay, but it can provide us with an upper limit of the pseudo-scalar mass. The second purpose is to compare the two-point functions to quenched results on similarly sized lattices and look for effects of the dynamical fermions in the scalar correlator. We compute zero-momentum correlators
$$C_{ij}(t)=\frac{1}{V}\underset{𝐱}{}\overline{\psi }(0,0)\mathrm{\Gamma }_i\psi (0,0)\overline{\psi }(𝐱,t)\mathrm{\Gamma }_j\psi (𝐱,t)$$
(15)
where the fermion fields $`\psi `$ are contracted with the appropriate flavor structure. We compute the quark propagators with Gaussian sources of radius $`2a`$ on gauge configurations in Coulomb gauge. We use point sinks and apply low-mode averaging using the four lowest eigenmodes of the Dirac operator DeGrand:2004qw ; DeGrand:2004wh .
In Fig. 7 we show the pseudo-scalar two-point function for our two smaller dynamical quark masses. Unfortunately, the corresponding masses, shown in Fig. 8, do not differ much. This can be attributed to the fact that $`T=8`$ is just too small. We thus see a superposition of excited states with the ground state which does not depend on the quark mass.
More impressive is the comparison of the scalar correlator from the dynamical and the quenched ensemble. It is known that this correlator turns negative in the quenched theory, which is a sign that quenched QCD is not unitary. In Fig. 9 (bottom) we show that this is the case even on a small lattice. The full two-point function is shown (connected by a line) as well as the contributions of the various topological sectors. Whereas the function is positive in $`\nu =0`$, the zero-modes in the sectors of non-trivial topology turn it negative. In the full theory, however, the fermion determinant suppresses the sectors with $`|\nu |>0`$ sufficiently for the whole two-point function to remain positive (Fig. 9 upper row).
## VII The condensate from Random Matrix Theory
It was proposed more than a decade ago that the distribution of the low-lying eigenvalues of the QCD Dirac operator in a finite volume can be predicted by random matrix theory (RMT) Shuryak:1992pi ; Verbaarschot:1993pm ; Verbaarschot:1994qf . Since then this hypothesis has received impressive support from lattice calculations, mainly quenched simulations Berbenni-Bitsch:1997tx ; Damgaard:1998ie ; Gockeler:1998jj ; Edwards:1999ra ; Giusti:2003gf , but also some dynamical ones using staggered quarks Berbenni-Bitsch:1998sy ; Damgaard:2000qt .
Typically, the predictions are made in the so-called epsilon regime, for which $`1/\mathrm{\Lambda }L1/m_\pi `$ with $`\mathrm{\Lambda }`$ a typical hadronic scale. However, it has been found that they describe the data in a wider range. In a recent large scale study, e.g., using the overlap operator on quenched configurations Giusti:2003gf , it could be shown that from lattices with a length larger than $`1.5\mathrm{fm}`$, the RMT predictions match the result of the simulation. Our dynamical lattices have a spatial extent of about $`1.3\mathrm{fm}`$. As we will see, random matrix theory describes our low-lying Dirac spectra quite well.
Our analysis is based on the distribution of the $`k`$-th eigenmode from RMT as presented in Ref. Damgaard:2000ah and successfully compared to simulation results in Ref. Damgaard:2000qt . The prediction is for the distribution of the dimensionless quantity $`\zeta =\lambda _k\mathrm{\Sigma }V`$ in each topological sector with $`\lambda _k`$ is the $`k`$-th eigenvalue of the Dirac operator, $`\mathrm{\Sigma }`$ the chiral condensate and $`V`$ the volume of the box. These distribution are universal and do not depend on additional parameters other than the number of flavors, the topological charge and the dimensionless quantity $`m_q\mathrm{\Sigma }V`$. (Note that topology effects the distributions, in contrast to the behavior of staggered fermions seen by Ref. Damgaard:2000qt .) By comparing the distribution of the eigenmodes with the RMT prediction one can thus measure the chiral condensate $`\mathrm{\Sigma }`$. The main advantage of this method is that it gives the zero quark mass, infinite volume condensate directly. The validity of the approach can be verified comparing the shape of the distribution for the various modes and topological sectors. The main uncertainty comes from a too small volume which causes deviations in the shape, particularly for the higher modes. As we will see below, the direct measurement of $`\mathrm{\Sigma }`$ from $`(\overline{\psi }\psi )(m)`$, which requires extrapolation into the chiral limit as well as correction for the finite volume, is unreliable for our data.
In Fig. 10 we show the distribution of the two lowest eigenmodes of the overlap operator (scaled by $`\mathrm{\Sigma }V`$) measured on the $`\nu =0`$ and $`\nu =\pm 1`$ parts of the $`am_q=0.03`$ and $`am_q=0.05`$ ensembles. To correct for the fact that our eigenvalues lie on a circle of radius $`R_0`$ instead of a straight line, we use the stereographic projection $`\stackrel{~}{\lambda }=|\lambda |\sqrt{1|\lambda |^2/2R_0}`$. We fit the RMT prediction from Ref. Damgaard:2000ah to these distributions. The chiral condensate $`\mathrm{\Sigma }`$ is the only free parameter in this fit; we get $`\mathrm{\Sigma }V/a=54(2)`$ and $`\mathrm{\Sigma }V/a=53(2)`$, respectively. This corresponds to $`r_0^3\mathrm{\Sigma }=0.46(2)`$ and 0.41(2). We studied the dependence of this result on the number of trajectories used in the thermalization and the separation of consecutive configurations. We found no systematic variation beyond omitting the first 100 trajectories and separating them by 5. The prediction agrees overall well with the measured distribution given the low statistics. However, the distribution of the lowest mode in the $`\nu =0`$ sector seems to have a tail at larger $`\lambda \mathrm{\Sigma }V`$ that does not match the prediction. This could be an effect of the small volume. We also show the prediction for the distribution of the third mode from our fitted values of $`\mathrm{\Sigma }V`$ in the third column of Fig. 10. The RMT curve and the data, again, agree quite well. However for the $`|\nu |=1`$ sector, the curve seems to be on the right of the data. This is probably a sign of the break-down of RMT for eigenvalues larger than the Thouless energy Osborn:1998nm ; Gockeler:1998jj .
The RMT predictions are made with the assumption that the volume is infinite. We (obviously) are not in that situation. In finite volume, in the epsilon-regime of chiral perturbation theory, finite volume modifies the formula for the condensate by multiplication by a shape factor, $`\mathrm{\Sigma }\rho \mathrm{\Sigma }`$, where
$$\rho =1+\frac{c(l_i/l)}{f_\pi ^2L^2}$$
(16)
and $`c(l_i/l)`$ depends on the geometryGasser:1986vb . We do not know $`\rho `$ since we have not measured $`f_\pi `$, but combining our lattice spacing and lattice size with $`f_\pi =93`$ MeV gives $`\rho 1.4`$. This is too large a correction to be trustworthy; again, we need a lattice with a larger physical size. Our two $`r_0^3\mathrm{\Sigma }`$ values must be multiplied by $`S_S`$ and $`\rho `$ to give a continuum value: with $`r_0=0.5`$ fm, $`\rho \mathrm{\Sigma }(\overline{MS})=`$ 0.032(1) or 0.029(1) GeV<sup>3</sup> or $`(\rho \mathrm{\Sigma }(\overline{MS}))^{1/3}=`$ 316 or 304 MeV. (Errors are from $`r_0/a`$ and the fit uncertainty in $`\mathrm{\Sigma }`$.)
We also attempted to measure the mass dependent quark condensate $`\overline{\psi }\psi (m)`$ using 12 random sources per lattice on subsets consisting of about 50 lattices per mass value, spaced by 5 trajectories. Our data is displayed in Fig. 11. We show the mass-dependent condensate summed over all topological sectors as well as its value in sectors of topological charge $`\nu =0`$ and $`|\nu |=1`$ (with the zero mode contribution removed for the latter case). As expected, there is no divergence in $`\overline{\psi }\psi (m)`$ at zero quark mass (as would occur in a quenched simulation).
The parameter $`\mathrm{\Sigma }`$ must be extracted from a fit of $`\overline{\psi }\psi (m)`$ to some functional form. For example, were we in a large volume, we would fit to the usual constant plus chiral logarithm plus $`Cm_q`$ formula.
We believe that we are, in fact, in a small volume, and that the appropriate fitting function is that for the so-called $`ϵ`$regime, given by Gasser and LeutwylerGasser:1986vb for the average-topology case, or by Leutwyler and SmilgaLeutwyler:1992yt , for the case of fixed topology. In either case condensate is given by
$$\mathrm{\Sigma }(V,\mu )=\frac{\rho \mathrm{\Sigma }}{Z}\frac{Z}{\mu }$$
(17)
with $`Z`$ the appropriate partition function written by the above authors or given equivalently by random matrix theory, and $`\mu =m_q\mathrm{\Sigma }V`$. We append a linear $`Cm_q`$ term to Eq. 17 and attempt to fit the parameters $`\mathrm{\Sigma }`$ and $`C`$ to the data. This was not successful.
First, we do not know which of our mass values lie in the region of validity of our fitting function. Fig. 12 shows a set of fits to the $`\nu =0`$ condensate. The fitting function is a sum of the fixed-topology expression plus a linear mass-dependent term. Curves (a) and (c) show the result of fits to the three or lowest two mass data points. Curves (b) and (d) are the part of the fit coming from the $`\mathrm{\Sigma }(V,\mu )`$ expression of Eq. 17. The three-mass fit gives $`r_0^3\mathrm{\Sigma }=0.53(4)`$ while the two-mass fit gives 0.67(8). These numbers are obviously unstable. We are currently generating data at smaller quark masses which should alleviate this problem by filling in the curve.
But the numbers we record are quite different from the RMT ones. We think (see Ref. Damgaard:2001ep for a useful discussion) that this is an artifact of the small simulation volume. The problem of principle that we face is that the formulas to which we fit $`\overline{\psi }\psi (m)`$ assume that the condensate is given by an integral over a known spectral density $`\rho ^{(\nu )}(\zeta ,\mu )`$ for all values of $`\zeta `$,
$$\frac{\mathrm{\Sigma }_\nu (\mu )}{\mathrm{\Sigma }}=2\mu _0^{\mathrm{}}𝑑\zeta \frac{\rho ^{(\nu )}(\zeta ,\mu )}{\zeta ^2+\mu ^2}.$$
(18)
We are doing simulations in small coarse lattices, and so this assumption is unwarranted. We expect that cutoff effects will alter the high eigenvalue part of the spectrum. Only a larger lattice will cure this problem.
We can check this assumption by fitting $`\overline{\psi }\psi (m)`$ in the $`|\nu |=1`$ sector. We get $`\mathrm{\Sigma }=0.70(3)`$ from a fit to the $`am=0.03`$, 0.05 and 0.1 data sets, and 0.88(10) from the $`am=0.03`$ and 0.05 sets. These are bigger discrepancies from the RMT results than the $`\nu =0`$ fits. Fig. 10 shows that the third eigenmode of the $`\nu =0`$ sector matches the RMT prediction much better than the third $`|\nu |=1`$ mode does.
## VIII Conclusion
We have presented results from simulations of two flavors of dynamical overlap fermions at sea quark masses of about 35, 55 and 100 MeV. For us, this is a second step in gaining experience with these simulations. We therefore focused on dynamical fermion effects in the results. We were able to show that the topological susceptibility is greatly reduced as compared to quenched simulations. This measurement also is a great strength of the use of overlap fermions for the sea quarks. The topological charge as defined by the index theorem has a direct impact on the update of the gauge fields during the simulation. However, we are still worried about the long auto-correlation time of the topological charge.
We also extracted the quark condensate from a comparison of the distribution of the lowest eigenvalues with random matrix theory. This method is much simpler than a direct fit to $`\overline{\psi }\psi (m)`$ and has the advantage of giving the infinite volume, zero mass condensate directly without need of extrapolation. The eigenvalue distributions fit consistently the distributions in the various topological sectors for the first three modes. However, due to our very small volume, there is a finite size correction to $`\mathrm{\Sigma }`$ which is not well under control. This correction is probably $`𝒪(40\%)`$ for our simulation. But it is expected to scale which $`1/L^2`$ and one thus only needs a moderately larger volume to make it small.
We were also able to demonstrate the effect of the dynamical quarks on the scalar meson two-point function. We showed that contrary to the quenched theory the scalar two-point function is not turned negative by the zero-modes in the sectors of non-trivial topology.
All this is possible through an efficient implementation of the overlap operator using fat links, multiple pseudo-fermion fields to decrease the auto-correlation time of the topological charge and improvements in the algorithm namely the computation of the height of the step in the fermionic action.
## Appendix A Z-factors from perturbation theory
A simple application of the techniques described in Ref. DeGrand:2002va gives us one loop predictions for the vector, axial vector, pseudo-scalar, and scalar currents for our action. They are shown in Table 4. The value of the momentum scale at which the strong coupling constant is evaluated (from the Lepage-Mackenzie conventionref:LM ) is also shown. At one-loop order there is no difference between the gauge propagator from a tadpole improved action and the tree-level one, so we use the tree-level propagator in our computation. In this order of perturbation theory, each step of stout smearing is equivalent to a step of unitarized APE-smearingref:APEblock with $`\rho =\alpha /6`$ in the terminology of Ref. Bernard:1999kc .
In practice, we define the coupling through the so-called “$`\alpha _V`$” scheme. We only know the one-loop expression relating the plaquette to the coupling; it is
$$\mathrm{ln}\frac{1}{3}\mathrm{Tr}U_p=\frac{8\pi }{3}\alpha _V(q^{})W$$
(19)
with $`W=0.366`$ and $`q^{}a=3.32`$ for the tree-level Lüscher Weisz action. In our calculation of the condensate, we take the lattice spacing from the Sommer parameter. The coupling from the plaquette is matched to its $`\overline{MS}`$ value and run to the needed value of $`q^{}`$, where we perform the match. Then the $`\overline{MS}`$ result is run to $`\mu =2`$ GeV using the usual two-loop formula. In our simulations $`\alpha _V(3.32/a)=0.192`$, 0.193, 0.193, and $`Z_s=1.19`$, 1.22 and 1.23 for the $`am_q=0.03`$, 0.05 and 0.10 data sets. Essentially all the difference from unity comes from the $`\overline{MS}`$ running from $`\mu =1/a`$ to 2 GeV. Using $`Z_m=1/Z_S`$ gives the $`\overline{MS}`$ quark masses quoted in the body of the paper.
## Acknowledgments
This work was supported by the US Department of Energy. It is a pleasure to thank A. Hasenfratz for conversations and to P. Damgaard for extensive correspondence about random matrix theory. Simulations were performed on the University of Colorado Beowulf cluster. We are grateful to D. Johnson for its construction and maintenance, and to D. Holmgren for advice about its design.
|
warning/0506/quant-ph0506001.html
|
ar5iv
|
text
|
# New binding-concealing trade-offs for quantum string commitment
## 1 Introduction
Commitment schemes are powerful cryptographic primitives. In a bit commitment scheme $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$, the committer is supposed to commit a bit $`b\{0,1\}`$ to $`\mathrm{𝖡𝗈𝖻}`$ in such a way that after the commit phase she cannot change her choice of the committed bit. This is referred to as the binding property. Also at this stage $`\mathrm{𝖡𝗈𝖻}`$ should not be able to figure out what the committed bit is. This is referred to as the concealing property. Later in the reveal phase $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ is supposed to reveal the bit $`b`$ and convince $`\mathrm{𝖡𝗈𝖻}`$ that this was indeed the bit which she committed earlier. Bit commitment schemes have been very well studied in both the classical and quantum models since existence of such schemes imply several interesting results in cryptography. It has been shown that bit commitment schemes imply existence of quantum oblivious transfer \[Yao95\] which in turn provides a way to do any two-party secure computation \[Kil88\]. They are also useful in constructing zero knowledge proofs \[Gol01\] and imply another very useful cryptographic primitive called secure coin tossing \[Blu83\]. But unfortunately strong negative results are known about them in case $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ are assumed to possess arbitrary computation power and information theoretic security is required. In this paper we are concerned with this setting of information theoretic security with unbounded computational resources with cheating parties. Classically bit commitment schemes are known to be impossible. In the quantum setting several schemes were proposed but later several impossibility results were shown \[May97, LC97, LC98, DKSW07\]. Negative results were also shown for approximate implementations of bit commitment schemes \[SR02, DKSW07\] in which trade-offs were shown for cheating probabilities of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$, referred to as binding-concealing trade-offs. Interestingly however Kent \[Ken04\] has exhibited that bit-commitment can be achieved using relativistic constraints. However we point out that in this work we do not keep considerations of relativity into picture and our setting is non-relativistic.
Now suppose instead of wanting to commit a bit $`b\{0,1\}`$, $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ wants to commit an entire string $`x\{0,1\}^n`$. One way to do this might be to commit all the bits of $`x`$ separately. Binding-concealing trade-offs of such schemes will be limited by the binding-concealing trade-offs allowable for bit commitment schemes. But it is conceivable that there might exist cleverer schemes which allow for better binding and concealing properties? This question was originally raised by Kent \[Ken03\]. Let us first begin by formally defining a quantum string commitment protocol. Our definition is similar to the one considered by Buhrman et al. \[BCH<sup>+</sup>07\]
###### Definition 1 (Quantum string commitment)
Let $`P=\{p_x:x\{0,1\}^n\}`$ be a probability distribution and let $`B`$ be a measure of information (we define several measures of information later). A $`(n,a,b)B\mathrm{𝖰𝖲𝖢}`$ protocol for $`P`$ is a quantum communication protocol \[Yao95, LC98\] between $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$. $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ gets an input $`x\{0,1\}^n`$ (chosen according to the distribution $`P`$), which is supposed to be the string to be committed. The starting joint state of the qubits of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ is some pure state. There are no intermediate measurements during the protocol and $`\mathrm{𝖡𝗈𝖻}`$ has a final checking $`\mathrm{𝖯𝖮𝖵𝖬}`$ measurement $`\{M_y|y\{0,1\}^n\}\{I_yM_y\}`$ (please see Sec. 2 for definition of $`\mathrm{𝖯𝖮𝖵𝖬}`$) to determine the value of the committed string by $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ or to detect her cheating. The protocol runs in two phases called the commit phase followed by the reveal phase. The following properties need to be satisfied.
1. (Correctness) Let $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ act honestly. Let $`\rho _x`$ be the state of $`\mathrm{𝖡𝗈𝖻}`$’s qubits at the end of the reveal phase of the protocol when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ gets input $`x`$. Then $`x,y\mathrm{𝖳𝗋}M_y\rho _x=1`$ iff $`x=y`$ and 0 otherwise.
2. (Concealing) Let $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ act honestly and $`\mathrm{𝖡𝗈𝖻}`$ be possibly cheating. Let $`\sigma _x`$ be the state of $`\mathrm{𝖡𝗈𝖻}`$’s qubits after the commit phase when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ gets input $`x`$. Then the $`B`$ information of the ensemble $`=\{p_x,\sigma _x\}`$ is at most $`b`$. In particular this is also true for both $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ acting honestly.
3. (Binding) Let $`\mathrm{𝖡𝗈𝖻}`$ act honestly and $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ be possibly cheating. Let $`c\{0,1\}^n`$ be a string in a special cheating register $`C`$ with $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ that she keeps independent of the rest of the registers till the end of the commit phase. Let $`\rho _c^{}`$ be the state of $`\mathrm{𝖡𝗈𝖻}`$’s qubits at the end of the reveal phase when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ has $`c`$ in the cheating register. Let $`\stackrel{~}{p}_c\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{𝖳𝗋}M_c\rho _c^{}`$. Then for all input strings $`x`$,
$$\underset{c\{0,1\}^n}{}p_c\stackrel{~}{p}_c2^{an}.$$
The idea behind the above definition is as follows. At the end of the reveal phase of an honest run of the protocol $`\mathrm{𝖡𝗈𝖻}`$ figures out $`x`$ from $`\rho _x`$ by performing the $`\mathrm{𝖯𝖮𝖵𝖬}`$ measurement $`\{M_x\}\{I_xM_x\}`$. He accepts the committed string to be $`x`$ iff $`M_x`$ succeeds and this happens with probability $`\mathrm{𝖳𝗋}M_x\rho _x`$. He declares $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ cheating if $`I_xM_x`$ succeeds. Thus due to the first condition, at the end of an honest run of the protocol, $`\mathrm{𝖡𝗈𝖻}`$ accepts the committed string to be exactly the input string of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ with probability 1. The second condition above takes care of the concealing property stating that the amount of $`B`$ information about $`x`$ that a possibly cheating $`\mathrm{𝖡𝗈𝖻}`$ gets is bounded by $`b`$. In bit-commitment protocols, the concealing property was quantified in terms of the probability with which $`\mathrm{𝖡𝗈𝖻}`$ can guess $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s bit. Buhrman et al. \[BCH<sup>+</sup>07\] in fact do consider $`\mathrm{𝖡𝗈𝖻}`$’s probability of guessing $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s input string as quantifying the concealing property. However in the proof of their trade-off result, they consider a related notion of information as a quantification of the concealing property. In this paper, we use various notions of information to quantify the concealing property of the protocol. The third condition guarantees the binding property. It makes sure that if a cheating $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ wants to postpone committing or wants to change her choice at the end of the commit phase, then she cannot succeed in making an honest $`\mathrm{𝖡𝗈𝖻}`$ accept her new choice with good probability, for a lot of different strings of her choice.
A few points regarding the above definition are important to note. We assume that the combined state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ at the beginning of the protocol is a pure state. Given this assumption, it can be assumed without loss of generality (due to the arguments of \[Yao95, LC98\]) that it remains a pure state till the end of the protocol (in an honest run). This is because $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ need not apply any intermediate measurements, before $`\mathrm{𝖡𝗈𝖻}`$ applies the final checking $`\mathrm{𝖯𝖮𝖵𝖬}`$ at the end of the protocol. Our impossibility result makes a critical use of this fact and fails to hold if the starting combined state is not a pure state. However, there are no restrictions on the starting pure state shared between $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$, it could even be an entangled state between them. The impossibility result in \[BCH<sup>+</sup>07\] has also been shown under this assumption. This assumption has also been made in showing impossibility results for bit-commitment schemes \[May97, LC97, LC98\]. The main reason why these arguments do not work, both for bit commitment and string commitment schemes, if the combined state is not a pure state is that the Local Transition Theorem (Thm. 2.2 mentioned later) fails to hold for mixed states. It is conceivable that, and will be interesting to see if better $`\mathrm{𝖰𝖲𝖢}`$ schemes exist when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ are forced (by some third party say) to start in some mixed state. Please look at \[DKSW07\] for extension of impossibility results for bit-commitment to a very large class of protocols.
### 1.1 Measures of information
As we will see later, the notion of information used in the above definition is very important and therefore let us briefly define various notions of information that we will be concerned with in this paper. The following notion of information, referred to as the quantum mutual information or the Holevo-$`\chi `$ information is one of the most commonly used.
###### Definition 2 (Holevo-$`\chi `$ information)
Given a quantum state $`\rho `$, the von-Neumann entropy of $`\rho `$ is defined as $`𝖲(\rho )\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{𝖳𝗋}\rho \mathrm{log}_2\rho `$. Given quantum states $`\rho ,\sigma `$, the Kullback-Leibler divergence or relative entropy between them is defined as $`𝖲(\rho \sigma )\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{𝖳𝗋}\rho (\mathrm{log}_2\rho \mathrm{log}_2\sigma )`$. Given an ensemble $`=\{p_x,\rho _x\}`$, let $`\rho \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_xp_x\rho _x`$, then its Holevo-$`\chi `$ information is defined as
$$\chi ()\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\underset{x}{}p_x(𝖲(\rho )𝖲(\rho _x))=\underset{x}{}p_x𝖲(\rho _x\rho ).$$
The following notion captures the amount of information that can be made available to the real world through measurements on the quantum encoding of a classical random variable.
###### Definition 3 (Accessible information)
Let $`=\{p_x,\rho _x\}`$ be an ensemble and let $`X`$ be a classical random variable such that $`\mathrm{Pr}(X=x)\stackrel{\mathrm{𝖽𝖾𝖿}}{=}p_x`$. Let $`Y^{}`$, correlated with $`X`$, be the classical random variable that represents the result of a $`\mathrm{𝖯𝖮𝖵𝖬}`$ measurement $``$ performed on $``$. The accessible information $`I_{\mathrm{acc}}()`$ of the ensemble $``$ is then defined to be
$$I_{\mathrm{acc}}()\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\underset{}{\mathrm{max}}I(X:Y^{}).$$
(1)
As mentioned before Buhrman et al. used $`\mathrm{𝖡𝗈𝖻}`$’s probability of guessing $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s input string as the measure of concealment of the protocol. However in the proofs of their impossibility result, they used the following notion of information.
###### Definition 4 ($`\xi `$ information \[BCH<sup>+</sup>07\])
The $`\xi `$ information of an ensemble $`=\{p_x,\rho _x\}`$ is defined as
$$\xi ()\stackrel{\mathrm{𝖽𝖾𝖿}}{=}n+\mathrm{log}_2\underset{x}{}\mathrm{𝖳𝗋}(p_x\rho ^{1/2}\rho _x)^2$$
where $`\rho =_xp_x\rho _x`$.
Let $`q_x`$ be the probability that $`\mathrm{𝖡𝗈𝖻}`$ correctly guesses $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s input string $`x`$ (with $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ honest) before the start of the reveal phase. \[BCH<sup>+</sup>07\] showed that any $`(n,a,b)\mathrm{𝖰𝖲𝖢}`$ protocol with $`_{x\{0,1\}^n}q_x2^b`$, is also a $`(n,a,b)\xi \mathrm{𝖰𝖲𝖢}`$ protocol. Hence their impossibility results for $`(n,a,b)\xi \mathrm{𝖰𝖲𝖢}`$ protocols implied same impossibility results for $`(n,a,b)\mathrm{𝖰𝖲𝖢}`$ protocols with $`_{x\{0,1\}^n}q_x2^b`$.
In this paper we also consider a notion of divergence information. It is based on the following notion of distance between two quantum states, considered by Jain, Radhakrishnan and Sen \[JRS02\].
###### Definition 5 (Observational divergence \[JRS02\])
Let $`\rho ,\sigma `$ be two quantum states. The observational divergence between them denoted $`𝖣(\rho \sigma )`$, is defined as,
$$𝖣(\rho \sigma )\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\underset{𝖬:\mathrm{𝖯𝖮𝖵𝖬}\mathrm{𝖾𝗅𝖾𝗆𝖾𝗇𝗍}}{\mathrm{max}}\mathrm{𝖳𝗋}M\rho \mathrm{log}_2\frac{\mathrm{𝖳𝗋}M\rho }{\mathrm{𝖳𝗋}M\sigma }.$$
The definition of divergence information of an ensemble is similar to the Holevo-$`\chi `$ information except the notion of distance between quantum states used is now observational divergence instead of relative entropy.
###### Definition 6 (Divergence information)
Let $`=\{p_x,\rho _x\}`$ be an ensemble and let $`\rho \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_xp_x\rho _x`$. Its divergence information is defined
$$𝒟()\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\underset{x}{}p_x𝖣(\rho _x\rho ).$$
### 1.2 Previous results
The impossibility of a strong string commitment protocol, in which both $`a,b`$ are required to be 0, is immediately implied by the impossibility of strong bit-commitment protocols. The question of a trade-off between $`a`$ and $`b`$ was studied by Buhrman et al. They studied this trade-off both in the scenario of single execution of the protocol and also in the asymptotic regime with several parallel executions of the protocol. In the scenario of single execution of the protocol they showed the following result.
###### Theorem 1.1 (\[BCH<sup>+</sup>07\])
For single execution of the protocol of a $`(n,a,b)`$-$`\xi `$-$`\mathrm{𝖰𝖲𝖢}`$, $`a+b+5\mathrm{log}_254n`$.
This then (as argued before) implied similar trade-off for a $`(n,a,b)`$-$`\mathrm{𝖰𝖲𝖢}`$ with $`_{x\{0,1\}^n}q_x2^b`$ (where $`q_x`$ be the probability that $`\mathrm{𝖡𝗈𝖻}`$ correctly guesses $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s input string $`x`$, with $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ honest, before the start of the reveal phase.) In the asymptotic regime they showed the following result in terms of the Holevo-$`\chi `$ information.
###### Theorem 1.2 (\[BCH<sup>+</sup>07\])
Let $`\mathrm{\Pi }`$ be a $`(n,,b)\chi \mathrm{𝖰𝖲𝖢}`$ scheme. Let $`\mathrm{\Pi }_m`$ represent $`m`$ parallel executions of $`\mathrm{\Pi }`$. Let $`a_m`$ represent the binding parameter of $`\mathrm{\Pi }_m`$ and let $`a\stackrel{\mathrm{𝖽𝖾𝖿}}{=}lim_m\mathrm{}\frac{a_m}{m}`$. Then, $`a+bn`$.
There are two reason why Thm. 1.2 may appear stronger than Thm. 1.1. One because there is no additive constant and the other because for many ensembles $``$, $`\chi ()\xi ()`$ as we show in Sec. 0.A. In fact, as we also show in Sec. 0.A, there exists ensembles $``$ for which $`\xi ()`$ is exponentially (in $`n`$) larger than $`\chi ()`$.
Along with these impossibility results Buhrman et al. interestingly also showed that if the measure of information considered is the accessible information, the above trade-offs no longer hold. For example there exists a $`\mathrm{𝖰𝖲𝖢}`$ scheme where $`a=4\mathrm{log}_2n+O(1)`$ and $`b=4`$ when measure of information is the accessible information. This therefore asserts that the choice of measure of information is crucial to (im)possibility. Previously Kent \[Ken03\] also exhibited trade-offs for some schemes on $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s probability of cheating and the amount of accessible information that $`\mathrm{𝖡𝗈𝖻}`$ gets about the committed string. However he did not allow $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ to be arbitrarily cheating, in particular $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ could not have started with a superposition of strings in the input register. Therefore the schemes that he considered were truly not $`\mathrm{𝖰𝖲𝖢}`$s as we have defined them.
### 1.3 Our results
We show the following binding-concealing trade-off for $`\mathrm{𝖰𝖲𝖢}`$s.
###### Theorem 1.3
For single execution of the protocol of a $`(n,a,b)𝒟\mathrm{𝖰𝖲𝖢}`$ scheme,
$$a+b+8\sqrt{b+1}+16n.$$
It was shown by Jain, Radhakrishnan and Sen \[JRS02\] that for any two states $`\rho ,\sigma `$, $`𝖣(\rho \sigma )𝖲(\rho \sigma )+1`$, which implies from Defn. 2 and 6 that for any ensemble $`,𝒟()\chi ()+1`$. This immediately gives us the following impossibility result in terms of Holevo-$`\chi `$ information.
###### Theorem 1.4
For single execution of the protocol of a $`(n,a,b)\chi \mathrm{𝖰𝖲𝖢}`$ scheme
$$a+b+8\sqrt{b+2}+17n.$$
We also consider the notion of maximum possible divergence information (similar to the notion of maximum possible Holevo-$`\chi `$ information considered by Jain \[Jai06\]) of an encoding $`E:x\rho _x`$. For a probability distribution $`\mu \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\{p_x\}`$ over $`\{0,1\}^n`$, let the ensemble $`_\mu (E)\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\{p_x,\rho _x\}`$. Let $`\rho _\mu \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_xp_x\rho _x`$.
###### Definition 7
(Maximum possible divergence information) Maximum possible divergence information of an encoding $`E:x\rho _x`$ is defined as $`\stackrel{~}{𝒟}(E)\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{max}_\mu 𝒟(_\mu (E))`$.
We show the following theorem which states that if the maximum possible divergence information in the qubits of $`\mathrm{𝖡𝗈𝖻}`$ at the end of the commit phase is small then $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ can actually cheat with good probability for any string $`x\{0,1\}^n`$ and not just on the average.
###### Theorem 1.5
For a $`\mathrm{𝖰𝖲𝖢}`$ scheme let $`\sigma _x`$ be as in Defn. 1 when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ act honestly in the commit phase. If for the encoding $`E:x\sigma _x,\stackrel{~}{𝒟}(E)b`$ then for all strings $`c\{0,1\}^n`$,
$$\stackrel{~}{p}_c2^{(b+8\sqrt{b+1}+16)},$$
where $`\stackrel{~}{p}_c`$ (as in Defn. 1) represents the probability of successfully revealing string $`c`$ (in the cheating string) by cheating $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$.
Again using the fact that for all ensembles $`𝖣(\rho \sigma )𝖲(\rho \sigma )+1`$ we immediately get the following theorem in terms of maximum possible Holevo-$`\chi `$ information $`\stackrel{~}{\chi }(E)`$ (which is similar to maximum possible divergence information and obtained by just replacing divergence with relative entropy.)
###### Theorem 1.6
For a $`\mathrm{𝖰𝖲𝖢}`$ scheme let $`\sigma _x`$ be as in Defn. 1 when $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ act honestly in the commit phase. If for the encoding $`E:x\sigma _x,\stackrel{~}{\chi }(E)b`$ then for all strings $`c\{0,1\}^n`$,
$$\stackrel{~}{p}_c2^{(b+8\sqrt{b+2}+17)},$$
where $`\stackrel{~}{p}_c`$ (as in Defn. 1) represents the probability of successfully revealing string $`c`$ (in the cheating string) by cheating $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$.
Now let us now discuss some aspects of our results.
1. In Thm. 1.4 the trade-off between $`a`$ and $`b`$ is similar (up to lower order terms of $`b`$) to the one shown by Buhrman et al. \[BCH<sup>+</sup>07\] as in Thm. 1.1. However the fact that $`b`$ in Thm. 1.4 represents the Holevo-$`\chi `$ information instead of the $`\xi `$-information (as in Thm. 1.1) makes it significantly stronger in certain cases as follows. We show in Sec. 0.A that for any ensemble $`\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\{2^n,\rho _x\}`$, where for all $`x`$, $`\rho _x`$ commutes with $`\rho \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_x2^n\rho _x`$, we have, $`\xi ()\chi ()`$. In fact, as we also show in Sec. 0.A, there exists ensembles $``$ for which $`\xi ()`$ is exponentially (in $`n`$) larger than $`\chi ()`$. Thm. 1.4 therefore becomes much stronger than Thm. 1.1 for ensembles where $`\xi ()\chi ()`$.
2. As mentioned before, Jain, Radhakrishnan and Sen \[JRS02\] have shown that for any ensemble $`,𝒟()\chi ()+1`$. However recently, Jain, Nayak and Su \[JNS08\] have shown that there exists ensembles $``$ such that $`\chi ()𝒟()`$ ($`\chi ()=\mathrm{\Omega }(\mathrm{log}_2n𝒟())`$ for some ensembles $``$ supported on $`\{0,1\}^n`$). For ensembles where this holds, Thm. 1.3 becomes much stronger than Thm. 1.4.
3. As we show in Sec. 3, our one shot result Thm. 1.4 immediately implies the asymptotic result Thm. 1.2 of Buhrman et al.
4. No counterparts of Thm. 1.5 and Thm. 1.6 were shown by Buhrman et al. and are therefore completely new.
5. If $`b`$ is large then the cheating attack (that we present) of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ would succeed with low probability (like $`2^b`$). However, as we show in a remark in Sec. 3, in case $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s cheating attack succeeds with low probability, she would still be able to ’reverse’ her cheating operations and reveal, with a high probability, at least some $`x^{}\{0,1\}^n`$ to $`\mathrm{𝖡𝗈𝖻}`$. That is, with a high probability, $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ will be able to prevent herself from being detected cheating by $`\mathrm{𝖡𝗈𝖻}`$.
6. It is easily seen that up to lower order terms in $`b`$, the above trade-offs are achieved by trivial protocols. For Thm. 1.3 above consider the following protocol. $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ in the concealing phase sends the first $`b`$ bits of the $`n`$-bit string $`x`$. In this case $`\mathrm{𝖡𝗈𝖻}`$ gets to know $`b`$ bits of divergence information about $`x`$. In the reveal phase a cheating $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ can now reveal any of the $`2^{nb}`$ strings $`x`$ (consistent with the first $`b`$ bits being the ones sent) with probability 1. Hence $`a=\mathrm{log}_22^{nb}=nb`$. For Thm. 1.5 above let $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ send one of the $`2^b`$ strings $`s\{0,1\}^b`$ uniformly to $`\mathrm{𝖡𝗈𝖻}`$ representing the first $`b`$ bits of $`x`$. The condition of Thm. 1.5 is satisfied. Now if in the reveal phase she wants to commit any $`x`$, she can do so with probability $`2^b`$ (in the event that the sent $`s`$ is consistent with $`x`$).
In the next section we state some quantum information theoretic facts that will be useful in the proofs of the impossibility results that we present in Sec. 3.
## 2 Preliminaries
All logarithms in this paper are taken with base 2 unless otherwise specified. Let $`,𝒦`$ be finite dimensional Hilbert spaces. For a linear operator $`A`$ let $`|A|=\sqrt{A^{}A}`$ and let $`\mathrm{𝖳𝗋}A`$ denote the trace of $`A`$. Given a state $`\rho `$ and a pure state $`|\varphi 𝒦`$, we call $`|\varphi `$ a purification of $`\rho `$ iff $`\mathrm{𝖳𝗋}_𝒦|\varphi \varphi |=\rho `$. A positive operator-valued measurement $`(\mathrm{𝖯𝖮𝖵𝖬})`$ element $`M`$ is a positive semi-definite operator such that $`IM`$ is also positive semi-definite, where $`I`$ is the identity operator. A $`\mathrm{𝖯𝖮𝖵𝖬}`$ is defined as follows.
###### Definition 8 ($`\mathrm{𝖯𝖮𝖵𝖬}`$)
An $`m`$ valued $`\mathrm{𝖯𝖮𝖵𝖬}`$ measurement $``$ on a Hilbert space $``$ is a set of operators $`\{M_i,i[m]\}`$ on $``$ such that $`i,M_i`$ is positive semi-definite and $`_{i[m]}M_i=I`$ where $`I`$ is the identity operator on $``$. A classical random variable $`Y^{}`$ representing the result of the measurement $``$ on a state $`\rho `$ is an $`m`$ valued random variable such that $`i[m],\mathrm{Pr}[Y^{}=i]\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{𝖳𝗋}M_i\rho `$.
Following fact follows easily from definition of von-Neumann entropy.
###### Lemma 1
Let $`\rho _1,\rho _2`$ be quantum states. Then $`𝖲(\rho _1\rho _2)=𝖲(\rho _1)+𝖲(\rho _2)`$.
We make a central use the following information-theoretic result called the substate theorem due to Jain, Radhakrishnan, and Sen \[JRS02\].
###### Theorem 2.1 (Substate theorem, \[JRS02\])
Let $`,𝒦`$ be two finite dimensional Hilbert spaces and $`dim(𝒦)dim()`$. Let $`^2`$ denote the two dimensional complex Hilbert space. Let $`\sigma ,\tau `$ be density matrices in $``$ such that $`𝖣(\sigma \tau )<\mathrm{}`$. Let $`|\overline{\sigma }`$ be a purification of $`\sigma `$ in $`𝒦`$. Then, for $`r>1`$, there exist pure states $`|\varphi ,|\theta 𝒦`$ and $`|\overline{\tau }𝒦^2`$, depending on $`r`$, such that $`|\overline{\tau }`$ is a purification of $`\tau `$ and $`\mathrm{𝖳𝗋}||\overline{\sigma }\overline{\sigma }||\varphi \varphi ||\frac{2}{\sqrt{r}}`$, where
$$|\overline{\tau }\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\sqrt{\frac{r1}{r2^{rk}}}|\varphi |1+\sqrt{1\frac{r1}{r2^{rk}}}|\theta |0$$
and $`k\stackrel{\mathrm{𝖽𝖾𝖿}}{=}𝖣(\sigma \tau )+6\sqrt{𝖣(\sigma \tau )+1}+4`$.
#### Remarks:
1. In the above theorem if the last qubit in $`|\overline{\tau }`$ is measured in the computational basis, then probability of obtaining 1 is $`(11/r)2^{rk}`$.
2. Later in a proof below we will let $`\sigma \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\rho _c`$ , $`\tau \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\rho _B`$ and $`|\overline{\sigma }\stackrel{\mathrm{𝖽𝖾𝖿}}{=}|\varphi _c`$ which will be explained later.
Following theorem is implicit in \[HJW93, May97, LC97, LC98\] although not called explicitly by the same name.
###### Theorem 2.2 (Local transition theorem)
Let $`\rho `$ be a quantum state in $`𝒦`$. Let $`|\varphi _1`$ and $`|\varphi _2`$ be two purification of $`\rho `$ in $`𝒦`$. Then there is a local unitary transformation $`U`$ acting on $``$ such that $`(UI)|\varphi _1=|\varphi _2`$.
We would also need the following theorem which follows from arguments similar to the one in Jain \[Jai06\] for a similar theorem about relative entropy.
###### Theorem 2.3
Let $`X`$ be a finite set. Let $`E:x\rho _x`$ be an encoding. Let $`\stackrel{~}{𝒟}(E)b`$, then there exists a distribution $`\mu \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\{q_x\}`$ on $`X`$ such that
$$xX,𝖣(\rho _x\rho )b,$$
where $`\rho \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_xq_x\rho _x`$.
The following theorem is shown by Helstrom \[Hel67\].
###### Theorem 2.4
Given two quantum states $`\rho `$ and $`\sigma `$, the probability of identifying the correct state is at most $`\frac{1}{2}+\frac{\mathrm{𝖳𝗋}|\rho \sigma |}{4}`$, or in other words the probability of distinguishing them is at most $`\frac{\mathrm{𝖳𝗋}|\rho \sigma |}{2}`$.
## 3 Proofs of impossibility
Proof of Thm. 1.3: Let us consider a $`\mathrm{𝖰𝖲𝖢}`$ scheme and let $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ get input $`x`$. After an honest run of the commit phase, let $`|\varphi _x`$ be the combined state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ and $`\rho _x`$ be the state of $`\mathrm{𝖡𝗈𝖻}`$’s qubits. Let $`=\{p_x,\rho _x\}`$. From the concealing property of the $`\mathrm{𝖰𝖲𝖢}`$ it follows $`𝖣()b`$. Let $`c`$ be the string in the cheating register $`C`$ of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$. Consider a cheating run of the protocol by $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ in which she starts with the superposition $`_x\sqrt{p_x}|x`$ in the input register and proceeds with the rest of the commit phase as before in the honest protocol. Let $`\mathrm{𝖡𝗈𝖻}`$ be honest all throughout our arguments. Since the input is classical and $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ can make its copy we can assume without loss of generality that the operations of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ in the honest run are such that they do not disturb the input register. Let $`|\psi `$ be the combined state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ in this cheating run at the end of the commit phase. Let $`A,B`$ correspond to $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$’s systems respectively. Now it can be seen that in the cheating run, at the end of the commit phase the qubits of $`\mathrm{𝖡𝗈𝖻}`$ are in the state $`\rho _B\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\mathrm{𝖳𝗋}_A|\psi \psi |=_xp_x\rho _x`$. Let $`r>1`$ to be chosen later. Let us now invoke substate theorem (Thm. 2.1) by putting $`\sigma \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\rho _c,|\overline{\sigma }\stackrel{\mathrm{𝖽𝖾𝖿}}{=}|\varphi _c`$, $`\tau \stackrel{\mathrm{𝖽𝖾𝖿}}{=}\rho _B`$ and $`r\stackrel{\mathrm{𝖽𝖾𝖿}}{=}r`$. Let $`|\psi _c\stackrel{\mathrm{𝖽𝖾𝖿}}{=}|\overline{\tau }`$ be obtained from Thm. 2.1 such that the extra single qubit register $`^2`$ is also with $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$. Since $`\mathrm{𝖳𝗋}_A|\psi _c\psi _c|=\mathrm{𝖳𝗋}_A|\psi \psi |=\rho _B`$, from Local transition theorem (Thm. 2.2) there exists a unitary transformation $`A_c`$ acting just on $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$’s system $`A`$ such that $`(A_cI_B)|\psi =|\psi _c`$, where $`I_B`$ is the identity transformation on $`\mathrm{𝖡𝗈𝖻}`$’s system. Now the cheating $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ (who’s intention is to reveal string $`c`$), applies the transformation $`A_c`$ to $`|\psi `$ and then continues with the rest of the reveal phase as in the honest run. Let $`|\varphi _c^{}\stackrel{\mathrm{𝖽𝖾𝖿}}{=}|\varphi `$ be obtained from Thm. 2.1 and hence, $`\mathrm{𝖳𝗋}||\varphi _c\varphi _c||\varphi _c^{}\varphi _c^{}||2/\sqrt{r}`$. Now it can be seen that when $`\mathrm{𝖡𝗈𝖻}`$ makes the final checking $`\mathrm{𝖯𝖮𝖵𝖬}`$, the probability of success $`\stackrel{~}{p}_c`$ for $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ is at least $`(11/r)2^{rk_c}(11/\sqrt{r})`$ where $`k_c=𝖣(\rho _c\rho _B)+6\sqrt{𝖣(\rho _c\rho _B)+1}+4`$. One way to see this is to imagine that $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ first measures the single qubit register $`^2`$ and then proceeds with the rest of the reveal phase. Now imagine that she obtains one on this measurement which from Thm. 2.1 has probability $`(11/r)2^{rk_c}`$. Also once she obtains one, the combined joint state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ is $`|\varphi _c^{}`$ whose trace distance with $`|\varphi _c`$ is at most $`2/\sqrt{r}`$. Since trace distance is preserved by unitary operations and is only smaller for subsystems and since after this $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ follows the rest of the reveal phase honestly, we can conclude the following: the final state resulting with $`\mathrm{𝖡𝗈𝖻}`$ will have trace distance at most $`2/\sqrt{r}`$ with the state with him at the end of a completely honest run of the protocol in which Alice starts with $`c`$ in the input register. Hence it follows from Thm. 2.4 that $`\mathrm{𝖡𝗈𝖻}`$ will accept at the end with probability at least $`11/\sqrt{r}`$ since he was accepting with probability 1 in the complete honest run of the protocol . Hence the overall cheating probability $`\stackrel{~}{p}_c`$ of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ is at least $`(11/r)2^{rk_c}(11/\sqrt{r})`$.
Although here we have imagined $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ doing an intermediate measurement on the single qubit register $`^2`$, it is not necessary and she will have the same cheating probability when she proceeds with the rest of the honest protocol after just applying the cheating transformation $`A_c`$ since the final qubits of $`\mathrm{𝖡𝗈𝖻}`$ will be in the same state in either case. Now,
$`2^{an}`$ $``$ $`{\displaystyle \underset{c}{}}p_c\stackrel{~}{p}_c`$
$``$ $`(11/r)(11/\sqrt{r})\left({\displaystyle \underset{c}{}}p_c2^{r(𝖣(\rho _c\rho _B)+6\sqrt{𝖣(\rho _c\rho _B)+1}+4)}\right)`$
$``$ $`(11/r)(11/\sqrt{r})2^{_crp_c(𝖣(\rho _c\rho _B)+6\sqrt{𝖣(\rho _c\rho _B)+1}+4)}`$
$``$ $`(11/r)(11/\sqrt{r})2^{r(b+6\sqrt{b+1}+4)}`$
The first inequality comes from definition of $`a`$ in Defn. 1. The third inequality comes from the convexity of the exponential function and the fourth inequality comes from definition of $`b`$ in Defn. 1, Defn. 6 and concavity of the square root function.
Now when $`b>15`$, we let $`r=1+\frac{1}{b}`$ and therefore,
$`(11/r)(11/\sqrt{r})2^{r(b+6\sqrt{b+1}+4)}`$ $``$ $`{\displaystyle \frac{0.5}{(b+1)^2}}2^{(b+6\sqrt{b+1}+7)}`$
$``$ $`2^{(b+8\sqrt{b+1}+8)}`$
When $`b15`$, we let $`r=1+1/15`$ and therefore,
$`(11/r)(11/\sqrt{r})2^{r(b+6\sqrt{b+1}+4)}`$ $``$ $`2^{(b+6\sqrt{b+1}+16)}`$
Therefore we get always, $`2^{an}2^{(b+8\sqrt{b+1}+16)}`$ which finally implies,
$$a+b+8\sqrt{b+1}+16n.$$
Proof of Thm. 1.2: Let $`b_m`$ represent the concealing parameter for $`\mathrm{\Pi }_m`$. It is easy to verify from Lem. 1 and definition of Holevo-$`\chi `$ information, Defn. 2, that $`b=b_m/m`$. Then Thm. 1.4 when applied to $`\mathrm{\Pi }_m`$ implies,
$``$ $`a_m+b_m+8\sqrt{b_m+2}+17mn`$
$``$ $`\underset{m\mathrm{}}{lim}{\displaystyle \frac{1}{m}}(a_m+b_m+8\sqrt{b_m+2}+17)n`$
$``$ $`a+bn`$
Proof of Thm. 1.5: Let $`\mu =\{\lambda _x\}`$ be the distribution on $`\{0,1\}^n`$ obtained from Thm. 2.3. Consider a cheating strategy of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ in which she puts the superposition $`_x\sqrt{\lambda _x}|x`$ in the register where she keeps the commit string. Let $`c`$ be the string in the cheating register of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$. Now by arguments as above probability of success $`\stackrel{~}{p}_c`$ for $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ is at least $`(11/\sqrt{r})(11/r)2^{rk_c}`$ where $`k_c,\rho _c,\rho `$ being as before. Since for all $`c,𝖣(\rho _c\rho )b`$ it implies (by setting $`r`$ appropriately) $`c,\stackrel{~}{p}_c2^{(b+8\sqrt{b+1}+16)}.`$
#### Remark:
Let us now see how, with a good probability overall, $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ will be able to prevent herself from being detected cheating by $`\mathrm{𝖡𝗈𝖻}`$. Let $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ have $`c`$ in the cheating register. Let $`r_c`$ be the probability of getting one on performing the two outcome measurement (obtained from Thm. 2.1) after the commit phase as in the cheating strategy described above in proof of Thm. 1.3. In case she gets one, she proceeds with the cheating strategy. In case she gets zero, she tries to rollback so that she can successfully reveal at least some string to $`\mathrm{𝖡𝗈𝖻}`$. For this she does the following.
1. She applies the transformation $`A_c^{}`$ (that is inverse of $`A_c`$).
2. She measures the input register in the computational basis and say she obtains $`x^{}`$.
3. She proceeds with the rest of the reveal phase as if her actual input was $`x^{}`$.
Assume that $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ obtains zero on performing the two-outcome measurement as in the cheating strategy described above which happens with probability $`1r_c`$. Now it can be verified that the trace distance between $`|\psi _c\psi _c|`$ and the combined state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ after obtaining zero on performing the measurement is at most $`2r_c`$. Since, $`A_c^{}`$ is unitary, this implies that the combined state of $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ and $`\mathrm{𝖡𝗈𝖻}`$ after applying $`A_c^{}`$, and $`|\psi \psi |`$ will be at most $`2r_c`$. Now we can argue as before that $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ can reveal some string successfully to $`\mathrm{𝖡𝗈𝖻}`$ with probability at least $`1r_c`$. Therefore overall, the probability that $`\mathrm{𝖠𝗅𝗂𝖼𝖾}`$ will be able to reveal some string is at least $`r_c+(1r_c)^21r_c`$. Now since typically $`r_c`$ is quite small (like $`2^b`$), $`1r_c`$ is quite close to 1.
### Acknowledgment
We thank Harry Buhrman, Matthias Christandl, Hoi-Kwong Lo, Jaikumar Radhakrishnan, and Pranab Sen for discussions. We also thank anonymous referees for suggestions on an earlier draft.
## Appendix 0.A Separations for $`\xi ()`$ and $`\chi ()`$
Let $`\stackrel{\mathrm{𝖽𝖾𝖿}}{=}\{1/2^n,\rho _x\}`$ be an ensemble with $`x\{0,1\}^n`$. Let $`\rho \stackrel{\mathrm{𝖽𝖾𝖿}}{=}_x2^n\rho _x`$. Lets assume that for all $`x`$, $`\rho _x`$ commutes with $`\rho `$ as is the case in classical ensembles. We show that in this case $`\xi ()\chi ()`$. Consider,
$`\xi ()`$ $`=`$ $`n+\mathrm{log}{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}(2^n\rho ^{1/2}\rho _x)^2`$
$`=`$ $`\mathrm{log}{\displaystyle \underset{x}{}}2^n\mathrm{𝖳𝗋}(\rho ^{1/2}\rho _x)^2`$
$``$ $`2^n{\displaystyle \underset{x}{}}\mathrm{log}\mathrm{𝖳𝗋}(\rho ^{1/2}\rho _x)^2\text{ (from concavity of }\mathrm{log}\text{ function)}`$
$`=`$ $`2^n{\displaystyle \underset{x}{}}\mathrm{log}\mathrm{𝖳𝗋}(\rho _x\rho ^1\rho _x)\text{ (since }\rho _x,\rho \text{ commute)}`$
$``$ $`2^n{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}\rho _x\mathrm{log}(\rho _x\rho ^1)\text{ (since }\mathrm{log}\mathrm{𝖳𝗋}BA\mathrm{𝖳𝗋}A\mathrm{log}B\text{, for }A,B\text{ quantum states)}`$
$`=`$ $`2^n{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}\rho _x(\mathrm{log}\rho _x\mathrm{log}\rho )\text{ (since }\rho _x,\rho \text{ commute)}`$
$`=`$ $`\chi ()`$
Next we show that there exists classical ensembles for which $`\xi ()`$ could be exponentially larger than $`\chi ()`$. Consider the ensemble of classical distributions $`\{2^n,P_x\}`$ for $`x\{0,1\}^n`$. Here each $`P_x`$ has support on $`\{0,1\}^n`$. Let $`ϵ(0,1)`$ be a constant. Let $`P_x(x)=2^{\frac{ϵn}{2}}`$ and let the other values for $`P_x(y),yx`$ be the same. Let $`P\stackrel{\mathrm{𝖽𝖾𝖿}}{=}_x2^nP_x`$. It is easy to verify that in this case $`P`$ is the uniform distribution on $`\{0,1\}^n`$. Now,
$`\xi ()`$ $`=`$ $`n+\mathrm{log}{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}(2^{2n}P^1P_x^2)`$
$`=`$ $`n+\mathrm{log}{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}(P^1P_x^2)`$
$``$ $`n+\mathrm{log}{\displaystyle \underset{x}{}}2^{n(1ϵ)}\text{ (since for all }x\text{}\mathrm{𝖳𝗋}P^1P_x^22^{n(1ϵ)}\text{ and since }\mathrm{log}\text{ is monotonic)}`$
$`=`$ $`n+\mathrm{log}2^{n(2ϵ)}`$
$`=`$ $`n(1ϵ)`$
Also we note that for all $`x`$, $`\mathrm{𝖳𝗋}P_x(\mathrm{log}P_x\mathrm{log}P)2^{\frac{ϵn}{2}}n(1ϵ/2)`$ and hence,
$`\chi ()`$ $`=`$ $`2^n{\displaystyle \underset{x}{}}\mathrm{𝖳𝗋}P_x(\mathrm{log}P_x\mathrm{log}P)`$
$``$ $`2^n{\displaystyle \underset{x}{}}2^{\frac{ϵn}{2}}n(1ϵ/2)`$
$`=`$ $`2^{\frac{ϵn}{2}}n(1ϵ/2)`$
Therefore by letting $`ϵ`$ to be a constant very close to $`0`$, we can let $`\xi ()`$ to be very close to $`n`$ whereas $`\chi ()`$ would still be exponentially small in $`n`$.
|
warning/0506/hep-ex0506003.html
|
ar5iv
|
text
|
# Nuclear effects in high-𝒑_𝑇 production of direct photons and neutral mesons
## I Introduction
The study of inclusive particle production at large transverse momenta ($`p_T`$) has yielded valuable information about perturbative quantum chromodynamics (PQCD), parton distribution functions (PDF), and fragmentation functions of partons Geist et al. (1990); McCubbin (1981); Owens (1987); Ferbel and Molzon (1984); Apanasevich et al. (1999). The use of nuclear targets provides, in addition, information on parton and hadron rescattering and explores the time evolution of the collision. Since the discovery of the nuclear enhancement of high-$`p_T`$ single-particle production Cronin et al. (1975); Antreasyan et al. (1979); Frisch et al. (1983), a large body of data has been accumulated to investigate nuclear-target effects in a wide variety of production processes, including those yielding single hadrons, dihadron pairs, Drell-Yan pairs, two-jet systems, and heavy flavors. Recent results from the RHIC program sta ; phe ; pho ; bra , in particular, have highlighted the differences between initial and final-state effects in the nuclear environment. Many approaches have been developed to explain these data, which have included models for multiple-scattering, Fermi motion, modification of parton densities in the nuclear medium, QCD higher-twist contributions, and new states of matter.
We present the results of a high-statistics study of nuclear effects in the inclusive production of direct photons, $`\pi ^0`$ and $`\eta `$ mesons at large $`p_T`$ using data from Fermilab experiment E706, and compare the results to predictions of a phenomenological model of nuclear effects wan .
## II Apparatus
### II.1 Meson West spectrometer
Fermilab E706 was a fixed-target experiment designed to measure the production of direct photons, neutral mesons, and associated particles at high-$`p_T`$ Apanasevich et al. (1998a); E70 (a); Apanasevich et al. (2004a, b); Apanasevich et al. (1997); E70 (b). The apparatus included a charged particle spectrometer and a large liquid argon calorimeter, as described below. Additional information about the Meson West spectrometer can be found in earlier papers Apanasevich et al. (1997); Apanasevich et al. (1998b).
This paper reports on data from the two primary data runs of the experiment. During the 1990 run, the target consisted of two 0.8 mm thick copper foils followed by two pieces of beryllium (Fig. 1:top). The upstream piece of beryllium was 3.7 cm long, while the length of the downstream piece was 1.1 cm. In the 1991–1992 run, the target consisted of two 0.8 mm thick copper foils immediately upstream of a liquid hydrogen target Allspach et al. (1991), followed by a 2.54 cm long beryllium cylinder (Fig. 1:bottom). The liquid hydrogen was contained in a 15.3 cm long mylar flask, which was supported in an evacuated volume with beryllium windows at each end (2.5 mm thickness upstream and 2.8 mm thickness downstream). The target material is detailed in Table 1.
The charged particle spectrometer consisted of silicon microstrip detectors (SSDs) in the target region and multiwire proportional chambers (PWCs) and straw tube drift chambers (STDCs) downstream of a large-aperture analysis magnet Apanasevich et al. (1997). Six 3$`\times `$3 cm<sup>2</sup> SSD planes were located upstream of the target region and used to reconstruct beam tracks. Two hybrid 5$`\times `$5 cm<sup>2</sup> SSD planes (25 $`\mu `$m pitch strips in the central 1 cm, 50 $`\mu `$m beyond) were located downstream of the target region. These were followed by eight 5$`\times `$5 cm<sup>2</sup> SSD planes of 50 $`\mu `$m pitch. The analysis dipole magnet imparted a $`0.45\mathrm{GeV}/c`$ $`p_T`$ impulse in the horizontal plane to charged particles. Downstream track segments were measured by means of four stations of four views ($`XYUV`$) of 2.54 mm pitch PWCs and two stations of eight (4$`X`$4$`Y`$) layers of STDCs with tube diameters 1.03 cm (upstream station) and 1.59 cm (downstream station) Bromberg et al. (1991).
Photons were detected in a large, lead and liquid-argon sampling electromagnetic calorimeter (EMLAC), located 9 m downstream of the target Apanasevich et al. (1998b). The EMLAC had a cylindrical geometry with an inner radius of 20 cm and an outer radius of 160 cm. The calorimeter had 33 longitudinal cells read out in two sections: an 11 cell front section (8.5 radiation lengths) and a 22 cell back section (18 radiation lengths). Each longitudinal cell consisted of a 2 mm thick lead cathode (the first cathode was constructed of aluminum), a double-sided copper-clad G-10 radial ($`R`$) anode board, a second 2 mm thick lead cathode, and a double-sided copper-clad G-10 azimuthal ($`\mathrm{\Phi }`$) anode board. The 2.5 mm gaps between these layers were filled with liquid argon. The physical layout is illustrated in Fig. 2.
The EMLAC readout was subdivided azimuthally into octants, each consisting of interleaved, finely segmented, radial and azimuthal views. This segmentation was realized by cutting the copper-cladding on the anode boards to form either radial or azimuthal strips. Signals from corresponding strips from all $`R`$ (or $`\mathrm{\Phi }`$) anode boards in the front (or back) section of a given octant were jumpered together. The copper-cladding on the radial anode boards was cut into concentric strips centered on the nominal beam axis. The width of the strips on the first $`R`$ board was 5.5 mm. The width of the strips on the following $`R`$ boards increased slightly so that the radial geometry was projective relative to the target region. The azimuthal strips were split at a radius of 40 cm into inner and outer segments; each inner strip subtended an azimuthal angle of $`\pi /192`$ radians, while outer strips covered $`\pi /384`$ radians.
The spectrometer was located at the end of the Meson West beamline. The design of the beamline, primary target, and primary beam dump were intended to minimize the rate of beam-halo muons incident upon the spectrometer. The beamline was capable of transporting either a primary (800 GeV/$`c`$) proton beam or unseparated secondary particle beams of either polarity to the experimental hall. The beamline Čerenkov detector was used to identify the secondary beam particles Striley (1996). This 43.4 m long helium-filled counter was located 100 m upstream of the experimental target. The positive secondary beam with mean momentum of 530 GeV/$`c`$ was 97% protons. The negative secondary beam with mean momentum of 515 GeV/$`c`$ was 99% pions.
At the end of the beamline was a 4.7 m long stack of steel surrounding the beam pipe and shadowing the EMLAC to absorb off-axis hadrons. A water tank was placed at the downstream end of this hadron shield to absorb low-energy neutrons. Surrounding the hadron shield and neutron absorber were walls of scintillation counters (VW) to identify penetrating muons. There was one wall at the upstream end and two walls at the downstream end of the hadron absorber during the 1990 run. An additional wall was added to the upstream end of the hadron absorber prior to the 1991–1992 run.
### II.2 Trigger
The E706 trigger selected interactions yielding high-$`p_T`$ showers in the EMLAC. The selection process involved several stages: beam and interaction definitions, a pretrigger, and high-$`p_T`$ trigger requirements E70 (c); Apanasevich et al. (1997); E70 (a). A scintillator hodoscope, located 2 m upstream of the target region, was used to detect beam particles, and reject interactions with more than one spacially isolated incident particle. Additional scintillator with a 1 cm diameter central hole was located just downstream of the beam hodoscope, and served to reject interactions initiated by particles in the beam halo BH . Two pairs of scintillator counters, mounted on the dipole analysis magnet, were used to identify interactions in the target. To minimize potential confusion in the EMLAC due to out-of-time interactions, a filter was employed to reject interactions that occurred within 60 ns of one another.
For those interactions that satisfied the beam and interaction requirements, the $`p_T`$ deposited in various regions of the EMLAC was evaluated by weighting the energy signals from the EMLAC $`R`$-channel amplifier fast outputs by a factor proportional to $`\mathrm{sin}\theta _i`$, where $`\theta _i`$ was the polar angle between the $`i^{th}`$ strip and the nominal beam axis. The pretrigger hi requirement for a given octant was satisfied when the $`p_T`$ detected in either the inner 128 $`R`$ channels or the outer $`R`$ channels of that octant was greater than a threshold value. A pretrigger signal was issued only when there was no evidence in that octant of substantial noise, significant $`p_T`$ attributable to an earlier interaction, or incident beam-halo muon detected by the VW.
Localized trigger groups were formed for each octant by clustering the $`R`$-channel fast-outputs into 32 groups of 8 channels. Each adjacent pair of 8 channel groups formed a group-of-16 strips. If the $`p_T`$ detected in any of these groups-of-16 was above a specified high (or low) threshold, then a local hi (or local lo) signal was generated for that octant. A single local hi (or single local lo) trigger was generated if a local hi (or local lo) signal was generated in coincidence with the pretrigger hi in the same octant.
Trigger decisions were also made based upon global energy depositions within an octant. A global lo signal was generated if the total $`p_T`$ in an octant exceeded a threshold value. The local$``$global lo trigger required a coincidence of the pretrigger hi signal with global lo and local lo signals from the same octant. The local lo requirement was included to suppress spurious global triggers due to coherent noise in the EMLAC.
The single local lo and local$``$global lo triggers were prescaled to keep them from dominating the trigger rate. Prescaled samples of beam, interaction, and pretrigger events were also recorded.
## III Analysis Methods
Data samples contributing to this analysis represent an integrated luminosity of 1.6 (6.8) pb<sup>-1</sup> and 1.6 (6.5) pb<sup>-1</sup> for 530 and 800 GeV/$`c`$ $`p`$Cu ($`p`$Be) interactions, respectively, as well as 0.3 (1.4) pb<sup>-1</sup> for 515 GeV/$`c`$ $`\pi ^{}`$Cu ($`\pi ^{}`$Be) interactions. These samples were accumulated during the 1991–1992 run. Results reported in this paper also use 0.9 (6.1) pb<sup>-1</sup> of $`\pi ^{}`$Cu ($`\pi ^{}`$Be) data recorded during the 1990 run. The following subsections describe the analysis procedures and methods used to correct the data for losses from inefficiencies and selection biases. Additional details can be found in our previous papers E70 (a); Apanasevich et al. (2004a); Apanasevich et al. (1997); Apanasevich et al. (2004b, 1998b).
### III.1 Charged-particle reconstruction
The two major aspects of the analysis procedure involved charged-particle and calorimeter-shower reconstruction (discussed in Sec. III.2). The charged-track reconstruction algorithm produced track segments upstream of the magnet using information from the SSDs, and downstream of the magnet using information from the PWCs and STDCs. These track segments were projected to the center of the magnet and linked to form final tracks whose calculated charges and momenta were used for the physics analysis. The charged track reconstruction is described in more detail elsewhere Apanasevich et al. (1997); Blusk (1995). The primary vertex reconstruction is described in Sec. III.5.
### III.2 Calorimeter shower reconstruction
The readout of each EMLAC quadrant was divided into four regions: left and right $`R`$, and inner and outer $`\mathrm{\Phi }`$. Strip energies from clusters in each region were fit to the shape of an electromagnetic shower determined from detailed Monte Carlo simulations and isolated-shower data. These fits were used to evaluate the positions and energies of the peaks in each region. Shower positions and energies were obtained by correlating peaks of approximately the same energy in the $`R`$ and $`\mathrm{\Phi }`$ regions within the same half octant. More complex algorithms were used to handle configurations with showers spanning multiple regions. The EMLAC readout was also subdivided longitudinally into front and back sections. This segmentation provided discrimination between showers generated by electromagnetically or hadronically interacting particles. Photons were defined as showers with at least 20% of the shower energy deposited in the front part of EMLAC, to reduce the backgrounds due to showers from hadronic interactions. Losses of photons due to this requirement were $`2`$%. A detailed event simulation was employed to correct for this and other effects including reconstruction smearing and losses. An expanded discussion of the EMLAC reconstruction procedures and performance can be found elsewhere Apanasevich et al. (1998b).
### III.3 Meson signals
For this study, $`\pi ^0`$ and $`\eta `$ mesons were reconstructed via their $`\gamma \gamma `$ decay modes. Only those $`\gamma \gamma `$ combinations with energy asymmetry $`A_{\gamma \gamma }=|E_{\gamma _1}E_{\gamma _2}|/(E_{\gamma _1}+E_{\gamma _2})<0.75`$ were considered to reduce uncertainties due to low energy photons. The meson signals have been corrected for losses due to the energy asymmetry cut and the branching fractions for the $`\gamma \gamma `$ decay modes Groom et al. (2000).
Photons were required to be reconstructed within the fiducial region of the EMLAC to exclude areas with reduced sensitivity. In addition, $`\gamma \gamma `$ combinations were restricted to the same octant to simplify the trigger analysis. A simple ray-tracing Monte Carlo program was employed to determine the correction for these fiducial requirements.
The correction for losses due to the conversion of photons into $`e^+e^{}`$ pairs was evaluated by projecting each reconstructed photon from the event vertex to the reconstructed position in the EMLAC. The radiation length of material traversed, up to the analysis magnet, was evaluated based upon detailed detector descriptions. The photon conversion probability was evaluated and used to account for conversion losses. The average correction for conversion losses was $`1.09`$ per photon for the Be target in the 1990 run ($`1.08`$ in 1991–1992 run) and $`1.19`$ per photon for the Cu target ($`1.16`$ in 1991–1992 run).
### III.4 Detector simulation
The Meson West spectrometer was modeled with a detailed geant gea simulation (DGS). A preprocessor was used to convert geant information into the simulated hits and strip energies associated with the various detectors. The preprocessor simulated hardware effects, such as channel noise and gain variations. Monte Carlo events were then processed through the same reconstruction software used for the analysis of the data. This technique accounted for inefficiencies and biases in the reconstruction algorithms. Reconstuction inefficiencies were relatively small over most of the kinematic range. More information on the detailed simulation of the Meson West spectrometer can be found elsewhere Apanasevich et al. (2004b); E70 (a); Apanasevich (2005).
As inputs to the geant simulation, we employed single particle distributions, reconstructed data events, and herwig-generated Marchesini et al. (1992) events. The herwig-generated $`\pi ^0`$, $`\eta `$, and direct-photon spectra were weighted in $`p_T`$ and rapidity to our measured results in an iterative fashion so that the final corrections were based on the data distributions rather than on the behavior of the physics generator. Figure 3 shows the $`\gamma \gamma `$ mass spectra in the $`\pi ^0`$ and $`\eta `$ mass regions in comparison to the DGS results for the $`\pi ^{}`$Cu data at 515 GeV/$`c`$, and Fig. 4 shows an analogous plot for our higher statistics $`p`$Be data at 530 GeV/$`c`$. In addition to providing evidence that the DGS simulated the EMLAC resolution well, the agreement between the levels of combinatorial background indicates that the DGS also provided a reasonable simulation of the underlying event structure. Since the DGS was tuned using our higher statistics Be data, Figs. 56 also show the level of agreement achieved for this target. Figure 5 shows a comparison between the DGS and the data for the sideband-subtracted energy asymmetry distribution for photons from $`\pi ^0`$ decays. This figure illustrates that the simulation accurately describes the losses of low-energy photons.
A second Monte Carlo simulation of the detector (PMC) was used to cross check the detailed simulation and for studies that required large statistics. This simulation employed parameterizations of physics cross sections and detector responses Apanasevich et al. (2004b); Begel (1999); Apanasevich (2005). The inclusive $`\pi ^0`$ and direct-photon cross sections were parameterized as two dimensional surfaces in $`p_T`$ and rapidity Apanasevich (2005). The $`\eta `$, $`\omega `$, and $`\eta ^{}`$ cross sections were parameterized using the measured $`\eta /\pi ^0`$ E70 (a); Apanasevich et al. (2004a), $`\omega /\pi ^0`$ E70 (b), and $`\eta ^{}/\eta `$ Diakonou et al. (1980); Geist et al. (1990) ratios. Generated mesons were decayed into final state particles; photons were smeared for energy and position resolution Apanasevich et al. (1998b). A vertex was generated in the simulated target for every event. Photons were allowed to convert into $`e^+e^{}`$ pairs; the energy of the resulting electrons was reduced using the geant function for bremsstrahlung radiation. Electron four-vectors were smeared for multiple scattering in the target and the resolution of the tracking system and adjusted for the magnet impulse. Figure 6 displays a comparison between the PMC and the data in the $`\pi ^0`$ and $`\eta `$ mass regions and for the $`\pi ^0`$ energy asymmetry. The PMC provides an adequate characterization of the data.
### III.5 Vertex reconstruction
The location of the interaction vertex was reconstructed using charged-particle tracks. Vertices were identified by means of an impact-parameter minimization technique Blusk (1995). A $`\chi ^2`$ was defined for a given vertex position using the impact parameters of the reconstructed tracks and their projection uncertainties. The vertex position was found by minimizing this $`\chi ^2`$. Vertices were found in $`X`$ and $`Y`$ independently and correlated based on the difference in their positions along the nominal beam direction ($`Z`$ axis). The $`Z`$ position of the matched vertex was the weighted average of the $`Z`$ positions found in $`X`$ and $`Y`$. The reconstructed vertex positions are presented in Fig. 7 as functions of $`Z`$ for the two target configurations. The beryllium, copper, and hydrogen targets are clearly visible, as are the SSDs and related support structures. The average resolution for the $`Z`$ location of the interaction vertex was $`300\mu `$Apanasevich et al. (1997); Blusk (1995).
The relative heights of the Cu and Be targets shown in Fig. 7 varies as a function of $`p_T`$. This is clearly evident in Fig. 8, which compares two $`\pi ^0`$ samples: one acquired using the highly prescaled interaction trigger, and the other using the single local hi trigger. The interaction triggered events are typically minimum-bias in character with low-$`p_T`$ $`\pi ^0`$’s. The number of primary vertices scales as $`A^{2/3}`$, where $`A`$ is the atomic weight of the target. The single local hi triggers are typically caused by hard-scatters that produce high-$`p_T`$ $`\pi ^0`$’s. The number of primary vertices in these events scale as $`A^1`$.
The DGS was used for detailed studies of the vertex reconstruction. The transverse positions of vertices were chosen according to beam profiles observed in the data. Longitudinal positions were determined using Monte Carlo methods based upon the interaction lengths of the materials in the target region (Table 1). DGS events were weighted to reproduce the relative number of vertices in the data. Results from the DGS compare favorably with the data in Fig. 7. This good agreement was particularly important for separating events with primary vertices in copper from those in the upstream piece of beryllium. Figure 9 displays the longitudinal vertex distribution, focussing on this region in the 1991–1992 target configuration. The shape of the tails in the data are well described by the DGS.
The vertex reconstruction efficiency was evaluated using the DGS Apanasevich (2005). Separate reconstruction efficiencies were evaluated for the Be, Cu, and H<sub>2</sub> targets. The reconstruction probability was defined as the number of vertices reconstructed in each target’s fiducial volume divided by the number of vertices generated in the fiducial volume. The reconstruction efficiency was the inverse of this probability. Defined in this manner, the reconstruction efficiency also corrected for the longitudinal resolution smearing of reconstructed vertices.
Additional beam particles occasionally interacted in the target material within the data-capture timing window of the tracking system. The extra tracks sometimes caused the vertex associated with the high-$`p_T`$ interaction to be misidentified. The bias introduced by these rare events favored configurations where the low-$`p_T`$ interaction took place within the downstream piece of Be. This primarily affected interactions in the Cu and upstream Be targets in the long 1991–1992 target configuration because of the relatively poor vertex resolution in those targets compared to the downstream Be. This bias was investigated by comparing the $`\pi ^0`$ cross sections measured in $`\pi ^{}`$Be interactions in the 1990 and the 1991–1992 runs, and by comparing $`\pi ^0`$ yields from the upstream and downstream Be pieces in the 1991–1992 target configuration. The number of Cu vertices were corrected for misidentifications arising from this source. The resulting correction was $`1.04`$ for the 1991–1992 $`\pi ^{}`$Cu sample at 515 GeV/c, $`1.06`$ for the 530 GeV/$`c`$ $`p`$Cu sample, and $`1.12`$ for the 800 GeV/$`c`$ $`p`$Cu sample hyd .
Each event in this analysis was required to have a reconstructed vertex in the target region. Longitudinal and transverse requirements were placed on vertices to define the data samples. The longitudinal cuts selected the target in which the incident beam interacted, while the transverse cuts ensured that the interaction occurred within the target material.
### III.6 Direct photons
The largest contribution to the direct-photon background comes from electromagnetic decays of neutral hadrons, particularly $`\pi ^0`$’s and $`\eta `$’s. For the purposes of the measurements reported here, a photon was a direct-photon candidate if it did not combine with another photon in the same octant to form a $`\pi ^0`$ with $`A_{\gamma \gamma }0.9`$ or an $`\eta `$ with $`A_{\gamma \gamma }0.8`$.
To suppress electrons, reconstructed showers were excluded from the sample when charged-particle tracks pointed to within 1 cm of shower center. The correction for this criterion in the direct-photon analysis is $`1.01`$ based upon studies of the impact of this requirement on reconstructed $`\pi ^0`$’s.
The residual background from $`\pi ^0`$’s and $`\eta `$’s, as well as from other sources of background, was calculated using DGS samples that contained no generated direct photons ($`\gamma _b`$). A smooth fit of the $`\gamma _b/\pi ^0`$ in $`p_T`$ and rapidity was used to extract the direct-photon cross sections. The systematic uncertainty in this background subtraction was estimated by varying the direct-photon definition as a function of the cut on $`A_{\gamma \gamma }`$. The direct-photon background was also investigated using the PMC. Figure 10 compares the direct-photon backgrounds estimated using the two Monte Carlo simulations. The close agreement provides additional confidence in our understanding of the background in our direct photon samples. Additional details are provided in Ref. Apanasevich et al. (2004b).
Fits to $`\gamma _b/\pi ^0`$ were only made for the beryllium target due to relatively poor DGS statistics in the other targets. However, $`\gamma _b/\pi ^0`$ is expected to be slightly different for each target due to the different amounts of target material the photons must traverse. Therefore, a correction to $`\gamma _b/\pi ^0`$ was calculated using the PMC. The differences in $`\gamma _b/\pi ^0`$ are shown for the 800 GeV/$`c`$ proton beam data in Fig. 11. The other data samples have similar behavior. The target differences were fit as functions of $`p_T`$ and rapidity for each incident beam and target configuration and applied as additive corrections to the nominal $`\gamma _b/\pi ^0`$ fit.
The impact of the background can be determined by normalizing the direct-photon candidate spectrum from the data and the simulation to the measured $`\pi ^0`$ cross section E70 (a); Apanasevich et al. (2004a). This $`\gamma /\pi ^0`$ ratio is displayed in Fig. 12 for all three incident beams for the copper target. The signal-to-background in all cases is large at high $`p_T`$.
### III.7 Summary of systematic uncertainties
The systematic uncertainties for the production of $`\pi ^0`$’s, $`\eta `$’s, and direct photons measured in $`\pi ^{}`$Cu and $`p`$Cu interactions are similar to those detailed in Refs. E70 (a); Apanasevich et al. (2004a, b). The principal contributions to the systematic uncertainty arose from the following sources: normalization, calibration of photon energy response and detector-resolution unsmearing, background subtraction, reconstruction efficiency, incident beam contamination (for the 515 GeV/$`c`$ and 530 GeV/$`c`$ secondary beams), beam halo muon rejection, geometric acceptance, photon conversions, trigger response, and vertex finding. The additional vertex-finding uncertainty associated with the confusion induced by multiple beam particles interacting in the target was $`2\%`$. The total systematic uncertainties, combined in quadrature, are quoted with the cross sections in the appropriate tables. Note that some of these contributions to the systematic uncertainty (e.g. normalization) are strongly correlated between bins.
Most of the experimental systematic uncertainties cancel in the ratio of cross sections measured on different targets. The residual uncertainties are due to target-related systematics associated mainly with vertex identification. The total systematic uncertainty in the ratio of cross sections measured on Cu to those on Be is $`\pm 3\%`$ for the 800 GeV/$`c`$ $`p`$ beam sample, and $`\pm 2\%`$ for the 530 GeV/$`c`$ $`p`$ and 515 GeV/$`c`$ $`\pi ^{}`$ beam samples. The systematic uncertainty associated with the ratios of Be to H and Cu to H are $`\pm 4\%`$ for all samples after correcting for the effects of vertex misidentification as discussed in Sec. III.5.
## IV Results and Discussion
### IV.1 Cross sections
The invariant differential cross sections per nucleon for direct-photon, $`\pi ^0`$, and $`\eta `$ production from 530 and 800 GeV/$`c`$ $`p`$ beams and 515 GeV/$`c`$ $`\pi ^{}`$ beam incident on copper are presented as functions of $`p_T`$ in Figs. 13 through 15. Results from 530 GeV/$`c`$ $`p`$ and 515 GeV/$`c`$ $`\pi ^{}`$ beams are averaged over the rapidity range $`0.75y_{\mathrm{cm}}0.75`$; results from the 800 GeV/$`c`$ $`p`$ beam are averaged over $`1.0y_{\mathrm{cm}}0.5`$. Data points are plotted at abscissa values that correspond to the average value of the cross section in each $`p_T`$ bin, assuming local exponential $`p_T`$ dependence Lafferty and Wyatt (1995). The inclusive cross sections are tabulated in the Appendix.
### IV.2 Ratios
E706 has previously reported results for direct-photon and $`\pi ^0`$ production on beryllium and hydrogen targets Apanasevich et al. (2004b); E70 (a); Apanasevich et al. (2004a). Since E706 is the only direct-photon experiment that used more than one nuclear target, our data provide a unique measurement of nuclear effects in direct-photon production. Figures 16 to 18 present the ratio of inclusive cross sections per nucleon measured on Cu to those measured on Be for $`\pi ^0`$ mesons and direct photons. (The Be/H ratios were presented in previous publications Apanasevich et al. (2004b); E70 (a); Apanasevich et al. (2004a), and are not reproduced here.) These ratios show clear evidence of nuclear enhancement in both $`\pi ^0`$ and direct-photon production and the $`\pi ^0`$ data exhibit the decrease at high $`p_T`$ first noted by Cronin et al. Cronin et al. (1975); Antreasyan et al. (1979). This behavior is generally attributed to the influence of multiple-parton scattering prior to the hard scatter Wang (1997).
Results from fits to constant ratios in restricted regions of $`p_T`$ have been overlaid on the data in Figs. 16 to 18, and summarized in Table 2. Fits were made for $`3.5<p_T<5.5`$ GeV/$`c`$ for the $`\pi ^0`$’s and over the entire $`p_T`$ range for the direct photons. The difference in the ratio for $`\pi ^0`$ mesons and direct photons is significant and may be indicative of differing roles of initial and final-state effects, since direct photons are not expected to be strongly impacted by final-state nuclear effects.
Expectations from a theoretical prediction for nuclear enhancement in the 515 GeV/$`c`$ $`\pi ^{}`$ direct-photon sample Guo and Qiu (1996) have been overlaid on Fig. 18. This calculation predicts a slight enhancement in direct-photon production and agrees with the data within the uncertainties.
### IV.3 Comparisons with hijing
The Cu to Be cross-section ratios are compared with results from the hijing Monte Carlo event generator in Figs. 19 and 20. hijing is a program designed to simulate particle production in $`pp`$, $`pA`$, and $`AA`$ collisions wan . It was necessary to normalize the hijing results in Figs. 19 and 20 to the data. However, the shapes of the curves, in the case of the proton beams, are in good agreement with the data. Renewed interest in the amount of nuclear enhancement as a function of rapidity has been generated by recent BRAHMS measurements Arsene et al. (2004) and corresponding results from the other RHIC experiments sta ; phe ; pho . Our high-statistics cross-section measurements in the central rapidity region can be used to tune theoretical models developed to describe the RHIC environment. The rapidity dependence of the Cu to Be cross-section ratios are compared to expectations from hijing in Figs. 21 to 23. hijing does not describe the rapidity dependence of the $`\pi ^0`$ data for the incident proton beams; the hijing results are generally peaked towards backward rapidities (like the BRAHMS data), whereas our data are relatively independent of rapidity. hijing provides a better description of our direct-photon data.
## V Conclusions
We have measured the invariant differential cross section per nucleon for direct-photon, $`\pi ^0`$, and $`\eta `$ production from 515 GeV/$`c`$ $`\pi ^{}`$ beam and 800 and 530 GeV/$`c`$ proton beams incident on copper as a function of $`p_T`$ and $`y_{\mathrm{cm}}`$. These data span the kinematic range $`1.0<p_T10`$ GeV/$`c`$ and central rapidities.
Ratios of these production cross sections to our previously published measurements on a beryllium target Apanasevich et al. (2004b); E70 (a); Apanasevich et al. (2004a) show a strong nuclear enhancement for $`\pi ^0`$ mesons and a smaller, but significant, enhancement for direct photons. We compare these measurements with expectations from a theoretical calculation (for direct-photon production for incident $`\pi ^{}`$ beam) and with the results of a Monte Carlo event generator, hijing, dedicated to the simulation of nuclear effects. hijing yields a good description of the shape of the $`p_T`$ dependence of the Cu to Be ratios for both $`\pi ^0`$’s and direct photons in the incident proton beam data samples. hijing also describes the rapidity dependence of direct photon production in those samples. However, hijing provides a relatively poor description of the rapidity dependence of $`\pi ^0`$ meson production.
###### Acknowledgements.
We thank the U. S. Department of Energy, the National Science Foundation, including its Office of International Programs, and the Universities Grants Commission of India, for their support of this research. The staff and management of Fermilab are thanked for their efforts in making available the beam and computing facilities that made this work possible. We are pleased to acknowledge the contributions of our colleagues on Fermilab experiment E672. We acknowledge the contributions of the following colleagues to the operation and upgrade of the Meson West spectrometer: W. Dickerson and E. Pothier from Northeastern University; J. T. Anderson, E. Barsotti Jr., H. Koecher, P. Madsen, D. Petravick, R. Tokarek, J. Tweed, D. Allspach, J. Urbin, and the cryo crews from Fermi National Accelerator Laboratory; T. Haelen, C. Benson, L. Kuntz, and D. Ruggiero from the University of Rochester; the technical staffs of Michigan State University and Pennsylvania State University for the construction of the straw tubes and of the University of Pittsburgh for the silicon detectors. We thank the following commissioning run collaborators for their invaluable contributions to the hardware and software infrastructure of the original Meson West spectrometer: G. Alverson, G. Ballocchi, R. Benson, D. Berg, D. Brown, D. Carey, T. Chand, C. Chandlee, S. Easo, W. Faissler, G. Glass, I. Kourbanis, A. Lanaro, C. Nelson Jr., D. Orris, B. Rajaram, K. Ruddick, A. Sinanidis, and G. Wu. We also thank X.-N. Wang for several helpful discussions.
*
## Appendix A Tabulated Cross Sections
In this appendix, we present tables of the measured invariant differential cross sections for direct photon, $`\pi ^0`$, and $`\eta `$ production on Cu targets as functions of $`p_T`$. In these tables, the first uncertainty is statistical and the second is systematic. In the case of the lowest two $`p_T`$ bins for the $`\pi ^0`$ measurement, the statistical and systematic uncertainties have been combined because of the large correlation between them.
|
warning/0506/math-ph0506070.html
|
ar5iv
|
text
|
# Untitled Document
Conformal Field Theories on the Two-Torus and Quotients of $`SL2(Z)`$
Antoine D. Coste\* , Giovanni Felder
* LPT-CNRS, UMR 8627,
University building 210 , F 91405 Orsay cedex, France
coste@math.uni-frankfurt.de
and FIM ETH Zentrum , CH-8092 Zurich
Switzerland
Abstract
We present remarkable properties of the groups $`SL2(Z/NZ)`$ which might be useful in detailed studies of some quotients appearing in Conformal Field Theories (CFTs).
Introduction
The main object underlying this study is a finite dimensional representation $`(\rho ,V)`$ of the group $`SL2(Z)`$ of matrices $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$. Our (writing) efforts are justified by the high interest for physics and mathematics of the space V; indeed there exist an infinity of such representations, a countable subset of explicit examples being provided by integrable representations of infinite dimensional affine algebras, (“Kac-Moody” algebras) and Virasoro algebra representations called “minimal models”. Any of the rational conformal field theories ( whenever their classification will be achieved) will provide such a representation.
In the following paragraphs we look at subgroups or quotients of $`SL2(Z)`$ expressing elements in terms of words in the two generators T and S. Physical motivation for this is that $`T`$ is represented by a unitary diagonal matrix whose eigenvalues are related to the physical dimension of fundamental excitations, whereas $`S`$ describes the effect of putting a box on one of its sides, or permuting the role of “space” and “time” in some hamiltonian description.
Since this field of research is very popular, we refer the reader to monographies (like the one by J.M. Drouffe and C. Itzykson) and go directly into:
Preliminary formulae:
The best to get acquainted with beauties of matrix groups over rings is maybe to let oneself play; therefore let us consider products of matrices of the form:
$$T^xS:=\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=\left(\begin{array}{cc}x& 1\\ 1& 0\end{array}\right)$$
$$T^{x_n}S\mathrm{}T^{x_1}S:=\left(\begin{array}{cc}A_n& B_n\\ C_n& D_n\end{array}\right)=\left(\begin{array}{cc}x_nA_{n1}C_{n1}& x_nB_{n1}D_{n1}\\ A_{n1}& B_{n1}\end{array}\right)$$
For arbitrary elements of $`SL2(Z)`$, we know the minimal n required is not bounded, according to the famous Farey’s enumeration of rational numbers between $`0`$ and $`1`$.
For $`n=3`$ we get:
$$T^xST^dST^yS=\left(\begin{array}{cc}(xdy+x+y)& 1+xd\\ (dy+1)& d\end{array}\right)$$
and also
$$T^uST^cST^v=\left(\begin{array}{cc}uc1& ucvuv\\ c& cv1\end{array}\right)$$
$$T^uST^cS^1T^v=\left(\begin{array}{cc}uc+1& ucv+u+v\\ c& cv+1\end{array}\right)$$
We will always denote integers with lower case letters: a,b,c,d ; and by A,B,C,D their residues modulo N , a fixed integer equal to the order of the matrix $`\rho (T)`$. Quite often $`adbc=1`$ and when $`(c,N)=1`$ we will denote by $`C^1`$ the inverse of C in $`Z/NZ`$. In the case $`(d,n)=1`$ we will denote its inverse also by $`𝐃^\mathrm{𝟏}`$.
In the next paragraphs, we will use some formulae valid for elements $`A,B,C,D`$ satisfying $`ADBC1`$ in a commutative ring. (in view of Z/NZ we will denote equality by $``$, but the formulae are also true in Z with equality of integers).
Proposition: If there exists $`U`$ such that $`UCA+1`$, then
$$T^UST^CST^{UDB}\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Proof: is straightforward from above formula. Similarly we also have:
Proposition: If there exists $`U`$ such that $`UCA1`$, then
$$T^UST^CS^1T^{BDU}\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Proposition: If there exists $`X`$ such that $`XDB1`$, then
$$T^XST^DST^{XCA}S\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Proposition: If there exists $`X`$ such that $`XDB+1`$, then
$$T^XST^DS^1T^{AXC}S\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Proposition: If there exists $`X`$ such that $`XA(1+C)`$, then
$$ST^XST^AST^{(D+BX)}\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Proposition: If there exists $`X`$ such that $`XA1C`$, then
$$ST^XST^AS^1T^{D+BX}\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
Examples:
$$\left(\begin{array}{cc}3& 4\\ 2& 3\end{array}\right)=T^2ST^2ST^2\text{ in }Z$$
$$\left(\begin{array}{cc}5& 8\\ 2& 3\end{array}\right)=T^2ST^2S^1T^2\text{ in }Z$$
If $`(c,N)=1`$ , we have both solutions $`(U_+,V_+)((A+1)C^1,(D+1)C^1)`$ ,
and $`(U_{},V_{})((A1)C^1,(D1)C^1)`$, where $`C^1`$ is the inverse of C mod N:
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)T^{AC^1}𝒮_CT^{DC^1}$$
$$\text{where }𝒮_C:=T^{C^1}ST^CST^{C^1}$$
$$T^{C^1}ST^CS^1T^{C^1}\left(\begin{array}{cc}0& C^1\\ C& 0\end{array}\right)(\sigma _{C^1}(S),\text{see below })$$
These equalities reflect some relations in the group $`SL2(Z/NZ)`$ :
$$T^{2C^1}ST^CS=ST^CS^1T^{2C^1}$$
$$\text{Since }S\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)S^1=\left(\begin{array}{cc}D& C\\ B& A\end{array}\right),$$
only when the four residues $`A,B,C,D`$ are non invertible mod $`N`$, could we have complicated expressions for elements of $`SL2(Z/NZ)`$. This could happen only when $`N`$ is not a prime power.
Results from Physical Approaches
In the sequel, we will use an abelian group of automorphisms (which we call the $`\sigma _L`$’s ) of the group $`SL2(Z/NZ)`$ defined by:
$$\sigma _L\left(\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\right):\left(\begin{array}{cc}A& BL\\ CL^1& D\end{array}\right)$$
Obviously, this group of automorphisms is isomorphic to the multiplicative group of invertible residues mod N:
Furthermore these morphisms satisfy:
$$\sigma _K(S)\sigma _L(S)\sigma _{K/L}(S)SS\sigma _{L/K}(S)\left(\begin{array}{cc}K/L& 0\\ 0& L/K\end{array}\right)$$
Theorems( de Boer, Goeree, A.C., Gannon, Lascoux, Bantay):
Any Rational Conformal Field Theory defines a representation $`\rho `$ of $`SL2(Z)`$ whose matrix elements $`\rho (M)_{pp^{}}`$ are in a cyclotomic field of N-th roots of unity, $`\rho (T^N)`$ is the identity, $`\mathrm{\Gamma }(N)`$ lies in its kernel , and for any matrix M of $`SL2(Z/NZ)`$, the above morphisms go into the cyclotomic characters:
$$\sigma _L(\rho (M)_{pp^{}})=\rho (\sigma _L(M))_{pp^{}}$$
$$\text{where on the l.h.s }\sigma _L(\xi _N)=\xi _N^L$$
In previous texts, explicit computations were given for the example of $`sl(2)`$ affine Lie algebra (so called Wess Zumino Witten models), Furthermore it was exposed how these theorems can lead to very compact formulae for the representation on Virasoro characters.
Relations
In this paragraph we give examples of constructions of quotient groups as anounced above:
We define the relations $`_N`$ to be:
$$_N:S^4=T^N=1,(ST)^3=S^2$$
The group generated by S, T and these relations, is for $`N6`$, associated to a tesselation by triangles. The question of which relations can be added in order to obtain a quotient group which is finite, or a triangulated surface of finite genus ( and finite area) naturally arises. Note that for some authors a Riemann surface is compact and this is not the case as long as we keep the hyperbolic metric structure for which the punctures are at ( logarithmically ) infinite distance.
Other questions are: what are the relations which lead to the group $`SL2(Z/NZ)`$ ?. At given N, what are the relations which give a surface of minimal genus?
Here we address only the first question for which we need the following:
Lemma: when $`BC2`$,
$$ST^CST^BT^BST^CS^1\left(\begin{array}{cc}1& B\\ C& 1\end{array}\right)$$
when $`BC0`$,
$$ST^CST^BT^BST^CS\left(\begin{array}{cc}1& B\\ C& 1\end{array}\right)$$
$$ST^CS^1T^BT^BST^CS^1\left(\begin{array}{cc}1& B\\ C& 1\end{array}\right)$$
Proposition:
$$ST^UST^AST^VT^XST^DST^YS\text{If and only if }$$
$$A(X+BY),D(V+BU),UXVY$$
$$\text{ where necessarily }B1+XD1+AV$$
Proposition: let p be a prime
$$\text{If }N=p^n>2,SL2(Z/NZ)\text{is presented with extra relations:}$$
$$H_A:=T^AST^{1/A}ST^AS^1\text{for invertible}A^{}s,H_AH_B=H_{AB}$$
$$H_AT=T^{A^2}H_A,H_AS=S^1H_{1/A}$$
Proof: define for each $`(C,D)`$ such that there exist $`A,B`$, $`ADBC1`$, a word $`X_{(C,D)}`$ in S and T which corresponds to both an SL2 matrix and an hyperbolic triangle. We can enumerate elements of the group mod N as words of the form $`T^xX_{(C,D)}`$.
Then one proves that $`SL2(Z/NZ)/\pm 1`$ is generated by relations given above (with $`S^2=1`$ ) by a careful study of glueing formulae at the boundary of the connected triangulated domain. This is achieved by checking by use of the above relations that for each $`X_{(C,D)}`$ , $`X_{(C,D)}T^{\pm 1}`$ and $`X_{(C,D)}S`$ are of the form $`T^LX_{(C^{},D^{})}`$ for some $`L,C^{},D^{}Z/NZ`$ . That the relations do not give a smaller quotient comes from the fact easily checked that the matrix group $`SL2(Z/NZ)`$ explicitly do satisfy these relations.
Since it deserves some time, let us give some explicit steps in a constructive proof of the above presentation:
$`X_{(0,1)}=`$ identity.
$`X_{(1,D)}=ST^D`$
For $`(c,N)=1`$, $`2c\frac{N1}{2}`$ and for any D: $`X_{(C,D)}=ST^CST^{(D+1)/C}`$. Note that for more general N one could also take
$`X_{(C,D)}=ST^CST^X\left(\begin{array}{cc}1& X\\ C& D\end{array}\right)`$ whenever exists $`X`$ such that $`CXD+1`$ ; note that this includes the cases $`X_{(C,1)}=ST^CS`$.
For c not coprime with p (thus with N ) , $`0c\frac{N}{2}`$ and $`(d,N)=1`$ , $`2d\frac{N1}{2}`$ we take $`X_{(C,D)}=ST^dST^{(1c)/d}S^1`$.
With the same c and $`(d,N)=1`$ , $`\frac{N}{2}<d2`$
$$X_{(C,D)}=ST^DST^{(1+C)/D}S\left(\begin{array}{cc}(1+C)/D& 1\\ C& D\end{array}\right)$$
.
Finally for C non invertible, $`2c<N/2`$: $`X_{(C,1)}=ST^CS^1=\left(\begin{array}{cc}1& 0\\ C& 1\end{array}\right)`$
This presentation by generators and relations being established, a few remarks are useful:
A dual point of view which is also useful is to construct the surface by glueing N-gons which are collections of triangles labelled by words $`Y_{(A,C)}T^z`$. Then the centers of the corresponding N-gon can be seen as having coordinate $`\tau =\frac{a}{c}`$ on the real axis ( boundary of the upper half plane). A more rigourous formulation is of course to identify the center of the N-gon to the orbit of $`\frac{a}{c}`$ under homographic $`\mathrm{\Gamma }(N)`$ tranformations.
Of course, the above relations in terms of the $`H_A^{}s`$ are redondant, it suffices to have them for generators of this abelian group (Cartan torus) isomorphic to the group of invertible residues mod N. Even one can find in literature various relations, which we already collected in a previous electronic text with T. Gannon (arXiv/math.QA/9909080 ). According to various authors their relations do in fact imply the above relations. (Such implication may come from interesting constructions and from the congruence subgroup property).
We give below for small values of N, such simple and compact looking presentations.
Note the genus of the Riemann surface is
$$\text{if }N=p^n,g=1+\frac{p^{3n}p^{3n2}6p^{2n}+6p^{2n2}}{24}$$
The above construction allows us to enumerate explicitly elements of $`SL_2(Z/NZ)`$, when N is a prime power. Of course an explicit use of the chinese remainder theorem gives us a description of the general group $`SL_2(Z/NZ)`$ as direct product of its primary factors.
Let us give explicitly formulae for the generators in terms of decomposition of $`1`$ (mod $`N`$ ) into orthogonal sum of idempotents: This is textbook result for commutative semisimple algebras, called in french “algèbres réduites” by Bourbaki:
$$1\mathrm{\Sigma }_{p,p|N}c_p\text{ mod }NT_p:=T^{c_p}$$
$$S_p:=S^2(ST^{1c_p})^3=S^2(T^{1c_p}S)^3$$
But a decomposition into primary factors gives a much too complicated description of elements of $`SL_2(Z/NZ)`$ as words in $`S`$ and $`T`$. There is a much smarter approach!:
Proposition:
Any element of $`SL_2(Z/NZ)`$ can be written with at most four powers of $`T`$, i. e. as a word like:
$$T^{x_3}ST^{x_2}ST^{x_1}ST^{x_0}$$
Proof: We start with the following lemma:
Let a,b,c,d, N be five integers satisfying $`adbc=1`$ , $`N>0`$ . Then there exists an integer m such that $`d^{}:=dmc`$ is coprime to $`N`$. Then denote $`𝐃^{}`$ its residue mod $`N`$ , and $`D^1`$ its inverse.
Dirichlet has even proven that one could find an infinity of values of m such that $`d^{}`$ is a prime number not dividing N; but here the requirement is much weaker, so that one can find a convenient m again with help of the chinese remainders, because that means there exist $`u,v`$ such that $`u(dmc)vN=1`$ which is equivalent to the existence of a residue $`M_p`$ modulo each $`p^\nu `$ factor of N such that the residue $`D_pM_pC_p`$ is invertible mod $`p^\nu `$. If $`D_p`$ is invertible, $`M_p=0`$ works, and if $`D_p`$ is not invertible , $`C_p`$ is and therefore any $`M_p`$ which is not divisible by $`p`$ will do the job.
Then we have ($`b^{}=bma`$ , $`d^{}=dmc`$ ):
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\left(\begin{array}{cc}a& bma\\ c& dmc\end{array}\right)\left(\begin{array}{cc}1& m\\ 0& 1\end{array}\right)$$
$$\text{Thus }\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)T^{(B^{}1)/D^{}}ST^D^{}ST^{(1+C^{})/D^{}}ST^M\text{ mod }N$$
Physicists’ intuition would find it natural since $`SL2(R)`$ is a three dimensional manifold, and indeed here we have succeeded in decomposing any group element into an expression with four parameters so there is some kind of one-parameter degree of arbitrariness. Nevertheless this is too naive a picture because $`SL2(Z)`$ is really exceptional: one needs unbounded words to express any matrix, the continued fraction of any rational $`a/c`$ can be of any arbitrary length. This is even exceptional, compared to $`SL3(Z)`$.
As a conclusion we could say that the groups $`SL2(Z/NZ)`$ appear in fact simpler than what could be anticipated from $`SL2(Z)`$. This allowed us to improve slightly upon prior works, bringing our little stone to the building.
Nevertheless there is a fascinating interplay between geometry and number theory, as usual due to the apparently chaotic occurence of primes and congruences when one decomposes an integer, N , into its prime factors. If we were considering strings, membranes or black holes we could claim that Conformal Field Theories bring more pieces into some cosmic hyperbolic puzzle! We prefer to let the reader appreciate the intrinsic beauty of mathematics and rigorous physics.
Examples:
We finally give explicit presentations for small values of $`N`$: $`N=5`$ is the famous F. Klein’s icosaedron (or dodecaedron). For $`N=6`$ we have a torus, which can be equivalently defined by the two presentations below. Another very interesting approach from a geometric point of view is to identify fundamental domains as done in places like Bonn by Kulkarni( see refs ). Examples of quotients from conformal theories are detailed in previous texts by A. C.
$$SL_2(Z/5Z)=<S,T|_5>$$
$$SL_2(Z/6Z)=<S,T|S^4=T^6=1,(ST)^3=S^2,ST^2ST^2=T^2ST^2S>$$
$$SL_2(Z/6Z)=<S,T|S^4=T^6=1,(ST)^3=S^2,ST^3ST^2=T^2ST^3S>$$
$$SL_2(Z/8Z)=<S,T|S^4=T^8=1,(ST)^3=S^2,ST^2ST^4=T^4ST^2S>$$
$$SL_2(Z/9Z)=<S,T|S^4=T^9=1,(ST)^3=S^2,ST^3ST^2S^1=T^4ST^2ST^2$$
$$,(ST^3)^2=(T^3S)^2,ST^4ST^4=T^4ST^4S^1,$$
$$ST^2ST^4ST^2=T^4ST^2ST^4S^1,$$
$$ST^2ST^2ST^4=T^2ST^4ST^3S^1>$$
$$SL_2(Z/10Z)=<S,T|S^4=T^{10}=1,(ST)^3=S^2,ST^2ST^5=T^5ST^2S$$
$$,ST^3ST^4=T^4S^1T^3S,ST^3ST^3ST=T^1ST^3ST^3S$$
$$,ST^4ST^5=T^5ST^4S>\text{ genus }13$$
References and Acknowledgements A. C. thanks J. Lascoux, H. Behr , T. Gannon , P. Ruelle, J. Wolfart, J-C. Schlage-Puchta, J. Froehlich, V. Pasquier, D. Zagier, for help and many conversations, M. Burger, M. Kraemer for very kind hospitality in FIM ETHZ.
P. Bantay, “The kernel of the modular representation and the Galois action in RCFT”, Communications in Math. Phys. 233 (2003) 423-438, and further texts.
A. Coste, “On Rational Conformal Field Theories: Explicit Modular Formulae” math-phys/0405, Letters in Mathematical Physics, Dijon (2005) vol72, 1-15.
A. Coste “Investigations sur les caracteres de Kac Moody et certains quotients de SL2(Z)” (in french) preprint IHES/P/1997/78, unpublished.
A. Coste, T. Gannon, Physics Letters B 323 (1994) 316.
J.M. Drouffe, C. Itzykson, “Statistical Theory of Fields (10 chapters)” , Ediscience, and references therein.
J. de Boer, J. Goeree, Communications in Math. Phys. 139, (1991) 267.
E. Hecke , Elementary number theory, (translated from German and reprinted) Springer GTM.
Ravi S. Kulkarni, “An arithmetic geometric study of the subgroups of the modular group”, American Journal of Math. 113 (1991), 1053-1133.
H. Rademacher, Analytic Number Theory, Springer.
|
warning/0506/math0506262.html
|
ar5iv
|
text
|
# Homological Properties of Color Lie Superalgebras
## 1. Introduction
Color Lie superalgebras are graded over an Abelian group $`G`$ and generalize Lie superalgebras. Background information on Lie superalgebras and color Lie superalgebras can be found in and and and , respectively. We will show, in a forthcoming paper , that the ideal structure of the enveloping algebra of a color Lie superalgebra can be very different from enveloping algebras of Lie superalgebras. By contrast, the results of this note illustrate that certain homological properties are the same. Our main strategy is to pass to the case where the grading group $`G`$ is finitely generated, so that the color Lie superalgebra is determined by a Lie superalgebra and a 2-cocycle defined on $`G`$.
In §2 we define color Lie superalgebras and state a theorem due to Scheunert (see ). In §3 we calculate the finitistic and injective dimensions of the enveloping algebra of a finite dimensional color Lie superalgebra over a field of characteristic $`0`$. The enveloping algebra of a finite dimensional color Lie superalgebra may have infinite global dimension (this is known for ordinary Lie superalgebras, see \[2, Proposition 5\]). However, for a finite dimensional color Lie superalgebra $``$ which is *positively graded*, i.e. $`=_+`$, we prove that, in analogy with the nongraded case, $`\mathrm{g}ldimU()=dim`$. In particular, theorem 3.1 generalizes \[8, Proposition 2.3\].
In theorem 3.2 we show that the enveloping algebra $`U()`$ of a color Lie superalgebra is Auslander-Gorenstein and Cohen-Macaulay. It thus follows from \[9, Theorem 1.4\] that $`U()`$ has a (right and left) quasi-Frobenius (QF) classical quotient ring.
## 2. Color Lie Superalgebras
Throughout $`k`$ denotes a field of characteristic $`2`$, $`G`$ denotes an Abelian group, and all algebras are associative $`k`$-algebras with 1. We call a map $`\gamma :G\times Gk^\times `$ ($`k^\times =k\backslash \{0\}`$) a *skew-symmetric bicharacter* on $`G`$ if it satisfies
$$\begin{array}{ccc}\text{(1)}\hfill & \text{ }\hfill & \gamma (f,gh)=\gamma (f,g)\gamma (f,h)\text{ and }\gamma (gh,f)=\gamma (g,f)\gamma (h,f)\text{ for any }f,g,hG\text{.}\hfill \\ \text{(2)}\hfill & \text{ }\hfill & \gamma (g,h)\gamma (h,g)=1\text{ for any }g,hG\text{.}\hfill \end{array}$$
Note that (2) implies that $`\gamma (g,g)=\pm 1`$ for any $`gG`$. Set $`G_\pm =\{gG:\gamma (g,g)=\pm 1\}`$, then $`G_+`$ is a subgroup of $`G`$ and $`[G:G_+]2`$. In \[14, Lemma 2\] it is shown that if $`G`$is finitely generated then the skew-symmetric bicharacters on $`G`$ are completely determined by the 2-cocycles on $`G`$ with coefficients in $`k^\times `$.
Generally, we shall be discussing objects which are graded over $`G`$ so let $`G`$-vec denote the category of $`G`$-graded vector spaces and graded linear maps. For any $`VG`$-vec, set $`x=g`$ if $`0xV_g`$. Let $`\gamma `$ be a skew-symmetric bicharacter on $`G`$. We will shorten our notation by writing $`\gamma (x,y)`$ instead of $`\gamma (x,y)`$ for homogeneous $`0xV`$, $`0yW`$ and $`V,WG`$-vec. Note that $`G`$-vec is naturally contained in $`_2`$-vec via the decomposition $`G=G_+\dot{}G_{}`$ and the group homomorphism $`\pi :G_2`$ given by $`\pi (G_+)=\overline{0}`$ and $`\pi (G_{})=\overline{1}`$.
Definition. A $`(G,\gamma )`$*-color Lie superalgebra* is a Lie algebra a pair $`(,,)`$ such that $`G`$-vec and $`,:_k`$ is a graded bilinear map which satisfies the following for any homogeneous $`x,y,z`$.
| $`\gamma \text{-skew-symmetry}`$ | $`x,y=\gamma (x,y)y,x`$ |
| --- | --- |
| $`\gamma \text{-Jacobi identity}`$ | $`\gamma (z,x),x,y,z+\gamma (y,z),z,x,y+\gamma (x,y),y,z,x=0`$ |
Example. (1) A *Lie superalgebra* is a $`(_2,\chi )`$-color Lie superalgebra where $`\chi (\overline{i},\overline{j})=`$
$`(1)^{ij}`$ for any $`i,j`$.
1. If $`A`$ is a $`G`$-graded algebra then there is a color Lie superalgebra structure $`A^{}`$ defined on $`A`$ subject to the condition that $`a,b=ab\gamma (a,b)ba`$ for any nonzero homogeneous $`a,bA`$.
We shall now describe how to construct color Lie superalgebras from Lie superalgebras as in . In as much as Lie superalgebras are $`_2`$-graded, we must consider $`G`$-gradings on Lie superalgebras which respect the $`_2`$-grading.
Definition. A Lie superalgebra $`(L,[,])`$, $`L=L_{\overline{0}}L_{\overline{1}}`$, is $`G`$*-graded* if $`L=_{gG}L_g`$, $`[L_f,L_g]L_{fg}`$ for any $`f,gG`$, and $`L_hL_{\overline{0}}`$ or $`L_hL_{\overline{1}}`$ for each $`hG`$.
Starting with a Lie superalgebra $`L`$ there are many choices of groups $`G`$ and $`G`$-gradings on $`L`$. Let $`G(L)`$ be the smallest subgroup of $`G`$ which grades $`L`$, that is, $`G(L)`$ is generated by $`\{gG:L_g0\}`$. We cannot always pass to the case that $`G=G(L)`$ since, for example, there are $`G`$-graded representations of $`L`$ which are not $`G(L)`$-graded.
Let $`G\left(L\right)_+`$ be the subgroup of $`G(L)`$ which is generated by $`\{gG:U(L)_gU(L)_{\overline{0}}\}`$. Then $`[G(L):G(L)_+]2`$. Set $`G\left(L\right)_{}=G\left(L\right)\backslash G\left(L\right)_+`$; then the pair $`(L,[,])`$ is a $`(G\left(L\right),\gamma _0)`$-color Lie superalgebra where $`\gamma _0`$ is the skew-symmetric bicharacter on $`G\left(L\right)`$ defined below.
$$\gamma _0(g,h)=\{\begin{array}{cc}\hfill 1& \text{if }g,hG\left(L\right)_{}\hfill \\ \hfill 1& \text{otherwise}\hfill \end{array}$$
To continue our construction of a color Lie superalgebra from a Lie superalgebra, we need the notion of a 2-cocycle.
Definition. A *2-cocycle* on $`G`$ is a map $`\sigma :G\times Gk^\times `$ which satisfies
$$\sigma (f,gh)\sigma (g,h)=\sigma (f,g)\sigma (fg,h)$$
for any $`f,g,hG`$.
If $`\sigma `$ is a 2-cocycle on $`G`$ then there is a skew-symmetric bicharacter $`\gamma `$ defined on $`G\left(L\right)`$ by $`\gamma (g,h)=\gamma _0(g,h)\sigma (g,h)\sigma (h,g)^1`$ for any $`g,hG\left(L\right)`$. Moreover, there is a $`(G\left(L\right),\gamma )`$-color Lie superalgebra $`(L^\sigma ,[,]^\sigma )`$ which has the same vector space structure as $`L`$ but bracket $`[,]^\sigma :L^\sigma \times L^\sigma L^\sigma `$ defined subject to the condition that $`[x,y]^\sigma =\sigma (x,y)[x,y]`$ for any homogeneous $`x,yL`$.
By setting $`f=h=e`$ in Definition 2 (where $`e`$ is the identity element of $`G`$) we obtain $`\sigma (e,e)=\sigma (g,e)=\sigma (e,g)`$ for any $`gG`$. Note that there is no loss in generality by assuming that $`\sigma (e,e)=1`$ since we can replace $`\sigma `$ with $`\sigma ^{}`$, where $`\sigma ^{}=\sigma (e,e)^1\sigma `$ and have $`L^\sigma L^\sigma ^{}`$ as $`(G\left(L\right),\gamma )`$-color Lie superalgebras.
Theorem 2.1 summarizes the results from we are interested in. For the reader’s convenience we will go over some basic definitions.
Definition. Let $`(,,)`$ be a $`(G,\gamma )`$*-color Lie superalgebra*.
1. A linear map $`\varphi :_1_2`$ between $`(G,\gamma )`$-color Lie superalgebras $`(_1,,_1)`$ and $`(_2,,_2)`$ is called a *homomorphism* if $`\varphi (x,y_1)=\varphi (x),\varphi (y)_2`$ for any $`x,y_1`$.
2. A $`G`$*-graded representation* of $``$ is a pair $`(V,\rho )`$ where $`VG`$-vec and $`\rho :End_k(V)^{}`$ is a homomorphism of color Lie superalgebras ($`End_k(V)`$ is a $`G`$-graded algebra since $`VG`$-vec$`).`$
3. For any $`G`$-graded algebra $`A`$, a graded map $`\varphi :A^{}`$ is called *compatible* if it is homomorphism of $`(G,\gamma )`$-color Lie superalgebras.
4. The *universal enveloping algebra* of $``$ is a $`G`$-graded algebra $`U()`$ and a compatible map $`\iota :U()`$ which satisfies the property that for any compatible map $`\varphi :A`$ there is a unique (graded) algebra homomorphism $`\mathrm{\Phi }:U()A`$ such that $`\mathrm{\Phi }\iota =\varphi `$.
###### Theorem 2.1 (Scheunert).
Let $``$ be a $`(G,\gamma )`$-color Lie superalgebra and $`L`$ and $`L^\sigma `$ as above.
1. If $`G`$ is finitely generated, then any color Lie superalgebra can be obtained as an $`L^\sigma `$ for appropriately chosen $`L`$ and 2-cocycle $`\sigma `$.
2. The enveloping algebra $`U(L^\sigma )`$ is obtained from $`U(L)`$ by defining a new multiplication $``$ on $`U`$ subject to the condition that for any homogeneous $`x,yU(L)`$ we have $`xy=\sigma (x,y)xy.`$
3. For a graded representation $`\rho :LEnd_k(V)`$ of $`L`$, there is a graded representation $`\rho ^\sigma :L^\sigma End_k(V)`$ of $`L^\sigma `$ which is obtained from $`\rho `$ subject to the condition that for any homogeneous $`xL^\sigma `$ and $`vV`$, $`\rho ^\sigma (x)(v)=\sigma (x,v)\rho (x)(v)`$. This defines a category equivalence between the categories of graded representations of $`L`$ and $`L^\sigma `$.
We are particularly interested in parts (2) and (3) of theorem 2.1. More generally, consider an arbitrary $`G`$-graded algebra $``$ and left $``$-modules $`V`$ and $`W`$. An $``$-module homomorphism $`\varphi :VW`$ is called *graded* if $`\varphi (V_h)W_h`$ for each $`hG`$. Let $`{}_{R}{}^{}_{}^{G}`$ denote the category whose objects are $`G`$-graded left $``$-modules and morphisms are all graded $``$-module homomorphisms from $`V`$ to $`W`$, denoted $`HOM_{}^0(V,W)`$, where $`V,W_R^G`$. In particular, $`{}_{k}{}^{}_{}^{G}=G`$-vec.
###### Lemma 2.2.
Let $`\sigma `$ be a 2-cocycle on $`G`$ and $`R`$ a $`G`$-graded $`k`$-algebra.
1. There is a $`G`$-graded $`k`$-algebra $`R^\sigma `$ with the same vector space structure as $`R`$ and whose multiplication $``$ is obtained from the multiplication of $`R`$ subject to the condition $`rs=\sigma (r,s)rs`$ for any homogeneous $`r,sR`$.
2. For any $`M_R^G`$ there corresponds $`M^\sigma _{R^\sigma }^G`$ such that $`M^\sigma `$ has the same vector space structure as $`M`$ but the $`R^\sigma `$-module structure . is obtained from the $`R`$-module structure of $`M`$ subject to the condition $`r.m=\sigma (r,m)rm`$ for any homogeneous $`rR`$ and $`mM`$.
3. The functor $`{}_{}{}^{\sigma }:_{R}^{}^G_{R^\sigma }^G`$ defined as in (2) is a category equivalence. In particular, $`HOM_R^0(V,W)=HOM_{R^\sigma }^0(V^\sigma ,W^\sigma )`$. Moreover, if $`V_R^G`$ is graded free, then so is $`V^\sigma _{R^\sigma }^G`$.
4. For any $`V,W_R^G`$ and $`\varphi HOM_R^0(V,W)`$ we have $`\mathrm{ker}(\varphi ^\sigma )=(\mathrm{ker}\varphi )^\sigma =\mathrm{ker}\varphi `$ and $`\mathrm{i}m(\varphi ^\sigma )=(\mathrm{i}m\varphi )^\sigma =\mathrm{i}m\varphi `$.
###### Proof.
It is easy to prove (1)-(3) directly. For (4), first note that $`\mathrm{ker}\varphi `$, $`\mathrm{i}m\varphi `$, $`\mathrm{ker}(\varphi ^\sigma )`$ and $`\mathrm{i}m(\varphi ^\sigma )`$ are $`G`$-graded submodules. For any homogeneous $`xV`$ we have $`x\mathrm{ker}(\varphi ^\sigma )x\mathrm{ker}(\varphi )x(\mathrm{ker}\varphi )^\sigma `$. This proves that $`\mathrm{ker}(\varphi ^\sigma )=(\mathrm{ker}\varphi )^\sigma =\mathrm{ker}\varphi `$. The proof that $`\mathrm{i}m(\varphi ^\sigma )=(\mathrm{i}m\varphi )^\sigma =\mathrm{i}m\varphi `$ is similar.
## 3. Homological Properties of $`U()`$
Throughout this section, $`\gamma `$ is a skew-symmetric bicharacter on $`G`$ and $`G_+`$=$`\{gG:\gamma (g,g)=\pm 1\}`$. For any $`VG`$-vec, set $`V_\pm =_{gG_\pm }V_g`$. If $`V=V_+`$ (respectively $`V=V_{}`$) then $`V`$ is called *positively* (respectively *negatively*) *graded*. We shall be primarily interested in the case where $``$ is an Abelian color Lie superalgebra. If $``$ is a positively graded, finite dimensional and Abelian $`(G,\gamma )`$-color Lie superalgebra, then $`U()k^\gamma [\theta _1,\theta _2,\mathrm{},\theta _n]`$, the *color polynomial ring* in $`n=dim`$ variables (see ). If $``$ is a negatively graded, finite dimensional and Abelian $`(G,\gamma )`$-color Lie superalgebra, then $`U()\mathrm{\Lambda }^\gamma ()`$, the *color exterior algebra (or color Grassmann algebra) of* $``$ (see ).
Example.
1. Choose $`qk^\times `$. Let $`\gamma `$ be the skew-symmetric bicharacter on $`G=^2`$ defined by $`\gamma (g,h)=q^{g_1h_2g_2h_1}`$ for any $`g=(g_1,g_2)`$, $`h=(h_1,h_2)G`$. Then $`G=G_+`$. Consider the Abelian $`(G,\gamma )`$-color Lie superalgebra with homogeneous basis $`\{x,y\}`$ such that $`x=(1,0)`$ and $`y=(0,1)`$. Then $``$ is positively graded and $`U()k[x,y:xy=qyx]`$, the so-called *quantum plane*.
2. Choose $`0qk^\times `$. Let $`\gamma `$ be the skew-symmetric bicharacter on $`G=^2`$ defined by $`\gamma (g,h)=(1)^{g_1h_1+g_2h_2}q^{g_1h_2g_2h_1}`$ for any $`g=(g_1,g_2)`$, $`h=(h_1,h_2)G`$. Then $`G_+=\{(i,j)G:i+j2\}`$ and $`G_{}=\{(i,j)G:i+j12\}`$. Consider the Abelian $`(G,\gamma )`$-color Lie superalgebra with homogeneous basis $`\{x,y\}`$ such that $`x=(1,0)`$ and $`y=(0,1)`$. Then $``$ is negatively graded and $`U()k[x,y:xy=qyx]/I`$, where $`I`$ is the ideal generated by $`\{x^2,y^2\}`$.
If $`V,WG`$-vec then the tensor product $`VWG`$-vec has grading defined by $`(VW)_g=_{hG}(V_hW_{gh^1})`$. For $`G`$-graded algebras $`A`$ and $`B`$ there is a $`\gamma `$*-graded tensor product* $`A\widehat{}_kB`$ with multiplication defined subject to the condition that $`(a\widehat{}b)(a^{}\widehat{}b^{})=\gamma (b,a^{})aa^{}\widehat{}bb^{}`$ for any nonzero homogeneous $`a,a^{}A`$ and $`b,b^{}B`$. The $`\gamma `$-graded tensor product is essential to the study of enveloping algebras of color Lie superalgebras. For example, the enveloping algebra of a color Lie superalgebra is not a Hopf algebra in the usual sense; it is only a Hopf algebra with respect to the $`\gamma `$-graded tensor product (see \[1, §3.2.9\] for details).
For any Noetherian ring $`R`$, the left (respectively right) finitistic dimension of $`R`$, denoted $`\mathrm{l}FPD\left(R\right)`$ (respectively $`\mathrm{r}FPD\left(R\right)`$), is the supremum of the projective dimensions of left (respectively right) $`R`$-modules of finite projective dimension. The following theorem generalizes \[8, Proposition 2.3\]. As in and , we assume that $`\mathrm{c}hark=0`$ in all that follows.
###### Theorem 3.1.
Let $`=_+_{}`$ be a finite dimensional color Lie superalgebra. Then $`\mathrm{g}ldim(U(_+))=\mathrm{l}FPD(U())=\mathrm{r}FPD(U())=\mathrm{i}njdim_{U()}(U())=dim(_+)`$.
###### Proof.
Let $`U=U()`$ denote the universal enveloping algebra of $``$ and set $`n=dim(_+)`$. We first show that $`\mathrm{l}FPD(U)=\mathrm{r}FPD(U))=\mathrm{i}njdim_U(U)=n\mathrm{g}ldim(U(_+))`$. Then we show that $`\mathrm{g}ldim(U(_+))n`$, that is, there exists a $`U(_+)`$-module with projective dimension at least $`n`$. As in the case of ordinary Lie algebras, the existence of such a module will be demonstrated using the *Chevalley-Eilenberg complex* (see Chapter 7 of for background).
1. With respect to the standard filtration, $`gr\left(U\right)`$ is Noetherian.
The filtration is defined by $`U^1=\{0\}`$, $`U^0=k`$ and, for $`m>0`$, $`U^m`$ is spanned by all monomials of length $`m`$. By \[1, Proposition III.2.8\] $`U`$ is Noetherian and the associated graded algebra $`gr\left(U\right)`$ is isomorphic to $`k^\gamma [\theta _1,\theta _2,\mathrm{},\theta _n]\widehat{}\mathrm{\Lambda }^\gamma (_{})`$, where $`\mathrm{\Lambda }^\gamma (_{})`$ is the color exterior algebra determined by the vector space $`_{}`$ with an Abelian color Lie super algebra structure. Thus $`gr\left(U\right)`$ is an iterated Ore extension of $`\mathrm{\Lambda }=\mathrm{\Lambda }^\gamma (_{})`$ without derivations, that is, $`gr\left(U\right)\mathrm{\Lambda }[\theta _1;\alpha _{\theta _1}][\theta _2;\alpha _{\theta _2}]\mathrm{}[\theta _n;\alpha _{\theta _n}]`$. Therefore, $`gr\left(U\right)`$ is Noetherian by \[11, Theorem I.2.10\].
1. $`\mathrm{i}njdim_UU\mathrm{i}njdim_{gr\left(U\right)}gr\left(U\right)=n`$ and $`\mathrm{g}ldimU(_+)\mathrm{g}ldimgr(U(_+))=n`$
Since $`U^m=0`$ when $`m<0`$, the filtration $`\{U^m\}`$ satisfies the *closure condition* of \[3, §1.2.11\]. Therefore, we may apply \[3, Theorem 1.3.12\] to obtain $`\mathrm{i}njdim_UU\mathrm{i}njdim_{gr\left(U\right)}gr\left(U\right)`$. Moreover, $`\mathrm{g}ldimU(_+)\mathrm{g}ldimgr(U(_+))`$ by \[11, 7.6.18\]. For any ring $`R`$ and automorphism $`\sigma `$ of $`R`$, one can show that $`\mathrm{g}ldim(R[t;\sigma ])=\mathrm{g}ldim(R)+1`$ and $`\mathrm{i}njdim_{R[t;\sigma ]}(R[t;\sigma ])=\mathrm{i}njdim_R(R)+1`$ by similar methods to the proofs in the case that $`\sigma `$ is the identity on $`R`$. Since $`gr(U(_+)k^\gamma [\theta _1,\theta _2,\mathrm{},\theta _n]`$ we have $`\mathrm{g}ldimgr(U(_+))=n=dim_+`$. By \[6, Corollary 6.3\] $`\mathrm{\Lambda }`$ is a Frobenius algebra; hence $`\mathrm{i}njdim_{gr(U)}gr(U)=n=dim_+`$ as $`gr(U)\mathrm{\Lambda }^\gamma [\theta _1,\theta _2,\mathrm{},\theta _n]`$.
1. $`\mathrm{l}FPD(U)=\mathrm{r}FPD(U))=\mathrm{i}njdim_U(U)\mathrm{g}ldim(U(_+))`$
Since $`\mathrm{g}ldimU(_+)`$ is finite and $`U`$ is free over $`U(_+)`$ on both sides, we may apply \[4, Theorem 1.4\] to obtain $`\mathrm{g}ldimU(_+)\mathrm{l}FPDU`$. By \[8, Proposition 2.1\] we have $`\mathrm{i}njdim_UU=\mathrm{l}FPD(U)=\mathrm{r}FPD(U)`$.
*The rest of the proof is dedicated to showing that if* $``$ *is a positively graded finite-dimensional* $`(G,\gamma )`$*-color Lie superalgebra then* $`n=dim\mathrm{g}ldimU()`$*.*
1. Pass to the case that $`=L^\sigma `$ for some Lie algebra $`L`$.
The modules we will consider are graded over the subgroup $`G()`$ of $`G`$ generated by $`\{gG:_g0\}`$. Therefore, we can pass to the case that $`G`$ is finitely generated by assuming that $`G=G()`$. By theorem 2.1(1), there is a $`G`$-graded Lie algebra $`L`$ with bracket $`[,]`$ and a 2-cocycle $`\sigma `$ on $`G`$ such that $`=L^\sigma `$.
1. The Chevalley-Eilenberg complex $`V_{}(L)^\underset{}{\epsilon }k`$ is a complex in $`{}_{U\left(L\right)}{}^{}_{}^{G}`$.
We must show that $`V_{}(L)^\underset{}{\epsilon }k`$ is a *graded complex*, i.e., the modules are graded and all differentiations are graded maps. Let $`\mathrm{\Lambda }L`$ denote the exterior algebra of $`L`$ and $`\mathrm{\Lambda }^iL`$ denote the $`i^{\text{th}}`$ exterior algebra of $`L`$. Then $`\mathrm{\Lambda }L=_{i=1}^n\mathrm{\Lambda }^iL`$ can be $`G`$-graded so that the homogeneous component of degree $`gG`$ is $`(\mathrm{\Lambda }L)_g=_{i=1}^n(\mathrm{\Lambda }^iL)_g`$ and $`(\mathrm{\Lambda }^iL)_g=_{g_1g_2\mathrm{}g_i=g}L_{g_1}L_{g_2}\mathrm{}L_{g_i}`$ where the sum is indexed by all appropriate $`g_1,g_2,\mathrm{},g_iG`$.
We have $`V_i(L)=U(L)_k\mathrm{\Lambda }^iL`$ which is a graded free left $`U(L)`$-module such that $`V_i(L)_g=_{hG}(U\left(L\right)_h)_k(\mathrm{\Lambda }^iL)_{hg^1}`$ so we only need to check that the *augmentation map* $`\epsilon :V_0(L)k`$ is graded and, for $`1in`$, $`d_i:V_i(L)V_{i1}(L)`$ is graded. The augmentation map $`\epsilon :V_0(L)=U\left(L\right)k`$ is graded since it is induced from the Lie algebra map $`Lk`$ which sends everything to zero, a graded map. It is also easy to see that $`d:V_1(L)V_0(L)`$ is graded since it is just the product map $`uxux`$. For $`1<in`$, the map $`d_i:V_i(L)V_{i1}(L)`$ is defined by the formula $`d_i(ux_1x_2\mathrm{}x_i)=\theta _1+\theta _2`$ where $`\theta _1`$ and $`\theta _2`$ are defined below.
$`\theta _1`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{i}{}}}(1)^{j+1}ux_jx_1x_2\mathrm{}\widehat{x}_j\mathrm{}x_i`$
$`\theta _2`$ $`=`$ $`{\displaystyle \underset{l<m}{}}(1)^{l+m}u[x_l,x_m]x_1x_2\mathrm{}\widehat{x}_l\mathrm{}\widehat{x}_m\mathrm{}x_i`$
Therefore $`d_i`$ is graded since if $`u,x_1,x_2,\mathrm{},x_i`$ are homogeneous elements then $`ux_1x_2\mathrm{}x_i,\theta _1`$ and $`\theta _2`$ are all homogeneous of the same degree.
1. The Chevalley-Eilenberg complex $`V_{}()^\underset{}{\epsilon }k`$ is a complex in $`{}_{U}{}^{}_{}^{G}`$.
By parts (2), (3) and (4) of Lemma 2.2, there is a projective resolution $`V_{}()`$ of $`k=k^\sigma `$ over $`U()`$ such that $`V^i()=V^i(L)^\sigma `$. For $`1in`$ the module $`V^i(L)`$ is a graded free $`U(L)`$-module hence the module $`V^i(L)^\sigma `$ is a graded free $`U(L^\sigma )`$-module by Lemma 2.2(3).
1. $`\mathrm{g}ldimU()n`$
As in \[16, Exercise 7.7.2\], there is a graded representation $`\rho :LEnd_k(\mathrm{\Lambda }^nL)`$ defined as below for $`0yL`$ and $`x_1,x_2,\mathrm{},x_n`$ a (homogeneous) basis of $`L`$.
$$\rho (y)(x_1x_2\mathrm{}x_n)=\underset{i=1}{\overset{n}{}}x_1x_2\mathrm{}[y,x_i]\mathrm{}x_n$$
Let $`M=\mathrm{\Lambda }^nL`$ be the (graded) $`U(L)`$-module defined by the above action. Then according to , $`\mathrm{E}xt_{U(L)}^n(k,M)\mathrm{ker}(d_n^{})/\mathrm{i}m(d_{n1}^{})k`$. By parts (2) and (4) of Lemma 2.2 and 6) above, we have $`\mathrm{ker}d_n^{}=(\mathrm{ker}d_n^{})^\sigma =\mathrm{ker}(d_n^{})^\sigma `$ and $`\mathrm{i}md_{n1}^{}=(\mathrm{i}md_{n1}^{})^\sigma =\mathrm{i}m(d_{n1}^{})^\sigma `$. This implies that
$$\underset{U}{\overset{n}{\mathrm{E}xt}}(k,M^\sigma )\left(\underset{U\left(L\right)}{\overset{n}{\mathrm{E}xt}}(k,M)\right)^\sigma k$$
therefore $`pd_U(M^\sigma )n`$ hence $`\mathrm{g}ldimUn`$ as desired.
The result follows from 3) and 7) above.
###### Theorem 3.2.
Let $``$ be a finite dimensional color Lie superalgebra. Then $`U()`$ is Auslander-Gorenstein and Cohen-Macaulay and thus has a quasi-Frobenius classical quotient ring.
###### Proof.
We have $`\mathrm{\Lambda }`$ is Auslander-Gorenstein and Cohen-Macaulay since $`\mathrm{\Lambda }`$ is a Frobenius algebra (see \[6, Corollary 6.3\]). As gr$`(U(L))`$ is an iterated Ore extension of $`\mathrm{\Lambda }`$ where each iteration is of the form $`R[x;\sigma ]`$, it follows from \[10, Lemma (ii), p. 184\] and \[5, Theorem 4.2\] that gr$`(U(L))`$ is Auslander-Gorenstein and Cohen-Macaulay. We thus conclude that $`U()`$ is Auslander-Gorenstein and Cohen-Macaulay by \[3, Theorem 1.4.1\] and \[15, Lemma 4.4\] (note that Theorem 3.1 is needed so that $`\mathrm{i}njdim_{U()}U()<\mathrm{}`$). The last statement follows from \[9, Theorem 1.4\].
Acknowledgements
The author is indebted to Jim Kuzmanovich for many helpful conversations concerning this material.
|
warning/0506/hep-th0506248.html
|
ar5iv
|
text
|
# More on the Triplet Killing Potentials of Quaternionic Kähler manifolds
## Abstract
We show the properties of the triplet Killing potentials of quaternionic Kähler manifolds which have been missing in the literature. It is done by means of the metric formula of the manifolds. We compute the triplet Killing potentials for the quaternionic Kähler manifold $`Sp(n+1)/Sp(n)Sp(1)`$ as an illustration.
PACS: 02.40.ky, 02.40.Tt, 04.65.+e
Keywords: Kähler manifold, Supergravity, Coset space
It is well-known that when hypermultiplets couple with $`𝒩`$=$`2`$ supergravity in four dimensions, the moduli space is required to be a quaternionic Kähler manifold. A quaternionic Kähler manifold is a $`4n`$-dimensional real manifold. The holonomy group of quaternionic Kähler manifolds is $`Sp(1)Sp(n)`$ and the $`Sp(1)`$ conection must not vanish. When the gravitational coupling $`\kappa 0`$, this $`Sp(1)`$ conection is required to vanish, the holonomy group gets contained in $`Sp(n)`$ and the quaternionic Kähler manifold approaches to a hyper Kähler manifold. Such a quaternionic or hyper Kähler manifold appears as a moduli space for type II superstring compactification on Calabi-Yau manifolds, together with a special Kähler manifold given by vector multiplets. The direct product of both moduli spaces is an arena for the mirror symmetry in the type II superstring.
Quaternionic Kähler manifolds are by now well-studied subjects-. The existence of the triplet Killing potentials $`\stackrel{}{M}^A`$ is a hallmark of quaternionic Kähler manifolds when they are group manifolds. The properties of the $`\stackrel{}{M}^A`$ were well studied in the literature-. But there are still important properties which are missing in generality of the arguments. In this letter we show them, explicitly constructing the metric of quaternionic Kähler manifolds. With the specific metric at hand the geometric quantities of the manifolds like the metric, the Riemann curvature etc, are given in terms of the triplet Killing potentials alone. Among them we stress on the relations implying that $`\stackrel{}{M}^A`$ is an eigen vector of the Beltrami-Laplace operator for the manifold and the quantity $`\stackrel{}{M}^A\stackrel{}{M}^A`$, which was called energy in , is a constant related with the Riemann scalar curvature. This is justified as long as the metric of quarternionic Kähler group manifolds manifolds is non-degenerate. Similar relations were known for the Killing potentials which exist for the ordinary Kähler group manifolds-. We develope our arguments noting a pararellism between both Kähler group manifolds. Finally we compute the triplet Killing potentials for the quarternionic Kähler coset space $`Sp(n+1)`$ $`/Sp(n)Sp(1)`$ to illustrate our general arguments.
We start with a summary of the geometric properties discussed in -. A quaternionic Kähler manifold is endowed with a triplet complex structure $`\stackrel{}{J}_a^b`$. Assume that it is a group manifold admitting an isometry $`G`$ realized by Killing vectors $`R^{Aa}`$ with $`A=1,2,\mathrm{},\mathrm{dim}G`$. For the quantity $`\stackrel{}{J}_a^bR_b^A(\stackrel{}{J}_a^bg_{bc}R^{Ac})`$ we have triplet Killing vectors $`\stackrel{}{M}^A`$ such that
$`_a\stackrel{}{M}^A=\nu \stackrel{}{J}_a^bR_b^A,\nu {\displaystyle \frac{R}{4n(n+2)}},`$ (1)
with the Riemann scalar curvature $`R`$. It is given by
$`\stackrel{}{M}^A=\stackrel{}{r}^A+R^{Aa}\stackrel{}{\omega }_a,`$ (2)
with a spin connection $`\stackrel{}{\omega }_a`$ and an appropriate $`\stackrel{}{r}^A`$ which will be explained soon later. $`\stackrel{}{M}^A`$ can be also written in the form
$`\stackrel{}{M}^A={\displaystyle \frac{1}{2n}}\stackrel{}{J}_a^b_bR^{Aa}`$ (3)
and satisfies the relation
$`\nu R_a^AR_b^B\stackrel{}{J}^{ab}=f^{ABC}\stackrel{}{M}^C+\stackrel{}{M}^A\times \stackrel{}{M}^B.`$ (4)
Here $`f^{ABC}`$ are the structure constants of $`G`$ and $`\stackrel{}{J}^{bc}g^{bd}\stackrel{}{J}_d^c`$.
We remember quite similar arguments for Kähler manifolds-. Namely a Kähler manifold is endowed with a singlet complex structure $`J_a^b`$. It may be locally set to be
$`J_a^b=\left(\begin{array}{cc}i\delta _\alpha ^\beta & 0\\ & \\ 0& i\delta _{\overline{\alpha }}^{\overline{\beta }}\end{array}\right).`$ (8)
If the Kähler manifold admits an isometry $`G`$, realized by (anti-)holomorphic Killing vectors $`R^{A\alpha }(\overline{R}^{A\overline{\alpha }})`$, there exist Killing potentials such that
$`_\alpha M^A=ig_{\alpha \overline{\beta }}\overline{R}^{A\overline{\beta }},_{\overline{\alpha }}M^A=ig_{\beta \overline{\alpha }}R^{A\beta }.`$ (9)
It was shown that such Killing potentials are given by
$`M^A`$ $`=`$ $`{\displaystyle \frac{i}{𝒩_{adj}}}_\alpha R^{A\alpha }\left(={\displaystyle \frac{i}{𝒩_{adj}}}_{\overline{\alpha }}\overline{R}^{A\overline{\alpha }}\right),`$ (10)
$`orM^A`$ $`=`$ $`i(R^{A\alpha }_\alpha KF^A),`$ (11)
and satisfy the relation
$`iR^{B\beta }\overline{R}^{C\overline{\gamma }}g_{\beta \overline{\gamma }}+iR^{C\beta }\overline{R}^{B\overline{\gamma }}g_{\beta \overline{\gamma }}=f^{ABC}M^A.`$ (12)
In (10) we have used the normalization $`f^{ABC}f^{ABD}=2𝒩_{adj}\delta ^{CD}`$. In (11) $`K`$ is the Kähler potential and $`F^A(\overline{F}^A)`$ are (anti-)holomorphic quantities that one may find from the Lie-variation
$`_{R^A}K=F^A+\overline{F}^A.`$ (13)
Eqs (1)$``$(4) of quaternionic Kähler group manifolds are so similar to (9)$``$(12) of Kähler group manifolds. About $`M^A`$ of Kähler group manifolds we know more properties than those. By the Lie-variation we have
$`_{R^A}M^B=f^{ABC}M^C.`$ (14)
From (9) and (10) it follows that $`M^A`$ is an eigen vector of the Beltrami-Laplace operator as
$`g^{\alpha \overline{\beta }}_\alpha _{\overline{\beta }}M^A=𝒩_{adj}M^A`$ (15)
Moreover, when Kähler manifolds are irreducible coset spaces, then the geometric quantities can be written in terms of the Killing potentials
$`g_{\alpha \overline{\beta }}`$ $`=`$ $`_\alpha M^A_{\overline{\beta }}M^A,`$ (16)
$`R_{\alpha \overline{\beta }\gamma \overline{\delta }}`$ $`=`$ $`_\alpha _{\overline{\beta }}M^A_\gamma _{\overline{\delta }}M^A=f^{ABE}f^{CDE}R_\alpha ^AR_{\overline{\beta }}^BR_\gamma ^CR_{\overline{\delta }}^D,`$ (17)
$`R`$ $`=`$ $`R_{\alpha \overline{\beta }\gamma \overline{\delta }}g^{\alpha \overline{\beta }}g^{\gamma \overline{\delta }}=𝒩_{adj}^2M^AM^A.`$ (18)
The geometrical properties of quaternionic Kähler manifolds which correspond to (14)$``$ (18) are missing in the literature-. In this letter we show them by completing the pararellism between both Kähler group manifolds. They are
$`_{R^A}\stackrel{}{M}^B`$ $`=`$ $`f^{ABC}\stackrel{}{M}^C+\stackrel{}{r}^A\times \stackrel{}{M}^B,`$ (19)
$`_a^a\stackrel{}{M}^A`$ $`=`$ $`2n\nu \stackrel{}{M}^A,`$ (20)
$`g_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{3\nu ^2}}_a\stackrel{}{M}^A_b\stackrel{}{M}^A,`$ (21)
$`R`$ $`=`$ $`{\displaystyle \frac{2}{3}}n(n+2)\stackrel{}{M}^A\stackrel{}{M}^A,`$ (22)
$`R_{abcd}`$ $`=`$ $`{\displaystyle \frac{1}{3\nu ^2}}[[_a,_b],_c]\stackrel{}{M}^A_d\stackrel{}{M}^A`$
$`=`$ $`f^{ABE}R_a^AR_b^Bf^{CDE}R_c^CR_d^D`$
$`{\displaystyle \frac{1}{2}}f^{ABC}R_a^AR_b^BR^{Ce}f^{DEF}R_c^DR_d^ER_e^F`$
$`{\displaystyle \frac{1}{4}}f^{ABC}R_d^AR_b^BR^{Ce}f^{DEF}R_c^DR_a^ER_e^F`$
$`+{\displaystyle \frac{1}{4}}f^{ABC}R_d^AR_a^BR^{Ce}f^{DEF}R_c^DR_b^ER_e^F.`$
Here $`\stackrel{}{r}^A`$ is the same quantity that appeared in (2). These formulae are useful when four-fermi couplings and scalar potentials in $`𝒩`$=$`2`$ supergravity are analyzed from a phenomelogical point of view, by identifying the isometry group $`G`$ or some subgroup with a grand unification gauge group.
We give a short review on quaternionic Kähler manifolds. Consider a real $`4n`$-dimensional Riemann manifold $``$ with local coordinates $`\varphi ^a=(\varphi ^1,\varphi ^2,\mathrm{},\varphi ^{4n})`$. The line element of the manifold is given by $`ds^2=g_{ab}d\varphi ^ad\varphi ^b.`$ If $``$ is a quaternionic manifold, there exists a triplet complex structure $`\stackrel{}{J}_a^b(J_a^{1b},J_a^{2b},J_a^{3b})`$ satisfying the property
$`J_a^{\alpha b}J_b^{\beta c}=\delta ^{\alpha \beta }\delta _a^c+ϵ^{\alpha \beta \gamma }J_a^{\gamma c}.\alpha =1,2,3,`$ (24)
and the holonomy group in the tangent space is $`Sp(1)Sp(n)`$. We define a set of vielbein $`1`$-forms
$`d\varphi ^ae_a^{i\mu },i=1,2,\mu =1,2,\mathrm{},2n,`$
which satisfy
$`e_a^{i\mu }e_{i\mu }^b=\delta _a^b,e_a^{i\mu }e_{j\nu }^a=\delta _j^i\delta _\nu ^\mu .`$
The holonomy groups $`Sp(1)`$ and $`Sp(n)`$ act on the indices $`i`$ and $`\mu `$ respectively. Using these vielbeins we can construct the triplet complex structure as
$`\stackrel{}{J}_a^b=ie_a^{i\mu }\stackrel{}{\sigma }_i^je_{j\mu }^b.`$ (25)
Lower or raise the indices as $`\stackrel{}{J}_a^bg_{bc}\stackrel{}{J}_{ac}`$ and $`g^{ab}\stackrel{}{J}_b^c=\stackrel{}{J}^{ac}`$. Then it is easy to show that $`\stackrel{}{J}_{ab}=\stackrel{}{J}_{ba},\stackrel{}{J}^{ab}=\stackrel{}{J}^{ba}.`$ We postulate that the vielbeins are covariantly constant:
$`_ae_b^{i\mu }=_ae_b^{i\mu }\mathrm{\Gamma }_{ab}^ce_c^{i\mu }+\omega _{aj\nu }^{i\mu }e_b^{j\nu }=0.`$ (26)
Here $`\omega _{aj\nu }^{i\mu }`$ is the spin connection Lie-valued in $`Sp(1)Sp(n)`$, so that
$`\omega _{aj\nu }^{i\mu }=\omega _{aj}^i\delta _\nu ^\mu +\omega _{a\nu }^\mu \delta _j^i.`$ (27)
We solve (26) for the spin conections $`\omega _{aj}^i`$ and $`\omega _{a\nu }^\mu `$ to find
$`\omega _{aj}^i={\displaystyle \frac{1}{2n}}e_{j\nu }^b\stackrel{}{}_ae_b^{i\nu },\omega _{a\nu }^\mu ={\displaystyle \frac{1}{2}}e_{j\nu }^b\stackrel{}{}_ae_b^{j\mu }.`$
Here $`\stackrel{}{}_c`$ is the covariant derivative which does not contain the spin connection. From the postulate it follows that the metric and the triplet complex structure are covariantly constant:
$`_ag_{bc}`$ $`=`$ $`_ag_{bc}\mathrm{\Gamma }_{ab}^eg_{ec}\mathrm{\Gamma }_{ac}^eg_{be}=0,`$ (28)
$`_a\stackrel{}{J}_b^c`$ $`=`$ $`_a\stackrel{}{J}_b^c\mathrm{\Gamma }_{ab}^e\stackrel{}{J}_e^c+\mathrm{\Gamma }_{ae}^c\stackrel{}{J}_b^e+\stackrel{}{\omega }_a\times \stackrel{}{J}_b^c=0,`$ (29)
with $`\stackrel{}{\omega }_ai\stackrel{}{\sigma }_i^j\omega _{aj}^i.`$ The Riemann curvature tensor is given by
$$R_{abc}^d=_b\mathrm{\Gamma }_{ac}^d_a\mathrm{\Gamma }_{bc}^d\mathrm{\Gamma }_{bc}^e\mathrm{\Gamma }_{ae}^d+\mathrm{\Gamma }_{ac}^e\mathrm{\Gamma }_{be}^d.$$
Correponding to the spin connection (27) it decomposes in the tangent space:
$`R_{abc}^de_d^{i\mu }e_{j\nu }^c=R_{abj}^i\delta _\nu ^\mu +R_{ab\nu }^\mu \delta _j^i.`$
We may write the $`Sp(1)`$ curveture $`R_{abj}^i`$ as $`\stackrel{}{R}_{ab}=i\stackrel{}{\sigma }_i^jR_{abj}^i,`$ or equivalently as
$`\stackrel{}{R}_{ab}={\displaystyle \frac{1}{2n}}R_{abc}^d\stackrel{}{J}_d^c={\displaystyle \frac{1}{2n}}R_{abcd}\stackrel{}{J}^{cd}={\displaystyle \frac{1}{n}}R_{cab}^d\stackrel{}{J}_d^c.`$ (30)
From the postulate (26) it is given by
$`\stackrel{}{R}_{ab}=_{[a}\stackrel{}{\omega }_{b]}+\stackrel{}{\omega }_a\times \stackrel{}{\omega }_b.`$ (31)
In they have shown various relations among $`\stackrel{}{R}_{ab},\stackrel{}{J}_a^b,R_{ab}(R_{cabd}g^{cd})`$ and $`R(`$ $`g^{ab}R_{ab})`$:
$`R_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(n+2)\stackrel{}{J}_a^c\stackrel{}{R}_{cb},`$ (32)
$`\stackrel{}{R}_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{n+2}}\stackrel{}{J}_a^cR_{cb}={\displaystyle \frac{1}{2}}\stackrel{}{J}_a^c\times \stackrel{}{R}_{cb},`$ (33)
$`R_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{4n}}Rg_{ab}.`$ (34)
Combining (33) and (34) gives
$`\stackrel{}{R}_{ab}=\nu \stackrel{}{J}_{ab}.`$ (35)
When the quaternionic Kähler manifold is a coset space $`G/H`$, there is a set of the Killing vectors $`R^{Aa}=(R^{1a},R^{2a},\mathrm{},R^{Na})`$ with $`N=\mathrm{dim}G`$ which represents the isometry group $`G`$. They are required to satisfy
$`_{R^A}R^{Ba}`$ $``$ $`R^{Ab}_bR^{Ba}R^{Bb}_bR^{Aa}=f^{ABC}R^{Ca},`$ (36)
$`_{R^A}g_{ab}`$ $``$ $`R^{Ac}_cg_{ab}+_aR^{Ac}g_{cb}+_bR^{Ac}g_{ca}=0,`$ (37)
$`_{R^A}\stackrel{}{J}_a^b`$ $``$ $`R^{Ac}_c\stackrel{}{J}_a^b+_aR^{Ac}\stackrel{}{J}_c^b_cR^{Ab}\stackrel{}{J}_a^c=\stackrel{}{r}^A\times \stackrel{}{J}_a^b,`$ (38)
in which $`_{R^A}`$ is the Lie-variation with respect to $`R^A`$. For covariant derivatives of the Killing vectors we have
$`_aR_b^A+_bR_a^A`$ $`=`$ $`0,`$ (39)
$`_a_bR_c^A`$ $`=`$ $`R_{bca}^dR_d^A,`$ (40)
$`_a_bR_c^A+_b_cR_a^A+_c_aR_b^A`$ $`=`$ $`0.`$ (41)
$`R_a^A_bR_c^A`$ $`=`$ $`{\displaystyle \frac{1}{2}}f^{ABC}R_a^AR_b^BR_c^C,`$ (42)
Eq. (39) is the Killing equation which is equivalent to (37). Eq. (40) follows by calculating the l.h.s. as
$`_a_bR_c^A`$ $`=`$ $`{\displaystyle \frac{1}{2}}_{\{a}_{b\}}R_c^A+{\displaystyle \frac{1}{2}}_{[a}_{b]}R_c^A`$
$`=`$ $`{\displaystyle \frac{1}{2}}(_{[a}_{c]}R_b^A_{[b}_{c]}R_a^A)+{\displaystyle \frac{1}{2}}_{[a}_{b]}R_c^A`$
$`=`$ $`{\displaystyle \frac{1}{2}}(R_{acb}^dR_{bca}^d+R_{abc}^d)R_d^A,`$
with (39) and the cyclic property of the Riemann curvature tensor
$`R_{abc}^d+R_{bca}^d+R_{cab}^d=0.`$ (43)
Combining (40) and (43) we obtain (41). When the coset space $`G/H`$ is irreducible, it follows from (37) or (39) that
$`g_{ab}=R_a^AR_b^A.`$ (44)
In using this metric formula our arguments differ from those in the literature. We will be able to explore more on the triplet Killing potentials $`\stackrel{}{M}^A`$. (See the later arguments for the reducible case, if there exists any quaternionic Kähler manifold for this case.) Put (36) in the covariant form
$`R^{Ab}_bR_c^BR^{Bb}_bR_c^A=f^{ABC}R_c^C.`$ (45)
Contract this equation with $`R_a^AR_b^B`$. Using the formula (44) and its consequence $`_c(R_a^A`$ $`R_b^A)=0`$ we then get (42).
Now we proceed to derive (19)$``$(More on the Triplet Killing Potentials of Quaternionic Kähler manifolds). We calculate the following the Lie-derivative
$`_{R^A}_cR^{Bb}=R^{Aa}_a_cR^{Bb}_aR^{Ab}_cR^{Ba}+_cR^{Aa}_aR^{Bb}.`$ (46)
On the other hand we have
$`_cR^{Aa}_aR^{Bb}+R^{Aa}_c_aR^{Bb}(AB)=f^{ABC}_cR^{Bb}.`$
from (36). Eliminating the last two terms in (46) by this formula and using (40) and (43) yields
$`_{R^A}_cR^{Bb}`$ $`=`$ $`R^{Aa}_{[a}_{c]}R^{Bb}+R^{Ba}_c_aR^{Ab}+f^{ABC}_cR^{Bb}`$
$`=`$ $`f^{ABC}_cR^{Bb}.`$
With this and (38) the triplet Killing potentials (3) satisfy the property (19). Note that
$`\stackrel{}{M}^A={\displaystyle \frac{1}{2n}}\stackrel{}{J}^{ab}_aR_b^A={\displaystyle \frac{1}{2n}}_a(\stackrel{}{J}^{ab}R_b^A)={\displaystyle \frac{1}{2n\nu }}_b^b\stackrel{}{M}^A,`$ (47)
by (1), (3) and (29). This implies that $`\stackrel{}{M}^A`$ is an eigen vector of the Beltrami-Laplace operator, i.e. (20). Note also that
$`_a\stackrel{}{M}^A_b\stackrel{}{M}^A=\nu ^2\stackrel{}{J}_a^cR_c^A\stackrel{}{J}_b^dR_d^A,`$ (48)
by (1). Calculate the r.h.s. with (24). Then (48) becomes (21) owing to (44). Eq. (22) can be shown by manipulating the formula
$`\stackrel{}{J}^{ab}\stackrel{}{J}^{cd}R_{abcd}=\stackrel{}{J}^{ab}\stackrel{}{J}^{cd}R_d^A_c_aR_b^A,`$ (49)
which is obvious from (40). By (1) and (3) the r.h.s. becomes
$`{\displaystyle \frac{\nu }{2n}}\stackrel{}{J}^{ab}\stackrel{}{J}^{cd}R_d^A_c_aR_b^A`$ $`=`$ $`^c\stackrel{}{M}^A_c\stackrel{}{M}^A`$ (50)
$`=`$ $`^c[_c({\displaystyle \frac{1}{2}}\stackrel{}{M}^A\stackrel{}{M}^A)]\stackrel{}{M}^A^c_c\stackrel{}{M}^A.`$
From the property (19) we have $`_{R^B}(\stackrel{}{M}^A\stackrel{}{M}^A)=0`$. Contracting this equation with $`R^{Bb}`$ and using (44) yields $`g^{ba}_a(\stackrel{}{M}^A\stackrel{}{M}^A)=0`$ . If the metric is non-degenerate, we get $`\stackrel{}{M}^A\stackrel{}{M}^A`$ to be a constant. Then (50) becomes
$`\stackrel{}{J}^{ab}\stackrel{}{J}^{cd}R_d^A_c_aR_b^A=4n^2\stackrel{}{M}^A\stackrel{}{M}^A,`$ (51)
by (47). On the other hand the l.h.s. of (49) is calculated as
$$\stackrel{}{J}^{ab}\stackrel{}{J}^{cd}R_{abcd}=2n\stackrel{}{J}^{ab}\stackrel{}{R}_{ab}=\frac{6n}{n+2}R,$$
by (24), (30) and (35). With this and (51) the formula (49) becomes (22). Finally we show (More on the Triplet Killing Potentials of Quaternionic Kähler manifolds). It consists of the two equalities. To show the equality in the first line we rewrite formula
$`_{[a}_{b]}\stackrel{}{J}_{cd}+_{[c}_{d]}\stackrel{}{J}_{ab}=0,`$
which is obvious from (29) as
$`R_{abce}\stackrel{}{J}_d^e+R_{abed}\stackrel{}{J}_c^e+R_{aecd}\stackrel{}{J}_b^e+R_{ebcd}\stackrel{}{J}_a^e=0,`$ (52)
by (31) and (35). Multiply both sides of (52) by $`\stackrel{}{J}_f^d`$. Using (24) and (40) we get
$`3R_{abcf}=R_d^A_e_aR_b^A\stackrel{}{J}_c^e\stackrel{}{J}_f^d+R_d^A_c_aR_e^A\stackrel{}{J}_b^e\stackrel{}{J}_f^d+R_d^A_c_eR_b^A\stackrel{}{J}_a^e\stackrel{}{J}_f^d.`$
This becomes
$`3R_{abcf}=R_d^A_{[a}_{b]}R_e^A\stackrel{}{J}_c^e\stackrel{}{J}_f^d+R_d^A_c_aR_e^A\stackrel{}{J}_b^e\stackrel{}{J}_f^dR_d^A_c_bR_e^A\stackrel{}{J}_a^e\stackrel{}{J}_f^d,`$
by (39 ) and (41). Using (1) in the r.h.s. of the last equation leads us to see that the equality in the first line of (More on the Triplet Killing Potentials of Quaternionic Kähler manifolds) indeed holds. To show the second equality in (More on the Triplet Killing Potentials of Quaternionic Kähler manifolds) we write the Riemann curvature tensor in the form
$`R_{abcd}=_c(R_d^A_aR_b^A)+(_cR_d^A)(_aR_b^A),`$ (53)
by (40). Calculate the first term of the r.h.s. by (42). We then see that $`R_{abcd}`$ is expressed in terms of $`_bR_c^A`$ alone. Note the formula
$`_bR_c^A=f^{ABC}R_b^BR_c^C+{\displaystyle \frac{1}{2}}R^{Aa}f_{abc},f_{abc}f^{ABC}R_a^AR_b^BR_c^C,`$
which can be shown by contracting (45) with $`R_b^B`$ and using (42) and (44). By this formula and the Jacobi identity for the structure constants $`f^{ABC}`$, the Riemann curvature tensor (53) gets expressed as given in (More on the Triplet Killing Potentials of Quaternionic Kähler manifolds).
We illustrate our general arguments for quaternionic Kähler group manifolds, taking the coset space $`Sp(n+1)/Sp(n)Sp(1)`$ as an example. The isometry group $`Sp(n+1)`$ is generated by operators $`X^{IJ}(=X^{JI}),I,J=1,2,\mathrm{},2(n+1)`$, satisfying the Lie-algebra
$`[X^{IJ},X^{KL}]=ϵ^{IK}X^{JL}+ϵ^{JL}X^{IK}+ϵ^{IL}X^{JK}+ϵ^{JK}X^{IL}.`$
Here $`ϵ^{IJ}`$ is a constant anti-symmetric tensor, which we take to be
$`ϵ^{IJ}=ϵ^{JI}=ϵ_{IJ}=ϵ_{JI}=\left(\begin{array}{ccccc}0& 1& & & \text{0}\\ 1& 0& & & \\ & & \mathrm{}& & \\ & & & 0& 1\\ \text{0}& & & 1& 0\end{array}\right),`$ (59)
for convenience. The generators are decomposed as
$$\{X^{IJ}\}=\{X^{i\mu },X^{\mu \nu },X^{ij}\},\mu ,\nu =1,2,\mathrm{},2n,i,j=1,2,$$
in which $`X^{\mu \nu }`$ and $`X^{ij}`$ are generators of the homogeneous group $`Sp(n+1)Sp(1)`$ and $`X^{i\mu }(=X^{\mu i})`$ are broken generators. The quadratic Casimir takes the form
$`{\displaystyle \frac{1}{2}}X^{IJ}X^{KL}ϵ_{IK}ϵ_{JL}=X^{i\mu }X^{j\nu }ϵ_{ij}ϵ_{\mu \nu }+{\displaystyle \frac{1}{2}}X^{\mu \nu }X^{\rho \sigma }ϵ_{\mu \rho }ϵ_{\nu \sigma }+{\displaystyle \frac{1}{2}}X^{ij}X^{kl}ϵ_{ik}ϵ_{jl}.`$
The coset space $`Sp(n+1)/Sp(n)Sp(1)`$ is locally parametrized by coordinates $`\varphi ^{𝒊𝝁}`$ corresponding to the broken generators $`X^{i\mu }`$. The line element is given by
$$ds^2=g_{𝒊𝝁\text{,}𝒋𝝂}d\varphi ^{𝒊𝝁}d\varphi ^{𝒋𝝂}.$$
We find the Killing vectors as non-linear realization of the Lie-algebra (36):
$`R^{(k\sigma )𝒊𝝁}`$ $``$ $`[X^{k\sigma },\varphi ^{𝒊𝝁\text{ }}]=ϵ^{k𝒊}ϵ^{\sigma 𝝁}+\varphi ^{k𝝁}\varphi ^{𝒊\sigma \text{ }}\text{,}`$
$`R^{(\rho \sigma )𝒊𝝁}`$ $``$ $`[X^{\rho \sigma },\varphi ^{𝒊𝝁\text{ }}]=ϵ^{\rho 𝝁}\varphi ^{𝒊\sigma }+ϵ^{\sigma 𝝁}\varphi ^{𝒊\rho }\text{,}`$ (60)
$`R^{(kl)𝒊𝝁}`$ $``$ $`[X^{kl},\varphi ^{𝒊𝝁\text{ }}]=ϵ^{k𝒊}\varphi ^{l𝝁}+ϵ^{l𝒊}\varphi ^{k𝝁}\text{.}`$
Then the metric is obtained from (44)
$`g^{𝒊𝝁\text{,}𝒋𝝂}`$ $`=`$ $`[ϵ(1ϵ\varphi ϵ\varphi )]^{𝒊𝒋}[ϵ(1ϵ\varphi ϵ\varphi )]^{𝝁𝝂}\text{,}`$
$`g_{𝒊𝝁\text{,}𝒋𝝂}`$ $`=`$ $`[(1ϵ\varphi ϵ\varphi )^1ϵ]_{𝒊𝒋}[(1ϵ\varphi ϵ\varphi )^1ϵ]_{𝝁𝝂}\text{.}`$
Here one should undestand matrix multiplication such that
$`(ϵ\varphi )_𝒊^{\text{ }𝝁}=ϵ_{𝒊𝒌}\varphi ^{𝒌𝝁}\text{,}`$ $`(ϵ\varphi )_𝝁^{\text{ }𝒊}=ϵ_{𝝁𝝆}\varphi ^{𝒊𝝁}\text{,}`$
$`(ϵ\varphi ϵ\varphi )_𝒊^{\text{ }𝒋}=ϵ_{𝒊𝒌}\varphi ^{𝒌𝝆}ϵ_{𝝆𝝈}\varphi ^{𝒋𝝈}\text{,}`$ $`(ϵ\varphi ϵ\varphi )_𝝁^{\text{ }𝝂}=ϵ_{𝝁𝝆}\varphi ^{𝒌𝝆}ϵ_{𝒌𝒍}\varphi ^{𝒍𝝂}\text{.}`$
We calculate the Affine connection as
$`\mathrm{\Gamma }_{𝒊𝝁\text{,}𝒋𝝂}^{\text{ }𝒌𝝀}=\delta _𝒋^𝒌\delta _𝝁^𝝀[(1ϵ\varphi ϵ\varphi )ϵ\varphi ϵ]_{𝒊𝝂}+\delta _𝒊^𝒌\delta _𝝂^𝝀[(1ϵ\varphi ϵ\varphi )ϵ\varphi ϵ]_{𝒋𝝁}\text{.}`$ (61)
The vielbeins are given by
$`e_{𝒊𝝁}^{i\mu }=[(1ϵ\varphi ϵ\varphi )^{\frac{1}{2}}]_𝒊^i[(1ϵ\varphi ϵ\varphi )^{\frac{1}{2}}]_𝝁^\mu \text{,}`$
so as to satisfy
$$g_{𝒊𝝁\text{,}𝒋𝝂}=e_{𝒊𝝁}^{i\mu }e_{𝒋𝝂}^{j\nu }ϵ_{ij}ϵ_{\mu \nu }\text{.}$$
Using the formula
$$(\varphi ϵ\varphi )^{𝒊𝒋}=\varphi ^{𝒊𝝆}ϵ_{𝝆𝝈}\varphi ^{𝒋𝝈}=\frac{1}{2}ϵ^{𝒊𝒋}\varphi ^2\text{,}$$
with $`\varphi ^2=ϵ_{𝒋𝒊}\varphi ^{𝒊𝝆}ϵ_{𝝆𝝈}\varphi ^{𝒋𝝈}`$, we rewrite the vielbeins in the form
$$e_{𝒊𝝁}^{i\mu }=\delta _𝒊^i[(1ϵ\varphi ϵ\varphi )]_𝝁^\mu (1\frac{1}{2}\varphi ^2)^{\frac{1}{2}}\text{.}$$
Then the triplet complex structure (25) becomes
$`\stackrel{}{J}_{𝒊𝝁}^{\text{ }𝒋𝝂}=i\stackrel{}{\sigma }_𝒊^{\text{ }𝒋}\delta _𝝁^𝝂\text{.}`$ (62)
Using (3) with (60), (61) and (62) we calculate the triplet Killing potentials to find
$`\stackrel{}{M}^{k\sigma }=2i{\displaystyle \frac{(\varphi \stackrel{}{\sigma })^{\rho k}}{1\frac{1}{2}\varphi ^2}},\stackrel{}{M}^{\rho \sigma }=2i{\displaystyle \frac{(\varphi \stackrel{}{\sigma }ϵ\varphi )^{\rho \sigma }}{1\frac{1}{2}\varphi ^2}},\stackrel{}{M}^{kl}=2i{\displaystyle \frac{(ϵ\stackrel{}{\sigma })^{kl}}{1\frac{1}{2}\varphi ^2}}.`$
The Casimir product of $`\stackrel{}{M}^A`$ is indeed constant:
$$\stackrel{}{M}^{IJ}\stackrel{}{M}^{KL}ϵ_{IK}ϵ_{JL}=12.$$
It is easy to check that the triplet Killing potentials have the property (19) with
$$\stackrel{}{r}^{k\sigma }=i(\varphi \stackrel{}{\sigma })^{k\sigma },\stackrel{}{r}^{\rho \sigma }=0,\stackrel{}{r}^{kl}=2i(ϵ\stackrel{}{\sigma })^{kl}.$$
On the other hand these $`\stackrel{}{r}`$’s may be obtained from (38) as well. Of course they coincide with each other.
The derivations of (21)$``$(More on the Triplet Killing Potentials of Quaternionic Kähler manifolds) relied on the metric formula (44). This formula holds at least for a class of quaternionic Kähler coset spaces, called the Wolf space, for which $`Sp(n+1)/Sp(n)Sp(1)`$ is the simplest example. To verify this we use the CCWZ formalism. The Wolf space is a compact coset space of type $`G/SSp(1)`$ or $`G/SSU(2)`$ with particularly chosen $`G`$ and $`S`$. The generators of $`G`$, denoted by $`T^A`$, are decomposed as
$`\{T^A\}`$ $`=`$ $`\{X^{i\mu },S^I,Q^\alpha \},`$
$`i=1,2,\mu =1,2,`$ $`\mathrm{}`$ $`2n,I=1,2,\mathrm{},\mathrm{dim}K,\alpha =1,2,3.`$
Here $`S^I`$ are generators of the homogeneous group $`S`$, while $`Q^\alpha `$ those of $`Sp(1)`$ or $`SU(2)`$ depending on the type. $`X^{i\mu }`$ are broken generators and transform as a tensor in the representation $`(\underset{¯}{2},\underset{¯}{2n})`$ under the homogeneous group $`Sp(1)S`$ or $`SU(2)S`$. In general it is said that the coset space is irreducible(reducible) when the representation of $`X^{i\mu }`$ is irreducible(reducible) under $`S`$. The Wolf space is irreducible. The quadratic Casimir of $`G`$ is given by
$`T^AT^A=X^{i\mu }X^{j\nu }\eta _{i\mu ,j\nu }+S^IS^I+Q^\alpha Q^\alpha ,`$
in which $`\eta _{i\mu ,j\nu }`$ is the Killing form for the coset part, i.e., $`\eta _{i\mu ,j\nu }=ϵ_{ij}ϵ_{\mu \nu }`$ for the previous example $`Sp(n+1)/Sp(n)Sp(1)`$. The coset space is locally parametrized by $`4n`$ real coordinates $`\varphi ^{𝒊𝝁}`$ corresponding to the broken generators $`X^{i\mu }`$. Cosider the quantity $`U=\mathrm{exp}(\varphi ^{𝒊𝝁}X^{j\nu }\eta _{𝒊𝝁\text{,}j\nu })`$. By left multiplication of $`\mathrm{exp}(ϵ^AT^A)G`$ we find
$`\mathrm{exp}(ϵ^AT^A)U(\varphi )=U(\varphi ^{})h(\varphi ,g),`$ (63)
appropriately choosing the conpensator $`h(\varphi ,g)Sp(1)S`$ or $`SU(2)S`$. Here $`ϵ^A`$ are global parameters. When they are infinitesimal (63) yields the Killing vectors $`R^{A𝒊𝝁}(\varphi )`$ as
$$\delta \varphi ^{𝒊𝝁}=\varphi _{}^{}{}_{}{}^{𝒊𝝁}\varphi ^{𝒊𝝁}=ϵ^AR^{A𝒊𝝁}(\varphi )\text{.}$$
The fundamental object to construct the metric is the Cartan-Maurer $`1`$-form $`U^1dU`$. It is valued in the Lie-algebra of $`G`$ as
$$U^1dU=e^{i\mu }X^{i\mu }\eta _{i\mu ,j\nu }+\omega ^IS^I+\omega ^\alpha Q^\alpha .$$
$`e^{i\mu }`$ is a vielbein $`1`$-form, while $`\omega ^I`$ and $`\omega ^\alpha `$ connection $`1`$-forms corresponding to the respective holonomy groups $`S`$ and $`Sp(1)`$. With this vielbein $`1`$-form the metric $`g_{𝒊𝝁\text{,}𝒋𝝂}`$ is given by
$`g_{𝒊𝝁\text{,}𝒋𝝂}=e_{𝒊𝝁}^{i\mu }e_{𝒋𝝂}^{j\nu }\eta _{i\mu ,j\nu }\text{.}`$ (64)
It is also given by (44) i.e., $`g_{𝒊𝝁\text{,}𝒋𝝂}=R_{𝒊𝝁}^AR_{𝒋𝝂}^A`$. Both metrics are equivalent because the value at the origin $`g_{𝒊𝝁\text{,}𝒋𝝂}|_{\varphi =0}`$ and the Lie-variation with respect to the Killing vectors $`_{R^A}g_{𝒊𝝁\text{,}𝒋𝝂}`$ are the same. Equivalently we can say that
$`e_{𝒊𝝁}^{i\mu }=R_{𝒊𝝁}^AU^{A,i\mu }\text{,}`$ (65)
in which $`U^{A,i\mu }`$ are matrix elements of $`U`$ in the adjoint representation of $`G`$. The point is that the Wolf space is irreducible. For reducible coset spaces the metric formula (64) should be generalized as
$`g_{𝒊𝝁\text{,}𝒋𝝂}=_{\mu =1}^Nc_\mu e_{𝒊𝝁}^{i\mu }e_{𝒋𝝂}^{j\nu }\eta _{i\mu ,j\nu }\text{,}`$ (66)
in which $`c_\mu `$ may take different values for each irreducible component of $`X^{i\mu }`$ such as $`c_\mu =(c_1,\mathrm{},c_1,c_2,\mathrm{},c_2,c_3\mathrm{}\mathrm{},c_N,\mathrm{},c_N)`$ with $`_{\mu =1}^N1=2n`$. Accordingly the formula (44) is generalized to the one obtained by putting the vielbeins (65) in (66). Presumably quaternionic Kähler manifolds could exist also in such reducible cases similarly to the ordinary Kähler manifold.
In this letter quaternionic Kähler manifolds were studied in view of an explicit construction of the metric. For the triplet Killing potentials we have shown the properties (19)$``$(More on the Triplet Killing Potentials of Quaternionic Kähler manifolds) which have been missing in the literature. Among them the properties (21)$``$(More on the Triplet Killing Potentials of Quaternionic Kähler manifolds) were derived with recourse to the metric formula (44). It seems that those properties have been overlooked in generality of the arguments with no specification of the metric. The metric can not be determined by the Killing equation (39) alone. One needs the initial condition at the origin of the manifold. By taking account of it the metric formula (44) was justified for the irreducible coset space. It is worth studying also the reducible case in the constructive approach of this letter.
Acknowledgements
I thank R. D’Auria for the hospitality at Politecnico di Torino. I also thank him and A. Van Proeyen for the discussions during the stay through which this work was started. The work was supported in part by the Grant-in-Aid for Scientific Research No. 13135212.
|
warning/0506/cond-mat0506060.html
|
ar5iv
|
text
|
# Friction and diffusion of matter-wave bright solitons
## Abstract
We consider the motion of a matter-wave bright soliton under the influence of a cloud of thermal particles. In the ideal one-dimensional system, the scattering process of the quasiparticles with the soliton is reflectionless, however, the quasiparticles acquire a phase shift. In the realistic system of a Bose-Einstein condensate confined in a tight waveguide trap, the transverse degrees of freedom generate an extra but small nonlinearity in the system which gives rise to finite reflection and leads to dissipative motion of the soliton. We calculate the velocity and temperature-dependent frictional force and diffusion coefficient of a matter wave bright soliton immersed in a thermal cloud.
Solitons are localized waves that propagate without spreading and attenuation. They appear from classical systems like ocean waves to optics and quantum systems like Bose-Einstein condensates (BEC) of atomic gases. A BEC of a dilute atomic gas with attractive two-body interactions in three dimensions (3D) is unstable and collapses gerton00 . In one dimension (1D), however, a BEC with attractive interaction is stable against collapse and forms a self-bound particle-like object known as a bright soliton. Recently, bright solitons of Bose-condensed <sup>7</sup>Li atoms were observed in quasi-1D waveguide traps at Rice University Strecker2002a and at ENS in Paris Khaykovich2002a .
One of the most important features of solitons is the dissipationless motion over long distances. Because of this property, optical solitons have important applications in transatlantic fiber-optic communication systems hasegawa02 . In this Letter we discuss how dissipative effects in the motion of a soliton in a thermal cloud can arise due to the 3D nature of the BEC in a tight waveguide. Although the dynamics is strictly one-dimensional, the transverse extent of the mean field generates extra nonlinear terms in the effective 1D equation, as first discussed in Ref. Muryshev2002a . This formalism can also be applied to other systems described by a cubic nonlinear Schrödinger equation perturbed by a quintic nonlinear term. An overview of related work in nonlinear optics can be found in Ref. kivshar03:book .
In this Letter we consider the scattering of quasiparticles by the quasi-1D matter-wave soliton Wynveen2000a similar to linear waves scattering on breathers flach:084101 . Muryshev et al. considered the interaction of quasiparticles with dark BEC solitons and conjectured that these lead to acceleration and eventually disintegration of the soliton in a thermal environment Muryshev2002a . Following a similar line of arguments, we show that the quasiparticles scattering on a bright soliton have a finite probability of reflection only due to the extra nonlinearity, which finally gives rise to dissipative effects. The bright soliton experiences friction and diffusive motion in a thermal cloud but maintains its integrity in contrast to dark solitons which disintegrate.
A BEC in a waveguide with a harmonic transverse confinement is well described by the Gross-Pitaevskii (GP) equation:
$$i\mathrm{}\frac{\psi }{t}=\left[\frac{\mathrm{}^2}{2m}^2+\frac{1}{2}m\omega ^2\rho ^2+\frac{4\pi \mathrm{}^2a}{m}|\psi |^2\right]\psi $$
(1)
where $`\psi `$ is the macroscopic wavefunction of the condensate, $`\omega `$ is the frequency of transverse trapping potential, and $`a`$ is the 3D scattering length. In the quasi-1D limit the effective dynamics of the system takes place along the free axis (x-axis) without exciting the transverse modes. The quasi-1D limit can be achieved when the mean field interaction is smaller than the radial excitation frequency, $`4\pi \mathrm{}^2|a||\psi |^2/m<\mathrm{}\omega `$. Aiming at an adiabatic separation of slow longitudinal and fast transverse motion we can write the full 3D wave function assuming cylindrical symmetry as, $`\psi (\stackrel{}{r},t)=\varphi (x,t)\chi (\rho ,x,t)`$. Here, $`\varphi `$ is the 1D (longitudinal) wavefunction and $`\chi `$ is the radial wavefunction with the normalisation convention $`|\chi |^22\pi \rho 𝑑\rho =1`$ and $`|\varphi |^2𝑑x=N`$, where $`N`$ is the number of bosons in the system. In the adiabatic or Born-Oppenheimer approximation we now assume that the radial wavefunction $`\chi `$ depends only weakly on the slow variables $`x`$ and $`t`$ and their derivatives of $`\chi `$ can be neglected. For this assumption to be correct, the time scale of dynamics should be longer than the inverse of the tranverse frequency $`\omega `$ and the relevant length scale should be significantly larger than the transverse oscillator length $`l=\sqrt{\mathrm{}/(m\omega )}`$. After substituting the ansatz for the wavefunction into Eq. (1) and neglecting the derivatives of $`\chi `$ with respect to $`x`$ and $`t`$, we obtain the following adiabatically decoupled equations for the longitudinal and the transverse wavefunctions:
$`i\mathrm{}_t\varphi `$ $`={\displaystyle \frac{\mathrm{}^2}{2m}}_x^2\varphi +\stackrel{~}{\mu }\varphi ,`$ (2)
$`[{\displaystyle \frac{\mathrm{}^2}{2m}}_\rho ^2+`$ $`{\displaystyle \frac{1}{2}}m\omega ^2\rho ^2+{\displaystyle \frac{4\pi \mathrm{}^2a}{m}}n|\chi |^2]\chi =\stackrel{~}{\mu }\chi ,`$ (3)
where we have introduced the transverse chemical potential $`\stackrel{~}{\mu }`$, which has to be found from the ground state solution of Eq. (3) as a function of the linear density $`n(x,t)=|\varphi (x,t)|^2`$. A simple scaling argument shows that $`\stackrel{~}{\mu }=\mathrm{}\omega f(an)`$, where $`f()`$ is a dimensionless function, which has been computed numerically in Ref. Berge2000a . Physical solutions of Eq. (3) are found only if $`an<0.47`$ weinstein83 ; Berge2000a , otherwise transverse collapse occurs carr:040401 . In the following we will be interested in the quasi-1D regime of small $`an`$ and expand $`\stackrel{~}{\mu }(an)`$ in a power series.
In the quasi-1D limit, when $`|a|n0.47`$, the radial wavefunction $`\chi `$ will be close to the ground state of the 2D harmonic oscillator with a Gaussian profile. We can expand $`\chi `$ in terms of the radial eigenmodes $`\phi _\nu (\rho )`$, $`\chi (\rho ,x)=\phi _0(\rho )+_\nu C_\nu (x)\phi _\nu (\rho )`$. The coefficients $`C_\nu `$ are small and can be calculated perturbatively. The transverse chemical potential $`\stackrel{~}{\mu }`$ can be obtained by using second order perturbation theory:
$$\stackrel{~}{\mu }=\mathrm{}\omega +gng_2n^2+\mathrm{},$$
(4)
where $`g=2a\mathrm{}\omega `$ and $`g_2=24\mathrm{ln}(4/3)a^2\mathrm{}\omega `$. A correction to the 1D coupling constant $`g`$ beyond the GP approach presented here has been found in Ref. Olshanii1998a . The constant $`g_2`$ was calculated first in Ref. Muryshev2002a and corrections beyond GP can be obtained by the self-consistent Hartree Fock Bogoliubov approach of Ref. cherny:043622 . We obtain the following effective equation describing the condensate in the quasi-1D limit:
$$[(\mathrm{}^2/2m)_x^2+g|\varphi |^2g_2|\varphi |^4]\varphi =\mu \varphi .$$
(5)
This is a nonlinear Schrödinger equation with a cubic and a quintic nonlinearity, as used before in Ref. Muryshev2002a . The possibility of collapse is inherent in this equation as the quintic nonlinear term is attractive. An estimate from the 3D GP equation (1) gives stability of a single soliton solution if $`N|a|/l<0.627`$ is fulfilled carr02 .
Without the extra nonlinearity associated with $`g_2`$, Eq. (5) is integrable. For attractive interactions at $`a<0`$, the bosons form a self-bound particle-like state known as a bright soliton with the wavefunction $`\varphi (x)=\sqrt{N/2b}\text{sech}(x/b)`$ and the chemical potential $`\mu =\mathrm{}^2/2mb^2`$, where $`b=l^2/(N|a|)`$. We notice that for a weak soliton parameter $`N|a|/l1`$, the system becomes quasi-one-dimensional ($`bl`$).
A soliton can be considered as a macroscopic particle of mass $`mN`$, moving in the bath of thermal excitations. Dissipative motion of the soliton arises due to the scattering of thermal atoms. Here we consider the interaction of thermally excited particles with the soliton within the Bogoliubov formalism Wynveen2000a ,
$$[H_0+H_1]\psi =ϵ(k)\psi $$
(6)
where, $`\psi =(u,v)`$ is a two component vector of particle ($`u`$) and hole ($`v`$) amplitudes, and $`ϵ`$ is the quasiparticle energy. The unperturbed Hamiltonian $`H_0`$ and the perturbation $`H_1`$ are given by
$`\widehat{H}_0`$ $`=`$ $`\left(\begin{array}{cc}\frac{\mathrm{}^2}{2m}_x^2\mu & 0\\ 0& +\frac{\mathrm{}^2}{2m}_x^2+\mu \end{array}\right)`$ (9)
$`\widehat{H}_1`$ $`=`$ $`\left(\begin{array}{cc}V_1(x)& V_2(x)\\ V_2^{}(x)& V_1(x)\end{array}\right),`$ (12)
where $`V_1=2g|\varphi |^23g_2|\varphi |^4`$ and $`V_2=g\varphi ^22g_2|\varphi |^2\varphi ^2`$. The scattering states have energy $`ϵ(k)=\frac{\mathrm{}^2k^2}{2m}+|\mu |`$. In one dimension, neglecting the extra nonlinearity ($`g_2=0`$), we obtain the exact solution of the scattering states:
$`u_k=`$ $`A(k)[kb+i\mathrm{tanh}(x/b)]^2e^{ikx}`$ (13)
$`v_k=`$ $`A(k)\text{sech}^2(x/b)e^{ikx},`$ (14)
where $`A(k)=1/(k^2b^21)`$ is a normalisation constant. The transmittance is given by
$$t=(kb+i)^2/(kbi)^2,$$
(15)
and the transmission probability is $`|t|^2=1`$. Hence, the quasiparticles scatter without reflection on the soliton but only aquire a phase shift and a time delay in the scattering process. Reflectionless scattering on a soliton in the integrable nonlinear Schrödinger equation (5) (with $`g_2=0`$) is a well-known result of mathematical soliton theory and is also found in an exact solution of the quantum many-body model in the limit of large particle number mcguire:622 . In the quasi-1D limit, the soliton thus becomes transparent and exhibits dissipationless motion in a thermal cloud.
Now we consider the scattering problem of quasiparticles in the presence of an extra nonlinearity that breaks the integrability. Assuming that the coupling constant $`g_2`$ of the extra nonlinear term is small, we can solve the scattering problem using Green’s function techniques. In order to solve Eq. (6) for the particle amplitude $`u`$ we construct the Green’s function for the $`u`$ component of $`H_0`$ satisfying
$$[(\mathrm{}^2/2m)_x^2\mu ϵ]G_1(xx^{})=\delta (xx^{}),$$
(16)
which is given by $`G_1(xx^{})=(m/\mathrm{}^2k)\mathrm{sin}(k|xx^{}|)`$. Since the potential is symmetric, the scattering states can be constructed with even or odd symmetry. The Lippmann-Schwinger equation for the particle channel can be written as
$`u_{e/o}`$ $`=`$ $`u_{e/o}^0+{\displaystyle G_1(xx^{})V_1(x^{})u_{e/o}(x^{})𝑑x^{}}`$ (17)
$`+{\displaystyle G_1(xx^{})V_2(x^{})v_{e/o}(x^{})𝑑x^{}},`$
where $`u_{e/o}`$ denotes even \[odd\] wave functions of the particle states and $`u_e^0=\mathrm{cos}(kx)`$, \[$`u_o^0=\mathrm{sin}(kx)`$\]. The most general wave function can be constructed from even and odd eigenstates: $`u_k=Au_k^e+Bu_k^o`$. Asymptotically this wave function becomes $`lim_x\mathrm{}u_k=e^{ikx}+re^{ikx}`$ and $`lim_x\mathrm{}u_k=te^{ikx}`$ where $`|t|^2`$ and $`R=|r|^2`$ are the transmission and the reflection coefficient, respectively. We obtain $`R(k)`$ by solving Eqs. (5, 6) numerically and matching with the asymptotic solutions, see Fig. 1.
An analytical estimate of the reflection coefficient can be obtained from Eq. (17) by approximating $`\varphi =\sqrt{N/2b}\text{sech}(x/b)`$ and $`u_{e/o}`$ and $`v_{e/o}`$ with the properly symmetrised solutions (13, 14). This approximation becomes exact for $`g_2=0`$ and relies on $`g_2`$ being a small parameter. The reflection coefficient is given by $`R=|r|^2`$ and
$$r(k)=i\frac{I_++I_{}}{(I_+i)(I_{}+i)},$$
(18)
where the terms $`I_+`$ and $`I_{}`$ are given by
$`I_\pm =2A(k)\left[kb+6\mathrm{ln}(4/3){\displaystyle \frac{N^2|a|^2}{l^2}}{\displaystyle \frac{Q_\pm (kb)}{kb}}\right]`$
with $`Q_\pm (x)=1/3+x^2\pm (1+x^2)^2\pi x/[3\mathrm{sinh}(\pi x)]`$. By using the Lippmann-Schwinger formalism instead of a simple Born approximation we obtain the correct limiting behaviour for small $`k`$ where $`R1`$. Total reflection is expected whenever the special resonant conditions leading to reflectionless scattering at $`g_2=0`$ are broken, as $`k0`$ implies a vanishing group velocity $`ϵ/(\mathrm{}k)=\mathrm{}k/m`$. This case is very different from phonons scattering on a perturbed dark soliton, which becomes transparent for small $`k`$ as found in Ref. Muryshev2002a . The approximation (18) reproduces the qualitative features but slightly overestimates the exact values of $`R`$ as seen in Fig. 1.
The reflection coefficient $`R`$ is a function of dimensionless momentum $`kb`$ and the soliton parameter $`N|a|/l`$. In the dissipative dynamics of a macroscopic object like a soliton, the microscopic parameter $`N|a|/l`$ enters through the reflection coefficient of the quasiparticles. Once we know the interaction of particles with a soliton from the microscopic theory, we can describe its motion in the bath of thermal particles at a given temperature. A bright soliton is a mesoscopic object with mass $`mN`$, and its dynamics is governed by classical motion. Therefore, we can define a phase space distribution function of soliton’s center of mass coordinate $`f(p,q,t)`$. When the soliton follows the classical trajectory then the distribution function takes a simple form $`f(p,q,t)=\delta (pp(t))\delta (qq(t))`$, where $`p(t),q(t)`$ are classical phase space trajectories. In the presence of a bath of thermal atoms, the atoms impart a momentum to the soliton in the scattering process. While the soliton is at rest, the force imparted on the soliton cancels on the average but, nevertheless, the stochastic nature of the force introduces a diffusive motion of the soliton. For a moving soliton, the average force imparted by the thermal particles does not vanish and gives rise to a frictional force on the soliton. To include the dissipative effects in the soliton’s motion we write down the kinetic equation for the phase space distribution function of the soliton landau81:kinetics :
$$\frac{f}{t}\frac{}{p}\left(\frac{H}{q}f\right)+\frac{}{q}\left(\frac{H}{p}f\right)=I_{\mathrm{coll}},$$
(19)
where, for small momentum transfer, the collision integral $`I_{\mathrm{coll}}`$ can be written as
$$I_{\mathrm{coll}}=\frac{}{p}\left[Af+\frac{}{p}(Bf)\right].$$
(20)
The terms $`A`$ and $`B`$ gives rise to friction and diffusion of the soliton respectively. The frictional force $`A`$ can be computed from the following expression,
$$A=\frac{dk}{2\pi }(2\mathrm{}k)R(k)\left|\frac{ϵ(k)}{\mathrm{}k}\right|N(E,k_BT)$$
(21)
where $`N(E,k_BT)`$ describes the distribution of thermal particles in the frame of the moving soliton with velocity $`v`$ and the energy $`E`$ takes the value $`E(k)=(\mathrm{}kmv)^2/2m`$. In each collision, the particle with momentum $`k`$ has a probability $`R`$ to reflect back and transfer the momentum $`2\mathrm{}k`$ to the soliton. This momentum transfer multiplied with the number of particles coming from each direction per unit time gives rise to a frictional force. When the soliton is at rest, the momentum transfer on each direction cancels on the average and as a result the friction vanishes.
At finite temperatures the thermal atoms are distributed according to the rules of quantum statistics. Although thermal equilibrium may be reached in an external trap dunjko03ep:solitonThermo , the subtle conditions of equilibrium are not necessarily fulfilled in a dynamical experimental situation. Here, we consider the motion of a soliton relative to a significantly warmer thermal cloud of atoms. We thus can assume a classical Boltzmann distribution of thermal atoms, $`N(E,k_BT)\mathrm{exp}(E/k_BT)`$. Dissipative effects of the soliton can be enhanced by increasing the density or the temperature of the thermal cloud. We consider the situation where $`10^4`$ thermal particles are confined within a length $`L=50b`$ ($`70\mu \mathrm{m}`$ for ENS soliton with $`b=1.4\mu \mathrm{m}`$), with a density of the thermal gas of $`nb=200`$ ($`10^{12}/\mathrm{cm}^3`$) and the velocity distribution of the thermal particles being controlled by changing the temperature. Within a certain range of the soliton velocity the frictional force increases linearly with velocity as seen in Fig. 2. When the velocity is increased further, nonlinear effects take over and the force decreases.
Now we can calculate the diffusion parameter of the transport equation:
$$B=\frac{dk}{2\pi }2(\mathrm{}k)^2R(k)\left|\frac{ϵ(k)}{\mathrm{}k}\right|N(E,k_BT).$$
(22)
This term describes the velocity fluctuations of the soliton and gives rise to a diffusion in the momentum space. A graph is shown in Fig. 3.
So far we considered only the low-energy elastic scattering of thermal quasiparticles on the soliton. To restrict our discussion to the quasi-1D case, we neglected the higher energy radial excitations $`\mathrm{}\omega `$. As an additional effect, the soliton can radiate particles if the colliding quasiparticle has higher energy than the binding energy $`|\mu |`$. Also nonlinear collective motion of the thermal cloud and the soliton is possible buljan05ep . However, the elastic scattering process discussed in this work will dominate if the condition $`k_BT<|\mu |<\mathrm{}\omega `$ is fulfilled. A tight radial trapping potential is suitable to avoid inelastic scattering processes. In the ENS experiment Khaykovich2002a , the oscillator length of radial confinement was $`l=1.4\mu \mathrm{m}`$. For a soliton parameter $`N|a|/l=0.4`$, the sound velocity at the center of the soliton becomes $`c_\mathrm{s}2.5\mathrm{mm}/\mathrm{s}`$. If a soliton with $`N10^3`$ particles moves with a velocity $`0.1c_\mathrm{s}`$, then it decelerates $`5.01\mathrm{mm}/\mathrm{s}^2`$ due to the frictional force. Finally it stops after $`0.05\mathrm{s}`$, travelling a distance of $`6\mu \mathrm{m}`$. The slowing down of a bright soliton can be observed experimentally by suitably manipulating the density and temperature of the thermal cloud.
Due to the friction force, the momentum $`\overline{p}`$ of the soliton changes as $`d\overline{p}/dt=A(v)`$. For small velocities, $`A=\gamma v`$ and the moving soliton stops after a time scale $`\tau =mN/\gamma `$. Due to the diffusion process the energy of a resting soliton changes as $`E=(B/2\gamma )[1e^{2\gamma t/mN}]`$. For $`t\tau `$, the energy of the soliton increases and finally it reaches a steady state with energy $`E=B/2\gamma `$.
In conclusion, we have investigated the effects of a thermal environment on the dynamics of bright matter-wave solitons and have calculated the frictional force and diffusion coefficient in a microscopic approach. Friction and diffusion effects occur due to the deviation from the quasi-one-dimensional limit. Both of them are generally small and can be controlled by the parameters of the system if unattenuated propagation of solitons is desired. However, the parameters can be chosen such that the dissipative effects become accessible to experimental observation with currently available techniques.
We acknowledge enlightening discussions with G. Shlyapnikov who suggested this problem.
|
warning/0506/quant-ph0506172.html
|
ar5iv
|
text
|
# Controlled Flow of Spin-Entangled Electrons via Adiabatic Quantum Pumping
## Abstract
We propose a method to dynamically generate and control the flow of spin-entangled electrons, each belonging to a spin-singlet, by means of adiabatic quantum pumping. The pumping cycle functions by periodic time variation of localized two-body interactions. We develop a generalized approach to adiabatic quantum pumping as traditional methods based on scattering matrix in one dimension cannot be applied here. We specifically compute the flow of spin-entangled electrons within a Hubbard-like model of quantum dots, and discuss possible implementations and identify parameters that can be used to control the singlet flow.
Entanglement is one of the most intriguing features that distinguish the quantum world from the classical. In recent years Nielsen-Chuang applications potential in quantum information and computation has injected renewed vigor into the study of entangled states. Since Bohm’s Bohm reformulation of the Einstein-Podolsky-Rosen (EPR) paradox EPR in terms of pairing of spins, spin singlets have become the canonical example of an elementary entangled state. From an experimental standpoint, manipulating singlet pairs of electron spins, in particular, is promising because of the vast expertise already available in solid-state electronics Awschalom .
Several proposals have emerged Oliver ; Saraga-Loss ; Hu-Sarma ; cooper-pairs that aim to generate a controlled flow of electron spin singlets using Coulomb blockade and tunnelling at pinched quantum dots that have weak coupling to leads. In these proposals, the natural tunnelling of unwanted single electrons can be difficult to suppress and/or singlet delivery time is constrained by tunnelling rates that are slow due to the necessity for pinched dots. Since manipulations in a quantum information device need to occur before decoherence sets in, fast delivery is essential, and any background of unpaired electrons reduces the purity of entanglement. It is therefore highly desirable to have a fast dynamic mechanism to generate a selective flow only of electrons belonging to singlets. In this paper we show that this can be accomplished through a generalized adiabatic quantum pumping that is induced by localized interactions and we discuss possible ways of experimental realization. Pinched dot tunnelling times and singlet formation times do not constrain delivery rates in our proposal. In addition, the absence of bias and the adiabatic nature of the pumping can reduce the production of heat in the apparatus.
The notion of quantum pumping has its roots in a speculative paper by Thouless Thouless in 1983, but advances in nanoscale transport have led to a renewed and growing interest in the phenomenon in recent years both theoretically and experimentally Brouwer ; turnstile ; nano-transport ; Chamon ; Watson ; Floquet-Buttiker ; Floquet-Kim . Quantum pumping is a coherent process that creates a direct current in the *absence* of any bias through a nanoscale device, by changing its scattering properties periodically through independent adiabatic variation of two or more physical parameters. Adiabatic quantum pumping of charge Brouwer , spin Chamon ; Watson and thermal Floquet-Buttiker currents have been considered. However previous studies have generally relied on a theoretical description based on transmission and reflection coefficients which cannot be used to describe the pumping of singlets due to localized interactions; therefore in this paper we also formulate a more generalized approach to the theory of quantum pumping.
*Singlet Current*: We consider a single available channel in a quasi one dimensional mesoscopic conductor connected to macroscopic contacts (Fig. 1(a)). In the presence of two body interactions the charge current can be defined in terms of the two particle reduced density matrix, $`\rho _2`$. If the interaction does *not* affect spins, the reduced density matrix separates into four independent spin subspaces, one singlet with a symmetric spatial part and three triplets with antisymmetric spatial parts; the singlet current is therefore equivalent to the charge current associated with the symmetric spatial part ($`\rho _2^S`$) of the two particle reduced density matrix
$`J_S(x_1,t)={\displaystyle \frac{e\mathrm{}}{m}}{\displaystyle }dx_2\mathrm{}\left\{_{x_1}\rho _2^S(x_1,x_2;x_1^{},x_2;t)\right\}_{x_1=x_1^{}}.`$ (1)
We consider a situation in which two-body interactions are localized, meaning that they are non-vanishing only when the particles are in a certain finite interval, as shown in Fig. 1(a). Such interactions lend themselves to a scattering description and the current can be approximately evaluated by expanding the density matrix in terms of two particle scattering states
$`J_S(x_1,t){\displaystyle \frac{e\mathrm{}}{2m}}{\displaystyle }dEF(E){\displaystyle }{\displaystyle \frac{dk_1}{2\pi }}{\displaystyle }{\displaystyle \frac{dk_2}{2\pi }}\delta (\frac{\mathrm{}^2k_1^2}{2m}+\frac{\mathrm{}^2k_2^2}{2m}E){\displaystyle }dx_2\mathrm{}\left\{_{x_1}\mathrm{\Psi }_{k_1,k_2}(x_1,x_2,t)\mathrm{\Psi }_{k_1,k_2}^{}(x_1^{},x_2,t)\right\}_{x_1=x_1^{}}.`$ (2)
Here $`E`$ denotes the energy required to remove a pair of particles from the many-body ground state, and $`F(E)`$ is the distribution of this pair energy. The effect of interaction on the current is determined completely by the two-particle *singlet* scattering states, $`\mathrm{\Psi }_{k_1,k_2}`$ arising from *free* singlet states $`\mathrm{\Phi }_{k_1,k_2}(x_1,x_2)=\frac{1}{\sqrt{2}}[\varphi _{k_1}(x_1)\varphi _{k_2}(x_2)+\varphi _{k_1}(x_2)\varphi _{k_2}(x_1)]`$ where $`\varphi _k(x)`$ denotes a single particle plane wave state with momentum $`\mathrm{}k`$.
*Pumped Current through Adiabatic Perturbation* : The two particle scattering states, and therefore the current, are determined by the interaction $`V(\overline{x},t)`$ between a pair of particles; we take it to be time-dependent and to occur only in a *finite* region $`|x_i|<l`$. Most importantly the interaction $`V(\overline{x},t)`$ is chosen to be localized so that it *only affects singlets thereby naturally eliminating the flow of triplets* in the absence of a bias. When the characteristic period $`\omega `$ of the time variation of the potential is slow compared to the time $`\delta t`$ the particles dwell in the scattering region Buttiker-Landauer , $`\omega \times \delta t1`$, we can apply adiabatic perturbation theory to express the scattering states of the time-dependent Hamiltonian in terms of the instantaneous states up to linear order
$`\mathrm{\Psi }_{\overline{k}}(\overline{x},t)\mathrm{\Psi }_{\overline{k}}^t(\overline{x})i\mathrm{}{\displaystyle 𝑑\overline{x}^{}G^t(\overline{x},\overline{x}^{};E)\frac{}{t}\mathrm{\Psi }_{\overline{k}}^t(\overline{x}^{})}.`$ (3)
We use the notation $`\overline{x}\{x_1,x_2\}`$ and $`\overline{k}\{k_1,k_2\}`$, so that $`G^t(\overline{x},\overline{x}^{};E)`$ is the *two*-particle *instantaneous* retarded Green’s function for the *full* Hamiltonian. The instantaneous state $`\mathrm{\Psi }_{\overline{k}}^t(\overline{x})`$ is a solution of the *time-independent* Lippmann-Schwinger equation for the potential $`V(\overline{x}^{},t)`$ at the specific time $`t`$,
$`\mathrm{\Psi }_{\overline{k}}^t(\overline{x})`$ $`=`$ $`\mathrm{\Phi }_{\overline{k}}(\overline{x})+{\displaystyle 𝑑\overline{x}^{}G_0(\overline{x},\overline{x}^{};E)V(\overline{x}^{},t)\mathrm{\Psi }_{\overline{k}}(\overline{x}^{})}.`$ (4)
where $`G_0`$ is the *free two-particle retarded* Green’s function. Taking the time derivatives of the defining equations for $`\mathrm{\Psi }^t(\stackrel{}{x}^{})`$ and $`G^t(\overline{x},\overline{x}^{};E)`$ enables us to express the second term in Eq.(3), that is linear in $`_t`$, as
$`\mathrm{\Delta }\mathrm{\Psi }_{\overline{k}}(\overline{x},t)=i\mathrm{}{\displaystyle 𝑑\overline{x}^{}𝑑\overline{x}^{\prime \prime }G^t(\overline{x},\overline{x}^{};E)}`$ (5)
$`\times G^t(\overline{x}^{},\overline{x}^{\prime \prime };E)\dot{V}(\stackrel{}{x}^{\prime \prime },t)\mathrm{\Psi }_{\overline{k}}^t(\stackrel{}{x}^{\prime \prime }).`$
If there is no bias or time-dependence, the laws of thermodynamics demand that there should be no current; we explicitly confirm that our expression for the current satisfies this essential physical requirement. The net current in the absence of time-dependence is evaluated by using the zeroeth order term from Eq. (3) for the scattering state in Eq. (2): $`\mathrm{\Psi }_{k_1,k_2}(x_1,x_2,t)\mathrm{\Psi }_{k_1,k_2}^t(x_1,x_2)`$. The resulting expression can be simplified by relating the imaginary part of the retarded Green’s function to the free singlet states: $`\mathrm{}\{G_0(\overline{x},\overline{x}^{};E)\}=\pi 𝑑k_1𝑑k_2\delta \left(\frac{\mathrm{}^2k_1^2}{2m}+\frac{\mathrm{}^2k_2^2}{2m}E\right)\mathrm{\Phi }_{\overline{k}}(\overline{x})\mathrm{\Phi }_{\overline{k}}^{}(\overline{x}^{})`$. Then repeated use of the Lippmann-Schwinger equation and properties of the Green’s functions shows that the net current corresponding to the zeroeth order $`\mathrm{\Psi }^t(\overline{x})`$ vanishes.
After confirming that our expression cannot produce spontaneous current, we evaluate the singlet current induced by the adiabatic time evolution. To linear order in the time dependence this involves the evaluation of
$`{\displaystyle }dx_2\mathrm{}\left\{_{x_1}\mathrm{\Psi }_{\overline{k}}^t(x_1,x_2)\mathrm{\Delta }\mathrm{\Psi }_{\overline{k}}^{}(x_1^{},x_2,t)\right\}_{x_1=x_1^{}}`$ (6)
within the expression for the current in Eq. (2). A calculation employing standard Green’s function identities, similar to that for the zeroeth order, leads to an expression for the net amount of singlet entangled electron pairs pumped in a complete cycle of period $`\tau `$, yield the main result of this paper:
$`Q_S(\tau )`$ $`=`$ $`{\displaystyle \frac{e\mathrm{}^2}{2\pi m}}{\displaystyle _0^\tau }𝑑t{\displaystyle 𝑑EF(E)\frac{}{E}\left[𝑑\overline{x}^{}\dot{V}(\overline{x}^{},t)𝑑x_2\mathrm{}\{G^t(\overline{x},\overline{x}^{};E)_{x_1}G^t(\overline{x},\overline{x}^{};E)\}\right]}.`$ (7)
*Singlet Pumping in a Turnstile Model*: We illustrate our results with a tight-binding model, (Fig. 1(b)) with two Hubbard impurities located at sites $`m,m`$
$`V(\overline{n},t)=U_{}(t)\delta _{n_1,m}\delta _{n_2,m}+U_+(t)\delta _{n_1,m}\delta _{n_2,m}.`$ (8)
An electron can interact with another only at those two sites therefore, due to the Pauli principle, only singlets are affected. The strength of the interactions, $`U_\pm (t)`$ are the two time-dependent pumping parameters. This concept is similar to a ‘turnstile model’ turnstile but differs significantly in that, instead of time-varying external potentials, the two-body interaction among electrons is varied in time. We separately derived a discrete version of Eq. (7); the end result amounts to replacing the coordinate arguments with site indices $`xn`$, integrals by sums and derivatives with a finite difference form. A lengthy calculation leads to an expression for the singlets pumped in a complete cycle in terms of the *free two-particle lattice* Green’s function, specifically two of its matrix elements $`G_0(0)=G_0(\overline{m},\overline{m};E)`$ and $`G_0(2\overline{m})=G_0(\overline{m},\overline{m};E)`$
$`Q_S(\tau )={\displaystyle \frac{e}{2\pi }}{\displaystyle _0^\tau }𝑑t{\displaystyle 𝑑EF(E)\frac{}{E}\underset{\pm }{}\dot{U}_\pm (t)\frac{|T_\pm (t)|^2\left[\mathrm{}\{G_0(0)\}(1+|T_{}(t)G_0(2\overline{m})|^2)\pm 2\mathrm{}\{T_{}(t)G_0(2\overline{m})G_0^\pm (2\overline{m})\}\right]}{U_\pm (t)^2|1T_{}(t)T_\pm (t)G_0(2\overline{m})G_0(2\overline{m})|^2}}.`$ (9)
Here $`\overline{m}\{m,m\}`$, and $`T_\pm (t)=1/(U_\pm ^1(t)+G_0(0))`$ is the T-matrix for a single Hubbard impurity, and $`G_0^\pm G_0(2\overline{m})`$ for $`\lambda =\pm i|\lambda |`$, when expressed in the form $`G_0(2\overline{m})=_0^\pi \frac{dk}{2\pi }\frac{e^{2m(ik\lambda )}}{\mathrm{sinh}(\lambda )}`$ with $`\mathrm{cosh}(\lambda )=E/2\mathrm{cos}(k)`$.
Exact analytical forms exist for the lattice Green’s functions, $`G_0`$, in terms of elliptic integrals Economou . For the purpose of numerical estimates, we assume a square-profile time dependence nano-transport in the plane of the parameters $`U_\pm (t)`$, shown in Fig. 2(a), where the two parameters change alternately between minimum value $`U_{min}`$ and a maximum value $`U_{max}`$. The pair distribution function is taken to be a Fermi function, $`F(E)1/[e^{\beta (E)}+1]`$. At low temperatures, an integration by parts with respect to energy yields an expression for the pumped singlets in terms of the maximum energy, $``$ available for a pair.
The onsite energy of each tight-binding site is taken to be zero and all energies are expressed in units of the nearest neighbor coupling strength. In Fig. 2(b) we plot the net singlets pumped in a single cycle as a function of the size and location of the square footprint of the time-cycle in the space of the parameters $`U_\pm `$; the flow depends on the enclosed region. Figure 3 shows the dependence of the singlet current on the parameter $``$ that measures the available energy for pairs of electrons determined by the chemical potential in the contacts. The two curves in the figure correspond to different locations of the Hubbard impurities, at lattice sites $`m=\pm 1`$ and $`m=\pm 2`$, illustrating the significant effect the spatial separation of two impurities have on the pumping rate. The direction of flow can also reverse for certain values of the various parameters; reversing the time-cycle is not the only way to reverse the direction of the current Brouwer . Quantum pumping has the intrinsic property that the magnitude of the pumped quantity, in this case singlets, is continuous in nature so that the delivery rate per cycle can be continuously adjusted. Thus there are several ways to precisely control the magnitude and direction of the flow of singlets dynamically.
*Discussion and Outlook*: The turnstile model could be implemented by taking the interaction sites to be quantum dots coupled to leads and varying the electron-electron interaction by decreasing the size of the dot periodically. Concurrent variation in the dot-lead coupling could in principle be compensated by counter-varying gate voltages. Another way to vary the interaction strength is to change the dielectric constant locally. This could be accomplished by a local shift of the electron wavefunction among layers of different materials in a heterostructure; this technique was recently used to control the electron $`g`$ factor Jiang . A shift between silicon and germanium layers or gallium arsenide (GaAs) and aluminium arsenide (AlAs) layers can in principle produce up to $`25\%`$ change in the dielectric constant, sufficient to see the effects described here.
The adiabatic condition requires the period, $`\tau \delta t`$, the dwell time given by $`\delta t=d/v`$, with $`d`$ being the size of the scattering region and $`v`$ the carrier velocity Buttiker-Landauer . Since our singlet pump requires no potential barriers that reduce kinetic energy, $`v`$ may be taken to be the Fermi velocity typically $`>10^5`$ m/s, so for a scattering region $`d10`$ nanometers the adiabatic condition would allow thousands of cycles per nanosecond. Our numerical estimates then give a pumping rate at the order of a thousand singlets per nanosecond. In typical Coulomb-blockade based schemes the most optimistic estimates yield a delivery rate at the order of one nanosecond per singlet Saraga-Loss ; Hu-Sarma . Thus our approach has the potential to be much faster.
Our result, Eq. (7), has the merit that it can also be applied to quantum pumping in systems that allow an independent particle description, as considered in previous studies. All the elements in Eq. (7) then reduce to single particle functions: $`F(E)f(E)`$ is the Fermi distribution function, the Green’s function is a *single* particle one with an asymptotic form $`i[m/(\mathrm{}^2k)]e^{ikx}\psi _k^t(x^{})`$ in terms of the scattering state $`\psi _k^t(x^{})`$ and wavevector $`k=[2mE/\mathrm{}^2]^{1/2}`$. This yields the pumped *charge*
$`Q(\tau )={\displaystyle \frac{em}{2\pi \mathrm{}^2}}{\displaystyle _0^\tau }𝑑t{\displaystyle 𝑑Ef(E)\frac{}{E}\left\{\frac{1}{k}\psi ^t|\dot{V}|\psi ^t\right\}},`$ (10)
which agrees with expressions derived in earlier works nano-transport ; Brouwer . But unlike all previous treatments, we never use 1D scattering matrix elements as they do not have useful generalizations when particles interact with each other in a region rather than scatter off an external potential.
To summarize, we have proposed a method based on two-body adiabatic quantum pumping for generating a dynamically controlled flow of spin-entangled electrons. The process is inherently coherent, potentially much faster than most current proposals, has reduced noise because of the lack of bias and the natural elimination of single electron flow, and allows for continuous adjustment of the flow through numerous physical parameters. All of these features can be developed and incorporated into a comprehensive scheme to generate a controlled flow of entangled electrons. Our goal here has been to present the basic idea, develop a theoretical framework for its description and discuss a possible physical model for implementation. We have in the process generalized the treatment of quantum pumping to incorporate interactions that cannot be treated in a scattering matrix approach. We anticipate that future studies can build upon the considerations in this work.
This work was supported by the Packard Foundation and NSF NIRT program grant DMR-0103068.
*Note added*: When we were finishing this paper we become aware of an excellent recent proposal by Beenakker *et al.* Beenakker to excite entangled electron-hole pairs using one-body potential in a scattering matrix approach. Our proposal is more difficult to implement, but it could be naturally incorporated into an electron-only spintronic device, it produces a current of singlets unadulterated by single electrons, and it will enjoy a longer decoherence time since hole decoherence times tend to be short.
|
warning/0506/math0506052.html
|
ar5iv
|
text
|
# Family of intersecting totally real manifolds of (𝐶^𝑛,0) and CR-singularities
## 1 Introduction
The aim of the article is to study the geometry of some germs of real analytic submanifolds of $`(C^n,0)`$. On the one hand, we shall study family of totally real submanifolds of $`(C^n,0)`$ intersecting at the origin. On the other hand, we shall study submanifolds having a CR-singularity at the origin. In both cases, we are primary interested in the holomorphic classification of such objects, that is the orbit of the action of the group of germs of holomorphic diffeomorphisms fixing the origin.
In this article, we shall mainly focus on the existence of complex analytic subsets intersecting such germs of real analytic manifolds. In the first problem, we shall also be interested in the problem of straightening holomorphically the family. We mean that we shall give sufficient condition which will ensure that, in a good holomorphic coordinates system, each submanifold of the family is an $`n`$-plane. In the case there are formal obstructions to straighten the family, we show the existence of a germ complex analytic variety which interesects the family along a set that can be straightened. The first part of this work takes its roots in and generalizes a recent work of Sidney Webster \[Web03\] from which it is very inspired. This part of the work start after having listen to Sidney Webster at the Partial Differential Equations and Several Complex Variables conference held in Wuhan University in June 2004.
The starting point of the first problem appeared already in the work of E. Kasner \[Kas13\] and was studied, from the formal view point, by G.A. Pfeiffer \[Pfe15\]. They were interested in pairs of real analytic curves in $`(C,0)`$ passing throught the origin. We shall not consider the case were some of the submanifolds are tangent to some others. We refer the reader to the works of I. Nakai \[Nak98\], J.-M. Trépreau \[Tré03\] and X. Gong \[AG05\] in this direction.
In the second part, we shall study $`(n+p1)`$-real analytic submanifolds of $`C^n`$ of the form
$$\{\begin{array}{cc}y_{p+1}=F_{p+1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ \mathrm{}\hfill & \\ y_{p+q}=F_{p+q}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ z_n=G(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \end{array}$$
where $`z^{}=(z_1,\mathrm{},z_p)C^p`$, $`z^{\prime \prime }=(z_{p+1},\mathrm{},z_{n1})C^{n1p}`$, $`pN^{}`$. The $`F`$’s and $`G`$ vanish at the second order at the origin. The origin is a singularity for the Cauchy-Riemann structures. The most studied case, up to now, is the case where $`p=1`$ ($`n`$-submanifold). Nevertheless, some work has been done for smaller dimensional submanifold by Adam Coffman (see for instance \[Cof04a, Cof04b\]). A nondegenerate real analytic surface in $`(C^2,0)`$ which is totally real except at the origin where it has a complex tangent can be regarded as a third order analytic deformation of the quadric
$$z_2=z_1\overline{z}_1+\gamma (z_1^2+\overline{z_1}^2).$$
This is due to Bishop \[Bis65\] and the nonnegative number $`\gamma `$ is called the Bishop invariant. When $`0<\gamma <1/2`$ (we say elliptic), J. Moser and S. Webster showed, in their pioneering work \[MW83\] that an analytic deformation of such a quadric could be transformed into a normal form (in fact a real algebraic variety) by the mean of a germ of holomorphic diffeomorphism preserving the origin. All the geometry a such a deformation is understood by the study of the normal form. When $`1/2<\gamma `$ (we say hyperbolic), it is known that such a statement doesn’t hold. So, what about the geometry ? Wilhelm Klingenberg Jr. showed \[Kli85\] that there exists a germ of complex curve passing throught the origin cutting the submanifold along two transversal real analytic curves which are tangent to $`\{\overline{z}_1=\lambda z_1\}`$ and $`\{\overline{z}_1=\lambda ^1z_1\}`$ respectively. Here, $`\lambda `$ is a solution of $`\gamma \lambda ^2+\lambda +\gamma =0`$ which is assumed to satify a diophantine condition : there exists $`M,\delta >0`$ such that, for all positive integer $`k`$, $`|\lambda ^k1|>Mk^\delta `$. We refer to \[BEPR00\] for a summary in this framework, to \[Hua04\] for a nice introduction and to \[Har84\] for another point of view.
We shall work with nondegenerate submanifolds (in some sense) and then define generalized Bishop invariants $`\{\gamma _i\}_{i=1,\mathrm{},p}`$ : There is a good holomorphic coordinates system at the origin in which the submanifold is defined by
$$(M_Q)\{\begin{array}{cc}y_\alpha =f_\alpha (z^{},\overline{z}^{},x^{\prime \prime })\alpha =p+1,\mathrm{},n1\hfill & \\ z_n=Q(z^{},\overline{z}^{})+g(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \end{array}$$
where the $`f_i`$’s and $`g`$ are germs of real analytic functions at the origin and of order greater than or equal to $`3`$ there. The quadratic polynomial $`Q`$ is of the form
$$Q(z^{},\overline{z}^{})=d_{i,l}z_i^{}\overline{z}_l^{}+\underset{i=1}{\overset{p}{}}\gamma _i((z_i^{})^2+(\overline{z}_i^{})^2),$$
the norm of the sesquilinear part of $`Q`$ being $`1`$. It is regarded as a perturbation of the quadric
$$(Q)\{\begin{array}{cc}y_{p+1}=\mathrm{}=y_{n1}=0\hfill & \\ z_n=Q(z^{},\overline{z}^{})\hfill & \end{array}$$
We shall consider the case where none of these invariants vanishes (see \[Mos85, HK95\] for results in this situation for $`p=1`$). Then, we begin a study à la Moser-Webster of such an object althought we shall not study here the normal form problem. We associate a pair $`(\tau _1,\tau _2)`$ of germs of holomorphic involutions of $`(C^{n+p1},0)`$ and also a germ of biholomorphism $`\mathrm{\Phi }=\tau _1\tau _2`$.
In the one hand, we shall show under some asumptions (in particular the diophantiness property of $`D\mathrm{\Phi }(0)`$), that if $`(M_Q)`$ is formally equivalent to its associated quadric $`(Q)`$ then it is holomorphically equivalent to it. This generalizes a result by X. Gong \[Gon94\](p=1).
On the other hand, the complex linear space $`\{z_{p+1}=\mathrm{}=z_n=0\}`$ intersects the quadric along the real linear set $`:=\{Q(z^{},\overline{z}^{})=0\}`$. We shall show that this situation survive under a small analytic perturbation. Namely, we shall show under some assuptions, that there exists a good holomorphic coordinates system at the origin in which the submanifold $`(M_Q)`$ intersects the complex linear space $`\{z_{p+1}=\mathrm{}=z_n=0\}`$ along a real analytic subset $`V`$ passing throught the origin. The latter is completely determined by the eigenvalues of $`D\mathrm{\Phi }(0)`$ and more precisely by the centralizer of the map $`vD\mathrm{\Phi }(0)v`$ in the space of non-linear formal maps. This generalizes the result of Wilhelm Klingenberg Jr.(see above) ($`p=1`$) in the sense that, there are holomorphic coordinates such that $`(M_Q)`$ intersects $`\{z_2=0\}`$ along $`\{(\zeta _1,\eta _1)R^2|\zeta _1\eta _1=0\}`$.
The core of these problems rests on geometric properties of dynamical systems associated to each situation. To be more specific, we shall deal, in the first part of this article, with germs of holomorphic diffeomorphisms of $`(C^n,0)`$ in a neighbourhood of the origin (a common fixed point). We shall consider those whose linear part at the origin is different from the identity. The main point is a result which gives the existence of germ of analytic subset of $`(C^n,0)`$ invariant by a abelian group of such diffeomorphisms under some diophantine condition. This kind of result was obtained by the author for a germ of holomorphic vector field at singular point \[Sto94\].
## 2 Abelian group of diffeomorphisms of $`(C^n,0)`$ and their invariant sets
The aim of this section is to prove the existence of complex analytic invariant subset for a commuting family of germs of holomorphic diffeomorphisms in a neighbourhood of a common fixed point. This is very inspired by a previous article of the author concerning holomorphic vector fields. Althought the objects are not the same, some of the computations are identical and we shall refer to them when possible.
Let $`D_1:=\text{diag}(\mu _{1,1},\mathrm{},\mu _{1,n}),\mathrm{},D_l:=\text{diag}(\mu _{l,1},\mathrm{},\mu _{l,n})`$ be diagonal invertible matrices. Let us consider a familly $`F:=\{F_i\}_{i=1,\mathrm{}l}`$ of commuting germs of holomorphic diffeomorphisms of $`(C^n,0)`$ which linear part, at the origin, is $`D:=\{D_ix\}_{i=1,\mathrm{}l}`$ :
$$F_i(x)=D_ix+f_i(x),\text{with}f_i(0)=0,Df_i(0)=0,f_i𝒪_n.$$
Let $``$ be an ideal of $`𝒪_n`$ generated by monomials of $`C^n`$. Let $`V()`$ be the germ at the origin, of the analytic subset of $`(C^n,0)`$ defined by $``$. It is left invariant by the familly $`D`$. Let us set $`\widehat{}:=\widehat{𝒪}_n`$. Here we denote $`𝒪_n`$ (resp. $`\widehat{𝒪}_n`$) the ring of germ of holomorphic function at the origin (resp. ring of formal power series) of $`C^n`$. Let $`Q=(q_1,\mathrm{},q_n)N^n`$ and $`x=(x_1,\mathrm{},x_n)C^n`$, we shall write
$$|Q|:=q_1+\mathrm{}+q_n,x^Q:=x_1^{q_1}\mathrm{}x_n^{q_n}.$$
Let $`\{\omega _k(D,)\}_{k1}`$ be the sequence of positive numbers defined by
$$\omega _k(D,)=inf\left\{\underset{1il}{\mathrm{max}}|\mu _i^Q\mu _{i,j}|0\right|\mathrm{\hspace{0.33em}2}|Q|2^k,1jn,QN^n,x^Q\}.$$
Let $`\{\omega _k(D)\}_{k1}`$ be the sequence of positive numbers defined by
$$\omega _k(D,)=inf\left\{\underset{1il}{\mathrm{max}}|\mu _i^Q\mu _{i,j}|0\right|\mathrm{\hspace{0.33em}2}|Q|2^k,1jn,QN^n\}.$$
###### Definition 2.1.
1. We shall say that the ideal $``$ is properly embedded if it has a set of monomial generators not involving a nonempty set $`𝒮`$ of variables.
2. We shall say that the familly $`D`$ is diophantine (resp. on $``$) if
$$\underset{k1}{}\frac{\mathrm{ln}\omega _k(D)}{2^k}<+\mathrm{}(resp.\underset{k1}{}\frac{\mathrm{ln}\omega _k(D,)}{2^k}<+\mathrm{}).$$
3. We shall say that the familly $`F`$ is formally linearizable on $`\widehat{}`$ if there exists a formal diffemorphism $`\widehat{\mathrm{\Phi }}`$ of $`(C^n,0)`$, tangent to the identity at the origin such that $`\widehat{\mathrm{\Phi }}_{}F_iD_ix(\widehat{})^n`$ for all $`1il`$.
4. A linear anti-holomorphic involution of $`C^n`$ is a map $`\rho (z)=P\overline{z}`$ where the matrix $`P`$ satisfies $`P\overline{P}=Id`$; $`\overline{z}`$ denotes the complex conjugate of $`z`$.
5. We shall say that $``$ is compatible with a anti-linear involution $`\rho `$ if the map $`\rho ^{}:\widehat{𝒪}_{n+p1}\overline{\widehat{𝒪}_{n+p1}}`$) defined $`\rho ^{}(f)=f\rho `$ maps $`\widehat{}`$ to $`\overline{\widehat{}}`$ and $`\widehat{CI}`$ to $`\overline{\widehat{CI}}`$.
Let $`\widehat{𝒪}_n^D`$ be the ring of formal invariant of the familly $`D`$, that is
$$\widehat{𝒪}_n^D:=\{f\widehat{𝒪}_n|f(D_ix)=f(x)i=1,\mathrm{},l\}.$$
It can be shown (as in proposition 5.3.2 of \[Sto00\]) that this ring is generated by a finite number of monomials $`x^{R_1},\mathrm{},x^{R_p}`$ and that the non-linear centralizer $`𝒞_D`$ of $`D`$ is a module over $`\widehat{𝒪}_n^D`$ of finite type. Let ResIdeal be the ideal generated by the monomials $`x^{R_1},\mathrm{},x^{R_p}`$ in $`𝒪_n`$.
###### Theorem 2.1.
Let $``$ be a monomial ideal (resp. properly embedded). Assume that the familly $`D`$ is diophantine (resp. on $``$). If the familly $`F`$ is formally linearizable on $`\widehat{}`$, then it is holomorphically linearizable on $``$. Moreover, there exists a unique such a diffeomorphism $`\mathrm{\Phi }`$ such that the projection of the Taylor expansion of $`\mathrm{\Phi }Id`$ onto $`𝒞_D`$ vanishes.
Moreover, let $`\rho `$ be a linear anti-holomorphic involution such that $`\rho 𝒞_D\rho =𝒞_D`$. We assume that $``$ is compatible with $`\rho `$. Assume that, for all $`1il`$, $`\rho F_i\rho `$ belongs to the group generated by the $`F_i`$’s. Then $`\mathrm{\Phi }`$ and $`\rho `$ commute with each other.
This theorem can be rephrased as follow : Under the afore-mentioned diophantine condition, then there exists a germ of holomorphic diffeomorphism $`\mathrm{\Phi }`$ such that $`\mathrm{\Phi }_{}F_iD_ix()^n`$ for all $`1il`$. As a consequence, in a good holomorphic coordinates system, the analytic subset $`V()`$ is left invariant by each $`F_i`$ and its restriction to it is the linear mapping $`xD_{i|V()}x`$.
###### Remark 2.1.
The familly $`D`$ can be diophantine while none of the $`D_i`$’s is.
The second part of the theorem will be used for applications in the third part of the article.
###### Corollary 2.1.
1. If the ring of invariant of $`D`$ reduces to the constants and if $`D`$ is diophantine, then $`F`$ is holomorphically linearizable in a neighbourhood of the origin. For one diffeomorphism, this was obtain by H. Rüssmann \[Rüs77, Rüs02\] and by T. Gramtchev and M.Yoshino \[GY99\] for an abelian group under a slightly coarser diophantine condition.
2. The existence of an invariant manifold for a germ of diffeomorphism was obtain by J. Pöschel \[Pös86\]. Despite the fact that we are dealing with a family of diffeomorphisms, the main difference is that we are able to linearize simultaneously on each irreducible component of analytic set.
According to M. Chaperon \[Cha86\]\[theorem 4, p.132\], if the family of diffeomorphisms is abelian then there exists a formal diffeomorphism $`\widehat{\mathrm{\Phi }}`$ such that
$$\widehat{\mathrm{\Phi }}_{}F_i(D_jz)=D_j\widehat{\mathrm{\Phi }}_{}F_i(z),1i,jl.$$
We call the family of $`\widehat{\mathrm{\Phi }}_{}F_i`$’s a formal normal form of the family $`F`$. Then we have the following corollary :
###### Corollary 2.2.
Let $`F`$ be an abelian family of germs of holomorphic diffeomorphisms of $`(C^n,0)`$. Let us assume that $`D`$ is diophantine on $`ResIdeal`$. If the non-linear centralizer of $`D`$ is generated by the $`x^{R_i}`$’s then $`F`$ is holomorphically linearizable on $``$.
###### Remark 2.2.
The condition that the non-linear centralizer of $`D`$ is generated by the $`x^{R_i}`$’s means: if $`\mu _i^Q=\mu _{i,j}`$ for some $`QN_2^n`$, $`1jn`$ and for all $`1il`$, then $`x^Q`$ belongs to the ideal generated by $`x^{R_1},\mathrm{},x^{R_p}`$. This a very weak condition since only all but a finite number of resonances satisfy this condition.
We shall prove that there exists a holomorphic map $`\varphi :(C^n,0)(C^n,0)`$, tangent to the identity at the origin, such that
$$\mathrm{\Phi }^1F_i\mathrm{\Phi }(y)=G_i(y):=D_iy+g_i(y)i=1,\mathrm{},l$$
where the components of $`g_i`$ are non-linear holomorphic functions and belong to the ideal $``$. It is unique if we require that its projection on $`\text{ResIdeal}`$ is zero.
Let us set $`x_j=\mathrm{\Phi }_j(y):=y_j+\varphi _j(y)`$, $`j=1,\mathrm{},n`$. Let us expand the equations $`F_i\mathrm{\Phi }(y)=\mathrm{\Phi }G_i`$, $`i=1,\mathrm{},l`$. For all $`1jn`$ and all $`i=1,\mathrm{},l`$, we have
$`\mu _{i,j}y_j+g_{i,j}(y)+\varphi _j(G_i(y))`$ $`=`$ $`\mu _{i,j}(y_j+\varphi _j(y))+f_{i,j}(\mathrm{\Phi }(y))`$
$`g_{i,j}(y)+\varphi _j(D_iy)`$ $`=`$ $`\mu _{i,j}\varphi _j(y)+f_{i,j}(\mathrm{\Phi }(y))`$
$`+(\varphi _j(G_i(y))\varphi _j(D_iy))`$
Let us expand the functions at the origin :
$$f_{i,j}(y)=\underset{QN_2^n}{}f_{i,j,Q}y^Q,g_{i,j}(y)=\underset{QN_2^n}{}g_{i,j,Q}y^Q\text{ and }\varphi _j(y)=\underset{QN_2^n}{}\varphi _{j,Q}y^Q.$$
Then we have
$$\underset{QN_2^n}{}\delta _{Q,j}^i\varphi _{j,Q}y^Q+g_{i,j}(y)=f_{i,j}(\mathrm{\Phi }(y))(\varphi _j(G_i(y))\varphi _j(D_iy))$$
(1)
where
$$\delta _{Q,j}^i:=\mu _i^Q\mu _{i,j},\mu _i:=(\mu _{i,1},\mathrm{},\mu _{i,n}).$$
Let $`\{f\}_Q`$ denotes the coefficient of $`x^Q`$ in the Taylor expansion at the origin of $`f`$. We compute $`\varphi _{j,Q}`$ and $`g_{i,j,Q}`$ by induction on $`|Q|2`$ in the following way :
* if $`y^Q`$ does not belongs to $``$ and $`\mathrm{max}_i|\delta _{Q,j}^i|0`$, then there exists $`1i_0l`$ such that $`|\delta _{Q,j}^{i_0}|=\mathrm{max}_i|\delta _{Q,j}^i|`$. We set
$`\varphi _{j,Q}`$ $`=`$ $`{\displaystyle \frac{1}{\delta _{Q,j}^{i_0}}}\left\{f_{i_0,j}(\mathrm{\Phi }(y))(\varphi _j(G_{i_0}(y))\varphi _j(D_{i_0}y))\right\}_Q`$
$`g_{i,j,Q}`$ $`=`$ $`0.`$
* If $`y^Q`$ does not belongs to $``$ and $`\mathrm{max}_i|\delta _{Q,j}^i|=0`$, then we have
$$\left\{f_{i_0,j}(\mathrm{\Phi }(y))(\varphi _j(G_{i_0}(y))\varphi _j(D_{i_0}y))\right\}_Q=0$$
and we set $`\varphi _{j,Q}=0=g_{i,j,Q}`$.
* If $`y^Q`$ belongs to $``$, we set
$`\varphi _{j,Q}`$ $`=`$ $`0`$
$`g_{i,j,Q}`$ $`=`$ $`\left\{f_{i,j}(\mathrm{\Phi }(y))(\varphi _j(G_i(y))\varphi _j(D_iy))\right\}_Q.`$
###### Lemma 2.1.
The formal diffeomorphism $`\mathrm{\Phi }`$ defined above linearizes simultanueously the family $`F`$ on $`\widehat{}`$ where $`\widehat{}:=\widehat{𝒪}_n`$.
###### Proof.
For all $`1i,jjl`$, we have
$$F_iF_j=F_jF_i\text{ thus }F_iF_j\mathrm{\Phi }=F_jF_i\mathrm{\Phi }.$$
Therefore, we have
$$D_iD_j\mathrm{\Phi }+D_i(f_j\mathrm{\Phi })+f_i(D_j\mathrm{\Phi }+f_j\mathrm{\Phi })=D_jD_i\mathrm{\Phi }+D_i(f_i\mathrm{\Phi })+f_j(D_i\mathrm{\Phi }+f_i\mathrm{\Phi }),$$
so that
$`f_j(\mathrm{\Phi }D_i)D_i(f_j\mathrm{\Phi })`$ $`=`$ $`f_i(\mathrm{\Phi }D_j)D_j(f_i\mathrm{\Phi })`$
$`+f_j(D_i\mathrm{\Phi }+f_i\mathrm{\Phi })f_j(\mathrm{\Phi }D_i)`$ $`+f_i(D_j\mathrm{\Phi }+f_j\mathrm{\Phi })f_i(\mathrm{\Phi }D_j).`$
Moreover, we have $`F_i\mathrm{\Phi }=\mathrm{\Phi }G_i`$. Hence, we have
$`f_i(D_j\mathrm{\Phi }+f_j\mathrm{\Phi })f_i(\mathrm{\Phi }D_j)`$ $`=`$ $`f_iF_j\mathrm{\Phi }f_i(\mathrm{\Phi }D_j)`$
$`=`$ $`f_i\mathrm{\Phi }G_jf_i\mathrm{\Phi }D_j`$
$`=`$ $`D(f_i\mathrm{\Phi })(D_jy)g_j+\mathrm{}`$
Assume the $`F_i`$’s are linearized on $`V()`$ up to order $`k2`$. This means that, for any $`1mn`$ and any $`1il`$, the $`k`$-jet $`J^k(g_{i,m})`$ belongs to $``$. The previous computation shows that the $`k+1`$-jet of $`f_i(D_j\mathrm{\Phi }+f_j\mathrm{\Phi })f_i(\mathrm{\Phi }D_j)`$ depends only on the $`k`$-jet of $`g_j`$ and belongs to $``$. The same is true for $`\varphi _j(G_i(y))\varphi _j(D_iy)`$. Therefore, if $`QN_2^n`$ with $`|Q|=k+1`$ is such that $`x^Q`$ does not belong to $``$, then we have
$$\{f_j(\mathrm{\Phi }D_i)D_i(f_j\mathrm{\Phi })\}_Q=\{f_i(\mathrm{\Phi }D_j)D_j(f_i\mathrm{\Phi })\}_Q;$$
that is, for all $`1mn`$,
$$(\mu _i^Q\mu _{i,m})\{f_{j,m}\mathrm{\Phi }\}_Q=(\mu _j^Q\mu _{j,m})\{f_{i,m}\mathrm{\Phi }\}_Q.$$
This means that equation $`(\text{1})`$ is solved by induction and that $`\mathrm{\Phi }`$ linearizes formally the $`F_i`$’s on $`V(\widehat{})`$. ∎
Let $`\rho `$ be a linear anti-holomorphic involution satisfying the assumptions of the theorem. We have $`F_i\mathrm{\Phi }=\mathrm{\Phi }G_i`$ where $`G_i`$ is linearized along $`\widehat{}`$. Hence, we have
$$(\rho F_i\rho )(\rho \mathrm{\Phi }\rho )=(\rho \mathrm{\Phi }\rho )(\rho G_i\rho ).$$
Let us set $`\stackrel{~}{F}_i:=\rho F_i\rho `$. By assusmptions, $`\stackrel{~}{F}_i`$ belongs to the group generated by the $`F_i`$’s. Since $`\rho ^{}\widehat{}\overline{\widehat{}}`$, then $`\rho G_i\rho `$ is a formal diffeomorphism which is linearized on $`\widehat{}`$. By assumptions, the projection of $`\rho \mathrm{\Phi }\rho Id`$ onto $`𝒞_D`$ vanishes identically. By uniqueness, we have $`\rho \mathrm{\Phi }\rho =\mathrm{\Phi }`$ since $`\mathrm{\Phi }`$ linearizes $`\stackrel{~}{F}_i`$ on $``$.
We shall prove, by using the majorant method, that $`\mathrm{\Phi }`$ actually converges on a polydisc of positive radius centered at the origin. Let us define $`N_2^n\widehat{}`$ to be the set of multiindices $`QN^n`$ such that $`|Q|2`$ and $`x^Q\widehat{}`$. Let $`f=_Qf_Qx^Q`$ and $`g=_Qg_Qx^Q`$ be formal power series. We shall say that $`g`$ dominates if $`|f_Q||g_Q|`$ for all multiindices $`Q`$.
First of all, for all $`1jn`$ and all $`QN_2^n\widehat{}`$ such that $`\mathrm{max}_{1il}|\delta _{j,Q}^i|0`$, we have
$$|\varphi _{j,Q}||\delta _{j,Q}|=|\{f_{i_0(Q),j}(\mathrm{\Phi })\}_Q|\{\overline{f}_{i_0(Q),j}(y+\overline{\varphi })\}_Q$$
where $`|\delta _{j,Q}|=\mathrm{max}_{1il}|\delta _{j,Q}^i|=|\delta _{i_0(Q,j),j,Q}|`$. In fact, $`\{f_i\mathrm{\Phi }G_jf_i\mathrm{\Phi }D_j\}_Q=0`$ whenever $`QN_2^n\widehat{}`$. This inequality still holds if $`\mathrm{max}_{1il}|\delta _{j,Q}^i|=0`$. Let us set
* $`\delta _Q:=\mathrm{min}\{|\delta _{j,Q}|,1jn`$ such that $`\delta _{j,Q}0\}`$,
* $`\delta _Q:=0`$ if $`\mathrm{max}_{1il}|\delta _{j,Q}^i|=0`$.
Let us sum over $`1jn`$ the previous inequalities. We obtain for all $`QN_2^n\widehat{}`$,
$$\delta _Q\underset{j=1}{\overset{n}{}}|\varphi _{j,Q}|\underset{j=1}{\overset{n}{}}|\varphi _{j,Q}||\delta _{j,Q}|\left\{\underset{j=1}{\overset{n}{}}\overline{f}_{i_0(Q,j),j}(y+\overline{\varphi })\right\}_Q\left\{\underset{i=1}{\overset{l}{}}\left(\underset{j=1}{\overset{n}{}}\overline{f}_{i,j}\right)(y+\overline{\varphi })\right\}_Q.$$
Since $`_{i=1}^l_{j=1}^nf_{i,j}`$ vanishes at the origin with its derivative as well, there exists positives constants $`a,b`$ such that
$$\underset{i=1}{\overset{l}{}}\underset{j=1}{\overset{n}{}}f_{i,j}\frac{a\left(_{j=1}^nx_j\right)^2}{1b\left(_{j=1}^nx_j\right)}.$$
Since the Taylor expansion of the right hand side has non-negative coefficients, we obtain
$$\delta _Q\stackrel{~}{\varphi }_Q\left\{\frac{a\left(_{j=1}^ny_j+\stackrel{~}{\varphi }\right)^2}{1b\left(_{j=1}^ny_j+\stackrel{~}{\varphi }\right)}\right\}_Q$$
where we have set $`\stackrel{~}{\varphi }_Q:=_{j=1}^n|\varphi _{j,Q}|`$ and $`\stackrel{~}{\varphi }=_{QN_2^n}\stackrel{~}{\varphi }_Qx^Q`$. Here, we have set $`\stackrel{~}{\varphi }_Q=0`$ whenever $`\delta _Q=0`$.
Let us define the formal power series $`\sigma (y)=_{QN_2^n}\sigma _Qy^Q`$ as follow :
$`QN_2^n(N_2^n\widehat{})\sigma _Q`$ $`=`$ $`0`$
$`QN_2^n\widehat{}\sigma _Q`$ $`=`$ $`\left\{{\displaystyle \frac{a\left(_{j=1}^ny_j+\sigma \right)^2}{1b\left(_{j=1}^ny_j+\sigma \right)}}\right\}_Q`$
###### Lemma 2.2.
\[Sto94\]\[Lemme 2.1\] The series $`\sigma `$ is convergent in a neighbourhood of the origin $`0C^n`$.
Let us define the sequence $`\{\eta _Q\}_{QN_1^n\widehat{}}`$ of positive number as follow :
1. $`PN_1^n\widehat{}`$ tel que $`|P|=1`$, $`\eta _P=1`$ ( such multiindice exists. ),
2. $`QN_2^n\widehat{}`$ with $`\delta _Q0`$
$$\delta _Q\eta _Q=\underset{\begin{array}{c}Q_jN_1^n,SN^n\\ Q_1+\mathrm{}+Q_p+S=Q\end{array}}{\mathrm{max}}\eta _{Q_1}\mathrm{}\eta _{Q_p},$$
the maximum been taken over the sets of $`p+1`$, $`1p|Q|`$, multiindices $`Q_1,\mathrm{},Q_p,S`$ such that $`1jp,Q_jN_1^n,|Q_j|<|Q|`$, $`SN^n`$. These sets are not empty.
3. $`QN_2^n\widehat{}`$ with $`\delta _Q=0`$, $`\eta _Q=0`$.
This sequence is well defined. In fact, if $`QN_2^n\widehat{}`$, then there exists multiindices $`Q_1,\mathrm{},Q_p,S`$ such that $`Q=Q_1+\mathrm{}+Q_p+S`$, $`1jp,Q_jN_1^n,|Q_j|<|Q|,SN^n`$. In this case, $`1jp,Q_jN_1^n\widehat{}`$.
The following lemmas are the key points.
###### Lemma 2.3.
\[Sto94\]\[Lemme 2.2\] For all $`QN_2^n\widehat{}`$, we have $`\stackrel{~}{\varphi }_Q\sigma _Q\eta _Q`$.
###### Lemma 2.4.
\[Sto94\]\[Lemme 2.3\] There exists a constant $`c>0`$ such that $`QN_2^n\widehat{},\eta _Qc^{|Q|}`$.
Let $`\theta >0`$ be such that $`4\theta :=\mathrm{min}_{i,j}|\lambda _{i,j}|1`$ (we can always assume this, even if this means using the inverse of one of the diffeomorphisms). If the ideal $``$ is properly embedded, then we shall set
$$4\theta :=\underset{1il,j𝒮}{\mathrm{min}}|\lambda _{i,j}|1$$
where $`𝒮`$ denotes the set of variables not involves in any generator. In particular, we have the property that if $`x^Q`$ then $`x_sx^Q`$ for all $`s𝒮`$.
By definition, $`\eta _Q`$ is a product of $`1/\delta _Q^{}`$ with $`|Q^{}||Q|`$. Let $`k`$ be a non-negative integer. Let us define $`\varphi ^{(k)}(Q)`$ (resp. $`\varphi _j^{(k)}(Q)`$) to be the number of $`1/\delta _Q^{}`$’s present in this product and such that $`0\delta _Q^{}<\theta \omega _k(D,)`$ (resp. and $`\delta _Q=\delta _{j,Q}`$). The lemma is a consequence of the following proposition
###### Proposition 2.1.
\[Sto94\]\[lemme 2.8\] For all $`QN_2^n\widehat{}`$, we have $`\varphi ^{(k)}(Q)2n\frac{|Q|}{2^k}`$ if $`|Q|2^k+1`$; and $`\varphi ^{(k)}(Q)=0`$ if $`|Q|2^k`$.
In fact, $`\varphi ^{(k)}(Q)`$ bounds the number of $`1/\delta _Q^{}`$’s appearing in the product defining $`\eta _Q`$ and such that $`\theta \omega _{k+1}(D,)\delta _Q^{}<\theta \omega _k(D,)`$.
###### Proof of lemma 2.4.
Let $`r`$ be the integer such that $`2^r+1|Q|<2^{r+1}+1`$. Then we have
$$\eta _Q\underset{k=0}{\overset{r}{}}\left(\frac{1}{\theta \omega _{k+1}(D,)}\right)^{\varphi ^{(k)}(Q)}.$$
By applying the Logarithm and proposition 2.1, we obtain
$`\mathrm{ln}\eta _Q`$ $``$ $`{\displaystyle \underset{k=0}{\overset{l}{}}}2n{\displaystyle \frac{|Q|}{2^k}}\left(\mathrm{ln}{\displaystyle \frac{1}{\theta \omega _{k+1}(D)}}\right)`$
$``$ $`|Q|\left(2n{\displaystyle \underset{k0}{}}{\displaystyle \frac{\mathrm{ln}\omega _{k+1}(D)}{2^k}}+2n\mathrm{ln}\theta ^1{\displaystyle \underset{k0}{}}{\displaystyle \frac{1}{2^k}}\right).`$
Since the familly $`D`$ is diophantine, we obtain $`\eta _Qc^{|Q|}`$ for some positive constant $`c`$. ∎
For any positive integer $`k`$, for any $`1jn`$, let us consider the function defined on $`N_2^n\widehat{}`$ to be
$$QN_2^n\widehat{},\psi _j^{(k)}(Q)=\{\begin{array}{c}1\text{ if }\delta _Q=|\delta _{j,Q}|0\text{ and }|\delta _{j,Q}|<\theta \omega _k(D,)\hfill \\ 0\text{ if }\delta _Q=0\text{ or }\delta _Q|\delta _{j,Q}|\text{ or }|\delta _{j,Q}|\theta \omega _k(D,)\hfill \end{array}$$
Then we have,
$$0\varphi _j^{(k)}(Q)=\psi _j^{(k)}(Q)+\underset{\begin{array}{c}Q_jN_1^n,SN^n\\ Q_1+\mathrm{}+Q_p+S=Q\end{array}}{\mathrm{max}}\left(\varphi _j^{(k)}(Q_1)+\mathrm{}+\varphi _j^{(k)}(Q_p)\right).$$
The proof of propositon 2.1 identitical to the proof of \[Sto94\]\[lemme 2.8\] except that we have to use the following version of \[Sto94\]\[lemme 2.7\].
###### Lemma 2.5.
Let $`QN_2^n\widehat{}`$ be such that $`\psi _j^{(k)}(Q)=1`$. If $`Q=P+P^{}`$ with $`(P,P^{})N_1^n\times N_2^n`$ and $`|P|2^k1`$, then $`(P,P^{})N_1^n\widehat{}\times N_2^n\widehat{}`$ and $`\psi _j^{(k)}(P^{})=0`$.
###### Proof.
Clearly, if $`Q=P+P^{}N_2^n\widehat{}`$ then $`(P,P^{})N_1^n\widehat{}\times N_2^n\widehat{}`$. There are two cases to consider :
1. if $`\delta _P^{}|\delta _{j,P^{}}|`$ or $`\delta _P^{}=0`$ then $`\psi _j^{(k)}(P^{})=0`$, by definition.
2. if $`\delta _P^{}=|\delta _{j,P^{}}|0`$, assume that $`\delta _P^{}<\theta \omega _k(D,)`$. Then, for all $`1il`$, we have
$$|\lambda _i^P^{}|>|\lambda _{i,j}|\theta \omega _k(D,)4\theta 2\theta =2\theta .$$
It follows that, for all $`1il`$,
$`2\theta \omega _k(D,)`$ $`>`$ $`|\lambda _i^Q\lambda _{i,j}|+|\lambda _i^P^{}\lambda _{i,j}|`$
$`>`$ $`|\lambda _i^Q\lambda _i^P^{}|=|\lambda _i^P^{}||\lambda _i^P1|.`$
If $``$ is properly embedded, for all $`a𝒮`$, we have $`x_ax^P`$. Therefore, for all $`1il`$, we have
$`2\theta \omega _k(D,)`$ $`>`$ $`2\theta |\lambda _{i,a}|^1|\lambda _i^{P+E_a}\lambda _{i,a}|`$
$`>`$ $`2\theta |\lambda _{i,a}|^1\omega _k(D,).`$
This contradicts the facts that $`\mathrm{min}_{1il,a𝒮}|\lambda _{i,a}|1`$. If $``$ is not properly embedded, then we obtain
$$2\theta \omega _k(D)>2\theta |\lambda _{i,a}|^1\omega _k(D)$$
for all $`1an`$. It is still a contradiction.
Hence, we have shown that $`\psi _j^{(k)}(P^{})=0`$. ∎
## 3 Family of totally real $`n`$-manifolds in $`(C^n,0)`$
Let us consider a family $`M:=\{M_i\}_{i=1,\mathrm{},m}`$ of real analytic totally real $`n`$-submanifold of $`C^n`$ passing throught the origin. Locally, each $`M_i`$ is the fixed point set of an anti-holomorphic involution $`\rho _i`$ : $`M_i=FP(\rho _i)`$ and $`\rho _i\rho _i=Id`$. This means that
$$\rho _i(z):=B_i\overline{z}+R_i(\overline{z})$$
where $`R_i`$ is a germ of holomorphic function at the origin with $`R_i(0)=0`$ and $`DR_i(0)=0`$. Each matrix $`B_i`$ is invertible and satisfies $`B_i\overline{B}_i=Id`$. The tangent space, at the origin, of $`M_i`$ is the totally real $`n`$-plane
$$\{z=B_i\overline{z}\}$$
We assume that there are all distinct one from another. Their intersection at the origin is the set
$$\left\{zC^n\right|B_i\overline{z}=z,i=1,\mathrm{},m\}\left\{zC^n\right|B_i\overline{B}_jz=z,i,j=1,\mathrm{},m\}.$$
It is contained in the common eigenspace of the $`B_i\overline{B}_j`$’s associated to the eigenvalue $`1`$. We shall not assume that this space is reduced to $`0`$.
Let us consider the group $`G`$ generated by the germs of holomorphic diffeomorphisms of $`(C^n,0)`$ $`F_{i,j}:=\rho _i\rho _j`$, $`1i,jm`$. Let $`D_{i,j}:=B_i\overline{B}_j`$ be the linear part at the origin of $`F_{i,j}`$. Let us set
$$F_{i,j}:=D_{i,j}z+f_{i,j}(z)$$
where $`f_{i,j}`$ is a germ of holomorphic function at the origin with $`f_{i,j}(0)=0`$ and $`Df_{i,j}(0)=0`$.
Let us write the relation $`F_{i,j}=\rho _i\rho _j`$ and $`\rho _i\rho _i=Id`$. We obtain
$`f_{i,j}(z)`$ $`=`$ $`B_i\overline{R}_j(z)+R_i(\overline{\rho }_j)`$ (2)
$`0`$ $`=`$ $`B_i\overline{R}_i(z)+R_i(\overline{\rho }_i).`$ (3)
By multiplying the first equation by $`\overline{B}_i`$, we obtain
$$\overline{R}_j(z)=\overline{B}_if_{i,j}(z)\overline{B}_iR_i(\overline{\rho }_j).$$
Hence,we have
$$0=B_j\overline{B}_if_{i,j}(z)B_j\overline{B}_iR_i(\overline{\rho }_j)+B_i\overline{f}_{i,j}(\overline{\rho }_j)B_i\overline{R}_i(\rho _j\rho _j).$$
Let us mupltiply by $`\overline{B}_i`$ on the left and take the conjugation. We obtain
$$0=D_{i,j}B_i\overline{f}_{i,j}(\overline{z})D_{i,j}B_i\overline{R}_i(\rho _j)+f_{i,j}(\rho _j)R_i(\overline{z}).$$
On the other hand, by evaluating equation $`(\text{3})`$ at $`\overline{\rho }_j`$, we obtain
$$0=B_i\overline{R}_i(\overline{\rho }_j)+R_i(\overline{F}_{i,j}).$$
At the end,we obtain
$$R_i(\overline{z})D_{i,j}R_i(\overline{F}_{i,j})=D_{i,j}B_i\overline{f}_{i,j}(\overline{z})+f_{i,j}(\rho _j).$$
(4)
###### Definition 3.1.
The $`\rho _i`$’s are simultaneously normalizable whenever $`R_i(\overline{z})D_{i,j}R_i(\overline{D}_{i,j}\overline{z})=0`$ for all $`1i,jl`$.
###### Remark 3.1.
If the group $`G`$ is holomorphically linearizable at the origin then the $`\rho _i`$’s are simultaneously normalizable.
Moreover, asssume the $`D_{i,j}`$’s are simultaneously diagonalizable and let us set $`D_{i,j}=\text{diag}(\mu _{i,j,k})`$. Then, for any $`1kn`$ and any $`1jm`$, the $`k`$-component $`\rho _{i,k}`$ of $`\rho _i`$ can be written as
$$\left(\rho _i(z)B_i\overline{z}\right)_k=\underset{\begin{array}{c}QN_2^n\\ j,\overline{\mu }_{i,j}^Q=\mu _{i,j,k}^1\end{array}}{}\rho _{i,k,Q}\overline{z}^Q.$$
Here, $`(f)_k`$ denotes the kth-component of $`f`$.
As a consequence, we have
###### Theorem 3.1.
Let us assume that the group $`G`$ associated to the family of totally real submanifolds $`M`$ is a semi-simple Lie group. Then the $`\rho _i`$’s are simultaneously and holomorphically normalizable in a neighbourhood of the origin.
###### Proof.
It is classical \[Kus67, GS68, CG97\] that if the Lie group $`G`$ of germs of difféomorphisms at a common fixed point is semi-simple then it is holomorphically linearizable in a neighbourhood of the origin. Then, apply the previous remark 3.1. ∎
###### Definition 3.2.
We shall say that such a family $`M=\{M_i\}_{i=1,\mathrm{},m}`$ of totally real $`n`$-submanifold of $`(C^n,0)`$ intersecting at the origin is commutative if the group $`G`$ is abelian.
¿From now on, we shall assume that $`M`$ is commutative and that the family $`D`$ of linear part of the group $`G`$ at the origin is diagonal. In other words, $`D_{i,j}=\text{diag}(\mu _{i,j,k})`$. Let $``$ be a monomial ideal of $`𝒪_n`$. It is genrated by some monomials $`x^{R_1},\mathrm{},x^{R_p}`$. We shall denote $`\overline{}`$ the ideal of $`C[[\overline{x}_1,\mathrm{},\overline{x}_n]]`$ generated by $`\overline{x}^{R_1},\mathrm{},\overline{x}^{R_p}`$.
###### Definition 3.3.
1. We shall say that the family $`M`$ of manifolds is non-resonnant whenever, for all $`1im`$, $`1kn`$ and for all $`QN_2^n`$, there exists a $`1jm`$ such that $`\overline{\mu }_{i,j}^Q\mu _{i,j,k}^1`$.
2. We shall say that the family $`M`$ of manifolds non-resonnant on $``$ whenever for all monomial $`z^Q`$ not belonging to $``$ and for all couple $`(i,k)`$, there exists $`j`$ such that $`\overline{\mu }_{i,j}^Q\mu _{i,j,k}^1`$.
###### Theorem 3.2.
Assume that the group $`G`$ is abelian. Let $``$ be a monomial ideal (resp. properly imbedded) left invariant by the family $`D:=\{D_{i,j}\}`$ and the $`B_i`$’s . Assume that $`D`$ is diophantine (resp. on $``$) and that $`M`$ is non-resonnant on $``$. Assume $`G`$ is formally linearizable on $``$. Then, the family $`F`$ is holomorphically linearizable on $``$. Moreover, in these coordinates, the $`\rho _i`$’s are anti-linearized on $`\overline{}`$.
###### Proof.
By theorem 2.1, the family $`F`$ is holomorphically linearized on $``$. Let us show that, in these coordinates, the $`\rho _i`$’s are anti-linearized on $`\overline{}`$.
Let us prove by induction on $`|Q|2`$ that $`\{\rho _{i,k}\}_Q=0`$ whenever $`z^Q`$ doesn’t belong to $``$ and $`\overline{\mu }_{i,j}^Q\mu _{i,j,k}^1`$. We recall that $`\{\rho _{i,k}\}_Q`$ denotes the coefficient of $`\overline{z}^Q`$ in the Taylor expansion of $`\rho _{i,k}`$. Assume it is case up to order $`k`$. Let $`QN_2^n`$ with $`|Q|=k+1`$. Let us compute $`\{\rho _{i,k}\}_Q`$. Using equation $`(\text{4})`$, we obtain
$`R_i(\overline{z})D_{i,j}R_i(\overline{D}_{i,j}\overline{z})`$ $`=`$ $`D_{i,j}B_i\overline{f}_{i,j}(\overline{z})+f_{i,j}(B_j\overline{z})`$
$`+D_{i,j}\left(R_i(\overline{D}_{i,j}\overline{z})R_i(\overline{F}_{i,j}\overline{z})\right)`$ $`+\left(f_{i,j}(\rho _j)f_{i,j}(B_j\overline{z})\right).`$
Moreover, $`F`$ is linearized on $`V()`$. Hence, both $`\{D_{i,j}B_i\overline{f}_{i,j}(\overline{z})+f_{i,j}(B_j\overline{z})\}_Q`$ and $`\{R_i(\overline{D}_{i,j}\overline{z})R_i(\overline{F}_{i,j}\overline{z})\}_Q`$ vanish when $`z^Q`$ doesn’t belong to $``$. Hence, if $`z^Q`$, then we have
$$(1\mu _{i,j,k}\overline{\mu }_{i,j}^Q)R_{Q,i,k}=\{(f_{i,j,k}(\rho _j)f_{i,j,k}(B_j\overline{z})\}_Q.$$
But by induction, we have
$$\{(f_{i,j,k}(\rho _j)f_{i,j,k}(B_j\overline{z})\}_Q=\{Df_{i,j,k}(B_j\overline{z})R_j+Df_{i,j,k}^2(B_j\overline{z})R_j^2+\mathrm{}\}_Q=0.$$
Therefore, since $`(1\mu _{i,j,k}\overline{\mu }_{i,j}^Q)0`$, then we have $`R_{Q,i,k}=0`$. That is,
$$\rho _i(z)=B_i\overline{z}mod\overline{}.$$
###### Corollary 3.1.
Under the assumptions of theorem 3.2, there exists a complex analytic subvariety $`𝒮`$ passing throught the origin and intersecting each totally real submanifold $`M_i`$. In good holomorphic coordinate system, $`𝒮`$ is a finite intersection of a finite union of complex hyperplane defined by complex coordinate subspaces :
$$𝒮=_i_j\{z_{i_j}=0\}.$$
The intersection $`M_k𝒮`$ is then given by
$$M_k𝒮=\{z_i_j\{z_{i_j}=0\}|B_k\overline{z}=z\}.$$
###### Proof.
The complex analytic subvariety $`𝒮`$ is nothing but $`V()`$. The trace of it on $`M_i`$ is the fixed points set of $`\rho _i`$ belonging to $`V()`$. It is non void since it contains the origin. According to the previous theorem, the $`\rho _i`$’s are holomorphically and simultaneously linearizable on $`V()`$. By assumptions, $``$ is a monomial ideal so $`V()`$ is a finite intersection of a finite union of hyperplane defined by coordinate subspaces :
$$𝒮=_i_j\{z_{i_j}=0\}.$$
###### Corollary 3.2.
Assume that the family $`M`$ is non-resonnant, $`G`$ is formally linearizable and $`D`$ is diophantine. Then, in a good holomorphic coordinates system, $`M`$ is composed of linear totally real subspaces
$$\underset{i}{}\left\{zC^n\right|B_i\overline{z}=z\}.$$
###### Remark 3.2.
If the family $`M`$ is non-resonnant and if for all $`(i,k)`$, one of the eigenvalues $`\mu _{i,j,k}`$’s belong to the unit circle, then $`G`$ is formally linearizable. In fact, for any $`QN_2^n`$, any $`1im`$, any $`1kn`$, there exists $`1jm`$ such that
$$\overline{\mu }_{i,j}^Q\mu _{i,j,k}^1=\overline{\mu }_{i,j,k}.$$
This means precisely that $`D`$ is non-resonnant in the classical sense. There is no obstruction to formal linearization.
###### Corollary 3.3.
Let $``$ be the ideal generated by the monomials $`x^{R_1},\mathrm{},x^{R_p}`$ generating the ring $`\widehat{𝒪}_n^D`$ of formal invariants of $`D`$. We assume that the non-linear centralizer of $`D`$ is generated by the same monomials. If $`D`$ is diophantine on $``$ then, in a good holomorphic coordinate system, we have
$$V()=\{z(C^n,0)|z^{R_1}=\mathrm{}=z^{R_p}=0\},$$
and
$$\rho _{i|V()}(z)=B_i\overline{z}.$$
###### Corollary 3.4.
Let us consider two totally real $`n`$-manifolds of $`(C^n,0)`$ not intersecting transversally at the origin. Assume that the $`l`$ first eignevalues of $`DF(0)`$ are one. Let $`\mu ^{R_1}=1,\mathrm{},\mu ^{R_p}=1`$ be the other (i.e. $`R_iN^n`$ and $`|R_i|>1`$) generators of resonnant relations. Let
$$V()=\{z(C^n,0)|z_1=\mathrm{}=z_l=z^{R_1}=\mathrm{}=z^{R_p}=0\}.$$
If $`DF(0)`$ is diophantine on $`V()`$, then in good holomorphic coordinate system,
$$M_iV()=\left\{zV()\right|B_i\overline{z}=z\},i=1,2.$$
## 4 Real analytic manifolds with CR singularities
Les us consider a $`(n+p1)`$-real analytic submanifold $`M`$ of $`C^n`$ of the form
$$\{\begin{array}{cc}z_1=x_1+iy_1\hfill & \\ \mathrm{}\hfill & \\ z_p=x_p+iy_p\hfill & \\ y_{p+1}=F_{p+1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ \mathrm{}\hfill & \\ y_{n1}=F_{n1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ z_n=G(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \end{array}$$
(5)
where we have set $`z^{}=(z_1,\mathrm{},z_p)`$, $`z^{\prime \prime }=(z_{p+1},\mathrm{},z_{n1})`$. The real analytic real (resp. complex) valued functions $`F_i`$ (resp. $`G`$) are assumed to vanish at the origin as well as its derivative. The tangent space at the origin contains the complex subspace defined by $`z^{}`$.
First of all, we shall show, under some assumptions, that there is a good holomorphic coordinate systems in which the $`F_i`$’s are of order greater than or equal to $`3`$ and the $`2`$-jet of the $`G`$’s depends only on $`z^{}`$ and its conjugate. When $`p=1`$, this was done by E. Bishop ($`n=2`$)\[Bis65\], and also by J. Moser and S. Webster ($`n2`$)\[MW83\].
### 4.1 Preparation
Let us consider the $`2`$-jet $`G^2`$ of $`G`$. It can be written as a sum of a quadratic polynomial
$$Q(z^{},\overline{z}^{}):=d_{i,l}z_i^{}\overline{z}_l^{}+\underset{1i,lp}{}e_{i,l}z_i^{}z_l^{}+\underset{1i,lp}{}f_{i,l}\overline{z}_i^{}\overline{z}_l^{}$$
and
$$\mathrm{\Sigma }:=a_{\alpha ,\beta }x_\alpha ^{\prime \prime }x_\beta ^{\prime \prime }+b_{\alpha ,i}x_\alpha ^{\prime \prime }z_i^{}+c_{\alpha ,i}x_\alpha ^{\prime \prime }\overline{z}_i^{}.$$
###### Definition 4.1.
Let
$$F:=\left(\begin{array}{ccc}f_{1,1}& \mathrm{}& f_{p,1}\\ \mathrm{}& & \mathrm{}\\ f_{1,p}& \mathrm{}& f_{p,p}\end{array}\right)$$
be the matrix associated to the so normalized $`Q`$. The symmetric matrix $`S=\frac{1}{2}(F+^tF)`$ will be called a Bishop matrix.
First of all, by a linear change of the coordinates $`z^{}`$, we can diagonalize the symmetric matrix $`S`$.
Our first non-degeneracy condition is that the sesquilinear part $`\stackrel{~}{H}:=d_{i,l}z_i^{}\overline{z}_l^{}`$ of $`Q`$ is not reduced to $`0`$. Let
$$\stackrel{~}{H}:=\underset{zC^p|z=1}{sup}\frac{|d_{j,i,l}z_i^{}\overline{z}_l^{}|}{z^{}^2}0.$$
to its norm. Let us set $`Z_n=z_n/\stackrel{~}{H}`$. Then the sequilinear part $`H`$ of the new $`Q`$ is of norm $`1`$. By setting
$$Z_n:=z_n+(f_{j,i,l}e_{j,i,l})z_i^{}z_l^{},$$
we transform $`Q`$ in the following form :
$$Q(z^{},\overline{z}^{})=H(z,^{}\overline{z}^{})+\underset{i=1}{\overset{p}{}}\gamma _i((z_i^{})^2+(\overline{z}_i^{})^2).$$
(6)
###### Definition 4.2.
The eigenvalues $`\gamma _1,\mathrm{},\gamma _n`$ of the so normalized Bishop matrix $`S`$ will be called the generalized Bishop invariants.
Let us show that, by an holomorphic change of coordinates, we can get rid of $`\mathrm{\Sigma }`$. First of all, let us get rid of the third member of $`\mathrm{\Sigma }`$ by a change of the form
$$\zeta _i^{}z_i^{}:=\zeta _i^{}+\underset{\gamma =p+1}{\overset{n1}{}}A_{i,\gamma }z_\gamma ^{\prime \prime },i=1,\mathrm{},p.$$
We have
$`G(z^{},\overline{z}^{},x^{\prime \prime })`$ $`=`$ $`{\displaystyle a_{\alpha ,\beta }x_\alpha ^{\prime \prime }x_\beta ^{\prime \prime }}+{\displaystyle b_{\alpha ,i}x_\alpha ^{\prime \prime }\left(\zeta _i^{}+\underset{\gamma =p+1}{\overset{n1}{}}A_{i,\gamma }z_\gamma ^{\prime \prime }\right)}`$
$`+{\displaystyle c_{\alpha ,i}x_\alpha ^{\prime \prime }\left(\zeta _i^{}+\underset{\gamma =p+1}{\overset{n1}{}}\overline{A}_{i,\gamma }\overline{z}_\gamma ^{\prime \prime }\right)}`$
$`+{\displaystyle d_{i,l}\left(\zeta _i^{}+\underset{\gamma =p+1}{\overset{n1}{}}A_{i,\gamma }z_\gamma ^{\prime \prime }\right)\left(\overline{\zeta }_l^{}+\underset{\gamma =p+1}{\overset{n1}{}}\overline{A}_{l,\gamma }\overline{z}_\gamma ^{\prime \prime }\right)}`$
$`+{\displaystyle f_{i,l}\left(\zeta _i^{}+\underset{\gamma =p+1}{\overset{n1}{}}A_{i,\gamma }z_\gamma ^{\prime \prime }\right)\left(\zeta _l^{}+\underset{\gamma =p+1}{\overset{n1}{}}A_{l,\gamma }z_\gamma ^{\prime \prime }\right)}`$
$`+{\displaystyle f_{i,l}\left(\overline{\zeta }_i^{}+\underset{\gamma =p+1}{\overset{n1}{}}\overline{A}_{i,\gamma }\overline{z}_\gamma ^{\prime \prime }\right)\left(\overline{\zeta }_l^{}+\underset{\gamma ^{}=p+1}{\overset{n1}{}}\overline{A}_{l,\gamma ^{}}\overline{z}_\gamma ^{}^{\prime \prime }\right)}`$
$`+\text{higher order terms.}`$
The coefficient of $`x_\alpha ^{\prime \prime }\overline{\zeta }_i`$ is
$$c_{\alpha ,i}+\underset{s}{}d_{s,i}A_{s,\alpha }+\underset{s}{}(f_{i,s}+f_{s,i})\overline{A}_{s,\alpha }$$
We assume that we can solve the set of equations
$`c_{\alpha ,i}`$ $`=`$ $`{\displaystyle \underset{s}{}}d_{s,i}A_{s,\alpha }+{\displaystyle \underset{s}{}}(f_{i,s}+f_{s,i})\overline{A}_{s,\alpha }`$
$`\overline{c}_{\alpha ,i}`$ $`=`$ $`{\displaystyle \underset{s}{}}\overline{d}_{s,i}\overline{A}_{s,\alpha }+{\displaystyle \underset{s}{}}(\overline{f}_{i,s}+\overline{f}_{s,i})A_{s,\alpha }.`$
For each $`p+1\alpha n1`$, let us set
$$A_\alpha :=\left(\begin{array}{c}A_{1,\alpha }\\ \mathrm{}\\ A_{p,\alpha }\end{array}\right),C_\alpha :=\left(\begin{array}{c}c_{\alpha ,1}\\ \mathrm{}\\ c_{\alpha ,p}\end{array}\right),D:=\left(\begin{array}{ccc}d_{1,1}& \mathrm{}& d_{p,1}\\ \mathrm{}& & \mathrm{}\\ d_{1,p}& \mathrm{}& d_{p,p}\end{array}\right).$$
The previous equations can be written as
$$\{\begin{array}{cc}C_\alpha =DA_\alpha +(F+^tF)\overline{A}_\alpha \hfill & \\ \overline{C}_\alpha =\overline{D}\overline{A}_\alpha +(\overline{F}+^t\overline{F})A_\alpha \hfill & \end{array}$$
(7)
These systems can be solved whenever
$$det\left(4^1S^1D\overline{S}^1\overline{D}I_p\right)0,$$
(8)
where $`S:=1/2(F+^tF)`$ is the Bishop matrix.
###### Remark 4.1.
When $`p=1`$, this condition reads $`4\gamma ^21`$ where $`\gamma `$ is the Bishop invariant \[Bis65\].
Now let us remove the two first terms of the $`\mathrm{\Sigma }`$’s (once we have done the previous change of coordinates, the coefficients appearing in $`\mathrm{\Sigma }`$ may have changed). Let us set
$$Z_n:=z_na_{\alpha ,\beta }z_\alpha ^{\prime \prime }z_\beta ^{\prime \prime }b_{\alpha ,i}z_\alpha ^{\prime \prime }z_i^{}.$$
Hence, we have
$`Z_n`$ $`=`$ $`G(z^{},\overline{z}^{},x^{\prime \prime }){\displaystyle a_{\alpha ,\beta }z_\alpha ^{\prime \prime }z_\beta ^{\prime \prime }}{\displaystyle b_{\alpha ,i}z_\alpha ^{\prime \prime }z_i^{}}`$
$`=`$ $`Q(z^{},\overline{z}^{})+\mathrm{\Sigma }{\displaystyle a_{\alpha ,\beta }z_\alpha ^{\prime \prime }z_\beta ^{\prime \prime }}{\displaystyle b_{\alpha ,i}z_\alpha ^{\prime \prime }z_i^{}}`$
$`+\text{higher order terms}.`$
But we have
$$z_\alpha ^{\prime \prime }=x_\alpha ^{\prime \prime }+iF_\alpha (z^{},\overline{z}^{},x^{\prime \prime })$$
where $`F_\alpha `$ is of order greater than or equal to two. Hence, the the 2-jet of $`Z_n`$ is precisely $`Q(z^{},\overline{z}^{},x^{\prime \prime })`$.
Let us consider the $`F_\alpha `$’s. As above, the $`2`$-jet of $`F_\alpha `$ can be written as the sum of
$$Q_\alpha (z^{},\overline{z}^{}):=d_{\alpha ,i,l}z_i^{}\overline{z}_l^{}+e_{\alpha ,i,l}z_i^{}z_l^{}+\overline{e}_{\alpha ,i,l}\overline{z}_i^{}\overline{z}_l^{}$$
and
$$\mathrm{\Sigma }_\alpha :=a_{\alpha ,\gamma ,\beta }x_\gamma ^{\prime \prime }x_\beta ^{\prime \prime }+2\text{Re}b_{\alpha ,\gamma ,i}x_\gamma ^{\prime \prime }z_i^{},$$
where $`a_{\alpha ,\gamma ,\beta }`$ is a real number, $`d_{\alpha ,i,l}=\overline{d}_{\alpha ,l,i}`$. Let us show that, under some other assumption, we can get rid of the $`z_i^{}\overline{z}_j^{}`$’s terms in the $`Q_\alpha ^{}s`$, $`\alpha =p+1,\mathrm{},n1`$. In fact, let us set
$$z_\alpha Z_\alpha :=z_\alpha +ib_\alpha z_n^{\prime \prime \prime },\alpha =p+1,\mathrm{},n1.$$
In the new coordinates, the $`2`$-jet of $`Z_\alpha `$ is
$$iQ_\alpha (z^{},\overline{z}^{})+ib_\alpha Q(z^{},\overline{z}^{})+i\mathrm{\Sigma }_\alpha .$$
Thus, in order that the coefficients of the $`z_i^{}\overline{z}_j^{}`$’s vanish, one must satisfies the second non-degeneracy condition :
each $`Q_\alpha `$ is proportional to $`Q`$. (9)
Let us show that we can get rid of the quadratic part of $`F_\alpha `$ by the using following change of coordinates :
$$z_\alpha Z_\alpha :=z_\alpha 2ib_{\alpha ,\gamma ,i}z_\gamma ^{\prime \prime }z_i^{}ia_{\alpha ,\gamma ,\beta }z_\gamma ^{\prime \prime }z_\beta ^{\prime \prime }2ie_{\alpha ,i,l}z_i^{}z_l^{}.$$
In fact, we have
$$\text{Im}(Z_\alpha )=\text{Im}\left(z_\alpha 2ib_{\alpha ,\gamma ,i}z_\gamma ^{\prime \prime }z_i^{}ia_{\alpha ,\gamma ,\beta }z_\gamma ^{\prime \prime }z_\beta ^{\prime \prime }2ie_{\alpha ,i,l}z_i^{}z_l^{}\right)$$
The $`2`$-jet of the right hand side is
$$Q_\alpha (z^{},\overline{z}^{},X^{\prime \prime })+\mathrm{\Sigma }_\alpha (z^{},\overline{z}^{},X^{\prime \prime })+\text{Im}\left(2ib_{\alpha ,\gamma ,i}X_\gamma ^{\prime \prime }z_i^{}ia_{\alpha ,\gamma ,\beta }X_\gamma ^{\prime \prime }X_\beta ^{\prime \prime }2ie_{\alpha ,i,l}z_i^{}z_l^{}\right).$$
It is equal to
$$d_{\alpha ,i,l}z_i^{}\overline{z}_l^{}$$
since
$$\text{Im}\left(2ie_{\alpha ,i,l}z_i^{}z_l^{}\right)=2\text{Re}e_{\alpha ,i,l}z_i^{}z_l^{}.$$
We summarize these results in the following lemma.
###### Lemma 4.1.
Under assumptions $`(\text{8})`$ and $`(\text{9})`$, the manifold $`(\text{5})`$ can be transformed, by an holomorphic change of variables, to
$$\{\begin{array}{cc}z_1=x_1+iy_1\hfill & \\ \mathrm{}\hfill & \\ z_p=x_p+iy_p\hfill & \\ y_{p+1}=f_{p+1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ \mathrm{}\hfill & \\ y_{n1}=f_{n1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ z_n=Q(z^{},\overline{z}^{})+g(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \end{array}$$
(10)
where the $`f_i`$’s and $`g`$ are germs of real analytic functions at the origin and of order greater than or equal to $`3`$ there. The quadratic polynomial $`Q`$ is of the form
$$Q(z^{},\overline{z}^{})=d_{i,l}z_i^{}\overline{z}_l^{}+\gamma _i((z_i^{})^2+(\overline{z}_i^{})^2),$$
the norm of the sesquilinear part of $`Q`$ being $`1`$.
In the next two sections, we shall adapt the constuction of Moser and Webster to our context.
### 4.2 Complexification
Let us complexify such a manifold $`M`$ by replacing $`\overline{z}_i`$ by $`w_i`$ in order to obtain a complex analytic $`(n+p1)`$-manifold $``$ of $`C^{2n}`$:
$$\{\begin{array}{cc}2x_\alpha =z_\alpha +w_\alpha ,\hfill & \\ z_\alpha w_\alpha =2iF_\alpha (z^{},w^{},x^{\prime \prime })=2i\overline{F}_\alpha (w^{},z^{},x^{\prime \prime }),\hfill & \\ z_n=G(z^{},w^{},x^{\prime \prime }),\hfill & \\ w_n=\overline{G}(w^{},z^{},x^{\prime \prime }),\hfill & \end{array}$$
(11)
where $`\alpha `$ ranges from $`p+1`$ to $`n1`$. In this situation, we have $`G(z^{},w^{},x\mathrm{"})=Q(z^{},w^{})+g(z^{},w^{},x\mathrm{"})`$ where $`g`$ as well as the $`F_\alpha `$’s are of order greater than or equal to $`3`$. As usual if $`G(y)=_Qg_Qy^Q`$ is a formal power series in $`C^k`$, then $`\overline{G}`$ denotes the formal power series $`_Q\overline{g}_Qy^Q`$. These equations imply that
$$z_\alpha =x_\alpha +iF_\alpha (z^{},w^{},x^{\prime \prime });w_\alpha =x_\alpha iF_\alpha (z^{},w^{},x^{\prime \prime }).$$
The variables $`(z^{},w^{},x^{\prime \prime })`$ may be used as complex coordinates. Let us define the anti-holomorphic involution $`\rho `$ of $`C^{2n}`$ by
$$\rho (z,w)=(\overline{w},\overline{z}).$$
A complex anaytic manifold $``$ of $`C^{2n}`$ comes from a manifold $`M`$ whenever it is preserved by $`\rho `$. In this case, $`M=\text{Fix}(\rho )`$. The restriction to $``$ of this map is the anti-holomorphic involution $`\rho (z^{},w^{},x^{\prime \prime })=(\overline{w}^{},\overline{z}^{},\overline{x}^{\prime \prime })`$.
The two projections $`\pi _1(z,w)=z`$ and $`\pi _2(z,w)=w`$, when restricted to $``$ have the form
$`\pi _1(z,w)`$ $`=`$ $`(z^{},x_\alpha +iF_\alpha (z^{},w^{},x^{\prime \prime }),G(z^{},w^{},x^{\prime \prime }))`$
$`\pi _2(z,w)`$ $`=`$ $`(w^{},x_\alpha iF_\alpha (z^{},w^{},x^{\prime \prime }),\overline{G}(w^{},z^{},x^{\prime \prime })).`$
Let us define the holomorphic involution $`\tau _1(z,w)=(\stackrel{~}{z},\stackrel{~}{w})`$ (resp.$`\tau _2(z,w)=(\stackrel{~}{z},\stackrel{~}{w})`$) on $``$ by $`w=\stackrel{~}{w}`$ (resp. $`z=\stackrel{~}{z}`$). This leads to the following equations :
$`\stackrel{~}{w}^{}`$ $`=`$ $`w^{}`$
$`\stackrel{~}{x}_\alpha +iF_\alpha (\stackrel{~}{z}^{},w^{},\stackrel{~}{x}^{\prime \prime })`$ $`=`$ $`x_\alpha +iF_\alpha (z^{},w^{},x^{\prime \prime }),\alpha =p+1,\mathrm{},n1`$
$`\overline{G}(w^{},\stackrel{~}{z}^{},\stackrel{~}{x}^{\prime \prime })`$ $`=`$ $`\overline{G}(w^{},z^{},x^{\prime \prime }).`$
By the implicit function theorem, there exist an holomorphic function $`\mathrm{\Gamma }_\alpha `$ such that
$$\stackrel{~}{x}_\alpha =\mathrm{\Gamma }_\alpha (z^{},\stackrel{~}{z}^{},w^{},\stackrel{~}{x}^{\prime \prime })=x_\alpha \text{mod}^2\alpha =p+1,\mathrm{},n1.$$
Here, $``$ denotes the maximal ideal of germ of holomorphic functions in $`C^{n+2p1}`$ at $`0`$. Hence, we have
$$\overline{G}(w^{},\stackrel{~}{z}^{},\mathrm{\Gamma }^{\prime \prime }(z^{},\stackrel{~}{z}^{},w^{},\stackrel{~}{x}^{\prime \prime }))=\overline{G}(w^{},z^{},x^{\prime \prime }).$$
(12)
Assume that the Bishop matrix is invertible.
Let us solve, in $`\stackrel{~}{z}^{}`$, the following equation :
$$Q(z^{},w^{})=Q(\stackrel{~}{z}^{},w^{}),$$
where $`Q`$ is the quadratic part of $`G`$. Therefore, we have
$$Q(\stackrel{~}{z}^{},w^{})Q(z^{},w^{})=D_zQ(z^{},w^{})(\stackrel{~}{z}^{}z^{})+\frac{1}{2}D_z^2Q(z^{},w^{})(\stackrel{~}{z}^{}z^{})^2.$$
Since $`\frac{^2Q}{z_jz_k}=0`$ if $`jk`$, we obtain that, for any $`i_0`$,
$$\frac{Q}{z_{i_0}}(z^{},w^{})+\frac{1}{2}\frac{^2Q}{z_{i_0}^2}(z^{},w^{})(\stackrel{~}{z}_{i_0}z_{i_0})=0.$$
But, for any $`1i_0p`$, we $`1/2\frac{^2Q}{z_{i_0}^2}(w^{},z^{})=\gamma _{i_0}0`$. Hence, we have
$$\stackrel{~}{z}_{i_0}=z_{i_0}\frac{1}{\gamma _{i_0}}\underset{l=1}{\overset{p}{}}d_{i_0,l}w_l.$$
Hence, by the implicit function theorem, equation $`(\text{12})`$ admit an holomorphic solution $`\tau _1(z^{},w^{},x^{\prime \prime })`$ which linear part at the origin is precisely $`T_1z^{}`$ :
$$\stackrel{~}{z}^{}=H(z^{},w^{},x^{\prime \prime })=T_1z^{}\text{mod}^2,\stackrel{~}{w}^{}=w^{},\stackrel{~}{x}^{\prime \prime }=\mathrm{\Gamma }(z^{},w^{},x^{\prime \prime }).$$
Here $``$ denotes the maximal ideal of germs of holomorphic functions of $`(C^{n+p1},0)`$ and the linear part $`T_1`$ is
$$T_1:=D\tau _1(0)=\left(\begin{array}{ccc}Id_p& S^1D& 0\\ 0& Id_p& 0\\ 0& 0& Id_q\end{array}\right),$$
where $`D`$ is the matrix $`(d_{i,j})_{1i,jp}`$ and $`I_q`$ stands for the identity matrix of dimension $`q:=np1`$. The map $`\pi _2`$ is a two-fold covering with covering transformation $`\tau _1`$.
Since $`\tau _2=\rho \tau _1\rho `$, $`\tau _2`$ restricted to $``$ is defined by
$$\stackrel{~}{z}^{}=z^{},\stackrel{~}{w}^{}=\overline{H}(w^{},z^{},x^{\prime \prime }),\stackrel{~}{x}^{\prime \prime }=\overline{\mathrm{\Gamma }}(w^{},z^{},x^{\prime \prime }).$$
The linear part at the origin of $`\tau _2`$ is
$$T_2:=\left(\begin{array}{ccc}Id_p& 0& 0\\ \overline{S}^1\overline{D}& Id_p& 0\\ 0& 0& Id_q\end{array}\right).$$
The map $`\pi _1`$ is a two-fold covering with covering transformation $`\tau _2`$.
Let us consider the germ of holomorphic diffeomorphism $`g`$ of $`(C^{n+p1},0)`$ defined to be $`g:=\tau _1\tau _2`$. It fixes the origin. The linear part of the diffeomorphism $`g`$ at the origin is
$$\mathrm{\Phi }:=Dg(0)=T_1T_2=\left(\begin{array}{ccc}Id_p+S^1D\overline{S^1D}& S^1D& 0\\ \overline{S^1D}& I_p& 0\\ 0& 0& I_q\end{array}\right).$$
### 4.3 Quadrics and linear involutions
In this section will shall study the relation between the linear part of the involutions, the linear anti-holomorphic involution and the quadric.
Under our assumptions, the set of fixed points of $`T_1`$ and $`T_2`$ is a $`q`$-dimensional vector space ($`q=np1`$) and these are the only common eigenvectors. In fact, let us try to solve $`T_1v=av`$ and $`T_2v=bv`$. Let us write $`v=(v_1,v_2,v_3)`$ with $`v_1,v_2`$ belong to $`C^p`$ while $`v_3`$ belongs to $`C^q`$. This leads to
$$()\{\begin{array}{c}v_1S^1Dv_2=av_1\hfill \\ v_2=av_2\hfill \\ v_3=av_3\hfill \end{array}()\{\begin{array}{c}v_1=bv_1\hfill \\ \overline{S}^1\overline{D}v_1v_2=bv_2\hfill \\ v_3=bv_3\hfill \end{array}$$
Assume $`v_3=0`$. If $`v_2=0`$ then $`a=1`$ (otherwise $`v_1=0`$ too). The second equation of $`()`$ gives $`\overline{S}^1\overline{D}v_1=0`$; that is $`v_1=0`$. This is not possible. Thus $`v_20`$ and $`a=1`$. The first equation of $`()`$ gives $`S^1Dv_2=2v_1`$. Moreover, we have $`b=1`$ since otherwise we would have $`v_1=0`$ and $`v_2=0`$ by the second equation of $`()`$. Therefore, using the same equation we obtain
$$\overline{S}^1\overline{D}S^1Dv_2=4v_2.$$
According to condition $`(\text{8})`$, we have $`v_2=0`$ and then $`v_1=0`$. Now, if $`v_30`$ then $`a=b=1`$, then we can apply the previous resonning to obtain $`v_1=v_2=0`$.
Let $`V_i`$ be the $`(1)`$-eigenspace of $`T_i`$, $`i=1,2`$. Let $`E`$ be their common $`(+1)`$-eigenspace. We assume that $`V_1,V_2`$ and $`E`$ span $`C^{n+p1}`$. Let $`F=V_1+V_2`$. We have $`C^{n+p1}=FE`$. Let us show that it is invariant under both $`T_1`$ and $`T_2`$. In fact, let $`v_2V_2`$. A priori, we have $`T_1v_2=\stackrel{~}{v}_1+\stackrel{~}{v}_2+e`$ where $`\stackrel{~}{v}_1`$ (resp. $`\stackrel{~}{v}_2`$, $`e`$) belongs to $`V_1`$ (resp. $`V_2`$, $`E`$). Since $`T_1^2=Id`$ and $`T_1v_1=v_1`$, we have $`v_2=\stackrel{~}{v}_1+T_1\stackrel{~}{v}_2+e`$. But, $`T_1\stackrel{~}{v}_2=\stackrel{~}{v}_1^{}+\stackrel{~}{v}_2^{}+e^{}`$. Hence, we have $`v_2=\stackrel{~}{v}_1+\stackrel{~}{v}_1^{}+\stackrel{~}{v}_2^{}+e^{}+e`$. Therefore, $`\stackrel{~}{v}_1+\stackrel{~}{v}_1^{}=0`$, $`\stackrel{~}{v}_2^{}v_2=0`$ and $`e^{}+e=0`$. It comes
$$T_1(v_2\stackrel{~}{v}_2e)=(v_2\stackrel{~}{v}_2e).$$
This means that $`v_2\stackrel{~}{v}_2e`$ belongs to $`V_1`$ and leads to $`e=0`$; that is $`F`$ is left invariant by $`T_1`$. The same argument applies to $`T_2`$. Let $`T_1^{}`$ (resp. $`T_2^{}`$, $`\mathrm{\Phi }^{}`$) be the restriction $`T_1`$ (resp. $`T_2`$, $`\mathrm{\Phi }`$) to $`F`$.
Let $`\mu `$ be an eigenvalue of $`\mathrm{\Phi }^{}`$ corresponding to an eigenvector $`f`$. Then, $`T_2^{}f=\mu T_1^{}f=\mu \mathrm{\Phi }^{}T_2^{}f`$. Hence, $`\mu ^1`$ is another eigenvalue of $`\mathrm{\Phi }^{}`$ associated to $`T_2^{}f`$. The vectors $`f`$ and $`T_2^{}f`$ are independant since otherwise we would have $`cf=T_2^{}f=\mu T_1^{}f`$ and $`T_1^{}`$ and $`T_2^{}`$ would have a common eigenvector. If $`\mu =\mu ^1`$, then the restriction of $`\mathrm{\Phi }^{}`$ to the span $`V_f`$ of $`f`$ and $`T_2^{}f`$ is $`\pm Id`$. This implies a common eigenvector to $`T_2^{}`$ and $`T_1^{}`$. As in \[MW83\]\[p.269\], we can choose $`\{f,T_2^{}f\}`$ as a basis of $`V_f`$. In this basis, the matrix of $`T_1^{}`$, $`T_2^{}`$ and $`\mathrm{\Phi }^{}`$ are
$$\mathrm{\Phi }_{|V_f}^{}=\left(\begin{array}{cc}\mu & 0\\ 0& \mu ^1\end{array}\right)T_{i|V_f}^{}=\left(\begin{array}{cc}0& \lambda _i\\ \lambda _i^1& 0\end{array}\right),$$
where $`\lambda _1=\mu `$ and $`\lambda _2=1`$.
Let us show that $`E`$ is invariant under the linear anti-holomorphic involution $`\rho `$. Let $`e`$ belongs to $`E`$. By definition, it is left invariant by both $`T_1`$ and $`T_2`$. Since, $`T_1\rho =\rho T_2`$, we obtain $`T_1\rho (e)=\rho (T_2e)=\rho (e)`$ and $`T_2\rho (e)=\rho (e)`$. Hence, $`\rho (E)=E`$.
Let us show that $`F`$ is invariant under the linear anti-holomorphic involution $`\rho `$. Let $`N`$ be the totally real fixed point set of $`\rho `$ on $`E`$. This means that $`E=N+iN`$ and that we may choose coordinates $`\zeta `$ on $`E`$ so that $`\rho :\zeta \overline{\zeta }`$. Let us show that $`\rho (F)`$ is also invariant by both $`T_1`$ and $`T_2`$. In fact, we have $`T_1\rho (F)=\rho T_2(F)=\rho (F)`$ and $`T_2\rho (F)=\rho (F)`$ as well. Hence, $`E\rho (F)=C^{n+p1}`$ is a decomposition preserved by both $`T_1`$ and $`T_2`$. The space $`\rho (F)`$ has to contain the $`(1)`$-eigenspace of both $`T_1`$ and $`T_2`$, that is $`F`$.
Let $`f`$ be a $`\mu `$-eigenvector of $`\mathrm{\Phi }`$. As above, $`T_2f`$ is a $`\mu ^1`$-eigenvector of $`\mathrm{\Phi }`$. Since $`\mathrm{\Phi }\rho \mathrm{\Phi }=T_1\rho T_2=\rho `$, we have
$$\rho (f)=\overline{\mu }\mathrm{\Phi }(\rho (f)).$$
Hence, $`\rho (f)`$ is a $`\overline{\mu }^1`$-eigenvector of $`\mathrm{\Phi }`$. Moreover, $`T_2\rho (f)=\rho (T_1f)=\overline{\mu }^1\rho (T_2f)`$ is a $`\overline{\mu }`$-eigenvector of $`\mathrm{\Phi }`$. This means that $`\rho (V_f)=V_{\rho (f)}`$.
Let us assume that the set $`\overline{\mu }^1`$ is different than $`\mu ^1`$ and $`\mu `$. Then, $`\rho (T_2\rho (f))=T_1f`$ is a $`\mu ^1`$-eigenvector of $`\mathrm{\Phi }`$. So, $`T_2\rho (T_2\rho (f))=T_2T_1f=\mathrm{\Phi }^1(f)=\mu ^1f`$ is a $`\mu `$-eigenvector of $`\mathrm{\Phi }`$. In this case, the matrices of $`\rho `$ and $`\mathrm{\Phi }`$ restricted to $`V_fV_{\rho (f)}`$ are
$$\rho _{|V_fV_{\rho (f)}}=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& \mu ^1\\ 1& 0& 0& 0\\ 0& \overline{\mu }& 0& 0\end{array}\right)\mathrm{\Phi }_{|V_fV_{\rho (f)}}=\left(\begin{array}{cccc}\mu & 0& 0& 0\\ 0& \mu ^1& 0& 0\\ 0& 0& \overline{\mu }^1& 0\\ 0& 0& 0& \overline{\mu }\end{array}\right).$$
We recall that $`V_f`$ denotes the span of $`f`$ and $`T_2f`$. We can do a similar analysis as in Moser-Webster article \[MW83\]\[p.269\] : let $`\mu _i`$ be an eigenvalue of $`\mathrm{\Phi }`$ of multiplicity $`n(i)`$ such that $`\mu _i=\overline{\mu }_i^1`$ or $`\mu _i=\overline{\mu }_i`$. Let $`\{f_1,\mathrm{},f_{n(i)}\}`$ be an associated basis. Let $`E_i`$ be the span of the basis $`\{f_1,\mathrm{},f_{n(i)},T_2f_1,\mathrm{},T_2f_{n(i)}\}`$. We have
* if $`\mu _i=\overline{\mu }_i^1`$, then $`\rho (f_k)=a_{j,k}f_k`$ so that
$$\rho (T_2f_k)=T_1\rho (f_k)=T_1(a_{j,k}f_k)=a_{j,k}T_1f_k=\overline{\mu }_ia_{j,k}T_2(f_k).$$
* if $`\mu _i=\overline{\mu }_i`$, then $`\rho (f_k)=a_{j,k}T_2f_k`$ so that
$$\rho (T_2f_k)=T_1\rho (f_k)=T_1(a_{j,k}T_2f_k)=\mu _ia_{j,k}f_k.$$
We have used the property that, if $`\mathrm{\Phi }(f)=\mu f`$ then $`T_2f=\mu T_1f`$. Hence we have proved the following
###### Lemma 4.2.
Let $`\mathrm{\Phi }`$, $`T_1`$, $`T_2`$ and $`\rho `$ as above. Then there exists a decomposition
$$C^{n+p1}=E_1\mathrm{}E_rE_{r+1}\mathrm{}E_sG$$
left invariant by $`\mathrm{\Phi }`$, $`T_1`$, $`T_2`$ and $`\rho `$ such that
* The $`E_i`$’s are complex vectorspaces of dimension $`2n(i)`$ if $`1ir`$ and $`4n(i)`$ otherwise. $`G`$ is a $`np1`$-dimensional vectorspace.
* The restrictions of $`\mathrm{\Phi }`$, $`T_1`$, $`T_2`$ to $`G`$ is the identity
* If $`1ir`$, then $`\mu _i\{\overline{\mu }_i,\overline{\mu }_i^1\}`$ and there exists coordinates $`(\zeta _i,\eta _i)`$ of $`E_i`$ ($`\zeta _i=(\zeta _{i,1},\mathrm{},\zeta _{i,n(i)})`$) such that
$$T_{k|E_i}=\left(\begin{array}{cc}0& \lambda _{k,i}I_{n(i)}\\ \lambda _{k,i}^1I_{n(i)}& 0\end{array}\right),k=1,2\mathrm{\Phi }_{|E_i}=\left(\begin{array}{cc}\mu _iI_{n(i)}& 0\\ 0& \mu _i^1I_{n(i)}0\end{array}\right),$$
where $`\lambda _{1,i}=\mu _i`$ and $`\lambda _{2,i}=1`$.
* If $`0<jsr`$, then $`\mu _{r+2j1}\{\overline{\mu }_{r+2j1},\overline{\mu }_{r+2j1}^1\}`$ and there exists coordinates $`\theta _j:=(\zeta _{r+2j1},\eta _{r+2j1},\zeta _{r+2j},\eta _{r+2j})`$ of $`E_{r+j}`$ such that
$$\mathrm{\Phi }_{|E_{r+j}}=\left(\begin{array}{cccc}\mu _{r+2j1}I_{n(r+j)}& 0& 0& 0\\ 0& \mu _{r+2j1}^1I_{n(r+j)}& 0& 0\\ 0& 0& \overline{\mu }_{r+2j1}^1I_{n(r+j)}& 0\\ 0& 0& 0& \overline{\mu }_{r+2j1}I_{n(r+j)}\end{array}\right).$$
and
$$T_{k|E_{r+j}}=\left(\begin{array}{cccc}0& \lambda _{k,j+2r1}I_{n(j+r)}& 0& 0\\ \lambda _{k,j+2r1}^1I_{m(j+r)}& 0& 0& 0\\ 0& 0& 0& \overline{\lambda }_{k,j+2r1}^1I_{m(j+r)}\\ 0& 0& \overline{\lambda }_{k,j+2r1}I_{m(j+r)}& 0\end{array}\right),$$
where $`k=1,2`$, $`\lambda _{1,r+2j1}=\mu _{r+2j1}`$, $`\mu _{r+2j}=\overline{\mu }_{r+2j1}^1`$ and $`\lambda _{2,r+2j1}=1`$.
* If $`\overline{\mu }_i=\mu _{i}^{}{}_{}{}^{1}`$ (we say hyperbolic), $`\rho _{|E_i}(\zeta _i,\eta _i)=\left(\begin{array}{cc}A_i& 0\\ 0& \overline{\mu }_iA_i\end{array}\right)\left(\begin{array}{c}\overline{\zeta }_i\\ \overline{\eta }_i\end{array}\right)`$ for some $`n(i)`$-square matrix $`A_i`$ such that $`A_i\overline{A}_i=I_{n(i)}`$.
* If $`\overline{\mu }_i=\mu _i`$ (we say elliptic), $`\rho _{|E_i}(\zeta _i,\eta _i)=\left(\begin{array}{cc}0& \mu _iA_i\\ A_i& 0\end{array}\right)\left(\begin{array}{c}\overline{\zeta }_i\\ \overline{\eta }_i\end{array}\right)`$ for some $`n(i)`$-square matrix $`A_i`$ such that $`\mu _iA_i\overline{A}_i=I_{n(i)}`$.
* If $`\mu _{r+2j1}\{\overline{\mu }_{r+2j1},\overline{\mu }_{r+2j1}^1\}`$, that is $`0<jsr`$, (we say complex),
$$\rho _{|E_{r+2j1}}(\theta _j)=\left(\begin{array}{cccc}0& 0& I_{n(j+r)}& 0\\ 0& 0& 0& \mu _{r+2j1}^1I_{n(j+r)}\\ I_{n(j+r)}& 0& 0& 0\\ 0& \overline{\mu }_{r+2j1}I_{n(j+r)}& 0& 0\end{array}\right)\left(\begin{array}{c}\overline{\zeta }_{r+2j1}\\ \overline{\eta }_{r+2j1}\\ \overline{\zeta }_{r+2j}\\ \overline{\eta }_{r+2j}\end{array}\right)$$
* $`\rho _{|G}(\upsilon )=\overline{\upsilon }`$.
###### Definition 4.3.
A coordinate or its associated eigenvalue will be called hyperbolic (resp. elliptic, complex) if $`\mu _i=\overline{\mu }_i^1`$, (resp. $`\mu _i=\overline{\mu }_i`$, $`\mu _i\{\overline{\mu }_i,\overline{\mu }_i^1\}`$).
###### Remark 4.2.
When $`p=1`$, the complex case doesn’t appear. In fact, in this case, the eigenspace associated to an eigenvalue $`\mu 1`$ is a $`2`$-dimensional vector space. Hence, we must have $`\mu \{\overline{\mu },\overline{\mu }^1\}`$.
### 4.4 Submanifolds with CR-singularities
Let us consider a real analytic $`(n+p1)`$-submanifolds $`M`$ of $`C^n`$ passing through the origin of the form
$$\{\begin{array}{cc}z_1=x_1+iy_1\hfill & \\ \mathrm{}\hfill & \\ z_p=x_p+iy_p\hfill & \\ y_{p+1}=F_{p+1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ \mathrm{}\hfill & \\ y_{n1}=F_{n1}(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \\ z_n=G(z^{},\overline{z}^{},x^{\prime \prime }):=Q(z^{},\overline{z}^{})+g(z^{},\overline{z}^{},x^{\prime \prime })\hfill & \end{array}$$
(13)
where the $`F_i`$’s and $`g`$ are real analytic in a neighbourhood of the origin and of order greater than or equal to $`3`$ at $`0`$. We assume that the quadratic polynomial $`Q`$ is normalized in the sense developped in the previous section. Hence, we have
$$Q=\underset{i,j}{}g_{i,j}z_i^{}\overline{z}_j^{}+\underset{i=1}{\overset{p}{}}\gamma _i((z_i^{})^2+(\overline{z}_i^{})^2),$$
where the sesquilinear part is of norm $`1`$, and the $`\gamma _i`$’s are the generalized Bishop invariants. The CR-structure is singular at the origin. For a $`n`$-variety, we have $`Q=z_1\overline{z}_1+\gamma (z_1^2+\overline{z}_1^2)`$. Let us complexify the equations as in $`(\text{11})`$:
$$:\{\begin{array}{cc}2x_\alpha =z_\alpha +w_\alpha ,\hfill & \\ z_\alpha w_\alpha =2iF_\alpha (z^{},w^{},x^{\prime \prime })=2i\overline{F}_\alpha (w^{},z^{},x^{\prime \prime }),\hfill & \\ z_n=G(z^{},w^{},x^{\prime \prime }),\hfill & \\ w_n=\overline{G}(w^{},z^{},x^{\prime \prime }),\hfill & \end{array}$$
where $`\alpha `$ ranges from $`p+1`$ to $`n1`$. Then there exists a pair $`(\tau _1,\tau _2)`$ of holomorphic involutions as defined above. It defines a germ of holomorphic diffeomorphism $`\mathrm{\Phi }:=\tau _1\tau _2`$ of $`(C^{n+p1},0)`$. We assume that it is not tangent to the identity at the origin.
###### Definition 4.4.
Let $``$ be monomial ideal of $`𝒪_{n+p1}`$. We shall say that $``$ is compatible with $`T_1`$ and $`T_2`$ (resp. with the anti-linear involution $`\rho `$) if the maps $`T_j^{}:\widehat{𝒪}_{n+p1}\widehat{𝒪}_{n+p1}`$ (resp. $`\rho ^{}:\widehat{𝒪}_{n+p1}\overline{\widehat{𝒪}_{n+p1}}`$) defined by $`T_j^{}(f)=fT_j`$ (resp. $`\rho ^{}(f)=f\rho `$ ) preserve the splitting $`\widehat{𝒪}_{n+p1}=\widehat{}\widehat{CI}`$ (resp. maps $`\widehat{}`$ to $`\overline{\widehat{}}`$ and $`\widehat{CI}`$ to $`\overline{\widehat{CI}}`$).
The main result of this section is the following.
###### Theorem 4.1.
Let $`M`$ be a third order analytic perturbation of the quadric $`\{y_{p+1}=\mathrm{}=y_{n1}=0,z_n=Q(z^{},\overline{z}^{})\}`$, $`z^{}=(z_1,\mathrm{},z_p)`$ as above. Let $`\tau _1`$, $`\tau _2`$ be the holomorphic involutions of $`(C^{n+p1},0)`$ associated to the complexified submanifold $``$ of $`C^{2n}`$. Let $``$ be a monomial ideal of $`𝒪_{n+p1}`$ compatible with $`D\tau _1(0)`$, $`D\tau _2(0)`$ and $`\rho `$. Assume that the involutions $`\tau _1`$ and $`\tau _2`$ are formally linearizable on the ideal $``$. Assume furthermore that $`D\mathrm{\Phi }(0)`$ is a diagonal matrix and is diophantine (resp. on $``$ whenever $``$ is properly embedded). Then the $`\tau _j`$’s are simultaneously and holomorphically linearizable on $``$. Moreover, the linearizing diffeomorphism can be chosen so that it commutes with the linear anti-holomorphic involution $`\rho `$.
###### Proof.
We can apply lemma 4.2 to the $`(D\mathrm{\Phi }(0),D\tau _1(0),D\tau _2(0))`$. We reorder the vectorspaces so that
$$C^{n+p1}=E_1\mathrm{}E_rE_{r+1}\mathrm{}E_sG$$
with $`E_1,\mathrm{},E_h`$ are hyperbolic, $`E_{h+1},\mathrm{},E_r`$ are elliptic and $`E_{r+1},\mathrm{},E_s`$ are complex. This means that $`\mu _i=\overline{\mu }_i^1`$ for $`1ih`$, $`\mu _i=\overline{\mu }_i`$ for $`h+1ir`$ and $`\mu _i\{\overline{\mu }_,\overline{\mu }_i^1\}`$ for $`r+1is`$.
Let us consider the ring $`\widehat{𝒪}_{n+p1}^{D\mathrm{\Phi }(0)}`$ of formal invariants of $`D\mathrm{\Phi }(0)`$. Since $`\mu _i\mu _i^1=1`$, the monomials $`\zeta _{i,u}\eta _{i,v}`$, $`1u,vn(i)`$ are invariants for $`1ir`$ as well as $`\zeta _{i,k,u}\eta _{i,k,v}`$, $`k=1,2`$, $`1u,vn(i)`$ and $`r+1is`$ and the $`\upsilon _j`$’s with $`j=1,\mathrm{},q`$ (their associated eigenvalues are $`1`$). More generally, let $`(p,q)N^{r+2s}\times N^{r+2s}`$ such that
$$\mathrm{\Pi }_{i=1}^r\mu _i^{p_iq_i}\mathrm{\Pi }_{j=r+1}^s\mu _j^{p_jq_r}\overline{\mu }_j^{p_{j+1}+q_{j+1}}=1.$$
then all the monomials
$$\mathrm{\Pi }_{i=1}^r\left(\mathrm{\Pi }_{k=1}^{n(i)}\zeta _{i,k}^{p_{i,k}}\eta _{i,k}^{q_{i,k}}\right)\mathrm{\Pi }_{j=r+1}^s\left(\mathrm{\Pi }_{k=1}^{n(j)}\zeta _{j,k}^{p_{j,k}}\eta _{j,k}^{q_{r,k}}\zeta _{j+1,k}^{p_{j+1,k}}\eta _{j+1,k}^{q_{r+1,k}}\right)$$
(14)
for which $`_{k=1}^{n(i)}p_{i,k}=p_i`$ and $`_{k=1}^{n(i)}q_{i,k}=q_i`$ ($`1ir+2s`$) are invariants (the $`p_{i,k}`$’s and $`q_{i,k}`$’s are non-negative integers). Such a monomial will be written as $`\zeta ^P\eta ^Q`$ where $`P=(\{p_{i,k}\}_{1kn(i)})_{1ir+2s}`$ and $`Q=(\{q_{i,k}\}_{1kn(i)})_{1ir+2s}`$. Let $``$ denotes the set of such $`(P,Q)`$ which generates, together with the $`\upsilon _j`$’s, the ring of invariants (in fact, it is a module of finite type; so just the generators can be selected).
Let $`ResIdeal`$ denote the ideal of $`𝒪_{np+1}`$ generated by the monomials of the set $``$.
Let us write that the $`\tau _i`$’s are simultaneously formally linearizable on $``$. To be more specific, let us write, for $`j=1,2`$
$$\tau _j(\zeta ^{},\eta ^{},\upsilon ^{})=\{\begin{array}{cc}\zeta _i^{\prime \prime }=\lambda _{j,i}\eta _i^{}+f_{j,i}(\zeta ^{},\eta ^{},\upsilon ^{})i=1,\mathrm{},r+2s\hfill & \\ \eta _i^{\prime \prime }=\lambda _{j,i}^1\zeta _i^{}+g_{j,i}(\zeta ^{},\eta ^{},\upsilon ^{})i=1,\mathrm{},r+2s\hfill & \\ \upsilon ^{\prime \prime }=\upsilon ^{}+h_j(\zeta ^{},\eta ^{},\upsilon ^{}).\hfill & \end{array},$$
and
$$\widehat{\mathrm{\Psi }}(\zeta ,\eta ,\upsilon )=\{\begin{array}{cc}\zeta _i^{}=\zeta _i+U_i(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \eta _i^{}=\eta _i+V_i(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \upsilon ^{}=\upsilon +W(\zeta ,\eta ,\upsilon ).\hfill & \end{array}.$$
Here we have used to following convention : if $`j1`$ then $`\mu _{r+2j}=\overline{\mu }_{r+2j1}^1`$ and $`\theta _j:=(\zeta _{r+2j1},\eta _{r+2j1},\zeta _{r+2j},\eta _{r+2j})`$ are coordinates on $`E_{r+j}`$. Let us write that $`\widehat{\mathrm{\Psi }}`$ conjugates $`\tau _j`$ to $`\stackrel{~}{\tau }_j`$ with
$$\stackrel{~}{\tau }_j(\zeta ,\eta ,\upsilon )=\{\begin{array}{cc}\zeta _i^{}=\lambda _{j,i}\eta _i+\stackrel{~}{f}_{j,i}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \eta _i^{}=\lambda _{j,i}^1\zeta _i+\stackrel{~}{g}_{j,i}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \upsilon ^{}=\upsilon +\stackrel{~}{h}_j(\zeta ,\eta ,\upsilon ).\hfill & \end{array},$$
where the $`\stackrel{~}{f}_{j,i}`$’s, the $`\stackrel{~}{g}_{j,i}`$’s and the components of $`\stackrel{~}{h}_j`$ belong to the ideal $``$. We have $`\widehat{\mathrm{\Psi }}\stackrel{~}{\tau }_j=\tau _j\widehat{\mathrm{\Psi }}`$; that is
$$()\{\begin{array}{cc}\lambda _{j,i}V_iU_i\stackrel{~}{\tau }_j=\stackrel{~}{f}_{j,i}f_{j,i}\widehat{\mathrm{\Psi }}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \lambda _{j,i}^1U_iV_i\stackrel{~}{\tau }_j=\stackrel{~}{g}_{j,i}g_{j,i}\widehat{\mathrm{\Psi }}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ WW\stackrel{~}{\tau }_j=\stackrel{~}{h}_jh_j\widehat{\mathrm{\Psi }}.\hfill & \end{array}$$
Let us find an equation involving only the unkown $`U_i`$ (resr. $`V_i`$, $`W`$). We recall that $`\mu _i:=\lambda _{1,i}\lambda _{2,i}^1`$. Since we have $`\lambda _{1,i}V_iU_i\stackrel{~}{\tau }_1=\stackrel{~}{f}_{1,i}f_{1,i}\widehat{\mathrm{\Psi }}`$, we have $`\lambda _{1,i}V_i\stackrel{~}{\tau }_2U_i\stackrel{~}{\mathrm{\Phi }}=\stackrel{~}{f}_{1,i}\stackrel{~}{\tau }_2f_{1,i}\widehat{\mathrm{\Psi }}\stackrel{~}{\tau }_2`$. Here, $`\stackrel{~}{\mathrm{\Phi }}`$ denotes $`\stackrel{~}{\tau }_1\stackrel{~}{\tau }_2`$. According to the second equation of $`()`$ for $`j=2`$, we have $`V_i\stackrel{~}{\tau }_2=\lambda _{2,i}^1U_i\stackrel{~}{g}_{2,i}+g_{2,i}\widehat{\mathrm{\Psi }}`$. Therefore, we obtain
$$\mu _iU_iU_i\stackrel{~}{\mathrm{\Phi }}=\left(\stackrel{~}{f}_{1,i}f_{1,i}\widehat{\mathrm{\Psi }}\right)\stackrel{~}{\tau }_2+\lambda _{1,i}\left(\stackrel{~}{g}_{2,i}g_{2,i}\widehat{\mathrm{\Psi }}\right).$$
Since the $`\tau _j`$’s are involutions, we have the following relations :
$`\lambda _{j,i}g_{j,i}+f_{j,i}\tau _j`$ $`=`$ $`0`$ (15)
$`\lambda _{j,i}^1f_{j,i}+g_{j,i}\tau _j`$ $`=`$ $`0`$ (16)
$`h_j\tau _j+h_j`$ $`=`$ $`0.`$ (17)
We have the following relations among the $`f`$’s and the $`g`$’s :
$$g_{1,i}\tau _2\mu _i^1g_{2,i}\tau _2=\lambda _{1,i}^1\left(f_{1,i}\tau _1f_{2,i}\tau _2\right)\tau _2.$$
(18)
Using the fact that $`f_{1,i}\widehat{\mathrm{\Psi }}\stackrel{~}{\tau }_2=f_{1,i}\tau _2\widehat{\mathrm{\Psi }}`$, we find the following relations :
$$\mu _iU_iU_i\stackrel{~}{\mathrm{\Phi }}=(\stackrel{~}{f}_{1,i}\mu _i\stackrel{~}{f}_{2,i})\stackrel{~}{\tau }_2(f_{1,i}\mu _if_{2,i})\tau _2\widehat{\mathrm{\Psi }}=:\alpha _i.$$
(19)
Similarly, we obtain
$$\mu _i^1V_iV_i\stackrel{~}{\mathrm{\Phi }}=(\stackrel{~}{g}_{1,i}\mu _i^1\stackrel{~}{g}_{2,i})\stackrel{~}{\tau }_2(g_{1,i}\mu _i^1g_{2,i})\tau _2\widehat{\mathrm{\Psi }}=:\beta _i,$$
(20)
as well as
$$WW\stackrel{~}{\mathrm{\Phi }}=(\stackrel{~}{h}_1\stackrel{~}{\tau }_2+\stackrel{~}{h}_2)(h_1\tau _2+h_2)\widehat{\mathrm{\Psi }}=:\gamma .$$
(21)
The germ of diffeomorphism $`\mathrm{\Phi }`$ is formally linearizable on $``$. If its linear part $`D\mathrm{\Phi }(0)`$ is diophantine (resp. on $``$ if $``$ is properly embedded), then, according to theorem 2.1, $`\mathrm{\Phi }`$ is actually holomorphically linearizable on $``$ and by a unique normalizing transformation which projection on the vector space spanned by $`\widehat{}`$ and the formal centralizer of $`D\mathrm{\Phi }(0)x`$ is zero. Let $`\mathrm{\Psi }^{}`$ be precisely this germ of holomorphic diffeomorphism of $`(C^{n+p1},0)`$ :
$$\mathrm{\Psi }^{}(\zeta ,\eta ,\upsilon )=\{\begin{array}{cc}\zeta _i^{}=\zeta _i+U_i^{}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \eta _i^{}=\eta _i+V_i^{}(\zeta ,\eta ,\upsilon )i=1,\mathrm{},r+2s\hfill & \\ \upsilon ^{}=\upsilon +W^{}(\zeta ,\eta ,\upsilon ).\hfill & \end{array}$$
Let us show that it commutes with $`\rho `$. First of all, we have $`\rho \mathrm{\Phi }\rho =\mathrm{\Phi }^1`$. Then, according to the second part of theorem 2.1, it is sufficient to show that $`\rho 𝒞_{D\mathrm{\Phi }(0)x}\rho =𝒞_{D\mathrm{\Phi }(0)x}`$. Let $`R𝒞_{D\mathrm{\Phi }(0)x}`$. We have $`D\mathrm{\Phi }(0)R(x)=RD\mathrm{\Phi }(0)(x)`$. We recall that $`D\mathrm{\Phi }(0)x=T_1T_2x`$. Therefore, we have
$$T_1T_2\rho R\rho =T_1\rho T_1R\rho =\rho T_2T_1R\rho =\rho RT_2T_1\rho =\rho R\rho T_1T_2.$$
Let us modify $`\mathrm{\Psi }^{}`$ in such a way that it linearizes simultaneously $`\tau _1,\tau _2`$ on $``$. Let $`pr(f)`$ (resp. $`pr_{}(f)`$) denote the projection of the formal power series $`f`$ on the vector space $`\widehat{CI}`$ spanned by the $`x^Q`$’s which doesn’t belong to $`\widehat{}`$ (resp. on $`\widehat{}`$). Let us project the equations $`()`$ onto this space. We find
$$()\{\begin{array}{cc}\lambda _{j,i}pr(V_i)pr(U_iT_j)=pr(f_{j,i}\widehat{\mathrm{\Psi }})i=1,\mathrm{},r+2s\hfill & \\ \lambda _{j,i}^1pr(U_i)pr(V_iT_j)=pr(g_{j,i}\widehat{\mathrm{\Psi }})i=1,\mathrm{},r+2s\hfill & \\ pr(W)pr(WT_j)=pr(h_j\widehat{\mathrm{\Psi }}).\hfill & \end{array}$$
On the other hand, projecting the equations $`(\text{19})`$,$`(\text{20})`$ and $`(\text{21})`$ on $`CI`$ we obtain
$`\mu _ipr(U_i)pr(U_i)D`$ $`=`$ $`pr\left(\left(f_{1,i}\mu _if_{2,i}\right)\tau _2\widehat{\mathrm{\Psi }}\right)`$
$`\mu _i^1pr(V_i)pr(V_i)D`$ $`=`$ $`pr\left(\left(g_{1,i}\mu _i^1g_{2,i}\right)\tau _2\widehat{\mathrm{\Psi }}\right)`$
$`pr(W)pr(WD)`$ $`=`$ $`pr\left(\left(h_1\tau _2+h_2\right)\widehat{\mathrm{\Psi }}\right).`$
Let $`𝒱_{\pm i}`$ (resp. $`𝒱_0`$), $`1ir`$, be the closed subspace of $`C[[\zeta ,\eta ,\upsilon ]]`$ generated by monomials $`\zeta ^{Q_1}\eta ^{Q_2}\upsilon ^{Q_3}`$ for which $`\mu ^{Q_1Q_2}=\mu _i^{\pm 1}`$ (resp. $`\mu ^{Q_1Q_2}=1`$). Let $`P_i`$ be the associated projection. Applying these various projections to the previous equations leads to the following compatibility equations :
$`P_i\left(pr\left(\left(f_{1,i}\mu _if_{2,i}\right)\tau _2\widehat{\mathrm{\Psi }}\right)\right)`$ $`=`$ $`0`$
$`P_i\left(pr\left(\left(g_{1,i}\mu _i^1g_{2,i}\right)\tau _2\widehat{\mathrm{\Psi }}\right)\right)`$ $`=`$ $`0`$ (22)
$`P_0\left(pr\left(\left(h_1\tau _2+h_2\right)\widehat{\mathrm{\Psi }}\right)\right)`$ $`=`$ $`0.`$
According to the properties of the eigenvalues of $`D\mathrm{\Phi }(0)`$ and its invariants, we have the following :
* If $`Q:=(Q_1,Q_2,Q_3)N^{p+p+q}`$ is such that $`\mu ^{Q_1Q_2}=1`$ then $`\zeta ^{Q_1}\eta ^{Q_2}\upsilon ^{Q_3}T_j`$ is a monomial $`\zeta ^{Q_1^{}}\eta ^{Q_2^{}}\upsilon ^{Q_3^{}}`$ such that $`\mu ^{Q_1^{}Q_2^{}}=1`$.
* If $`Q:=(Q_1,Q_2,Q_3)N^{2p+q}`$ and $`1ip`$ is such that $`\mu ^{Q_1Q_2}=\mu _i^{\pm 1}`$ then $`\zeta ^{Q_1}\eta ^{Q_2}\upsilon ^{Q_3}T_j`$ is a monomial $`\zeta ^{Q_1^{}}\eta ^{Q_2^{}}\upsilon ^{Q_3^{}}`$ such that $`\mu _k^{Q_1^{}Q_2^{}}=\mu _i^1`$.
Since $``$ is compatible with the $`T_j`$’s, we have $`pr(fT_j)=pr(f)T_j`$ for any formal power series $`f`$. Moreover, we have
$$P_{\pm i}(fT_j)=P_i(f)T_j,P_0(fT_j)=P_0(f).$$
We have some freedom in the choice of the normalizing transformation of the $`\tau _j`$’s. Let us show that there is a unique solution $`U_i`$, $`V_i`$ and $`W`$ such that
$$pr_{}U_i=pr_{}V_i=pr_{}W=0,$$
and
$$P_i(pr(V_i))=0;P_0(pr(W))=0.$$
(23)
###### Remark 4.3.
If the formal centralizer $`\widehat{𝒞}_{D\varphi (0)x}`$ is contained in $`\widehat{}`$ then the last condition is always satisfied.
In fact, let apply $`P_i`$ (resp. $`P_i`$, $`P_0`$) to the first (resp. second, last) equation of $`()`$ and let us show that system obtained is solvable. We have
$`P_i(pr(U_i))T_j`$ $`=`$ $`P_i(pr(f_{j,i}\widehat{\mathrm{\Psi }}))`$
$`\lambda _{j,i}^1P_i(pr(U_i))`$ $`=`$ $`P_i(pr(g_{j,i}\widehat{\mathrm{\Psi }}))`$
$`P_0(pr(h_j\widehat{\mathrm{\Psi }}))`$ $`=`$ $`0.`$
The last equation is obtained from equation $`(\text{17})`$. In fact, we have
$$0=h_j\tau _j\widehat{\mathrm{\Psi }}+h_j\widehat{\mathrm{\Psi }}=(h_j\widehat{\mathrm{\Psi }})\stackrel{~}{\tau }_j+(h_j\widehat{\mathrm{\Psi }}).$$
Since $`\stackrel{~}{\tau }_j`$ is linear on $``$, we have
$$pr((h_j\widehat{\mathrm{\Psi }})\stackrel{~}{\tau }_j)=pr((h_j\widehat{\mathrm{\Psi }})T_j)=pr(h_j\widehat{\mathrm{\Psi }})T_j.$$
Therefore, we have $`P_0(pr(h_j\widehat{\mathrm{\Psi }}))=0`$. The compatibility condition of the the first equation is
$$P_i(pr(f_{j,i}\widehat{\mathrm{\Psi }}))T_j+\lambda _{j,i}P_i(pr(g_{j,i}\widehat{\mathrm{\Psi }}))=0.$$
Let us show that it can be obtained from equation $`(\text{15})`$. In fact, let us first compose on the right by $`\widehat{\mathrm{\Psi }}`$ and then project onto $`CI`$. We have
$$\lambda _{j,i}pr(g_{j,i}\widehat{\mathrm{\Psi }})+pr(f_{i,j}\widehat{\mathrm{\Psi }}\stackrel{~}{\tau }_j)=0.$$
As above, this leads to
$$\lambda _{j,i}pr(g_{j,i}\widehat{\mathrm{\Psi }})+pr(f_{i,j}\widehat{\mathrm{\Psi }}T_j)=0,$$
and then to
$$\lambda _{j,i}P_i\left(pr(g_{j,i}\widehat{\mathrm{\Psi }})\right)+P_i\left(pr(f_{i,j}\widehat{\mathrm{\Psi }})\right)T_j=0.$$
Let us consider the following equations
$$u_i=\lambda _{j,i}P_i(pr(g_{j,i}(\zeta +U^{}+u,\eta +V^{},\upsilon +W^{}))),i=1,\mathrm{},r+2s,$$
(24)
where the $`u_i`$’s are the unknowns and which components belong to the range of $`P_ipr`$. Here, the maps $`U^{}`$, $`V^{}`$ and $`W^{}`$ are holomorphic maps defined by the $`\mathrm{\Psi }^{}`$, the normalized transformation which linearizes $`\mathrm{\Phi }`$ on $``$. Since the $`g_{j,i}`$’s are non-linear holomorphic functions, we can apply the implicit function theorem to obtain holomorphic $`u_i`$’s in a neighbourhood of the origin satisfying equation $`(\text{24})`$. The map $`u`$ is alspo solution of the equations $`u_iT_1=P_i(pr(f_{1,i}(\mathrm{\Psi }^{}+u)))`$. Then, the germ of holomorphic diffeomorphism at the origin of $`C^{r+n1}`$ defined to be
$$\mathrm{\Psi }(\zeta ,\eta ,\upsilon )=\{\begin{array}{cc}\zeta _i^{}=\zeta _i+U_i^{}(\zeta ,\eta ,\upsilon )+u_i(\zeta ,\eta ,\upsilon ),i=1,\mathrm{},r+2s\hfill & \\ \eta _i^{}=\eta _i+V_i^{}(\zeta ,\eta ,\upsilon ),i=1,\mathrm{},r+2s\hfill & \\ \upsilon ^{}=\upsilon +W^{}(\zeta ,\eta ,\upsilon )\hfill & \end{array}$$
linearizes simultaneously and holomorphically the $`\tau _j`$’s on $``$. It is the unique germ of diffeomorphism which linearizes the $`\tau _j`$’s and satisfies to $`(\text{23})`$ and $`(\text{24})`$. It remains to show that$`\mathrm{\Psi }`$ commutes whith the anti-holomorphic involution $`\rho `$. In fact, we have $`\widehat{\mathrm{\Psi }}\stackrel{~}{\tau }_2=\tau _2\widehat{\mathrm{\Psi }}`$. By composition on the right and on the left by $`\rho `$, we obtain
$$\left(\mathrm{\Psi }^{}+(\rho u\rho )\right)(\rho \stackrel{~}{\tau }_2\rho )=\tau _1\left(\mathrm{\Psi }^{}+(\rho u\rho )\right).$$
This is due to the fact that $`\rho `$ is an involution which commutes with $`\mathrm{\Psi }^{}`$. It is to be noticed that the map $`\rho \stackrel{~}{\tau }_2\rho `$ is equal to $`T_1`$ modulo the ideal $``$. Therfore, if we project outside $``$, we obtain
$$pr(\mathrm{\Psi }^{})T_1+pr(\rho u\rho )T_1=pr\left(\tau _1\left(\mathrm{\Psi }^{}+(\rho u\rho )\right)\right).$$
According to the properties of $`\mathrm{\Psi }^{}`$, we have obtain
$$P_i(pr(\rho u\rho )_i)T_1=P_i(pr\left(f_{1,i}\left(\mathrm{\Psi }^{}+(\rho u\rho )\right)\right)).$$
Therefore, $`\rho u\rho `$ is solution of the same equation as $`u`$. By uniqueness, we obtain $`\rho u\rho =u`$. Hence, $`\mathrm{\Psi }`$ commutes with $`\rho `$. This end the proof of the theorem. ∎
Let us express a direct geometric implication.
###### Corollary 4.1.
Under the assumptions of the theorem and in the new holomorphic coordinates system, the complexified submanifold $``$ intersects each irreducible component of the zero locus $`V()`$ which is invariant under $`\rho `$, $`T_1`$ and $`T_2`$ along a complex submanifold which is a complexified quadric of $`C^k`$, $`kn`$ (if the $`T_j`$’s are still involutions) .
###### Proof.
Under the assumptions of the theorem, the $`\tau _j`$’s are holomorphically linearizable on $``$. Hence, in the new coordinates, their restriction to any irreducible component of $`V()`$ is equal to their restriction of their linear part at the origin. Assume that $`V_i`$ is an irreductible component of $`V()`$ which is invariant under $`\rho `$, $`T_1`$ and $`T_2`$. Since $``$ is a monomial ideal, its zero locus is an intersection of unions of complex hyperplanes $`\{\zeta _{i_j}=0\}`$, $`\{\eta _{i_j}=0\}`$ and $`\{\upsilon _{k_l}=0\}`$. Hence it is a linear complex manifold. Then, the submanifold $`V_iFix(\rho )`$ of $`M`$ is a quadric associated to the linear involutions $`T_{j|V_j}`$. ∎
### 4.5 Holomorphic equivalence to quadrics
As a direct consequence, we have the following result :
###### Theorem 4.2.
Assume that $`M`$ is formally equivalent to the quadric defined by its quadratic part. Let $`\mathrm{\Phi }`$ be the associated germ of holomorphic diffeomorphism. Assume that its linear part $`D\mathrm{\Phi }(0)`$ is semi-simple and diophantine, then $`M`$ is biholomophic to the quadric
$$(Q)\{\begin{array}{cc}z_1=x_1+iy_1\hfill & \\ \mathrm{}\hfill & \\ z_p=x_p+iy_p\hfill & \\ y_{p+1}=0\hfill & \\ \mathrm{}\hfill & \\ y_{n1}=0\hfill & \\ z_n=Q(z^{},\overline{z}^{})\hfill & \end{array}$$
(25)
When $`p=1`$, $`n=2`$, this result is due to X. Gong \[Gon94\] in the hyperbolic case.
###### Proof.
Let us apply the main theorem with the zero ideal $`:=(0)`$. In fact, it is compatible with the $`D\tau _j(0)`$’s and $`\rho `$. That $`M`$ is formally equivalent to the quadric $`(Q)`$ means precisely that the $`\tau _j`$’s are simultaneously formally linearizable. According to the assumption, the $`\tau _j`$’s are simultaneously and holomorphically linearizable. Since the linearizing diffeomorphism $`\mathrm{\Psi }`$ commutes with the linear anti-holomorphic involution $`\rho `$, $`\mathrm{\Psi }`$ is the complexified of an holomorphic diffeomorphism $`\psi `$ of $`(C^n,0)`$ which maps $`M`$ to the quadric $`(Q)`$. ∎
### 4.6 Cutting varieties
We want to apply the theorem to the ideal $`:=ResIdeal`$, the ideal generated by the set $``$ of resonnant monomials $`(\text{14})`$. Almost all pair of involution $`(\tau _1,\tau _2)`$ are formally linearizable along the resonant ideal $``$. In fact, the normal form theory for the diffeomorphism $`\mathrm{\Phi }`$ tells us that all but a finite number of the monomials appearing in the Taylor expansion of a normal form belongs to the resonannt ideal.
###### Lemma 4.3.
Assume that the non-linear centralizer $`𝒞_{D\mathrm{\Phi }(0)}`$ is included in the resonnant ideal $`ResIdeal`$ and that $`D\mathrm{\Phi }(0)`$ has distinct eigenvalues. Then the $`\tau _j`$’s are formally and simultaneously linearizable on $`=ResIdeal`$. Moreover, $``$ is compatible with the $`D\tau _j(0)`$’s and $`\rho `$.
###### Proof.
Let us write
$$\mathrm{\Phi }=\{\begin{array}{cc}\mu _i\zeta _i+\varphi _i,i=1,\mathrm{},r+2s\hfill & \\ \mu _i^1\eta _i+\psi _i,i=1,\mathrm{},r+2s\hfill & \\ \upsilon +\theta .\hfill & \end{array}$$
Since $`\mathrm{\Phi }=\tau _1\tau _2`$, we have
$$\{\begin{array}{cc}\lambda _{1,i}g_{2,i}+f_{1,i}\tau _2=\varphi _i\hfill & \\ \lambda _{1,i}^1f_{2,i}+g_{1,i}\tau _2=\psi _i\hfill & \\ h_2+h_1\tau _2=\theta .\hfill & \end{array}$$
But according to relations $`(\text{15})`$, $`(\text{16})`$ and $`(\text{17})`$, we have
$$\{\begin{array}{cc}g_{2,i}=\lambda _{2,i}^1f_{2,i}\tau _2\hfill & \\ g_{1,i}=\lambda _{1,i}^1f_{1,i}\tau _2\hfill & \\ h_2+h_2\tau _2=0.\hfill & \end{array}$$
(26)
Therefore, we have
$$\{\begin{array}{cc}(\lambda _{1,i}\lambda _{2,i}^1)f_{2,i}\tau _2+f_{1,i}\tau _2=\varphi _i\hfill & \\ \lambda _{1,i}^1f_{2,i}\lambda _{1,i}^1f_{1,i}\mathrm{\Phi }=\psi _i\hfill & \\ h_2\tau _1+h_1\tau _2=\theta .\hfill & \end{array}$$
We find that
$$\{\begin{array}{cc}\mu _i^1f_{1,i}f_{1,i}\mathrm{\Phi }=\mu _i\psi _i+\mu _i^1\varphi _i\tau _2\hfill & \\ f_{2,i}=f_{1,i}\mathrm{\Phi }+\mu _i\psi _i\hfill & \end{array}$$
(27)
Assume that the $`\varphi _i`$’s, $`\psi _i`$’s and $`\theta `$ belong to $``$. Assume that the $`k`$-jet of $`\tau _1T_1`$ and $`\tau _2T_2`$ belong to also the ideal $``$. Then, by equalities $`(\text{26})`$ and $`(\text{27})`$, we have that $`\mu _i^1f_{1,i,k+1}f_{1,i,k+1}D\mathrm{\Phi }(0)`$ belong to $``$. Here, $`f_{1,i,k+1}`$ denotes the homogenous polynomial of order $`k+1`$ in the Taylor expansion of $`f_{1,i}`$ at the origin. Since the centralizer of $`D\mathrm{\Phi }(0)`$ is contained in $``$ then $`f_{1,i,k+1}`$ belongs to $``$. It follows that $`f_{2,i,k+1}`$ also belongs to $``$. So, are the $`g_{j,i,k+1}`$’s and the $`h_{j,k+1}`$’s. Therefore, under the assuptions, the $`\tau _i`$’s are formally linearizable on the ideal $``$ since $`\mathrm{\Phi }`$ is.
About the second point, the $`D\tau _j(0)`$’s leave each monomial $`\zeta _i\eta _i`$ invariant. Moreover, if a monomial $`\zeta ^Q\eta ^R`$ is first integral of $`D\mathrm{\Phi }(0)`$, then so is $`\eta ^Q\zeta ^R`$. this is due to the fact that if $`\mu ^Q(\mu ^1)^R=1`$ then $`(\mu ^1)^Q(\mu )^R=1`$. As a consequence, each first integral of $`D\mathrm{\Phi }(0)`$ is sent to a first integral of $`D\mathrm{\Phi }(0)`$ by the $`D\tau _j(0)`$’s. Furthermore, $`\overline{\rho }^{}`$ maps a monomial first integral to another. This is due to the fact that the eigenvalues are distinct. ∎
The assumption of the lemma means that the resonnances are generated by the generators of the ring of invariants of the linear part of $`\mathrm{\Phi }`$ at the origin.
###### Theorem 4.3.
Assume that the non-linear centralizer $`𝒞_{D\mathrm{\Phi }(0)x}`$ is included in the resonnant ideal $`ResIdeal`$. Moreover let us assume that $`D\mathrm{\Phi }(0)`$ has distinct hyperbolic eigenvalues and is diophantine (on $`ResIdeal`$ if the latter is properly embedded). Then, there are good holomorphic coordinates $`(v_1,\mathrm{},v_n)`$ of $`C^n`$, in which the submanifold $`M`$ intersects the complex linear manifold $`\{v_{p+1}=\mathrm{}=v_n=0\}`$ along a real analytic subset $`V`$ whose complexified is nothing but the zero locus of the resonnant ideal $`ResIdeal`$ : $`V=V(ResIdeal)Fix(\rho )`$, where $`Fix(\rho )`$ denotes the fixed point set of $`\rho `$.
###### Proof.
First of all, we make a change of coordinates as in lemma 4.2. Then, the quadratic functions which are invariant under both $`T_1`$ and $`T_2`$ are linear combinations of monomials of the form $`\zeta _i\eta _i`$ and $`\upsilon _j\upsilon _k`$. Therefore, in these coordinates, the functions defining the complexified quadric are functions $`f,g`$ of the $`\zeta _i\eta _i`$’s and $`\upsilon _j\upsilon _k`$’s. Since these monomials vanish on the zero locus of $`ResIdeal`$, so does $`f`$ and $`g`$. We apply theorem 4.1 : $`V(ResIdeal)`$ is invariant under the $`\tau _j`$’s and their restriction to it are equal to $`T_{j|V(ResIdeal)}`$. Because of hyperbolicity, each irreductible component of $`V(ResIdeal)`$ is invariant under $`\rho `$. Hence, the functions which are invariants by the $`\tau _j`$’s vanish on zero on the zero locus of $`Resideal`$. ∎
When $`p=1`$, $`n2`$ and in the hyperbolic case, this result is due to Wilhelm Klingenberg Jr. \[Kli85\]. He found, for instance in the case $`n=2`$, that, in some suitable holomorphic coordinates, the complex analytic set $`\{v_2=0\}`$ intersects $`M`$ along $`\{(\zeta _1,\eta _1)R^2,\zeta _1\eta _1=0\}`$. In fact, in this case, the ring of formal invariants of $`\mathrm{\Phi }`$ is generated by only one element: $`\zeta _1\eta _1`$.
###### Remark 4.4.
When the spectrum of $`D\mathrm{\Phi }(0)`$ contains elliptic eigenvalues and complex eigenvalues, then one has to use the set
$$\stackrel{~}{V}=V()\left(\underset{i}{}\{\zeta _i=\eta _i=0\}\right))(\underset{j𝒞}{}\{\zeta _{r+2j1}=\eta _{r+2j1}=\zeta _{r+2j}=\eta _{r+2j}=0\})$$
which is invariant by $`\rho `$. Then, the conclusion of the previous theorem holds with $`\stackrel{~}{V}Fix(\rho )`$.
|
warning/0506/hep-th0506062.html
|
ar5iv
|
text
|
# Impact of supersymmetry on the nonperturbative dynamics of fuzzy spheres
## 1 Introduction
Matrix models are considered as one of the most promising candidates for a nonperturbative formulation of string theories. Indeed some concrete models are proposed as constructive definitions of superstring and M theories . These models are obtained from the dimensional reduction of super Yang-Mills theory in ten dimensions. In a broad sense, such models belong to the class of the so-called dimensionally reduced models (or large-$`N`$ reduced models), which were studied intensively in the eighties as an equivalent description of large-$`N`$ gauge theories. Unlike the old models, however, the new models are written in terms of Hermitian matrices, and they have manifest supersymmetry, which is expected to have crucial effects on their dynamics <sup>1</sup><sup>1</sup>1See ref. for a comprehensive review on these issues..
An important feature of these matrix models is that the space-time is not introduced from the outset, but it emerges dynamically as the eigenvalue distribution of the bosonic matrices. In fact there are certain evidences in the IIB matrix model that four-dimensional space-time is generated dynamically . In refs. the free energy of space-time with various dimensionality has been calculated using the gaussian expansion method, and the free energy turned out to take the minimum value for the four-dimensional space-time. In ref. it was found that the fuzzy $`\mathrm{S}^2\times \mathrm{S}^2`$ (but not the fuzzy S<sup>2</sup>) is a solution to the 2-loop effective action. See also refs. for related works on this issue.
By adding a Chern-Simons term to the matrix models, one obtains fuzzy spheres as classical solutions , and their dynamical properties have been studied in refs. . This provides a matrix description of the so-called Myers effect in string theory . The emergence of a fuzzy sphere in matrix models may be regarded as a prototype of the dynamical generation of space-time since it has lower dimensionality than the original dimensionality that the model can actually describe. When $`k`$ fuzzy spheres coincide, the gauge symmetry enhances from U(1)<sup>k</sup> to U($`k`$). By expanding the theory around such a solution, one obtains a U($`k`$) gauge theory on a noncommutative geometry . Therefore the model may also serve as a toy model for the dynamical generation of gauge group, which is expected to occur in the IIB matrix model .
In fact one can use the above matrix models to define a regularized field theory on the fuzzy sphere as has been done on a noncommutative torus , which enables nonperturbative studies of such theories from first principles . This is motivated from the general expectation that noncommutative geometry provides a crucial link to string theory and quantum gravity . Yet another motivation is to use the fuzzy sphere (or its generalization ) as a regularization scheme alternative to the lattice regularization . Unlike the lattice, fuzzy spheres preserve the continuous symmetries of the space-time considered, and hence it is expected to ameliorate the well-known problem concerning chiral symmetry and supersymmetry. A challenge in this direction is to remove the effects of noncommutativity of the space-time in the “continuum limit”. The fuzzy sphere is also useful in the Coset Space Dimensional Reduction , where one can take the compact part of space-time to be a fuzzy coset .
Whatever the motivation is, the stability of fuzzy-sphere-like solutions is clearly one of the most important issues. In cases when there are more than one stable solutions, one can identify the true vacuum by comparing the corresponding free energy. This will be important in the dynamical determination of the space-time dimensionality and the gauge group in superstring theory. In the series of papers , we addressed such issues in various kinds of models using both perturbative calculations and Monte Carlo simulations. In ref. we have studied the dimensionally reduced 3d Yang-Mills models with the Chern-Simons term, which has the fuzzy 2-sphere (S<sup>2</sup>) as a classical solution . We have found a first-order phase transition as we vary the coefficient ($`\alpha `$) of the Chern-Simons term. For small $`\alpha `$ the large-$`N`$ behavior of the model is the same as in the pure Yang-Mills model, whereas for large $`\alpha `$ a single fuzzy S<sup>2</sup> appears dynamically. In addition we find that the $`k`$ coincident fuzzy spheres, which are also classical solutions of the same model, cannot be realized as the true vacuum in this model even in the large-$`N`$ limit. This implies that the dynamical gauge group is U(1) in this model. In refs. we have extended this work to various matrix models, which incorporate four-dimensional fuzzy manifolds as classical solutions. While the fuzzy S<sup>4</sup> turned out to be unstable , we find that the fuzzy CP<sup>2</sup> and the fuzzy $`\text{S}^2\times \text{S}^2`$ are stable at large $`N`$ although the true vacuum is actually given by the fuzzy S<sup>2</sup>. In the latter two cases the gauge group generated dynamically turned out to be U(1) as well. In ref. , on the other hand, it has been shown for the first time that gauge groups of higher rank can be realized in the true vacuum by adding a mass term to the 3d Yang-Mills-Chern-Simons model.
The aim of the present paper is to study the impact of supersymmetry on the dynamics of the fuzzy spheres. The simplest 3d supersymmetric model is problematic nonperturbatively since the partition function is divergent . This leads us to study the 4d supersymmetric model with a cubic term instead. Indeed it turns out that the supersymmetry has a striking impact. Unlike the bosonic models, the fuzzy sphere is always stable if the large-$`N`$ limit is taken in such a way that various correlation functions scale. We also observe an interesting phenomenon that the power-law tail of the eigenvalue distribution, which exists in the supersymmetric models without the Chern-Simons term , disappears in the presence of the fuzzy sphere in the large-$`N`$ limit. Coincident fuzzy spheres turn out to be unstable, which implies that the dynamically generated gauge group is U(1) in the present model.
This paper is organized as follows. In section 2 we define the model and discuss its fuzzy sphere solutions. In section 3 we study the phase diagram of the model. In section 4 we study the geometrical structure of the dominant configurations. In section 5 we study coincident fuzzy spheres and discuss the dynamical gauge group. Section 6 is devoted to summary and discussions. In appendix A we explain the algorithm for our Monte Carlo simulations. In appendices B and C we provide the details of perturbative calculations.
## 2 The model and the fuzzy sphere
The model we study in this paper is defined by the action
$`S`$ $`=`$ $`S_\mathrm{b}+S_\mathrm{f},`$ (1)
$`S_\mathrm{b}`$ $`=`$ $`N\text{tr}\left({\displaystyle \frac{1}{4}}[A_\mu ,A_\nu ]^2+{\displaystyle \frac{2}{3}}i\alpha {\displaystyle \underset{i,j,k=1}{\overset{3}{}}}ϵ_{ijk}A_iA_jA_k\right),`$ (2)
$`S_\mathrm{f}`$ $`=`$ $`N\text{tr}\left(\overline{\psi }_\alpha (\mathrm{\Gamma }_\mu )_{\alpha \beta }[A_\mu ,\psi _\beta ]\right),`$ (3)
where $`A_\mu `$ ($`\mu =1,2,3,4`$) are $`N\times N`$ traceless hermitian (bosonic) matrices, and $`\psi _\alpha ,\overline{\psi }_\alpha `$ ($`\alpha =1,2`$) are $`N\times N`$ traceless complex (fermionic) matrices. Here and henceforth we assume that repeated Greek indices are summed over all possible integers. The $`ϵ_{ijk}`$ ($`i,j,k=1,2,3`$) is a totally anti-symmetric tensor with $`ϵ_{123}=1`$. The $`2\times 2`$ matrices $`\mathrm{\Gamma }_\mu `$ are Weyl-projected gamma matrices in four dimensions, and they are given explicitly as
$`\mathrm{\Gamma }_1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\mathrm{\Gamma }_2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\mathrm{\Gamma }_3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\mathrm{\Gamma }_4=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right).`$ (12)
The convergence of the integration over $`A_\mu `$ is a non-trivial issue since the integration region is non-compact. At $`\alpha =0`$ the partition function is finite for arbitrary $`N`$, as first conjectured by ref. and proved later by ref. , and this remains to be the case also for $`\alpha 0`$ . Moreover, since the fermion determinant of this model is positive semi-definite (See ref. for a proof), the model can be studied by Monte Carlo simulations without confronting the so-called sign problem. The pure super Yang-Mills model ($`\alpha =0`$), which may be regarded as the 4d version of the IIB matrix model , has been studied intensively . The sign problem does not occur even if one includes the cubic term, which is real <sup>2</sup><sup>2</sup>2This is in contrast to the Chern-Simon term in ordinary gauge theories in euclidean space-time, which is purely imaginary. Note that the coefficient $`\alpha `$ in (2) should be chosen to be real in order for fuzzy sphere solutions to exist..
For $`\alpha =0`$ the model has manifest SO(4) symmetry and supersymmetry. The cubic term in (2) obviously breaks the SO(4) symmetry down to SO(3). It also breaks supersymmetry, but the effects of breaking is “soft” since the power of $`A_\mu `$ is lower than the quartic term . Therefore one may still anticipate to see peculiar effects of supersymmetry. We repeat that the 3d supersymmetric model, which has been studied perturbatively , is actually problematic nonperturbatively since the partition function is divergent . Therefore, the present 4d model is the simplest model that can be studied in order to examine the impact of supersymmetry on the fuzzy sphere dynamics.
Let us then consider the classical solutions of this model. For $`\psi =0`$ the equation of motion reads
$`[A_\nu ,[A_\nu ,A_i]]+i\alpha {\displaystyle \underset{j,k=1}{\overset{3}{}}}ϵ_{ijk}[A_j,A_k]=0\text{ for }i=1,2,3,`$ (13)
$`[A_\nu ,[A_\nu ,A_4]]=0.`$
Apart from the solution given by commuting matrices, which exists also for $`\alpha =0`$, we have the fuzzy $`\mathrm{S}^2`$ solution given by
$`\{\begin{array}{cc}A_i^{(\mathrm{S}^2)}=\alpha L_i^{(N)}\hfill & \text{ for }i=1,2,3,\hfill \\ A_4^{(\mathrm{S}^2)}=0,\hfill & \end{array}`$ (16)
where $`L_i^{(r)}`$ ($`i=1,2,3`$) represents the $`r`$-dimensional irreducible representation of the SU$`(2)`$ Lie algebra
$$[L_i^{(r)},L_j^{(r)}]=i\alpha ϵ_{ijk}L_k^{(r)}.$$
(17)
The solution $`A_\mu ^{(\mathrm{S}^2)}`$ satisfies
$`{\displaystyle \underset{i=1}{\overset{3}{}}}\left(A_i^{(\mathrm{S}^2)}\right)^2={\displaystyle \frac{1}{4}}(N^21)\alpha ^2\mathbf{\hspace{0.17em}1}_N,`$ (18)
which implies that the “radius” of the fuzzy sphere is given by
$$\rho =\frac{1}{2}\alpha \sqrt{N^21}.$$
(19)
We consider more general solutions in section 5.
## 3 Phase diagram
### 3.1 Monte Carlo results
In this section we calculate various quantities by Monte Carlo simulation and study the phase diagram of the model (1). We show results obtained by using the fuzzy sphere $`A_\mu ^{(\mathrm{S}^2)}`$ as the initial configuration, but we have checked that the result is the same for other initial configurations such as $`A_\mu =0`$ or some randomly generated configurations. For brevity we introduce the notation
$`F_{\mu \nu }`$ $`=`$ $`i[A_\mu ,A_\nu ],`$ (20)
$`M`$ $`=`$ $`{\displaystyle \frac{2}{3N}}i{\displaystyle \underset{i,j,k=1}{\overset{3}{}}}ϵ_{ijk}\text{tr}(A_iA_jA_k).`$ (21)
We note that there is an exact result
$`{\displaystyle \frac{1}{N}}\text{tr}(F_{\mu \nu })^2+3\alpha M=6\left(1{\displaystyle \frac{1}{N^2}}\right),`$ (22)
which can be derived as in the bosonic case . This result has been used to check our code for the simulation.
By performing one-loop calculation around the fuzzy sphere $`A_\mu ^{(\mathrm{S}^2)}`$, we obtain the leading large-$`N`$ behaviors as (See appendix C for the details)
$`{\displaystyle \frac{1}{N}}\text{tr}(F_{\mu \nu })^2`$ $``$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\alpha }^4+6,`$ (23)
$`{\displaystyle \frac{1}{\sqrt{N}}}M`$ $``$ $`{\displaystyle \frac{1}{6}}\stackrel{~}{\alpha }^3+0,`$ (24)
$`{\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{N}}\text{tr}(A_\mu )^2`$ $``$ $`{\displaystyle \frac{1}{4}}\stackrel{~}{\alpha }^2+0,`$ (25)
where we have introduced the rescaled parameter
$`\stackrel{~}{\alpha }=\alpha \sqrt{N}.`$ (26)
In the r.h.s. of eqs. (23) $``$ (25) the first term represents the classical result, and the second term represents the one-loop correction.
In figure 1 we plot the results for $`\frac{1}{N}\text{tr}(F_{\mu \nu })^2`$ and $`\frac{1}{\sqrt{N}}M`$ obtained by Monte Carlo simulations. We find that Monte Carlo data agree with the one-loop results even at $`\stackrel{~}{\alpha }=0`$. This is rather surprising since the expansion parameter in the perturbative calculation (at finite $`N`$) is $`\frac{1}{\alpha ^4}`$. As we will see shortly, however, the system actually changes its behavior at $`\stackrel{~}{\alpha }\frac{1}{\sqrt{N}}`$, and the agreement in figure 1 below that point should rather be considered as accidental.
Let us then consider the quantity $`\frac{1}{N}\text{tr}(A_\mu )^2`$, which we have postponed since it involves a subtle issue. At $`\alpha =0`$ this quantity is actually divergent even for finite $`N`$, as first observed in numerical studies and further confirmed by ref. . (Ref. provides some analytical explanation.) On the other hand, from perturbative calculations around the fuzzy sphere, we obtain a finite result for finite $`N`$ (See eqs. (100) and (101)). In order to clarify the situation, let us first look at the history of $`\frac{1}{N}\text{tr}(A_\mu )^2`$, which is plotted in figure 2 for various $`\alpha `$ at $`N=4`$. The horizontal axis represents the number of “trajectories” in the hybrid Monte Carlo algorithm (See appendix A.) with the parameters $`\nu =100`$ and $`\mathrm{\Delta }\tau =0.01`$. For $`\alpha =0`$ the history has a lot of spikes, and these spikes are responsible for the divergence of $`\frac{1}{N}\text{tr}(A_\mu )^2`$. As we increase $`\alpha `$ the spikes become less frequent, and their height gets lowered. At $`\alpha 1.1`$ the history looks quite regular. We also looked at the history at larger $`N`$, and find that it becomes regular (no spikes) for $`\alpha 0.5`$ at $`N=8`$ and for $`\alpha 0.3`$ at $`N=16`$. The transition point $`\alpha _{\mathrm{tr}}`$ is roughly consistent with $`\alpha _{\mathrm{tr}}\frac{1}{N}`$ (i.e., $`\stackrel{~}{\alpha }_{\mathrm{tr}}\frac{1}{\sqrt{N}}`$).
According to ref. , switching on $`\alpha `$ does not change the convergence properties of the matrix integrals. Therefore, unless there is some special mechanism for canceling the leading divergence, we get $`\frac{1}{N}\text{tr}(A_\mu )^2=\mathrm{}`$ even for $`\alpha 0`$. The reason why we get finite results by perturbative calculations should then be that such a calculation only include the region of the configuration space near the fuzzy sphere solution so that it does not take into account the configurations that have large $`\frac{1}{N}\text{tr}(A_\mu )^2`$, which may actually contribute crucially to the vacuum expectation value of that quantity.
One can, however, define a finite quantity
$`R=\sqrt{{\displaystyle \frac{1}{N}}\text{tr}(A_\mu )^2},`$ (27)
which is finite at $`\alpha =0`$ and behaves as O(1) at large $`N`$ <sup>3</sup><sup>3</sup>3Strictly speaking, ref. studies a slightly different observable, but the large-$`N`$ behavior is expected to be qualitatively the same. in the present parametrization of the action. Due to the argument of ref. , this quantity is finite also for $`\alpha 0`$. We plot the Monte Carlo results for $`R/\sqrt{N}`$ in figure 3 on the left as a function of $`\stackrel{~}{\alpha }`$ for $`N=4`$, $`8`$ and $`16`$. At large $`\stackrel{~}{\alpha }`$ the data agree very well with the classical result $`R=\frac{1}{2}\stackrel{~}{\alpha }\sqrt{N}`$ for the fuzzy sphere solution.
In figure 3 on the right we plot $`R`$ against $`\alpha `$. We expect that $`R`$ approaches a finite value in the large-$`N`$ limit for each $`\alpha `$ in the small-$`\alpha `$ regime, and our data are roughly consistent with this picture. If we assume the transition to take place at the point where the fuzzy sphere result $`R\frac{1}{2}\stackrel{~}{\alpha }\sqrt{N}`$ becomes comparable to the pure super Yang-Mills behavior $`R\mathrm{O}(1)`$, the transition point should be $`\stackrel{~}{\alpha }\frac{1}{\sqrt{N}}`$, which roughly agrees with the point $`\stackrel{~}{\alpha }_{\mathrm{tr}}`$, where the spikes in figure 2 get suppressed.
Thus we conclude that the system actually undergoes some transition at $`\stackrel{~}{\alpha }\frac{1}{\sqrt{N}}`$, below which the system behaves similarly to the pure super Yang-Mills model ($`\alpha =0`$). The quantities in figure 1 are insensitive to the transition since their behavior in the pure super Yang-Mills phase just happens to be the same as in the fuzzy sphere phase. Note that $`M_{\alpha =0}=0`$ due to parity symmetry $`A_iA_i`$, whereas from the perturbative expansion around the fuzzy sphere, one finds that the one-loop contribution to $`M`$ is absent due to supersymmetry <sup>4</sup><sup>4</sup>4Note that the cubic term, which breaks supersymmetry softly, does not appear in the relevant one-loop calculation. This is not the case, however, at higher loop calculations., and as a consequence the result (103) has a smooth extrapolation to $`\alpha =0`$, which agrees with $`M_{\alpha =0}=0`$. Since $`M`$ and $`\frac{1}{N}\text{tr}F^2`$ are related to each other through the exact result (22), the agreement of the former propagates to that of the latter. Thus the success of one-loop results for these quantities in the small-$`\stackrel{~}{\alpha }`$ regime does not mean that we are still in the fuzzy sphere phase, but it simply means that these quantities are insensitive to the transition from the fuzzy sphere phase to the pure super Yang-Mills phase.
The above results are in striking contrast to those obtained in the bosonic model , where we observed that the fuzzy sphere becomes unstable at some finite $`\stackrel{~}{\alpha }`$ and various quantities show a hysteresis behavior, implying a first order phase transition.
### 3.2 Theoretical understanding based on the effective action
In the previous section we observed that the fuzzy sphere is stable in the large-$`N`$ limit at any finite $`\stackrel{~}{\alpha }`$ unlike the bosonic model. Here we would like to provide some theoretical understanding of the striking difference between the bosonic and supersymmetric cases based on the one-loop effective action. For that purpose let us consider a one-parameter family of configurations given by
$`\{\begin{array}{cc}A_i=\beta L_i^{(N)}\hfill & \text{ for }i=1,2,3,\hfill \\ A_4=0,\hfill & \end{array}`$ (30)
where the fuzzy sphere solution (16) corresponds to $`\beta =\alpha `$. The one-loop effective action around (30) can be calculated along the line described in appendix B, and we get the result at large $`N`$ as
$`{\displaystyle \frac{1}{N^2}}W_{1\mathrm{loop}}^{(\beta )}=\left({\displaystyle \frac{1}{8}}\stackrel{~}{\beta }^4{\displaystyle \frac{1}{6}}\stackrel{~}{\alpha }\stackrel{~}{\beta }^3\right)\mathrm{log}N,`$ (31)
where $`\stackrel{~}{\beta }=\beta \sqrt{N}`$. The one-loop effective action has a minimum at $`\stackrel{~}{\beta }=\stackrel{~}{\alpha }`$ for arbitrary $`\stackrel{~}{\alpha }`$.
In analogous calculations in the bosonic models , the one-loop contribution gives rise to a term proportional to $`\mathrm{log}\stackrel{~}{\beta }`$. Due to this term the (local) minimum disappears below some critical $`\stackrel{~}{\alpha }`$, which indeed agrees well with the Monte Carlo results. In the present supersymmetric case, the $`\stackrel{~}{\beta }`$-dependent one-loop term is absent due to supersymmetry. Thus we can understand the qualitative difference between the bosonic case and the supersymmetric case observed in Monte Carlo simulations.
## 4 Geometrical structure
In this section we study the geometrical structure of the dominant configurations in the supersymmetric model. For that purpose we consider the “Casimir operator”
$`Q=(A_\mu )^2`$ (32)
and define its eigenvalue distribution $`f(x)`$ as
$`f(x)={\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}\delta (x\lambda _j),`$ (33)
where $`\lambda _j`$ $`(j=1,2,\mathrm{},N)`$ represent the eigenvalues of $`Q`$. Let us note that $`\frac{1}{N}\text{tr}(A_\mu )^2`$ discussed in the previous section is related to $`f(x)`$ as
$`{\displaystyle \frac{1}{N}}\text{tr}(A_\mu )^2={\displaystyle \frac{1}{N}}\text{tr}Q={\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}\lambda _j={\displaystyle _0^{\mathrm{}}}xf(x)𝑑x.`$ (34)
In figure 4 we plot the eigenvalue distribution for the same set of $`\alpha `$ and $`N`$ as in figure 2.
### 4.1 Power-law tail
The results for $`\alpha =0`$ reproduce the power-law behavior
$$f(x)x^2,$$
(35)
at large $`x`$, which has been first discovered in ref. and studied also in ref. . This is related to the divergence of $`\frac{1}{N}\text{tr}(A_\mu )^2`$ discussed in section 3.1. In fact this quantity is known to diverge logarithmically , which explains the power in (35) due to (34).
At $`\alpha =0.7`$ the magnitude of the power-law tail becomes much weaker, but our data are consistent with the existence of the power-law tail. At $`\alpha 1.1`$ the power-law tail becomes hardly visible, which corresponds to the disappearance of the spikes in the history of $`\frac{1}{N}\text{tr}(A_\mu )^2`$ seen in Fig. 2. If we assume that $`\frac{1}{N}\text{tr}(A_\mu )^2=\mathrm{}`$ for $`\alpha 0`$, as we argued in the previous section, the power-law tail should be still there, but simply hidden by the main contribution coming from the fuzzy-sphere-like configurations <sup>5</sup><sup>5</sup>5Within perturbative expansion around the fuzzy sphere configuration, we expect to obtain a distribution which decays faster since the $`n`$-th moment $`𝑑xx^nf(x)=\frac{1}{N}\text{tr}\{(A_\mu )^2\}^n`$ can be calculated as a finite quantity to all orders. . Considering that those configurations are enhanced by the Boltzmann weight $`e^{\mathrm{const}.\stackrel{~}{\alpha }^4N^2}`$ at large $`\stackrel{~}{\alpha }`$ compared with the configurations that give the power-law tail, we expect that the magnitude of the power-law tail decreases as $`e^{\mathrm{const}.\stackrel{~}{\alpha }^4N^2}`$. Therefore the power-law tail is expected to disappear completely if we take the large-$`N`$ limit at fixed $`\stackrel{~}{\alpha }`$.
### 4.2 Spherical geometry
In order to clarify the geometrical structure of the fuzzy-sphere-like configurations, we decompose the Casimir operator $`Q`$ as $`Q=Q^{(123)}+Q^{(4)},`$ where
$`Q^{(123)}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}(A_i)^2,`$ (36)
$`Q^{(4)}`$ $`=`$ $`(A_4)^2,`$ (37)
and calculate the eigenvalue distribution for $`Q^{(123)}`$ and $`Q^{(4)}`$, which we denote as $`f^{(123)}(x)`$ and $`f^{(4)}(x)`$, respectively. Figure 5 shows the result for $`\alpha =1.5`$ at $`N=4`$.
The figure on the right shows that the eigenvalues of $`Q^{(4)}`$ is quite small, which can be understood from the one-loop result (101) for $`\frac{1}{N}\text{tr}(A_4)^2_{1\mathrm{loop}}`$, which vanishes as O($`\frac{1}{N^2}\mathrm{log}N`$) in the large-$`N`$ limit with fixed $`\stackrel{~}{\alpha }`$. As a consequence the distribution $`f^{(123)}(x)`$ shown on the left of the figure 5 is almost identical to $`f(x)`$ shown on the bottom right of figure 4. We also find that $`f^{(123)}(x)`$ is peaked around the classical result. Thus we confirm that the dominant configurations indeed have the geometry of a 2-sphere.
At $`\alpha =0`$ the distribution $`f(x)`$ has an empty region around $`x=0`$. Similar behavior has been observed also in the bosonic model , and it can be understood by the uncertainty principle. Therefore the geometrical structure of the dominant configurations at $`\alpha =0`$ should rather be considered as that of a solid ball.
## 5 Dynamical gauge group
In fact the equation of motion (13) has a class of solutions of the form
$`\{\begin{array}{cc}A_i^{(k\mathrm{S}^2)}=\alpha (L_i^{(n)}\mathrm{𝟏}_k)\hfill & \text{ for }i=1,2,3,\hfill \\ A_4^{(k\mathrm{S}^2)}=0,\hfill & \end{array}`$ (40)
where $`N=nk`$. These solutions represent the $`k`$ coincident fuzzy $`\mathrm{S}^2`$, and the Casimir operator $`Q`$ takes the value
$`Q={\displaystyle \frac{1}{4}}(n^21)\alpha ^2\mathbf{\hspace{0.17em}1}_N,`$ (41)
meaning that the “radius” of the fuzzy spheres is given by $`\rho =\frac{1}{2}\alpha \sqrt{n^21}\frac{1}{2k}\alpha N`$, which becomes smaller as $`k`$ increases. By expanding the theory around such a configuration, one obtains noncommutative Yang-Mills theory with the U($`k`$) gauge group .
In order to study the stability of such a configuration, we perform Monte Carlo simulation for $`N=8`$ and $`\alpha =1.0`$ using the $`k=2`$ coincident fuzzy spheres as the initial configuration. In figure 6 we plot the history of the eigenvalues of $`Q`$. The horizontal axis represents the number of “trajectories” in the hybrid Monte Carlo algorithm (See appendix A.) with the parameters $`\nu =100`$ and $`\mathrm{\Delta }\tau =0.0001`$. Thus the $`k=2`$ multi fuzzy sphere configuration is unstable and decays into the single fuzzy sphere.
This phenomenon can be understood by considering the free energy, which is calculated in appendix B up to one-loop. At large $`N`$ the result reads
$`{\displaystyle \frac{1}{N^2}}W_{1\mathrm{loop}}^{(k)}={\displaystyle \frac{1}{24k^2}}\stackrel{~}{\alpha }^4\mathrm{log}N.`$
Since the one-loop term represented by the second term is independent of $`k`$, we find that the free energy takes the smallest value for $`k=1`$.
Although the conclusion concerning the dynamical gauge group is the same as in the bosonic models , we note that the reasoning is different. In the bosonic models the one-loop term in the free energy has the $`\mathrm{log}k^2`$ term, which actually favors large $`k`$. However, if one decreases $`\stackrel{~}{\alpha }`$ so that the one-loop term becomes more important, the fuzzy sphere solutions disappear before the free energy starts to favor $`k>1`$. In the present supersymmetric case, the fuzzy sphere solutions remain to be there, but due to the absence of the $`\mathrm{log}k^2`$ term, the $`k=1`$ solution is always favored.
## 6 Conclusion
In this paper we have studied the dimensionally reduced 4d super Yang-Mills model with an extra Chern-Simons term, which incorporates fuzzy spheres as classical solutions. We have found that supersymmetry indeed has substantial effects on the dynamics of fuzzy spheres.
While the observables that appear in the action change continuously as we vary $`\alpha `$, the model actually possesses two distinctive phases, which is demonstrated by a well-defined observable $`\sqrt{\frac{1}{N}\text{tr}(A_\mu )^2}`$. The tail of the eigenvalue distribution changes drastically as one moves from one phase to the other. In the pure super Yang-Mills phase we observe the same power-law tail as the one known for $`\alpha =0`$, but it disappears in the fuzzy sphere phase in the large-$`N`$ limit. This allows us to identify the critical point quite accurately. From our Monte Carlo data up to $`N=16`$ and some theoretical considerations, we speculate that the transition point $`\stackrel{~}{\alpha }_{\mathrm{tr}}`$ goes to zero as $`\stackrel{~}{\alpha }_{\mathrm{tr}}=\mathrm{O}\left(\frac{1}{\sqrt{N}}\right)`$.
Our results are in sharp contrast to the results obtained for analogous bosonic models , where the fuzzy sphere becomes unstable at some finite critical $`\stackrel{~}{\alpha }`$. A strong first-order phase transition has been observed in the bosonic models, and the (lower) critical point obtained by Monte Carlo simulation can be reproduced very well from the one-loop effective action. In the present paper we have shown that the one-loop term in the effective action vanishes up to an irrelevant constant due to supersymmetry. This explains our observation that the fuzzy sphere is stable down to vanishingly small $`\stackrel{~}{\alpha }`$ at large $`N`$.
We have also studied the dynamical generation of the gauge group. In the bosonic case the quantum instability of the fuzzy spheres was an obstacle in obtaining a non-trivial gauge group in the true vacuum. In the present supersymmetric case, this instability is gone, but we also lost the one-loop term in the free energy for the coincident fuzzy spheres, which favors higher multiplicity. As a result, we obtain the U$`(1)`$ gauge group again. We note, however, that this conclusion applies only to the fuzzy sphere phase, and whether we can obtain a nontrivial gauge group in supersymmetric models without a Chern-Simons term such as the IIB matrix model is still an interesting open question.
###### Acknowledgments.
We thank Hajime Aoki, Subrata Bal, Satoshi Iso, Kazuyuki Kanaya, Yoshihisa Kitazawa and Dan Tomino for valuable discussions. This work was partially funded by the ”Pythagoras” and ”Pythagoras II” project, which is co-funded by the European Social Fund (75$`\%`$) and Greek National Resources (25$`\%`$). The work of T.A. and J.N. was supported in part by Grant-in-Aid for Scientific Research (Nos. 03740 and 14740163, respectively) from the Ministry of Education, Culture, Sports, Science and Technology.
## Appendix A Details of the Monte Carlo simulation
In this section we explain the algorithm used for our simulation. Our algorithm is similar to the one adopted in ref. , but the crucial difference is that we make the Metropolis reject/accept procedure at the end of each trajectory. In the previous algorithm there was a systematic error due to discretization required for solving the Hamiltonian equation, and the step size $`\mathrm{\Delta }\tau `$ for the discretization had to be sent to zero. In the present algorithm we do not need such an extrapolation. Another difference is that we do not use the noisy estimator for estimating the r.h.s. of the Hamilton equation since it causes some systematic error. Instead we invert the Dirac operator directly using the LU decomposition. Each of the two modifications increases the computational effort for making one trajectory (for fixed parameters in the algorithm) from O($`N^5`$) to O($`N^6`$), which is the price we have to pay to make the algorithm “exact”.
The “exact” algorithm is essentially the hybrid Monte Carlo (HMC) algorithm, which is used in studying the large-$`N`$ behavior of the phase quenched version of the IIB matrix model . In that case the one-loop approximation has been used to decrease the computational effort from O($`N^6`$) to O($`N^3`$). The hybrid algorithms are standard in full QCD simulations, and it is useful also in simulating matrix models as demonstrated in refs. . If we had used the Metropolis algorithm in the present model, for instance, the computational effort would have been O($`N^8`$). Note, however, that simulating matrix models is generally harder than simulating field theories due to the non-local nature of the interaction. Even in the bosonic case, the computational effort is at least O($`N^3`$), which grows faster than the number of d.o.f., which is O($`N^2`$).
Let us first recall an explicit form of the fermion determinant derived in ref. . We define a complete basis for the general complex $`N\times N`$ matrices as
$`(t^a)_{ij}=\delta _{ii_a}\delta _{jj_a}(a=1,2,\mathrm{},N^2),`$ (42)
where $`i_a,j_a`$ are integers within the range $`1i_a,j_aN`$ satisfying
$`a=N(i_a1)+j_a.`$ (43)
By taking into account that the fermionic matrices $`\psi _\alpha `$ and $`\overline{\psi }_\alpha `$ are traceless, integration over the fermionic matrices yields the fermion determinant $`det`$, where the $`2(N^21)\times 2(N^21)`$ matrix $``$ is given by
$`_{a\alpha ,b\beta }`$ $`=`$ $`_{a\alpha ,b\beta }^{}_{N^2\alpha ,b\beta }^{}\delta _{i_aj_a}_{a\alpha ,N^2\beta }^{}\delta _{i_bj_b}.`$ (44)
Here the $`2N^2\times 2N^2`$ matrix $`^{}`$ is defined as
$`_{a\alpha ,b\beta }^{}`$ $`=`$ $`(\mathrm{\Gamma }_\mu )_{\alpha \beta }\text{tr}(t^a[A_\mu ,t^b]).`$ (45)
The effective action for $`A_\mu `$ can be written as
$`S_{\mathrm{eff}}[A]=S_\mathrm{b}[A]\mathrm{log}\text{Det}[A].`$ (46)
Following the idea of the hybrid Monte Carlo algorithm, we introduce auxiliary bosonic hermitian matrices $`P_\mu `$ and consider the action
$`S_{\mathrm{HMC}}[P,A]={\displaystyle \frac{1}{2}}\text{tr}(P_\mu )^2+S_{\mathrm{eff}}[A].`$ (47)
Since $`P_\mu `$ does not couple to $`A_\mu `$, we retrieve the original model trivially by integrating out $`P_\mu `$. We regard the action $`S_{\mathrm{HMC}}[P,A]`$ as the Hamiltonian of a classical system described by $`A_\mu (\tau )`$ and its conjugate momentum $`P_\mu (\tau )`$, where $`\tau `$ denotes the fictitious time of the classical system. Then as an update procedure, we may take the old configuration $`(P,A)`$ as the initial configuration $`(P(0),A(0))`$ and solve the Hamiltonian equation
$`{\displaystyle \frac{d(A_\mu )_{ij}}{d\tau }}`$ $`=`$ $`{\displaystyle \frac{S_{\mathrm{HMC}}}{(P_\mu )_{ij}}}=(P_\mu )_{ji},`$ (48)
$`{\displaystyle \frac{d(P_\mu )_{ij}}{d\tau }}`$ $`=`$ $`{\displaystyle \frac{S_{\mathrm{HMC}}}{(A_\mu )_{ij}}}={\displaystyle \frac{S_{\mathrm{eff}}}{(A_\mu )_{ij}}}`$ (49)
$`=`$ $`N\left([A_\nu ,[A_\mu ,A_\nu ]]+2i\alpha ϵ_{\mu \nu \rho }A_\nu A_\rho \right)_{ji}\text{Tr }\left(^1{\displaystyle \frac{}{(A_\mu )_{ij}}}\right)`$
for a finite fictitious time $`T`$ (this defines “one trajectory”) to obtain $`(P(T),A(T))`$. The symbol Tr in (49) denotes a trace over the $`2(N^21)`$-dimensional index, and the derivative $`\frac{}{(A_\mu )_{ij}}`$ is given explicitly by
$`{\displaystyle \frac{_{a\alpha ,b\beta }}{(A_\mu )_{ij}}}=(\mathrm{\Gamma }_\mu )_{\alpha \beta }\left([t^b,t^a]\right)_{ji}.`$ (50)
Since the trace of $`A_\mu `$ is not conserved during the evolution, we subtract the trace part $`A_\mu ^{}=A_\mu (T)\left\{\frac{1}{N}\text{tr}A_\mu (T)\right\}\mathrm{𝟏}`$. Thus we obtain the updated configuration $`(P^{},A^{})`$, where $`P_\mu ^{}=P_\mu (T)`$. Due to the Hamiltonian conservation, this update procedure preserves the action $`S_{\mathrm{HMC}}[P,A]`$. Using also the fact that the transition between $`(P,A)`$ and $`(P^{},A^{})`$ is reversible, one can readily verify the detailed balance. After each trajectory, we update the momentum $`P_\mu `$ fixing $`A_\mu `$, which can be done by simply generating gaussian variables since the $`P_\mu `$-dependent part of the action (47) is gaussian. This procedure is necessary to avoid the ergodicity problem.
In actual calculations we have to discretize the Hamiltonian equation (49). The reversibility of the time evolution can be preserved by using the so-called leap-frog discretization, but the Hamiltonian conservation is inevitably violated. However, we may accept the configuration $`(P^{},A^{})`$ as the updated configuration with the probability $`\mathrm{max}(1,e^{\mathrm{\Delta }S_{\mathrm{HMC}}})`$, where $`\mathrm{\Delta }S_{\mathrm{HMC}}=S_{\mathrm{HMC}}[P^{},A^{}]S_{\mathrm{HMC}}[P,A]`$, and duplicate the old configuration when rejected. By adding such a Metropolis accept/reject procedure, we can preserve the detailed balance. The step size $`\mathrm{\Delta }\tau `$ for the time evolution should be small enough to keep the acceptance rate reasonably high. The discretized Hamiltonian equation is given by
$`(P_\mu ^{(1/2)})_{ij}`$ $`=`$ $`(P_\mu ^{(0)})_{ij}{\displaystyle \frac{\mathrm{\Delta }\tau }{2}}{\displaystyle \frac{dS_{\mathrm{eff}}}{d(A_\mu )_{ij}}}(A_\mu ^{(0)}),`$
$`(A_\mu ^{(1)})_{ij}`$ $`=`$ $`(A_\mu ^{(0)})_{ij}+\mathrm{\Delta }\tau (P_\mu ^{(1/2)})_{ji},`$ (51)
$`(P_\mu ^{(n+1/2)})_{ij}`$ $`=`$ $`(P_\mu ^{(n1/2)})_{ij}\mathrm{\Delta }\tau {\displaystyle \frac{dS_{\mathrm{eff}}}{d(A_\mu )_{ij}}}(A_\mu ^{(n)}),`$
$`(A_\mu ^{(n+1)})_{ij}`$ $`=`$ $`(A_\mu ^{(n)})_{ij}+\mathrm{\Delta }\tau (P_\mu ^{(n+1/2)})_{ji},`$ (52)
$`(P_\mu ^{(\nu )})_{ij}`$ $`=`$ $`(P_\mu ^{(\nu 1/2)})_{ij}{\displaystyle \frac{\mathrm{\Delta }\tau }{2}}{\displaystyle \frac{dS_{\mathrm{eff}}}{d(A_\mu )_{ij}}}(A_\mu ^{(\nu )}),`$ (53)
where $`n=1,2,\mathrm{},\nu 1`$ and $`T=\nu \mathrm{\Delta }\tau `$, and we have introduced the short-hand notation $`P_\mu ^{(r)}=P_\mu (r\mathrm{\Delta }\tau )`$ and $`A_\mu ^{(r)}=A_\mu (r\mathrm{\Delta }\tau )`$. At each step of the “Molecular Dynamics”, we have to calculate the inverse $`^1`$, and at the end of each trajectory, we have to calculate $`det`$. These are the dominant part of the numerical calculation, and it requires a CPU time of the order of O$`(N^6)`$.
The hybrid Monte Carlo algorithm involves two parameters $`T`$ and $`\mathrm{\Delta }\tau `$, which can be optimized in such a way that the computational effort for obtaining one statistically independent configuration is minimized. The optimization can be done in a standard way . First we fix $`T`$ and optimize $`\mathrm{\Delta }\tau `$ so that the effective speed of motion in the configuration space, which is given by the acceptance rate times $`\mathrm{\Delta }\tau `$, is maximized. Using the $`\mathrm{\Delta }\tau `$ optimized for each $`T`$, we minimize the autocorrelation time (in units of “Molecular Dynamics step”) with respect to $`T`$. For instance, at $`N=16`$ and $`\alpha =0.0`$ we obtain the optimal values $`\mathrm{\Delta }\tau 0.006`$ and $`T1.0`$.
## Appendix B One-loop free energy
In this section we formulate the perturbation theory around fuzzy sphere solutions, and derive the one-loop free energy. We decompose $`A_\mu `$, $`\psi `$ and $`\overline{\psi }`$ into the classical background and the fluctuation as
$`A_\mu `$ $`=`$ $`X_\mu +\stackrel{~}{A}_\mu ,`$ (54)
$`\psi `$ $`=`$ $`\chi +\stackrel{~}{\psi },\text{ }\overline{\psi }=\overline{\chi }+\stackrel{~}{\overline{\psi }},`$ (55)
and obtain the free energy around the classical solutions by integrating over $`\stackrel{~}{A}_\mu `$, $`\stackrel{~}{\psi }`$ and $`\stackrel{~}{\overline{\psi }}`$ perturbatively. Here we take the classical solution to be the $`k`$ coincident fuzzy spheres $`X_\mu =A_\mu ^{(k\mathrm{S}^2)}`$, $`\chi =\overline{\chi }=0`$, which includes the single fuzzy sphere as a special case $`k=1`$.
In order to remove the zero modes associated with the SU($`N`$) invariance, we introduce the gauge fixing term and the corresponding ghost term
$`S_{\mathrm{g}.\mathrm{f}.}={\displaystyle \frac{N}{2}}\text{tr}[X_\mu ,A_\mu ]^2,`$ (56)
$`S_{\mathrm{ghost}}=N\text{tr}\left([X_\mu ,\overline{c}][A_\mu ,c]\right)=N\text{tr}\left(\overline{c}[X_\mu ,[A_\mu ,c]]\right),`$ (57)
where $`c`$ and $`\overline{c}`$ are the ghost and anti-ghost fields respectively.
The total action $`S_{\mathrm{total}}=S+S_{\mathrm{g}.\mathrm{f}.}+S_{\mathrm{ghost}}`$ can be written as
$`S_{\mathrm{total}}`$ $`=`$ $`S_{\mathrm{cl}}+S_{\mathrm{kin}}+S_{\mathrm{int}},`$ (58)
$`S_{\mathrm{cl}}`$ $`=`$ $`N\text{tr}\left({\displaystyle \frac{1}{4}}[X_\mu ,X_\nu ]^2+{\displaystyle \frac{2}{3}}i\alpha {\displaystyle \underset{i,j,k=1}{\overset{3}{}}}ϵ_{ijk}X_iX_jX_k\right),`$ (59)
$`S_{\mathrm{kin}}`$ $`=`$ $`N\text{tr}([\stackrel{~}{A}_\mu ,\stackrel{~}{A}_\nu ][X_\mu ,X_\nu ]+i\alpha {\displaystyle \underset{i,j,k=1}{\overset{3}{}}}ϵ_{ijk}[\stackrel{~}{A}_i,\stackrel{~}{A}_j]X_k`$ (60)
$`{\displaystyle \frac{1}{2}}[X_\mu ,\stackrel{~}{A}_\nu ]^2+\overline{c}[X_\mu ,[X_\mu ,c]]\stackrel{~}{\overline{\psi }}\mathrm{\Gamma }_\mu [X_\mu ,\stackrel{~}{\psi }]),`$
$`S_{\mathrm{int}}`$ $`=`$ $`N\text{tr}([\stackrel{~}{A}_\mu ,\stackrel{~}{A}_\nu ][X_\mu ,\stackrel{~}{A}_\nu ]{\displaystyle \frac{1}{4}}[\stackrel{~}{A}_\mu ,\stackrel{~}{A}_\nu ][\stackrel{~}{A}_\mu ,\stackrel{~}{A}_\nu ]`$ (61)
$`+{\displaystyle \frac{2}{3}}i\alpha {\displaystyle \underset{i,j,k=1}{\overset{3}{}}}ϵ_{ijk}\stackrel{~}{A}_i\stackrel{~}{A}_j\stackrel{~}{A}_k+\overline{c}[X_\mu ,[\stackrel{~}{A}_\mu ,c]]\stackrel{~}{\overline{\psi }}\mathrm{\Gamma }_\mu [\stackrel{~}{A}_\mu ,\stackrel{~}{\psi }]).`$
The linear terms in $`\stackrel{~}{A}_\mu `$ cancel since $`X_\mu `$ satisfies the classical equation of motion.
Noting that the background configuration $`X_\mu `$ includes a factor of $`\alpha `$, we can rescale the fluctuations as $`\stackrel{~}{A}_\mu \alpha \stackrel{~}{A}_\mu `$, $`c\alpha c`$, $`\overline{c}\alpha \overline{c}`$, $`\stackrel{~}{\psi }\alpha ^{\frac{3}{2}}\stackrel{~}{\psi }`$, $`\stackrel{~}{\overline{\psi }}\alpha ^{\frac{3}{2}}\stackrel{~}{\overline{\psi }}`$ so that all the terms in the total action $`S_{\mathrm{total}}`$ become proportional to $`\alpha ^4`$. This means that the expansion parameter of the present perturbation theory is $`\frac{1}{\alpha ^4}`$.
The free energy $`W`$ is defined by
$`\text{e}^W`$ $`=`$ $`{\displaystyle \text{d}\stackrel{~}{A}\text{d}c\text{d}\overline{c}\text{d}\stackrel{~}{\psi }\text{d}\stackrel{~}{\overline{\psi }}\text{e}^{S_{\mathrm{total}}}},`$ (62)
which can be calculated as a perturbative expansion $`W=_{j=0}^{\mathrm{}}W_j`$, where $`W_j=\text{O}(\alpha ^{4(1j)})`$. The classical part is simply given by $`W_0=S_\mathrm{b}[X]`$. In order to evaluate the one-loop term $`W_1`$, we note that the kinetic terms can be written as
$$S_{\mathrm{kin}}=N\text{tr}\left(\frac{1}{2}\stackrel{~}{A}_\nu (𝒫_\lambda )^2\stackrel{~}{A}_\nu +\overline{c}(𝒫_\lambda )^2c\right)N\text{tr}\left(\stackrel{~}{\overline{\psi }}\mathrm{\Gamma }_\mu 𝒫_\mu \stackrel{~}{\psi }\right),$$
(63)
where we have introduced the operator $`𝒫_\mu `$ which acts on a traceless $`N\times N`$ matrix $`M`$ as
$$𝒫_\mu M[X_\mu ,M].$$
(64)
Then the one-loop term can be expressed as
$`W_1`$ $`=`$ $`W_{1,\mathrm{b}}+W_{1,\mathrm{f}},`$ (65)
$`W_{1,\mathrm{b}}`$ $`=`$ $`𝒯r\mathrm{log}\left\{N(𝒫_\mu )^2\right\},`$ (66)
$`W_{1,\mathrm{f}}`$ $`=`$ $`𝒯r^{}\mathrm{log}(N\mathrm{\Gamma }_\mu 𝒫_\mu ),`$ (67)
where the symbol $`𝒯r`$ denotes the trace in the $`(N^21)`$-dimensional linear space which consists of traceless $`N\times N`$ matrices, and $`𝒯r^{}`$ includes the trace over spinor indices as well.
### B.1 Single fuzzy sphere
Let us first consider the single fuzzy sphere $`X_\mu =A_\mu ^{(\mathrm{S}^2)}`$. The classical part is given by
$$W_0=\frac{1}{24}N^2\alpha ^4(N^21),$$
(68)
and the one-loop terms can be written as
$`W_{1,\mathrm{b}}`$ $`=`$ $`𝒯r\mathrm{log}(N\alpha ^2𝒬),`$ (69)
$`W_{1,\mathrm{f}}`$ $`=`$ $`𝒯r^{}\mathrm{log}(N\alpha 𝒟).`$ (70)
The operators $`𝒬`$ and $`𝒟`$ are defined as
$$𝒬=\underset{i=1}{\overset{3}{}}(_i)^2,𝒟=\underset{i=1}{\overset{3}{}}\sigma _i_i,$$
(71)
where $`_i`$ acts on a traceless $`N\times N`$ matrix $`M`$ as $`_iM[L_i^{(N)},M]`$. In order to evaluate the one-loop terms, we need to solve the eigenvalue problem of the operators $`𝒬`$ and $`𝒟`$.
The eigenvectors of the operator $`𝒬`$ are given by the “matrix spherical harmonics” $`Y_{lm}`$ ($`l=0,1,\mathrm{},N1`$ and $`m=l,\mathrm{},l`$), which span a complete basis of the space of $`N\times N`$ matrices and have the properties analogous to the usual spherical harmonics such as
$`{\displaystyle \frac{1}{N}}\text{tr}\left(Y_{lm}^{}Y_{l^{}m^{}}\right)`$ $`=`$ $`\delta _{ll^{}}\delta _{mm^{}},`$ (72)
$`Y_{lm}^{}`$ $`=`$ $`(1)^mY_{l,m}.`$ (73)
The corresponding eigenvalues are given by $`l(l+1)`$; i.e.,
$$𝒬Y_{lm}=l(l+1)Y_{lm}.$$
(74)
Thus the one-loop term from the bosonic contribution is obtained as
$$W_{1,\mathrm{b}}=\underset{l=1}{\overset{N1}{}}(2l+1)\mathrm{log}\left[N\alpha ^2l(l+1)\right].$$
(75)
Here $`l=0`$ has been omitted from the sum since the trace $`𝒯r`$ in (66) should be taken in the space of traceless $`N\times N`$ matrices.
In order to solve the eigenvalue problem of the operator $`𝒟`$, we note that
$`𝒟={\displaystyle \underset{i=1}{\overset{3}{}}}(𝒥_i)^2𝒬{\displaystyle \frac{3}{4}},`$ (76)
where we have defined the “total angular momentum” operator
$`𝒥_i=_i+{\displaystyle \frac{\sigma _i}{2}}.`$ (77)
By making a linear combination of eigenvectors of $`𝒬`$ with the eigenvalue $`l(l+1)`$, we can construct the eigenvectors of both $`_{i=1}^3(𝒥_i)^2`$ and $`𝒥_3`$ with the eigenvalues $`j(j+1)`$ and $`m`$, respectively, where $`j`$ can be either $`j=l+\frac{1}{2}`$ ($`l=0,\mathrm{},N1`$) or $`j=l\frac{1}{2}`$ ($`l=1,\mathrm{},N1`$), and $`m`$ takes half-integer values in the range $`|m|j`$. Explicitly, the eigenvectors are given by the “matrix spinorial-spherical harmonics”
$`𝒴_{l+\frac{1}{2},m}`$ $`=`$ $`\sqrt{{\displaystyle \frac{l+\frac{1}{2}+m}{2l+1}}}Y_{l,m\frac{1}{2}}|+\sqrt{{\displaystyle \frac{l+\frac{1}{2}m}{2l+1}}}Y_{l,m+\frac{1}{2}}|,`$ (78)
$`𝒴_{}^{}{}_{l\frac{1}{2},m}{}^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{l+\frac{1}{2}m}{2l+1}}}Y_{l,m\frac{1}{2}}|\sqrt{{\displaystyle \frac{l+\frac{1}{2}+m}{2l+1}}}Y_{l,m+\frac{1}{2}}|,`$ (79)
for the cases $`j=l+\frac{1}{2}`$ and $`j=l\frac{1}{2}`$, respectively. Here the symbol $`|`$ and $`|`$ denotes the two-dimensional eigenvectors of $`\sigma _3`$ corresponding to the eigenvalues $`1`$ and $`1`$, respectively. From eq. (76), the “matrix spinorial-spherical harmonics” are also eigenvectors of $`𝒟`$ with the eigenvalues
$`𝒟`$ $`=`$ $`j(j+1)l(l+1){\displaystyle \frac{3}{4}}`$ (80)
$`=`$ $`\{\begin{array}{cc}l\hfill & \text{ for}j=l+\frac{1}{2}\hfill \\ (l+1)\hfill & \text{ for}j=l\frac{1}{2}.\hfill \end{array}`$ (83)
Namely we have the relation
$`𝒟𝒴_{l+\frac{1}{2},m}`$ $`=`$ $`l𝒴_{l+\frac{1}{2},m},`$ (84)
$`𝒟𝒴_{}^{}{}_{l\frac{1}{2},m}{}^{}`$ $`=`$ $`(l+1)𝒴_{}^{}{}_{l\frac{1}{2},m}{}^{}.`$ (85)
Thus the one-loop term from the fermionic contribution is obtained as
$`W_{1,\mathrm{f}}=\left[{\displaystyle \underset{l=1}{\overset{N1}{}}}\mathrm{\hspace{0.17em}2}(l+1)\mathrm{log}(N\alpha l)+{\displaystyle \underset{l=1}{\overset{N1}{}}}\mathrm{\hspace{0.17em}2}l\mathrm{log}\{N\alpha (l+1)\}\right],`$ (86)
where $`l=0`$ has been omitted from the first sum since the trace $`𝒯r^{}`$ in (67) should be taken in the space of traceless $`N\times N`$ matrices (and over the spinor indices). Let us rewrite the above expression as
$`W_{1,\mathrm{f}}=W_{1,\mathrm{b}}(N^21)\mathrm{log}N+\mathrm{log}N.`$ (87)
Thus the fermionic contribution cancel the bosonic contribution up to the $`\alpha `$-independent constant. From (68) and (87) the one-loop free energy for the single fuzzy sphere is obtained at large $`N`$ as
$$W_{1\mathrm{loop}}N^2\left(\frac{1}{24}\stackrel{~}{\alpha }^4\mathrm{log}N\right).$$
(88)
### B.2 $`k`$ coincident fuzzy spheres
Next we consider the $`k`$ coincident fuzzy spheres $`X_\mu =A_\mu ^{(k\mathrm{S}^2)}`$. The classical part of the free energy is given by
$$W_0=\frac{1}{24}\stackrel{~}{\alpha }^4(n^21).$$
(89)
In order to calculate the one-loop term, we consider the $`n\times n`$ version of the matrix spherical harmonics $`Y_{lm}^{(n)}`$, and define
$$Y_{lm}^{(a,b)}Y_{lm}^{(n)}𝐞^{(a,b)},$$
(90)
where $`𝐞^{(a,b)}`$ denotes a $`k\times k`$ matrix whose ($`a,b`$) element is 1 and all the other elements are zero. The $`N\times N`$ matrices $`Y_{lm}^{(a,b)}`$ form a complete basis of the space of $`N\times N`$ matrices, and they are the eigenvectors of the operator $`(𝒫_\mu )^2`$ for the present background; i.e.,
$$(𝒫_\mu )^2Y_{lm}^{(a,b)}=\alpha ^2l(l+1)Y_{lm}^{(a,b)}.$$
(91)
Let us note that the operator $`(𝒫_\mu )^2`$ has $`k^2`$ zero modes corresponding to $`l=m=0`$ with arbitrary $`(a,b)`$. The mode $`_{a=1}^kY_{00}^{(a,a)}`$ corresponds to the trace mode, which should be omitted due to the traceless condition. Here we omit the other zero modes by hand, and simply consider the non-zero modes. Then the one-loop term $`W_{1,\mathrm{b}}`$ from the bosonic contribution is obtained as
$$W_{1,\mathrm{b}}=k^2\underset{l=1}{\overset{n1}{}}(2l+1)\mathrm{log}\left[N\alpha ^2l(l+1)\right].$$
(92)
The calculation of the fermionic one-loop term proceeds in the same way except that we have to use the matrix spinorial-spherical harmonics for each of $`k^2`$ blocks. In this case we have $`2k^2`$ zero modes, and two of them correspond to the trace mode. We omit the other zero modes by hand and obtain
$`W_{1,\mathrm{f}}=k^2\left[{\displaystyle \underset{l=1}{\overset{n1}{}}}\mathrm{\hspace{0.17em}2}(l+1)\mathrm{log}(N\alpha l)+{\displaystyle \underset{l=1}{\overset{n1}{}}}\mathrm{\hspace{0.17em}2}l\mathrm{log}\{N\alpha (l+1)\}\right].`$ (93)
The cancellation between the bosonic contribution and the fermionic one occurs here as well, and the one-loop free energy is given by
$$W_{1\mathrm{loop}}N^2\left(\frac{1}{24k^2}\stackrel{~}{\alpha }^4\mathrm{log}N\right).$$
(94)
## Appendix C One-loop calculation of various observables
In this section we apply the perturbation theory discussed in the previous section to the one-loop calculation of various observables which are studied by Monte Carlo simulations in this paper. We take the background to be $`k`$ coincident fuzzy spheres $`X_\mu =A_\mu ^{(k\mathrm{S}^2)}`$, but the results for the single fuzzy sphere can be readily obtained by setting $`k=1`$. As in appendix B.2, we omit the zero modes for $`k2`$.
We note that the number of loops in the relevant diagrams can be less than the order of $`1/\alpha ^4`$ in the perturbative expansion since we are expanding the theory around a nontrivial background. At the one-loop level, the only nontrivial task is to evaluate the tadpole $`(\stackrel{~}{A}_\mu )_{ij}`$ explicitly, which, however, turns out to vanish due to supersymmetry.
### C.1 Propagators and the tadpole
The propagators for $`\stackrel{~}{A}_\mu `$, the ghosts and the fermion fields are given respectively as
$`(\stackrel{~}{A}_\mu )_{ij}(\stackrel{~}{A}_\nu )_{kl}_0`$ $`=`$ $`\delta _{\mu \nu }{\displaystyle \frac{1}{n}}{\displaystyle \underset{ab}{}}{\displaystyle \underset{l=1}{\overset{n1}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \frac{1}{N\alpha ^2l(l+1)}}\left(Y_{lm}^{(a,b)}\right)_{ij}\left(Y_{lm}^{(a,b)}\right)_{kl},`$ (95)
$`(c)_{km}(\overline{c})_{pq}_0`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{ab}{}}{\displaystyle \underset{l=1}{\overset{n1}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \frac{1}{N\alpha ^2l(l+1)}}\left(Y_{lm}^{(a,b)}\right)_{ij}\left(Y_{lm}^{(a,b)}\right)_{kl},`$ (96)
$`(\psi )_{ij}(\overline{\psi })_{kl}_0`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{ab}{}}{\displaystyle \underset{l=0}{\overset{n1}{}}}{\displaystyle \underset{m=l\frac{1}{2}}{\overset{l+\frac{1}{2}}{}}}{\displaystyle \frac{1}{N\alpha l}}\left(𝒴_{l+\frac{1}{2},m}^{(a,b)}\right)_{ij}\left(𝒴_{l+\frac{1}{2},m}^{(a,b)}\right)_{kl}`$ (97)
$`+{\displaystyle \frac{1}{n}}{\displaystyle \underset{ab}{}}{\displaystyle \underset{l=1}{\overset{n1}{}}}{\displaystyle \underset{m=l+\frac{1}{2}}{\overset{l\frac{1}{2}}{}}}{\displaystyle \frac{1}{N\alpha (l+1)}}\left(𝒴_{}^{}{}_{l\frac{1}{2},m}{}^{(a,b)}\right)_{ij}\left(𝒴_{}^{}{}_{l\frac{1}{2},m}{}^{(a,b)}\right)_{kl},`$
where the symbol $`_0`$ refers to the expectation value calculated using the kinetic term $`S_{\mathrm{kin}}`$ only. The tadpole $`\stackrel{~}{A}_i_{1\mathrm{loop}}`$ ($`i=1,2,3`$) at the one-loop level can be calculated as
$`\stackrel{~}{A}_i_{1\mathrm{loop}}`$ $`=`$ $`N\stackrel{~}{A}_i\text{tr}\left([\stackrel{~}{A}_\nu ,\stackrel{~}{A}_\rho ][X_\nu ,\stackrel{~}{A}_\rho ]\right)_0N\stackrel{~}{A}_i\text{tr}\left(\overline{c}[X_\nu ,[\stackrel{~}{A}_\nu ,c]]\right)_0`$ (98)
$`N\stackrel{~}{A}_i\text{tr}\left(\stackrel{~}{\overline{\psi }}\mathrm{\Gamma }_\nu [\stackrel{~}{A}_\nu ,\stackrel{~}{\psi }]\right)_0.`$
By redoing the calculation in ref. in the present model, we find that the bosonic contribution and the fermionic contribution cancel each other even at finite $`N`$. We also find that $`\stackrel{~}{A}_4=0`$ to all orders in perturbation theory due to parity invariance $`A_4A_4`$.
### C.2 One-loop results for various observables
Using the propagator and the tadpole obtained in the previous section, we can evaluate various observables easily at the one-loop level. For instance the two-point function $`\frac{1}{N}\text{tr}(A_\mu )^2`$ can be evaluated as follows. Let us decompose it as
$$\frac{1}{N}\text{tr}(A_\mu )^2=\frac{1}{N}\underset{i=1}{\overset{3}{}}\text{tr}(A_i)^2+\frac{1}{N}\text{tr}(A_4)^2.$$
(99)
Each term on the r.h.s. can be calculated as
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{3}{}}}\text{tr}(A_i)^2_{1\mathrm{loop}}`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{3}{}}}\left[\text{tr}(X_iX_i)+2\text{tr}\left(X_i\stackrel{~}{A}_i_{1\mathrm{loop}}\right)+\text{tr}(\stackrel{~}{A}_i)^2_0\right]`$ (100)
$`=`$ $`\alpha ^2\left[{\displaystyle \frac{1}{4}}(n^21)+0+{\displaystyle \frac{3}{\alpha ^4n^2}}{\displaystyle \underset{l=1}{\overset{n1}{}}}{\displaystyle \frac{2l+1}{l(l+1)}}\right],`$
$`{\displaystyle \frac{1}{N}}\text{tr}(A_4)^2_{1\mathrm{loop}}`$ $`=`$ $`{\displaystyle \frac{1}{N}}\text{tr}(\stackrel{~}{A}_4)^2_0={\displaystyle \frac{1}{\alpha ^2n^2}}{\displaystyle \underset{l=1}{\overset{n1}{}}}{\displaystyle \frac{2l+1}{l(l+1)}}.`$ (101)
At large $`N`$ with fixed $`\stackrel{~}{\alpha }=\alpha \sqrt{N}`$, we get
$`{\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{N}}\text{tr}(A_\mu )^2_{1\mathrm{loop}}{\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{3}{}}}\text{tr}(A_i)^2_{1\mathrm{loop}}{\displaystyle \frac{1}{4k^2}}\stackrel{~}{\alpha }^2.`$ (102)
The Chern-Simons term $`M`$ can be evaluated as
$`M_{1\mathrm{loop}}`$ $`=`$ $`{\displaystyle \frac{2i}{3N}}ϵ_{ijk}\left[\text{tr}(X_iX_jX_k)+3\text{tr}\left(X_iX_j\stackrel{~}{A}_k_{1\mathrm{loop}}\right)\right]`$ (103)
$`=`$ $`{\displaystyle \frac{1}{6}}\alpha ^3(n^21).`$
At large $`N`$ with fixed $`\stackrel{~}{\alpha }=\alpha \sqrt{N}`$, we obtain
$$\frac{1}{\sqrt{N}}M_{1\mathrm{loop}}\frac{1}{6k^2}\stackrel{~}{\alpha }^3.$$
(104)
The observable $`\frac{1}{N}\text{tr}F^2`$ can be calculated in a similar manner, but we can also obtain it from the exact result (22) using (104) as
$`{\displaystyle \frac{1}{N}}\text{tr}(F_{\mu \nu })^2_{1\mathrm{loop}}`$ $`=`$ $`6\left(1{\displaystyle \frac{1}{N^2}}\right)3\alpha M_{1\mathrm{loop}}`$ (105)
$``$ $`{\displaystyle \frac{1}{2k^2}}\stackrel{~}{\alpha }^4+6.`$ (106)
### C.3 Alternative derivation
Since $`\text{tr}F^2`$ and $`M`$ are the operators that appear in the action $`S`$, we can obtain their expectation values easily by using the free energy calculated for the $`k`$ coincident fuzzy sphere in Appendix B. Let us deform the bosonic action as
$`S_\mathrm{b}(\beta _1,\beta _2,\alpha )=N^2\left[{\displaystyle \frac{1}{4}}\beta _1\text{tr}(F_{\mu \nu })^2+\beta _2\alpha M\right]`$ (107)
with two free parameters $`\beta _1`$, $`\beta _2`$, and define the corresponding free energy as
$`\text{e}^{W(\beta _1,\beta _2,\alpha )}={\displaystyle \text{d}A\text{d}\psi \text{d}\overline{\psi }\text{e}^{S_\mathrm{b}(\beta _1,\beta _2,\alpha )S_\mathrm{f}}}.`$ (108)
Then $`\text{tr}(F_{\mu \nu })^2`$ and $`M`$ can be obtained by
$`{\displaystyle \frac{1}{N}}\text{tr}(F_{\mu \nu })^2`$ $`=`$ $`{\displaystyle \frac{4}{N^2}}{\displaystyle \frac{W}{\beta _1}}|_{\beta _1=\beta _2=1},`$ (109)
$`M`$ $`=`$ $`{\displaystyle \frac{1}{\alpha N^2}}{\displaystyle \frac{W}{\beta _2}}|_{\beta _1=\beta _2=1}.`$ (110)
By rescaling the integration variables as $`A_\mu \beta _1^{\frac{1}{4}}A_\mu `$, $`\psi \beta _1^{\frac{1}{8}}\psi `$ and $`\overline{\psi }\beta _1^{\frac{1}{8}}\overline{\psi }`$, we get
$$W(\beta _1,\beta _2,\alpha )=\frac{3}{2}(N^21)\mathrm{log}\beta _1+W(1,1,\alpha \beta _1^{\frac{3}{4}}\beta _2).$$
(111)
Using the one-loop result
$$W(1,1,\alpha )_{1\mathrm{loop}}=\frac{N^2}{24}\alpha ^4(n^21)k^2(n^21)\mathrm{log}N+k^2\mathrm{log}n,$$
(112)
we can reproduce (103) and (105).
|
warning/0506/astro-ph0506427.html
|
ar5iv
|
text
|
# The theory of canonical perturbations applied to attitude dynamics and to the Earth rotation. Osculating and nonosculating Andoyer variables.
## 1 The Hamiltonian approach to rotational dynamics
### 1.1 Historical preliminaries
The perturbed rotation of a rigid body has long been among the key topics of both spacecraft engineering (Giacaglia & Jefferys 1971; Zanardi & Vilhena de Moraes 1999) and planetary astronomy (Kinoshita 1977; Laskar & Robutel 1993; Touma & Wisdom 1994; Mysen 2004, 2006). While free spin (the Euler-Poinsot problem) permits an analytic solution in terms of the elliptic Jacobi functions, perturbed motion typically requires numerical treatment, though sometimes it can be dealt with by analytical means (like, for example, in Kinoshita 1977 ). Perturbation may come from a physical torque, or from an inertial torque caused by the frame noninertiality, or from nonrigidity (Getino & Ferrándiz 1990; Escapa, Getino & Ferrándiz 2001, 2002). The free-spin Hamiltonian, expressed through the Euler angles and their conjugate momenta, is independent of one of the angles, which reveals an internal symmetry of the problem. In fact, this problem possesses an even richer symmetry (Deprit & Elipe 1993), whose existence indicates that the unperturbed Euler-Poinsot dynamics can be reduced to one degree of freedom. The possibility of such reduction is not readily apparent and can be seen only under certain choices of variables. These variables, in analogy with the orbital mechanics, are called rotational elements. It is convenient to treat the forced-rotation case as a perturbation expressed through those elements.
The Andoyer variables are often chosen as rotational elements (Andoyer 1923, Giacaglia & Jefferys 1971, Kinoshita 1972), though other sets of canonical elements have appeared in the literature (Richelot 1850; Serret 1866; Peale 1973, 1976; Deprit & Elipe 1993; Fukushima & Ishizaki 1994)<sup>1</sup><sup>1</sup>1 Some authors use the term “Serret-Andoyer elements.” This is not correct, because the set of elements introduced by Richelot (1850) and Serret (1862) differs from the one employed by Andoyer (1923).. After a transition to rotational elements is performed within the undisturbed Euler-Poinsot setting, the next step is to extend this method to a forced-rotation case. To this end, one will have to express the torques via the elements. On completion of the integration, one will have to return back from the elements to the original, measurable, quantities – i.e., to the Euler angles and their time derivatives.
### 1.2 The Kinoshita-Souchay theory of rigid-Earth rotation
The Hamiltonian approach to spin dynamics has found its most important application in the theory of Earth rotation. A cornerstone work on this topic was carried out by Kinoshita (1977) who switched from the Euler angles defining the Earth orientation to the Andoyer variables, and treated their dynamics by means of the Hori (1966) and Deprit(1969) technique.<sup>2</sup><sup>2</sup>2 For an introduction into the Hori-Deprit method see Boccaletti & Pucacco (2002) and Kholshevnikov (1975, 1985). Kinoshita (1977) referred only to the work by Hori (1966). Then he translated the results of this development back into the language of Euler’s angles and provided the precessional and nutation spectrum. Later his approach was extended to a much higher precision by Kinoshita & Souchay (1990) and Souchay, Losley, Kinoshita & Folgueira (1999).
### 1.3 Subtle points
When one is interested only in the orientation of the rotator, it is sufficient to have expressions for the Euler angles as functions of the elements. However, when one needs to know also the instantaneous angular velocity, one needs expressions for the Euler angles’ time derivatives. This poses the following question: if we write down the expressions for the Euler angles’ derivatives via the canonical elements in the free-spin case, will these expressions stay valid under perturbation? In the parlance of orbital mechanics, this question may be formulated like this: are the canonical elements always osculating? As we shall demonstrate below, under angular-velocity-dependent disturbances the condition of osculation is incompatible with that of canonicity, and therefore expression of the angular velocity via the canonical elements will, under such types of perturbations, become nontrivial.
In 2004 the question acquired a special relevance to the Earth-rotation theory. While the thitherto available observations referred to the orientation of the Earth figure (Kinoshita et al. 1978), a technique based on ring laser gyroscope provided a direct measurement of the instantaneous angular velocity of the Earth relative to an inertial frame (Schreiber et al. 2004, Petrov 2007).
Normally, rotational elements are chosen to have evident physical interpretation. For example, the Andoyer variable $`G`$ coincides with the absolute value of the body’s spin angular momentum, while two other variables, $`H`$ and $`L`$, are chosen to coincide, correspondingly, with the $`Z`$-component of the angular momentum in the inertial frame, and with its $`z`$-component in the body frame. The other Andoyer elements, $`g,l,h`$, too, bear some evident meaning. Hence another important question: will the canonical rotational elements preserve their simple physical meaning also under disturbance?
## 2 The canonical perturbation theory <br>in orbital and attitude dynamics
### 2.1 Kepler and Euler
In orbital dynamics, a Keplerian conic, emerging as an undisturbed two-body orbit, is regarded to be a “simple motion,” so that all the other available motions are conveniently considered as distortions of such conics, distortions implemented through endowing the orbital constants $`C_j`$ with their own time dependence. Points of the orbit can be contributed by the “simple curves” either in a nonosculating fashion, as in Fig. 1, or in the osculating manner, as in Fig. 2.
The disturbances, causing the evolution of the motion from one instantaneous conic to another, are the primary’s nonsphericity, the gravitational pull of other bodies, the atmospheric
Fig. 1. The perturbed orbit is a set of points belonging to a sequence of confocal instantaneous ellipses that are not supposed to be tangent or even coplanar to the orbit. As a result, the physical velocity $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ (tangent to the orbit) differs from the Keplerian velocity $`\stackrel{}{𝐠}`$ (tangent to the ellipse). To parameterise the depicted sequence of nonosculating ellipses, and to single it out of the other sequences, it is suitable to employ the difference between $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ and $`\stackrel{}{𝐠}`$, expressed as a function of time and six (nonosculating) orbital elements: $`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)=\dot{\stackrel{\mathbf{}}{𝒓}}(t,C_1,...,C_6)\stackrel{}{𝐠}(t,C_1,...,C_6).`$ In the literature, $`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)`$ is called gauge function or gauge velocity or, simply, gauge.
Fig. 2. The orbit is represented by a sequence of confocal instantaneous ellipses that are tangent to the orbit, i.e., osculating. Now, the physical velocity $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ (tangent to the orbit) coincides with the Keplerian velocity $`\stackrel{}{𝐠}`$ (tangent to the ellipse), so that their difference vanishes everywhere: $`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)=\mathrm{\hspace{0.17em}0}`$. This is the so-called Lagrange constraint or Lagrange gauge. Orbital elements obeying it are called osculating.
and radiation-caused drag, the relativistic corrections, and the non-inertiality of the reference system.
On Fig. 1 the orbit consists of points, each of which is donated by a representative of a certain family of “simple” curves (confocal ellipses). These instantaneous ellipses are not supposed to be tangent or even coplanar to the orbit. As a result, the physical velocity $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ (tangent to the orbit) differs from the Keplerian velocity $`\stackrel{}{𝐠}`$ (tangent to the ellipse). To parameterise the depicted sequence of nonosculating ellipses, and to single it out of the other sequences, it is suitable to employ the difference between $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ and $`\stackrel{}{𝐠}`$, expressed as a function of the time and the orbital elements: $`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)=\dot{\stackrel{\mathbf{}}{𝒓}}(t,C_1,...,C_6)\stackrel{}{𝐠}(t,C_1,...,C_6).`$ Evidently,
$`\dot{\stackrel{\mathbf{}}{𝒓}}={\displaystyle \frac{\stackrel{\mathbf{}}{𝒓}}{t}}+{\displaystyle \underset{j=1}{\overset{6}{}}}{\displaystyle \frac{C_j}{t}}\dot{C}_j=\stackrel{}{𝐠}+\stackrel{\mathbf{}}{𝚽},`$
i.e., the unperturbed Keplerian velocity is $`\stackrel{}{𝐠}\stackrel{\mathbf{}}{𝒓}/t`$, while the said difference $`\stackrel{\mathbf{}}{𝚽}`$ is the convective term that emerges when the instantaneous ellipses are being gradually altered by the perturbation (and when the orbital elements become time-dependent): $`\stackrel{\mathbf{}}{𝚽}=\left(\stackrel{\mathbf{}}{𝒓}/C_j\right)\dot{C}_j`$. When one fixes a particular functional dependence of $`\stackrel{\mathbf{}}{𝚽}`$ upon time and the elements, this function, $`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)`$, is called gauge function or gauge velocity or, simply, gauge.
On Fig. 2, the perturbed orbit is represented with a sequence of confocal instantaneous ellipses that are tangent to the orbit, i.e., osculating. Under this choice, the physical velocity $`\dot{\stackrel{\mathbf{}}{𝒓}}`$ (tangent to the orbit) will coincide with the Keplerian velocity $`\stackrel{}{𝐠}`$ (tangent to the ellipse), so that their difference $`\stackrel{\mathbf{}}{𝚽}(tC_1,...,C_6)`$ will vanish everywhere:
$`\stackrel{\mathbf{}}{𝚽}(t,C_1,...,C_6)\dot{\stackrel{\mathbf{}}{𝒓}}(t,C_1,...,C_6)\stackrel{}{𝐠}(t,C_1,...,C_6)={\displaystyle \underset{j=1}{\overset{6}{}}}{\displaystyle \frac{C_j}{t}}\dot{C}_j=0.`$
This, so-called Lagrange constraint or Lagrange gauge, is the necessary and sufficient condition of osculation of the orbital elements $`C_j`$ (Brouwer & Clemence 1961). Historically, the first attempt of using nonosculating elements dates back to Poincare (1897), though he never explored them from the viewpoint of a non-Lagrange constraint choice. (See also Abdullah & Albouy (2001), p. 430.) Parameterisation of nonosculation through a non-Lagrange constraint was offered in Efroimsky (2002a,b).
Similarly to orbital dynamics, in attitude dynamics, a complex spin can be presented as a sequence of instantaneous configurations borrowed from a family of some “simple rotations”. (Efroimsky 2004) It is convenient to employ in this role the motions exhibited by an undeformable free top experiencing no torques.<sup>3</sup><sup>3</sup>3 Here one opportunity will be to utilise in the role of “simple” motions the non-circular Eulerian cones described by the actual triaxial top, when it is unforced. Another opportunity will be to use, as “simple” motions, the circular Eulerian cones described by a dynamically symmetrical top (and to treat its actual triaxiality as another perturbation). The main result of our paper will be invariant under this choice. Each such undisturbed “simple motion” will be a trajectory on the three-dimensional manifold of the Euler angles (Synge & Griffith 1959). For the lack of a better term, we shall call these unperturbed motions “Eulerian cones,” implying that the loci of the rotational axis, which correspond to each such non-perturbed spin state, make closed cones (circular, for an axially symmetrical rotator; and elliptic for a triaxial one). Then, to implement a perturbed motion, we shall have to go from one Eulerian cone to another, just as in Fig. 1 and 2 we go from one Keplerian ellipse to another. Hence, similar to those pictures, a smooth “walk” over the instantaneous Eulerian cones may be osculating or nonosculating.
The torques, as well as the actual triaxiality of the top and the non-inertial nature of the reference frame, will then act as perturbations causing this “walk.” Perturbations of the latter two types depend not only upon the rotator’s orientation but also upon its angular velocity.<sup>4</sup><sup>4</sup>4 When we study the Earth rotation relative to the precessing plane of the Earth orbit about the Sun, the frame precession gives birth to a fictitious torque (sometimes called “inertial torque”) that depends upon the Earth’s angular velocity.
### 2.2 Delaunay and Andoyer
In orbital dynamics, we can express the Lagrangian of the reduced two-body problem via the spherical coordinates $`q_j=\{r,\phi ,\theta \}`$, then derive their conjugated momenta $`p_j`$ and the Hamiltonian $`(q,p)`$, and then carry out the Hamilton-Jacobi procedure (Plummer 1918), to arrive at the Delaunay variables
$`\{Q_1,Q_2,Q_3;P_1,P_2,P_3\}\{L,G,H;l_o,g,h\}=`$
(1)
$`\{\sqrt{\mu a},\sqrt{\mu a\left(1e^2\right)},\sqrt{\mu a\left(1e^2\right)}\mathrm{cos}i;M_o,\omega ,\mathrm{\Omega }\},`$
where $`\mu `$ denotes the reduced mass.
Similarly, in rotational dynamics one can define a state of a spinning top by the three Euler angles $`q_j=\{\phi ,\theta ,\psi \}`$ and their canonical momenta $`p_j=\{p_\phi ,p__\theta ,p_\psi \}`$; and then carry out a canonical transformation to the Andoyer elements<sup>5</sup><sup>5</sup>5 In attitude dynamics, the Andoyer elements $`l,g,h`$ play the role of coordinates, while $`L,G,H`$ are their conjugate momenta. In the orbital case, the Delaunay variables $`L,G,H`$ play the role of coordinates, while $`l,g,h`$ defined as in (1) act as momenta. Needless to say, this is merely a matter of convention. (See formulae (9.31 - 9.32) in Goldstein 1981.) For example in some textbooks on orbital mechanics the Hamiltonian perturbation is deliberately introduced with an opposite sign, while the Delaunay elements $`l,g,h`$, too, are defined with signs opposite to given in (1). Under such a convention, the Delaunay elements $`l,g,h`$ become coordinates, while $`L,G,H`$ act as momenta. $`\{l,g,h;L,G,H\}`$. By definition, the element $`G`$ is the magnitude of the angular-momentum vector, $`L`$ is the projection of the angular-momentum vector on the principal axis $`\widehat{𝐛}_3`$ of the body, while $`H`$ is the projection of the angular-momentum vector on the $`\widehat{𝐬}_3`$ axis of the inertial coordinate system. The variable $`h`$ conjugate to $`H`$ is the angle from the inertial reference longitude to the ascending node of the invariable plane (the one perpendicular to the angular momentum). The variable $`g`$ conjugate to $`G`$ is the angle from the ascending node of the invariable plane on the reference plane to the ascending node of the equator on the invariable plane. Finally, the variable conjugate to $`L`$ is the angle $`l`$ from the ascending node of the equator on the invariable plane to the the $`\widehat{𝐛}_1`$ body axis. Two auxiliary quantities defined through
$`\mathrm{cos}I={\displaystyle \frac{H}{G}},\mathrm{cos}J={\displaystyle \frac{L}{G}},`$
have obvious meaning: $`I`$ is the angle between the angular-momentum vector and the $`\widehat{𝐬}_3`$ space axis, while $`J`$ is the angle between the angular-momentum vector and the $`\widehat{𝐛}_3`$ principal axis of the body, as depicted on Fig. 3.
Andoyer (1923) introduced his variables in a manner different from canonical constants: while his variables $`G,H,h`$ are constants (for a free triaxial rotator), the other three, $`L,l,g`$, do evolve in time, because the Andoyer Hamiltonian of a free top
Fig. 3. A reference coordinate system (inertial or precessing) is constituted by axes $`𝐬_1,𝐬_2,𝐬_3`$. A body-fixed frame is defined by the principal axes $`𝐛_1,𝐛_2,𝐛_3`$. The third frame is constituted by the angular-momentum vector $`\stackrel{}{\text{G}}`$ and a plane orthogonal thereto (the so-called invariable plane). The lines of nodes are denoted with $`𝐢,𝐥,𝐣`$. The attitude of the body relative to the reference frame is given by the Euler angles $`h_f,I_f,\varphi _f`$. The orientation of the invariable plane with respect to the reference frame is determined by the angles $`h`$ and $`I`$. The inclination $`I`$ is equal to the angle that the angular-momentum vector $`\stackrel{}{\text{G}}`$ makes with the reference axis $`𝐬_3`$. The angle $`J`$ between the invariable plane and the body equator coincides with the angle that $`\stackrel{}{\text{G}}`$ makes with the major-inertia axis $`𝐛_3`$ of the body. The projections of the angular momentum toward the reference axis $`𝐬_3`$ and the body axis $`𝐛_3`$ are $`H=G\mathrm{cos}I`$ and $`L=G\mathrm{cos}J`$.
$`(g,h,l,G,H,L)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{sin}^2l}{A}}+{\displaystyle \frac{\mathrm{cos}^2l}{B}}\right)\left(G^2L^2\right)+{\displaystyle \frac{L^2}{2C}}`$
is a nonvanishing function of the variables $`l,L`$ and $`G`$. (Notations $`A,B,C`$ stand for the inertia matrix’ principal values that are assumed, without loss of generality, to obey the inequality $`ABC`$.) So, to make our analogy complete, we may carry out one more canonical transformation, from the regular Andoyer set $`\{l,g,h,L,G,H\}`$ to the “modified Andoyer set” $`\{l_o,g_o,h;L_o,G,H\}`$, where $`L_o,l_o,g_o`$ are the initial values of $`L,l,g`$. The modified set consists only of constants of integration, wherefore the appropriate Hamiltonian becomes nil.<sup>6</sup><sup>6</sup>6 We shall not write down the explicit form of this transformation, because it is sufficient for us to know that it is canonical. This follows from the group property of canonical transformations and from the fact that the transformations from $`\{l,g,h,L,G,H\}`$ to $`\{\phi ,\theta ,\psi ,p_\phi ,p_\theta ,p_\psi \}`$ and from $`\{l_o,g_o,h;L_o,G,H\}`$ to $`\{l,g,h,L,G,H\}`$ are canonical. The latter transformation is canonical, for it is simply the time evolution. This canonical transition from Andoyer-type variables to their initial values is not new – see Fukushima & Ishizaki (1994). Historically, the first set of rotational elements was constituted by constants (Richelot 1850). Serret (1866) found the generating function of a canonical transformation from $`\{\phi ,\theta ,\psi ,p_\phi ,p_\theta ,p_\psi \}`$ to that set. His development was polished by Radau (1869) and Tisserand (1889). The Serret-Richelot set consisted of the following constants: $`\{g_o,h,t_o;G,H,T_{kin}\}`$, where $`h,G`$ and $`H`$ coincide with the appropriate Andoyer elements, $`T_{kin}`$ is the rotational kinetic energy, $`t_o`$ is the initial moment of time, and $`g_o`$ is the initial value of the Andoyer element $`g`$. Therefore, the modified Andoyer set of variables is analogous to the Delaunay set with $`l_o=M_o`$, while the regular Andoyer elements are analogous to the Delaunay elements with $`l=M`$ used instead of $`l_o=M_o`$. We would stress that, in analogy with the orbital case, the variables $`h,G,H`$ are constants (and $`h_o=h,G_o=G,H_o=H`$) only in the unperturbed, free-spin, case.
To summarise this section, in both cases we start out with
$`\dot{q}={\displaystyle \frac{^{(o)}}{p}},\dot{p}={\displaystyle \frac{^{(o)}}{q}},`$ (2)
$`q`$ and $`p`$ being the coordinates and their conjugated momenta, in the orbital case, or the Euler angles and their momenta, in the rotation case. Then we switch, via a canonical transformation
$`q=f(Q,P,t),`$
(3)
$`p=\chi (Q,P,t),`$
to
$`\dot{Q}={\displaystyle \frac{^{}}{P}}=0,\dot{P}={\displaystyle \frac{^{}}{Q}}=\mathrm{\hspace{0.33em}0},^{}=\mathrm{\hspace{0.33em}0},`$ (4)
where $`Q`$ and $`P`$ denote the set of Delaunay elements, in the orbital case, or the (modified, as explained above) Andoyer set $`\{l_o,g_o,h;L_o,G,H\}`$, in the case of rigid-body rotation.
This scheme relies on the fact that, for an unperturbed Keplerian orbit (and, similarly, for an undisturbed Eulerian cone), its six-constant parameterisation may be chosen so that:
1. the parameters are constants and, at the same time, are canonical variables $`\{Q,P\}`$ with a zero Hamiltonian: $`^{}(Q,P)=\mathrm{\hspace{0.17em}0}`$;
2. for constant $`Q`$ and $`P`$, the transformation equations (3) are mathematically equivalent to the dynamical equations (2).
In practice, this scheme is implemented via the Hamilton-Jacobi procedure.
### 2.3 Canonical perturbation theory: canonicity versus osculation
Under perturbation, the “constants” $`Q,P`$ begin to evolve so that, after their insertion into
$`q=f(Q(t),P(t),t),`$
(5)
$`p=\chi (Q(t),P(t),t)`$
($`f`$ and $`\chi `$ being the same functions as in (3) ), the resulting motion obeys the disturbed equations
$`\dot{q}={\displaystyle \frac{\left(^{(o)}+\mathrm{\Delta }\right)}{p}},\dot{p}={\displaystyle \frac{\left(^{(o)}+\mathrm{\Delta }\right)}{q}}.`$ (6)
We also want our “constants” $`Q`$ and $`P`$ to remain canonical and to obey
$`\dot{Q}={\displaystyle \frac{\left(^{}+\mathrm{\Delta }^{}\right)}{P}},\dot{P}={\displaystyle \frac{\left(^{}+\mathrm{\Delta }^{}\right)}{Q}},`$ (7)
where
$`^{}=\mathrm{\hspace{0.33em}0}\text{and}\mathrm{\Delta }^{}(Q,Pt)=\mathrm{\Delta }(q(Q,P,t),p(Q,P,t),t).`$ (8)
Above all, we wish that the perturbed “constants” $`C_jQ_1,Q_2,Q_3,P_1,P_2,P_3`$ (the Delaunay elements, in the orbital case, or the modified Andoyer elements, in the rotation case) remain osculating. This means that the perturbed velocity will be expressed by the same function of $`C_j(t)`$ and $`t`$ as the unperturbed one used to. Let us check to what extent this optimism is justified. The perturbed velocity reads
$`\dot{q}=\text{g}+\mathrm{\Phi },`$ (9)
where
$`\text{g}(C(t),t){\displaystyle \frac{q(C(t),t)}{t}}`$ (10)
is the functional expression for the unperturbed velocity; and
$`\mathrm{\Phi }(C(t),t){\displaystyle \underset{j=1}{\overset{6}{}}}{\displaystyle \frac{q(C(t),t)}{C_j}}\dot{C}_j(t)`$ (11)
is the convective term. Since we chose the “constants” $`C_j`$ to make canonical pairs $`(Q,P)`$ obeying (7 \- 8) with vanishing $`^{}`$, then insertion of (7) into (11) will result in
$`\mathrm{\Phi }={\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{q}{Q_n}}\dot{Q}_n(t)+{\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{q}{P_n}}\dot{P}_n(t)={\displaystyle \frac{\mathrm{\Delta }(q,p)}{p}}.`$ (12)
We see that in some situations the canonicity requirement is incompatible with osculation.<sup>7</sup><sup>7</sup>7 For the first time, this observation was made in Efroimsky & Goldreich (2003). To be specific, under a momentum-dependent perturbation we still can use ansatz (5) for calculation of the coordinates and momenta, but cannot impose the osculation condition $`\mathrm{\Phi }=0`$ (i.e., we cannot use $`\dot{q}=\text{g}`$ for calculating the velocities). Instead, we must use (9) with the substitution (12). This generic rule applied both to orbital and rotational motions. Its application to the orbital case is illustrated by Fig. 2. There, the constants $`C_j=(Q_n,P_n)`$ parameterise instantaneous ellipses which, for nonzero $`\mathrm{\Phi }`$, are not tangent to the trajectory. In orbital mechanics, the variables preserving canonicity at the cost of osculation are called “contact elements” (term coined by Victor Brumberg). The osculating and contact variables coincide when the disturbance is velocity-independent. Otherwise, they differ already in the first order of the time-dependent perturbation (Efroimsky & Goldreich 2003, 2004). Luckily, in some situations, their secular parts differ only in the second order (Efroimsky 2005), a fortunate circumstance anticipated by Goldreich (1965), who came across these elements in a totally different context unrelated to canonicity.
The case of rotational motion will parallel the theory of orbits. Now, instead of the instantaneous Keplerian conics, one will deal with instantaneous Eulerian cones (i.e., with the loci of the rotational axis, corresponding to non-perturbed spin states). Indeed, the situation of an axially symmetric unsupported top at each instant of time is fully defined by the three Euler angles $`q_n=\varphi ,\theta ,\psi `$ and their time derivatives $`\dot{q}_n=\dot{\varphi },\dot{\theta },\dot{\psi }`$. The evolution of these six quantities is governed by three dynamical equations of the second order (the three projections of $`d\stackrel{}{𝐋}/dt=\stackrel{\mathbf{}}{𝝉}`$, where $`\stackrel{}{𝐋}`$ is the angular momentum and $`\stackrel{\mathbf{}}{𝝉}`$ is the torque) and, therefore, this evolution will depend upon the time and the six integration constants:
$`q_n=f_n(C_1,...,C_6,t),`$
(13)
$`\dot{q}_n=\text{g}_n(C_1,...,C_6,t),`$
where the functions $`\text{g}_n`$ and $`f_n`$ are interconnected via $`\text{g}_nf_n/t`$, for $`n=\psi ,\theta ,\varphi `$.
Under disturbance, the motion will be altered:
$`q_n=f_n(C_1(t),...,C_6(t),t),`$
(14)
$`\dot{q}_n=\text{g}_n(C_1(t),...,C_6(t),t)+\mathrm{\Phi }_n(C_1(t),...,C_6(t),t),`$
where
$`\mathrm{\Phi }_n(C_1(t),...,C_6(t),t){\displaystyle \underset{j=1}{\overset{6}{}}}{\displaystyle \frac{f_n}{C_j}}\dot{C}_j.`$ (15)
If we want the “constants” $`C_j`$ to constitute canonical pairs $`(Q,P)`$ obeying (7 \- 8), then insertion of (7) into (15) will result in
$`\mathrm{\Phi }_n(C_1(t),...,C_6(t),t){\displaystyle }{\displaystyle \frac{f_n}{Q}}\dot{Q}+{\displaystyle }{\displaystyle \frac{f_n}{P}}\dot{P}={\displaystyle \frac{\mathrm{\Delta }(q,p)}{p_n}},`$ (16)
so that the canonicity requirement (7 \- 8) violates the gauge freedom in a non-Lagrange fashion.
To draw this subsection to a close, let us sum up two facts. First, no matter what the Hamiltonian perturbation is to be, the Delaunay (in the orbital case) or the modified Andoyer (in the attitude case) variables $`Q,P`$ always remain canonical. They do so simply because they are *a priori* defined to be canonical – see equations (4) and (7 \- 8) above. Second, as we have seen from (15 \- 16), the osculating character of the $`Q,P`$ variables is lost under momentum-dependent perturbations of the Hamiltonian.<sup>8</sup><sup>8</sup>8 It is possible, of course, to choose the other way and preserve osculation at the cost of canonicity. In the orbital case, one should simply set $`\stackrel{\mathbf{}}{𝚽}=\mathrm{\hspace{0.17em}0}`$ in equations (52 - 57) of Efroimsky (2006). In the attitude case, though, this will be a more cumbersome construction, never implemented in the literature hitherto.
### 2.4 From the modified Andoyer elements to the regular ones
So far our description of perturbed spin, (13 \- 16), has merely been a particular case of the general development (5 \- 12). The sole difference was that the role of canonically-conjugated integration constants $`C=(Q,P)`$ in (13 \- 16) should be played not by the Delaunay variables (as in the orbital case) but by some rotational elements – like, for example, the Richelot-Serret variables (see the footnote in subsection 2.2 above) or by the modified Andoyer set $`(l_o,g_o,h;L_o,G,H)`$ consisting of the initial values of the regular Andoyer elements. The developments conventionally used in the theory of Earth rotation, as well as in spacecraft attitude engineering, are almost always set out in terms of the regular Andoyer elements, not in terms of their initial values (the paper by Fukushima & Ishizaki (1994) being a unique exception). Fortunately, all our gadgetry, developed above for the modified Andoyer set, stays applicable for the regular set. To prove this, let us consider the unperturbed parameterisation of the Euler angles $`q_n=(\varphi ,\theta ,\psi )`$ via the regular Andoyer elements $`A_j=(l,g,h;L,G,H)`$:
$`q_n=f_n(A_1(C,t),...,A_6(C,t)),`$ (17)
each element $`A_i`$ being a function of time and of the initial values $`C_j=(l_o,g_o,h;L_o,G,H)`$. When a perturbation gets turned on, the parameterisation (17) stays, while the time evolution of the elements $`A_i`$ changes: beside the standard time-dependence inherent in the free-spin Andoyer elements, the perturbed elements acquire an extra time-dependence through the evolution of their initial values.<sup>9</sup><sup>9</sup>9 This is fully analogous to the transition from the unperturbed mean longitude, $`M(t)=M_o+n\left(tt_o\right),\text{with}M_o,n,t_o=const,`$ to the perturbed one, $`M(t)=M_o(t)+{\displaystyle _{t_o}^t}n(t^{})𝑑t^{},\text{with}t_o=const,`$ in orbital dynamics. Then the time evolution of an Euler angle $`q_n=(\phi ,\theta ,\psi )`$ will be given by a sum of two items: (1) the angle’s unperturbed dependence upon time and time-dependent elements; and (2) an appropriate addition $`\mathrm{\Phi }_n`$ that arises from a perturbation-caused alteration of the elements’ dependence upon the time:
$`\dot{q}_n=\text{g}_n+\mathrm{\Phi }_n.`$ (18)
The unperturbed part is
$`\text{g}_n={\displaystyle \underset{i=1}{\overset{6}{}}}{\displaystyle \frac{f_n}{A_i}}\left({\displaystyle \frac{A_i}{t}}\right)_C,`$ (19)
while the convective term is given by
$`\mathrm{\Phi }_n={\displaystyle \underset{i=1}{\overset{6}{}}}{\displaystyle \underset{j=1}{\overset{6}{}}}\left({\displaystyle \frac{f_n}{A_i}}\right)_t\left({\displaystyle \frac{A_i}{C_j}}\right)_t\dot{C}_j={\displaystyle \underset{j=1}{\overset{6}{}}}\left({\displaystyle \frac{f_n}{C_j}}\right)_t\dot{C}_j`$
$`={\displaystyle \underset{j=1}{\overset{3}{}}}\left({\displaystyle \frac{f_n}{Q_j}}\right)_t\dot{Q}_j+{\displaystyle \underset{j=1}{\overset{3}{}}}\left({\displaystyle \frac{f_n}{P_j}}\right)_t\dot{P}_j={\displaystyle \frac{\mathrm{\Delta }(q,p)}{p_n}},`$ (20)
where the set $`C_j`$ is split into canonical coordinates and momenta like this: $`Q_j=(l_o,g_o,h)`$ and $`P_j=(L_o,G,H)`$. In the case of free spin they obey the Hamilton equations with a vanishing Hamiltonian and, therefore, are all constants. In the case of disturbed spin, their evolution is governed by (7 \- 8), substitution whereof in (20) will once again take us to (16). This means that the non-osculation-caused convective corrections to the velocities stay the same, no matter whether we parameterise the Euler angles through the modified Andoyer elements (variable constants) or through the regular Andoyer elements. This invariance will become obvious if, once again, we consider the analogy with orbital mechanics: on Fig. 1, the correction $`\stackrel{\mathbf{}}{𝚽}`$ is independent of how we choose to parameterise the nonosculating instantaneous ellipse – through the Delaunay set with $`M_o`$ or through the one containing $`M`$.
This consideration yields the following consequences:
(a) Under momentum-dependent perturbations, calculation of the angular velocities via the elements must be performed not through the second equation of (13) but through the second equation of (14), with (16) substituted therein. The convective term given by (16) is nonzero when the perturbation is angular-velocity-dependent. In other words, under such type of perturbations, the canonicity condition imposed upon the Richelot-Serret or the Andoyer elements is incompatible with osculation. An example of such perturbation shows itself in the theory of planetary rotation, when we switch to a coordinate system associated with the orbit plane. Precession of this plane makes the frame noninertial, and the appropriate Lagrangian perturbation depends upon the planet’s angular velocity. The corresponding Hamiltonian perturbation (denoted in the Kinoshita-Souchay theory by $`E`$) comes out momentum-dependent. In this theory the Andoyer elements are introduced in the precessing frame, and since the precession-caused perturbation is momentum-dependent, these elements come out nonosculating. For this reason, their substitution into the undisturbed expressions (2.6) and (6.26 - 6.27) in Kinoshita (1977) will not render the angular velocity relative to the precessing frame wherein the elements were introduced. To furnish the angular velocity relative to that frame, these expressions must be amended with the appropriate convective terms.
(b) The above circumstance, instead of being a flaw of the Kinoshita-Souchay theory, turns out to be its strong point. It can be shown that Kinoshita’s undisturbed expressions for the angular velocity via the elements keep rendering the *inertial* angular velocity, even when the elements defined in a precessing frame are plugged therein. Briefly speaking, we first introduce the Andoyer elements in an unperturbed setting (inertial frame) and write down the expressions, via these elements, for the Euler angles and velocities relative to the inertial frame. Then we introduce a momentum-dependent perturbation, i.e., switch to a precessing frame, and in that frame we introduce the Andoyer elements. Insertion thereof into the unperturbed expressions for the Euler angles and angular velocities gives us the Euler angles relative to the precessing frame and (due to the nonosculating nature of the elements) the angular velocity relative to the *inertial* frame, not to the precessing one.<sup>10</sup><sup>10</sup>10 This mishap is an example of osculation loss. We introduce the elements in a certain frame (the precessing frame of the orbit), plug them into the unperturbed expressions for the Euler angles and for the angular velocities, and here comes the result: while we obtain the correct values of the Euler angles relative to the said frame, we do not get the right values for the angular velocity relative to that frame. (Instead, our formulae return the values of the angular velocity relative to another, inertial, frame.) This happens because the disturbance, associated with a transition to the precessing frame, depends not only upon the Earth’s orientation but also upon its angular velocity. Or, stated alternatively, because the appropriate Hamiltonian variation $`\mathrm{\Delta }`$ depends upon the momenta $`p`$ canonically conjugated to the Euler angles $`q`$: $`\mathrm{\Phi }{\displaystyle \frac{\mathrm{\Delta }}{p}}\mathrm{\hspace{0.17em}0}.`$ A proof of this fact will be presented in Appendix 1.3.
This fact should not be regarded as a disadvantage of the Kinoshita-Souchay theory, because in some situations it is the inertial, not the relative, angular velocity that is measured (Schreiber et al. 2004, Petrov 2007). Under these circumstances, Kinoshita’s expressions for the angular velocity should be employed (as long as they are correctly identified as the formulae for the *inertial* angular velocity).
## 3 The angular velocity relative to the precessing frame
In the theory of Earth rotation, three angular velocities emerge:
$`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}\text{the relative angular velocity,}`$
i.e., the body’s angular velocity relative to a precessing orbital frame;
$`\stackrel{\mathbf{}}{𝝁}\text{the precession rate of the orbital frame with respect to an inertial one;}`$
$`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\text{the inertial angular velocity,}`$
i.e., the body’s angular velocity with respect to the inertial frame.
Evidently, the latter is the sum of the two former ones: $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}=\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}+\stackrel{\mathbf{}}{𝝁}`$.
If some day we develop an experimental technique for measuring the Earth’s angular velocity relative to the precessing plane of its orbit, we shall have to compare the observations with the theoretical predictions for the directional angles of this, *relative*, angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$.
The Kinoshita (1977) theory was created with intention to furnish the precessing-frame-related directional angles<sup>11</sup><sup>11</sup>11 Here and hereafter the term “directional angles” will stand for the longitude of the node and the inclination of the plane perpendicular to the Earth figure. An analogous meaning is understood for the directional angles of the angular velocity. of the Earth figure (formulae (2.3) and (6.24 - 6.25) in Kinoshita’s paper). This theory also provides precessing-frame-related directional angles of the Earth’s angular-velocity vector (formulae (2.6) and (6.26 - 6.27) in *Ibid.*). We prove in Appendix 1.3 below that, contrary to the expectations, the latter expressions render the directional angles not of the relative but of the *inertial* angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$:
$`I_r^{^{(inert)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}ge\mathrm{cos}\left(2l+g\right)\right],`$ (21)
and
$`h_r^{^{(inert)}}=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{sin}ge\mathrm{sin}\left(2l+g\right)\right].`$ (22)
where the angles $`I`$ and $`J`$ are as on Fig. 3, while $`e`$ is introduced as a measure of triaxiality of the rotator:<sup>12</sup><sup>12</sup>12 For the Earth, $`J\mathrm{\hspace{0.17em}10}^6`$ rad, which justifies the common approximation to write all formulae up to the first order in $`J`$ (Kinoshita 1977). The value of the triaxiality parameter is: $`e=\mathrm{\hspace{0.17em}3.3646441}\times \mathrm{\hspace{0.17em}10}^3`$ (Escapa, Getino & Ferrándiz 2002).
$`e{\displaystyle \frac{\left[\left(1/B\right)\left(1/A\right)\right]/2}{\left(1/C\right)\left[\left(1/A\right)+\left(1/B\right)\right]/2}}`$ (23)
$`A,B,C`$ being the principal moments of inertia.<sup>13</sup><sup>13</sup>13 The Earth is assumed to be rigid, and its body axes are chosen to diagonalise its inertia matrix.
This is a very nontrivial and counterintuitive fact. On introducing the Andoyer variables in the precessing frame of orbit, we plug them into the standard expressions for the orientation angles and the angular velocity. Doing so, we naturally expect to obtain the orientation and the spin rate relative to that precessing frame. We indeed get the body orientation relative to that frame, but the rendered angular velocity turns out to be not the one relative to the precessing frame wherein the Andoyer elements were introduced. Instead, the standard formulae give us the angular velocity relative to some other frame, the inertial one (as if we had used the Andoyer variables defined in the inertial frame). This is an interesting (and still underappreciated by mathematicians) internal symmetry instilled into the Andoyer construction: we can go through a continuum of Andoyer sets (each set introduced in a different precessing frame), but their substitution into the standard formulae for the angular velocity will always return the angular velocity relative to the inertial frame.
A proof of this fact begins with a study of the physical meaning of the Andoyer elements introduced in a precessing frame (presented in Appendix A.1.1). Completion of the proof demands a sequence of calculations so laborious that we chose to put them into the Appendix (see Appendices A.1.2 - A.1.3). This entire situation remarkably parallels a similar episode from the theory of Delaunay elements in orbital dynamics (see the end of Appendix A.1.3).
Now, what if we want to know the angular velocity relative to the precessing frame, i.e., $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$? The precessing-frame-related directional angles of this angular velocity will look as
$`I_r^{^{(rel)}}=I_r^{^{(inert)}}+I_{r}^{}{}_{}{}^{^{(\mathrm{\Phi })}},`$ (24)
and
$`h_r^{^{(rel)}}=h_r^{^{(inert)}}+h_r^{^{(\mathrm{\Phi })}},`$ (25)
the extra terms $`I_{r}^{}{}_{}{}^{^{(\mathrm{\Phi })}},h_{r}^{}{}_{}{}^{^{(\mathrm{\Phi })}}`$ emerging because, as explained above, one has to add the convective term $`\mathrm{\Phi }`$ to the unperturbed velocity g , in order to obtain the full velocity $`\dot{q}`$ under disturbance. Here $`q`$ stands for the three Eulerian angles<sup>14</sup><sup>14</sup>14 Be mindful that in the physics and engineering literature the Euler angles are traditionally denoted with $`(\varphi ,\theta ,\psi )`$. In the literature on the Earth rotation, very often the inverse convention, $`(\psi ,\theta ,\varphi )`$, is employed. In the Kinoshita-Souchay theory, these angles are denoted with $`(h,I,\varphi )`$. The angles defining orientation of the Earth’s figure are accompanied with the subscript $`f`$ and are termed as: $`(h_f,I_f,\varphi _f)`$. The directional angles of the Earth’s angular-velocity vector are equipped with subscript $`r`$. For the relative and the inertial angular velocities the angles are denoted with $`(h_r^{^{(rel)}},I_r^{^{(rel)}},\varphi _r^{^{(rel)}})`$ and $`(h_r^{^{(inert)}},I_r^{^{(inert)}},\varphi _r^{^{(inert)}})`$, accordingly. $`q_n=\{h_f,I_f,\varphi _f\}`$ defining the the orientation of the principal axes of the Earth, relative to the precessing frame, so $`\dot{q}_n`$ will signify time derivatives of Euler angles relative to this precessing frame. The convective terms, entering the expressions for $`\dot{q}_n=\{\dot{h}_f,\dot{I}_f,\dot{\varphi }_f\}`$, can be calculated using formula (16) – see the Appendices 3 and 4 below. The ensuing corrections to the Euler angles determining the orientation of the instantaneous spin axis will look as
$`I_r^{^{(\mathrm{\Phi })}}=\dot{\pi }_1{\displaystyle \frac{C}{L}}\mathrm{cos}I\mathrm{sin}(h\mathrm{\Pi }_1)+\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}\left[\mathrm{sin}\pi _1\mathrm{cos}I\mathrm{cos}(h\mathrm{\Pi }_1)+\mathrm{cos}\pi _1\mathrm{sin}I\mathrm{sin}I\right]`$
$`+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2)`$ (26)
and
$`h_r^{^{(\mathrm{\Phi })}}=\dot{\pi }_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{cos}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{sin}\pi _1\mathrm{sin}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega ).`$ (27)
These corrections depend upon two angles that define the orientation of the precessing orbit with respect to an inertial frame – the inclination $`\pi _1`$ and the node $`\mathrm{\Pi }_1`$.
Let us make rough numerical estimates for the case of the rigid Earth. Putting together (21), (24), and (26), we see that in the resulting expression for $`I_r^{^{(rel)}}`$
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}ge\mathrm{cos}\left(2l+g\right)\right]\dot{\pi }_1{\displaystyle \frac{C}{L}}\mathrm{cos}I\mathrm{sin}(h\mathrm{\Pi }_1)`$
$`+\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}\left[\mathrm{sin}\pi _1\mathrm{cos}I\mathrm{cos}(h\mathrm{\Pi }_1)+\mathrm{cos}\pi _1\mathrm{sin}I\mathrm{sin}I\right]+O(J^2)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2)`$ (28)
we have the leading term, $`I`$, and three additions – of order $`J\mathrm{\hspace{0.17em}10}^6`$, of order $`Je\mathrm{\hspace{0.17em}10}^9`$, and of order<sup>15</sup><sup>15</sup>15 This estimate ensues from the trivial observation that $`C/L\omega ^1`$. Regarding the numbers: according to Lieske et al. (1977) and Seidelmann (1992), $`\dot{\pi }_147\mathrm{"}/century`$, while $`\dot{\mathrm{\Pi }}_1870\mathrm{"}/century`$ $`\mathrm{\hspace{0.17em}2.4}\times 10^3deg/yr`$. On the other hand, $`\omega 360deg/day1.3\times 10^5deg/yr`$ whence $`\mathrm{\Phi }/\omega \dot{\pi }_1/\omega \mathrm{\hspace{0.17em}10}^9`$. (We could as well have used the IERS value of $`\omega \mathrm{\hspace{0.17em}7.3}\times \mathrm{\hspace{0.17em}10}^5rad/s\mathrm{\hspace{0.17em}1.3}\times \mathrm{\hspace{0.17em}10}^7deg/century\mathrm{\hspace{0.17em}4.7}\times \mathrm{\hspace{0.17em}10}^{10}\mathrm{"}/century`$.) $`\mathrm{\Phi }/\omega \dot{\pi }_1/\omega \mathrm{\hspace{0.17em}10}^9`$. We see that *the nonosculation-caused convective terms $`\mathrm{\Phi }`$ provide an effect on the spin-axis orientation, which is of the same order as the $`Je`$ term stemming from triaxiality*.<sup>16</sup><sup>16</sup>16 The nutational spectra of these two contributions are, however, quite different (secular vs. periodic). As $`\mathrm{\hspace{0.33em}1}rad\mathrm{\hspace{0.17em}0.2}\times \mathrm{\hspace{0.17em}10}^{\mathrm{\hspace{0.17em}6}}^{\prime \prime }`$ and $`J\mathrm{\hspace{0.17em}10}^6`$, then the $`J`$ term brings into $`I_r`$ a contribution of an arcsecond order, while the $`\mathrm{\Phi }`$ and $`Je`$ terms give corrections of order milliacrseconds. We also see that the terms of order $`J\mathrm{\Phi }/\omega `$ and those of $`(\mathrm{\Phi }/\omega )^2`$ are much less than one percent of a microarcsecond and may be neglected.
Numerical estimates for the expression
$`h_r^{^{(rel)}}=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{sin}ge\mathrm{sin}\left(2l+g\right)\right]`$
$`\dot{\pi }_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{cos}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{sin}\pi _1\mathrm{sin}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega )`$ (29)
will be similar.
Formulae (26 \- 27) constitute the main result of this paper.<sup>17</sup><sup>17</sup>17 Through the medium of equations (69 \- 71) it is also possible to express these corrections via the Euler set, instead of the Andoyer variables. In Appendices 2 - 4 we present their derivation based on formulae (18) and (20). It would be important to note that the resulting corrections acquire the form (25 \- 27) provided the coordinate system co-precessing with the orbit is chosen as in Kinoshita (1977), i.e., by three consecutive Euler rotations $`\widehat{𝐑}_3(\mathrm{\Pi }_1)\widehat{𝐑}_N(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)`$, letter $`Z`$ standing for an inertial axis orthogonal to the ecliptic of epoch, $`N`$ denoting the line of nodes, and $`\mathrm{\hspace{0.17em}3}`$ being a precessing axis perpendicular to the ecliptic of date – see Appendix A.2.2.3. Under an alternative choice of axes within the co-precessing frame, expressions for $`I_r^{^{(rel)}}`$ and $`h_r^{^{(rel)}}`$ will look differently. For example, a transition carried out by only two Euler rotations, $`\widehat{𝐑}_N(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)`$, as in Appendix A.2.2.2, will yield expression (136) instead of (25), and (146) instead of (27).
In principle, (25 \- 27) might as well be derived by purely geometrical means, i.e., from the formula $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}=\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}+\stackrel{\mathbf{}}{𝝁}`$. We however chose the method based on (18) and (20), because this method is fundamental and applicable to *any kind* of momentum-dependent perturbations of the Hamiltonian – for example, to the perturbations caused by deviations from rigidity, as studied by Getino & Ferrándiz (1994). A similar situation will emerge in the (yet to be built) relativistic theory of the Earth rotation.
## 4 Conclusions
In this article we explained that the unperturbed spin states (“Eulerian cones”) play in the attitude dynamics the same role as the unperturbed two-body orbits (“Keplerian conics”) play in the orbital mechanics. Just as the orbital elements parameterising Keplerian conics, the rotational elements parameterising Eulerian cones may be either osculating or nonosculating. If the perturbation depends upon the velocity (in the orbital case) or upon the angular velocity (in the attitude case), the condition of osculation is incompatible with the condition of canonicity. In these situations the standard equations furnish the Delaunay (in the orbital case) or Andoyer (in the attitude case) elements, which are not osculating, – circumstance important when the elements are employed for calculation of the velocity or angular velocity. The functional form of the expression for a velocity or an angular velocity through elements depends upon whether these elements are osculating or not.
A remarkable peculiarity is shared by the Delaunay and Andoyer elements. Suppose the perturbation is caused by a transition to a precessing frame of reference, and the elements are introduced in this noninertial frame. Their substitution into the unperturbed expressions for the Cartesian coordinates (or the Euler angles) will render the right position (or the attitude) relative to the precessing frame wherein these elements were defined. Now, suppose that we impose on our elements the condition of canonicity. Since the frame-precession-caused perturbation is momentum-dependent, the canonicity condition is incompatible with the osculation one. Hence, when our elements are inserted into the unperturbed expressions for the velocity or angular velocity, they will NOT return the velocity with respect to the precessing frame. It turns out, though, that they will render the velocity relative to the inertial frame. While for the orbital case this was proven in Efroimsky & Goldreich (2003, 2004), in the current paper we proved this fact also for the attitude case.
This has ramifications for the Kinoshita-Souchay theory of the Earth rotation. In this theory, the Andoyer elements are defined in a precessing frame of the Earth orbit. In Kinoshita (1977) these elements were *ab initio* canonical – simply because Kinoshita obtained them via a canonical transformation (see section 3 of his work). As demonstrated in sections 2.3 - 2.4 of our paper, the by-default-imposed canonicity condition made the elements nonosculating. Insertion of such elements into the unperturbed equations for the angular velocity (formulae (2.6) and (6.26 - 6.27) in Kinoshita 1977) does not yield the angular velocity relative to the frame wherein the elements were defined (the precessing frame). Rather, the equations will still furnish the angular velocity relative to the inertial frame of reference. This way of osculation loss might be a flaw of the Kinoshita-Souchay theory, had we expected it to render the angular velocity with respect to a precessing frame. In reality, the osculation loss is an advantage of the theory, because the presently available experimental technique (Schreiber et al. 2004, Petrov 2007) provides for the measurement of the angular velocity relative to the inertial frame – the velocity furnished by the Kinoshita-Souchay theory.
In the final section we provide expressions (26 \- 27) for the body-frame-related directional angles of the planet’s angular velocity relative to a frame coprecessing with the planet’s orbit. The method wherewith we calculate these angles is general and applicable to to *any kind* of momentum-dependent perturbations of the Hamiltonian – for example, to the perturbations caused by deviations from rigidity.
Acknowledgments
The authors would like to deeply thank Hiroshi Kinoshita for his extremely valuable consultations on some subtleties of his theory. ME is also grateful to Pini Gurfil, George Kaplan, Jean Souchay, and Jack Wisdom for fruitful and stimulating conversations on the subject. AE’s contribution was partially supported by Spanish Projects AYA2004-07970 and AYA2005-08109.
Appendix 1.
The Andoyer variables introduced in a precessing frame
A 1.1 Formalism
Let us consider an unsupported rigid body whose spin is to be studied in a coordinate system, which itself is precessing relative to some inertial frame. The said system is assumed to precess at a rate $`\stackrel{\mathbf{}}{𝝁}`$ so the kinetic energy of rotation, in the inertial frame, is given by
$`T_{kin}={\displaystyle \frac{1}{2}}\stackrel{\mathbf{}}{𝝎}_{}^{^{(inert)}}{}_{}{}^{^T}𝕀\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}={\displaystyle \frac{1}{2}}\left(\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}+\stackrel{\mathbf{}}{𝝁}\right)^^T𝕀\left(\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}+\stackrel{\mathbf{}}{𝝁}\right)={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=x,y,z}{}}I_i\left(\omega _i^{^{(rel)}}+\mu _i\right)^2`$
(30)
$`={\displaystyle \frac{1}{2}}A\left(\omega _x^{^{(rel)}}+\mu _x\right)^2+{\displaystyle \frac{1}{2}}B\left(\omega _y^{^{(rel)}}+\mu _y\right)^2+{\displaystyle \frac{1}{2}}C\left(\omega _z^{^{(rel)}}+\mu _z\right)^2,`$
where $`I_i(A,B,C)`$ are the principal values of the inertia matrix of the body, $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$ is the inertial angular velocity (i.e., the one with respect to an inertial frame), $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$ is the relative angular velocity (i.e., the one with respect to a precessing coordinate system), while $`\stackrel{\mathbf{}}{𝝁}`$ is the rotation rate of the precessing frame with respect to the inertial frame. In (30), both $`\stackrel{\mathbf{}}{𝝎}`$ ’ and $`\stackrel{\mathbf{}}{𝝁}`$ are resolved into their components along the principal axes $`\widehat{𝐛}_1\widehat{𝐱},\widehat{𝐛}_2\widehat{𝐲},\widehat{𝐛}_3\widehat{𝐳}`$ of the rotating body. Expression (30) is fundamental and stays, no matter whether $`\stackrel{\mathbf{}}{𝝁}`$ depends on the rotator’s orientation, or whether it carries a direct time dependence.
The role of canonical coordinates will be played the Euler angles<sup>18</sup><sup>18</sup>18 We would once again remind that the Euler angles, though normally termed $`(\varphi ,\theta ,\psi )`$, in the astronomical literature are often denoted as $`(\psi ,\theta ,\varphi )`$. In the Kinoshita-Souchay theory notations $`(h,I,\varphi )`$ are employed. The angles defining orientation of the Earth’s figure and of the Earth’s angular-velocity vector are accompanied with the subscripts $`f`$ and $`r`$, correspondingly: $`(h_f,I_f,\varphi _f)`$ and $`(h_r,I_r,\varphi _r)`$. The directional angles of the inertial angular velocity will be denoted with $`(h_r^{^{(inert)}},I_r^{^{(inert)}},\varphi _r^{^{(inert)}})`$. Those of the relative velocity will be called $`(h_r^{^{(rel)}},I_r^{^{(rel)}},\varphi _r^{^{(rel)}})`$.
$`q_n=(h_f,I_f,\varphi _f)`$ (31)
that map the *precessing* coordinate basis into the principal body basis. To compute their conjugate momenta, let us assume that noninertiality of the precessing coordinate system is the only angular-velocity-dependent perturbation. Then the momenta are simply the derivatives of the kinetic energy. With aid of the formulae for the body-frame components of the relative angular velocity,<sup>19</sup><sup>19</sup>19 It should be emphasised, that the components of the angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$ are related to the body axes, but the angular velocity itself is the relative one (i.e., that with respect to the precessing coordinate system). Our formulae (32 \- 34) are analogous to equations (2.4) in Kinoshita (1977). At the initial step of his development, Kinoshita used his equations (2.4) to express the inertial angular velocity (what we call $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$) via the Euler angles introduced in an inertial frame. Then, on having introduced a precessing orbital frame, Kinoshita employed these equations for expressing the angular velocity through the Euler angles introduced in a precessing frame. Kinoshita did not explore whether this operation would furnish the relative angular velocity (what we call $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$) or still the inertial one.
$`\omega _x^{^{(rel)}}=\dot{h_f}\mathrm{sin}I_f\mathrm{sin}\varphi _f+\dot{I_f}\mathrm{cos}\varphi _f,`$ (32)
$`\omega _y^{^{(rel)}}=\dot{h_f}\mathrm{sin}I_f\mathrm{cos}\varphi _f\dot{I_f}\mathrm{sin}\varphi _f,`$ (33)
$`\omega _z^{^{(rel)}}=\dot{h_f}\mathrm{cos}I_f+\dot{\varphi _f},`$ (34)
we obtain:
$`p_{_{h_f}}={\displaystyle \frac{T_{kin}}{\dot{h}_f}}=A\left(\omega _x^{^{(rel)}}+\mu _x\right)\mathrm{sin}I_f\mathrm{sin}\varphi _f+B\left(\omega _y^{^{(rel)}}+\mu _y\right)\mathrm{sin}I_f\mathrm{cos}\varphi _f+C\left(\omega _z^{^{(rel)}}+\mu _z\right)\mathrm{cos}I_f,`$ (35)
$`p_{_{I_f}}={\displaystyle \frac{T_{kin}}{\dot{I}_f}}=A\left(\omega _x^{^{(rel)}}+\mu _x\right)\mathrm{cos}\varphi _fB\left(\omega _y^{^{(rel)}}+\mu _y\right)\mathrm{sin}\varphi _f.`$ (36)
$`p_{_{\varphi _f}}={\displaystyle \frac{T_{kin}}{\dot{\varphi }_f}}=C\left(\omega _z^{^{(rel)}}+\mu _z\right),`$ (37)
These formulae enable one to express the angular-velocity components $`\omega _i`$ and the derivatives $`\dot{q}_n=(\dot{h}_f,\dot{I}_f,\dot{\varphi }_f)`$ via the momenta $`p_n=(p_{_{h_f}},p_{_{I_f}},p_{_{\varphi _f}})`$. Insertion of (35 \- 37) into
$`={\displaystyle \underset{n=1}{\overset{3}{}}}\dot{q}_np_n=\dot{h}_fp_{_{h_f}}+\dot{I}_fp_{_{I_f}}+\dot{\varphi }_fp_{_{\varphi _f}}T+V(h_f,I_f,\varphi _f;t)`$ (38)
results, after some tedious algebra, in
$`=T+\mathrm{\Delta },`$ (39)
the perturbation $`\mathrm{\Delta }`$ consisting of a potential term $`V`$ (presumed to depend only upon the time and the angular coordinates, not upon the momenta) and a precession-generated inertial term $`E`$:
$`\mathrm{\Delta }=V(h_f,I_f,\varphi _f;t)+E,`$
where
$`E=\mu _x\left[{\displaystyle \frac{\mathrm{sin}\varphi _f}{\mathrm{sin}I_f}}\left(p_{_{h_f}}p_{_{\varphi _f}}\mathrm{cos}I_f\right)+p_{_{I_f}}\mathrm{cos}\varphi _f\right]`$
(40)
$`\mu _y\left[{\displaystyle \frac{\mathrm{cos}\varphi _f}{\mathrm{sin}I_f}}\left(p_{_{h_f}}p_{_{\varphi _f}}\mathrm{cos}I_f\right)p_{_{I_f}}\mathrm{sin}\varphi _f\right]\mu _zp_{_{\varphi _f}},`$
expression equivalent to formulae (24 - 25) in Giacaglia & Jefferys (1971).<sup>20</sup><sup>20</sup>20 It can be shown (Gurfil, Elipe, Tangren & Efroimsky 2007) that the body-frame-related components $`g_i`$ of the angular momentum are connected with the Euler angles and their conjugate momenta through $`g_1=p_{_{h_f}}{\displaystyle \frac{\mathrm{sin}\varphi }{\mathrm{sin}I}}+p_{_{I_f}}\mathrm{cos}\varphi p_{_{\varphi _f}}\mathrm{sin}\varphi \mathrm{cot}I,g_2=p_{_{h_f}}{\displaystyle \frac{\mathrm{cos}\varphi }{\mathrm{sin}I}}p_{_{I_f}}\mathrm{sin}\varphi p_{_{\varphi _f}}\mathrm{cos}\varphi \mathrm{cot}I,g_3=p_{_{\varphi _f}},`$ whence it can be seen that (40) is merely another form of the relation $`\mathrm{\Delta }=\stackrel{\mathbf{}}{𝝁}\stackrel{}{\text{G}}`$. The latter can also be expressed via the Andoyer variables: $`\mathrm{\Delta }=\stackrel{\mathbf{}}{𝝁}\stackrel{}{\text{G}}=\mu _1\sqrt{G^2L^2}\mathrm{sin}l\mu _2\sqrt{G^2L^2}\mathrm{cos}l\mu _3L.`$ Now let us employ the machinery set out in subsection 2.2. The fact that the Andoyer elements are introduced in a noninertial frame is accounted for by the emergence of the $`\mu `$-terms in the expression (38) for the disturbance $`\mathrm{\Delta }`$. Insertion of (40) into (20) entails:
$`\dot{q}_n=\text{g}_n+{\displaystyle \frac{\mathrm{\Delta }}{p_n}}`$ (41)
where $`q_nh_f,I_f,\varphi _f`$, and the convective terms are given by
$`{\displaystyle \frac{\mathrm{\Delta }}{p_{_{h_f}}}}={\displaystyle \frac{\mu _x\mathrm{sin}\varphi _f+\mu _y\mathrm{cos}\varphi _f}{\mathrm{sin}I_f}},`$ (42)
$`{\displaystyle \frac{\mathrm{\Delta }}{p_{_{I_f}}}}=\mu _x\mathrm{cos}\varphi _f+\mu _y\mathrm{sin}\varphi _f,`$ (43)
$`{\displaystyle \frac{\mathrm{\Delta }}{p_{_{\varphi _f}}}}=\left(\mu _x\mathrm{sin}\varphi _f+\mu _y\mathrm{cos}\varphi _f\right)\mathrm{cot}I_f\mu _z,`$ (44)
where $`\dot{q}_n`$ stand for $`\dot{h}_f,\dot{I}_f,\dot{\varphi }_f`$, and $`p_n`$ signify the corresponding momenta, while $`\mu _x,\mu _y,\mu _z`$ are the components of $`\stackrel{\mathbf{}}{𝝁}`$ in the principal axes of the body.
A 1.2 The physical interpretation of the Andoyer variables
defined in a precessing frame
The physical content of the Andoyer construction built in an inertial frame is transparent: see Fig. 3 and explanation thereto. Will all the Andoyer variables and the auxiliary angles $`I`$ and $`J`$ retain the same physical meaning if we re-introduce the Andoyer construction in a noninertial frame? The answer is affirmative, because a transition to a noninertial frame is no different from any other perturbation: precession of the fiducial frame $`(\widehat{𝐬}_1,\widehat{𝐬}_2,\widehat{𝐬}_3)`$ is equivalent to emergence of an extra perturbing torque, one generated by the inertial forces (i.e., by the fictitious forces emerging in the noninertial frame of references). In the original Andoyer construction assembled in an inertial space, the invariable plane was orthogonal to the instantaneous direction of the angular-momentum vector: if the perturbing torques were to instantaneously vanish, the angular-momentum vector (and the invariable plane orthogonal thereto) would freeze in their positions relative to the fiducial axes $`(\widehat{𝐬}_1,\widehat{𝐬}_2,\widehat{𝐬}_3)`$ (which were inertial and therefore indifferent to vanishing of the perturbation). Now, that the Andoyer construction is built in a precessing frame, the fiducial plane is no longer inertial. Nevertheless if the inertial torques were to instantaneously vanish, then the invariable plane would still freeze relative to the fiducial plane (because the fiducial plane would seize its precession). Therefore, all the variables retain their initial meaning. In particular, the variables $`I`$ and $`J`$ defined as above will be the angles that the angular-momentum makes, correspondingly, with the precessing $`\widehat{𝐬}_3`$ space axis and with the $`\widehat{𝐛}_3`$ principal axis of the body.<sup>21</sup><sup>21</sup>21 On the interrelation between the Andoyer variables, referred to an inertial frame, and those referred to a moving frame see equation (3.3) in Kinoshita (1977). Among other things, this explains why Laskar & Robutel (1993) and Touma & Wisdom (1993, 1994), who explored the history of the Martian obliquity, arrived to very close results. Both groups rightly used the angle $`I`$ as an approximation for the obliquity. While Touma & Wisdom (1993, 1994) employed (a somewhat simplified version of) the Kinoshita formalism in an inertial frame, Laskar & Robutel (1993) used this machinery in a precessing frame of the orbit. Now we understand why they obtained so close results, with minor differences stemming, most likely, from averaging-caused error accumulation in the latter paper. (The computation by Touma and Wisdom was based on unaveraged equations of motion, while Laskar and Robutel employed orbit-averaged equations.)
A 1.3 Calculation of the angular velocities via the Andoyer
variables introduced in a precessing frame of reference
Let us now have a look at the well-known expressions
$`\omega _x^{^{(rel)}}=\dot{h}_f\mathrm{sin}I_f\mathrm{sin}\varphi _f+\dot{I}_f\mathrm{cos}\varphi _f,`$ (45)
$`\omega _y^{^{(rel)}}=\dot{h}_f\mathrm{cos}\varphi _f\mathrm{sin}I_f\dot{I}_f\mathrm{sin}\varphi _f,`$ (46)
$`\omega _z^{^{(rel)}}=\dot{\varphi }_f+\dot{h}_f\mathrm{cos}I_f,`$ (47)
for the principal-axes components of the precessing-frame-related angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$. These formulae render this angular velocity as a function of the rates of Euler angle’s evolution, so one can symbolically denote the functional dependence (45 \- 47) as $`\stackrel{\mathbf{}}{𝝎}=\stackrel{\mathbf{}}{𝝎}(\dot{q})`$. This dependence is linear, so insertion of (41) therein will yield:
$`\stackrel{\mathbf{}}{𝝎}\left(\dot{q}(A)\right)=\stackrel{\mathbf{}}{𝝎}(\text{g}(A))+\stackrel{\mathbf{}}{𝝎}\left(\mathrm{\Delta }/p\right),`$ (48)
with $`A`$ denoting the set of Andoyer variables, and $`p`$ signifying the canonical momenta corresponding to the Euler angles. Direct substitution of (42 \- 44) into (45 \- 47) will then show that the second term on the right-hand side in (48) is exactly $`\stackrel{\mathbf{}}{𝝁}`$:
$`\stackrel{\mathbf{}}{𝝎}(\dot{q}(A))=\stackrel{\mathbf{}}{𝝎}(\text{g(A)})\stackrel{\mathbf{}}{𝝁}.`$ (49)
Since the $`\stackrel{\mathbf{}}{𝝎}(\dot{q})`$ is, *ab initio*, the relative angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$ (i.e., that of the body frame relative to the precessing frame), and since $`\stackrel{\mathbf{}}{𝝁}`$ is the precession rate of that frame with respect to the inertial one, then $`\stackrel{\mathbf{}}{𝝎}(\text{g}(A))`$ will always return the inertial angular velocity of the body, $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$ (i.e., the angular velocity relative to the inertial frame). It will do so even despite the fact that now the Andoyer parameterisation is introduced in a precessing coordinate frame!
In brief, the above line of reasoning may be summarised as:
$`\begin{array}{ccc}\stackrel{\mathbf{}}{𝝎}\left(\dot{q}(A)\right)=\stackrel{\mathbf{}}{𝝎}(\text{g}(A))+\stackrel{\mathbf{}}{𝝎}\left(\mathrm{\Delta }/p\right),& & \\ & & \\ \stackrel{\mathbf{}}{𝝎}(\dot{q}(A))=\stackrel{\mathbf{}}{𝝎}^{^{(rel)}},& & \\ & & \\ \stackrel{\mathbf{}}{𝝎}(\mathrm{\Delta }/p)=\stackrel{\mathbf{}}{𝝁},& & \\ & & \\ \stackrel{\mathbf{}}{𝝎}^{^{(rel)}}=\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\stackrel{\mathbf{}}{𝝁}.& & \end{array}\}\stackrel{\mathbf{}}{𝝎}(\text{g}(A))=\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$ (57)
where the entities are defined as follows:
$`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}\text{the relative angular velocity,}`$
i.e., the body’s angular velocity relative to a precessing orbital frame;
$`\stackrel{\mathbf{}}{𝝁}\text{the precession rate of that frame with respect to an inertial one;}`$
$`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\text{the inertial angular velocity,}`$
i.e., the body’s angular velocity with respect to the inertial frame.
This development parallels a situation in orbital dynamics. There the role of canonical elements is played by the Delaunay set $`C=(Q;P)=(L,G,H;M_o,\omega ,\mathrm{\Omega })`$ with
$`L\mu ^{1/2}a^{1/2},G\mu ^{1/2}a^{1/2}\left(1e^2\right)^{1/2},H\mu ^{1/2}a^{1/2}\left(1e^2\right)^{1/2}\mathrm{cos}i,`$
the parameters $`a,e,i,\omega ,\mathrm{\Omega },M_o`$ being the Kepler orbital elements. In the unperturbed setting (the two-body problem in inertial axes), the Cartesian coordinates $`\stackrel{\mathbf{}}{𝒓}(x_1,x_2,x_3)`$ and velocities $`(\dot{x}_1,\dot{x}_2,\dot{x}_3)`$ are expressed via the time and the Delaunay constants by means of the following functional dependencies:
$`\stackrel{\mathbf{}}{𝒓}=\stackrel{\mathbf{}}{𝒇}(C,t)\text{and}\stackrel{}{\text{v}}=\stackrel{}{\text{g}}(C,t),\text{where}\stackrel{}{\text{g}}\stackrel{\mathbf{}}{𝒇}/t.`$ (58)
If we want to describe a satellite orbiting a precessing oblate planet, we may fix our reference frame on the precessing equator of date. Then the two-body problem will get amended with two disturbances. One, $`\mathrm{\Delta }_{oblate}`$, caused by the presence of the equatorial bulge of the planet, will depend only upon the satellite’s position. Another one, $`\mathrm{\Delta }_{precess}`$, will stem from the noninertial nature of our frame and, thus, will give birth to velocity-dependent inertial forces. Under these perturbations, the Delaunay constants (now introduced in the precessing frame) will become canonical variables evolving in time. As explained in subsection 2.3, the velocity-dependence of one of the perturbations involved will make the Delaunay variables nonosculating (provided that we keep them canonical). On the one hand, the expression $`\stackrel{\mathbf{}}{𝒓}=\stackrel{\mathbf{}}{𝒇}(C(t),t)`$ will return the correct Cartesian coordinates of the satellite in the precessing equatorial frame, i.e., in the frame wherein the Delaunay variables were introduced. On the other hand, the expression $`\stackrel{}{\text{g}}(C,t)`$ will no longer return the correct velocities in that frame. Indeed, according to (9 \- 10), the Cartesian components of the velocity in the precessing equatorial frame will be given by $`\stackrel{}{\text{g}}(C,t)+\mathrm{\Delta }_{precess}/\stackrel{}{\text{p}}`$. However, since the second term of this sum is equal to $`\stackrel{\mathbf{}}{𝝁}\times \stackrel{\mathbf{}}{𝒓}`$, then $`\stackrel{}{\text{g}}(C,t)`$ turns out to always render the velocity with respect to the inertial frame of reference (Efroimsky & Goldreich 2004, Efroimsky 2005).
Appendix 2.
The instantaneous angular velocity
in the Kinoshita-Souchay theory
The main burden of subsection 2.3 in the text above was to highlight the need to add the convective term $`\mathrm{\Phi }`$ to the unperturbed velocity g , in order to obtain the full velocity $`\dot{q}`$ under disturbance. Here $`q`$ stands for a vector consisting of the three Eulerian angles $`q_n=h_f,I_f,\varphi _f`$ defining the orientation of the principal axes of the Earth relative to the precessing frame. The corresponding convective terms, entering the expressions for $`\dot{q}_n=\dot{h}_f,\dot{I}_f,\dot{\varphi }_f`$, are given by formula (20). Our eventual goal will be to calculate the corresponding corrections to the Euler angles determining the instantaneous axis of rotation in a precessing frame of reference.
A 2.1 The unperturbed velocities
In this subsection we shall write the unperturbed Euler angles’ partial time derivatives $`\text{g}_nq_n/t`$ as functions of these angles and of the Andoyer variables.
In the Kinoshita-Souchay theory of Earth rotation, the Euler angles defining the figure of the Earth are denoted with
$`q_n=(h_f,I_f,\varphi _f),`$ (59)
the subscript standing for “figure.”
Now, let us denote the principal body axes with $`\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}3}`$ and the appropriate moments of inertia with $`A,B,C`$ (so that $`ABC`$). The angular momentum $`\stackrel{}{𝐋}`$ is connected with the Earth-figure Euler angles via the body-frame components (45 \- 47) of the inertial-frame-related<sup>22</sup><sup>22</sup>22 Be mindful that in formulae (45 \- 47) the notations $`h_f,I_f,\varphi _f`$ stood for the Euler angles defining the body’s orientation relative to a precessing frame. For this reason, (45 \- 47) furnished the relative angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}`$. In this subsection we are beginning with the unperturbed situation, when the orbit frame is yet assumed to be inertial. Hence, at this moment, $`h_f,I_f,\varphi _f`$ yet denote the angles relative to the inertial frame, and hence the same formulae render the inertial angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$. angular velocity $`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}`$:
$`L_x=A\omega _x^{^{(inert)}}=A\left(\dot{h_f}\mathrm{sin}I_f\mathrm{sin}\varphi _f+\dot{I_f}\mathrm{cos}\varphi _f\right),`$ (60)
$`L_y=B\omega _y^{^{(inert)}}=B\left(\dot{h_f}\mathrm{sin}I_f\mathrm{cos}\varphi _f\dot{I_f}\mathrm{sin}\varphi _f\right),`$ (61)
$`L_z=C\omega _z^{^{(inert)}}=C\left(\dot{h_f}\mathrm{cos}I_f+\dot{\varphi _f}\right).`$ (62)
On the other hand, the body-frame components of the angular momentum will be related to the Andoyer elements through<sup>23</sup><sup>23</sup>23 At this point, we are discussing the unperturbed case, with no frame precession. However, as explained in subsection A1.2, the interconnection (63 \- 65) between the Andoyer elements and the components of $`\stackrel{}{𝐋}`$ will stay valid also when the precession is “turned on” (and both the elements and the components of $`\stackrel{}{𝐋}`$ are introduced in a precessing frame of reference).
$`L_x=\sqrt{G^2L^2}\mathrm{sin}l,`$ (63)
$`L_y=\sqrt{G^2L^2}\mathrm{cos}l,`$ (64)
$`L_zL.`$ (65)
Substituting (63 \- 65) into (60 \- 62) and solving for the rates of change of the Euler angles will entail:
$`{\displaystyle \frac{h_f}{t}}={\displaystyle \frac{1}{\mathrm{sin}I_f}}\left[{\displaystyle \frac{L_x}{A}}\mathrm{sin}\varphi _f+{\displaystyle \frac{L_y}{B}}\mathrm{cos}\varphi _f\right]={\displaystyle \frac{1}{\mathrm{sin}I_f}}\sqrt{G^2L^2}\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right],`$ (66)
$`{\displaystyle \frac{I_f}{t}}={\displaystyle \frac{L_x}{A}}\mathrm{cos}\varphi _f{\displaystyle \frac{L_y}{B}}\mathrm{sin}\varphi _f=\sqrt{G^2L^2}\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{cos}\varphi _f}{A}}{\displaystyle \frac{\mathrm{cos}l\mathrm{sin}\varphi _f}{B}}\right],`$ (67)
$`{\displaystyle \frac{\varphi _f}{t}}={\displaystyle \frac{L_z}{C}}\mathrm{cot}I_f\left[{\displaystyle \frac{L_x}{A}}\mathrm{sin}\varphi _f+{\displaystyle \frac{L_y}{B}}\mathrm{cos}\varphi _f\right]`$
$`={\displaystyle \frac{L}{C}}\sqrt{G^2L^2}\mathrm{cot}I_f\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right],`$ (68)
where we deliberately replaced $`\dot{h_f},\dot{I_f},\dot{\varphi _f}`$ with $`h_f/t,I_f/t,\varphi _f/t`$, because so far we have been considering the situation of no disturbances turned on (i.e., the case when the full derivatives coincide with the partial ones, and lack convective terms). Our next step will be to turn on the disturbance $`\mathrm{\Delta }`$, which will include a transition from an inertial frame to the precessing frame of the Earth’s orbit. In accordance with formulae (14) and (16), this transition will generate additions to the derivatives (66 \- 68), the additions that make the difference between a total and a partial derivative.
A 2.2 Turning on the perturbation – switching to a precessing frame
Our goal here is to derive the convective terms $`\mathrm{\Phi }_n=(\mathrm{\Phi }_{_{h_f}},\mathrm{\Phi }_{_{I_f}},\mathrm{\Phi }_{_{\varphi _f}})`$ that are to be added to the partial derivatives (66 \- 68), to get the full time derivatives $`\dot{q}_n=(\dot{h}_f,\dot{I}_f,\dot{\varphi }_f)`$.
A 2.2.1 Generalities
As explained by Kinoshita (1977), the undisturbed dependence of the Euler angles of the Earth’s figure upon the Andoyer elements can be approximated with
$`h_f=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\mathrm{sin}g+O(J^2),`$ (69)
$`I_f=I+J\mathrm{cos}g+O(J^2),`$ (70)
$`\varphi _f=l+gJ\mathrm{cot}I\mathrm{sin}g+O(J^2),`$ (71)
$`J`$ and $`I`$ being the angles that the invariable plane (the one orthogonal to the angular momentum $`\stackrel{}{𝐆}`$) makes with the body equator and with the ecliptic plane of date, correspondingly. (For the Earth, $`J`$ is of order $`10^6.`$) As evident from Fig. 3, these angles are interconnected with the Andoyer variables $`L`$ and $`G`$ through formulae
$`L=G\mathrm{cos}J`$ (72)
and
$`H=G\mathrm{cos}I.`$ (73)
Under perturbations, formulae (69 \- 71) will stay valid. However, the expressions for the angles’ evolution rate, (66 \- 68), will acquire convective additions (20) caused by the loss of osculation. These additions, entering the expressions for $`\dot{q}_n=(\dot{h}_f,\dot{I}_f,\dot{\varphi }_f)`$, will read, accordingly, as
$`\mathrm{\Phi }_{h_f}={\displaystyle \frac{\mathrm{\Delta }}{p_{h_f}}},\mathrm{\Phi }_{I_f}={\displaystyle \frac{\mathrm{\Delta }}{p_{I_f}}},\mathrm{\Phi }_{\varphi _f}={\displaystyle \frac{\mathrm{\Delta }}{p_{\varphi _f}}}.`$ (74)
So our next step will be to calculate these three terms.
Among the perturbations entering the Kinoshita theory, there is a so-called “E term.” It emerges due to a transition from an inertial frame to a noninertial one, i.e., from a coordinate system associated with the ecliptic of epoch to the one associated with the ecliptic of date. Simply speaking, in the Kinoshita theory the Earth rotation is considered in a noninertial frame of the terrestrial orbit precessing about the Sun. In Kinoshita (1977), the $`xy`$ plane of this noninertial frame is referred to as the *moving plane*. In his theory, this “E term” is the only one dependent not only upon the instantaneous orientation but also upon the angular velocity of the Earth (or, in the Hamiltonian formulation, upon the momenta conjugate to the Euler angles of the Earth’s figure). Hence, in this situation $`\mathrm{\Delta }/p_j=E/p_j`$. The expressions for $`\mathrm{\Delta }`$ and the “E term” are rendered by formulae (39 \- 40) where $`p_{_{h_f}},p_{_{I_f}},p_{_{\varphi _f}}`$ denote the canonical momenta, while $`\mu _x`$, $`\mu _y`$, $`\mu _z`$ signify the body-frame components of the angular rate at which the orbit plane is precessing relative to an inertial coordinate system.<sup>24</sup><sup>24</sup>24 We would point out that in Kinoshita’s theory the origin of both $`\mathrm{\Pi }_1`$ and $`h`$ is the mean equinox, whereas in our formalism the origin simply coincides with the $`x`$ axis. For our $`\mathrm{\Pi }_1`$ and $`h`$ to coincide with those of Kinoshita, not only must we choose our inertial coordinate system with its $`xy`$ plane being within the ecliptic of epoch, but we should also choose the $`x`$ axis to coincide with the mean equinox of epoch. Similarly, not only should our precessing frame to be associated with the ecliptic of date, but the precessing $`x`$ axis whence we reckon the angles should be placed exactly at the angular distance of $`\mathrm{\Pi }_1`$ from the node – see Fig. 2 in Kinoshita (1977). (Mind that in the presence of precession Kinoshita employs notation $`h^{}`$ instead of $`h`$.)
In order to continue, we need the expressions for the body-frame components $`\mu _x,\mu _y,\mu _z`$. These can be obtained from the precessing-frame components $`\mu _1,\mu _2,\mu _3`$ by means of the appropriate rotation matrix:
$`\left[\begin{array}{ccc}\mu _x& & \\ \mu _y& & \\ \mu _z& & \end{array}\right]=`$ (78)
$`\left[\begin{array}{ccc}\mathrm{cos}\varphi _f\mathrm{cos}h_f\mathrm{sin}\varphi _f\mathrm{cos}I_f\mathrm{sin}h_f& \mathrm{cos}\varphi _f\mathrm{sin}h_f+\mathrm{sin}\varphi _f\mathrm{cos}I_f\mathrm{cos}h_f& \mathrm{sin}\varphi _f\mathrm{sin}I_f\\ \mathrm{sin}\varphi _f\mathrm{cos}h_f\mathrm{cos}\varphi _f\mathrm{cos}I_f\mathrm{sin}h_f& \mathrm{sin}\varphi _f\mathrm{sin}h_f+\mathrm{cos}\varphi _f\mathrm{cos}I_f\mathrm{cos}h_f& \mathrm{cos}\varphi _f\mathrm{sin}I_f\\ \mathrm{sin}I_f\mathrm{sin}h_f& \mathrm{sin}I_f\mathrm{cos}h_f& \mathrm{cos}I_f\end{array}\right]\left[\begin{array}{ccc}\mu _1& & \\ \mu _2& & \\ \mu _3& & \end{array}\right]`$ (85)
Since no other contributions in $`\mathrm{\Delta }`$ other than $`E`$ depend upon the momenta, then
$`\mathrm{\Phi }_{h_f}={\displaystyle \frac{\mathrm{\Delta }}{p_{_{h_f}}}}={\displaystyle \frac{E}{p_{_{h_f}}}}={\displaystyle \frac{\mathrm{sin}\varphi _f}{\mathrm{sin}I_f}}\mu _x{\displaystyle \frac{\mathrm{cos}\varphi _f}{\mathrm{sin}I_f}}\mu _y=\mu _1\mathrm{cot}I_f\mathrm{sin}h_f\mu _2\mathrm{cot}I_f\mathrm{cos}h_f\mu _3,`$ (86)
$`\mathrm{\Phi }_{I_f}={\displaystyle \frac{\mathrm{\Delta }}{p_{_{I_f}}}}={\displaystyle \frac{E}{p_{_{I_f}}}}=\mu _x\mathrm{cos}\varphi _f+\mu _y\mathrm{sin}\varphi _f=\mu _1\mathrm{cos}h_f\mu _2\mathrm{sin}h_f,`$ (87)
$`\mathrm{\Phi }_{\varphi _f}={\displaystyle \frac{\mathrm{\Delta }}{p_{_{\varphi _f}}}}={\displaystyle \frac{E}{p_{_{\varphi _f}}}}={\displaystyle \frac{\mathrm{sin}\varphi _f\mathrm{cos}I_f}{\mathrm{sin}I_f}}\mu _x+{\displaystyle \frac{\mathrm{cos}\varphi _f\mathrm{cos}I_f}{\mathrm{sin}I_f}}\mu _y\mu _z={\displaystyle \frac{\mathrm{sin}h_f}{\mathrm{sin}I_f}}\mu _1+{\displaystyle \frac{\mathrm{cos}h_f}{\mathrm{sin}I_f}}\mu _2.`$ (88)
Naturally, none of the $`\mathrm{\Phi }`$ terms bears dependence upon $`\varphi __f`$.
To write down the $`\mathrm{\Phi }`$ terms as functions of the longitude $`\mathrm{\Pi }_1`$ and inclination $`\pi _1`$ of the ecliptic of date on that of epoch, we shall insert into (86 \- 88) the appropriate expressions for $`\mu _1,\mu _2,\mu _3`$. However, at this point care is needed, because of the freedom of choice of a coordinate system co-precessing with the orbital plane.<sup>25</sup><sup>25</sup>25 We are grateful to Hiroshi Kinoshita who explained to us the choice accepted in his works.
A 2.2.2 The precession rate $`\stackrel{\mathbf{}}{\mu }`$ as seen in a certain coordinate system
associated with the precessing equator of date
Let the inertial axes $`(X,Y,Z)`$ be fixed in space so that $`X`$ and $`Y`$ belong to the ecliptic of epoch. A rotation within the ecliptic-of-epoch plane by longitude $`\mathrm{\Pi }_1`$, from the axis
$`X`$, will define the line of nodes. A rotation about this line by an inclination angle $`\pi _1`$ will give us the ecliptic of date. The line of nodes, $`\mathrm{\hspace{0.33em}1}`$, along with axis $`\mathrm{\hspace{0.33em}2}`$ naturally chosen within the ecliptic-of-date plane, and with axis $`\mathrm{\hspace{0.33em}3}`$ orthogonal to this plane, will constitute the precessing coordinate system, with the appropriate basis denoted by $`(\widehat{𝐞}_1,\widehat{𝐞}_2,\widehat{𝐞}_3)`$. For example, the unit vector $`\widehat{𝐞}_3`$ reads in the inertial axes $`(X,Y,Z)`$ as
$`\widehat{𝐞}_3=(\mathrm{sin}\pi _1\mathrm{sin}\mathrm{\Pi }_1,\mathrm{sin}\pi _1\mathrm{cos}\mathrm{\Pi }_1,\mathrm{cos}\pi _1)^^T.`$ (89)
The Earth’s angular velocity relative to the inertial and precessing axes obey
$`\stackrel{\mathbf{}}{𝝎}^{(inert)}=\stackrel{\mathbf{}}{𝝎}^{(rel)}+\stackrel{\mathbf{}}{𝝁},`$ (90)
$`\stackrel{\mathbf{}}{𝝁}`$ being the precession rate of the precessing axes $`\widehat{𝐞}_j`$ relative to the inertial axes $`(X,Y,Z)`$. In the inertial axes, this rate is given by
$`\stackrel{\mathbf{}}{𝝁}{}_{}{}^{}{}_{}{}^{}=(\dot{\pi }_1\mathrm{cos}\mathrm{\Pi }_1,\dot{\pi }_1\mathrm{sin}\mathrm{\Pi }_1,\dot{\mathrm{\Pi }}_1)^^T,`$ (91)
because this expression satisfies the equality $`\stackrel{\mathbf{}}{𝝁}{}_{}{}^{}{}_{}{}^{}\times \widehat{𝐞}_3=\dot{\widehat{𝐞}}_3`$, as can be easily seen from (89) and (91).
In a frame precessing with the ecliptic, the precession rate will be represented by the vector
$`\stackrel{\mathbf{}}{𝝁}=\widehat{𝐑}_{ed}\stackrel{\mathbf{}}{𝝁}{}_{}{}^{}{}_{}{}^{},`$ (92)
where
$`\widehat{𝐑}_{ed}=\widehat{𝐑}_1(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)=\left[\begin{array}{ccc}\mathrm{cos}\mathrm{\Pi }_1& \mathrm{sin}\mathrm{\Pi }_1& 0\\ & & \\ \mathrm{cos}\pi _1\mathrm{sin}\mathrm{\Pi }_1& \mathrm{cos}\pi _1\mathrm{cos}\mathrm{\Pi }_1& \mathrm{sin}\pi _1\\ & & \\ \mathrm{sin}\pi _1\mathrm{sin}\mathrm{\Pi }_1& \mathrm{sin}\pi _1\mathrm{cos}\mathrm{\Pi }_1& \mathrm{cos}\pi _1\end{array}\right]`$ (98)
is the matrix of rotation from the ecliptic of epoch to that of date. From (92 \- 98) we get the components of the precession rate,<sup>26</sup><sup>26</sup>26 Equivalently, one can find the components of $`\stackrel{\mathbf{}}{𝝁}`$ as the elements of the skew-symmetric matrix $`\dot{\widehat{𝐑}}_{ed}\widehat{𝐑}_{ed}^^1`$. as seen in the co-precessing coordinate frame $`(1,\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}3})`$:
$`\stackrel{\mathbf{}}{𝝁}=(\mu _1,\mu _2,\mu _3)^^T=(\dot{\pi }_1,\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1,\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1)^^T.`$ (99)
Substitution of these components into (86 \- 88) entails:
$`\mathrm{\Phi }_{h_f}=\dot{\pi }_1\mathrm{cot}I_f\mathrm{sin}h_f\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{cot}I_f\mathrm{cos}h_f\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1,`$ (100)
$`\mathrm{\Phi }_{I_f}=\dot{\pi }_1\mathrm{cos}h_f\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{sin}h_f,`$ (101)
$`\mathrm{\Phi }_{\varphi _f}={\displaystyle \frac{\mathrm{sin}h_f}{\mathrm{sin}I_f}}\dot{\pi }_1+{\displaystyle \frac{\mathrm{cos}h_f}{\mathrm{sin}I_f}}\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1.`$ (102)
A 2.2.3 The precession rate $`\stackrel{\mathbf{}}{\mu }`$ as seen in a different coordinate system
associated with the precessing equator of date
(the system used by Kinoshita 1977)
In the preceding subsection the transition from the ecliptic of epoch to the one of date was implemented by two Euler rotations: $`\widehat{𝐑}_{ed}=\widehat{𝐑}_N(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)=\widehat{𝐑}_1(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)`$. The axis $`\mathrm{\hspace{0.17em}1}`$ of the precessing frame was assumed to coincide with the line of nodes, $`N`$. Evidently, this choice was just one out of an infinite multitude. An alternative option was employed by Kinoshita (1977), who used a sequence of three Eulerian rotations: $`\widehat{𝐑}_{ed}^K=\widehat{𝐑}_3(\mathrm{\Pi }_1)\widehat{𝐑}_N(\pi _1)\widehat{𝐑}_Z(\mathrm{\Pi }_1)`$. Specifically, having performed the two rotations described above, Kinoshita then rotated the axis $`\mathrm{\hspace{0.17em}1}`$ within the ecliptic of date by angle $`\mathrm{\Pi }_1`$ away from the line of nodes $`N`$. Due to reasoning analogous to what was presented in the subsection above, the sequence of three rotations gives, instead of (99), the following expression:<sup>27</sup><sup>27</sup>27 In the precessing frame, the angular momentum reads: $`(G\mathrm{sin}I\mathrm{sin}h,G\mathrm{sin}I\mathrm{cos}h,H)^^T`$, quantities $`I`$, $`h`$, and $`H`$ being as in Fig. 3. This, together with (103) and the formula $`\mathrm{\Delta }=E=\stackrel{\mathbf{}}{𝝁}\stackrel{}{𝐆}`$, will entail: $`\mathrm{\Delta }=E=\dot{\mathrm{\Pi }}_1H\left(1\mathrm{cos}\pi _1\right)\dot{\pi }_1G\mathrm{sin}I\mathrm{sin}(h\mathrm{\Pi }_1)+\dot{\mathrm{\Pi }}_1G\mathrm{sin}I\mathrm{cos}(h\mathrm{\Pi }_1)\mathrm{sin}\pi _1,`$ which coincides with expression (3.4) in Kinoshita (1977).
$`\stackrel{\mathbf{}}{𝝁}=(\mu _1,\mu _2,\mu _3)^^T`$
(103)
$`=(\dot{\pi }_1\mathrm{cos}\mathrm{\Pi }_1\dot{\mathrm{\Pi }}\mathrm{sin}\mathrm{\Pi }_1\mathrm{sin}\pi _1,\dot{\pi }_1\mathrm{sin}\mathrm{\Pi }_1+\dot{\mathrm{\Pi }}_1\mathrm{cos}\mathrm{\Pi }_1\mathrm{sin}\pi _1,\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1\dot{\mathrm{\Pi }}_1)^^T.`$
Insertion thereof into (86 \- 88) will yield:
$`\mathrm{\Phi }_{h_f}=\dot{\pi }_1\mathrm{cot}I_f\mathrm{sin}(h_f\mathrm{\Pi }_1)\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{cot}I_f\mathrm{cos}(h_f\mathrm{\Pi }_1)\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1+\mathrm{\Pi }_1,`$ (104)
$`\mathrm{\Phi }_{I_f}=\dot{\pi }_1\mathrm{cos}(h_f\mathrm{\Pi }_1)\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{sin}(h_f\mathrm{\Pi }_1),`$ (105)
$`\mathrm{\Phi }_{\varphi _f}={\displaystyle \frac{\mathrm{sin}(h_f\mathrm{\Pi }_1)}{\mathrm{sin}I_f}}\dot{\pi }_1+{\displaystyle \frac{\mathrm{cos}(h_f\mathrm{\Pi }_1)}{\mathrm{sin}I_f}}\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1.`$ (106)
A 2.3 The perturbed velocities
According to (18), to get the full evolution of the figure-axis Euler angles under perturbation, one should sum the unperturbed velocities, given by the partial derivatives (66 \- 68), with the appropriate convective terms (86 \- 88):
$`\dot{h}_f=\left({\displaystyle \frac{h_f}{t}}\right)_C+\mathrm{\Phi }_{_{h_f}}=`$ (107)
$`{\displaystyle \frac{1}{\mathrm{sin}I_f}}\sqrt{G^2L^2}\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right]+\dot{\pi }_1\mathrm{cot}I_f\mathrm{sin}h_f\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{cot}I_f\mathrm{cos}h_f\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1,`$
$`\dot{I}_f=\left({\displaystyle \frac{I_f}{t}}\right)_C+\mathrm{\Phi }_{_{I_f}}`$
$`=\sqrt{G^2L^2}\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{cos}\varphi _f}{A}}{\displaystyle \frac{\mathrm{cos}l\mathrm{sin}\varphi _f}{B}}\right]\dot{\pi }_1\mathrm{cos}h_f\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{sin}h_f,`$ (108)
$`\dot{\varphi }_f=\left({\displaystyle \frac{\varphi _f}{t}}\right)_C+\mathrm{\Phi }_{_{\varphi _f}}`$
$`={\displaystyle \frac{L}{C}}\sqrt{G^2L^2}\mathrm{cot}I_f\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right]{\displaystyle \frac{\mathrm{sin}h_f}{\mathrm{sin}I_f}}\dot{\pi }_1+{\displaystyle \frac{\mathrm{cos}h_f}{\mathrm{sin}I_f}}\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1.`$ (109)
These expressions will help us to determine the instantaneous orientation of the spin axis.
A 2.4 The precessing-frame-related angular velocity
expressed through the Andoyer elements
introduced in the precessing frame.
Let angles $`h_r^{^{(rel)}}`$ and $`I_r^{^{(rel)}}`$ be the precessing-frame-related node and inclination of vector of the angular velocity relative to the precessing frame, and let $`\omega ^{^{(rel)}}|\stackrel{\mathbf{}}{𝝎}^{(rel)}|`$ be the relative spin rate. Our next step will be to derive expressions for $`h_r^{^{(rel)}}`$ and $`I_r^{^{(rel)}}`$ as functions of $`h_f,I_f,\varphi _f,\dot{h}_f,\dot{I}_f,\dot{\varphi }_f`$, and then to substitute the expressions (107 \- 109) for $`\dot{h}_f,\dot{I}_f,\dot{\varphi }_f`$ therein, in order to express $`h_r^{^{(rel)}}`$ and $`I_r^{^{(rel)}}`$ via $`h_f,I_f,\varphi _f`$ only. To accomplish this step, let us begin with the formulae interconnecting the precessing-frame components of the angular velocity relative to the precessing frame with the figure-axis Euler angles and with the spin-axis Euler angles. These formulae are fundamental and perturbation-invariant:
$`\omega _1^{(rel)}=\dot{I}_f\mathrm{cos}h_f+\dot{\varphi }_f\mathrm{sin}I_f\mathrm{sin}h_f`$ (110)
$`\omega _2^{(rel)}=\dot{I}_f\mathrm{sin}h_f\dot{\varphi }_f\mathrm{sin}I_f\mathrm{cos}h_f`$ (111)
$`\omega _3^{(rel)}=\dot{h}_f+\dot{\varphi }_f\mathrm{cos}I_f`$ (112)
and
$`\omega _1^{(rel)}=\omega ^{^{(rel)}}\mathrm{sin}I_r^{^{(rel)}}\mathrm{sin}h_r^{^{(rel)}}`$ (113)
$`\omega _2^{(rel)}=\omega ^{^{(rel)}}\mathrm{sin}I_r^{^{(rel)}}\mathrm{cos}h_r^{^{(rel)}}`$ (114)
$`\omega _3^{(rel)}=\omega ^{^{(rel)}}\mathrm{cos}I_r^{^{(rel)}}.`$ (115)
Both the inertial and relative spin rates, $`\omega ^{^{(inert)}}`$ and $`\omega ^{^{(rel)}}`$, can be most conveniently calculated in the body frame. In that frame, the inertial angular velocity can be written as
$`\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}=({\displaystyle \frac{L_x}{A}},{\displaystyle \frac{L_y}{B}},{\displaystyle \frac{L_z}{C}})^^T=({\displaystyle \frac{G\mathrm{sin}J\mathrm{sin}l}{A}},{\displaystyle \frac{G\mathrm{sin}J\mathrm{cos}l}{B}},{\displaystyle \frac{G\mathrm{cos}J}{C}})^^T,`$ (116)
whence its absolute value turns out to be
$`\omega ^{^{(inert)}}=\sqrt{\left({\displaystyle \frac{L_x}{A}}\right)^2+\left({\displaystyle \frac{L_y}{B}}\right)^2+\left({\displaystyle \frac{L_z}{C}}\right)^2}=`$
$`\sqrt{\left({\displaystyle \frac{G\mathrm{sin}J\mathrm{sin}l}{A}}\right)^2+\left({\displaystyle \frac{G\mathrm{sin}J\mathrm{cos}l}{B}}\right)^2+\left({\displaystyle \frac{G\mathrm{cos}J}{C}}\right)^2}={\displaystyle \frac{G}{C}}\left[1+O\left(J^2\right)\right]={\displaystyle \frac{L}{C}}\left[1+O\left(J^2\right)\right].`$ (117)
To derive the relative rate, square the obvious equality $`\stackrel{\mathbf{}}{𝝎}^{^{(rel)}}=\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\stackrel{\mathbf{}}{𝝁}`$, to obtain:
$`(\stackrel{\mathbf{}}{𝝎}^{^{(rel)}})^^2=(\stackrel{\mathbf{}}{𝝎}^{^{(inert)}})^^2\mathrm{\hspace{0.33em}2}\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\stackrel{\mathbf{}}{𝝁}+\stackrel{\mathbf{}}{𝝁}^^2.`$ (118)
Hence,
$`\omega ^{^{(rel)}}=\omega ^{^{(inert)}}\left[\mathrm{\hspace{0.33em}1}\alpha +O\left(\left(\mathrm{\Phi }/\omega \right)^2\right)\right]={\displaystyle \frac{L}{C}}\left[\mathrm{\hspace{0.33em}1}\alpha +O\left(\left(\mathrm{\Phi }/\omega \right)^2+O(J^2)\right)\right]`$ (119)
and
$`{\displaystyle \frac{1}{\omega ^{^{(rel)}}}}={\displaystyle \frac{1}{\omega ^{^{(inert)}}}}\left[\mathrm{\hspace{0.33em}1}+\alpha +O\left(\left(\mathrm{\Phi }/\omega \right)^2\right)\right]={\displaystyle \frac{C}{L}}\left[\mathrm{\hspace{0.33em}1}+\alpha +O\left(\left(\mathrm{\Phi }/\omega \right)^2\right)+O(J^2)\right],`$ (120)
where
$`\alpha {\displaystyle \frac{\stackrel{\mathbf{}}{𝝎}^{^{(inert)}}\stackrel{\mathbf{}}{𝝁}}{(\stackrel{\mathbf{}}{𝝎}^{^{(inert)}})^^2}}.`$ (121)
is of order $`\mathrm{\Phi }/\omega `$. Dot-multiplying (116) by (78), we arrive at:
$`\alpha ={\displaystyle \frac{G\mathrm{cos}J}{C}}{\displaystyle \frac{\mu _z}{(\stackrel{\mathbf{}}{𝝎}^{^{(inert)}})^^2}}+O(J\mathrm{\Phi }/\omega )`$
$`={\displaystyle \frac{C}{L}}\left[\mu _1\mathrm{sin}I_f\mathrm{sin}h_f\mu _2\mathrm{sin}I_f\mathrm{cos}h_f+\mu _3\mathrm{cos}I_f\right]+O(J\mathrm{\Phi }/\omega )+O(J^2).`$ (122)
In the coordinate system described in A.2.2.2, substitution of (99) makes the latter read
$`\alpha ={\displaystyle \frac{C}{L}}\left(\dot{\pi }_1\mathrm{sin}I_f\mathrm{sin}h_f\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{sin}I_f\mathrm{cos}h_f+\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1\mathrm{cos}I_f\right)+O(J\mathrm{\Phi }/\omega )+O(J^2).`$ (123)
In Kinoshita’s coordinate system described in subsection A.2.2.3, one should use for the components of $`\stackrel{\mathbf{}}{𝝁}`$ not expression (99) but (103), insertion whereof into (122) results in:
$`\alpha ={\displaystyle \frac{C}{L}}(\dot{\pi }_1\mathrm{sin}I_f\mathrm{sin}(h_f\mathrm{\Pi }_1)\dot{\mathrm{\Pi }}_1\mathrm{sin}\pi _1\mathrm{sin}I_f\mathrm{cos}(h_f\mathrm{\Pi }_1)+\dot{\mathrm{\Pi }}_1\mathrm{cos}\pi _1\mathrm{cos}I_f`$
(124)
$`\dot{\mathrm{\Pi }}_1\mathrm{cos}I_f)+O(J\mathrm{\Phi }/\omega )+O(J^2).`$
This formula will enable us to derive approximate (valid up to $`O\left(J^2\right)+O\left(J\mathrm{\Phi }/\omega \right)+O\left((\mathrm{\Phi }/\omega )^2\right)`$) expressions for $`h_r^{^{(rel)}}`$ and $`I_r^{^{(rel)}}`$ expressed as functions of the Andoyer variables.
Appendix 3. Expression for $`I_r^{^{(rel)}}`$
Together, (112) and (115) will give:
$`\omega ^{^{(rel)}}\mathrm{cos}I_r^{^{(rel)}}=\dot{h}_f+\dot{\varphi }_f\mathrm{cos}I_f`$ (125)
or, equivalently,
$`\omega ^{^{(rel)}}\left(\mathrm{cos}I_r^{^{(rel)}}\mathrm{cos}I_f\right)=\dot{h}_f+\left(\dot{\varphi }_f\omega ^{^{(rel)}}\right)\mathrm{cos}I_f.`$ (126)
Since we are planning to carry out all calculations neglecting terms $`O\left(J^2\right)`$, and since the three inclinations $`I_f,I_r^{^{(rel)}},I`$ differ from one another by quantities of order $`O(J)`$, we can approximate the left-hand side of (126) with the first-order terms of its Taylor expansion:
$`\omega ^{^{(rel)}}\left(\mathrm{sin}I_f\right)\left(I_r^{^{(rel)}}I_f\right)=\dot{h}_f+\left(\dot{\varphi }_f\omega ^{^{(rel)}}\right)\mathrm{cos}I_f+O\left(J^2\right),`$ (127)
wherefrom
$`I_r^{^{(rel)}}I_f={\displaystyle \frac{\dot{h}_f}{\omega ^{^{(rel)}}}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}+\left[{\displaystyle \frac{\dot{\varphi }_f}{\omega ^{^{(rel)}}}}\mathrm{\hspace{0.33em}1}\right]\left(\mathrm{cot}I_f\right)+O\left(J^2\right)`$
$`=\left({\displaystyle \frac{h_f}{t}}+\mathrm{\Phi }_{h_f}\right){\displaystyle \frac{1}{\omega ^{^{(rel)}}}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}+\left[\left({\displaystyle \frac{\varphi _f}{t}}+\mathrm{\Phi }_{\varphi _f}\right){\displaystyle \frac{1}{\omega ^{^{(rel)}}}}\mathrm{\hspace{0.33em}1}\right]\left(\mathrm{cot}I_f\right)+O\left(J^2\right).`$ (128)
To get rid of the time derivatives, employ formulae (107 \- 109) or (66 \- 68):
$`I_r^{^{(rel)}}I_f=\mathrm{cot}I_f\left\{\mathrm{\hspace{0.33em}1}{\displaystyle \frac{1}{\omega ^{^{(rel)}}}}\left[{\displaystyle \frac{L}{C}}\sqrt{G^2L^2}\mathrm{cot}I_f\left({\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right)\right]\right\}`$
(129)
$`{\displaystyle \frac{1}{\omega ^{^{(rel)}}\mathrm{sin}^2I_f}}\sqrt{G^2L^2}\left({\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right){\displaystyle \frac{1}{\omega ^{^{(rel)}}\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)+O\left(J^2\right)`$
For $`\varphi _f`$ we can use the approximation $`\varphi _f=l+\text{g}J\mathrm{cot}I\mathrm{sin}\text{g}+O(J^2)`$. Besides, in the terms of order $`J`$ and of order $`\mathrm{\Phi }/\omega `$ we can substitute $`\omega ^{^{(rel)}}`$ with $`\omega ^{^{(inert)}}=L/C+O(J^2)`$. Such substitutions will entail errors of orders $`O(J\mathrm{\Phi }/\omega )`$ and $`O((\mathrm{\Phi }/\omega )^2)`$. However, in the leading term we must use (119 \- 123). Thus we get:
$`I_r^{^{(rel)}}I_f=\mathrm{cot}I_f\left(\mathrm{\hspace{0.17em}1}{\displaystyle \frac{1}{\omega ^{^{(rel)}}}}{\displaystyle \frac{L}{C}}\right)+{\displaystyle \frac{\sqrt{G^2L^2}}{\omega ^{^{(rel)}}}}\left(\mathrm{cot}^2I_f{\displaystyle \frac{1}{\mathrm{sin}^2I_f}}\right)\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{sin}\varphi _f}{A}}+{\displaystyle \frac{\mathrm{cos}l\mathrm{cos}\varphi _f}{B}}\right]`$
$`{\displaystyle \frac{1}{\omega ^{^{(rel)}}}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)+O(J^2)=`$
$`\mathrm{cot}I_f\left(1{\displaystyle \frac{1+\alpha }{\omega ^{^{(inert)}}}}{\displaystyle \frac{L}{C}}\right){\displaystyle \frac{G\mathrm{sin}J}{G/C}}\left[{\displaystyle \frac{\mathrm{sin}l}{A}}\mathrm{sin}\left(l+\text{g}J\mathrm{cot}I\mathrm{sin}\text{g}\right)+{\displaystyle \frac{\mathrm{cos}l}{B}}\mathrm{cos}\left(l+\text{g}J\mathrm{cot}I\mathrm{sin}\text{g}\right)\right]`$
$`{\displaystyle \frac{1}{L/C}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)+O\left(J^2\right)+O\left(J\mathrm{\Phi }/\omega \right)+O((\mathrm{\Phi }/\omega )^2)`$
$`=\alpha \mathrm{cot}I_fCJ\left[{\displaystyle \frac{\mathrm{sin}l}{A}}\mathrm{sin}\left(l+\text{g}\right)+{\displaystyle \frac{\mathrm{cos}l}{B}}\mathrm{cos}\left(l+\text{g}\right)\right]`$
$`{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)+O\left(J^2\right)+O\left(J\mathrm{\Phi }/\omega \right)+O\left((\mathrm{\Phi }/\omega )^2\right),`$ (130)
whence, by using (123) and the formula $`I_f=I+J\mathrm{cos}g+O(J^2)`$, we arrive at:
$`I_r^{^{(rel)}}=I+J\left\{\mathrm{cos}\text{g}{\displaystyle \frac{C}{A}}\mathrm{sin}l\mathrm{sin}\left(l+\text{g}\right){\displaystyle \frac{C}{B}}\mathrm{cos}l\mathrm{cos}\left(l+\text{g}\right)\right\}\alpha \mathrm{cot}I_f`$
$`{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)+O\left(J^2\right)+O\left(J\mathrm{\Phi }/\omega \right)+O\left((\mathrm{\Phi }/\omega )^2\right)`$
Via trigonometric transformations, the second term gets simplified as:
$`\mathrm{cos}\text{g}{\displaystyle \frac{C}{A}}\mathrm{sin}l\mathrm{sin}\left(l+\text{g}\right){\displaystyle \frac{C}{B}}\mathrm{cos}l\mathrm{cos}\left(l+\text{g}\right)=\mathrm{cos}\text{g}{\displaystyle \frac{C}{A}}{\displaystyle \frac{\mathrm{cos}\text{g}\mathrm{cos}(2l+\text{g})}{2}}`$ (131)
$`{\displaystyle \frac{C}{B}}{\displaystyle \frac{\mathrm{cos}\text{g}+\mathrm{cos}(2l+\text{g})}{2}}=\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right].`$
Insertion of (86 \- 88) into the fourth term will make it look:
$`{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left(\mathrm{\Phi }_{h_f}+\mathrm{\Phi }_{\varphi _f}\mathrm{cos}I_f\right)={\displaystyle \frac{C}{L}}{\displaystyle \frac{\mu _3}{\mathrm{sin}I_f}}`$ (132)
Hence, we have for $`I_r^{^{(rel)}}`$:
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right]+{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mu _3}{\mathrm{sin}I_f}}\alpha \mathrm{cot}I_f`$
$`+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O\left((\mathrm{\Phi }/\omega )^2\right)`$ (133)
where the parameter $`e`$, given by (23), is the measure of triaxiality of the rotator.
In the precessing coordinate system obtained from the inertial one by two Euler rotations, as in subsection A.2.2.2, we must now substitute (123) for $`\alpha `$ and (99) for $`\mu _3`$, to get:
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right]{\displaystyle \frac{C}{L}}\dot{\pi }_1\mathrm{cos}I_f\mathrm{sin}h`$
(134)
$`+{\displaystyle \frac{C}{L}}\dot{\mathrm{\Pi }}_1\left(\mathrm{sin}\pi _1\mathrm{cos}I_f\mathrm{cos}h+\mathrm{cos}\pi _1\mathrm{sin}I_f\right)+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2).`$
In the Kinoshita precessing axes obtained from the inertial ones by three rotations, as in subsection A.2.2.3, we should substitute (124) for $`\alpha `$ and (103) for $`\mu _3`$, to get:
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right]{\displaystyle \frac{C}{L}}\dot{\pi }_1\mathrm{cos}I_f\mathrm{sin}(h\mathrm{\Pi }_1)`$
(135)
$`+{\displaystyle \frac{C}{L}}\dot{\mathrm{\Pi }}_1\left(\mathrm{sin}\pi _1\mathrm{cos}I_f\mathrm{cos}(h\mathrm{\Pi }_1)+\mathrm{cos}\pi _1\mathrm{sin}I_f\mathrm{sin}I_f\right)+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2).`$
To arrive to the final expression for $`I_r^{^{(rel)}}`$, we shall make use of (69 \- 71). These formulae will enable us to substitute, in the above expression, $`I_f`$ and $`h_f`$ with $`I`$ and $`h`$, correspondingly. All in all, in the coordinate system as in A.2.2.2 we have:
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right]{\displaystyle \frac{C}{L}}\dot{\pi }_1\mathrm{cos}I\mathrm{sin}h`$
(136)
$`+{\displaystyle \frac{C}{L}}\dot{\mathrm{\Pi }}_1\left(\mathrm{sin}\pi _1\mathrm{cos}I\mathrm{cos}h+\mathrm{cos}\pi _1\mathrm{sin}I\right)+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2).`$
In the coordinate system as in A.2.2.3, we obtain:
$`I_r^{^{(rel)}}=I+J\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{cos}\text{g}e\mathrm{cos}\left(2l+\text{g}\right)\right]{\displaystyle \frac{C}{L}}\dot{\pi }_1\mathrm{cos}I\mathrm{sin}(h\mathrm{\Pi }_1)`$
(137)
$`+{\displaystyle \frac{C}{L}}\dot{\mathrm{\Pi }}_1\left[\mathrm{sin}\pi _1\mathrm{cos}I\mathrm{cos}(h\mathrm{\Pi }_1)+\mathrm{cos}\pi _1\mathrm{sin}I\mathrm{sin}I\right]+O\left(J^2\right)+O(J\mathrm{\Phi }/\omega )+O((\mathrm{\Phi }/\omega )^2).`$
In (136) and (137), the first two terms coincide with those given by the second of formulae (2.6) in Kinoshita (1977). They make $`I_r^{^{(inert)}}`$, while the third term is $`I_r^{^{(\mathrm{\Phi })}}`$.
Appendix 4. Expression for $`h_r^{^{(rel)}}`$
Expressions (110) and (113) result in
$`\omega ^{^{(rel)}}\mathrm{sin}I_r^{^{(rel)}}\mathrm{sin}h_r^{^{(rel)}}=\dot{I}_f\mathrm{cos}h_f+\dot{\varphi }_f\mathrm{sin}I_f\mathrm{sin}h_f,`$ (138)
while (111) and (114) entail
$`\omega ^{^{(rel)}}\mathrm{sin}I_r^{^{(rel)}}\mathrm{cos}h_r^{^{(rel)}}=\dot{I}_f\mathrm{sin}h_f\dot{\varphi }_f\mathrm{sin}I_f\mathrm{cos}h_f.`$ (139)
Multiplying the former with $`\mathrm{cos}h_f`$ and the latter with $`\mathrm{sin}h_f`$, and then summing up the two results, we arrive at
$`\dot{I}_f=\omega ^{^{(rel)}}\mathrm{sin}I_r^{^{(rel)}}\mathrm{sin}\left(h_r^{^{(rel)}}h_f\right).`$ (140)
Since the difference $`h_r^{^{(rel)}}h_f`$ is expected to be of order $`O(J)+O(\mathrm{\Phi }/\omega )`$, the above formula may be rewritten as
$`h_r^{^{(rel)}}h_f={\displaystyle \frac{\dot{I}_f}{\omega ^{^{(rel)}}\mathrm{sin}I_f}}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega )`$ (141)
or, according to (120), as
$`h_r^{^{(rel)}}h_f={\displaystyle \frac{1+\alpha }{\omega ^{^{(inert)}}}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\left({\displaystyle \frac{I_f}{t}}+\mathrm{\Phi }_{I_f}\right)+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega )`$ (142)
where $`\alpha `$ is of order $`O(\mathrm{\Phi }/\omega )`$ and therefore may be dropped. Recall that, according to (117), the absolute value of the angular-velocity vector can be approximated, up to $`O(J^2)`$, with $`\omega ^{^{(inert)}}L/CG/C`$, while $`\sqrt{G^2L^2}`$ can be expressed as $`\sqrt{G^2L^2}=G\mathrm{sin}JGJLJ`$. Together with approximations $`\varphi _f=l+\text{g}J\mathrm{cot}I\mathrm{sin}\text{g}+O(J^2)`$ and $`I_f=I+J\mathrm{cos}\text{g}+O(J^2)`$ , it will enable us to rewrite (67) as
$`{\displaystyle \frac{I_f}{t}}=LJ\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{cos}(l+g)}{A}}{\displaystyle \frac{\mathrm{cos}l\mathrm{sin}(l+g)}{B}}\right]+O(J^2),`$ (143)
insertion whereof in (143) will then entail
$`h_r^{^{(rel)}}h_f=`$
$`{\displaystyle \frac{JC}{\mathrm{sin}I_f}}\left[{\displaystyle \frac{\mathrm{sin}l\mathrm{cos}(l+g)}{A}}{\displaystyle \frac{\mathrm{cos}l\mathrm{sin}(l+g)}{B}}\right]+{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I_f}}\mathrm{\Phi }_{I_f}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega )`$ (144)
This, along with $`h_f=h+J\mathrm{sin}g/\mathrm{sin}I+O(J^2)`$ and $`I_f=I+J\mathrm{cos}\text{g}+O(J^2)`$, yields:
$`h_r^{^{(rel)}}=h+J{\displaystyle \frac{\mathrm{sin}\text{g}}{\mathrm{sin}I}}+{\displaystyle \frac{JC}{\mathrm{sin}I}}\{{\displaystyle \frac{\mathrm{sin}l}{A}}\mathrm{cos}(l+\text{g}){\displaystyle \frac{\mathrm{cos}l}{B}}\mathrm{sin}(l+\text{g})\}+{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I}}\mathrm{\Phi }_{I_f}+O(...)`$
$`=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left\{\mathrm{sin}\text{g}+{\displaystyle \frac{C}{A}}{\displaystyle \frac{\mathrm{sin}\left(2l+\text{g}\right)\mathrm{sin}\text{g}}{2}}{\displaystyle \frac{C}{B}}{\displaystyle \frac{\mathrm{sin}\left(2l+\text{g}\right)+\mathrm{sin}\text{g}}{2}}\right\}`$
$`+{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I}}\left(\mu _1\mathrm{cos}h_f\mu _2\mathrm{sin}h_f\right)+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega )`$
$`=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{sin}\text{g}e\mathrm{sin}\left(2l+\text{g}\right)\right]`$
$`+{\displaystyle \frac{C}{L}}{\displaystyle \frac{1}{\mathrm{sin}I}}\left(\mu _1\mathrm{cos}h_f\mu _2\mathrm{sin}h_f\right)+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega ).`$
To get the answer in the precessing coordinate axes obtained from the inertial ones by two Euler rotations, as in A.2.2.2, we substitute (99) for $`\mu _1`$ and $`\mu _2`$. It yields:
$`h_r^{^{(rel)}}=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{sin}\text{g}e\mathrm{sin}\left(2l+\text{g}\right)\right]`$ (145)
$`\dot{\pi }_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{cos}h}{\mathrm{sin}I}}\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{sin}\pi _1\mathrm{sin}h}{\mathrm{sin}I}}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega ).`$ (146)
To obtain the answer in the Kinoshita precessing coordinate system obtained from the inertial one by three Euler rotations, as in A.2.2.3, we should substitute (103) for $`\mu _1`$ and $`\mu _2`$. This will entail:
$`h_r^{^{(rel)}}=h+{\displaystyle \frac{J}{\mathrm{sin}I}}\left(1{\displaystyle \frac{C}{2A}}{\displaystyle \frac{C}{2B}}\right)\left[\mathrm{sin}\text{g}e\mathrm{sin}\left(2l+\text{g}\right)\right]`$ (147)
$`\dot{\pi }_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{cos}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}\dot{\mathrm{\Pi }}_1{\displaystyle \frac{C}{L}}{\displaystyle \frac{\mathrm{sin}\pi _1\mathrm{sin}(h\mathrm{\Pi }_1)}{\mathrm{sin}I}}+O(J^2)+O((\mathrm{\Phi }/\omega )^2)+O(J\mathrm{\Phi }/\omega ),`$ (148)
the triaxiality parameter $`e`$ being rendered by (23). In (146) and (148), the first two terms coincide with those given by the first expression of (2.6) in Kinoshita (1977). They constitute $`h_r^{^{(inert)}}`$, while the third term is $`h_r^^\mathrm{\Phi }`$.
|
warning/0506/physics0506078.html
|
ar5iv
|
text
|
# Cherenkov Radiation from 𝑒⁺𝑒⁻ Pairs and Its Effect on 𝜈_𝑒 Induced Showers
## I Introduction
Cherenkov radiation from relativistic particles has been known for over 70 years . However, to date, almost all studies have concentrated on the radiation from individual particles. Frank , Eidman and Balazs considered the Cherenkov radiation from electric and magnetic dipoles, but only in the limit of vanishing separations $`d`$. Their work was nicely reviewed by Jelley .
Several more recent calculations have considered Cherenkov radiation from entire electromagnetic showers, in the coherent or almost coherent limit . The fields from the $`e^+`$ and $`e^{}`$ largely cancel, and the bulk of the coherent radiation is due to the net excess of $`e^{}`$ over $`e^+`$ (the Askaryan effect) . Hadronic showers produce radiation through the same mechanism . Coherent radiation occurs when the wavelength of the radiation is large compared to the radial extent of the shower; for real materials, this only occurs for radio waves.
Here, we consider another case, the reduction of radiation from slightly-separated oppositely-charged co-moving pairs. This includes $`e^+e^{}`$ pairs produced by photon conversion. When high-energy photons convert to $`e^+e^{}`$ pairs, the pair opening angle is small and the $`e^+`$ and $`e^{}`$ separate slowly.
Near the pair, the electric and magnetic fields from the $`e^+`$ and $`e^{}`$ must be considered separately. However, for an observer far away from the pair (compared to the pair separation $`d`$), the electric and magnetic fields from the $`e^+`$ and $`e^{}`$ largely cancel. Cherenkov radiation is produced at a distance of the order of the photon wavelength $`\mathrm{\Lambda }`$ from the charged particle trajectory. So, for $`d<\mathrm{\Lambda }`$, cancellation reduces the Cherenkov radiation from a pair to below that for two independent particles. For a typical pair opening angle $`m/k`$, where $`k`$ is the photon energy and $`m`$ the electron mass, without multiple scattering, $`\mathrm{\Lambda }>d`$ for a distance $`k\mathrm{\Lambda }/m`$. For blue light ($`\mathrm{\Lambda }=400`$ nm) from a 1 TeV pair, the radiation is reduced until the pair travels a distance of 40 cm (neglecting multiple scattering).
A similar cancellation effect was observed for energetic ($`100`$ GeV) $`e^+e^{}`$ pairs in nuclear emulsions . Ionization from newly created $`e^+e^{}`$ pairs is reduced when the pair separation is less than the screening distance for ionization in the target.
In this paper, we calculate the Cherenkov radiation from $`e^+e^{}`$ pairs, simulate optical radiation from pairs follow realistic trajectories, and consider the radiation from electromagnetic showers. We consider two classes of experiments: underwater/in-ice neutrino observatories and air Cherenkov telescopes.
## II Cherenkov radiation from pairs
Cherenkov radiation from closely spaced $`e^+e^{}`$ pairs can be derived by extending the derivation for point charges, by replacing a point charge with an oppositely charged, separated pair. We sketch the derivation for radiation from point charges, review previous work on radiation from infinitesimal dipoles, and derive the expression for Cherenkov radiation from a closely-spaced co-moving pair.
We follow the notation and derivation from Ref. . In Fourier space, the charge density $`\rho `$ and current density $`\stackrel{}{J}`$ from a point charge $`ze`$ propagating with speed $`v`$ in the $`x_1`$ direction can be written as
$`\rho (\stackrel{}{k},\omega )`$ $`={\displaystyle \frac{ze}{2\pi }}\delta (\omega k_1v)`$ (1)
$`\stackrel{}{J}(\stackrel{}{k},\omega )`$ $`=\stackrel{}{v}\rho (\stackrel{}{k},\omega )`$
where $`\stackrel{}{k}`$ is the wave vector and $`\omega `$ the photon energy. This current deposits energy into the medium through electromagnetic interactions. We use Maxwell’s equations beyond a radius $`a`$ around the particle track, where $`a`$ is comparable to the average atomic separation. Then, by conservation of energy, the Cherenkov radiation power is equal to the the energy flow through a cylinder of this radius, giving
$$\left(\frac{dE}{dx}\right)=caRe_0^{\mathrm{}}B_3^{}(\omega )E_1(\omega )𝑑\omega .$$
(2)
$`E_1`$ is the component of $`\stackrel{}{E}`$ parallel to the particle track, and $`B_3`$ is the component of $`\stackrel{}{B}`$ in the $`x_3`$ direction, evaluated at an impact parameter $`b`$ at a point with $`x_2=b`$, $`x_3=0`$. We omit the time-phase factors for brevity.
Using the wave equations in a dielectric medium and the definition of fields, then integrating over momenta, which eliminates the space-phase factors, one finds
$$E_1(\omega )=\frac{ize\omega }{v^2}\left(\frac{2}{\pi }\right)^{1/2}\left[\frac{1}{ϵ(\omega )}\beta ^2\right]K_0(\lambda b)$$
(3)
where
$`\lambda ^2={\displaystyle \frac{\omega ^2}{v^2}}[1\beta ^2ϵ(\omega )].`$
Similarly,
$`E_2(\omega )`$ $`={\displaystyle \frac{ze}{v}}\left({\displaystyle \frac{2}{\pi }}\right){\displaystyle \frac{\lambda }{ϵ(\omega )}}K_1(\lambda b)`$ (4)
$`B_3(\omega )`$ $`=ϵ(\omega )\beta E_2(\omega ).`$
$`K_0`$ and $`K_1`$ are the zeroth and first order modified Bessel functions of the second kind.
The far-field radiation depends on the asymptotic form of the energy deposition at $`|\lambda a|1`$. For $`\beta >1/\sqrt{ϵ(\omega )}`$ for real $`ϵ(\omega )`$, $`\lambda `$ is completely imaginary. The asymptotic contribution of the Bessel functions in the integrand of $`dE/dx`$ is finite, giving the well-known expression for the Cherenkov radiation
$`\left({\displaystyle \frac{dE}{dx}}\right)`$ $`={\displaystyle \frac{(ze)^2}{c^2}}`$ (5)
$`\times {\displaystyle _{ϵ(\omega )>1/\beta ^2}}\omega (1{\displaystyle \frac{1}{\beta ^2ϵ(\omega )}})d\omega .`$
Note how $`a`$ has dropped out \[10, Ch. 13\]. The derivation of this Cherenkov radiation may be expanded to give the field from a pair.
The radiation from an $`e^+e^{}`$ pairs depends on two parameters: the separation $`d`$ and the angle between the direction of motion and the orientation of the pair. For relativistic pairs created by photon conversion, the transverse (to the direction of motion) separation is important; the longitudinal separation of a highly relativistic pair can be neglected, due to Lorentz length contraction.
Balazs provided an expression for Cherenkov radiation from an infinitesimal dipole $`D`$ oriented transverse to its momentum. The fields are approximated by by a linear Taylor expansion of the corresponding point-charge fields:
$`E_1^{(D)}(\omega )=d{\displaystyle \frac{E_1(\omega )}{x_2}};B_3^{(D)}(\omega )=d{\displaystyle \frac{B_3(\omega )}{x_2}}`$
where $`d`$ is the effective pair separation, so $`D=zed`$. Then, following the same steps as in the point-charge case, Balazs finds
$`\left({\displaystyle \frac{dE}{dx}}\right)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{D^2}{c^4}}`$ (6)
$`\times {\displaystyle _{ϵ(\omega )>1/\beta ^2}}ϵ(\omega )\omega ^3(1{\displaystyle \frac{1}{\beta ^2ϵ(\omega )}})^2d\omega .`$
For a point dipole oriented parallel to its direction of motion, the radiation is negligible for $`\beta 1`$ .
To compute the Cherenkov radiation for finite separations $`d`$, let us consider a pair moving in the $`+x`$ direction. The pair lies entirely in the transverse plane $`y`$-$`z`$, with the line between them making an angle $`\alpha `$ with respect to the $`y`$-axis. Then, generalizing Eq. (1), the charge density from the pair is
$$\rho (\stackrel{}{k},\omega )=\frac{ze}{2\pi }\delta (\omega k_1v)\left[e^{i(k_2y_+k_3z_+)}e^{i(k_2y_{}k_3z_{})}\right].$$
The two charges have positions relative to the center of mass
$`y_+={\displaystyle \frac{d}{2}}\mathrm{cos}\alpha `$ $`z_+={\displaystyle \frac{d}{2}}\mathrm{sin}\alpha `$
$`y_{}={\displaystyle \frac{d}{2}}\mathrm{cos}\alpha `$ $`z_{}={\displaystyle \frac{d}{2}}\mathrm{sin}\alpha .`$
The angle $`\alpha `$ is the relative azimuth between the line connecting the two charges and the azimuth of observation.
The generalization of $`E_1(\omega )`$ of Eq. 3 is
$`E_1(\omega )=`$ $`{\displaystyle \frac{ize\omega }{v^2}}\sqrt{{\displaystyle \frac{\pi }{2}}}\left({\displaystyle \frac{1}{ϵ(\omega )}}\beta ^2\right)`$ (7)
$`\times \left[K_0(\lambda b_{})K_0(\lambda b_+)\right]`$ (8)
where
$$b_\pm =\sqrt{\frac{d^2}{4}\mathrm{sin}^2\alpha +(b\pm \frac{d}{2}\mathrm{cos}\alpha )^2}.$$
As before, we take $`|\lambda a|1`$ and $`a<b`$, so we need only consider $`db`$; there is little interference for $`db`$. Therefore, we can simplify using
$$b_\pm b\pm \frac{d}{2}\mathrm{cos}\alpha .$$
Then, as before, considering completely imaginary $`\lambda `$ and $`|\lambda a|1`$,
$`E_1(\omega )=`$ $`{\displaystyle \frac{2ze\omega }{c^2}}\left(1{\displaystyle \frac{1}{\beta ^2ϵ(\omega )}}\right)`$ (9)
$`\times \sqrt{{\displaystyle \frac{i}{|\lambda |}}}{\displaystyle \frac{e^{i|\lambda |b}}{b}}\mathrm{sin}\left[{\displaystyle \frac{d}{2}}|\lambda |\mathrm{cos}\alpha \right]`$
and a similar expression for $`B_3(\omega )`$. Here we have taken $`b_{}b_+b`$ in the denominator.
At $`\alpha =\pm \pi /2`$, $`E_1(\omega )=0`$. The Cherenkov radiation is no longer symmetric about the direction of motion, and vanishes at right angles to the direction of the dipole. As the charge separation increases (or the wavelength decreases), the angular distribution evolves from two wide lobes into a many-lobed structure, as shown in Fig. 1. After integration over even a narrow range of $`\omega `$ or $`d`$, the angular distribution becomes an almost-complete disk, with two narrow zeroes remaining at a direction perpendicular to the dipole vector.
After assembling the pieces, and averaging over $`\alpha `$, we find the generalization of Eq. (5),
$`\left({\displaystyle \frac{dE}{dx}}\right)=`$ $`{\displaystyle \frac{(ze)^2}{c^2}}{\displaystyle _{ϵ(\omega )>1/\beta ^2}}𝑑\omega `$ (10)
$`\omega \left(1{\displaystyle \frac{1}{\beta ϵ(\omega )}}\right)\times 2\left[1J_0(\lambda d)\right].`$
Here $`J_0`$ is the zeroth order Bessel function of the first kind. For $`\lambda d1`$, this reproduces Eq. (6). For $`\lambda d1`$, the $`dE/dx`$ is twice that expected for an independent particle (Eq. (5)). The transition is shown in Fig. 2. As the emission wavelength $`\mathrm{\Lambda }`$ approaches $`d`$, the pair spectrum converges to the point-charge spectrum in an oscillatory fashion, characteristic of the Bessel function. For certain values of $`\lambda d`$, the radiation exceeds that of two independent charged particles.
For the remainder of the paper, we assume that media satisfy $`\sqrt{ϵ(\omega )}=n`$, where $`n`$ is independent of frequency. In realistic detection media, any variation of $`n`$ with frequency is small, and would have little effect on Cherenkov radiation from relativistic particles.
With real $`e^+e^{}`$ pairs, two effects should be considered. Electromagnetic radiation is not emitted instantaneously, but occurs while the radiating particles travel a distance known as the formation length, $`l_f`$. For Cherenkov radiation, $`l_f=\mathrm{\Lambda }/\mathrm{sin}^2(\theta _C)=\mathrm{\Lambda }ϵ\beta ^2/(ϵ\beta ^21)`$ , depends only on the Cherenkov emission angle, $`\theta _C`$ and the photon wavelength; $`l_f`$ depends only slightly (through $`\beta `$) on the electron energy.
While the pair is covering the distance $`l_f`$, the pair separation will change by an amount $`\mathrm{\Delta }d=l_f\mathrm{sin}(\theta )`$, where $`\theta `$ is the angle between the $`e^+`$ and $`e^{}`$ velocity vectors. Since $`\theta `$ is of order $`1/\gamma `$, $`\mathrm{\Delta }d/d1`$, so the change in separation is not significant.
Second, the Cherenkov radiation produced at a point ($`x`$-coordinate) depends on the fields emitted by the charged particles at earlier times, when $`d`$ may be different than at the point of radiation. For full rigor, these retarded separations should be used in the calculation. Again, this has a negligible effect on the results.
## III Radiation from $`e^+e^{}`$ pairs in showers
Many experiments study Cherenkov radiation from large electromagnetic showers. The radiation from a shower may be less than would be expected if every particle were treated as independent. We use a simple simulation to consider 300 to 800 nm radiation from electromagnetic showers. This frequency range is typical for photomultiplier based Cherenkov detectors; at longer wavelength, there is little radiation, while shorter wavelength light is absorbed by the glass in the phototube.
We simulated 1000 $`\gamma `$ conversions to $`e^+e^{}`$ pairs with total energies from $`10^8`$ to $`10^{20}`$ eV. Pairs were produced with the energy partitioned between the $`e^+`$ and $`e^{}`$ following the Bethe-Heitler differential cross section $`d\sigma E_\pm (1E_\pm )`$, where $`E_\pm `$ is the electron (or positron energy). At high energies in dense media (above $`10^{16}`$ eV in water or ice), the LPM effect becomes important, and more asymmetric pairs predominate . The pairs are generated with initial opening angle of $`m/k`$; the fixed angle is a simplification, but the pair separation is dominated by multiple scattering, so it has little effect on our results.
The $`e^{}`$ and $`e^+`$ are tracked through a water medium (with $`n=\sqrt{ϵ}=1.3`$) in steps of $`0.02X_0`$, where $`X_0`$ is the radiation length, 36.1 cm in water. At each step, the particles multiple-scatter, following a Gaussian approximation \[14, Ch. 27\]. The particles radiate bremsstrahlung photons, using a simplified model where photon emission follows a Poisson distribution, with mean free path $`X_0`$. Although this model has almost no soft bremsstrahlung, soft emission has little effect on Cherenkov radiation, since the electron or positron velocity is only slightly affected.
At each step, we compute the Cherenkov radiation for each pair. They are treated coherently when $`d<2\mathrm{\Lambda }`$; at larger separations the particles radiate independently.
As shown in Fig. 3, the particles in lower energy pairs ($`<10^{10}`$ eV) radiate almost independently. In contrast, the radiation from very high energy pairs ($`>10^{15}`$ eV) is largely suppressed. The broad excursions slightly above unity occur when $`J_0(\lambda d)>1`$ for many of the scattered pairs.
## IV Implications for experiments
At least two types of astrophysical observatories depend on Cherenkov radiation. Water and ice based neutrino observatories observe Cherenkov radiation from the charged particles produced in neutrino interactions, and air Cherenkov telescopes look for $`\gamma `$-ray induced electromagnetic showers in the Earth’s atmosphere.
Current neutrino observatories can search for electron neutrinos with energies above 50 TeV (for $`\nu _\mu `$, the threshold is much lower) . They use large arrays of photomultiplier tubes to observe the Cherenkov radiation from $`\nu _e`$ induced showers. For water, $`n1.3`$, Fig. 3 shows that $`\lambda d<1`$ while the pair travels significant distances. Ice is similar to water, with a slightly lower density; $`n`$ of ice depends on its structure, and is typically $`1.29`$ .
To quantify the effect of Cherenkov radiation from $`\nu _e`$ interactions, we use a toy model of an electromagnetic shower. The shower evolves through generations, with each generation having twice as many particles as the preceding generation, with half the energy. Each generation evolves over a distance of $`X_0`$; other simulations have evolved generations over a shorter distance $`(\mathrm{ln}2)X_0`$, leading to a more compact shower . In these showers, most of the particles are produced in the last radiation lengths.
Fig. 4 shows the Cherenkov radiation expected from a model $`10^{20}`$ eV shower with coherent Cherenkov radiation (solid line) and in a model where all particles radiate independently (dotted line). This model does not include the LPM effect, so it should be considered only illustrative. The LPM effect lengthens the high-energy (above a few $`10^{15}`$ eV) portion of the shower. By spreading the shower longitudinally, the LPM effect will give the electrons and positrons more time to separate, and so will somewhat lessen the difference between the two results. However, it is clear from Fig. 4 that coherence has a significant effect for the first $`22`$ generations. Since the front of the shower contains relatively few particles, it will not affect the measured energy; the change in number of radiated photons (and hence on the energy measurement) should be less than 1%.
However, the suppression will affect the apparent length of the shower. For the first $`8`$ generations, the shower will emit less light than a single charged particle. Because of the LPM effect, each of these generations (with mean particle energy $`E_g`$ above a few $`10^{15}`$ eV) develop over a distance $`X=X_0\sqrt{E_g/5E_{LPM}}`$, where $`E_{LPM}=278`$ TeV for water is the effective LPM energy , greatly elongating the shower. So, the first 8 generations include most of the length of the shower. So, the suppression of Cherenkov radiation hides the initial shower development, making the shower appear considerably more compact. The reduction in early-stage radiation should help in separating electron cascades from muon-related backgrounds, especially muons that undergo hard interactions, and lose a large fraction of their energy.
Atmospheric Cherenkov telescopes like the Whipple observatory study astrophysical $`\gamma `$-rays with energies from 100 GeV to 10 TeV. These telescopes observe Cherenkov radiation from pairs in the upper atmosphere; for a 1 TeV shower, the maximum particle density occurs at an altitude of 8 km above sea level (asl) , where the density is about $`1/3`$ that at sea level. Since $`n1`$ depends linearly on the density, at 8 km asl $`n11\times 10^4`$, so for 500 nm photons radiated from ultra-relativistic particles, $`\lambda d<1`$ only for $`d<6\mu \text{m}`$. In this low-density medium, the effect of the pair opening angle is significant and multiple scattering is less important. Pairs with $`k<1`$ TeV will separate by 30 $`\mu `$m in a distance less than 30 meters; at 8 km asl, this is 3% of a radiation length. This distance is too short to affect the radiation pattern from the shower.
Cherenkov radiation is also used in lead-glass block calorimetry, and in Cherenkov counters for particle identification; their response to photon conversions may be affected by this coherence.
Although the reactions are slightly different, a similar analysis applies to the reduction of ionization by $`e^+e^{}`$ pairs. Perkins observed that the ionization from pairs with mean energy 180 GeV in emulsion was surppressed for the first $`250`$ $`\mu `$m after the pairs were created . With $`X_0=3`$ cm (typical for emulsion), the $`e^+`$ and $`e^{}`$ trajectories will be about 4 nm apart after travelling $`250`$ $`\mu `$m. For relativistic particles, the screening distance (effective range for $`dE/dx`$) is determined by the plasma frequency of the medium, $`\omega _p`$. For silver bromide, the dominant component of emulsion, $`\mathrm{}\omega _p=48`$ eV (in a complete emulsion, $`\mathrm{}\omega _p`$ will be slightly lower). This yields a screening distance $`c/\omega _p=4`$ nm, which is very close to the calculated separation.
## V Conclusion
We have calculated the Cherenkov radiation from $`e^+e^{}`$ pairs as a function of the pair separation $`d`$. When $`d^2<v^2/(\omega ^2[1\beta ^2ϵ(\omega )])`$, the radiation is suppressed compared to that from two independent particles.
This suppression affects the radiation from electromagnetic showers in dense media. Although the total radiation from a shower is not affected, emission from the front part of the shower is greatly reduced; this will affect studies of the shower development, and may affect measurements of the position of the shower.
This work was funded by the U.S. National Science Foundation under Grant number OPP-0236449 and the U.S. Department of Energy under contract number DE-AC-76SF00098.
|
warning/0506/astro-ph0506564.html
|
ar5iv
|
text
|
# Spiders in Lyot Coronagraphs
## 1. Introduction
In the last few years several exciting new coronagraphic designs have been proposed. These designs focus on suppressing starlight within a few resolution elements of bright stars in order to detect planetary companions of these stars fractions of an arcsecond from the star, a task which requires a contrast ratio of more than a million to one for extrasolar jovian planets (Burrows et al., 2004; Burrows, 2005). An even more ambitious goal, that of detecting and characterizing extra-solar terrestrial planets, demands contrast ratios of a billion to one or higher.
Recent progress in coronagraphic concepts has yielded several classes of coronagraphs that achieve the requisite suppression in the ideal case (e.g., Kuchner & Traub, 2002; Soummer et al., 2003a, b; Kasdin et al., 2003; Guyon, 2003), which generally requires zero wavefront error, and unobscured apertures. The fully-optimized, diffraction-limited Lyot project coronagraph (Lyot, 1939; Sivaramakrishnan et al., 2001; Oppenheimer et al., 2003, 2004; Digby et al., 2004; Makidon et al., 2005), which is deployed at the Air Force AEOS 3.6m telescope (Roberts & Neyman, 2002) is the first coronagraph to operate in the regime of ‘extreme adaptive optics’ (ExAO) with Strehl ratios around 90% under the best seeing conditions in the $`H`$-band. This instrument, and all future ground and space coronagraphs, will have to work with non-zero wavefront error, and on aperture geometries subject to real-world engineering constraints. In Lloyd & Sivaramakrishnan (2005) and Sivaramakrishnan et al. (2005) we analyzed the effects of wavefront errors on Lyot coronagraphs. In this paper, we focus on the effects of spider diffraction in a Lyot coronagraph.
Next generation AO systems, designed for Gemini and an ESO Very Large Telescope, are likely to be dedicated to coronagraphic imaging and spectroscopy in the search for exo-jupiters. It is therefore important to understand and quantify the effects spiders have on coronagraphic imaging as these systems are designed and constructed over the next five years. This is particularly relevant in light of recent work on high dynamic range coronagraphy on obscured apertures (Soummer, 2005).
We focus on an analytical treatment of this problem with the goal of producing some general understanding of the problem, not just a calculation for a specific case at a specific telescope. Over and above the narrow application of specific numerical calculations, there is a computational difficulty in modelling coronagraphs on obstructed apertures when taking spider vanes into account. For example, coronagraphic image simulations require about 6 to 8 samples across each resolution element in order to accurately model the effects of a $`4\lambda /D`$ image plane occulting stop (‘focal plane mask’) in the focal plane ($`\lambda `$ being the central wavelength of the observing bandpass, and $`D`$ the telescope diameter). This entails optical calculations using arrays 6 to 8 times the diameter of the entrance pupil. The 8 m Gemini telescope’s spider vanes are 1 cm thick. Thus the numerical arrays required to model the optical train can be 12000–16000 elements on a side. Furthermore, these numerical calculations can require extra-ordinary resolution and dynamic range to converge on the correct answer, and the effects of aliasing in an FFT can be quite severe at the $`10^9`$ contrast level \[e.g., even very high resolution numerical calculations of the Four Quadrant Phase Mask coronagraph (Rouan et al., 2000) are inadequate to correctly calculate the coronagraphic point-spread function (PSF)\], requiring an analytical solution (Lloyd et al., 2003). Fourier transform routines of optical modelling programs will need more memory than can be put in most of today’s computers. We are therefore forced to fall back on developing analytical estimates of the effect of secondary mirror support vanes on coronagraphic images. We treat classical Lyot coronagraphs (Lyot, 1939) with their hard-edged, opaque focal plane masks; modified Lyot coronagraphs with graded focal plane masks described by Gaussian<sup>1</sup><sup>1</sup>1Ftaclas, private communication or ‘band-limited’ (Kuchner & Traub, 2002) functions; as well as apodized pupil Lyot coronagraphs (Soummer et al., 2003a; Soummer, 2005). Guyon (2003) discusses pure pupil apodized high dynamic range imaging on arbitrarily-shaped apertures, Soummer et al. (2003b) treats dual-zone phase-mask coronagraphy on arbitrary apertures, and Lloyd et al. (2003) presents ways of dealing with four-quadrant phase mask coronagraphy on centrally-obscured apertures with spiders.
## 2. Monochromatic coronagraphic theory with spider supports
Here we briefly recapitulate our basic monochromatic Fourier optics formalism. A more detailed treatment can be found in Sivaramakrishnan et al. (2001); Lloyd & Sivaramakrishnan (2005). We recollect that a plane monochromatic wave travelling in the $`z`$ direction in a homogenous medium without loss of energy can be characterized by a complex amplitude $`E`$ representing the transverse (e.g., electric) field strength of the wave. The full spatio-temporal expression for the field strength is $`Ee^{(i\kappa z\omega t)}`$, where $`\omega /\kappa =c`$, the speed of the wave. We do not use the term field to denote image planes — the traditional optics usage — we always use the term to denote electromagnetic fields or scalar simplifications of them. The wavelength of the wave is $`\lambda =2\pi /\kappa `$. The time-averaged intensity of a wave at a point is proportional to $`EE^{}`$, where the average is taken over one period, $`T=2\pi /\omega `$, of the harmonic wave. The phase of the complex number $`E`$ represents a phase difference from the reference phase associated with the wave. A real, positive $`E`$ corresponds to an electric field oscillating in phase with our reference wave. A purely imaginary positive value of $`E`$ indicates that the electric field lags by a quarter cycle from the reference travelling wave. Transmission through passive, linear filters such as apertures, apodizers, and so forth, is represented by multiplying the field strength by the transmission of these objects which modify the wave. Again, such multiplicative modification is accomplished using complex numbers to represent phase changes forced on the wave incident on such objects.
We assume that Fourier optics describes our imaging system: image field strengths are given by the Fourier transform of aperture (or pupil — we use the two terms interchangeably) illumination functions, and vice versa.
A telescope aperture is described by a transmission function pattern $`A(x)`$, where $`x=(x_1,x_2)`$ is the location in the aperture, in units of the wavelength of the light (see Fig. 1). The corresponding aperture illumination describing the electric field strength in the pupil (in response to an unaberrated, unit field strength, monochromatic incident wave) is $`E_A=A(x)`$. From this point onwards we drop the common factor $`Ee^{(i\kappa z\omega t)}`$ when describing fields. The aperture intensities ($`E_AE_A^{}`$) for two coronagraphic designs are shown in Fig. 2 (top row). The field strength in the image plane, $`E_B=a(k)`$, is the Fourier transform of $`E_A`$, where $`k=(k_1,k_2)`$ is the image plane coordinate in radians. Because of the Fourier relationship between pupil and image fields, $`k`$ is also a spatial frequency vector for a given wavelength of light. We refer to this complex-valued field $`a`$ as the ‘amplitude-spread function’ (ASF), by analogy with the PSF of an optical system. The PSF is $`aa^{}`$. Our convention is to change the case of a function to indicate its Fourier transform. We multiply the image field $`E_B`$ by a mask function $`m(k)`$ to model the focal plane mask of a coronagraph. The image field immediately after this mask is $`E_C=m(k)E_B`$. The electric field in the re-imaged pupil after the focal plane mask (the Lyot pupil, see Fig. 2, middle rows) is $`E_D`$, which is the Fourier transform of $`E_C`$. We use the fact that the transform of the image plane field $`E_B`$ is just the aperture illumination function $`E_A`$ itself, so the Lyot pupil field is $`E_D=M(x)E_A`$, where $``$ is the convolution operator.
If the Lyot pupil stop transmission is $`N(x)`$, the electric field after the Lyot stop is $`E_E=N(x)E_D`$. The transform of this last expression is the final coronagraphic image field strength: $`E_F=n(k)[m(k)E_B]`$. In this paper we look at the effects of secondary support spiders on the final coronagraphic PSF corresponding to the field strength $`E_F`$. We take into account the fact that $`A(x)`$ may be apodized (Fig. 2, right), so $`A`$ is a graded function rather than a function that takes values of either 0 or 1. Understanding high dynamic range Lyot coronagraphy hinges on understanding the structure of the field strength $`E_D`$ in the Lyot plane located at D, as well as the repercussions of such structure in the final image plane $`E_F`$.
## 3. Spider diffraction in a Lyot coronagraph
The mask function in a Lyot coronagraph is best expressed as $`m(k)=1w(k)`$, where $`w(k)`$ is the ‘focal plane mask shape’ function. For a hard-edged mask $`w(k)=\mathrm{\Pi }(D|k|/s)`$, where $`s`$ is the mask diameter in units of the resolution of the optical system, $`\lambda /D`$. The function $`\mathrm{\Pi }(\alpha )`$ takes the value unity for $`|\alpha |<0.5`$, and is zero elsewhere. We note that $`w(0,0)=1`$ (which constrains $`W(x)`$ to have unit area). The Fourier transform of the focal plane mask transmission function $`m(k)`$ is $`M(x)=\delta (x)W(x)`$, so the Lyot pupil electric field of a Lyot coronagraph can be expressed as
$`E_D`$ $`=`$ $`[\delta (x)W(x)]A(x).`$ (1)
(An example of the mask transmission function for a simple band-limited coronagraph is shown in Fig. 3.)
We must understand the morphology of the Lyot pupil field in order to understand the extent to which secondary spider supports reduce the final coronagraphic image’s dynamic range.
We model a single spider vane across the entrance pupil of a telescope by writing
$`E_A`$ $`=`$ $`A(x)(1\mathrm{\Pi }(x_1/ϵ))`$ (2)
($`ϵ`$ being the width of the spider vane). Secondary obstructions result in different effects. We do not treat them here (see e.g., Sivaramakrishnan et al., 2001; Lloyd et al., 2003, for details on this topic). Equation (1) with this aperture function produces a Lyot pupil field
$`E_L`$ $`=`$ $`A(x)(1\mathrm{\Pi }(x_1/ϵ)\delta (x_2))`$ (3)
$`W(x))[A(x)(1\mathrm{\Pi }(x_1/ϵ)\delta (x_2))].`$
We are only concerned with the ‘interior’ of the aperture in the Lyot plane, where, by design, we are satisfied with coronagraphic performance of our aperture without spider support vanes. This means that in our estimation, $`A(x)WA(x)`$ is sufficiently small for our scientific purposes (e.g., Soummer et al., 2003a; Soummer, 2005) or zero (Rouan et al., 2000; Aime et al., 2001, 2002; Kuchner & Traub, 2002; Soummer et al., 2003b). We therefore drop this component of equation (3), so in the interior of the Lyot pupil we obtain
$`E_{L,\mathrm{interior}}`$ $`=`$ $`A(x)[\mathrm{\Pi }(x_1/ϵ)\delta (x_2)]`$ (4)
$`W(x)(A(x)\mathrm{\Pi }(x_1/ϵ)\delta (x_2)).`$
This is the contribution (in the Lyot pupil) due to diffraction from a long thin obstruction such as a spider vane in the entrance aperture.
Fig. 2 (left) shows the intensity of the Lyot pupil field using a perfect theoretical coronagraph, the band-limited coronagraph, where, in the absence of spider vanes and optical aberrations, the interior of the Lyot field is identically zero.
There are two components to the field strength in the interior region of the Lyot pupil. The bright central stripe is exactly the spider vane width $`ϵ`$ — its brightness is very close to that of the brightness of the clear or apodized aperture in the entrance pupil. This is the first term in equation (4), viz., $`A(x)[\mathrm{\Pi }(x/ϵ)\delta (x_2)]`$. For high dynamic range applications it can be masked out with a thin opaque strip in the Lyot pupil stop. This strip can be oversized for practical reasons without noticeably affecting throughput at the Lyot stop. We call this kind of Lyot stop a ‘Lyot spider stop’. It resembles the ‘reticulated Lyot stop’ of Sivaramakrishnan & Yaitskova (2005), which masks out bright inter-segment gaps in the coronagraphic Lyot plane of extremely large, highly-segmented telescopes.
The extended low-intensity ‘aura’ of the bright spider vane is described by the second term in equation (4), $`W(x)(A(x)\mathrm{\Pi }(x_1/ϵ)\delta (x_2))`$ (Fig. 2, bottom row). For the applications we consider here, the equivalent width of the function $`W(x)`$ is of order $`D/s`$, since we wish to search for faint companions and structure outside an inner working angle of about $`s\lambda /2D`$ of a bright, on-axis star. Typically $`s`$ will lie between 4 and 10.
When the width of the focal plane mask becomes larger than $`10\lambda /D`$, this strip of dimmer light can be be removed by a Lyot pupil stop with an oversized spider vane obstruction without sacrificing Lyot stop throughput too much. We point out that the exact geometry of the spiders is not relevant — our approach can deal with non-orthogonal spiders as easily as with perfectly aligned spiders.
### 3.1. The coronagraphic PSF with spiders
We now estimate the field strength in the final coronagraphic image with a Lyot spider stop. The coronagraphic ASF without a Lyot stop of any kind is the Fourier transform of the field in the interior of the Lyot pupil described by equation (4):
$`a_c(k)`$ $`=`$ $`ϵa(k)\mathrm{sinc}(ϵk_1)`$ (5)
$`ϵw(k)\left(a(k)\mathrm{sinc}(ϵk_1)\right).`$
The first term of the ASF $`a_c(k)`$ is $`ϵa(k)[\mathrm{sinc}(ϵk_1)]`$. This has the shape of a regular spider diffraction spike, but is down a factor $`ϵ^2`$ in intensity from the direct image’s spider spike. Masking out the bright spider vane removes this first term in the expression for $`a_c(k)`$. The remaining term is modulated by the mask shape function $`m(k)`$ itself. If the mask shape is the circular top hat function $`\mathrm{\Pi }(D|k|/s)`$, spider diffraction in the coronagraphic image plane is confined to the region behind the mask when a Lyot spider stop is used. The $`\mathrm{sinc}(ϵk_1)`$ function describes the direct image’s spider profile along the length of the spider ‘spike’ in the image plane. This function is close to unity in the region behind or just around the focal plane mask in the direct or coronagraphic image planes. $`w(k)`$ is also unity at scales where $`a(k)`$ is significant in size if the focal plane mask is a few to several resolution elements wide. Thus, where $`w(k)`$ is non-zero (for hard-edged masks), $`a(k)\mathrm{sinc}(ϵk_1)`$ is very close to the integral of $`a(k)`$ over its entire domain. The value of this integral is $`A(0)`$ because of the Fourier relation between $`a`$ and $`A`$. We stress that we are trying to estimate the size of the spider effects here rather than calculate them exactly.
However, we should consider an appropriately optimized Lyot stop — on a traditional or ‘classical’ Lyot coronagraph (or a band-limited coronagraph) this stop is undersized relative to the entrance pupil. In apodized pupil Lyot coronagraphs, the Lyot stop is not undersized; it lets the entire graded entrance pupil through. In the former case we would need to convolve $`a_c(k)`$ by the Fourier transform of this undersized pupil. If the undersized classical Lyot stop is described by the function $`A^{}(x)`$, the classical or band-limited coronograph’s ASF is
$`a_{c,LC}(k)`$ $`=`$ $`ϵa^{}(k)\mathrm{sinc}(ϵk_1)`$ (6)
$`ϵa^{}(k)\left[w(k)\left(a(k)\mathrm{sinc}(ϵk_1)\right)\right]`$
(we note that $`a^{}(k)a(k)=a^{}(k)`$ if the Lyot stop and entrance pupil are not apodized, and the Lyot stop support is a subset of the entrance pupil support).
With a Lyot spider stop the first term disappears, leaving only the second term in equation (6). We see from this that the effects of spider diffraction are now concentrated behind the mask, as in the case of $`a_c`$, but because of a convolution with $`a^{}(k)`$, the diffracted light leaks out about a diffraction width outside the actual mask (in the final coronagraphic image plane). Thus we can immediately conclude that for a hard-edged focal plane mask with diameter $`s\lambda /D`$, combined with an optimized Lyot spider stop, residual spider diffraction is very small outside a circle of diameter $`s\lambda /D+\lambda /D^{}`$ around the on-axis star ($`D^{}`$ being the Lyot stop outer diameter, projected back to the primary mirror). Extending this argument to obstructed apertures is straightforward: the spider diffraction will fall drastically at about a PSF equivalent width away from the geometric shadow of the focal plane mask in the final image plane. The secondary obstruction will make for a larger distance scale for the fall-off of residual spider diffraction outside the focal plane mask’s edge. For graded focal plane masks (such as Gaussian or band-limited coronagraphs), residual spider diffraction after using an optimized Lyot spider stop extends everywhere that the mask shape function is non-zero.
Monochromatic coronagraphic images using a simple and a spider Lyot stop in a band-limited coronagraph are shown in Fig. 4, cases BL-a and BL-b respectively. Here the perfect coronagraph on a clear unobstructed aperture with no phase or amplitude errors produces no on-axis light whatsoever in the Lyot pupil after it is stopped down by the Lyot stop (Fig. 2, in the ‘no spider’ column’s Lyot stop intensity), or in the following coronagraphic image plane. A Lyot stop optimized without consideration of spider vanes lets through a horizontally-oriented ‘spider spike’ of light (Fig. 4, BL-a). If there were two crossed spider vanes across the aperture, this image would show the usual cross observers associate with secondary support spiders. The ‘ringing’ in the vertical direction has a period of the resolution element induced by the Lyot stop: approximately twice the size of the non-coronagraphic resolution $`\lambda /D`$ for this design. Removal of the bright strip of light from the Lyot plane (where the geometrical image of the spider vane is located) by using a Lyot spider stop removes much of the light from the coronagraphic image (Fig. 4, BL-b). The dark vertical stripes are located at multiples of $`8\lambda /D`$ in this panel. The focal plane mask structure of the band-limited coronagraph design in this example is seen by comparing this residual coronagraphic image with the focal plane mask itself (shown in Fig. 3).
For the apodized pupil Lyot coronagraph the convolution is with the Fourier transform of the full entrance pupil without apodization. This coronagraphic design does not undersize the Lyot pupil, so $`D^{}=D`$. For the unobscured aperture with a spider, this Fourier transform is the $`\mathrm{jinc}`$ function, $`2J_1(x)/x`$ (where $`J_1(x)`$ is the Bessel function of the first kind, with index $`1`$) —
$`a_{c,APLC}(k)`$ $`=`$ $`ϵ\mathrm{jinc}(Dk)a(k)\mathrm{sinc}(ϵk_1)`$ (7)
$`ϵ\mathrm{jinc}(Dk)\left[w(k)\left(a(k)\mathrm{sinc}(ϵk_1)\right)\right].`$
We can treat centrally obscured apertures by using the difference of two $`\mathrm{jinc}`$ functions, although for simplicity we show the unobscured aperture case here. Once again, an optimized Lyot spider stop will remove the first term in equation (7), leaving the second term. This term is reduced in field strength by a factor $`ϵ`$ from what is essentially $`A(0,0)`$, and it is also restricted to a diffraction-width (or equivalent width of the PSF for obscured apertures) around the the projection of the focal plane mask on the final coronagraphic image plane.
### 3.2. Apodized pupil and apodized occulter coronagraphs
On a telescope with a secondary mirror obstructing the entrance aperture, it is possible to design Lyot coronagraphs such as apodized pupil Lyots, and Gaussian or band-limited coronagraphs. The presence of spider vanes makes the apodized pupil design preferable to the apodized occulter designs. If we examine Figs. 2 and 4 in the case of the theoretically perfect band-limited design we can understand why this is the case.
The focal plane mask used in this example possesses a $`1\mathrm{jinc}(\alpha k)`$ transmission function. $`\alpha `$ is chosen to produce a first zero in the $`\mathrm{jinc}`$ function at $`k=8\lambda /D`$. The width of the residual broad swath of low intensity light is twice the bandwidth or equivalent width of the function describing the focal plane mask (Fig. 2, Lyot plane on the left).
In the image plane of the band-limited coronagraph with the Lyot spider stop (Fig. 4 BL-b), we see that light diffracted by the spider vane spills into the coronagraphic image plane wherever the focal plane mask has a transmission that is not unity, i.e., wherever our focal plane mask had any opacity whatsoever. In fact the light occupies a slightly larger area (by about one resolution element) of the image after passage through an undersized Lyot stop.
With an apodized pupil design that utilizes a hard-edged focal plane mask, the on-axis coronagraphic image (which is due almost entirely to spider vane diffraction when we use our analytical, circularly-symmetric apodizer) is localized to a circular area with a radius one resolution element larger than the original focal plane mask (Fig. 4 AP-b). Thus, for the 8 m Gemini telescope with an apodized pupil Lyot coronagraph with a $`4\lambda /D`$ diameter focal plane mask, all diffraction from the spider vanes is restricted to a disk $`6\lambda /D`$ in diameter. Furthermore, the contrast ratio between this diffracted light and a PSF taken with the same apodizer and Lyot stop, but no focal plane stop, is of the order of $`10^6`$ with a simple analytical apodizer design.
This is visible in Fig. 4, where we see the structure of the focal plane mask in residual spider diffraction after using a Lyot spider stop. We simulated a single spider vane three pixels wide across a 170 pixel diameter aperture. We calculated the PSFs with the band-limited design described here, as well as with a numerical approximation of an analytical apodized pupil Lyot coronagraph optimized for a circular pupil (Soummer et al., 2003a). While Soummer (2005) shows that it is possible to refine the apodized pupil coronagraph design to accomodate spiders, we used the circular aperture apodizer design because an azimuthally symmetric apodizer is easier to align in practice, and the analytical apodizing function is easy to generate numerically. The particular example we present is taken from Soummer et al. (2003a). It has a 19% throughput apodizer matched to a $`3.74\lambda /D`$ focal plane mask diameter. Figs. 5 and 6 show cuts through the PSFs of both designs, with and without Lyot spider stops. These cuts show that the apodized pupil coronagraphic design outperforms the band-limited design away from residual spider spikes in the coronagraphic image when the bright spiders are blocked in the Lyot pupil, in spite of the fact that the apodized pupil design’s focal plane mask equivalent width is smaller than that of our band-limited example. This localization of diffracted light due to spiders suggests that apodized pupil coronagraphic designs suit future ground-based adaptive optics coronagraphic instruments on existing telescopes.
## 4. Discussion
We have demonstrated that even after masking out bright spiders in the Lyot plane of a Lyot coronagraph, some residual spider diffraction will be seen around the focal plane mask in the coronagraphic image. These effects cause a brightness proportional to the square of the spider thickness, but are localized to a PSF-width around the focal plane mask. Hard-edged masks show a stronger localization of the light diffracted by spiders than Gaussian or band-limited masks. This localization is desirable behavior for an extreme adaptive optics coronagraph on existing 8–10 m class telescopes, all of which possess secondary mirror support spider vanes.
The residual brightness due to spider vane diffraction will affect speckle statistics in this region (Aime & Soummer, 2004), inducing a larger variance in intensity there. However, on hard-edged focal plane mask coronagraphs, masking out the bright spiders in the Lyot plane does enable good suppression of spider diffraction just one resolution element away from the mask edge. Opto-mechanical tolerances for on-axis telescope designs are looser than for extreme off-axis designs. The advent of on-axis apodized pupil Lyot coronagraphs with good suppression of on-axis sources (Soummer, 2005) makes it important to understand the effects of spiders in Lyot coronagraphs and their modern variants when designing coronagraphs to search for extrasolar jovian companions using ground-based, next-generation adaptive optics systems on today’s 8–10 m class telescopes.
We acknowledge frequent helpful discussions with R. Soummer, and his generous contribution of a numerical realization of the apodized pupil Lyot coronagraph on an unobstructed circular aperture. We are grateful to the referee for useful comments, and to P. E. Hodge, P. Greenfield, J. T. Miller, and N. Dencheva for their role in developing and supporting the Python Numarray module (Greenfield et al., 2002, 2003), wrapping the numerical Fourier transform library FFTW (Frigo & Johnson, 1997) for Numarray, and providing support for matplotlib (Hunter, 2005). We also thank the Space Telescope Science Institute’s Research Programs Office and Director’s Discretionary Research Fund for support. This work has also been supported by the National Science Foundation Science and Technology Center for Adaptive Optics, managed by the University of California at Santa Cruz under cooperative agreement No. AST-9876783, and by the National Science Foundation under Grant No. AST-0215793 and Grant No. AST-0334916.
|
warning/0506/cond-mat0506505.html
|
ar5iv
|
text
|
# Tunneling resonances in quantum dots: Coulomb interaction modifies the width
## Abstract
Single-electron tunneling through a zero-dimensional state in an asymmetric double-barrier resonant-tunneling structure is studied. The broadening of steps in the $`I`$$`V`$ characteristics is found to strongly depend on the polarity of the applied bias voltage. Based on a qualitative picture for the finite-life-time broadening of the quantum dot states and a quantitative comparison of the experimental data with a non-equilibrium transport theory, we identify this polarity dependence as a clear signature of Coulomb interaction.
Single-electron tunneling through zero-dimensional states has been observed in a wide variety of systems, including metallic islands, lateral quantum dots in gated semiconductor devices, vertical dots in double-barrier resonant tunneling structures, and molecular systems such as carbon nanotubes review\_exp . The $`I`$$`V`$ characteristics has the shape of a staircase, in which each step is associated with the opening of a new transport channel through the system. Many features observed in the transport measurements deshnull ; su ; equ ; thomas ; icps can be explained either within a single-particle picture for non-interacting electrons or by the orthodox theory of sequential tunneling beenakker ; averin ; glazman valid for weak dot-lead tunnel coupling. This is no longer the case for interaction effects on transport beyond the weak-tunneling limit. A famous example is the zero-bias anomaly of Kondo-assisted tunneling Kondo\_theory ; Kondo\_exp .
In this work, we report on a new clear signature of Coulomb interaction beyond weak tunneling, that is achieved under much less stringent experimental conditions than required for the Kondo effect to occur. This signature is contained in the width of the first step of the $`I`$$`V`$ characteristics. We observe that the width strongly depends on the polarity of the applied bias voltage. This behavior can neither be explained within a single-particle picture nor by sequential-tunneling theory beenakker ; averin ; ralph . Due to Coulomb interaction, the finite-life-time broadening of the dot levels becomes energy dependent. As a consequence, the values for the broadening at the two considered steps of opposite polarity can differ by up to a factor of two for strongly asymmetric coupling strengths of the two tunnel barriers. We use the results of a diagrammatic real-time transport theory koenig that includes the above described physics to find reasonable quantitative agreement with the experimental data.
The experiment was performed with a highly asymmetric double-barrier resonant-tunneling device grown by molecular beam epitaxy on an n<sup>+</sup>–type GaAs substrate. An undoped 10 nm wide GaAs quantum well is sandwiched between 5 and 8 nm thick Al<sub>0.3</sub>Ga<sub>0.7</sub>As tunneling barriers separated from highly-doped GaAs contacts (Si-doped with $`n_{\text{Si}}=4\times 10^{17}\text{cm}^3`$) by 7 nm thick undoped GaAs spacer layers. The sample was fabricated as a pillar of $`2`$ $`\mu `$m diameter. Two-terminal dc-measurements of the $`I`$$`V`$ characteristics were performed in a dilution refrigerator at temperatures between 20 mK and 1 K. The studied GaAs quantum well embedded between two AlGaAs barriers can be viewed as a two-dimensional system with the edges and residual impurities confining the lateral electron motion and thus forming dots. Tunneling through the energetically lowest state of the dot, at the energy $`E_0=33`$ meV prl and with a lateral extent of $`10`$ nm, produces the lowest resonance peak in the differential conductance $`G=dI/dV`$.
The system considered is sketched in Fig. 1(a). Electrons tunnel from the heavily-doped emitter through the spin-degenerate quantum-dot level embedded inside the quantum well. The level comes to resonance with the emitter’s electro-chemical potential at a finite bias voltage, indicated by a step in the current, as shown in Fig. 1(b).
Coulomb interaction prevents double occupation of the level for the considered range of bias voltage. This gives rise to charging effects that are differently important for the two polarities. The bottleneck of transport is provided by the thicker tunnel barrier. For $`V<0`$, the “charging direction”, as sketched in Fig. 1 (a), the dot is predominantly singly occupied, and, therefore, the two spin channels effectively block each other. This reduces the current-step height to one half of the value for the opposite polarity, $`V>0`$, the “non-charging direction”, for which the dot is predominantly empty averin ; glazman ; ralph ; equ ; thomas ; icps . This is well understood within the sequential-tunneling picture.
The focus of this paper, however, is the finite broadening of the step edge. This broadening can be measured as the full width at half maximum (FWHM)-value of the differential conductance ($`dI/dV`$) peaks associated with the current steps. Figure 1(c) shows our experimental low-temperature differential conductance peaks for both bias polarities. As a result, the width $`\mathrm{\Delta }V_{}=152`$ $`\mu `$V for negative bias has roughly twice the value than that of $`\mathrm{\Delta }V_+=71`$ $`\mu `$V. In previous experiments the resonance width in asymmetric structures has been already studied and a polarity-dependence has been observed, but their explanations referred to polarity-dependent leverage-factors su or internal saturation processes equ .
The width of the current step edge at low temperature reflects the finite-life-time broadening of the zero-dimensional state due to tunneling in and out of the dot. The golden-rule rates for an electron tunneling in and out of the dot are given by $`\mathrm{\Gamma }^+(\omega )=_{r=\text{L},\text{R}}\mathrm{\Gamma }_rf_r(\omega )`$ and $`\mathrm{\Gamma }^{}(\omega )=_{r=\text{L},\text{R}}\mathrm{\Gamma }_r[1f_r(\omega )]`$, respectively. The tunnel-coupling strength is characterized by the constant $`\mathrm{\Gamma }_r=2\pi \nu _r|t_r|^2`$, where $`\nu _r`$ is the density of states in lead $`r`$, and $`t_r`$ is the tunneling amplitude, and $`f_r(\omega )=f(\omega \mu _r)`$ is the Fermi function of lead $`r`$ with electro-chemical potential $`\mu _r`$. We note that the presence of strong Coulomb interaction introduces an asymmetry between tunneling-in and tunneling-out processes: while for an empty dot there are two possibilities to choose the spin state of the incoming electron, the spin state of an electron leaving the dot is fixed. It is, therefore, the combination $`2\mathrm{\Gamma }^+(\omega )+\mathrm{\Gamma }^{}(\omega )`$, i.e.,
$$\underset{r}{}\mathrm{\Gamma }_r\left[1+f_r(\omega )\right],$$
(1)
that determines the finite-life-time broadening. This expression is energy dependent, i.e., it depends on the relative position of the relevant transport channels to the Fermi energy of the leads. At low temperature, the broadening probed by the electrons within the transport window set by the Fermi energies of emitter and collector is either $`2\mathrm{\Gamma }_\text{L}+\mathrm{\Gamma }_\text{R}`$ or $`\mathrm{\Gamma }_\text{L}+2\mathrm{\Gamma }_\text{R}`$, depending on whether the left lead serves as emitter or collector. This qualitatively explains the polarity-dependent step width. It is an interaction effect since in the absence of Coulomb interaction the asymmetry between tunneling-in and tunneling-out processes is lifted by processes involving double occupancy of the dot, and the finite-life-time broadening is energy independent, given by $`\mathrm{\Gamma }_\text{L}+\mathrm{\Gamma }_\text{R}`$.
For a more quantitative analysis we employ a diagrammatic non-equilibrium transport theory. In particular, we make use of the results for the current obtained within the so-called resonant-tunneling approximation for transport through a zero-dimensional state in presence of strong Coulomb interaction such that double occupancy is prohibited. The technique, the approximation scheme, and the steps of the calculation are presented in Ref. koenig . The result for the current is
$$I=\frac{e}{h}_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{2\mathrm{\Gamma }_\text{L}\mathrm{\Gamma }_\text{R}[f_\text{L}(\omega )f_\text{R}(\omega )]}{\left[\omega \epsilon \text{Re}\mathrm{\Sigma }(\omega )\right]^2+\left[\text{Im}\mathrm{\Sigma }(\omega )\right]^2},$$
(2)
where $`\epsilon `$ is the dot level energy \[with the appropriate incorporation of a bias voltage influencing level energy $`\epsilon =\epsilon _0+(1\alpha _+)eV`$ and Fermi-energies of the leads: $`\mu _\text{R}=\text{c}onst.<\epsilon _0,\mu _\text{L}=\mu _\text{L}^0+eV`$ for the non-charging direction\] and the self energy (see Fig. 2)
$$\mathrm{\Sigma }(\omega )=\underset{r}{}\frac{\mathrm{\Gamma }_r}{2\pi }𝑑\omega ^{}\frac{1+f_r(\omega ^{})}{\omega \omega ^{}+i0^+}.$$
(3)
Evaluating the integral leads to $`\text{Re}\mathrm{\Sigma }(\omega )=_r\frac{\mathrm{\Gamma }_r}{2\pi }\left[\mathrm{ln}\left(\frac{\beta E_C}{2\pi }\right)\text{Re}\psi \left(\frac{1}{2}+i\frac{\beta (\omega \mu _r)}{2\pi }\right)\right]`$ and
$$\text{Im}\mathrm{\Sigma }(\omega )=\underset{r}{}\frac{\mathrm{\Gamma }_r}{2}[1+f_r(\omega )].$$
(4)
The real part is weakly dependent on a high-energy cutoff $`E_C`$ given by the smaller of the charging energy for double occupancy or the band width of the leads. In the imaginary part, we recover the structure of the finite-life-time broadening as postulated in the qualitative discussion above. Deep in the Kondo-regime the approximation above is no longer valid. The two spin channels become independent from each other Kondo\_theory and $`\text{Im}\mathrm{\Sigma }(\omega )=_r\mathrm{\Gamma }_r/2`$.
We remark that for the *step height* the well-known (sequential-tunneling) result $`\mathrm{\Delta }I_+=(2e/\mathrm{})\mathrm{\Gamma }_\text{L}\mathrm{\Gamma }_\text{R}/(2\mathrm{\Gamma }_\text{L}+\mathrm{\Gamma }_\text{R})`$ for the non-charging direction \[and $`\text{L}\text{R}`$ for the charging direction\] is reproduced.
The focus of this paper, however, is on the *width of the step edge*. At *low temperature*, Eq. (2) simplifies to
$$I_+=\frac{e}{h}\underset{\mu _\text{R}}{\overset{\mu _\text{L}^0+eV}{}}𝑑\omega \frac{2\mathrm{\Gamma }_\text{L}\mathrm{\Gamma }_\text{R}}{\left[\omega \epsilon \text{Re}\mathrm{\Sigma }(\omega )\right]^2+\left[\mathrm{\Gamma }_\text{L}+\frac{\mathrm{\Gamma }_\text{R}}{2}\right]^2}$$
(5)
with $`\epsilon =\epsilon _0+(1\alpha _+)eV`$ and $`\alpha _+`$ being the leverage factor for positive bias \[and $`\alpha _{}=1\alpha _+`$ being the leverage-factor for negative bias\], denoting the voltage drop over the left barrier. Neglecting the real part of the self-energy for the moment, we find that the differential conductance as a function of $`V`$ is a Lorentzian with a FWHM of:
$$\alpha _+e\mathrm{\Delta }V_+=2\mathrm{\Gamma }_\text{L}+\mathrm{\Gamma }_\text{R}\text{and}\alpha _{}e\mathrm{\Delta }V_{}=\mathrm{\Gamma }_\text{L}+2\mathrm{\Gamma }_\text{R}.$$
For $`\alpha _+\alpha _{}`$ and strongly asymmetric tunnel-coupling strengths, $`\mathrm{\Gamma }_\text{L}\mathrm{\Gamma }_\text{R}`$, we get the relation $`\mathrm{\Delta }V^{}2\mathrm{\Delta }V^+`$.
At *high temperature* we find for the FWHM $`e\mathrm{\Delta }V_\pm =3.525k_\mathrm{B}T/\alpha _\pm +e\mathrm{\Delta }V_\pm ^0`$, i.e. the temperature broadening of the Fermi-function. As predicted from sequential-tunneling theory beenakker the width increases linearly with temperature. The constant term $`\mathrm{\Delta }V_\pm ^0`$ is of the order of $`\mathrm{\Gamma }`$, and, in general, also polarity dependent remark1 .
For a detailed comparison between theory and experiment we need first to determine the system parameters. The factor $`\alpha `$ determining the bare level shift with bias voltage is gained from the linear high-temperature dependence of $`\mathrm{\Delta }V_\pm (T)`$ as $`\alpha _+=0.53`$ and $`\alpha _{}=0.5`$ so that $`\alpha _++\alpha _{}1.`$ The coupling constants $`\mathrm{\Gamma }_{\text{L}/\text{R}}`$ could, in principle, both be determined from the current steps $`\mathrm{\Delta }I_\pm `$. For strong asymmetry, however, as is the case here, the maximum current is limited entirely by the bottleneck of the smaller coupling $`\mathrm{\Gamma }_\text{L}`$. As a consequence, the step heights only fix $`\mathrm{\Gamma }_\text{L}=\mathrm{\Delta }I^+\mathrm{}/(2e)=2\mathrm{\Delta }I^{}\mathrm{}/(2e)=0.64`$ $`\mu `$eV.
Considering the results derived above for zero temperature FWHM, we note that this width is essentially the sum of different couplings and, hence, is dominated by the larger coupling $`\mathrm{\Gamma }_\text{R}`$. Therefore we can find $`\mathrm{\Gamma }_\text{R}=(1\alpha ^+)e\mathrm{\Delta }V^+(1\alpha ^{})\mathrm{\Delta }V^{}/2.`$ In fact, we gain better accuracy by fitting the full $`dI/dV`$–peak for charging polarity and lowest temperature to pinpoint $`\mathrm{\Gamma }_\text{R}=40`$ $`\mu `$eV. The high-energy cut-off $`E_C=30`$ meV is given by the bandwidth, i.e. the value of the Fermi energy, being of the same magnitude as the charging energy.
Figure 3 shows the comparison between experimental data and theoretical calculations for the differential conductance. For the charging direction (negative-bias), Fig. 3(a), we find a good agreement between experiment and theory for both the resonance width and amplitude. In the non-charging direction (positive bias), Fig. 3(b), the peak width is reduced by about a factor of two for both experiment and theory. The experimental data show some extra features. First there is an oscillatory fine structure on the positive-voltage side of the resonance attributed to the fluctuations of the local density of states of the emitter, see, e.g., Ref. ldos . Moreover, we observe an enhanced resonance amplitude as compared to the theoretical calculation. This effect may be related to an additional many-body phenomenon at the Fermi-edge int .
The temperature dependence of the resonance width in the broad range between $`20`$ mK and $`400`$ mK is shown in Fig. 4(a). The polarity dependence of the width is clearly visible. At low temperature experimental and theoretical data nicely match. For temperatures above 200 mK the width increases linearly with $`e\alpha _\pm \mathrm{\Delta }V_\pm =3.525k_\mathrm{B}T+𝒪(\mathrm{\Gamma })`$. This experimental result is again adequately reproduced by theory.
Further support for our explanation of the polarity dependence as an interaction effect is given by the magnetic-field dependence. The data, shown in Fig. 4(b), reveal that polarity dependence is strongly reduced by a Zeeman-splitting (on a scale of $`1`$ T, which is of the order of the magnitude of the coupling). In this case, only one spin state contributes to transport, and the system is equivalent to a noninteracting one, for which theory predicts $`\text{Im}\mathrm{\Sigma }(\omega )=_r\mathrm{\Gamma }_r/2`$ and, thus, a low-temperature width $`\alpha _\pm e\mathrm{\Delta }V_\pm =\mathrm{\Gamma }_\text{L}+\mathrm{\Gamma }_\text{R}`$ independent of the polarity. This is in accordance with the trend seen in the experimental data.
The discrepancies remaining between experimental data and theoretical simulation are consistent with the expected range of accuracy. The main source of experimental errors are systematic fluctuations in the current, in particular the local DOS fluctuations prominent in non-charging direction which become even more pronounced at high temperatures. This artefact particularly limits the precision to which we can determine the current step height $`\mathrm{\Delta }I_+`$.
We study the width of the current step of single-electron tunneling through a zero-dimensional state in a double-barrier resonant tunneling structure and find that it is strongly polarity dependent. This is interpreted as a clear signature of Coulomb interaction. The latter introduces an asymmetry between possible tunneling-in and tunneling-out processes, which gives rise to an energy-dependent finite-life-time broadening. The bias-polarity dependent step width, observed in our experiment, where strengths of the tunnel coupling between dot and the two leads were highly asymmetric, can be simulated with reasonable agreement by applying a diagrammatic non-equilibrium transport theory for interacting quantum dots.
We acknowledge sample growth by A. Förster and H. Lüth and discussions with H. Schoeller. Financial assistance was granted by BMBF and by DFG via SFB491 and GRK726.
|
warning/0506/astro-ph0506672.html
|
ar5iv
|
text
|
# Stellar Cluster Fiducial Sequences with the Advanced Camera for Surveys1
## 1. Introduction
The Advanced Camera for Surveys (ACS; Ford et al. 1998) on the Hubble Space Telescope (HST) provides far greater broad-band optical sensitivity than its predecessor, the Wide Field Planetary Camera 2 (WFPC2). Two of the most widely used filters on the ACS Wide Field Camera (WFC) are the F606W (broad $`V`$) and F814W ($`I`$); the former provides three times the throughput of the analogous filter on WFPC2, while the latter provides five times the throughput of the corresponding filter on WFPC2. Although the F555W on ACS more closely approximates the Johnson $`V`$ bandpass, the F606W has far more grasp, making it the filter of choice for deep imaging programs.
This enormous advance in sensitivity allows the HST to resolve the old main sequence in populations out to the edge of the Local Group, given a reasonable investment of exposure time ($`100`$ orbits) and a sufficiently sparse field ($`\mu _V26`$ mag arcsec<sup>-2</sup>). In order to measure the star formation history in the Andromeda halo, we used this capability to obtain extremely deep images of Andromeda in a field 51 from the nucleus on the southeast minor axis (Brown et al. 2003). In addition to
<sup>1</sup>Based on observations made with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS 5-26555. These observations are associated with proposals 9453 and 10265.
these deep images, our program included brief exposures of five Galactic globular clusters, using the same camera and filters: NGC 6341, NGC 6752, NGC 104, NGC 5927, and NGC 6528. These cluster observations serve two purposes: they provide empirical isochrones for old simple stellar populations spanning a wide range of metallicity, and they serve as calibrators for the transformation of theoretical isochrones to the ACS bandpasses. Subsequent to our first ACS program, we obtained additional observations of the Andromeda tidal stream (discovered by Ibata et al. 2001) and outer disk; this second program included an observation of the old open cluster NGC 6791 in order to expand the range of metallicities sampled in the empirical isochrones. Because these data would be useful references for other ACS programs, we tabulate here, for each cluster, the ridge lines tracing the main sequence (MS), subgiant branch (SGB), and red giant branch (RGB) stars, plus the horizontal branch (HB) loci. In addition, we describe our transformation of the Victoria-Regina isochrones (hereafter VRI; Bergbusch & VandenBerg 2001) to the ACS bandpasses. The VRI color-temperature relations have been revised recently (VandenBerg & Clem 2003), and the VRI grid has been extended to higher metallicities and lower ages (VandenBerg, Bergbusch, & Dowler 2005). The isochrone interpolation code and the models should be available from the Canadian Astronomy Data Centre later this year.
Table 1:Parameters<sup>a</sup> of Galactic clusters observed with ACS
| | exposures | | | | | |
| --- | --- | --- | --- | --- | --- | --- |
| | F606W | F814W | $`(mM)_V`$ | $`E(BV)`$ | | age |
| Name | (sec) | (sec) | (mag) | (mag) | \[Fe/H\] | (Gyr) |
| NGC 6341 (M92) | 0.5079,5,90 | 0.5079,6,100 | 14.60 | 0.023 | $`2.14`$ | 14.5 |
| NGC 6752 | 0.5079,4,40 | 0.5079,4,45 | 13.17 | 0.055 | $`1.54`$ | 14.5 |
| NGC 104 (47 Tuc) | 0.5079,6,70 | 0.5079,5.5,72 | 13.27 | 0.024 | $`0.70`$ | 12.5 |
| NGC 5927 | 2,30,500 | 0.6934,15,340 | 15.85 | 0.42 | $`0.37`$ | 12.5 |
| NGC 6528 | 4,50,450 | 1,20,350 | 16.31 | 0.55 | $`+0.00`$ | 12.5 |
| NGC 6791 | 0.5079,5,50 | 0.5079,5,50 | 13.50 | 0.14 | $`+0.30`$ | 9.0 |
<sup>a</sup>See §4 for discussion and references for these parameters.
## 2. Observations and Data Reduction
The Galactic clusters in our programs are listed in Table 1. The globular clusters were observed as part of HST program GO-9453, while NGC 6791 was observed as part of HST program GO-10265. Although the enormous sensitivity of ACS has enabled great strides in the deep imaging of faint targets, ironically, the camera sensitivity makes it challenging to observe relatively bright star clusters in our own Galaxy. We observed each cluster for one orbit, staggering the exposure times by an order of magnitude to increase the dynamic range. Due to the exposure overheads (CCD readout, buffer dumps, etc.), only six ACS/WFC images can be taken in a single orbit outside of the HST continuous viewing zone, allowing three images in each bandpass (F606W and F814W). For the three relatively distant globular clusters (NGC 6341, NGC 5927, and NGC 6528), we roughly centered the WFC on the cluster core, maximizing the number of stars in the samples. For the two relatively nearby globular clusters (NGC 104 and NGC 6752), we offset the ACS images from the cluster core, to include regions of lower background and higher signal-to-noise ratio for the photometry of the relatively faint main sequence stars. The nearby open cluster NGC 6791 is relatively sparse and subtends an area much larger than the ACS field of view, so we centered the WFC roughly in the cluster core.
We observed these clusters to create empirical isochrones and to calibrate the transformation of theoretical isochrones to the ACS bandpasses. The observing strategy was designed to efficiently achieve these purposes, but the data are not optimal for detailed studies of the clusters themselves. In particular, the exposures are brief (sometimes less than 1 sec), they are not split into two subexposures, nor are they dithered, thus foregoing the traditional method of cosmic ray rejection and precluding a better sampled point spread function (PSF). However, by comparing each set of cluster images, even with their different exposure times, we were able to create adequate cosmic ray masks. Note that, in 2003, the effective exposure times for commanded exposures of less than 1 sec were remeasured and updated in the HST calibration pipeline, such that a 0.5 sec exposure is really 0.5079 sec, and a 0.7 sec exposure is really 0.6934 sec (Gilliland & Hartig 2003). Because the images were not dithered, no large changes in plate scale were applied when the masked images were registered and coadded, as often done to better sample the PSF.
All of the images of a given cluster in a given bandpass were coadded using the DRIZZLE package (Fruchter & Hook 2002), with masking of cosmic rays, saturated pixels, and bad pixels. Although no gross changes to the plate scale were applied, this step does correct for the geometric distortion, small temporal changes in plate scale due to velocity aberration and telescope breathing, and a small difference in plate scale
between the two bandpasses. Because the images are filled with thousands of stars, we used the stellar positions in each exposure to accurately determine the shifts and scales needed to register the exposures. The DRIZZLE package also provides software for the masking of cosmic rays. We tuned the software parameters to aggressively mask cosmic rays without masking the cores of unsaturated bright stars; these masks were then confirmed by visual inspection. In the medium and long exposures, the mask for all saturated pixels (either due to a bright star or cosmic ray) was enlarged with a border of 7 pixels to mask the bleeding of saturated pixels into neighboring pixels; the short exposure suffers from very few cosmic-ray hits and very little bleeding of saturated pixels.
The five globular clusters were all observed using a CCD gain of 1 e<sup>-</sup> per data number (DN), in order to match exactly the observing mode used in our deep imaging of the Andromeda halo. The observations were planned before the installation of ACS, when it was unclear how much uncertainty there would be in the relative gain corrections on the camera. Subsequent calibration programs have accurately measured the gain correction, and showed that a gain of 2 e<sup>-</sup>/DN offers significant advantages when observing bright stars, with little increase in quantization noise. With a gain of 1 e<sup>-</sup>/DN, the analog-to-digital converter (ADC) saturates at 65,535 DN, prior to the CCD full well of 84,700 e<sup>-</sup>, making it difficult to recover the flux from a saturated star. With a gain of 2 e<sup>-</sup>/DN, fluxes can be measured not only to the CCD full well, but beyond, via sampling the neighboring pixels where charge “bleeds” from the saturated pixel (i.e., e<sup>-</sup> are conserved in the conversion to DN). Because a gain of 2 e<sup>-</sup>/DN can significantly extend the dynamic range of the CCD, we observed NGC 6791 with a gain of 2 e<sup>-</sup>/DN. As it turns out, the scarcity of giant stars in our chosen field made this choice somewhat moot.
The undersampled images do not lend themselves to accurate PSF-fitting; therefore, we obtained aperture photometry using the DAOPHOT package (Stetson 1987). We experimented with a range of aperture sizes, and ultimately chose two different apertures for the unsaturated and saturated stars. For the unsaturated stars, we chose a circular aperture of radius 2.5 pixels (0.125$`\mathrm{}`$) and a sky annulus of radii 7–18 pixels. With a gain of 1 e<sup>-</sup>/DN, saturated stars cannot be accurately photometered in a circular aperture because charge is not conserved, so instead we measured the flux in a circular annulus of radii 2.5–3.5 pixels, again with a 7–18 pixel sky annulus. We discarded from the catalog saturated stars that bled beyond a 2.5 pixel radius, or if they were adversely affected by bad pixels or were blended with bright neighbors. The photometry of the saturated stars was normalized to the same zeropoint as the unsaturated stars by comparing photometry of the unsaturated stars in the circular aperture and the annular aperture. The photometric catalog was then corrected to true apparent magnitudes using TinyTim models of the HST PSF (Krist 1995) and observations of the standard star EGGR 102 (a $`V=12.8`$ mag DA white dwarf) in the same filters, with agreement at the 1% level.
Charge transfer inefficiency (CTI) can be a problem for aging large-format CCDs in the space radiation environment, causing stars to appear fainter than they actually are. The ACS detector consists of two chips, $`4144\times 2068`$ pixels each, separated by a small horizontal gap. Stars which fall closer to the gap undergo more parallel transfers when the detector is read, and thus suffer from more charge loss due to CTI. The CTI correction is approximately linear with the position of a star relative to the gap, and approximately linear with the age of the detector. The correction is larger for faint stars and smaller when there is a significant background. Because the globular clusters were all observed shortly after the ACS launch, the CTI correction for isolated stars would be negligible ($`0.01`$ mag), even at the faint end of the ridge lines we define for each cluster color-magnitude diagram (CMD), as discussed below. Because the globular cluster images are significantly crowded, the CTI is, in practice, even smaller than it would be for isolated stars. We thus applied no CTI correction to our globular cluster photometry. Unlike the globular clusters, NGC 6791 was observed 2.6 years after launch, and its ACS images are relatively sparse; because the CTI correction is small but not completely negligible ($`0.01`$ mag near the faint limit at chip center), we applied a CTI correction to our NGC 6791 photometry, using the algorithm of Riess & Mack (2005).
Our photometry is in the STMAG system: $`m=2.5\times `$ log$`{}_{10}{}^{}f_{\lambda }^{}21.1`$ mag, where $`f_\lambda =`$ e$`{}_{}{}^{}\times \mathrm{PHOTFLAM}/\mathrm{EXPTIME}`$, EXPTIME is the exposure time, and PHOTFLAM is $`7.906\times 10^{20}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup> / (e<sup>-</sup> s<sup>-1</sup>) for the F606W filter and $`7.072\times 10^{20}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup> / (e<sup>-</sup> s<sup>-1</sup>) for the F814W filter. The STMAG system is a convenient system because it is referenced to an unambiguous flat $`f_\lambda `$ spectrum; an object with $`f_\lambda =3.63\times 10^9`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup> has a magnitude of 0 in every filter. Another convenient and unambiguous system that is widely used is the ABMAG system: $`m=2.5\times `$ log$`{}_{10}{}^{}f_{\nu }^{}48.6`$ mag; it is referenced to a flat $`f_\nu `$ spectrum, such that an object with $`f_\nu =3.63\times 10^{20}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> has a magnitude of 0 in every filter. It is thus trivial and unambiguous to convert any of the data presented herein from STMAG to ABMAG: for F606W, ABMAG = STMAG $`0.169`$ mag, and for F814W, ABMAG = STMAG $`0.840`$ mag. Although our photometry could be transformed to ground magnitude systems (e.g., Johnson $`V`$) for comparison to theoretical isochrones as well as other data in the literature, such transformations always introduce significant systematic errors (see Sirianni et al. 2005). Instead of converting HST data to ground bandpasses so that they can be compared to models in ground bandpasses, it is preferable to produce models in one of the HST instrument magnitude systems, in either STMAG or ABMAG.
The CMDs of each cluster are shown in Figures 1 through 6, along with the ridge line spanning the MS, SGB, and RGB. Each ridge line was created by defining regions along the MS-SGB-RGB locus, and taking the median color and magnitude in each region. The size of the region was varied along the locus to allow clipping of outliers while including most of the stars appropriate for defining the ridge line. Larger regions were defined in parts of the CMD where the locus of stars is relatively linear, where the photometric scatter is significant (at the faint end), and where stars are scarce (at the bright end). Smaller regions were needed where the locus curves significantly (otherwise the ridge line would smooth over these features), but fortunately this is also where the photometric errors are relatively reasonable and the CMD is relatively well-populated (between the turnoff and the base of the RGB). Near the tip of the RGB, where the CMD is very sparse, a mean was used instead of a median if less than 5 stars fell in the region. Figures 1–6 show, at representative locations, horizontal and vertical bars spanning the sizes of these regions used for defining the ridge lines. We have also highlighted the HB locus in each cluster; these stars were simply selected by eye from the obvious overdensity of points in the vicinity of the HB. The observed fiducials (ridge line and HB locus) are difficult to determine accurately in NGC 6528, because the cluster suffers from high, spatially variable reddening (Heitsch & Richtler 1999), and in NGC 6791, due to the scarcity of stars. A dotted line in each figure indicates where stars can become saturated if they are well-centered on an ACS pixel. Note that when these ridge lines were presented previously (Brown 2005; Brown et al. 2003, 2004), the ridge lines were not shown above the saturation point, due to uncertainties in the correction for saturated stars with a gain of 1 e<sup>-</sup>/DN. Although we have taken care to correct the saturated stars, the correction is not as precise as the correction that can be done with a gain of 2 e<sup>-</sup>/DN, because one is effectively extrapolating from the wings of the PSF. The correction is also quite large at the tip of the RGB; for example, at m$`{}_{F814W}{}^{}=10.5`$ mag in the NGC 104 CMD, a star clipped at the limit of the ADC can exhibit $``$25% less DN than the number of e<sup>-</sup> actually generated on the ACS detector.
Although we have shown the globular cluster fiducials in several papers and conference proceedings, we have provided the tabular data to ACS observers upon request only. Publication of the tables was delayed until we obtained our NGC 6791 observations so that we could present a complete cluster dataset that extended to super-solar metallicity. It is worth noting that, in the interim, a separate group published ridge lines for these globular cluster data (Bedin et al. 2005). Besides our addition of NGC 6791, there are several differences between their work and that presented here, related to the photometric methods, handling of saturated stars, and reddening. More significantly, Bedin et al. (2005) used a magnitude system referenced to a spectrum of Vega that is older than the recent one of Bohlin & Gilliland (2004). They also compared the globular cluster data to a different isochrone set (Pietrinferni et al. 2004), and only did so for three of the globular clusters (NGC 6341, NGC 6752, and NGC 104), which is understandable, given the larger uncertainties associated with the clusters at higher metallicity. The work we present here will aid in the interpretation of past and future papers in our study of the star formation history in the Andromeda halo (HST program GO-9453), disk, and tidal stream (HST program GO-10265). It will also facilitate the analyses of other groups working with CMDs derived from the most popular ACS bandpasses (F606W and F814W). These tables are given in the following sections, where we also discuss how the ridge lines can be transformed to account for different levels of reddening.
## 3. Transformation of the Victoria-Regina Isochrones
We briefly described our transformation of the VRI in Brown et al. (2004) but elaborate here. As distributed, the VRI provide, for a simple stellar population, the physical parameters (effective temperature, surface gravity, luminosity, and mass)
Fig. 1– Left panel: The CMD for NGC 6341, along with its ridge line (grey curve). The HB locus is highlighted in grey. A dashed line shows where stars can become saturated if well-centered on an ACS pixel. At representative points along the ridge line, the size of the regions used to define the ridge line are shown by horizontal and vertical lines. Right panel: The CMD for NGC 6341, along with an isochrone and ZAHB sequence, shown with the empirical color correction (solid curve) and without this correction (dashed curve). Note that the axes are not the same as those in the left panel, to better show the level of agreement between the models and data. The ZAHB sequence shows better agreement with the data when no empirical color correction is applied.
and the observed magnitudes in ground bandpasses ($`B`$, $`V`$, $`R`$, and $`I`$). The adopted color-temperature relations for these ground bandpasses are the ones described by VandenBerg & Clem (2003)<sup>2</sup> and, as shown by VandenBerg (2000), they yield synthetic CMDs that agree very well with the observed CMDs for a number of well-studied Galactic clusters. For this reason, Brown et al. (2003) used the synthetic spectra of Lejeune, Cuisinier, & Buser (1997) to calculate a differential transformation between the ground bandpasses and the corresponding ACS bandpasses, $`Vm_{F606W}`$ and $`Im_{F814W}`$, for all points on the VRI, and then applied those differences to the isochrones, thus producing a set of isochrones in the ACS bandpasses. Comparison to the ACS observations of Galactic
globular clusters showed that a small empirical color correction ($`0.05`$ mag) was required to force agreement between the transformed isochrones and the observations. In a subsequent analysis of the Andromeda globular cluster SKHB-312, Brown et al. (2004) found that a direct transformation from physical parameters to ACS bandpasses (instead of the
<sup>2</sup>These transformations are nearly identical with those reported by Bell & Gustafsson (1989) for turnoff stars. At cooler temperatures, redward adjustments to the synthetic colors were applied in order that the predicted MS and RGB slopes agreed well with those observed. The color-temperature relations derived from Kurucz model atmospheres (e.g., see Castelli 1999) were used for stars hotter than $`7000`$ K.
Fig. 2– The same as Figure 1, but for NGC 6752.
differential transformation above) was slightly preferable, when done in conjunction with a more extensive grid of synthetic spectra (Castelli & Kurucz 2003). The Castelli & Kurucz (2003) grid provides spectra over a wide range of metallicity, with and without alpha-enhancement, so that the chemical compositions in the spectra can be well-matched to those in the isochrones (Lejeune et al. 1997 provide only scaled-solar models). Isochrones transformed in this new direct manner, when compared to the ACS Galactic cluster data, still demonstrated the need for an empirical color correction, but it was somewhat smaller and the functional form was simpler than that required in Brown et al. (2003). In the end, the result was very similar to that obtained by Brown et al. (2003), because in both cases the isochrones were forced to agree with the same set of observational data.
To transform the isochrones to the ACS bandpasses, we first interpolate the synthetic spectra grid of Castelli & Kurucz (2003) in metallicity, effective temperature, and surface gravity
to produce a spectrum at each point on the isochrone, and then redden that spectrum using the curve of Fitzpatrick (1999). Although extinction is often handled in the literature by a “reddening vector” that is constant over the full range of a CMD, in reality the reddening produces a change in flux in each bandpass that depends upon the spectral energy distribution of the star; thus, it is more accurate to calculate the reddening at each point on the isochrone. The reddened spectrum is then converted from an energy spectrum (erg cm<sup>-2</sup> s<sup>-1</sup> Å<sup>-1</sup>) to a photon spectrum (photon cm<sup>-2</sup> s<sup>-1</sup> Å<sup>-1</sup>), multiplied by the throughput of each bandpass, and integrated over wavelength, to produce the expected count rate on the detector (e<sup>-</sup> s<sup>-1</sup>). Finally, this count rate is converted to STMAG, using the same zeropoints given in §2. Because the synthetic spectra grid of Castelli & Kurucz (2003) does not extend below 3500 K, the isochrone is truncated for any points below this temperature (near the RGB tip).
Note that the VRI do not include He diffusion, which would
Fig. 3– The same as Figure 1, but for NGC 104.
decrease their ages at a given turnoff luminosity by $`10`$%, thus avoiding discrepancies with the age of the Universe (VandenBerg et al. 2002). Although the ages of isochrones with He diffusion are likely more accurate, if diffusion is allowed to act efficiently on other elements in the surface layers, such models show significant discrepancies when compared to observed CMDs. For example, they fail to explain either the Li abundance versus effective temperature relationship obeyed by field Population II dwarfs (see Richard et al. 2002 and references therein) or the lack of any detectable difference in the derived abundances between globular cluster stars at the turnoff and on the lower giant branch (Gratton et al. 2001; James et al. 2004). Apparently, there must be some competing processes at work (e.g., turbulence at the base of the convective envelope, as invoked by Richard et al. 2002) that reduces the efficacy of diffusion in the surface layers of metal-deficient stars. However, the age effect is mainly due to the settling of He in the stellar
core, and presumably this still occurs at close to expected rates. Although the best available models to use in comparisons with stellar data are arguably those by Richard et al. (2002), which take diffusion and turbulence into account, they have (so far) been computed for only a few values of \[Fe/H\] and only as far as the lower giant branch. In view of these considerations, it seems advisable to fit non-diffusive models to observed CMDs, and to reduce the ages so obtained by $`10`$% in order to provide the best estimates of cluster ages.
## 4. Cluster – Isochrone Comparison
The comparison of transformed isochrones to observed clusters is not completely straightforward, because for even the most well-studied clusters, there are still significant uncertainties in their parameters: age, chemical composition, reddening, and distance. We evaluated the various values for these parameters in the literature, and settled upon those (Table 1)
Fig. 4– The same as Figure 1, but for NGC 5927.
that minimized the empirical correction of the transformed isochrones and allowed a uniform correction for all clusters. If the same empirical correction can be applied to the transformed isochrones associated with each of these clusters, then this systematic offset can be reasonably attributed to systematic errors in the isochrones, synthetic spectra, and calibration of the bandpasses. In general, our values for the distance, \[Fe/H\], and reddening come from the literature, while the age for each cluster is that which provides the best agreement between the ACS data and the isochrone in the vicinity of the turnoff. As distributed, the VRI interpolation code produces isochrones over a continuous range of age at discrete metallicities, but the sampling of the metallicity grid is fine enough that isochrones can be matched to clusters well within 0.1 dex. For each of the globular clusters, we compare to an isochrone with \[$`\alpha `$/Fe\] = +0.3 and a metallicity closest to that in Table 1 (see Maraston et al. 2003 for observational evidence of $`\alpha `$-enhancement in Galactic globular
clusters over the full range of metallicity). For the open cluster NGC 6791, we interpolate the two isochrones nearest in metallicity with \[$`\alpha `$/Fe\] = 0. Note that, to first order, isochrones without alpha-enhancement look similar to isochrones with alpha-enhancement at lower metallicity; this approximation is sometimes used in the literature, although it is more accurate at lower metallicities than at metallicities near the solar value. Before comparing these isochrones to the observed CMDs, we briefly discuss the parameters we adopted for each cluster.
NGC 6341. For NGC 6341, we assumed the same distance, \[Fe/H\], and reddening as used in Brown et al. (2003), but increased the age from 14 Gyr to 14.5 Gyr, which fits better under the new direct transformation method. VandenBerg & Clem (2003) found good agreement between the VRI and a $`BV`$ CMD of M92 (Stetson & Harris 1988) when assuming \[Fe/H\] = $`2.14`$, $`E(BV)=0.023`$ mag, $`(mM)_V=14.60`$ mag, and an age of 15 Gyr. This value of \[Fe/H\] is very
Fig. 5– The same as Figure 1, but for NGC 6528.
close to that found spectroscopically by Zinn & West (1984; $`2.24\pm 0.08`$) and Carretta & Gratton (1997; $`2.16\pm 0.02`$). For our comparison between the ACS data and VRI (Figure 1), we used the isochrones at \[Fe/H\] = $`2.14`$. The extinction comes from the dust maps of Schlegel, Finkbeiner, & Davis (1998), for a sight-line through the Galaxy, which is reasonable for a distant halo cluster. The distance comes from Grundahl et al. (2000), based upon the metal-poor subgiant HD 140283; Grundahl et al. (2000) also found an age of 14.5 Gyr when performing a distance-independent fit of the VRI to Str$`\ddot{\mathrm{o}}`$mgren photometry of NGC 6341.
NGC 6752. For NGC 6752, we assumed the same distance, \[Fe/H\], and reddening as used in Brown et al. (2003), but again increased the age from 14 Gyr to 14.5 Gyr, which fits better under the new direct transformation method. VandenBerg (2000) found good agreement between the VRI and the $`BV`$ CMD of NGC 6752 (Penny & Dickens 1986) when
assuming \[Fe/H\] = $`1.54`$, which is the same value found spectroscopically by Zinn & West (1984; $`1.54\pm 0.09`$). For our comparison between the ACS data and VRI (Figure 2), we used the isochrones at \[Fe/H\] = $`1.54`$. Renzini et al. (1996) found a true distance modulus of $`(mM)_0=13.05`$ mag by comparison of the cluster white dwarf sequence to local white dwarfs with accurate parallaxes; they assumed an extinction of $`E(BV)=0.04`$ mag (Penny & Dickens 1986), and thus found an apparent distance modulus of $`(mM)_V=13.17`$. We have updated the extinction to $`E(BV)=0.055`$ mag, using the dust maps of Schlegel et al. (1998). Our adopted age falls in the range found in the recent literature; Renzini et al. (1996) found an age of 15.5 Gyr using non-diffusive models and 14.5 Gyr using diffusive models, while VandenBerg (2000) found an age of 12.5 Gyr, on the assumption of a somewhat larger distance modulus, using non-diffusive models.
NGC 104. For NGC 104, we assumed the same parameters
Fig. 6– The same as Figure 1, but for NGC 6791.
as those used in Brown et al. (2004). The metallicity comes from the high-resolution spectroscopy of Carretta & Gratton (1997; $`0.70\pm 0.03`$), which agrees well with the \[Fe/H\]<sub>II</sub> fits of Kraft & Ivans (2003; $`0.70\pm 0.09`$) and with the spectroscopic fits of Zinn & West (1984; $`0.71\pm 0.08`$). For our comparison between the ACS data and VRI (Figure 3), we used isochrones at \[Fe/H\] = $`0.705`$. From fits to the cluster white dwarf sequence, Zoccali et al. (2001) found an apparent distance modulus of $`(mM)_V=13.27\pm 0.14`$ mag. Gratton et al. (2003) present several different determinations for $`E(BV)`$, with an average value of 0.024 mag. Our adopted age of 12.5 Gyr is midway between that found by VandenBerg & Clem (2003; 12 Gyr), who adopted a slightly larger distance and smaller metal abundance, and Zoccali et al. (2001; 12.9 Gyr with diffusive models and 13.5 Gyr with non-diffusive models).
NGC 5927. For NGC 5927, we assumed the same parameters as those used in Brown et al. (2004). The metallicity is
that adopted by Harris (1996), and represents a combination of various spectroscopic measurements in the literature (Zinn 1985; Armandroff & Zinn 1988; François 1991). For our comparison between the ACS data and VRI (Figure 4), we used isochrones at \[Fe/H\] = $`0.397`$. Brown et al. (2004) adopted values for distance and reddening that were slightly different from those of Harris (1996), who gives $`(mM)_V=15.81`$ mag and $`E(BV)=0.45`$ mag, but the Brown et al. (2004) values (15.85 mag and 0.42 mag, respectively) are well within the uncertainties, given the high, spatially-variable reddening (Heitsch & Richtler 1999). Our adopted age of 12.5 Gyr is also well within the wide range of ages in the literature for this cluster (see Feltzing & Gilmore 2000 for a summary of recent estimates for the age, metallicity, distance, and reddening).
NGC 6528. For NGC 6528, we assumed the same distance, reddening, and age as used in Brown et al. (2003), but we assumed a slightly higher metallicity, \[Fe/H\] = 0.0, which is still well within the range of spectroscopic metallicities in the literature (Momany et al. 2003). For our comparison between the ACS data and VRI (Figure 5), we used isochrones at \[Fe/H\] = 0. As with the case of NGC 5927, NGC 6528 suffers from high, spatially-variable reddening (Heitsch & Richtler 1999 and references therein), resulting in considerable uncertainties in its fundamental parameters (see Feltzing & Gilmore 2000). Brown et al. (2003) found that if NGC 6528 is assumed to lie at an apparent distance modulus of $`(mM)_V=16.15`$ mag (Momany et al. 2003), the cluster, when shifted to the Andromeda distance and reddening, is too faint by 0.16 mag, implying $`(mM)_V=16.31`$ mag, which is the distance we have assumed here. We assumed the same extinction, $`E(BV)=0.55`$ mag, as found by Momany et al. (2003) in their analysis of optical and near-infrared imaging. Our adopted age of 12.5 Gyr is nearly the same as that found by Momany et al. (2003; 12.6 Gyr). Because the metallicity and reddening are very uncertain for NGC 6528, it is difficult to evaluate the empirical corrections needed to achieve agreement between the isochrone and observed CMD, but those corrections are the same as those applied to the other clusters.
NGC 6791. For NGC 6791, the Carney et al. (2005) analysis of their infrared HB photometry implies $`(mM)_V=13.50`$ mag and $`E(BV)=0.14`$ mag, and we have assumed those values here. The metallicity for NGC 6791 is quite uncertain (see Taylor 2001 for an extensive review), but Carney et al. (2005) find acceptable fits to their infrared photometry if they assume \[Fe/H\] = +0.5 (and an age of 7.5 Gyr) or \[Fe/H\] = +0.3 (and and age of 9 Gyr); we find good agreement with the latter pair of parameters. This metallicity falls midway between two metallicities in the VRI, so we interpolated between isochrones at \[Fe/H\] = +0.23 and \[Fe/H\] =+.37; the comparison between the isochrone and ACS data is shown in Figure 6.
The comparisons between the data and isochrones are shown in Figures 1–6 (right-hand panels). The dashed line shows the transformed isochrone when the above steps are performed, but no empirical correction is applied. The isochrones are generally too red over much of the CMD, except for the bottom of the main sequence. Furthermore, the color difference between the turnoff and the base of the RGB is larger than observed, which may indicate, e.g., a preference either for models that treat diffusive processes and turbulence, or for models with higher oxygen abundances, as both possibilities are ways of reducing the turnoff temperatures without affecting the location of the RGB significantly. While adopting older isochrones could also alleviate this problem, the required ages are implausibly large. Because an accurate determination of the age and metallicity in a CMD generally comes from a fit of the upper MS, turnoff, SGB, and RGB, we apply a small empirical color correction, specified as a first-order polynomial:
$`(m_{F606W}m_{F814W})_{new}=0.92\times (m_{F606W}m_{F814W})_{old}0.06.`$
This correction gives good agreement from the upper MS to a point on the RGB slightly above the HB, but then the upper RGB increasingly deviates to the blue. Thus, we change the correction for stars at log $`g<2`$ to:
$`(m_{F606W}m_{F814W})_{new}=0.92\times (m_{F606W}m_{F814W})_{old}0.06+0.05\times (2`$ log $`g)`$.
This correction gives good agreement to the end of the isochrone (which, as stated above, can end before the RGB tip, due to the 3500 K limit in the grid of synthetic spectra). In general, the transformed isochrones with the empirical correction (shown as a solid curve in the right-hand panels of Figures 1–6) agree at the $`0.02`$ mag level with the observed ridge lines from a point 1.5 mag below the MS turnoff through the upper RGB, which thus allows the isochrones to be used for fitting the age and metallicity in ACS CMDs. Over this range, the empirical correction is small ($`0.05`$ mag) but not insignificant. In the two most metal-rich clusters (NGC 6528 and NGC 6791), the agreement on the upper RGB might not be as good as that in the more metal-poor clusters, but this is difficult to evaluate, given the scarcity of stars in NGC 6791 and the high differential reddening in NGC 6528. We make no attempt to improve the agreement between the isochrones and data on the lower MS, because there is too much variation in the disagreement, from cluster to cluster, to suggest a suitable correction. Furthermore, photometry on the lower MS is not required to measure ages in CMDs (and it is beyond the reach of current instrumentation for stellar populations beyond our Galaxy and its satellites).
We also show models of the zero-age HB (ZAHB) for each cluster, again transformed using the method above. For consistency with the isochrones, the ZAHB loci were taken from VandenBerg et al. (2000); their $`M_V`$ values in the center of the instability strip agree well with those derived from studies of RR Lyrae stars in Galactic globular clusters spanning a wide range in metallicity (see De Santis & Cassisi 1999; Cacciari et al. 2005). The ZAHB should trace the bottom of the HB locus, because luminosity increases during HB evolution. In general, the ZAHB models show better agreement with the data when no color correction is applied (dashed) than when the color correction is applied (solid). Although we do not use HB models in our fits of the Andromeda star formation history, it appears that a transformation of HB models to the ACS bandpasses should not employ the empirical color correction we employ in the transformation of the isochrones. Because the uncorrected ZAHB agrees well across the entire range of color, especially at the blue end, this might be indicating that a significant part of the correction required for the isochrones is driven by limitations in the synthetic spectra of stars at low effective temperature and low surface gravity, where one must account for large atomic and molecular opacities.
The transformed models for each cluster can also be used to estimate the effects of reddening on the observed fiducials, allowing them to be transformed to any given distance and reddening, for comparison to other CMDs. For example, in Brown et al. (2004), we transformed the fiducials of NGC 5927 and NGC 104 to the Andromeda distance and reddening. At each point on the fiducial, we took the closest point on the associated model, calculated the position in the CMD at both the cluster reddening and the Andromeda reddening, applied this difference to the fiducial, and then shifted by the difference in distance modulus. The absorption in each of the ACS bandpasses ($`A_{F606W}`$ and $`A_{F814W}`$) is a function of both the effective temperature of the source and the foreground reddening, as parameterized by $`E(BV)`$. The ratio of absorption in an ACS band to $`E(BV)`$ is not constant: $`R_{F606W}A_{F606W}/E(BV)`$ and $`R_{F814W}A_{F814W}/E(BV)`$ are both very slowly varying functions of $`E(BV)`$. For example, if one assumed that $`A_{F606W}`$ at the NGC 104 turnoff was ten times larger for $`E(BV)=1.0`$ mag than it was at $`E(BV)=0.1`$ mag, this would introduce an error of 0.08 mag. Thus, the absorption in HST bandpasses has sometimes been published in tables giving absorption in the bandpasses at specific values of effective temperature and $`E(BV)`$ (e.g., see Bedin et al. 2005 for absorption in ACS bands and Holtzman et al. 1995 for absorption in WFPC2 bands, both specified at two distinct effective temperatures). However, using the isochrone associated with each ridge line, one can parameterize the absorption along the ridge line more generally so that it can be transformed to any reference frame. The parameterization is straightforward because $`R_{F606W}`$ and $`R_{F814W}`$ are nearly linear functions of $`E(BV)`$, so we provide, in Tables 2–7, the ridge line for each cluster (as observed) plus the coefficients ($`\alpha `$, $`\beta `$) needed to calculate $`R_{F606W}`$ and $`R_{F814W}`$ along the ridge line:
$`R_{F606W}=\alpha _{F606W}+\beta _{F606W}\times E(BV)`$
$`R_{F814W}=\alpha _{F814W}+\beta _{F814W}\times E(BV)`$.
The absorption is then the product of $`R`$ and $`E(BV)`$. We derive these coefficients by transforming the corresponding isochrone to different values of $`E(BV)`$, fitting $`R_{F606W}`$ and $`R_{F814W}`$ as functions of $`E(BV)`$, and then matching points on the isochrone to points on the ridge line (with the isochrone at the cluster extinction). Where the ridge lines extend beyond the isochrone at the faint and bright ends of the ridge line, we simply take the closest isochrone point in the CMD. Because $`R_{F606W}`$ and $`R_{F814W}`$ are well-approximated by the coefficients in Table 2–7, there is little discernible difference ($`0.01`$ mag) between a transformation of the ridge line using these coefficients and the more exact method of Brown et al. (2004). In Tables 8–13, we provide the analogous data for the HB loci, again by calculating the ZAHB model at different values of $`E(BV)`$, fitting the dependence of $`R_{F606W}`$ and $`R_{F814W}`$ on $`E(BV)`$, and matching stars in the HB locus to points in the ZAHB model at the same color. We stress that the ridge lines and HB loci in Tables 2–13 present the fiducials as observed, and thus they include the reddening for each cluster (Table 1); to transform these fiducials to a different reference frame with a distinct reddening, one must first subtract the reddening intrinsic to each cluster and then add the reddening for the new reference frame, using the coefficients in Tables 2–13.
Table 2:NGC 6341 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 24.346 | 2.628 | -0.077 | 24.279 | 1.709 | -0.025 |
| 23.731 | 2.628 | -0.077 | 23.724 | 1.709 | -0.025 |
| 23.103 | 2.628 | -0.077 | 23.178 | 1.709 | -0.025 |
| 22.570 | 2.628 | -0.077 | 22.726 | 1.709 | -0.025 |
| 22.090 | 2.639 | -0.078 | 22.327 | 1.712 | -0.025 |
| 21.592 | 2.654 | -0.080 | 21.909 | 1.715 | -0.025 |
| 21.194 | 2.666 | -0.082 | 21.572 | 1.718 | -0.025 |
| 20.712 | 2.679 | -0.083 | 21.159 | 1.721 | -0.025 |
| 20.248 | 2.689 | -0.084 | 20.747 | 1.723 | -0.025 |
| 19.947 | 2.695 | -0.084 | 20.475 | 1.725 | -0.025 |
Note: Table 2 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 3:NGC 6752 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 22.610 | 2.613 | -0.074 | 22.432 | 1.707 | -0.025 |
| 22.178 | 2.613 | -0.074 | 22.066 | 1.707 | -0.025 |
| 21.831 | 2.613 | -0.074 | 21.773 | 1.707 | -0.025 |
| 21.398 | 2.613 | -0.074 | 21.424 | 1.707 | -0.025 |
| 20.970 | 2.620 | -0.075 | 21.084 | 1.708 | -0.025 |
| 20.609 | 2.633 | -0.077 | 20.799 | 1.711 | -0.025 |
| 20.295 | 2.644 | -0.079 | 20.544 | 1.713 | -0.025 |
| 19.874 | 2.658 | -0.080 | 20.194 | 1.716 | -0.025 |
| 19.526 | 2.669 | -0.082 | 19.897 | 1.719 | -0.025 |
| 19.169 | 2.679 | -0.083 | 19.585 | 1.721 | -0.025 |
Note: Table 3 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 4:NGC 104 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 23.966 | 2.567 | -0.067 | 23.323 | 1.692 | -0.024 |
| 23.332 | 2.567 | -0.067 | 22.836 | 1.692 | -0.024 |
| 22.784 | 2.567 | -0.067 | 22.415 | 1.692 | -0.024 |
| 22.077 | 2.567 | -0.067 | 21.867 | 1.692 | -0.024 |
| 21.519 | 2.578 | -0.068 | 21.436 | 1.698 | -0.025 |
| 21.028 | 2.594 | -0.071 | 21.067 | 1.704 | -0.025 |
| 20.466 | 2.618 | -0.075 | 20.641 | 1.710 | -0.025 |
| 20.022 | 2.636 | -0.077 | 20.293 | 1.713 | -0.025 |
| 19.771 | 2.647 | -0.079 | 20.088 | 1.715 | -0.025 |
| 19.467 | 2.658 | -0.080 | 19.832 | 1.717 | -0.025 |
Note: Table 4 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 5:NGC 5927 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 26.339 | 2.560 | -0.066 | 25.207 | 1.679 | -0.024 |
| 25.678 | 2.560 | -0.066 | 24.710 | 1.679 | -0.024 |
| 25.018 | 2.560 | -0.066 | 24.229 | 1.679 | -0.024 |
| 24.284 | 2.564 | -0.067 | 23.688 | 1.685 | -0.024 |
| 23.600 | 2.579 | -0.068 | 23.188 | 1.696 | -0.024 |
| 22.972 | 2.601 | -0.072 | 22.719 | 1.706 | -0.025 |
| 22.554 | 2.618 | -0.075 | 22.396 | 1.710 | -0.025 |
| 22.204 | 2.633 | -0.077 | 22.112 | 1.713 | -0.025 |
| 21.873 | 2.647 | -0.079 | 21.837 | 1.716 | -0.025 |
| 21.546 | 2.658 | -0.080 | 21.552 | 1.718 | -0.025 |
Note: Table 5 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 6:NGC 6528 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 24.847 | 2.559 | -0.066 | 23.908 | 1.680 | -0.024 |
| 24.554 | 2.564 | -0.066 | 23.710 | 1.686 | -0.024 |
| 24.253 | 2.567 | -0.067 | 23.508 | 1.690 | -0.024 |
| 23.961 | 2.575 | -0.068 | 23.309 | 1.695 | -0.024 |
| 23.660 | 2.583 | -0.069 | 23.104 | 1.700 | -0.025 |
| 23.365 | 2.594 | -0.071 | 22.900 | 1.704 | -0.025 |
| 23.088 | 2.604 | -0.072 | 22.696 | 1.707 | -0.025 |
| 22.823 | 2.618 | -0.074 | 22.493 | 1.710 | -0.025 |
| 22.566 | 2.629 | -0.076 | 22.288 | 1.712 | -0.025 |
| 22.316 | 2.640 | -0.078 | 22.088 | 1.714 | -0.025 |
Note: Table 6 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 7:NGC 6791 ridge line and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 22.838 | 2.552 | -0.067 | 22.014 | 1.670 | -0.023 |
| 22.333 | 2.559 | -0.067 | 21.687 | 1.680 | -0.024 |
| 21.640 | 2.569 | -0.067 | 21.238 | 1.690 | -0.024 |
| 21.071 | 2.581 | -0.069 | 20.880 | 1.700 | -0.025 |
| 20.667 | 2.595 | -0.071 | 20.600 | 1.706 | -0.025 |
| 20.002 | 2.621 | -0.075 | 20.102 | 1.713 | -0.025 |
| 19.647 | 2.634 | -0.077 | 19.828 | 1.716 | -0.025 |
| 19.327 | 2.645 | -0.078 | 19.567 | 1.718 | -0.025 |
| 18.905 | 2.657 | -0.080 | 19.202 | 1.721 | -0.025 |
| 18.480 | 2.666 | -0.081 | 18.822 | 1.723 | -0.025 |
Note: Table 7 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 8:NGC 6341 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 16.858 | 2.801 | -0.092 | 18.074 | 1.752 | -0.025 |
| 16.838 | 2.801 | -0.092 | 17.993 | 1.752 | -0.025 |
| 16.245 | 2.801 | -0.092 | 17.397 | 1.752 | -0.025 |
| 16.769 | 2.801 | -0.092 | 17.905 | 1.752 | -0.025 |
| 16.789 | 2.801 | -0.092 | 17.916 | 1.752 | -0.025 |
| 16.943 | 2.801 | -0.092 | 18.068 | 1.752 | -0.025 |
| 16.791 | 2.801 | -0.092 | 17.916 | 1.752 | -0.025 |
| 16.739 | 2.801 | -0.092 | 17.862 | 1.752 | -0.025 |
| 16.722 | 2.799 | -0.092 | 17.841 | 1.751 | -0.025 |
| 16.586 | 2.799 | -0.092 | 17.703 | 1.751 | -0.025 |
Note: Table 8 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 9:NGC 6752 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 17.674 | 2.802 | -0.092 | 18.954 | 1.752 | -0.025 |
| 17.740 | 2.802 | -0.092 | 18.979 | 1.752 | -0.025 |
| 17.271 | 2.802 | -0.092 | 18.507 | 1.752 | -0.025 |
| 18.010 | 2.802 | -0.092 | 19.244 | 1.752 | -0.025 |
| 17.569 | 2.802 | -0.092 | 18.797 | 1.752 | -0.025 |
| 17.696 | 2.802 | -0.092 | 18.923 | 1.752 | -0.025 |
| 17.744 | 2.802 | -0.092 | 18.967 | 1.752 | -0.025 |
| 17.315 | 2.802 | -0.092 | 18.536 | 1.752 | -0.025 |
| 16.533 | 2.802 | -0.092 | 17.741 | 1.752 | -0.025 |
| 17.530 | 2.802 | -0.092 | 18.737 | 1.752 | -0.025 |
Note: Table 9 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 10:NGC 104 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 14.075 | 2.678 | -0.083 | 14.484 | 1.722 | -0.025 |
| 13.874 | 2.678 | -0.083 | 14.279 | 1.722 | -0.025 |
| 14.028 | 2.678 | -0.083 | 14.432 | 1.722 | -0.025 |
| 14.056 | 2.678 | -0.083 | 14.460 | 1.722 | -0.025 |
| 13.998 | 2.678 | -0.083 | 14.400 | 1.722 | -0.025 |
| 14.028 | 2.678 | -0.083 | 14.430 | 1.722 | -0.025 |
| 14.085 | 2.678 | -0.083 | 14.485 | 1.722 | -0.025 |
| 13.940 | 2.678 | -0.083 | 14.339 | 1.722 | -0.025 |
| 14.068 | 2.678 | -0.083 | 14.466 | 1.722 | -0.025 |
| 14.074 | 2.678 | -0.083 | 14.472 | 1.722 | -0.025 |
Note: Table 10 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 11:NGC 5927 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 16.407 | 2.671 | -0.082 | 16.402 | 1.721 | -0.025 |
| 16.430 | 2.671 | -0.082 | 16.421 | 1.721 | -0.025 |
| 16.370 | 2.671 | -0.082 | 16.361 | 1.721 | -0.025 |
| 16.506 | 2.671 | -0.082 | 16.490 | 1.721 | -0.025 |
| 16.403 | 2.671 | -0.082 | 16.384 | 1.721 | -0.025 |
| 16.534 | 2.671 | -0.082 | 16.515 | 1.721 | -0.025 |
| 16.407 | 2.671 | -0.082 | 16.386 | 1.721 | -0.025 |
| 16.503 | 2.671 | -0.082 | 16.479 | 1.721 | -0.025 |
| 16.452 | 2.671 | -0.082 | 16.428 | 1.721 | -0.025 |
| 16.520 | 2.671 | -0.082 | 16.495 | 1.721 | -0.025 |
Note: Table 11 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 12:NGC 6528 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 16.784 | 2.669 | -0.081 | 16.654 | 1.720 | -0.025 |
| 16.766 | 2.669 | -0.081 | 16.636 | 1.720 | -0.025 |
| 16.728 | 2.669 | -0.081 | 16.597 | 1.720 | -0.025 |
| 16.954 | 2.669 | -0.081 | 16.818 | 1.720 | -0.025 |
| 16.981 | 2.669 | -0.081 | 16.841 | 1.720 | -0.025 |
| 16.813 | 2.669 | -0.081 | 16.670 | 1.720 | -0.025 |
| 16.835 | 2.669 | -0.081 | 16.691 | 1.720 | -0.025 |
| 16.927 | 2.669 | -0.081 | 16.778 | 1.720 | -0.025 |
| 16.876 | 2.669 | -0.081 | 16.727 | 1.720 | -0.025 |
| 16.935 | 2.669 | -0.081 | 16.784 | 1.720 | -0.025 |
Note: Table 12 is presented in its entirety in the electronic edition of the Astronomical Journal. A portion is shown here for guidance regarding its form and content.
Table 13:NGC 6791 HB locus and reddening coefficients
| $`m_{F606W}`$ | $`\alpha _{F606W}`$ | | $`m_{F814W}`$ | $`\alpha _{F814W}`$ | |
| --- | --- | --- | --- | --- | --- |
| (mag) | (mag) | $`\beta _{F606W}`$ | (mag) | (mag) | $`\beta _{F814W}`$ |
| 14.455 | 2.623 | -0.074 | 14.533 | 1.719 | -0.025 |
| 14.528 | 2.620 | -0.073 | 14.597 | 1.719 | -0.025 |
| 14.524 | 2.620 | -0.073 | 14.583 | 1.719 | -0.025 |
| 14.480 | 2.615 | -0.073 | 14.502 | 1.717 | -0.025 |
| 14.537 | 2.612 | -0.073 | 14.501 | 1.713 | -0.025 |
Note: Table 13 is presented in its entirety here, but for completeness a machine-readable table is available in the electronic edition of the Astronomical Journal, matching the format of Tables 2–12.
In Figure 7, we show all of the ridge lines and HB loci transformed to a distance of 10 pc and $`E(BV)=0`$ mag, using the information in Tables 1–13. Specifically, for each cluster, we calculated $`M_{F606W}m_{F606W}A_{F606W}(mM)_0`$ and $`M_{F814W}m_{F814W}A_{F8146W}(mM)_0`$, with $`A=\alpha \times E(BV)+\beta \times [E(BV)]^2`$ and $`(mM)_0=(mM)_V3.1\times E(BV)`$. When placed in the same reference frame, the clusters fall at their expected relative locations, given their relative ages and metallicities (Table 1). Although the distances for these clusters were determined via a variety of methods, there is clear consistency in the relative distances and reddenings, given the excellent agreement at the HB. Note that the NGC 6528 and NGC 6791 photometry reaches only $`3`$ mag below the turnoff, and so their fiducials are truncated above the point where the main sequence begins to steepen (in contrast with the other clusters, which have fiducials reaching $`5`$ mag below the turnoff); this variation in depth exaggerates the impression that the fiducials are diverging at the faint limit.
## 5. Summary
Using observations of Galactic clusters spanning a wide range in metallicity, we provide ridge lines and HB loci in the two most popular ACS/WFC bandpasses, along with coefficients for transforming these fiducials to a reference frame with an arbitrary reddening and distance. We also provide the algorithm for transforming the VRI to the ACS bandpasses, which provides good agreement ($`0.02`$ mag) with the cluster data from a point 1.5 mag below the main sequence turnoff to a point near the tip of the RGB. The empirical fiducials and the isochrone transformations should be useful in the analysis of stellar population data obtained using ACS.
Support for proposals 9453 and 10265 was provided by NASA through a grant from STScI, which is operated by AURA, Inc., under NASA contract NAS 5-26555. We are grateful to P. Stetson for providing his codes and assistance. We are also grateful to F. Castelli for making her grids of synthetic spectra available to us. We thank the members of the scheduling and operations teams at STScI (especially P. Royle, D. Taylor, and D. Soderblom) for their efforts in executing our large HST programs.
Fig. 7– The ridge lines and HB loci of Figures 1–6, shifted to a distance of 10 pc and no reddening, using the reddening coefficients of Tables 2–13 and the distance and reddening specified in Table 1. The globular cluster fiducials (color points and solid curves) and open cluster (grey points and dashed curve) fall where one would expect relative to each other, given their relative ages and metallicities (Table 1). The HB loci also show good agreement, which suggests the relative distance moduli and extinctions are consistent, even though they were derived via a variety of techniques in the literature (fitting subgiants, white dwarfs, and HB stars).
|
warning/0506/math0506506.html
|
ar5iv
|
text
|
# Rankin-Cohen brackets and formal quantization
## 1 Introduction
In the study of transversal index theory, Connes and Moscovici introduced a Hopf algebra, $`_1`$, which governs the local symmetry in calculating the index of a transversal elliptic operator. Interestingly, Connes and Moscovici discovered an action of $`_1`$ on the modular Hecke algebras.
Inspired by this action, Connes and Moscovici found many similarities between the theory of codimension one foliations and the theory of modular forms. For example, they showed that the Hopf cyclic version of the Godbillon-Vey cocycle gives rise to a 1-cocycle on $`PSL(2,)`$ with values in Eisenstein series of weight 2, and that the Schwarzian 1-cocycle corresponds to an inner derivation implemented by a level 1 Eisenstein series of weight 4. In particular, inspired by Zagier’s Rankin-Cohen deformation on modular forms, Connes and Moscovici constructed a universal deformation formula for an action of $`_1`$ with a projective structure. In this paper, we aim to reconstruct this deformation formula using noncommutative Poisson geometry as developed by the second author and .
The origin of the Rankin-Cohen deformation is a work of Rankin. Rankin in 1956 described all polynomials in the derivatives of modular forms with values again in modular forms. Based on Rankin’s work, in 1977, Cohen defined a sequence of bilinear operations on modular forms indexed by nonnegative integer $`n`$, which assigns to two modular forms, $`f`$ of weight $`k`$ and $`g`$ of weight $`l`$, a modular form of weight $`k+l+2n`$. Their results showed that for any given integer $`n0`$, there is essentially(up to a constant) only one bilinear operator mapping<sup>1</sup><sup>1</sup>1$`_p`$ is the space of modular forms of weight $`p`$. $`_p_q`$ to $`_{p+q+2n}`$ $`p,q_0`$. They are later called Rankin-Cohen brackets and usually denoted by $`RC_n`$. These operators were further studied and played an important role in the theory of modular forms. Zagier observed that the sum of Rankin-Cohen brackets defines an associative product on the algebra $`:=_{l0}_l`$. Zagier’s proof of the associativity of this product, which involves infinitely many equalities, was rather combinatoric. Cohen, Manin, and Zagier explained this deformation using the theory of automorphic pseudo differential operators. The calculation still involves many interesting and complicated combinatoric identities. In this paper, we will first reconstruct Zagier’s Rankin-Cohen deformation using the methods of deformation quantization of symplectic manifolds developed by Fedosov . In particular, we will show that this deformation is isomorphic to the standard Moyal product. The calculation involved in our construction is easier and more transparent than those and .
To reconstruct Connes-Moscovici’s Rankin-Cohen deformation for $`_1`$ action, we need to first understand the projective structure introduced by Connes and Moscovic . The notion of a projective structure of $`_1`$ is a generalization of the projective structure on an elliptic curve ( see ). Our idea to understand this structure is to look at the defining action of $`_1`$ on a groupoid algebra associated to a codimension one foliation. In this case, we discovered that the existence of a projective structure is equivalent to the existence of a certain type of invariant symplectic connection. This geometric explanation provides a natural connection to the results in Tang , where he studied the deformation quantization of a groupoid algebra. The existence of an invariant symplectic connection is a sufficient condition for the existence of a deformation quantization of a groupoid algebra. Therefore, in the case of a codimension one foliation, Tang’s construction implies that with a projective structure, one can construct a deformation quantization (a star product) of the corresponding foliation groupoid algebra. Furthermore, our calculation in Section 5 exhibits that when the symplectic connection is flat, the star product on the groupoid algebra can be expressed by an element $`RC`$ in $`_1_1[[\mathrm{}]]`$. To obtain a universal deformation for a $`_1`$ action with a projective structure as Connes and Moscovici , we construct a fully injective $`_1`$ action on the union of groupoid algebras of those foliation groupoids with a fixed type of invariant symplectic connections. Therefore, we are able to reconstruct the universal deformation formula on $`_1`$ by pulling back the star products on the groupoid algebras.
All the above deformations, including , , and , are all formal deformation, which means that the deformation parameter $`t`$ is a formal variable. It is more interesting to ask whether one can make a deformation strict in the sense of Rieffel. This will be studied in the next paper .
Acknowledgment: We would like to thank Alain Connes and Henri Moscovici for explaining the results and asking us interesting questions. Tang would like to thank Max Karoubi and Richzard Nest for their hosts of his visit of Institute de Henri Poincaré during summer 2004, where the paper started. Yao wants to thank Don Zagier for his inspiring course given at Collège de France.
## 2 Prerequisites
In this section, we review the materials needed for this paper.
### 2.1 Codimension one foliations and the Hopf algebra
For a constant rank foliation on $`M`$, we choose a complete flat transversal $`X`$. We look at the oriented frame bundle $`FX`$ of $`X`$ with the lifted holonomy foliation groupoid action, which defines an étale groupoid $`𝒢FX`$. Connes and Moscovici found a Hopf algebra $`_k`$ acting on the smooth groupoid algebra $`C_c^{\mathrm{}}(𝒢)`$, where $`k`$ is the codimension of the foliation. We exhibit this Hopf algebra in the case of $`k=1`$.
In the case of a codimension one foliation, the complete transversal $`X`$ is a flat 1-dim manifold, and $`FX`$ is isomorphic to $`X\times ^+`$ by fixing a flat connection on $`FXX`$. We introduce coordinates $`x`$ on the $`X`$ component and $`y`$ on the $`^+`$ component. Let $`\mathrm{\Gamma }`$ be a pseudogroup associated to the foliation acting on $`X`$. The lifted action of $`\mathrm{\Gamma }`$ on $`FX`$ is
$$(x,y)(\varphi (x),\varphi ^{}(x)y),\varphi \mathrm{\Gamma }.$$
We look at the groupoid $`FX\mathrm{\Gamma }FX`$. It is an étale groupoid with a natural symplectic form $`\omega =\frac{\mathrm{d}x\mathrm{d}y}{y^2}`$.
On $`FX`$, we consider vector fields $`X=y_x`$ and $`Y=y_y`$. It is easy to check that $`Y`$ is invariant under the $`\mathrm{\Gamma }`$ action, but $`X`$ is not, and has the following commutation relation,
$$U_\varphi XU_\varphi ^1=Xy\frac{\varphi _{}^{1}{}_{}{}^{\prime \prime }(x)}{\varphi _{}^{1}{}_{}{}^{}(x)}Y.$$
We introduce the following operators on $`𝒜`$.
$$\begin{array}{cc}X(fU_\varphi )\hfill & =X(f)U_\varphi ,\hfill \\ Y(fU_\varphi )\hfill & =Y(f)U_\varphi ,\hfill \\ \delta _1(fU_\varphi )\hfill & =\mu _{\varphi ^1}fU_\varphi ,\hfill \\ \delta _n(fU_\varphi )\hfill & =X^{n1}(\mu _{\varphi ^1})fU_\varphi ,\hfill \end{array}$$
(1)
where $`\mu _{\varphi ^1}(x,y)=y\frac{\varphi _{}^{1}{}_{}{}^{\prime \prime }(x)}{\varphi _{}^{1}{}_{}{}^{}(x)}`$.
The commutation relation among the above operators are
$$\begin{array}{cc}[Y,X]=X,\hfill & [X,\delta _n]=\delta _{n+1},\hfill \\ [Y,\delta _n]=n\delta _n,\hfill & [\delta _n,\delta _m]=0.\hfill \end{array}$$
The operators $`X,Y,\delta _n,n`$ form an infinite dimensional Lie algebra $`H_1`$, and the Hopf algebra $`_1`$ is defined to be the universal enveloping algebra of $`H_1`$.
We define the following operations on $`_1`$:
1. product $`:_1_1_1`$ by the product on $`_1`$ as the universal enveloping algebra of $`H_1`$.
2. coproduct $`\mathrm{\Delta }:_1_1_1`$ by
$$\begin{array}{c}\mathrm{\Delta }Y=Y1+1Y,\hfill \\ \mathrm{\Delta }\delta _1=\delta _11+1\delta _1,\hfill \\ \mathrm{\Delta }X=X1+1X+\delta _1Y,\hfill \\ \mathrm{\Delta }\delta _n=[\mathrm{\Delta }X,\mathrm{\Delta }\delta _{n1}].\hfill \end{array}$$
3. counit $`ϵ:_1`$ by taking the value of the identity component.
4. antipode $`S:_1_1`$ by
$$S(X)=X+\delta _1Y,S(Y)=Y,S(\delta _1)=\delta _1.$$
It is straightforward to check that $`(_1,,\mathrm{\Delta },S,ϵ,id)`$ defines a Hopf algebra.
### 2.2 Deformation quantization a la Fedosov
Fedosov’s construction of deformation quantizations of a symplectic manifold can be formulated as follows.
Let $`(M,\omega )`$ be a $`2n`$ dimensional symplectic manifold. At each fiber $`T_xM`$ of the tangent bundle, which is a symplectic vector space, we define a Weyl algebra $`W_x`$ to be an associative algebra over $``$ with a unit, whose elements are of the form
$$a(y,\mathrm{})=\underset{k,|\alpha |0}{}\mathrm{}^ka_{k,\alpha }y^\alpha ,$$
where $`\mathrm{}`$ is a formal parameter and $`y=(y^1,\mathrm{},y^{2n})T_xM`$ is a tangent vector, $`\alpha =(\alpha _1,\mathrm{},\alpha _{2n})`$ is a multi-index, $`y^\alpha =(y^1)^{\alpha _1}\mathrm{}(y^{2n})^{\alpha _{2n}}`$.
The product of elements $`a,bW_x`$ is defined as follows:
$$\begin{array}{cc}ab\hfill & =\mathrm{exp}(\frac{i\mathrm{}}{2}\omega ^{ij}\frac{}{y^i}\frac{}{z^j})a(y,\mathrm{})b(z,\mathrm{})|_{z=y}\hfill \\ & =_{k=0}^{\mathrm{}}(\frac{i\mathrm{}}{2})^k\frac{1}{k!}\omega ^{i_1j_1}\mathrm{}\omega ^{i_kj_k}\frac{^ka}{y^{i_1}\mathrm{}y^{i_k}}\frac{^kb}{y^{j_1}\mathrm{}y^{j_k}}.\hfill \end{array}$$
We consider the Weyl algebra bundle $`W`$ over $`(M,\omega )`$ for which the fiber at the point $`x`$ is $`W_x`$, and denote $`C^{\mathrm{}}(W)`$ to be the algebra of smooth sections of $`W`$ with pointwise multiplication $``$. To introduce the Fedosov connection, we look at the algebra $`C^{\mathrm{}}(W\mathrm{\Lambda })=_{q=0}^{2n}\mathrm{\Gamma }^{\mathrm{}}(W\mathrm{\Lambda }^q)`$, where $`\mathrm{\Lambda }^q`$ is set of smooth $`q`$forms.
We introduce several operations on $`C^{\mathrm{}}(W\mathrm{\Lambda })`$.
1. commutator, i.e. $`[a,b]=ab(1)^{deg(a)deg(b)}ba`$.
2. $`\delta ,\delta ^{}:C^{\mathrm{}}(W\mathrm{\Lambda })C^{\mathrm{}}(W\mathrm{\Lambda })`$, i.e.
$$\delta a=dx^k\frac{a}{y^k},\delta ^{}a=y^ki(\frac{}{x^k})a.$$
A Fedosov connection on the Weyl algebra bundle $`W`$ is a connection $`D`$ such that for any section $`aC^{\mathrm{}}(W\mathrm{\Lambda })`$,
$$D^2a=\frac{i}{\mathrm{}}[\mathrm{\Omega },a]=0.$$
Fedosov in showed that given a torsion free symplectic connection $``$ on $`M`$ with Christoffel $`\mathrm{\Gamma }_{ijk}`$, one can construct an abelian connection on $`W`$ of the following form
$$D=\delta ++\frac{i}{\mathrm{}}[r,],$$
where $`a:=da+\frac{i}{\mathrm{}}[\mathrm{\Gamma },a]`$, with $`\mathrm{\Gamma }=\frac{1}{2}\mathrm{\Gamma }_{ijk}y^iy^jdx`$, and $`r`$ is a local 1-form with values in $`W`$.
We look at the subalgebra $`W_DC^{\mathrm{}}(W)`$ consisting of flat sections of $`D`$. The main theorem that we will use is the following:
###### Theorem 2.1
For any $`a_0C^{\mathrm{}}(M)[[\mathrm{}]]`$, there exists a unique section $`aW_D`$, which is denoted by $`\sigma ^1(a_0)`$, such that $`\sigma (a)=a_0`$, where $`\sigma (a)`$ means the projection onto the center: $`\sigma (a)=a(x,0,h)`$.
This implies that there is a one-to-one correspondence between $`W_D`$ and $`C^{\mathrm{}}(M)[[\mathrm{}]]`$. Accordingly we can define on $`C^{\mathrm{}}(M)[[\mathrm{}]]`$ an associative star product
$$ab=\sigma (\sigma ^1(a)\sigma ^1(b)).$$
(2)
### 2.3 Deformation Quantization of Groupoids
The second named author considered deformation quantization of the groupoid algebra of a pseudo étale groupoid and proved that one can construct star products on such groupoids. As a special case, we have that for an étale groupoid with an invariant symplectic structure and an invariant symplectic connection on the base, the groupoid algebra can be formally deformation quantized. In this subsection, we recall the basic concepts and constructions from Tang .
###### Definition 1
(Block, Getzler and Xu) A Poisson structure on an associative algebra $`A`$ is an element $`[\mathrm{\Pi }]`$ of the Hochschild cohomology group $`H^2(A,A)`$ such that the cohomology class of the Gerstenhaber bracket $`[\mathrm{\Pi },\mathrm{\Pi }]`$ vanishes.
###### Definition 2
Let $`(A,[\mathrm{\Pi }])`$ be a noncommutative Poisson algebra, and $`A[[\mathrm{}]]`$ the space of formal power series with coefficients in $`A`$. A formal deformation quantization of $`(A,[\mathrm{\Pi }])`$ (or in other words star product) is an associative product
$$:A[[\mathrm{}]]\times A[[\mathrm{}]]A[[\mathrm{}]],(a_1,a_2)a_1a_2=\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{}^kc_k(a_1,a_2)$$
satisfying the following properties:
1. Each one of the maps $`c_k:A[[\mathrm{}]]A[[\mathrm{}]]A[[\mathrm{}]]`$ is $`[[\mathrm{}]]`$-bilinear;
2. One has $`c_0(a_1,a_2)=a_1a_2`$ for all $`a_1,a_2A`$;
3. The relation
$$a_1a_2c_0(a_1,a_2)\frac{i}{2}\mathrm{}\mathrm{\Pi }(a_1,a_2)\mathrm{}^2A[[\mathrm{}]]$$
holds true for some representative $`\mathrm{\Pi }Z^2(A,A)`$ of the Poisson structure and all $`a_1,a_2A`$.
For an étale groupoid $`𝒢`$ with an invariant symplectic form $`\omega `$ and a invariant symplectic connection $``$ on the base, we define a Hochschild 2-cochain on $`C^{\mathrm{}}(𝒢)`$ by
$$\mathrm{\Pi }(a_1,a_2)(g)=\underset{g_1g_2=g}{}\pi (g)(da_1(g_1),da_2(g_2)),g𝒢a_1,a_2C^{\mathrm{}}(𝒢),$$
(3)
where $`da_1(g_1)`$ and $`da_2(g_2)`$ have been pulled back to $`g`$ along the maps $`t`$ and $`s`$, and $`\pi `$ is the Poisson structure associated to the symplectic form $`\omega `$. This definition is legitimate because $`t`$ and $`s`$ are local diffeomorphisms. It was proved that this Hochschild 2-cochain gives rise to a Poisson structure on $`C^{\mathrm{}}(𝒢)`$ if there is an invariant symplectic connection.
Tang showed that the above noncommutative Poisson structure $`\mathrm{\Pi }`$ on the groupoid algebra admits a formal deformation quantization. Such a deformation can be constructed as follows: first using Fedosov’s construction , given an invariant symplectic connection, we construct an invariant star product on the algebra of smooth functions on the unit space $`𝒢^{(0)}`$. The deformation of the groupoid algebra $`C^{\mathrm{}}(𝒢)`$ is a crossed product algebra of the above deformation on the base $`C^{\mathrm{}}(𝒢^{(0)})`$ and the associated pseudogroup $`𝒢`$ action.
### 2.4 Rankin-Cohen deformation
It is well known that if $`f(z)`$ is a modular form, $`\frac{1}{2\pi i}\frac{d}{dz}f`$ is not a modular form any more. Following , we introduce a differential operator $`X`$ as
$$X\stackrel{def}{=}\frac{1}{2\pi i}\frac{d}{dz}\frac{1}{12\pi i}\frac{d}{dz}(\mathrm{log}\mathrm{\Delta })Y,$$
where $`\mathrm{\Delta }(z)=(2\pi )^{12}\eta ^{24}(z)=(2\pi )^{12}q_{n=1}^{\mathrm{}}(1q^n)^{24},q=e^{2\pi z}`$ and $`Y(f)=\frac{k}{2}f,f_k`$, the space of modular forms of weight $`k`$.
It is straightforward to check that $`X`$ and $`Y`$ acts on $`=_k_k`$ satisfying $`[Y,X]=X`$. Under these two operators, the Rankin-Cohen bracket $`RC_n`$ can be written as follows, for $`f_k,g_l`$
$$\begin{array}{cc}RC_n(f,g)\hfill & =_{r+s=n}(1)^r\left(\begin{array}{c}n+k1\\ s\end{array}\right)\left(\begin{array}{c}n+l1\\ r\end{array}\right)f^{(r)}g^{(s)},\hfill \end{array}$$
where $`f^{(r)}`$ (or $`g^{(s)})`$ is the $`r`$-th (or $`s`$-th) derivative of $`f`$ (or $`g)`$, and $`(\alpha )_k\stackrel{def}{=}\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$.
In , Zagier observed that $`_nRC_n`$ defines an associative product on $``$. This product actually defines a universal deformation formula of the Lie algebra $`h_1`$, consisting of $`X,Y`$ with $`[Y,X]=X`$, since $`h_1`$ acts on $``$ injectively. It is worth mentioning that $`h_1`$ is the Lie algebra of the $`\mathrm{`}\mathrm{`}ax+b\mathrm{"}`$ group.
Inspired by the Rankin-Cohen brackets, Connes and Moscovici introduced a family of Rankin-Cohen type elements in $`(_1_1)[[\mathrm{}]]`$ as follows.
###### Definition 2.2
()Let $`_1`$ act on an algebra $`A`$. This action is called projective if $`\delta _2^{}\stackrel{def}{=}\delta _1^2\frac{1}{2}\delta _2`$ is inner implemented by an element $`\mathrm{\Omega }A`$, so that
$$\delta _2^{}(a)=[\mathrm{\Omega },a],aA,$$
and
$$\delta _k(\mathrm{\Omega })=0,k.$$
Assume that the action of $`_1`$ action an algebra $`A`$ is projective. Define
$$\begin{array}{cc}RC\hfill & =_{n=0}^{\mathrm{}}\mathrm{}^n_{k=0}^n\frac{A_k}{k!}(2Y+k)_{nk}\frac{B_{nk}}{(nk)!}(2Y+nk)_k\hfill \\ A_{m+1}\hfill & =S(X)A_mm\mathrm{\Omega }^0(Y\frac{m1}{2})A_{m1},\hfill \\ B_{m+1}\hfill & =XB_mm\mathrm{\Omega }(Y\frac{m1}{2})B_{m1},\hfill \end{array}$$
(4)
where $`\mathrm{\Omega }^0`$ is the right multiplication of $`\mathrm{\Omega }`$.
Connes and Moscovici proved that $`RC`$ defines a universal deformation formula of a projective $`_1`$ action.
## 3 Universal deformation of $`h_1`$
If we set all $`\delta _n`$ to be $`0`$, the Lie algebra $`H_1`$ is reduced to $`h_1`$, the Lie algebra of the “$`ax+b`$” group, and $`_1`$ becomes $`𝒰(h_1)`$, the universal enveloping algebra of $`h_1`$. In this case, $`RC`$ defined by (4) is simplified to the following universal deformation formula of $`h_1`$,
$$RC_n(a,b)\stackrel{def}{=}\underset{k=0}{\overset{n}{}}\left[\frac{(1)^k}{k!}X^k(2Y+k)_{nk}(a)\frac{1}{(nk)!}X^{nk}(2Y+nk)_k(b)\right],$$
(5)
where $`X,Yh_1`$ are such that $`[Y,X]=X`$, $`(\alpha )_k\stackrel{def}{=}\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$, and $`a,bA`$.
We spend this section studying this universal deformation.
### 3.1 Giaquinto-Zhang’s deformation of $`h_1`$
A nice deformation formula for $`h_1`$ has already been given by Giaquinto and Zhang \[Thm 2.20\]: Given two elements $`X,Y`$ with $`[Y,X]=X`$, the following expression defines a universal deformation formula(UDF) of the Hopf algebra associated to $`h_1`$
$$F=\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^n}{n!}F_n=1\times 1+tXY+\frac{t^2}{2!}\left(X^2Y_22XY_1XY_1+Y_2X^2\right)+\mathrm{},$$
where $`F_n`$ is defined to be $`F_n=_{r=0}^n(1)^r\left(\genfrac{}{}{0pt}{}{n}{r}\right)X^{nr}Y_rX^rY_{nr}`$.
###### Proposition 3.1
The above defined $`F`$ can be realized by the standard Moyal product.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$ We consider the space $`\times _+`$ on which $`X`$ and $`Y`$ act as $`Y=y\frac{}{y}`$, and $`X=\frac{1}{y}\frac{}{x}`$. It is obvious that the action of $`X`$ and $`Y`$ on $`\times _+`$ is injective.
With the following identity,
$$Y_r=Y(Y+1)\mathrm{}(Y+r1)=(y)^r\frac{^r}{y^r},$$
it is straightforward to check that the above defined $`F`$ in this representation is equal to the Moyal product. $`\mathrm{}`$
### 3.2 Rankin-Cohen deformation of $`h_1`$
We should point out that the above universal deformation formula of $`h_1`$ is not equal to the one induced from $`RC`$ in Equation (5). However, we will show that it is equivalent to the Giaquinto-Zhang’s deformation.
We set $`(V,\omega ):=(^2=\{(p,q)\},dpdq)`$ and denote by $`𝔥=𝔥(V,\omega ):=V\times `$ the associated Heisenberg algebra. Setting $`𝔤:=𝔰𝔩_2()=\mathrm{𝚜𝚙𝚊𝚗}_{}\{H,E,F\}`$, $`([H,E]=2E,[H,F]=2F,[E,F]=H)`$, we form the natural semi-direct product $`\stackrel{~}{𝔤}:=𝔤\times 𝔥`$. The (infinitesimal) affine linear action $`\stackrel{~}{\gamma }\mathrm{\Gamma }(T(V))`$ is then strongly hamiltonian. We let $`\lambda :\stackrel{~}{𝔤}C^{\mathrm{}}(V)`$ denote the corresponding moment map. Explicitly, denoting fundamental vector fields by $`A_x^{}:=\frac{d}{dt}|_0\mathrm{exp}(tA)xA\stackrel{~}{𝔤}`$, one has
$$\begin{array}{ccccc}H^{}=p_p+q_q;\hfill & E^{}=q_p;\hfill & F^{}=p_q;\hfill & P^{}=_p;\hfill & Q^{}=_q;\hfill \\ \lambda _H=pq;\hfill & \lambda _E=\frac{1}{2}q^2;\hfill & \lambda _F=\frac{1}{2}p^2;\hfill & \lambda _P=q;\hfill & \lambda _Q=p.\hfill \end{array}$$
We have that $`[A^{},B^{}]=[A,B]^{}`$ and $`\lambda _{[A,B]}=\{\lambda _A,\lambda _B\}`$ where $`\{u,v\}=_pu_qv_pv_qu`$, and $`A,B\stackrel{~}{𝔤}`$.
Let $`S:=AN=\mathrm{exp}(\mathrm{𝚜𝚙𝚊𝚗}\{H,E\})`$ denote the Iwasawa component in $`SL(2,)`$, which is the $`\mathrm{`}\mathrm{`}ax+b\mathrm{"}`$ group. We consider the open orbit $`𝒪\stackrel{def}{=}S(0,1)`$ in $`V`$, which is equal to the set $`[q>0]`$. Since $`S`$ acts simply transitively on $`𝒪`$, we have the identification $`\varphi :S𝒪:gg(1,0)`$. We still denote by $`\lambda :\stackrel{~}{𝔤}C^{\mathrm{}}(S)`$ the transported restricted moment map, that is:
$$\lambda _A:=\varphi ^{}(\lambda _A|_𝒪)(A\stackrel{~}{𝔤}).$$
(6)
###### Lemma 3.2
Denoting by $`\stackrel{~}{X}_g:=\frac{d}{dt}|_0g\mathrm{exp}(tX)`$ the left-invariant vector field associated to $`Xh_1=\text{Lie }(S)`$, one has:
1. $`\stackrel{~}{H}.\lambda _{X+v}=(2)\lambda _X+(1)\lambda _v\text{for }X𝔤\text{ and }vV;`$
2. $`\stackrel{~}{E}^r.\lambda _X=\mathrm{\hspace{0.33em}0}\text{ for }r3,`$ for all $`X𝔤`$;
3. $`\stackrel{~}{E}^r.\lambda _v=\mathrm{\hspace{0.33em}0}\text{ for }r2,`$ for all $`vV`$.
Proof. A convenient parametrization of the group manifold $`S`$ is given by:
$$^2S:(a,\mathrm{})\mathrm{exp}(aH)\mathrm{exp}(\mathrm{}E).$$
In these coordinates, the group law reads $`(a,\mathrm{})(a^{},\mathrm{}^{})=(a+a^{},e^{2a^{}}\mathrm{}+\mathrm{}^{})`$. We deduce the expressions for the left-invariant vector fields:
$$\stackrel{~}{H}=_a2\mathrm{}_{\mathrm{}};\stackrel{~}{E}=_{\mathrm{}}.$$
The corresponding chart on the orbit $`𝒪S`$ is given by
$$p=e^a\mathrm{};q=e^a.$$
Note that this is a global Darboux chart on $`𝒪`$ as for $`dad\mathrm{}=\pm \varphi ^{}\omega |_𝒪`$. The corresponding (uncomplete) moment map reads as
$$\lambda _H=\mathrm{};\lambda _E=\frac{1}{2}e^{2a};\lambda _F=\frac{1}{2}\mathrm{}^2e^{2a};\lambda _P=e^a;\lambda _Q=e^a\mathrm{}.$$
A straightforward computation then yields the lemma. $`\mathrm{}`$
From (5), for any left $`𝒰(h_1)`$ action on an algebra $`A`$, the Rankin-Cohen brackets on $`𝒰(h_1)`$ is defined by,
$$RC_n(a,b):=\underset{k=0}{\overset{n}{}}\left[\frac{(1)^k}{k!}X^k(2Y+k)_{nk}(a)\frac{1}{k!}X^{nk}(2Y+nk)_k(b)\right],$$
where $`X,Yh_1`$ are such that $`[Y,X]=X`$, $`(\alpha )_k\stackrel{def}{=}\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$, and $`a,bA`$.
Since $`h_1`$ acts as left invariant vector fields on $`S`$, $`𝒰(h_1)`$ acts as left invariant differential operators on $`C^{\mathrm{}}(S)`$, and $`RC_n`$, an element of $`𝒰(h_1)𝒰(h_1)`$, acts as a left invariant bidifferential operator on $`C^{\mathrm{}}(S)`$. Since $`[H,E]=2E`$, we set
$$\stackrel{~}{H}=2Y\text{ and }\stackrel{~}{E}=X.$$
###### Lemma 3.3
For all $`A`$ in $`\stackrel{~}{𝔤}`$, we have
$$[\lambda _A,u]_n\stackrel{def}{=}RC_n(\lambda _A,u)RC_n(u,\lambda _A)=\mathrm{\hspace{0.33em}0}\text{for }n1.$$
(7)
Proof. For $`X𝔤`$ and $`vV`$, Lemma 3.2 implies that $`X^k(2Y+r)_s.\lambda _{X+v}=(2+r)_sX^k\lambda _X+(1+r)_sX^k\lambda _v=0`$ if $`k>2`$. Therefore, in the expression (5) of $`RC_n(\lambda _{X+v},u)`$ only the first three terms corresponding to $`k=0,1,2`$ contribute. In each of them the following (left hand side) factor occurs:
$$\text{ for }k=0:(2)_n\lambda _X+(1)_n\lambda _v;$$
(8)
$$\text{ for }k=1:\stackrel{~}{E}.[(1)_{n1}\lambda _X+(0)_{n1}\lambda _v];$$
(9)
$$\text{ for }k=2:\stackrel{~}{E}^2.[(0)_{n2}\lambda _X+(1)_{n2}\lambda _v].$$
(10)
1. The first expression (8) vanishes identically for $`n3`$. Indeed, $`(2)_n=(2)(2+1)(2+2)\mathrm{}(2+n1)`$ is zero as soon as $`n12`$; and similarly for $`(1)_n`$;
2. In the same way, the second expression (9) vanishes for $`n21`$, i.e. $`n3`$;
3. At last, the third expression (10) is equal to $`(n2)!\stackrel{~}{E}^2(\lambda _v)`$ which is identically zero by Lemma 3.2 item (iii). We conclude by observing that $`RC_0`$ and $`RC_2`$ are symmetric. $`\mathrm{}`$
By Lemma 3.3, the Rankin-Cohen deformation (4) defines a $`\stackrel{~}{𝔤}`$ invariant star product on $`(V,\omega )`$. In Corollary 2, Section 2.7 of , Gutt showed that there is a unique $`\stackrel{~}{𝔤}`$-invariant star product on $`(V,\omega )`$, which is the standard Moyal product. We conclude that the Rankin-Cohen deformation on $`C^{\mathrm{}}(S)`$ is identical to the Moyal product.
###### Proposition 3.4
The reduced Rankin-Cohen deformation realized on $`𝒪V`$ coincides with the restriction to $`𝒪`$ of the standard Moyal product on $`(V,\mathrm{\Omega })`$.
To generalize the construction in Proposition 3.4, we explain its relation to Fedosov’s construction of deformation quantization of symplectic manifolds.
The natural action of $`S\mathrm{`}\mathrm{`}ax+b\mathrm{"}`$ on $``$,
$$\mathrm{exp}(aH+nE)x_1:=e^{2a}x_1+ne^a,$$
lifts to $`T^{}()=^2`$ as
$$\mathrm{exp}(aH+nE)(x_1,x_2):=(e^{2a}x_1+ne^a,e^{2a}x_2).$$
The $`S`$-orbit $`\stackrel{~}{𝒪}`$ of point $`\stackrel{~}{o}:=(0,1)=dx_1|_0T^{}(^2)`$ is then naturally isomorphic as $`S`$-homogeneous space to $`𝒪V`$; namely one has the identification:
$$\phi :𝒪\stackrel{~}{𝒪}:ge_2g\stackrel{~}{o}.$$
In $`(p,q)`$-coordinates on $`𝒪`$, this reads:
$$\phi (p,q)=(\frac{p}{2q},q^2).$$
Identifying $`\stackrel{~}{𝒪}`$ with $`S`$ (via $`\phi \varphi `$), we obtain the expressions for the left invariant vector fields:
$$\stackrel{~}{H}=2x_2_{x_2};\stackrel{~}{E}=\frac{1}{x_2}_{x_1}.$$
In particular, we set
$$\stackrel{~}{H}=2Y\text{ and }\stackrel{~}{E}=X.$$
By letting $`^𝒪`$ denote the restriction to $`𝒪`$ of the standard symmetric flat connection on $`V`$ ($`^𝒪{}_{_p}{}^{}_{p}^{}=^𝒪{}_{_q}{}^{}_{p}^{}=^𝒪{}_{_q}{}^{}_{q}^{}=0`$), and setting
$$^{\stackrel{~}{𝒪}}:=\phi (^𝒪),$$
we obtain a symplectic connection on $`\stackrel{~}{𝒪}`$,
$$_{_{x_1}}^{\stackrel{~}{𝒪}}_{x_1}=0;_{_{x_1}}^{\stackrel{~}{𝒪}}_{x_2}=\frac{1}{2x_2}_{x_1};_{_{x_2}}^{\stackrel{~}{𝒪}}_{x_2}=\frac{1}{2x_2}_{x_2}.$$
(11)
We identify $`\stackrel{~}{𝒪}`$ with $`\times ^+`$, and use $`^{\stackrel{~}{𝒪}}`$ to construct deformation quantization of $`(\times ^+,\omega \stackrel{def}{=}dxdy)`$ as described in Section 2.2.
###### Corollary 3.5
The reduced Rankin-Cohen deformation on $`\stackrel{~}{𝒪}`$ is identical to Fedosov’s construction of the star product on $`(\stackrel{~}{𝒪},\omega )`$ using the connection $`^{\stackrel{~}{𝒪}}`$ with the characteristic form equal to $`\frac{1}{i\mathrm{}}\omega `$.
## 4 Projective structures
To reconstruct Connes-Moscovici’s Rankin-Cohen deformation, we need to understand the geometric meaning of their Definition 2.2, a projective structure.
### 4.1 The flat case
We look at the connection $`^{\stackrel{~}{𝒪}}`$ considered in Section 3, (11).
###### Proposition 4.1
The connection $`^{\stackrel{~}{𝒪}}`$ (11) is invariant under the local diffeomorphism $`\varphi :x_1\stackrel{~}{x_1}\stackrel{def}{=}\varphi (x_1),x_2\stackrel{~}{x_2}\stackrel{def}{=}\frac{x_2}{\varphi ^{}(x_1)}`$ if and only if $`\delta _2^{}(\varphi )=0`$. Here $`_1`$ acts on $`\varphi `$ as in Section 2.1.
Notation:We use $``$ to replace $`^{\stackrel{~}{𝒪}}`$ in the rest of the paper.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$We have the following transformation rules of vector fields.
$$\begin{array}{c}\frac{}{\stackrel{~}{x}_1}=\frac{1}{\varphi ^{}(x_1)}\frac{}{x_1}+\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{x_2},\hfill \\ \frac{}{\stackrel{~}{x}_2}=\varphi ^{}\frac{}{x_2}.\hfill \end{array}$$
The invariance of $``$ implies that we should have
$$\begin{array}{cc}_{\varphi _{}(\frac{}{x_1})}\varphi _{}(\frac{}{x_1})\hfill & =_{\varphi ^{}(x_1)\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}}(\varphi ^{}(x_1)\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2})\hfill \\ & =\varphi ^2_{\frac{}{\stackrel{~}{x}_1}}\frac{}{\stackrel{~}{x}_1}+\varphi ^{}\frac{}{\stackrel{~}{x}_1}(\varphi ^{})\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^{}}x_2_{\frac{}{\stackrel{~}{x}_1}}\frac{}{\stackrel{~}{x}_2}\varphi ^{}\frac{}{\stackrel{~}{x}_1}(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)\frac{}{\stackrel{~}{x}_2}\hfill \\ & \frac{\varphi ^{\prime \prime }}{\varphi ^{}}x_2_{\frac{}{\stackrel{~}{x}_2}}\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}(\varphi ^{})\frac{}{\stackrel{~}{x}_1}+(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)^2_{\frac{}{\stackrel{~}{x}_2}}\frac{}{\stackrel{~}{x}_2}\hfill \\ & +\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)\frac{}{\stackrel{~}{x}_2}\hfill \\ & =\varphi ^{}\frac{1}{\varphi ^{}}(\varphi ^{\prime \prime })\frac{}{\stackrel{~}{x}_1}\varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}\varphi ^{}[\frac{1}{\varphi ^{}}\frac{\varphi ^{\prime \prime \prime }\varphi ^22\varphi ^{\prime \prime 2}\varphi ^{}}{(\varphi ^2)^2}x_2+(\frac{\varphi ^{\prime \prime }}{\varphi ^2})^2x_2]\frac{}{\stackrel{~}{x}_2}\hfill \\ & \varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}+0+(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)^2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_2}+\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}\frac{}{\stackrel{~}{x}_2}\hfill \\ & =\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^3}x_2\frac{}{\stackrel{~}{x}_2},\hfill \end{array}$$
$$\begin{array}{c}\begin{array}{cc}_{\varphi _{}(\frac{}{x_1})}\varphi _{}(\frac{}{x_2})\hfill & =_{\varphi ^{}(x_1)\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}}(\frac{1}{\varphi ^{}}\frac{}{\stackrel{~}{x}_2})\hfill \\ & =\varphi ^{}\frac{1}{\varphi ^{}}_{\frac{}{\stackrel{~}{x}_1}}\frac{}{\stackrel{~}{x}_2}+\varphi ^{}\frac{}{\stackrel{~}{x}_1}(\frac{1}{\varphi ^{}})\frac{}{\stackrel{~}{x}_2}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{\varphi ^{}}_{\frac{}{\stackrel{~}{x}_2}}\frac{}{\stackrel{~}{x}_2}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}(\frac{1}{\varphi ^{}})\frac{}{\stackrel{~}{x}_2}\hfill \\ & =\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}+\varphi ^{}\frac{1}{\varphi ^{}}(\frac{\varphi ^{\prime \prime }}{\varphi ^2})\frac{}{\stackrel{~}{x}_2}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{\varphi ^{}}(\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_2})0\hfill \\ & =\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}\frac{1}{2}\frac{\varphi ^{\prime \prime }}{\varphi ^2}\frac{}{\stackrel{~}{x}_2}=\varphi _{}(\frac{1}{2x_2}\frac{}{x_1}),\hfill \end{array}\hfill \\ \\ \begin{array}{cc}_{\varphi _{}(\frac{}{x_2})}\varphi _{}(\frac{}{x_2})\hfill & _{\frac{1}{\varphi ^{}}\frac{}{\stackrel{~}{x}_2}}(\frac{1}{\varphi ^{}}\frac{}{\stackrel{~}{x}_2})=\frac{1}{\varphi ^2}_{\frac{}{\stackrel{~}{x}_2}}(\frac{}{\stackrel{~}{x}_2})+\frac{1}{\varphi ^{}}\frac{}{\stackrel{~}{x}_2}(\frac{1}{\varphi ^{}})\frac{}{\stackrel{~}{x}_2}\hfill \\ & =\frac{1}{\varphi ^2}(\frac{1}{2\stackrel{~}{x}_2})\frac{}{\stackrel{~}{x}_2}+0=\varphi _{}(\frac{1}{2x_2}\frac{}{x_2})|_{(\stackrel{~}{x}_1,\stackrel{~}{x}_2)}.\hfill \end{array}\hfill \end{array}$$
We see easily that the invariance of the connection under $`\varphi `$ is equivalent to $`\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}=0`$, i.e. $`\delta _2^{}(\varphi )=0`$. $`\mathrm{}`$
### 4.2 The general case
For the general case of nontrivial $`\delta _2^{}`$, we look at the following connection.
$$\begin{array}{cc}_{\frac{}{x_1}}\frac{}{x_1}=\mu (x_1,x_2)\frac{}{x_2},\hfill & _{\frac{}{x_1}}\frac{}{x_2}=\frac{1}{2x_2}\frac{}{x_1},\hfill \\ & \\ _{\frac{}{x_2}}\frac{}{x_1}=\frac{1}{2x_2}\frac{}{x_1},\hfill & _{\frac{}{x_2}}\frac{}{x_2}=\frac{1}{2x_2}\frac{}{x_2}.\hfill \end{array}$$
(12)
Here $`\mu `$ is a suitable function.
###### Theorem 4.2
Let $`\mathrm{\Gamma }`$ be a pseudogroup generated by local diffeomorphisms on $``$ acting on $`\times ^+`$ by $`\varphi :x_1\varphi (x_1),x_2\frac{x_2}{\varphi ^{}(x_1)}`$, $`\varphi \mathrm{\Gamma }`$. Assume that the dimension of the fixed point set of each element $`\varphi \mathrm{\Gamma }`$ is strictly less than 2. The connection $``$ in (12) is invariant under $`\mathrm{\Gamma }`$ if and only if the $`_1`$ action on the corresponding groupoid algebra $`\mathrm{\Gamma }C_c^{\mathrm{}}(\times ^+)`$ is projective.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$ Given a local diffeomorphism $`\varphi `$, we have the following quantity different from the proof of Proposition 4.1. All the others are same.
$$\begin{array}{cc}_{\varphi _{}(\frac{}{x_1})}\varphi _{}(\frac{}{x_1})\hfill & =_{\varphi ^{}(x_1)\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}}(\varphi ^{}(x_1)\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2})\hfill \\ & =\varphi ^2_{\frac{}{\stackrel{~}{x}_1}}\frac{}{\stackrel{~}{x}_1}+\varphi ^{}\frac{}{\stackrel{~}{x}_1}(\varphi ^{})\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^{}}x_2_{\frac{}{\stackrel{~}{x}_1}}\frac{}{\stackrel{~}{x}_2}\hfill \\ & \varphi ^{}\frac{}{\stackrel{~}{x}_1}(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)\frac{}{\stackrel{~}{x}_2}\frac{\varphi ^{\prime \prime }}{\varphi ^{}}x_2_{\frac{}{\stackrel{~}{x}_2}}\frac{}{\stackrel{~}{x}_1}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}(\varphi ^{})\frac{}{\stackrel{~}{x}_1}\hfill \\ & +(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)^2_{\frac{}{\stackrel{~}{x}_2}}\frac{}{\stackrel{~}{x}_2}+\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{}{\stackrel{~}{x}_2}(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)\frac{}{\stackrel{~}{x}_2}\hfill \\ & =\varphi ^2\mu (\stackrel{~}{x}_1,\stackrel{~}{x}_2)\frac{}{\stackrel{~}{x}_2}+\varphi ^{}\frac{1}{\varphi ^{}}(\varphi ^{\prime \prime })\frac{}{\stackrel{~}{x}_1}\varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}\hfill \\ & \varphi ^{}[\frac{1}{\varphi ^{}}\frac{\varphi ^{\prime \prime \prime }\varphi ^22\varphi ^{\prime \prime 2}\varphi ^{}}{(\varphi ^2)^2}x_2+(\frac{\varphi ^{\prime \prime }}{\varphi ^2})^2x_2]\frac{}{\stackrel{~}{x}_2}\varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_1}\hfill \\ & +(\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2)^2\frac{1}{2\stackrel{~}{x}_2}\frac{}{\stackrel{~}{x}_2}+\frac{\varphi ^{\prime \prime }}{\varphi ^2}x_2\varphi ^{}\frac{\varphi ^{\prime \prime }}{\varphi ^2}\frac{}{\stackrel{~}{x}_2}\hfill \\ & =[\varphi ^2\mu (\stackrel{~}{x}_1,\stackrel{~}{x}_2)\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^3}x_2]\frac{}{\stackrel{~}{x}_2}.\hfill \end{array}$$
By the invariance of $``$, we have
$$[\varphi ^2\mu (\stackrel{~}{x}_1,\stackrel{~}{x}_2)\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^3}x_2]\frac{}{\stackrel{~}{x}_2}=\varphi _{}(\mu (x_1)\frac{}{x_2})=\mu (x_1,x_2)\frac{1}{\varphi ^{}}\frac{}{\stackrel{~}{x}_2},$$
and
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^3}x_2=\varphi ^2\mu (\varphi (x_1),\frac{x_2}{\varphi ^{}})\frac{1}{\varphi ^{}}\mu (x_1,x_2).$$
(13)
By Equation (13), we have
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^2}x_2^2=\varphi ^4\stackrel{~}{x}_2\mu (\varphi (x_1),\frac{x_2}{\varphi ^{}})x_2\mu (x_1,x_2).$$
(14)
1. $``$. Let $`\varphi `$ be an element in $`\mathrm{\Gamma }`$.
We introduce $`\nu =\frac{\mu (x_1,x_2)}{x_2}`$, and Equation (14) is equivalent to
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^2}=\varphi ^2\nu (\varphi (x_1),\frac{x_2}{\varphi ^{}})\nu (x_1,x_2).$$
Define $`\omega (x_1,x_2)=\nu (x_1,\frac{1}{x_2})`$, and we have
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^2}=\varphi ^2\nu (\varphi (x_1),\frac{x_2}{\varphi ^{}})\nu (x_1,x_2)=\varphi ^2\omega (\varphi (x_1),\frac{\varphi ^{}}{x_2})\omega (x_1,\frac{1}{x_2}).$$
Introduce $`y=\frac{1}{x_2}`$, the above equation gives
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^2}=\varphi ^2\omega (\varphi (x_1),\varphi ^{}y)\omega (x,y).$$
(15)
Finally, letting $`\mathrm{\Omega }(x,y)=y^2\omega (x,y),x_1=x`$, we see that Equation (15) implies
$$\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi ^{\prime \prime 2}}{\varphi ^2}y^2=\varphi ^2y^2\omega (\varphi (x_1),\varphi ^{}y)\omega (x,y)y^2=(\varphi ^1)^{}(\mathrm{\Omega })(x,y)\mathrm{\Omega }(x,y).$$
The left hand side of the above equation is equal to the expression of $`\delta _2^{}(\varphi ^1)`$. The above equality shows that $`\delta _2^{}`$ is inner when we consider the $`_1`$ action on the foliation groupoid $`FX𝒢`$ as in Section 2.1.
2. $``$. Suppose that the $`_1`$ action on $`\mathrm{\Gamma }C_c^{\mathrm{}}(\times ^+)`$ is projective.
We first show that if the $`_1`$ action is projection on $`\mathrm{\Gamma }C_c^{\mathrm{}}(\times ^+)`$, the support of $`\mathrm{\Omega }`$ has to be on the unit space. We write $`\mathrm{\Omega }=_{\alpha \mathrm{\Gamma }}\mathrm{\Omega }_\alpha U_\alpha `$ and $`\delta _2^{}(U_\varphi )U_\varphi =[\mathrm{\Omega },U_\varphi ]`$, and have the following observations.
1. From $`\delta _i(\mathrm{\Omega })=0,i>0`$, we know that $`\delta _i(U_\alpha )\mathrm{\Omega }_\alpha =0,\alpha `$.
2. From $`\delta _i(f)=0`$ for any $`fC_c^{\mathrm{}}(\times ^+)`$, we have that $`[\mathrm{\Omega },f]=_{\alpha \mathrm{\Gamma }}(\alpha ^{}(f)f)\mathrm{\Omega }_\alpha U_\alpha `$. Therefore $`(\alpha ^{}(f)f)\mathrm{\Omega }_\alpha =0`$, for all $`\alpha \mathrm{\Gamma }`$.
For a given $`\alpha \mathrm{\Gamma }`$ not equal to identity, we have that $`\delta _i(U_\alpha )\mathrm{\Omega }_\alpha =0,i>0`$ and $`(\alpha ^{}(f)f)\mathrm{\Omega }_\alpha =0`$. If there is $`x_0\times ^+`$ such that $`\mathrm{\Omega }_\alpha (x_0)0`$, then at $`x_0`$, there is a neighborhood $`N`$ of $`x_0`$ on which $`\delta _i(U_\alpha )=0`$. In particular $`\delta _1(U_\alpha )=\mathrm{log}((\alpha ^1)^{^{}})^{}=0`$. Solving this differential equation, we know that $`\alpha `$ on $`N`$ must act like $`\alpha :(x_1,x_2)(ax_1+b,ax_2)`$. By the fact that $`(\alpha ^{}(f)f)\mathrm{\Omega }_\alpha (x_0)=0`$ on $`N`$, for any smooth function, we know that $`\alpha (x_0)=x_0`$. The same argument show that all $`xN`$ has to be fixed by $`\alpha `$, since $`\mathrm{\Omega }_\alpha (x)0`$. But this contradicts our assumption that the fixed point set of $`\alpha `$ is at most 1 dimensional. This shows that $`\mathrm{\Omega }_\alpha =0`$.
From the above argument, we know that $`\mathrm{\Omega }`$ has to be supported on the unit space. At this time, the projective condition is equivalent to
$$\delta _2(\varphi ^1)=y^2\frac{\varphi ^{\prime \prime \prime }\varphi ^{}\frac{3}{2}\varphi _{}^{\prime \prime }{}_{}{}^{2}}{\varphi ^2}U_\varphi =(\mathrm{\Omega }\varphi ^{}(\mathrm{\Omega }))U_\varphi .$$
From (15) and the transformation there, we know that the existence of $`\mathrm{\Omega }`$ implies the existence of an invariant connection like (12). $`\mathrm{}`$
###### Remark 4.3
Here, for calculation convenience, we have identified the Frame bundle $`F`$ with the cotangent bundle $`T^{}`$ by $`\tau :(x,y)(x,\frac{1}{y})`$. The connection $``$ is defined on $`T^{}`$. By $`\tau `$, it is also defined on $`F`$.
In Theorem 4.2, the assumption that the fixed point set of any element in $`\mathrm{\Gamma }`$ is at most one dimensional is only used in the sufficient part of the proof. Generally, $`\mathrm{\Omega }`$ is supported on the fixed point set $`B^{(0)}`$ of $`\mathrm{\Gamma }`$, i.e. $`\{(\gamma ,x)|\gamma \mathrm{\Gamma },\gamma (x)=x\}`$. $`\mathrm{\Gamma }`$ acts on $`B^{(0)}`$, by conjugation action. The similar result of Theorem 4.2 is extended to this general situation without any extra effort.
Theorem 4.2’ Let $`\mathrm{\Gamma }`$ be a pseudogroup generated by local diffeomorphisms on $``$ and $`B^{(0)}=\{(\gamma ,x)\mathrm{\Gamma }\times \times ^+|\gamma x=x\}`$ be the fixed point set. The projective action $`(\rho ,\mathrm{\Omega })`$ of $`_1`$ on $`\mathrm{\Gamma }C_c^{\mathrm{}}(\times ^+)`$ is one to one correspondent to a $`\mathrm{\Gamma }`$ invariant connection $``$ on $`\times ^+`$ of form (12) and a smooth function $`f`$ on $`\mathrm{\Gamma }\times \times ^+`$, which is supported on $`B^{(0)}\{(id,x)|x\times ^+\}`$ and invariant under $`\mathrm{\Gamma }`$ conjugation action.
## 5 Universal deformation formula for $`_1`$
In this section, we will use a Fedosov type construction to reconstruct the universal deformation formula of $`_1`$ originally constructed by Connes and Moscovici .
### 5.1 Zagier’s deformation
In this subsection, we discuss the influence of the above new connection (12) on the star product (2).
###### Corollary 5.1
The connection $``$ (12) is flat if and only if $`\mu (x_1,x_2)=x_2\nu (x_1)`$, where $`\nu (x_1)`$ is an arbitrary smooth function on $``$.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$ The curvature of $``$ can be directly calculated to be equal to
$$\begin{array}{cc}R(\frac{}{x_1},\frac{}{x_2})(\frac{}{x_1})\hfill & =(\frac{\mu }{x_2}\frac{\mu }{x_2})\frac{}{x_2}\hfill \\ R(\frac{}{x_1},\frac{}{x_2})(\frac{}{x_2})\hfill & =0.\hfill \end{array}$$
Therefore, $`R=0`$ if and only if $`\frac{\mu }{x_2}\frac{\mu }{x_2}=0`$. The solution of this first order differential equation is that $`\mu =x_2\nu (x_1)`$, where $`\nu (x_1)`$ is an arbitrary smooth function on $``$. $`\mathrm{}`$
In this section, we restrict ourselves to the case that the connection (12) is flat, which means that $`\mu (x_1,x_2)=x_2\nu (x_1)`$. We consider the deformation quantization of $`(\times ^+,dx_1dx_2)`$ using this connection.
The Christoffel symbols of the connection $`^{\stackrel{~}{𝒪}}`$ are calculated as follows,
$$\mathrm{\Gamma }_{11}^1=\mathrm{\Gamma }_{12}^2=\mathrm{\Gamma }_{21}^2=\mathrm{\Gamma }_{22}^1=0,\mathrm{\Gamma }_{11}^2=\mu ,\mathrm{\Gamma }_{12}^1=\mathrm{\Gamma }_{21}^1=\frac{1}{2x_2},\mathrm{\Gamma }_{22}^2=\frac{1}{2x_2}.$$
Taking the Formula (5.1.8) in with the same notations, we have
$$\begin{array}{cc}\mathrm{\Gamma }_{111}=\omega _{11}\mathrm{\Gamma }_{11}^1+\omega _{12}\mathrm{\Gamma }_{11}^2=\omega _{12}\mu ,\hfill & \mathrm{\Gamma }_{211}=\omega _{21}\mathrm{\Gamma }_{11}^1+\omega _{22}\mathrm{\Gamma }_{11}^2=0,\hfill \\ \mathrm{\Gamma }_{112}=\omega _{11}\mathrm{\Gamma }_{12}^1+\omega _{12}\mathrm{\Gamma }_{12}^2=0,\hfill & \mathrm{\Gamma }_{121}=\omega _{11}\mathrm{\Gamma }_{21}^1+\omega _{12}\mathrm{\Gamma }_{21}^2=0,\hfill \\ \mathrm{\Gamma }_{212}=\omega _{21}\mathrm{\Gamma }_{12}^1+\omega _{22}\mathrm{\Gamma }_{12}^2=\frac{1}{2x_2}\omega _{21},\hfill & \mathrm{\Gamma }_{221}=\omega _{21}\mathrm{\Gamma }_{21}^1+\omega _{22}\mathrm{\Gamma }_{21}^2=\frac{1}{2x_2}\omega _{21},\hfill \\ \mathrm{\Gamma }_{122}=\omega _{11}\mathrm{\Gamma }_{22}^1+\omega _{12}\mathrm{\Gamma }_{22}^2=\frac{1}{2x_2}\omega _{12},\hfill & \mathrm{\Gamma }_{222}=\omega _{21}\mathrm{\Gamma }_{22}^1+\omega _{22}\mathrm{\Gamma }_{22}^2=0.\hfill \end{array}$$
We have the following expression for $`\mathrm{\Gamma },\mathrm{\Gamma }a,a\mathrm{\Gamma }`$, and $`[\mathrm{\Gamma },a]`$.
$$\mathrm{\Gamma }=\frac{1}{2}\omega _{21}\{[\mu (u^1)^2+\frac{1}{2}(2u^2)^2]dx_1+\frac{1}{2}2u^1u^2dx_2\},$$
and
$$\begin{array}{cc}\frac{i}{h}[\mathrm{\Gamma },a]\hfill & =(\frac{1}{2}(\mu )2a_{m,n}(u^1)^mn(u^2)^{n1}\frac{1}{4x_2}2a_{m,n}m(u^1)^{m1}(u^2)^{n+1})dx_1\hfill \\ & +\frac{1}{4x_2}(2a_{m,n}(u^1)^mn(u^2)^n2a_{m,n}m(u^1)^m(u^2)^n)dx_2.\hfill \end{array}$$
It is a direct check that when $`\mu =x_1\nu (x_2)`$, $`^2`$ and $`D^2`$ are both 0. By Theorem 2.1, for each $`fC^{\mathrm{}}(\times _+)[[\mathrm{}]]`$, there is a unique solution of the equation $`Da=0`$ with $`a_{0,0}=f`$. In the following, we calculate the explicit expression of $`a`$.
The expression of $`Da`$ is calculated as follows.
$$\begin{array}{cc}Da\hfill & =a\delta a=\delta a+da+\frac{i}{h}[\mathrm{\Gamma },a]\hfill \\ & =a_{m,n}m(u^1)^{m1}(u^2)^ndx_1a_{m,n}(u^1)^mn(u^2)^{n1}dx_2\hfill \\ & +\frac{a_{m,n}}{x_1}(u^1)^m(u^2)^ndx_1+\frac{a_{m,n}}{x_2}(u^1)^m(u^2)^ndx_2\hfill \\ & +[\mu a_{m,n}n(u^1)^{m+1}(u^2)^{n1}\frac{a_{m,n}}{2x_2}m(u^1)^{m1}(u^2)^{n+1}]dx_1\hfill \\ & +\frac{a_{m,n}}{2x_2}(nm)(u^1)^m(u^2)^ndx_2.\hfill \end{array}$$
The equation $`Da=0`$ gives the following system of differential equations:
$$a_{m+1,n}(m+1)+\frac{a_{m,n}}{x_1}(n+1)\mu a_{m1,n+1}\frac{a_{m+1,n1}}{2x_2}(m+1)=0,$$
and
$$a_{m,n+1}(n+1)+\frac{a_{m,n}}{x_2}+\frac{a_{m,n}}{2x_2}(nm)=0.$$
Given $`a_{0,0}=f`$, we solve the system of equations by induction.
$$\begin{array}{cc}a_{m,0}\hfill & =\frac{1}{m}(\frac{a_{m1,0}}{x_1}\mu a_{m2,1})=\frac{1}{m}(\frac{a_{m1,0}}{x_1}\mu (\frac{}{x_2}\frac{m2}{2x_2})a_{m2,0}),\hfill \\ a_{m,n}\hfill & =\frac{1}{n!}(\frac{}{x_2}\frac{m}{2x_2})\mathrm{}(\frac{}{x_2}+\frac{nm1}{2x_2})a_{m,0}.\hfill \end{array}$$
If we set
$$\begin{array}{cc}X\hfill & =\frac{1}{x_2}\frac{}{x_1},\hfill \\ Y\hfill & =x_2\frac{}{x_2},\hfill \end{array}$$
it is direct check that
$$\begin{array}{cc}A_{m+1}\hfill & =XA_mm\frac{\mu }{x_2^3}(Y\frac{m1}{2})A_{m1},\hfill \\ B_{m+1}\hfill & =XB_mm\frac{\mu }{x_2^3}(Y\frac{m1}{2})B_{m1},\hfill \\ a_{m,n}\hfill & =\frac{(1)^nx_2^{mn}}{n!}\frac{A_m}{m!}(Y+\frac{m}{2})\mathrm{}(Y+\frac{m+n1}{2})a,\hfill \\ b_{n,m}\hfill & =\frac{(1)^mx_2^{nm}}{m!}\frac{B_n}{n!}(Y+\frac{n}{2})\mathrm{}(Y+\frac{m+n1}{2})b.\hfill \end{array}$$
The above expression of $`A_m,B_m`$ is exactly identical to the recurrence relation as described in (2.9) of of Connes and Moscovici with $`S(X)=X`$, and $`\mathrm{\Omega }=\frac{\mu }{x_2^3}=\frac{\nu }{x_2^2}`$. The star product constructed in this way defines the Zagier’s deformation for $`h_1`$ constructed from Rankin-Cohen brackets on modular forms with a forth degree element.
###### Remark 5.2
For computation reasons, we have chosen that a special form of connections defined by Equation (12), which is flat. Because of the flatness, the calculation is quite simple and transparent. When the connection is not flat, Fedosov’s construction still works, but the calculation is much more complicated. However, the star product should be able to be expressed by the same formula.
###### Remark 5.3
As explained in Remark 4.3, the connection and the star product discussed in this subsection are both on the cotangent bundle $`T^{}`$. However, all these constructions can be pulled back to the frame bundle by $`\tau `$(See Remark 4.3) without any difficulty.
### 5.2 Full injectivity
We have shown in the last subsection that the deformation quantization of the standard symplectic structure on the upper half plane using the connection (12) with $`\mu (x_1,x_2)=x_2\nu (x_1)`$ gives rise to Zagier’s deformation formula on modular forms. To generalize this deformation to a universal deformation formula of a projective $`_1`$ action, we adapt the method used by Connes and Moscovici \[Sec. 3\] to our situation. We briefly recall their construction in the following, and refer to for the detail.
Firstly, we introduce a free abelian algebra $`P`$ with a set of generators indexed by $`_0`$, $`Z_0,Z_1,\mathrm{},Z_n,\mathrm{}`$. On $`P`$, we define a $`_1`$ action as follows,
$$Y(Z_j)\stackrel{def}{=}(j+2)Z_j,X(Z_j)\stackrel{def}{=}Z_{j+1},\delta _k(p)=0,pP,j0.$$
Secondly, we consider the crossed product algebra $`\stackrel{~}{_1}\stackrel{def}{=}P_1P`$, which is equal to $`P_1P`$ as a vector space. Denote this algebra by $`\stackrel{~}{_1}`$. Connes and Moscovici defines on $`\stackrel{~}{_1}`$ an Hopf algebra structure over $`P`$, with $`\alpha ,\beta :P\stackrel{~}{_1}`$ defined by
$$\alpha (p)=p11,\beta (q)=11q,p,qP.$$
Thirdly, to deal with the projective structure, we define $`\stackrel{~}{\delta _2^{}}\stackrel{def}{=}\delta _2\frac{1}{2}\delta _2\alpha (Z_0)+\beta (Z_0)`$, $`\stackrel{~}{_s}`$ as the quotient of $`\stackrel{~}{_1}`$ by the ideal generated by $`\stackrel{~}{\delta _2}^{}`$. $`\stackrel{~}{}_s`$ is still a Hopf algebra over $`P`$ because $`\mathrm{\Delta }(\stackrel{~}{\delta _2^{}})=\stackrel{~}{\delta _2}^{}1+1\stackrel{~}{\delta _2}^{}`$.
Fixing a function $`\mu (x_1,x_2)`$, we consider a pseudogroup $`\mathrm{\Gamma }`$ action on $``$ whose lifting onto $`T^{}`$ preserves the connection $``$ (12) defined by $`\mu `$. By Theorem 4.2, the $`_1`$ action on the corresponding groupoid algebra $`𝒜_{\mu ,\mathrm{\Gamma }}\stackrel{def}{=}C_c^{\mathrm{}}(\times ^+)\mathrm{\Gamma }`$ is projective with $`\mathrm{\Omega }`$ defined in the proof.
We define $`\rho _{\mu ,\mathrm{\Gamma }}:P𝒜_{\mu ,\mathrm{\Gamma }}`$ by $`\rho (Z_k)=X^k(\mathrm{\Omega })`$ and make $`𝒜_{\mu ,\mathrm{\Gamma }}`$ into a module algebra over $`\stackrel{~}{}_1|P`$ by
$$\chi _{\mu ,\mathrm{\Gamma }}(phq)(U_\gamma f)\stackrel{def}{=}\rho _{\mu ,\mathrm{\Gamma }}(p)h(U_\gamma f)\rho _{\mu ,\mathrm{\Gamma }}(q).$$
One easily checks that $`𝒜_{\mu ,\mathrm{\Gamma }}`$ becomes a module algebra over $`\stackrel{~}{}_s|P`$ because when the $`_1`$ action is projective, $`\stackrel{~}{\delta _2^{}}`$ acts as $`0`$.
We define action $`\chi _{\mu ,\mathrm{\Gamma }}^n`$,
$$\chi _{\mu ,\mathrm{\Gamma }}^{(n)}:\underset{n}{\underset{}{\stackrel{~}{}_s_P\mathrm{}\stackrel{~}{}_s}}(\underset{n}{\underset{}{𝒜_{\mu ,\mathrm{\Gamma }}\mathrm{}𝒜_{\mu ,\mathrm{\Gamma }}}},𝒜_{\mu ,\mathrm{\Gamma }})$$
by means of acting on each components, where $``$ means the set of linear maps.
We fix $`\mu =x_1\nu (x_1)`$, and have the following Proposition analogous to \[Prop. 12\].
###### Proposition 5.4
For each $`n`$, $`_{\nu (x_1),\mathrm{\Gamma }}Ker\chi _{x_2\nu (x_1),\mathrm{\Gamma }}^{(n)}=0`$.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$ There is no difference between the proofs for different $`n`$. Therefore, for simplicity, we only prove the proposition for $`n=1`$.
Following the proof of \[Prop. 12\], an arbitrary element of $`\stackrel{~}{}_s`$ can be written uniquely as a finite sum of the form
$$H=\underset{j,k,l,m}{}\alpha (p_{jklm})\beta (q_{jklm})\delta _1^jX^kY^l,$$
where $`p,qP`$.
Let $`\chi _{x_2\nu (x_1),\mathrm{\Gamma }}(H)=0`$, for arbitrary $`\nu (x_1)`$ and pseudogroup $`\mathrm{\Gamma }`$ preserving the connection defined by $`x_2\nu (x_1)`$. From the proof of Theorem 4.2, we know that in this case, $`\mathrm{\Omega }=x_2^2\nu (x_1)`$.
If $`U_\gamma f𝒜_{x_2\nu (x_2),\mathrm{\Gamma }}`$, then
$$\underset{j,k,l,m}{}\rho _{x_1\nu (x_2),\mathrm{\Gamma }}(p_{jklm})\gamma ^{}(\rho _{x_1\nu (x_2)}(q_{jklm}))\delta _1(\gamma )^jX^kY^l(f).$$
We notice that $`f`$ can be arbitrary smooth function on $`\times ^+`$, and $`X^kY^l=x_2^{m+l}\frac{d^k}{dx_1^m}\frac{d^l}{dx_2^l}`$. This implies that
$$\underset{j,m}{}\rho _{x_1\nu (x_2),\mathrm{\Gamma }}(p_{jklm})\gamma ^{}(\rho _{x_1\nu (x_2)}(q_{jklm}))\delta _1(\gamma )^j=0,$$
for any $`l,m`$.
To prove the Proposition, we consider the following family of algebras, $`𝒜_{x_2\nu (x_2),\mathrm{\Gamma }}`$.
Fix a diffeomorphism $`\varphi _{O_1,O_2}`$ from an open set $`O_1`$ to the other open set $`O_2`$, with $`O_1`$ disjoint from $`O_2`$. The disjointness between $`O_1`$ and $`O_2`$ makes the set $`\mathrm{\Gamma }_\varphi \stackrel{def}{=}\{id|_{},id|_{O_1},id|_{O_2},\varphi ,\varphi ^1,\}`$ into a pseudogroup. Starting with any connection $`_1`$ of the form (12) with $`\mu =x_2\nu (x_1)`$ on $`O_1`$, we first push forward this connection to $`O_2`$ by $`\varphi `$, and then extend the connections defined on $`O_1`$ and $`O_2`$ to a global connection $`\stackrel{~}{}`$ on $`\times ^+`$. The extension of the connection is well defined because $`O_1`$ is disjoint from $`O_2`$, (we may need to restrict to a smaller open subset $`O_2^{}`$ of $`O_2`$ by a cutoff function) and is $`\mathrm{\Gamma }_\varphi `$ invariant by its definition. According to our construction, we have that $`\stackrel{~}{_s}`$ act on the corresponding groupoid algebra $`𝒜_{\varphi _{O_1,O_2},\stackrel{~}{}}`$.
Now at any $`x`$, we fix $`O_1`$ containing $`x`$, and let $`O_2`$, $`\varphi `$, $`_1`$ vary. It is not hard to see that if $`H`$ vanishes on this family of algebra $`𝒜_{\varphi _{O_1,O_2},\stackrel{~}{}}`$, we must have that $`H`$ vanishes at $`x`$, because $`H`$ has only finite number of terms but this family of algebras has infinitely many freedoms. Hence $`H`$ has to be equal to 0. $`\mathrm{}`$
### 5.3 Universal deformation $`_1`$ with a projective structure
We consider the groupoid algebra $`𝒜_{x_2\nu (x_1),\mathrm{\Gamma }}`$. Because the connection defined by $`x_1\nu (x_1)`$ in (12) is $`\mathrm{\Gamma }`$ invariant, the results in Section 2.3 implies that the symplectic form $`\frac{\mathrm{d}x\mathrm{d}y}{y^2}`$ on $`\times _+`$, which is invariant under any $`\mathrm{\Gamma }`$, defines a noncommutative Poisson structure on $`C_c^{\mathrm{}}(\times ^+)\mathrm{\Gamma }`$. Furthermore, we extend this Poisson structure to a deformation of $`C_c^{\mathrm{}}(\times ^+)\mathrm{\Gamma }`$. This deformation can be realized by the crossed product of the star product constructed in Section 5.1 with $`\mathrm{\Gamma }`$.
In Section 5.1, the $``$ product is expressed as follows: for $`f,gC_c^{\mathrm{}}(\times _+)`$,
$$\begin{array}{cc}fg\hfill & =_{n=0}^{\mathrm{}}\mathrm{}^n_{k=0}^n\frac{A_k}{k!}(2Yk)_{nk}(a)\frac{B_{nk}}{(nk)!}(2Yn+k)_k(b)\hfill \\ A_{m+1}\hfill & =XA_mmx_2\mu (Y\frac{m1}{2})A_{m1}=XA_mm\mathrm{\Omega }(Y\frac{m1}{2})A_{m1},\hfill \\ B_{m+1}\hfill & =XB_mmx_2\mu (Y\frac{m1}{2})B_{m1}=XB_mm\mathrm{\Omega }(Y\frac{m1}{2})B_{m1}.\hfill \end{array}$$
The crossed product of $``$ with $`\mathrm{\Gamma }`$ is written as $`f_\gamma U_\gamma g_\beta U_\beta \stackrel{def}{=}f_\gamma \gamma ^{}(g_\beta )U_{\gamma \beta }`$ defines a deformation quantization of $`C_c^{\mathrm{}}(\times ^+)\mathrm{\Gamma }`$.
According to the formulas of $``$ and the $`\mathrm{\Gamma }`$ crossed product, the deformed product $``$ on $`C_c^{\mathrm{}}(\times ^+)\mathrm{\Gamma }`$ can be expressed by $`\stackrel{~}{}_s`$ as follows,
$$\begin{array}{cc}RC\hfill & =_{n=0}^{\mathrm{}}\mathrm{}^n_{k=0}^n\frac{A_k}{k!}(2Y+k)_{nk}\frac{B_{nk}}{(nk)!}(2Y+nk)_k\hfill \\ A_{m+1}\hfill & =S(X)A_mm\mathrm{\Omega }^0(Y\frac{m1}{2})A_{m1},\hfill \\ B_{m+1}\hfill & =XB_mm\mathrm{\Omega }(Y\frac{m1}{2})B_{m1},\hfill \end{array}$$
where $`\mathrm{\Omega }^0`$ is the right multiplication of $`\mathrm{\Omega }`$.
By Proposition 5.4, we conclude $`RC`$ can be pulled back to $`\stackrel{~}{}_s`$ and defines an associative universal deformation for any projective $`_1`$ actions.
## 6 Deformation without Projective structures—noncommutative Poisson structure
In the above deformation (4), we have assumed the action to be projective. One can ask whether one can go beyond this. Recently, a construction of Bressler, Gorokhovsky, Nest, and Tsygan strongly suggests that this general RC deformation may still exist.
In this section, we look at the first order approximation of the general deformation. We prove that $`RC_1`$ generally defines a noncommutative Poisson structure without any assumptions.
###### Proposition 6.1
For an $`_1`$ action on an $`A`$, $`RC_1=X2Y+2YX+\delta _1Y2Y`$ defines a noncommutative Poisson structure on $`A`$.
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}.`$ The proof of this proposition is calculation. We need to find an element $`B`$ in $`_1_1`$, such that for any $`a,b,cA`$,
$$aB(b,c)B(ab,c)+B(a,bc)B(a,b)c=RC_1(RC_1(a,b),c)RC_1(a,RC_1(b,c)).$$
In order to find such a $`B`$, we first look at the special case where the Hopf algebra action is projective. In this case, the associativity of the Connes-Moscovici’s universal deformation formula of $`_1`$ implies that $`RC_2`$ is a right choice of $`B`$.
For a general $`_1`$ action, we first look at the following term
$$B^{}=S(X)^2Y(2Y+1)+S(X)(2Y+1)X(2Y+1)+Y(2Y+1)X^2.$$
We calculate the difference between the Hochschild coboundary of $`B^{}`$ and $`[RC_1,RC_1]`$.
$$\begin{array}{cc}& (b(B^{})[RC_1,RC_1])(a,b,c)\hfill \\ =\hfill & 4Ya\delta _2^{}YbYc+2Y^2a\delta _2^{}bYc+2Ya\delta _2^{}bYc+2Ya\delta _2^{}bY^2c\hfill \\ & \\ =\hfill & 2[a\delta _2^{}Y^2bYc\delta _2^{}Y^2(ab)Yc+\delta _2^{}Y^2aY(bc)\delta _2^{}Y^2a(Yb)c]4\delta _2^{}YaYbYc\hfill \\ & 2\delta _2^{}aY^2bYc+2Ya\delta _2^{}aY^2c+2Ya\delta _2^{}aYc\hfill \\ & \\ =\hfill & 2[a\delta _2^{}Y^2bYc\delta _2^{}Y^2(ab)Yc+\delta _2^{}Y^2aY(bc)\delta _2^{}Y^2a(Yb)c]\hfill \\ & 2[a\delta _2^{}YbY^2c\delta _2^{}Y(ab)Y^2c+\delta _2^{}YaY^2(bc)\delta _2^{}Ya(Y^2b)c]\hfill \\ & 2\delta _2^{}aY^2bYc2\delta _2^{}aYbY^2c+2Ya\delta _2^{}bYc\hfill \\ & \\ =\hfill & 2[a\delta _2^{}Y^2bYc\delta _2^{}Y^2(ab)Yc+\delta _2^{}Y^2aY(bc)\delta _2^{}Y^2a(Yb)c]\hfill \\ & 2[a\delta _2^{}YbY^2c\delta _2^{}Y(ab)Y^2c+\delta _2^{}YaY^2(bc)\delta _2^{}Ya(Y^2b)c]\hfill \\ & \frac{2}{3}[a\delta _2^{}bY^3c\delta _2^{}(ab)Y^3c+\delta _2^{}aY^3(bc)+\delta _2^{}aY^3(b)c]+2Ya\delta _2^{}bYc\hfill \\ =\hfill & 2[a\delta _2^{}Y^2bYc\delta _2^{}Y^2(ab)Yc+\delta _2^{}Y^2aY(bc)\delta _2^{}Y^2a(Yb)c]\hfill \\ & 2[a\delta _2^{}YbY^2c\delta _2^{}Y(ab)Y^2c+\delta _2^{}YaY^2(bc)\delta _2^{}Ya(Y^2b)c]\hfill \\ & \frac{2}{3}[a\delta _2^{}bY^3c\delta _2^{}(ab)Y^3c+\delta _2^{}aY^3(bc)+\delta _2^{}aY^3(b)c]\hfill \\ & 2[a\delta _2^{}YbYc\delta _2^{}Y(ab)Yc+\delta _2^{}YaY(bc)\delta _2^{}YaY(b)c]\hfill \\ & [a\delta _2^{}bY^2c\delta _2^{}(ab)Y^2c+\delta _2^{}aY^2(bc)\delta _2^{}aY^2(b)c],\hfill \end{array}$$
where $`b(B^{})`$ is the Hochschild coboundary of $`B^{}`$ and $`\delta _2^{}=\delta _2\frac{1}{2}\delta _1^2`$.
It is straightforward to check the following identities.
$$\begin{array}{cc}\hfill b(\delta _2^{}Y^2Y)(a,b,c)& =a\delta _2^{}Y^2bYc\delta _2^{}Y^2(ab)Yc+\delta _2^{}Y^2aY(bc)\delta _2^{}Y^2a(Yb)c,\hfill \\ \hfill b(\delta _2^{}Y^3)(a,b,c)& =a\delta _2^{}bY^3c\delta _2^{}(ab)Y^3c+\delta _2^{}aY^3(bc)\delta _2^{}a(Y^3b)c,\hfill \\ \hfill b(\delta _2^{}YY)(a,b,c)& =a\delta _2^{}YbYc\delta _2^{}Y(ab)Yc+\delta _2^{}YaY(bc)\delta _2^{}Ya(Yb)c,\hfill \\ \hfill b(\delta _2^{}Y^2)(a,b,c)& =a\delta _2^{}bY^2c\delta _2^{}(ab)Y^2c+\delta _2^{}aY^2(bc)\delta _2^{}aY^2bc\hfill \\ \hfill b(\delta _2^{}YY^2)(a,b,c)& =a\delta _2^{}YbY^2c\delta _2^{}Y(ab)Y^2c+\delta _2^{}YaY^2(bc)\delta _2^{}YaY^2(b)c.\hfill \end{array}$$
Therefore, the calculation suggests the introduction of $`B^{\prime \prime }=2\delta _2^{}+Y^2Y+\frac{2}{3}\delta _2^{}Y^3+2\delta _2^{}YY^2+2\delta _2^{}YY+\delta _2^{}Y^2`$ and $`B=B^{}+B^{\prime \prime }`$. And we have $`b(B)=b(B^{}+B^{\prime \prime })=[RC_1,RC_1]`$. $`\mathrm{}`$
Pierre Bieliavsky, Département de Mathématique, Université Catholique de Louvain, Chemin du cyclotron, 2, 1348 Louvain-La-Neuve, Belgium. Email: bieliavsky@math.ucl.ac.be.
Xiang Tang, Department of Mathematics, Washington University, St. Louis, MO, 63130, U.S.A., Email:xtang@math.wustl.edu.
Yijun Yao, Centre de Mathématiques École Polytechnique, Ecole Polytechnique, 91128 Palaiseau Cedex, France, Email: yao@math.polytechnique.fr.
|
warning/0506/hep-ex0506041.html
|
ar5iv
|
text
|
# LATEST RESULTS FROM PHOBOS AT RHIC
## 1 Introduction
The study of heavy ion collisions constitutes an important part of the recent experimental and theoretical effort to understand the strong interaction which binds quarks and gluons into hadrons. Predictions based on QCD indicate the existence of a new state of matter dominated by the strong interaction, where bound hadrons no longer exist, provided the energy density is sufficiently high (greater than a GeV/fm<sup>3</sup>). Heavy ion collisions are the only method to create such a high energy density in the laboratory. The current understanding of the phase structure of strongly interacting matter, the properties of the matter in the various phases and the nature of the transition between them is, to a large extent, driven by experiment. One of the important discoveries at the Relativistic Heavy Ion Collider (RHIC) is that an extremely high energy density system is created, where hadronic degrees of freedom cannot be relevant any more. There is evidence for a very significant level of interaction between the constituents of this system, as opposed to earlier expectations.
In this experimental talk we will emphasize a few basic scaling and factorization features of the data collected by the PHOBOS experiment at RHIC, in comparison with earlier measurements. These simple rules highlight common features of collisions of heavy ions (Au+Au) and simpler systems (d+Au, p+p) in a broad range of collision energies ($`\sqrt{s_{_{NN}}}`$=6.7 to 200 GeV). Details of these findings can be found in a volume summarizing the results of the four experiments from the first three years of RHIC, in which our contribution $`^\mathrm{?}`$ is titled ‘The PHOBOS Perspective on Discoveries at RHIC’. To collect the data presented here, we used the magnetic spectrometer and the multiplicity arrays (covering the $`|\eta |<5.4`$ pseudo-rapidity region) of the PHOBOS experiment (described in detail elsewhere $`^\mathrm{?}`$).
## 2 Scaling features in heavy ion data
We concentrate on three topics: transverse momentum ($`p_T`$) and pseudo-rapidity ($`\eta `$) distributions of charged hadrons, and the azimuthal asymmetry of their production which is called elliptic flow. The consensus in the heavy ion community is that in the early stage of high energy heavy ion collisions a new state of strongly interacting matter has been created. The above observables are related to relevant physical characteristics of this newly created matter, such as the suppression of high-$`p_T`$ particles (jet quenching); the initial energy density and entropy after the collision and the boost-invariance of particle production along the colliding beam direction; and the collective motion of particles resulting from secondary interactions and the properties of the equation of state.
Determining the centrality of a heavy ion collision is essential to provide the basis of comparison with more elementary (p+p, p+$`\overline{\mathrm{p}}`$) processes. Instead of the impact parameter, two different quantities are used to quantify the centrality: the number of participant nucleon pairs ($`N_{part}/2`$) and the number of binary nucleon-nucleon collisions ($`N_{coll}`$). Since the nucleon-nucleon cross section increases with collision energy, so does the $`N_{coll}/N_{part}`$ ratio, which also grows with decreasing impact parameter for simple geometrical reasons. These two quantities are calculated by measuring the charged hadron multiplicity in various regions of $`\eta `$, combined with a comprehensive Monte-Carlo simulation including the Glauber model. $`^\mathrm{?}`$
If one normalizes the transverse momentum distribution of charged hadrons measured in different centrality bins to the $`p_T`$ distribution measured in p+p($`\overline{\mathrm{p}}`$) collisions and also divides by $`N_{coll}`$, one observes a gradual decrease of this quantity ($`R_{AA}`$) with decreasing impact parameter (Fig. 1). Note that the expectation for ‘hard’ collisions with large momentum transfer and no re-interaction with the created medium would be a constant $`R_{AA}`$ equal to unity. However, the $`R_{AA}^{N_{part}}`$ quantity (where we replaced $`N_{coll}`$ by $`N_{part}`$ in the denominator) scales with centrality much more precisely in Au+Au collisions at $`\sqrt{s_{_{NN}}}`$=62.4 and 200 GeV, $`^\mathrm{?}`$ pointing to new physical interpretations. $`^\mathrm{?}`$ More strikingly, the $`p_T`$ spectra normalized to the most central bin at each energy, and also normalized by $`N_{part}`$, as illustrated in the bottom row of Fig. 1, are identical within errors in each of our centrality bins. This is a clear factorization of collision energy and centrality. It will be interesting to compare these conclusions to the data from half a billion Cu+Cu events taken in the present RHIC run, since, for the same number of participants, the collision zone will have a very different geometry in the Cu+Cu and the Au+Au events.
Another simple feature of the data is the extended longitudinal scaling of pseudo-rapidity ($`\eta `$) density distributions of charged particles, in p+Emulsion and d+Au, $`^\mathrm{?}`$ as shown in Fig. 2. By plotting these distributions as a function of $`\eta ^{}=\eta \pm y_{target}`$, thus effectively looking at them from the rest frame of one of the colliding beams, we observe that the data at various energies fall on a common limiting curve, in both reference frames. The longitudinal scaling extends over more than an order of magnitude in beam energy.
A similar scaling was observed earlier in heavy ion (Au+Au) collisions by PHOBOS, illustrated on the right panel of Fig. 2 for the 6% most central data. In addition, applying the fact that $`dN/d\eta `$ at $`\eta 0`$ in central collisions scales logarithmically with $`\sqrt{s_{_{NN}}}`$, we can extrapolate the $`dN/d\eta `$ distribution to LHC energies and give a simple experiment-based prediction at $`\sqrt{s_{_{NN}}}`$=5500 GeV. The prediction gives $`dN/d\eta |_{\eta =0}`$ 1100 and 14000 charged particles in total. $`^\mathrm{?}`$
Longitudinal scaling has also been recently observed in the elliptic flow of particles produced in heavy ion collisions. $`^\mathrm{?}`$ The elliptic flow parameter ($`v_2`$) is a sensitive probe of the properties of the newly created, very dense and hot matter in the early stage of the collision. In Fig. 3, the pseudo-rapidity dependence of this $`v_2`$ parameter is plotted for semi-central Au+Au events and for various energies, where we use the $`\eta ^{}=\eta y_{beam}`$ parameter again. There seems to be a universal curve governing $`v_2`$ over a broad range of $`\eta ^{}`$. This is shown more precisely on the right panel of Fig. 3, where we used the symmetry and plotted $`v_2`$ as a function of $`|\eta |y_{beam}`$.
## 3 Conclusions
We reported a few simple scaling and factorization properties of charged hadron production in heavy ion collisions at RHIC energies. We observed that the transverse momentum spectra approximately scale with the number of participant nucleons, and the energy and centrality dependence of these spectra precisely factorize in the range of centrality and $`p_T`$ we studied. Longitudinal scaling was seen in an extended pseudo-rapidity region in Au+Au (and d(p)+A) collisions, both in case of the $`dN/d\eta `$ pseudo-rapidity density and in case of the $`v_2`$ elliptic flow parameter. These simple features of the data impose constraints on models attempting to describe and understand the basic particle production mechanism in heavy ion collisions.
## Acknowledgments
This work was partially supported by U.S. DOE grants DE-AC02-98CH10886, DE-FG02-93ER40802, DE-FC02-94ER40818, DE-FG02-94ER40865, DE-FG02-99ER41099, and W-31-109-ENG-38, by U.S. NSF grants 9603486, 0072204, and 0245011, by Polish KBN grant 1-P03B-062-27(2004-2007), by NSC of Taiwan Contract NSC 89-2112-M-008-024, and by the Hungarian Scientific Research Fund (OTKA F049823) and the János Bolyai Research Grant.
## References
|
warning/0506/math0506474.html
|
ar5iv
|
text
|
# Examples of quasi-hyperbolic dynamical systems with slow decay of correlations
## 1 Decay of correlation
Let $`\sigma (f)`$ denote the non-negative quantity defined by
$$\sigma ^2(f)=_Xf^2𝑑\mu +2\underset{k1}{}<f,A^kf>.$$
Theorem :If the function $`f`$ is aperiodic and if $`\phi :X\times Y`$ is a centered hölder-continuous function such that the function $`y_X\phi (a,y)𝑑\mu (a)`$ is not a coboundary in the sense of the flow the the number
$$\mathrm{\Sigma }^2(\phi )=_{}_Y\left(_X\phi (a,g_by)𝑑\mu (a)\right)\left(_X\phi (c,y)𝑑\mu (c)\right)𝑑\nu (y)𝑑b$$
is positive and the $`<\phi T^k,\phi >`$ is equivalent to
$$\frac{1}{\sqrt{2\pi k}\sigma (f)}\stackrel{2}{}(\phi ).$$
Proof : We’re interested in
$`<\phi T^k,\phi >`$ $`=`$ $`𝔼(\phi (T^k.,g_{S_kf(.)}.)\phi (.,.))`$
$`=`$ $`{\displaystyle _{X\times Y}}\phi (T^kx,g_{S_{k1}f(x)}y)\phi (x,y)𝑑\mu (x)𝑑\nu (y).`$
Let us consider the function $`F`$ :
$`X\times \times X`$ $``$ $``$
$`(a,b,c)`$ $``$ $`{\displaystyle _Y}\phi (a,g_by)\phi (c,y)𝑑\nu (y)({\displaystyle _Y}\phi (a,y)𝑑\nu (y))({\displaystyle _Y}\phi (c,y)𝑑\nu (y))`$
The correlation $`<\phi T^k,\phi >`$ is
$$_XF(T^kx,S_{k1}f(x),x)𝑑\mu (x)+_X(_Y\phi (T^kx,y)𝑑\nu (y))(_Y\phi (x,y)𝑑\nu (y))𝑑\mu (x).$$
The action of $`A`$ is exponentially mixing thus the second element of this sum tends exponentially fast toward 0.
The action of the flow is exponentially mixing. The function $`F`$ decrease exponentially fast when $`b`$ grows. The local limit theorem of Guivarc’h and Hardy assure that the integral
$$_XF(T^kx,S_{k1}f(x),x)𝑑\mu (x)$$
is equivalent to,
$`{\displaystyle \frac{1}{\sqrt{2k\pi }\sigma (f)}}{\displaystyle _{X\times \times X}}F(a,b,c)𝑑\mu (a)𝑑b𝑑\mu (c)`$
$`={\displaystyle \frac{1}{\sqrt{2\pi k}\sigma (f)}}{\displaystyle _{}}({\displaystyle _X}{\displaystyle _X}{\displaystyle _Y}\phi (a,g_by)\phi (c,y)d\nu (y)d\mu (a)d\mu (c)`$
$`\underset{0}{\underset{}{{\displaystyle _X}{\displaystyle _X}{\displaystyle _Y}\phi (a,y)𝑑\nu (y){\displaystyle _Y}\phi (c,y)𝑑\nu (y)𝑑\mu (a)𝑑\mu (c)}})db`$
$`={\displaystyle \frac{1}{\sqrt{2\pi k}\sigma (f)}}\underset{^2(\phi ).}{\underset{}{{\displaystyle _{}}{\displaystyle _Y}\left({\displaystyle _X}\phi (a,g_by)d\mu (a)\right)\left({\displaystyle _X}\phi (c,y)d\mu (c)\right)d\nu (y)db.}}`$
The quantity $`^2(\phi )`$ is zero if and only if the function $`y_X\phi (a,y)𝑑\mu (a)`$ is a coboundary in the sense of the flow (see for more details) that is if and only if there exists a function $`\phi :Y`$ such that for almost every $`yY`$ one has
$$_X\phi (a,y)𝑑\mu (a)=\underset{t0}{lim}\frac{\psi (g+y)\phi (y)}{t}.$$
$`\mathrm{}`$
## 2 Variances of the ergodic sums
Corollary :If $`^2(\phi )`$ is not zero the variance of $`_{k=0}^{n1}\phi T^k)`$ is equivalent to
$$\frac{8}{3}\frac{^2(\phi )}{\sqrt{2\pi }\sigma (f)}n^{3/2}.$$
Proof : This is a direct computation using the theorem. $`\mathrm{}`$
Remark : An analogous construction is possible in continuous time. Let $`(Z,\psi _t,\rho )`$ be another Anosov flow. Define
$`\chi _t:Z\times Y`$ $``$ $`Z\times Y`$
$`(z,y)`$ $``$ $`(\chi _tz,\varphi _{_0^tf(\psi _sz)𝑑s}y).`$
One obtain a quasi-hyperbolic flow with decay of correlations in $`t^{1/2}`$. To show this one has to use the results of Waddington .
## 3 Convergence in law
Once we know the behaviour of the variances, we can ask the question of the convergence in law. We treat this problem on an example. Many technical difficulties arise when we try to generalise our result.
Our example is the following one. Let us consider
\- the group $`G=PSL(2,)`$,
\- $`\mathrm{\Gamma }`$ a cocompact lattice in $`G`$,
\- the one-parameter group $`g_t:=\left(\begin{array}{cc}e^{t/2}& 0\\ 0& e^{t/2}\end{array}\right)`$,
\- the two dimensional torus $`2,𝕋^2`$,
\- a centered function $`f`$ from $`𝕋^2`$ to $``$,
\- the matrix $`A:=\left(\begin{array}{cc}2& 1\\ 1& 1\end{array}\right)`$.
In 1988 Rudolph has introduced the application
$`T:𝕋^2\times G/\mathrm{\Gamma }`$ $``$ $`𝕋^2\times G/\mathrm{\Gamma }`$
$`(x,y)`$ $``$ $`(Ax,g_{f(x)}y).`$
An example of concrete function for which the system has good properties is the function :
$`f:𝕋^2`$ $``$ $``$
$`(x_1,x_2)`$ $``$ $`sin(2\pi x_1).`$
The transformation $`T`$ has two non zero Lyapounoff exponents, $`\lambda =\frac{3+\sqrt{5}}{2}`$ and $`\frac{1}{\lambda }`$, associated to two directions one expanded the other contracted uniformly by $`T`$ and a zero exponent in the tangent spaces to the submanifolds $`\{x\}\times G/\mathrm{\Gamma }`$. Let $`\mu `$ and $`\nu `$ be the Haar measure and the measure coming from the Haar measure respectively.
The system $`(𝕋^2\times G/\mathrm{\Gamma },T,\mu \nu )`$ is a regular version of the $`T,T^1`$-transformation. Both are examples of $`K`$-systems that do not have the Benoulli property (, ) .
In $`^2`$ let $`\stackrel{~}{x_0}`$ (resp.$`\stackrel{~}{x_1}`$ denote the intersection point of the contracted right line passing through the point $`(1,0)`$ (resp. $`(1,1)`$) and the expanded right line passing through the point $`(1,1)`$. These are a homoclinic points : in the torus $`T^k\stackrel{~}{x_0}`$ tend and $`T^k\stackrel{~}{x_1}`$ to zero when $`k`$ tends toward $`+\mathrm{}`$ and toward $`\mathrm{}`$.
Theorem : Let $`f`$ be a hölder continuous function such that
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}(f(T^k\stackrel{~}{x_0})f(T^k\stackrel{~}{x_1}))0.$$
Let $`\phi :𝕋^2\times G/\mathrm{\Gamma }`$ be a hölder continuous function such that $`^2(\phi )`$ is positive. There exist three brownian motion (non necessarily reduced) $`W,W_+,W_{}`$ such that, if $`L_t(x)`$ denotes the local time of $`W`$ in $`x`$, one has :
$$\frac{1}{n^{3/4}}\underset{k=0}{\overset{n1}{}}\phi T^k\stackrel{}{}_0^+\mathrm{}L_1(x)𝑑W_+(x)+_0^+\mathrm{}L_1(x)𝑑W_{}(x)$$
Remark : One easily verify that if $`f`$ is a coboundary the condition on $`f`$ isn’t satisfied. For regular function and for the automorphism $`A`$, to be aperiodic simply means not to be a coboundary. The condition of the theorem is a priori stronger than aperiodicity. Is it really stronger ? Another question : for which systems and which functions the notions of periodicity and coboundaricity do coincide ?
One always has $`\mathrm{\Sigma }^2(\phi _{𝕋^2}\phi (a,.)d\mu (a))=0`$. The study of the convergence in law of $`\phi T^k`$ is thus reduced to the one of $`(_{𝕋^2}\phi (a,.)d\mu (a))T^k`$. So one can restrict ourselves to the case where $`\phi `$ depends uniquely on the second coordinate. That what we’re doing from now on.
We’ll use the following result of Kesten and Spitzer.
Theorem : (Kesten, Spitzer) Let $`X_i`$ an i.i.d sequence of centered random variables with values in $``$. Let $`(\xi _i)`$ an i.i.d sequence of centered random variables independent of the $`X_i`$ variables. There exist three brownian motion (non necessarily reduced) $`W,W_+,W_{}`$ such that, if $`L_t(x)`$ denotes the local time of $`W`$ in $`x`$, one has :
$$\frac{1}{n^{3/4}}\underset{k=0}{\overset{n1}{}}\xi _{X_i+\mathrm{}+X_k}\stackrel{}{}_0^+\mathrm{}L_1(x)𝑑W_+(x)+_0^+\mathrm{}L_1(x)𝑑W_{}(x)$$
The theorem is a direct consequence of the two following propositions and of the result of Kesten and Spitzer.
Define
$$N(n,p)=\mathrm{}\{k\{0,\mathrm{},n1\}|/S_kf[p,p+1[\}.$$
Proposition : Under the hypotheses of the theorem, one has
$$\frac{1}{n^{3/4}}\underset{p}{}[N(n,p)_p^{p+1}\phi (g_t)dt\underset{k:S_kf[p,p+1[}{}\phi (g_{S_kf})]^{}0$$
Let $`\mathrm{\Phi }`$ denote the function $`\mathrm{\Phi }(y)=_0^1\phi (g_ty)𝑑t`$. One has $`_p^{p+1}\phi (g_ty)𝑑t=\mathrm{\Phi }g_p(y)`$.
Proposition : Under the hypotheses of the theorem, there exists two independent sequences of i.i.d centered random variables with values in $``$ , $`(X_i),(\xi _i)`$ such that :
$$𝔼(exp(it\frac{1}{n^{3/4}}\underset{p}{}N(n,p)\mathrm{\Phi }g_p))𝔼(exp(it\frac{1}{n^{3/4}}\underset{k=0}{\overset{n1}{}}\xi _{X_1+\mathrm{}+X_2}))\underset{n+\mathrm{}}{}0.$$
To prove these two propositions we use the three following lemmas.
Lemma 1 : (Moderated deviation)There exist $`C>0,c>0`$ such that, for very $`n`$, every $`\beta ]0,\frac{1}{2}[`$, one has
$$(S_nf>n^{1\beta })Ce^{cn^{12\beta }}$$
(Here the probability $``$ is the Lebesgue measure on $`𝕋^2`$).
Lemma 2 : (Multiple mixing for the geodesic flow) Let $`m`$ and $`m^{}`$ be two integers, $`(\mathrm{\Phi }_i)_{i=i}^{m+m^{}}`$ hölder-continuous functions defined on $`G/\mathrm{\Gamma },t_1,\mathrm{}t_m0s_1\mathrm{}s_m^{},T>0`$ real numbers. There exists $`C>0`$, and $`\delta ]0,1[`$,such that
Cov $`({\displaystyle \underset{i=1}{\overset{m}{}}}\mathrm{\Phi }_ig_{t_i},{\displaystyle \underset{j=1}{\overset{m^{}}{}}}\mathrm{\Phi }_jg_{s_i+T})C({\displaystyle \underset{i=1}{\overset{m+m^{}}{}}}\mathrm{\Phi }_i_{\mathrm{}}+{\displaystyle \underset{j}{}}[\mathrm{\Phi }_j]{\displaystyle \underset{ij}{}}\varphi _j_{\mathrm{}})\delta ^T.`$
Lemma 3 : Let
\- $`I`$ be an interval of length 1,
\- $`\epsilon >0`$ a real number,
\- $`J,K`$ be two subintervals of $`I`$ of length $`\frac{1}{[n^\epsilon ]}`$.
Let $`N(n,I)`$ denote the quantity :
$$N(n,I)=\mathrm{}\{k\{0,n1\}/S_kfI\}.$$
We define in the same way $`N(n,K)`$ and $`N(n,J)`$.
Under the hypotheses of the theorem there exist $`\xi >0,C>0`$ such that :
$`𝔼(N(n,I)^2)[n^\epsilon ]𝔼(N(n,I)N(n,J))|Cn^{1\xi }`$
$`𝔼(N(n,J)N(n,K))[n^\epsilon ]^2𝔼(N(n,I))|Cn^{1\xi }`$
$`𝔼(N(n,I)^3)|Cn^{3/2}`$
For the first and second lemmas see et . The third one is deduced from a result of speed of convergence in the local limit theorem for the sums $`S_nf`$.
|
warning/0506/hep-th0506152.html
|
ar5iv
|
text
|
# High-energy effective theory for matter on close Randall Sundrum branes
## I Introduction
Advances in string/M-theory have recently motivated the study of new cosmological scenarios for which our Universe would be embedded in compactified extra dimensions where one extra dimension could be very large relative to the Planck scale. Although the notion of extra dimensions is not new Kaluza:1921tu ; Klein:1926tv ; Antoniadis:1990ew , braneworld scenarios offer a new approach for realistic cosmological models. In some of these models, spacetime is effectively five-dimensional and gauge and matter fields are confined to three-branes while gravity and bulk fields propagate in the whole spacetime Gibbons:1986wg ; Brax:2004xh ; Davis:2005au ; Langlois:2002bb . Playing the role of a toy model, the Randall Sundrum (RS) scenario is of special interest Randall:1999ee . In the RS model, the extra dimension is compactified on an $`S_1/_2`$ orbifold, with two three-branes (or boundary branes) at the fixed points of the $`_2`$ symmetry. In the model no bulk fields are present and only gravity propagates in the bulk which is filled with a negative cosmological constant. In the low-energy limit, an effective four-dimensional theory can be derived on the branes Mendes:2000wu ; Khoury:2002jq ; Kanno:2002ia ; Shiromizu:2002qr . However, beyond this limit, braneworlds models differ remarkably from standard four-dimensional models and have some distinguishing elements which could either generate cosmological signatures or provide alternative scenarios to standard cosmology. This has been pointed out in many publications Langlois:2003yy ; Maartens:1999hf ; Langlois:2000ns ; Copeland:2000hn ; Liddle:2001yq ; Sami:2003my ; Maartens:2003tw ; Koyama:2004ap ; deRham:2004yt ; Calcagni:2003sg ; Calcagni:2004bh ; Papantonopoulos:2004bm ; Kunze:2003vp ; Liddle:2003gw ; Shiromizu:2002ve and in particular in deRham:2005xv where the characteristic features of the model are pointed out in the limit that two such three-branes are close to each other. In particular the effective four-dimensional theory was derived in the limit when the distance between the branes is much smaller than the length scale characteristic for the five-dimensional Anti-de Sitter (AdS) bulk. The effective four-dimensional Einstein equations are affected by the braneworld nature of the model and new terms in the Einstein equations contain arbitrary powers of the first derivative of the brane distance. In deRham:2005xv the main results have been obtained for the simple case where no matter is present on the branes.
The main purpose of this work is to extend this analysis and to derive an effective theory in the presence of matter on the branes. At high energy, matter couples to gravity in a different way to what is usually expected in a standard four-dimensional scenario. In particular gravity couples quadratically to the stress-energy tensor of matter fields on the brane as well as to the electric part of the five-dimensional Weyl tensor, which encodes information about the bulk geometry Shiromizu:1999wj ; Binetruy:1999hy ; Flanagan:1999cu ; Mukohyama:1999qx . Consequently we expect the covariant theory in presence of matter to be genuinely different than normal four-dimensional gravity and to bring some interesting insights on the way matter might have coupled to gravity at the beginning of a hypothetical braneworld Universe.
Our analysis relies strongly on the assumption that the brane separation is small, so that the results would only be valid just after or just before a collision. However, it is precisely this regime that is of great importance if one is to interpret the Big Bang as a brane collision Khoury:2001wf ; Khoury:2001bz ; Webster:2004sp ; Gibbons:2005rt ; Jones:2002cv or as a collision of bubbles Gen:2001bx . In particular, we may point out Blanco-Pillado:2003hq where it is shown that bubbles collision could lead to a Big Crunch. The authors show that, close to the collision, the bubbles could be treated as branes. Their collision would lead to a situation where the branes are sticking together, creating a spatially-flat expanding Universe, where inflation could take place. In that model, the collision will be well defined and not lead to any five-dimensional singularities.
In order to study the presence of matter on the branes in a model analogue to RS, we first derive, from the five-dimensional theory, the exact Friedmann equations on the branes for the background. This is done in section II, in the limit where the branes are close together, i.e. either just before or just after a brane collision. We then give in section III an overview of the effective four-dimensional theory in the limit of small brane separation as presented in deRham:2005xv and show how the theory can be formally extended in order to accommodate the presence of matter on the boundary branes. We then check that this theory gives a result that agrees perfectly with the five-dimensional solution for the background. Having checked the consistency of this effective theory for the background, we use it in order to study the effect of matter perturbations about an empty background (i.e. a ‘stiff source’ approximation) in section V. For this we consider the branes to be empty for the background and introduce matter on the positive-tension brane only at the perturbed level. We then study with more detail the propagation of tensor and scalar perturbations. Although the five-dimensional nature of the theory does not affect the way tensor and scalar perturbations propagate in a given background, the coupling to matter is indeed affected. In particular, we show that the effective four-dimensional Newtonian constant depends both on the distance between the branes and their rate of separation. We then extend the analysis in order to have a better insight of what might happen when the small brane separation condition is relaxed. The implications of our results are discussed in section VI. Finally, in appendix B, we present the technical details for the study of scalar perturbations within this close-brane effective theory.
## II Five-dimensional background behaviour
We consider a Randall-Sundrum two brane model allowing the presence of general stress-energies on each brane. Specifically, we assume an action of the form
$`S`$ $`=`$ $`{\displaystyle \mathrm{d}^5x\sqrt{g}\left(\frac{1}{2\kappa ^2}R\mathrm{\Lambda }\right)}`$
$`+{\displaystyle \underset{i=\pm }{}}{\displaystyle _^i}\mathrm{d}^4x\sqrt{g^i}\left(\lambda _i+_i\right),`$
where the two four-dimensional integrals run over the positive- and negative-tension branes $`^\pm `$ respectively and $`g_{\mu \nu }^\pm `$ are the induced metrics. We assume a $`_2`$ symmetry across the branes.
The five-dimensional bulk is filled with a negative cosmological constant $`\mathrm{\Lambda }=6/\kappa ^2L^2`$, where $`L`$ is the associated AdS radius and $`\kappa ^2`$ the five-dimensional Newtonian constant. The tensions $`\lambda _\pm `$ are, without loss of generality, assumed to take their standard fine-tuned values
$$\lambda _\pm =\pm \frac{6}{\kappa ^2L},$$
(2)
with any deviations from these absorbed into the matter Lagrangians $`_\pm `$. The resulting four-dimensional stress-energy tensors on the brane can be written as
$$T_{\mu \nu }^\pm =\frac{6}{\kappa ^2L}g_{\mu \nu }^\pm +\tau _{\mu \nu }^{(\pm )}.$$
In this paper we use the index conventions that Greek indices run from 0 to 3 and Latin from 1 to 3, referring to the Friedmann-Robertson-Walker (FRW) coordinate systems defined below in (4).
The point of this section is to extract as much information as possible about the dynamics of the system in the case of cosmological symmetry in order to obtain a result against which the effective theory can be checked. Therefore, we both assume the bulk and the brane stress-energies $`\tau _{\mu \nu }^\pm `$ to have the required symmetries. Generalising the analysis of deRham:2005xv , we work again in the stationary Birkhoff frame:
$`\mathrm{d}s^2`$ $`=`$ $`\mathrm{d}Y^2n^2(Y)\mathrm{d}T^2+a^2(Y)\mathrm{d}𝐱^2`$
$`\text{with }a^2(Y)`$ $`=`$ $`e^{2Y/L}+{\displaystyle \frac{𝒞}{4}}e^{2Y/L}`$ (3)
$`n^2(Y)`$ $`=`$ $`L^2a^{}(Y)^2=a^2{\displaystyle \frac{𝒞}{a^2}},`$
with flat spatial geometry for simplicity. The trajectories of the branes are $`Y=Y_\pm (T)`$ giving the induced line elements
$`\mathrm{d}s_\pm ^2`$ $`=`$ $`(n_\pm ^2\dot{Y}_\pm ^2)\mathrm{d}T^2+a_\pm ^2\mathrm{d}𝐱^2`$ (4)
$``$ $`\mathrm{d}t_\pm ^2+a_\pm ^2(t_\pm )\mathrm{d}𝐱^2,`$
where $`a_\pm (T)=a\left(Y=Y_\pm (T)\right)`$ and similarly for $`n_\pm (T)`$. The velocities of the branes are completely prescribed by the Israël junction conditions Israel:1966rt :
$`{\displaystyle \frac{\mathrm{d}Y_\pm }{\mathrm{d}T}}^2`$ $`=`$ $`n_\pm ^2\left(1{\displaystyle \frac{n_\pm ^2}{a_\pm ^2}}F_\pm \left(\rho _\pm \right)\right)`$ (5)
$`F_\pm \left(\rho _\pm \right)`$ $`=`$ $`\left(1\pm {\displaystyle \frac{\kappa ^2L}{6}}\rho _\pm \right)^2,`$ (6)
where $`\rho _\pm `$ are the brane energy densities $`\tau _0^{(\pm )0}`$. The Hubble parameter on each brane then follows as
$`H_\pm ^2`$ $`=`$ $`{\displaystyle \frac{1}{L^2}}\left({\displaystyle \frac{1}{F_\pm \left(\rho _\pm \right)}}{\displaystyle \frac{n_\pm ^2}{a_\pm ^2}}\right)`$ (7)
$`=`$ $`{\displaystyle \frac{𝒞}{L^2a_\pm ^4}}\pm {\displaystyle \frac{\kappa ^2}{3L}}\rho _\pm +{\displaystyle \frac{\kappa ^4}{36}}\rho _\pm ^2.`$ (8)
As in deRham:2005xv we now consider the limit of small brane separation by replacing $`n_\pm `$ and $`a_\pm `$ with their values $`n_0`$ and $`a_0`$ at the collision (equivalent to taking the leading order in $`d/L`$ where $`d`$ is related to the radion, as defined below).
To this level of approximation the brane position are then given by
$`Y_\pm (T)`$ $``$ $`Y_0v_\pm \left(TT_0\right)`$ (9)
$`v_\pm `$ $`=`$ $`n_0\sqrt{1{\displaystyle \frac{n_0^2}{a_0^2}}F_\pm \left(\rho _\pm (T=T_0)\right)},`$ (10)
where here and subsequently we take $``$ to denote the leading order in $`d/L`$, and we have chosen to consider the motion of the branes immediately after a collision at $`T=T_0`$ and $`Y=Y_0`$ when the branes are moving apart. Note firstly that the branes will in general be moving with different velocities. Secondly, the limit of large energy density $`\rho _\pm \mathrm{}`$ corresponds to $`v_\pm n_0`$, i.e. the limit of null brane velocity.
The transformation
$`TT_0`$ $`=`$ $`{\displaystyle \frac{t}{n_0}}\mathrm{cosh}\alpha (y)`$
$`YY_0`$ $`=`$ $`t\mathrm{sinh}\alpha (y)`$ (11)
$`\alpha (y)`$ $`=`$ $`(y1)\mathrm{tanh}^1\left({\displaystyle \frac{v_+}{n_0}}\right)+y\mathrm{tanh}^1\left({\displaystyle \frac{v_{}}{n_0}}\right),`$
brings the brane loci (9) to the fixed positions $`y=0,1`$, with line element
$`\mathrm{d}s^2`$ $``$ $`d(t)^2\mathrm{d}y^2\mathrm{d}t^2+a(y,t)^2\mathrm{d}𝐱^2`$ (12)
$`d(t)`$ $`=`$ $`t\left(\mathrm{tanh}^1\left({\displaystyle \frac{v_+}{n_0}}\right)+\mathrm{tanh}^1\left({\displaystyle \frac{v_{}}{n_0}}\right)\right).`$
Note as a consistency check that the global coordinate $`t`$ coincides for $`y=0,1`$ with the proper times $`t_\pm `$ on the two branes (in the small $`d`$ limit) as defined in (4), e.g.
$$\frac{\mathrm{d}t}{\mathrm{d}T}|_{y=0}=\sqrt{n_0^2v_+^2}\frac{\mathrm{d}t_+}{\mathrm{d}T}.$$
A generalisation of this metric to
$$\mathrm{d}s^2=A(x,y)^2\mathrm{d}y^2+q_{\mu \nu }(x,y)\mathrm{d}x^\mu \mathrm{d}x^\nu ,$$
(13)
for branes fixed at $`y=0,1`$ is the starting point for the derivation of the effective theory in the next section. There, the proper distance between the two branes is measured along a trajectory of constant $`x^\mu `$, i.e. it is taken to be
$$d(x)=_0^1dyA(x,y)^2.$$
(14)
In particular, if we choose a specific gauge for which $`A(x,y)`$ is independent of $`y`$, the metric (13) is simply
$$\mathrm{d}s^2=d(x)^2\mathrm{d}y^2+q_{\mu \nu }(x,y)\mathrm{d}x^\mu \mathrm{d}x^\nu .$$
(15)
As discussed in deRham:2005xv , it is unclear whether such a gauge may be fixed in general, but it can be shown that the resulting effective theory is not sensitive to the $`y`$ dependence of $`A`$.
From (7) and (10) it can then be shown that the Hubble parameter at the time of collision is related to the rate of expansion of the fifth dimension with respect to proper time $`t`$ by
$`H_+(0)`$ $`=`$ $`{\displaystyle \frac{1}{L}}\mathrm{tanh}{\displaystyle \frac{\dot{d}(0)}{2}}`$
$`+{\displaystyle \frac{\kappa ^2}{6}}\left(\rho _+(0)\mathrm{coth}\dot{d}(0)+\rho _{}(0)\mathrm{cosech}\dot{d}(0)\right).`$
Note that this is in general not the same as other definitions of the radion, more common ones being the distance along integral curves of the normal (lines of constant $`x`$ are not in general geodesics in this metric, see Fig. 1) or, different again, $`Y_{}(T)Y_+(T)`$. When the effective theory is defined from a moduli space approximation, the radion often enters via a ratio of the conformal factors on the branes, but this is not meaningful apart from in the small-velocity limit. It is however of note that all these definitions are proportional in the special case of cosmological symmetry and small brane separation.
## III Close branes effective theory description
### III.1 Formalism
We work in a frame where the branes are assumed to be exactly static at $`y=0,1`$ with metric (15) in order to simplify the implementation of the Israël junction conditions, which would otherwise be difficult. From the Gauss equations, the Einstein tensor on a $`y=\text{const}`$ hypersurface is given by Binetruy:1999ut ; Shiromizu:1999wj :
$`G_\nu ^\mu (y)`$ $`=`$ $`{\displaystyle \frac{3}{L^2}}\delta _\nu ^\mu +KK_\nu ^\mu K_\alpha ^\mu K_\nu ^\alpha `$
$`{\displaystyle \frac{1}{2}}\delta _\nu ^\mu \left(K^2K_\alpha ^\beta K_\beta ^\alpha \right)E_\nu ^\mu .`$
The unknown quantity in (III.1) is the electric part of the projected Weyl tensor $`E_\nu ^\mu `$ which is traceless, enabling us to write the Ricci scalar purely in terms of the extrinsic curvature:
$`R={\displaystyle \frac{12}{L^2}}K^2+K_\alpha ^\beta K_\beta ^\alpha .`$ (18)
The Weyl tensor $`E_\nu ^\mu `$ can be expressed in terms of more recognisable quantities as
$`E_\nu ^\mu ={\displaystyle \frac{1}{d}}{\displaystyle \frac{}{y}}K_\nu ^\mu {\displaystyle \frac{1}{d}}D^\mu D_\nu dK_\alpha ^\mu K_\nu ^\alpha +{\displaystyle \frac{1}{L^2}}\delta _\nu ^\mu ,`$ (19)
where $`D_\mu `$ is the covariant derivative with respect to $`q_{\mu \nu }(y)`$, implying from (III.1) that:
$`G_\nu ^\mu ={\displaystyle \frac{1}{d}}`$ $`[D^\mu D_\nu d+{\displaystyle \frac{}{y}}K_\nu ^\mu `$
$`+{\displaystyle \frac{2d}{L^2}}\delta _\nu ^\mu +dKK_\nu ^\mu {\displaystyle \frac{d}{2}}\delta _\nu ^\mu (K^2K_\alpha ^\beta K_\beta ^\alpha )],`$
where the second line is of higher order in the small distance limit $`dL`$.
In order to find the derivative of the extrinsic curvature $`\frac{}{y}K_\nu ^\mu `$ on the brane, we consider the Taylor expansion:
$`K_\nu ^\mu (y=1)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}_y^{(n)}K_\nu ^\mu |_{y=0}.`$ (21)
We expand the $`n^{\text{th}}`$ derivative of the extrinsic curvature in powers of $`d/L`$, keeping only the leading term:
$`{\displaystyle \frac{^n}{y^n}}K_\nu ^\mu |_{y=0}K_\nu ^{\mu (n)}=𝒦_\nu ^{\mu (n)}\left(1+𝒪\left({\displaystyle \frac{d}{L}}\right)\right)`$
and, as shown in Appendix A, one can obtain the recurrence relation
$`𝒦_\nu ^{\mu (n)}=\widehat{O}𝒦_\nu ^{\mu (n2)},`$ (22)
where the operator $`\widehat{O}`$ is defined by
$`\widehat{O}Z_\nu ^\mu =\left[d^{,\mu }Z_\nu ^\alpha +d_{,\nu }Z^{\alpha \mu }d^{,\alpha }Z_\nu ^\mu \right]d_{,\alpha },`$ (23)
giving
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}𝒦_\nu ^{\mu (n)}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n+1)!}}\widehat{O}^n𝒦_\nu ^{\mu (1)}+{\displaystyle \frac{1}{(2n)!}}\widehat{O}^n𝒦_\nu ^{\mu (0)}`$ (24)
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}\sqrt{\widehat{O}}}{\sqrt{\widehat{O}}}}𝒦_\nu ^{\mu (1)}+\mathrm{cosh}\sqrt{\widehat{O}}𝒦_\nu ^{\mu (0)}.`$
We then have a formal expression for the first derivative of the extrinsic curvature in terms of the radion and stress-energy:
$`{\displaystyle \frac{}{y}}K_\nu ^\mu |_{y=0}{\displaystyle \frac{\sqrt{\widehat{O}}}{\mathrm{sinh}\sqrt{\widehat{O}}}}\left[\mathrm{cosh}\sqrt{\widehat{O}}K_\nu ^{(+)\mu }K_\nu ^{()\mu }\right],`$ (25)
where $`K_\nu ^{(\pm )\mu }=K_\nu ^\mu (y=0,1)`$. It is straightforward then to obtain the corresponding result for the negative-tension brane:
$`{\displaystyle \frac{}{y}}K_\nu ^\mu |_{y=1}`$ $`=`$ $`{\displaystyle \frac{}{y}}K_\nu ^\mu |_{y=0}+{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{K_\nu ^{\mu (n)}}{\left(n1\right)!}}`$ (26)
$``$ $`{\displaystyle \frac{}{y}}K_\nu ^\mu |_{y=0}+\left(\mathrm{cosh}\sqrt{\widehat{O}}1\right){\displaystyle \frac{}{y}}K_\nu ^\mu |_{y=0}`$
$`+\mathrm{sinh}\sqrt{\widehat{O}}K_\nu ^{(+)\mu }`$
$``$ $`{\displaystyle \frac{\sqrt{\widehat{O}}}{\mathrm{sinh}\sqrt{\widehat{O}}}}\left[\mathrm{cosh}\sqrt{\widehat{O}}K_\nu ^{()\mu }K_\nu ^{(+)\mu }\right].`$
### III.2 Einstein equations on the branes
The next step is to use the Israël junction conditons to rewrite the extrinsic curvatures of the two branes in terms of the stress-energy tensors and the tensions:
$`K_\nu ^{(+)\mu }`$ $`=`$ $`{\displaystyle \frac{1}{L}}\delta _\nu ^\mu {\displaystyle \frac{\kappa ^2}{2}}\widehat{\tau }_\nu ^{(+)\mu }`$ (27)
$`K_\nu ^{()\mu }`$ $`=`$ $`{\displaystyle \frac{1}{L}}\delta _\nu ^\mu +{\displaystyle \frac{\kappa ^2}{2}}\widehat{\tau }_\nu ^{()\mu }`$ (28)
$`\widehat{\tau }_\nu ^{(\pm )\mu }`$ $``$ $`\tau _\nu ^{(\pm )\mu }{\displaystyle \frac{1}{3}}\tau _{}^{(\pm )}\delta _\nu ^\mu .`$ (29)
This gives us both the value of the Ricci scalar on the branes (18)
$`R^{(\pm )}={\displaystyle \frac{\kappa ^2}{L}}\tau ^{(\pm )}+{\displaystyle \frac{\kappa ^4}{4}}\left(\widehat{\tau }^{(\pm )\mathrm{\hspace{0.17em}2}}\widehat{\tau }_\beta ^{(\pm )\alpha }\widehat{\tau }_\alpha ^{(\pm )\beta }\right)`$ (30)
and, substituting (25) (or (26)) into (III.1), the effective Einstein equations:
$`G_\nu ^{(\pm )\mu }`$ $``$ $`{\displaystyle \frac{1}{d}}D^{(\pm )\mu }D_\nu ^{(\pm )}d\pm {\displaystyle \frac{1}{Ld}}\left|d\right|\mathrm{tanh}{\displaystyle \frac{\left|d\right|}{2}}\delta _\nu ^\mu `$
$`\pm {\displaystyle \frac{1}{Ld}}{\displaystyle \frac{1}{\left|d\right|}}\left(\mathrm{tanh}{\displaystyle \frac{\left|d\right|}{2}}+\mathrm{tan}{\displaystyle \frac{\left|d\right|}{2}}\right)^\mu d_\nu d`$
$`+{\displaystyle \frac{\kappa ^2}{2d}}𝒜_\nu ^{(\pm )\mu }`$
$`+\left[\pm {\displaystyle \frac{2\kappa ^2}{L}}\left(\tau _\nu ^{(\pm )\mu }{\displaystyle \frac{\tau ^{(\pm )}}{6}}\delta _\nu ^\mu \right)+\kappa ^4\mathrm{\Pi }_\nu ^{(\pm )\mu }\right],`$
where
$$𝒜_\nu ^{(\pm )\mu }=\sqrt{\widehat{O}}\left(\mathrm{coth}\sqrt{\widehat{O}}\widehat{\tau }_\nu ^{(\pm )\mu }+\mathrm{cosech}\sqrt{\widehat{O}}\widehat{\tau }_\nu ^{()\mu }\right),$$
(32)
and
$$\mathrm{\Pi }_\nu ^{(\pm )\mu }=\frac{1}{8}\tau _\beta ^{(\pm )\alpha }\tau _\alpha ^{(\pm )\beta }\delta _\nu ^\mu \frac{1}{12}\tau ^{(\pm )}\tau _\nu ^{(\pm )\mu }\frac{1}{72}\tau ^{(\pm )\mathrm{\hspace{0.17em}2}}\delta _\nu ^\mu .$$
(33)
From the tracelessness of $`E_{\mu \nu }^{(\pm )}`$ we obtain two equivalent equations of motion for the radion,
$`\mathrm{}^{(\pm )}d`$ $``$ $`{\displaystyle \frac{\left|d\right|}{L}}\left(3\mathrm{tanh}{\displaystyle \frac{\left|d\right|}{2}}\mathrm{tan}{\displaystyle \frac{\left|d\right|}{2}}\right){\displaystyle \frac{\kappa ^2}{2}}𝒜^{(\pm )}`$
$`+\left[\pm {\displaystyle \frac{\kappa ^2d}{3L}}\tau ^{(\pm )}{\displaystyle \frac{\kappa ^4d}{4}}\widehat{\tau }_\nu ^{(\pm )\mu }\widehat{\tau }_\mu ^{(\pm )\nu }\right].`$
The terms in square brackets in (III.2) and (III.2) will turn out to be of higher order as $`d0`$ and so should not strictly be included. However, for exact cosmological symmetry, they are the only higher order terms and we shall keep them for the time being. Later on they shall be discarded.
The conservation of the stress-energy tensor on both branes follows directly from the Codacci equation Binetruy:1999ut ; Shiromizu:1999wj :
$`D_\mu K_\nu ^\mu D_\nu K=0,`$ (35)
which, evaluated on the branes implies
$`D_\mu ^{(+)}\tau _\nu ^{(+)\mu }=D_\mu ^{()}\tau _\nu ^{()\mu }=0.`$ (36)
### III.3 Low-energy limit
As a first consistency check of this close-brane theory, we consider its low-energy limit and compare it with the effective four-dimensional low-energy theory Mendes:2000wu ; Khoury:2002jq ; Kanno:2002ia ; Shiromizu:2002qr for small brane separation. In that common limit, both theories should agree.
In the low-energy limit, the magnitude of the stress-energy tensor on the brane is small compared to the brane tension. Any quadratic term $`\kappa ^4\tau ^{(\pm )\mathrm{\hspace{0.17em}2}}`$ is negligible compared to $`\frac{\kappa ^2}{L}\tau ^{(\pm )}`$, so that $`\mathrm{\Pi }_\nu ^{(\pm )\mu }`$ may be dropped in (III.2) and, from (30), the Ricci tensor on the brane is:
$`R^{(\pm )}={\displaystyle \frac{\kappa ^2}{L}}\tau ^{(\pm )}.`$ (37)
Furthermore in the low-energy limit, the branes are moving slowly, $`\left|d\right|1`$, to linear order in $`\left|d\right|`$, we have:
$`𝒜_\nu ^{(\pm )\mu }=\widehat{\tau }_\nu ^{(+)\mu }+\widehat{\tau }_\nu ^{()\mu }.`$ (38)
The effective Einstein equation on the brane at low energy is therefore
$`G_\nu ^{(\pm )\mu }`$ $``$ $`{\displaystyle \frac{1}{d}}D^{(\pm )\mu }D_\nu ^{(\pm )}d+{\displaystyle \frac{\kappa ^2}{2d}}\left(\widehat{\tau }_\nu ^{(+)\mu }+\widehat{\tau }_\nu ^{()\mu }\right)`$
$`\pm {\displaystyle \frac{1}{Ld}}\left(^\mu d_\nu d{\displaystyle \frac{1}{2}}\left(d\right)^2\delta _\nu ^\mu \right),`$
with the equation of motion for $`d`$:
$`\mathrm{}^{(\pm )}d=\pm {\displaystyle \frac{1}{L}}\left(d\right)^2+{\displaystyle \frac{\kappa ^2}{6}}\left(\tau ^{(+)}+\tau ^{()}\right).`$ (40)
We can therefore write (III.3) in the more common form:
$`G_\nu ^{(\pm )\mu }`$ $``$ $`{\displaystyle \frac{\kappa ^2}{2d}}\left(\tau _\nu ^{(+)\mu }+\tau _\nu ^{()\mu }\right)`$
$`+{\displaystyle \frac{1}{d}}\left(D^{(\pm )\mu }D_\nu ^{(\pm )}d\mathrm{}^{(\pm )}d\right)`$
$`\pm {\displaystyle \frac{1}{Ld}}\left(^\mu d_\nu d+{\displaystyle \frac{1}{2}}\left(d\right)^2\delta _\nu ^\mu \right),`$
which is precisely the result we get from the effective low-energy theory in the close brane limit Mendes:2000wu ; Khoury:2002jq ; Kanno:2002ia ; Shiromizu:2002qr ; deRham:2004yt .
## IV Cosmological Symmetry
In the most general case, the coupling of the radion to matter on the branes given by (32) is intractable. However, we are concerned here with the case of cosmological symmetry as a check on the validity of the theory. We take (12) as our metric and notice that
$`\widehat{O}\left(\begin{array}{cc}A& 0\\ 0& B\delta _j^i\end{array}\right)=\dot{d}^2\left(\begin{array}{cc}\hfill A& 0\\ \hfill 0& B\delta _j^i\end{array}\right).`$ (46)
We can then obtain the coupling tensors (32) in closed form:
$`𝒜_0^{(\pm )0}`$ $`=`$ $`\dot{d}\mathrm{cot}\dot{d}\left({\displaystyle \frac{2}{3}}\rho _\pm +p_\pm \right)`$
$`\dot{d}\mathrm{cosec}\dot{d}\left({\displaystyle \frac{2}{3}}\rho _{}+p_{}\right)`$
$`𝒜_j^{(\pm )i}`$ $`=`$ $`{\displaystyle \frac{\dot{d}}{3}}\left(\rho _\pm \mathrm{coth}\dot{d}+\rho _{}\mathrm{cosech}\dot{d}\right)\delta _j^i.`$ (48)
The resulting equations of motion follow from (III.2),(III.2) and (36):
$`H_\pm ^2`$ $``$ $`{\displaystyle \frac{1}{d}}\left[\dot{d}H_\pm \pm {\displaystyle \frac{\dot{d}}{L}}\mathrm{tanh}{\displaystyle \frac{\dot{d}}{2}}+{\displaystyle \frac{\kappa ^2}{6}}𝒜_i^{(\pm )i}\right]`$
$`\pm {\displaystyle \frac{2\kappa ^2}{3L}}\rho _\pm +{\displaystyle \frac{1}{18}}\kappa ^4\rho _\pm ^2`$
$`\dot{H}_\pm +2H_\pm ^2`$ $``$ $`\pm {\displaystyle \frac{\kappa ^2}{6L}}\left(\rho _\pm 3p_\pm \right){\displaystyle \frac{\kappa ^4}{36}}\rho _\pm \left(\rho _\pm +3p_\pm \right)`$ (50)
$`\ddot{d}+3H_+\dot{d}`$ $``$ $`{\displaystyle \frac{\dot{d}}{L}}\left(3\mathrm{tanh}{\displaystyle \frac{\dot{d}}{2}}\mathrm{tan}{\displaystyle \frac{\dot{d}}{2}}\right)`$
$`+{\displaystyle \frac{\kappa ^2d}{L}}\left({\displaystyle \frac{\rho _+}{3}}p_+\right)+{\displaystyle \frac{\kappa ^2}{2}}\left(𝒜_0^{(+)0}+𝒜_i^{(+)i}\right)`$
$`+{\displaystyle \frac{\kappa ^4d}{4}}\left({\displaystyle \frac{7}{9}}\rho _+^2+p_+^2+{\displaystyle \frac{4}{3}}\rho _+p_+\right)`$
$`\dot{\rho }_\pm `$ $`=`$ $`3H_\pm \left(\rho _\pm +p_\pm \right).`$ (52)
Equations (52) and (50) together imply
$$H_\pm ^2\frac{\stackrel{~}{𝒞}}{La_\pm ^4}\pm \frac{\kappa ^2}{3L}\rho _\pm +\frac{\kappa ^4}{36}\rho _\pm ^2,$$
(53)
where $`\stackrel{~}{𝒞}`$ is an integration constant which, at this order, can be identified with the bulk parameter $`𝒞`$ via (7). The system is now manifestly finite as $`d0`$. Note that, apart from the presence of quadratic terms, (50) is the same result as that obtained from the moduli space approximation and is, in fact, exact ($`d`$ decouples as a consequence of the simplicity of the Weyl tensor for an AdS-Schwarzschild bulk). However, the additional information from (IV) gives additional information not obtainable from a low-energy effective theory. Since $`H_+`$ takes a finite value at the collision, the coefficient of $`d^1`$ in (IV) must vanish at $`d=0`$; this implies then that
$`H_+(0)`$ $`=`$ $`{\displaystyle \frac{1}{L}}\mathrm{tanh}{\displaystyle \frac{\dot{d}(0)}{2}}`$
$`+{\displaystyle \frac{\kappa ^2}{6}}\left(\rho _+(0)\mathrm{coth}\dot{d}(0)+\rho _{}(0)\mathrm{cosech}\dot{d}(0)\right),`$
in agreement with the exact result (II).
## V Effective theory for perturbations
More interesting is the study of cosmological perturbations, for which a relatively straightforward solution of the above system is also available. We shall give a few examples here and point out some interesting features. Throughout we work only with the positive-tension brane, assuming the negative-tension brane to be empty, and drop the $`\pm `$ signs. Also, we shall assume for simplicity that matter on the positive-tension brane is only introduced at the perturbed level, i.e. the background solution $`a(t),d(t)`$ is that obtained from (50) and (IV) in the absence of matter. We therefore have
$$\begin{array}{ccccc}\tau _\nu ^\mu (𝐱,t)& =& 0\hfill & +& \delta \tau _\nu ^\mu \hfill \\ d(𝐱,t)& =& d(t)\hfill & +& \delta d(𝐱,t)\hfill \\ g_{\mu \nu }(𝐱,t)& =& \overline{g}_{\mu \nu }\hfill & +& \delta g_{\mu \nu },\hfill \end{array}$$
where $`\overline{g}_{\mu \nu }`$ is the usual flat FRW metric with scale factor $`a(t)`$. However, in the following we will set $`\delta d=0`$, either because we are considering tensor perturbations only or because we choose to work in such a gauge. Hence we shall assume that $`d`$ takes its background value.
### V.1 Tensor Perturbations
As the simplest starting point we consider perturbations using the above formalism and we choose to work in conformal time. We take the metric to be
$$\mathrm{d}s^2=a(\eta )^2\left(\mathrm{d}\eta ^2+\left(\delta _{ij}+h_{ij}\right)\mathrm{d}x^i\mathrm{d}x^j\right),$$
(55)
with the usual transverse traceless conditions
$$h_i^i=h_{j,i}^i=0$$
on the perturbation, spatial indices being raised by $`\delta ^{ij}`$. The resulting Ricci tensor perturbation is then
$`\delta R_{00}`$ $`=`$ $`\delta R_{0i}=0,`$
$`\delta R_{ij}`$ $`=`$ $`\left(^{}+2^2\right)h_{ij}+h_{ij}^{}{\displaystyle \frac{1}{2}}\overline{\mathrm{}}h_{ij},`$ (56)
where $`=a^{}/a=aH`$, primes denote differentiation with respect to conformal time $`\eta `$ and $`\overline{\mathrm{}}=_\eta ^2+_i^i`$ is the Minkowski space wave operator. We assume that these gravity waves are sourced by tensor matter at the perturbative level, i.e.
$$\delta \tau _\nu ^\mu =\delta \widehat{\tau }_\nu ^\mu =\left(\begin{array}{cc}0& 0\\ 0& \tau _j^i\end{array}\right),\tau _i^i=0.$$
The perturbed Klein Gordon equation for tensor matter just reduces to $`\delta \left(\mathrm{}d\right)0`$, so it is consistent to set the scalar perturbation $`\delta d`$ to zero, i.e. to study purely tensor fluctuations. In this case, the equation of motion for the perturbations follows from (III.2):
$`\delta R_{ij}`$ $``$ $`{\displaystyle \frac{1}{d}}\delta \left(D_iD_jd\right)+{\displaystyle \frac{\dot{d}}{Ld}}\mathrm{tanh}{\displaystyle \frac{\dot{d}}{2}}a^2h_{ij}`$
$`+{\displaystyle \frac{\kappa ^2}{2d}}\sqrt{\widehat{O}}\mathrm{coth}\sqrt{\widehat{O}}\delta \widehat{\tau }_{ij},`$
where we have now dropped the sub-dominant matter terms. It is straightforward to obtain
$`\delta \left(D_iD_jd\right)`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}h_{ij}^{}+h_{ij}\right)d^{}`$
$`\sqrt{\widehat{O}}\mathrm{coth}\sqrt{\widehat{O}}\delta \widehat{\tau }_{ij}`$ $`=`$ $`\dot{d}\mathrm{coth}\dot{d}\tau _{ij}.`$
Equations (IV) and (50) then imply the relation for the background Hubble parameter
$$^{}+^2=0,^2=\frac{d^{}}{d}+a^2\frac{\dot{d}}{L}\mathrm{tanh}\frac{\dot{d}}{L}.$$
Putting this all together we obtain, to leading order in $`d`$:
$$\widehat{}h_{ij}\left[\overline{\mathrm{}}\frac{d^{}}{d}\frac{}{\eta }\right]h_{ij}=\kappa ^2\frac{\dot{d}}{d}\mathrm{coth}\dot{d}\delta \tau _{ij}.$$
(58)
The same calculation repeated subject to the low-energy approximation, not assuming small $`d`$, is straightforward. Since the matter is traceless, the standard equations at low energy Mendes:2000wu ; Khoury:2002jq ; Kanno:2002ia ; Shiromizu:2002qr ; deRham:2004yt give
$$R_{ij}\frac{1\psi }{L^2\psi }\left(2LD_i^{(+)}_jdg_{ij}d^2+2_id_jd\right)+\frac{\kappa ^2}{L\psi }\tau _{ij}^{(+)},$$
where $`\psi =1\mathrm{exp}\left(2d/L\right)`$. Perturbing this gives
$$\overline{\mathrm{}}h_{ij}2\left[+\frac{\left(1\psi \right)}{\psi }\frac{d^{}}{L}\right]h_{ij}^{}\frac{2\kappa ^2}{L\psi }\delta \tau _{ij},$$
(59)
using the equations
$$^{}+^2=0,^2+2\frac{\left(1\psi \right)}{L\psi }d^{}\frac{d^2}{L^2}\frac{\left(1\psi \right)}{\psi },$$
for the background. The small-$`d`$ limit of (59), where $`\psi 2d/L`$, is then
$$\overline{\mathrm{}}h_{ij}\frac{d^{}}{d}h_{ij}^{}\frac{\kappa ^2}{d}\delta \tau _{ij}.$$
The operator $`\widehat{}`$ defined in (58) is therefore the same as one would find in the low-energy theory. The difference lies in the source term; in the high-energy theory, the effective four-dimensional Newton constant on the positive-tension brane is related to the five-dimensional one by
$`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}={\displaystyle \frac{\dot{d}}{d}}\mathrm{coth}\dot{d}\kappa ^2,`$ (60)
whereas the low-energy result has
$`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{\kappa ^2}{d_0}}.`$ (61)
As is the case in the low-energy effective theory, the coupling to matter is different for the background as it is for the perturbations - for the background, the coupling can be identified from (8) or (53) as $`\kappa ^2/L`$, as opposed to (60). When either the branes are stabilised, and $`d`$ is not treated as a dynamical variable, $`dd_0=\text{const}`$, or the velocity is small $`\dot{d}1`$ (which is the case in the low-energy limit), it is easy to see that (60) and (61) agree. However, for arbitrary brane velocities, when the radion is not stabilised, the exact result for small $`d`$ is given by (60). As expected, the effective Newton constant picks up a dependence on $`d`$, as it does in the low-energy theory, but more unexpected, it also contains some degree of freedom: the brane separation velocity. Whilst this is not expected to be relevant today, since one would assume the radion is stabilised in the present Universe, it would be extremely important near the brane collision. As discussed in section II, $`\dot{d}`$ would be approximately constant, $`\dot{d}v`$, leading to
$`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{\mathrm{coth}v}{t}}\kappa ^2,`$ (62)
where the coefficient $`\mathrm{coth}v`$ could take any value greater than $`1`$ depending on the matter content of the branes.
### V.2 Scalar Perturbations
We now consider scalar metric perturbations on the brane sourced by a perfect fluid at the perturbative level (again, the background geometry is taken to be empty). We choose to work in a gauge where $`\delta d=0`$, i.e. to evaluate the perturbations on hypersurfaces of constant $`d`$, in which the metric perturbation can be taken as
$`\mathrm{d}s^2`$ $`=`$ $`a(\eta )^2((1+2\mathrm{\Phi })\mathrm{d}\eta ^2+4E_{,i}\mathrm{d}\eta \mathrm{d}x^i`$
$`+(1+2\mathrm{\Psi })\delta _{ij}\mathrm{d}x^i\mathrm{d}x^j).`$
The calculations are not nearly so straightforward as for tensors and have therefore been relegated to Appendix B. The result is the following equation of motion for the curvature perturbation
$$\widehat{}\mathrm{\Psi }\left[\overline{\mathrm{}}\frac{d^{}}{d}\frac{}{\eta }\right]\mathrm{\Psi }=a^2\frac{\kappa ^2}{6}\frac{\dot{d}}{d}\mathrm{coth}\dot{d}\delta \rho ,$$
(64)
giving rise to the same relation between the four-dimensional Newtonian constant and the five-dimensional as in (60). Here again we may check that, apart from the modification of the effective Newtonian constant on the brane, the perturbations propagate in the given background exactly the same way as they would if the theory were genuinely four-dimensional. This is a very important result for the propagation of scalar perturbations if they are to generate the observed large-scale structure. The five-dimensional nature of the theory does affect the background behaviour but on this background the perturbations behave exactly the same way as they would in the four-dimensional theory.
This result is of course only true in the close-brane limit, for which the theory contains no higher than second derivatives, only powers of first derivatives. When the branes are no longer very close to each other, the theory will become higher-dimensional (in particular the theory becomes non-local in the one brane limit). The presence of these higher-derivative corrections (not expressible as powers of first derivatives) is expected to modify the way perturbations propagate in a given background, mainly because extra Cauchy data would need to be specified, making the perturbations non adiabatic deRham:2004yt . However if we consider a scenario for which the large-scale structure is generated just after the brane collision, the mechanism for the production of the scalar perturbations will be very similar to the standard four-dimensional one.
### V.3 Relation between the four- and five-dimensional Newtonian constant
The relation (60) between $`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}`$ and the five-dimensional constant $`\kappa ^2`$ is formally only valid for small distance between the branes. However if we consider the analysis of deRham:2005xv , we may have some insights of what will happen if we had not stopped the expansion to leading order in $`d`$. Here, terms of the form $`d\ddot{d}`$ and more generally any term of the form $`d^nd^{(n+1)}`$ have been considered as negligible in comparison to $`\dot{d}`$ and therefore only the terms of the form $`\dot{d}^n`$ have been considered in the expansion. In a more general case, when the branes are not assumed to be very close to each other, any term of the form $`d^nd^{(n+1)}`$ should be considered and would affect the relation between the four-dimensional Newtonian constant and the five-dimensional one. For moving branes, we therefore expect the relation between $`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}`$ and $`\kappa ^2`$ to be:
$$\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}=\frac{\kappa ^2}{L}\mathrm{\Omega }(\frac{d}{L},\dot{d},d\ddot{d},\mathrm{},d^nd^{(n+1)}).$$
(65)
The relation is therefore a functional of $`d`$: $`\mathrm{\Omega }\left[d(t)\right]`$ has an infinite number of independent degree of freedom.
In the low-energy limit, or when the radion is stabilised, $`dd_0=\text{const}`$, the exact expression of $`\mathrm{\Omega }`$ is Garriga:1999yh :
$$\mathrm{\Omega }\mathrm{\Omega }\left[d(t)=d_0\right]=\frac{e^{d_0/L}}{2\mathrm{sinh}d_0/L}$$
(66)
For close branes, another limit is now known: when $`dL`$,
$$\mathrm{\Omega }\left[dL\right]=\frac{\dot{d}}{d}\mathrm{coth}\dot{d}.$$
(67)
But in a general case, $`\mathrm{\Omega }`$ (and therefore $`\kappa _{4d}^{(+)\mathrm{\hspace{0.17em}2}}`$) is expected to be a completely dynamical degree of freedom. For the present Universe the radion is supposed to be stabilised, but in early-Universe cosmology, the effective four-dimensional Newton constant could be very different from its present value. It might therefore be interesting to understand what the constraints on such time-variation of the Newtonian constant would be and how it would constrain the brane velocity Clifton:2005xr ; Barrow:1996kc , or whether such a time variation could act as a signature for the presence of extra dimensions.
## VI Conclusions
In the first part of this work, we derived the exact behaviour of FRW branes in the presence of matter. The characteristic features come from the presence of the $`\rho ^2`$ terms in the Friedmann equation and from the ‘dark energy’ Weyl term. In the limit of close separation we related the contribution of the Weyl term to the expansion of the fifth dimension. We then used this result to test the close-brane effective theory that was first derived in deRham:2005xv but now with matter introduced on the branes. For this we have shown how matter can be included using a formal sum of operators acting on the stress-energy tensors for matter fields on both branes. In the general case the action of this sum of operators on an arbitrary stress-energy tensor would not be available in closed form, although one could in principle proceed perturbatively. When a specific scenario is chosen, however, one can make considerable analytical progress. Assuming cosmological symmetry, the action of the operators on the stress-energy tensor is remarkably simple and the sum can be evaluated analytically. We then compared the result with the exact five-dimensional result in the limit of small brane separation. As expected both results agree perfectly. Furthermore we have checked that, in the low-energy limit, our close-brane effective theory agrees perfectly with the effective four-dimensional low-energy theory, giving another consistency check.
We then used this close-brane effective theory in order to understand the way matter couples to gravity at the perturbed level. In order to do so, we considered a scenario in the stiff source approximation for which the background is supposed to be unaffected by the presence of matter and considered the production of curvature and tensor perturbation sourced by the presence of matter fields on the brane. Although the five-dimensional nature of the theory does affect the background behaviour, we have shown that for a given background the perturbations propagate the same way as they would in a standard four-dimensional theory. This is only true in the limit of small brane separation and is not expected to be valid outside this regime. However, since the large-scale structure of the Universe might have been produced in a period for which the branes could have been close together (for instance just after a brane collision initiating the Big Bang), this regime is of special interest. The fact that the perturbations behave the same way, for a given background, as they would in a four-dimensional theory is a remarkable result for the production of the large scale structure which could be almost unaffected by the presence of the fifth dimension. On the other hand, the relation between the four-dimensional Newtonian constant and the five-dimensional one is however affected by the expansion of the fifth-dimension. It has been shown in the literature Randall:1999ee ; Garriga:1999yh that four-dimensional Newtonian constant was dependant on the distance between the branes, giving a possible explanation of the hierarchy problem. In this paper we show that the four-dimensional Newtonian constant also has some dependence on the brane velocity which we computed exactly in the small-distance limit, which might be able to provide an observational signature for the presence of extra dimensions. Outside the small $`d`$ regime, we expect the four-dimensional Newtonian constant to depend on the five-dimensional one not only through the brane separation velocity $`\dot{d}`$ but also on higher derivatives of the distance between the branes $`d^{(n)}`$, making the requirement for moduli stabilisation even more fundamental for any realistic cosmological setup within braneworld cosmology.
## VII Acknowledgements
The authors would like to thank Anne Davis for her supervision and comments on the manuscript and Andrew Tolley for useful discussions. SLW is supported by PPARC and CdR by DAMTP.
## Appendix A Leading order derivative of the extrinsic curvature
In this appendix, we shall derive an expression for the derivative of the extrinsic curvature on he branes. We will not be able to calculate these quantities exactly, but will be able to obtain a relatively simple expression for its leading-order contribution. We will focus on the positive-tension brane first, and our starting point shall be the Taylor series
$$K_\nu ^{()\mu }=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}K_\nu ^{\mu (n)},$$
(68)
where we are defining
$$K_\nu ^{\mu (n)}\frac{^n}{y^n}K_\nu ^\mu |_{y=0}.$$
We are interested only in the leading order contribution to $`K_\nu ^{\mu (n)}`$,
$$K_\nu ^{\mu (n)}=𝒦_\nu ^{\mu (n)}\left(1+𝒪\left(d/L\right)\right),$$
and the aim of this section will be to establish that
$$𝒦_\nu ^{\mu (n)}=\widehat{O}𝒦_\nu ^{\mu (n2)},$$
(69)
where the operator $`\widehat{O}`$ is defined by
$`\widehat{O}Z_\nu ^\mu `$ $``$ $`\left[d^{,\mu }Z_\nu ^\alpha +d_{,\nu }Z^{\alpha \mu }d^{,\alpha }Z_\nu ^\mu \right]d_{,\alpha }.`$
This implies that $`𝒦_\nu ^{\mu (n)}`$ is of the same order as $`𝒦_\nu ^{\mu (n2)}`$, and will allow us to produce a simple, albeit formal, expression for this sum, which will be the starting point for writing down the small-$`d`$ effective theory in the next section. We will proceed by induction, and throughout make the following assumptions about the order of terms:
* $`_\mu _\nu d_\nu dd^0`$
* $`\tau _\nu ^{(+)\mu }`$, $`D_\alpha ^{(+)}\tau _\nu ^{(+)\mu }`$ are at worst as divergent as the geometry
* $`E_\nu ^{(+)\mu }d^1`$.
The assumption on the order of magnitude of the matter terms is reasonable, since the matter introduced on the brane is expected to scale as the scale factor for the background and as the curvature perturbation for general perturbations. Since the curvature perturbations is in general expected to diverge logarithmically at the collision, we can hence assume that $`\tau _\nu ^{(+)\mu }`$, is, at worst, logarithmically divergent in $`d`$. This implies that the extrinsic curvature on the brane is itself at worse logarithmically divergent in $`d`$. Similarly, we know that $`E_\nu ^{(+)\mu }d^1`$ for the low-energy theory. Although we have argued that at high energy the moduli space approximation does not give the exact expression for the Weyl tensor, we have seen that (at least for the background) the behaviour is the same, differing only in corrections at higher order in the velocity. In particular $`E_\nu ^\mu `$ should go as $`d^1`$ at high energies as well (we will see later that this is indeed the case). From (19) we have
$`K_\nu ^{\mu (1)}`$ $`=`$ $`dE_\nu ^\mu \left(y=0\right)D^\mu _\nu d|_{y=0}`$ (70)
$`{\displaystyle \frac{\kappa ^2d}{L}}\tau _\nu ^{(+)\mu }{\displaystyle \frac{\kappa ^4}{4}}d\widehat{\tau }_\alpha ^{(+)\mu }\widehat{\tau }_\nu ^{()\alpha }`$
which, from the above assumptions, gives us
$$𝒦_\nu ^{\mu (1)}d^0.$$
(71)
For the second derivative of the extrinsic curvature, i.e. for $`n=2`$, we need expressions for the derivatives of the Weyl tensor and the Christoffel symbols. It is straightforward to show that
$$\mathrm{\Gamma }_{\mu \nu }^\alpha =D_\mu (dK_\nu ^\alpha )+D_\nu (dK_\mu ^\alpha )D^\alpha (dK_{\mu \nu }),$$
(72)
and the derivative of the Weyl tensor is deRham:2005xv
$`E_\nu ^\mu `$ $`=`$ $`d(2K_\nu ^\alpha E_\alpha ^\mu {\displaystyle \frac{3}{2}}KE_\nu ^\mu {\displaystyle \frac{1}{2}}K_\beta ^\alpha E_\alpha ^\beta \delta _\nu ^\mu +C_{\alpha \nu \beta }^\mu K^{\alpha \beta }`$
$`2\widehat{K}_\mu ^\alpha \widehat{K}_{\alpha \beta }\widehat{K}_\nu ^\alpha {\displaystyle \frac{7}{6}}\widehat{K}_{\alpha \beta }\widehat{K}^{\alpha \beta }\widehat{K}_{\mu \nu }{\displaystyle \frac{1}{2}}q_{\mu \nu }\widehat{K}_{\alpha \beta }\widehat{K}_\rho ^\beta \widehat{K}^{\alpha \rho })`$
$`{\displaystyle \frac{1}{2d}}D^\alpha \left[d^2D^\mu K_{\alpha \nu }+d^2D_\nu K_\alpha ^\mu 2d^2D_\alpha K_\nu ^\mu \right],`$
where $`\widehat{K}_\nu ^\mu =K_\nu ^\mu \frac{1}{4}K\delta _\nu ^\mu `$. On the brane, from the Israël matching conditions, the trace of the extrinsic curvature is $`K\tau `$, hence $`\widehat{K}_\nu ^\mu \tau `$ also. So the cubic terms in $`\widehat{K}`$ will be of higher order than the $`KE`$ terms, as will the $`CK`$ term. The leading terms will, in fact, just be the first three, giving
$$E_{\nu }^{\mu }{}_{}{}^{}(0)\frac{4d}{L}E_\nu ^\mu (0).$$
(74)
On the brane,
$`D_\alpha ^{(+)}\left(dK_\nu ^{(+)\mu }\right)`$ $`=`$ $`\left(_\alpha d\right)K_\nu ^{(+)\mu }+dD_\alpha ^{(+)}K_\nu ^{(+)\mu }`$ (75)
$``$ $`\left(_\alpha d\right)K_\nu ^{(+)\mu },`$
the second term being subdominant from the assumption that $`D_\alpha ^{(+)}K_\nu ^{(+)\mu }D_\alpha ^{(+)}\widehat{\tau }_\nu ^{(+)\mu }`$ is of higher order than $`d^1`$. The derivative of the Christoffel symbol will similarly be of the same order as the extrinsic curvature on the brane:
$$\mathrm{\Gamma }_{\mu \nu }^\alpha (0)\left(d_{,\mu }K_\nu ^{(+)\alpha }+d_{,\nu }K_\mu ^{(+)\alpha }d^{,\alpha }K_{\mu \nu }^{(+)}\right).$$
(76)
Taking the derivative of (19) gives
$`K_{\nu }^{\mu }{}_{}{}^{\prime \prime }(y)`$ $`=`$ $`dE_{\nu }^{\mu }{}_{}{}^{}+2dK^{\mu \beta }D_\beta _\nu d+q^{\mu \beta }\mathrm{\Gamma }_{\beta \nu }^\alpha _\alpha d`$ (77)
$`d_y\left(K_\alpha ^\mu K_\nu ^\alpha \right),`$
in the bulk. Evaluated on the brane using (74) and (76), the dominant term (of order $`d^0`$) is the one containing the derivative of the Christoffel symbol
$`K_\nu ^{\mu (2)}`$ $``$ $`q^{(+)\mu \beta }\mathrm{\Gamma }_{\beta \nu }^\alpha (0)_\alpha dK.`$ (78)
Since we have shown in (76) that $`\mathrm{\Gamma }_{\beta \nu }^\alpha (0)d_{,\beta }K_\nu ^{(+)\alpha }`$, on the brane, the second derivative of the extrinsic curvature is hence of the same order as the extrinsic curvature itself $`K_\nu ^{\mu (2)}K_\nu ^\mu `$.
Using (76), we have proved the result for $`n=2`$:
$`𝒦_\nu ^{\mu (2)}`$ $`=`$ $`d_{,\alpha }\left[d^{,\mu }K_\nu ^{(+)\alpha }+d_{,\nu }K^{(+)\alpha \mu }d^{,\alpha }K_{\mu \nu }^{(+)}\right]`$ (79)
$``$ $`\widehat{O}K_\nu ^{(+)\mu }.`$
The second derivative of the Christoffel symbol follows from (72):
$`\mathrm{\Gamma }_{\mu \nu }^{\alpha \prime \prime }(y)`$ $`=`$ $`D_\mu \left(dK_\nu ^\alpha \right)+D_\nu \left(dK_\mu ^\alpha \right)D^\alpha \left(dK_{\mu \nu }^{}\right)`$
$`+d\left(\mathrm{\Gamma }_{\mu \rho }^\alpha K_\nu ^\rho +\mathrm{\Gamma }_{\nu \rho }^\alpha K_\mu ^\rho 2\mathrm{\Gamma }_{\mu \nu }^\rho K_\rho ^\alpha \right)`$
$`dq^{\alpha \beta }q_{\mu \sigma }\left(\mathrm{\Gamma }_{\beta \rho }^\sigma K_\nu ^\rho \mathrm{\Gamma }_{\beta \nu }^\rho K_\rho ^\sigma \right)`$
$`d\left(q^{\alpha \beta }q_{\mu \sigma }\right)^{}\left(\mathrm{\Gamma }_{\beta \rho }^\sigma K_\nu ^\rho \mathrm{\Gamma }_{\beta \nu }^\rho K_\rho ^\sigma \right)`$
where, recall, $`K_{\nu }^{\mu }{}_{}{}^{}=d_nK_\nu ^\mu `$ is a tensor, hence the use of covariant derivatives. These $`D(dK^{})`$ terms are of order $`d^0`$, whilst the others are all of higher order, when evaluated on the brane. The leading term is
$`\mathrm{\Gamma }_{\mu \nu }^{\alpha \prime \prime }(0)`$ $``$ $`_\mu dK_\nu ^\alpha (0)+_\nu dK_\mu ^\alpha (0)^\alpha dK_{\mu \nu }^{}(0)`$ (81)
$``$ $`d^0.`$
Substituting (A) into (77) gives a complicated second-order differential equation for $`K_\nu ^\mu `$. Taking repeated $`y`$-derivatives of this equation would be impractical, but to start with all we want to do is to work to leading order. We will first identify which term is dominant, before actually evaluating it. We therefore drop all indices and numerical factors for the time being, writing $`q`$ for the metric (with indices in any position), $`K`$ for $`K_\nu ^\mu `$ and $``$ for $`_\mu `$. For example, $`q_{\mu \nu }^{}dK_{\mu \nu }=dq_{\mu \sigma }K_\nu ^\sigma `$ and so we would write $`q^{}=dqK`$. The equation for $`K^{\prime \prime }`$ can then be written symbolically as
$`K^{\prime \prime }`$ $`=`$ $`d\left(^2+\mathrm{\Gamma }+\mathrm{\Gamma }+\mathrm{\Gamma }^2\right)q^{}+dK^3+dK`$
$`+dKK^{}+qd\left(\mathrm{\Gamma }^{}+d\mathrm{\Gamma }K+dK\right)+^2dq^{},`$
and, from (78), we already know that the dominant term is $`qd\mathrm{\Gamma }^{}K`$. We know that $`K^{(m)}(0)`$ and $`\mathrm{\Gamma }^{(m)}(0)`$ are all of order $`d^0`$ or $`K(0)`$ for $`m=0,1,2`$. Recalling that the extrinsic curvature on the brane is at worse logarithmically divergent in $`d`$, terms of the form $`dK(0)`$ will hence be negligible compared to terms of order $`d^0`$ (and of course compared to terms of order $`K(0)`$). Compared to the $`d^1`$ divergence, $`K(0)`$ is hence still negligible. In what follows, terms of order $`d^0`$ (such as $`K^{}(0)`$, $`\mathrm{\Gamma }(0)`$ and $`\mathrm{\Gamma }^{\prime \prime }(0)`$) and terms of order $`K(0)`$ (such as $`K(0)`$, $`K^{\prime \prime }(0)`$ and $`\mathrm{\Gamma }^{}(0)`$) can hence be treated in a similar way. Since they are all at worse going as $`K(0)`$, we shall use in what follows the notation $`K^{(m)}(0)\mathrm{\Gamma }^{(m)}(0)K(0)`$ for $`m=0,1,2`$. We shall hence take as the inductive hypothesis that this result is true for all $`mn`$. In particular,
$`\mathrm{\Gamma }^{(m)}(0)`$ $``$ $`K(0)0mn`$
$`q^{(m+1)}(0)`$ $``$ $`dK(0)0mn`$
$`q(0)`$ $``$ $`d^0.`$
Writing $`l=n1`$, we have
$`K^{(n+1)}=_y^{(l)}K^{\prime \prime }`$ (83)
$`=_y^{(l)}[d(^2+\mathrm{\Gamma }+\mathrm{\Gamma }+\mathrm{\Gamma }^2)q^{}`$
$`+dK^3+dK+dKK^{}+^2dq^{}`$
$`+qd(\mathrm{\Gamma }^{}+d\mathrm{\Gamma }K+dK)].`$
Now, evaluating on the brane, we examine the order of each of these terms to find which is the dominant. For example, remembering that $``$ and $`_y`$ commute, we have for $`l0`$,
$$_y^{(l)}\left(d^2q^{}\right)|_{y=0}=d^2\left(q^{(n)}(0)\right)dK(0)$$
and similarly
$`_y^{(l)}\left(d\left(\mathrm{\Gamma }+\mathrm{\Gamma }+\mathrm{\Gamma }^2\right)q^{}+dK^3+dK\right)|_{y=0}`$ $``$ $`dK(0)`$
$`_y^{(l)}\left(dKK^{}+qd\left(d\mathrm{\Gamma }K+dK\right)+^2dq^{}\right)|_{y=0}`$ $``$ $`dK(0).`$
Finally,
$$_y^{(l)}\left(q\mathrm{\Gamma }^{}d\right)|_{y=0}=d\underset{m=0}{\overset{l}{}}\left(\begin{array}{c}m\\ l\end{array}\right)q^{(lm)}(0)\mathrm{\Gamma }^{(l+1)}(0),$$
and the $`l=0`$ term dominates this last sum, being of order $`K(0)`$. Hence the dominant term in the expression for $`K_\nu ^{\mu (n+1)}`$ is, as in the $`K_\nu ^{\mu (2)}`$ case, the one with the derivative of the Christoffel symbol, of the same order as $`K(0)`$.
We have now proved half of the inductive hypothesis, but still need to show that $`\mathrm{\Gamma }^{(n+1)}(0)K(0)`$. From (72), we have
$`\mathrm{\Gamma }_{\mu \nu }^{\alpha (n+1)}(y)=`$ $`_\mu (dK_\nu ^{\alpha (n)})+_\nu (dK_\mu ^{\alpha (n)})^\alpha (dK_{\mu \nu }^{(n)})`$ (87)
$`+d_y^{(n)}[\mathrm{\Gamma }K\mathrm{terms}]`$
$`{\displaystyle \underset{m=0}{\overset{n1}{}}}\left(\begin{array}{c}m\\ n1\end{array}\right)\left(dK_\nu ^{\rho (m)}\right)_{,\beta }_y^{(nm)}\left(q^{\alpha \beta }q_{\mu \rho }\right).`$
By the inductive hypothesis
$$d_y^{(n)}\left(\mathrm{\Gamma }K\right)dK(0),_y^{(nm)}\left(q_{\mu \rho }^{\alpha \beta }\right)dK(0)n>m$$
and
$$_\alpha \left(dK_\nu ^{\mu (n)}\right)=d_{,\alpha }K_\nu ^{\mu (n)}+𝒪(dK(0)).$$
Therefore we can now read off from (87) the leading order contribution to $`\mathrm{\Gamma }^{(n)}`$,
$`\mathrm{\Gamma }_{\mu \nu }^{\alpha (n+1)}(0)`$ $`=`$ $`d_{,\mu }K_\nu ^{\alpha (n)}+d_{,\nu }K_\mu ^{\alpha (n)}d^{,\alpha }K_{\mu \nu }^{(n)}`$ (88)
$``$ $`K(0),`$
which agrees with (81) and proves the inductive hypothesis.
We are now finally in a position to calculate $`𝒦_\nu ^{\mu (n)}`$ for general $`n`$. We know that the leading contribution obtained from repeatedly differentiating the right-hand side of (77) is
$$K_\nu ^{\mu (n)}q^{\mu \beta }\mathrm{\Gamma }_{\beta \nu }^{\alpha (n1)}_\alpha d$$
which, from (88), immediately gives (69) and the result is proved.
## Appendix B Scalar Perturbations in an FRW background
In this appendix, we shall present some of the details for the calculation of scalar perturbations of section V.2, with the metric perturbation as given in (V.2). We recall that we picked the comoving gauge for which $`\delta d=0`$. In that gauge we then have:
$$\delta \left|d\right|=\dot{d}\mathrm{\Phi }$$
Terms that appear to be sub-dominant will only be dropped at the end. Using (30), we get:
$$\delta R=\frac{\kappa ^2}{L}\left(\delta \rho 3\delta p\right).$$
(89)
Since $`a^{\prime \prime }=0`$ for the background (we assume the brane to be empty), (89) implies
$`\mathrm{\Psi }^{\prime \prime }`$ $`+`$ $`{\displaystyle \frac{a^{}}{a}}\left(2k^2E+\mathrm{\Phi }^{}+3\mathrm{\Psi }^{}\right)`$
$`+`$ $`{\displaystyle \frac{k^2}{3}}\left(2\mathrm{\Psi }+2E^{}\mathrm{\Phi }\right)={\displaystyle \frac{\kappa ^2}{6L}}a^2\left(\delta \rho 3\delta p\right).`$
We now perturb (III.2), writing
$`z`$ $`=`$ $`\left|d\right|`$
$`f(z)`$ $`=`$ $`z\mathrm{tanh}(z/2)`$
$`g(z)`$ $`=`$ $`{\displaystyle \frac{1}{z}}\left(\mathrm{tanh}(z/2)+\mathrm{tan}(z/2)\right)`$
for simplicity. The $`ij`$ (with $`ij`$) component of the Einstein equation, to first order in the perturbations, reduces to:
$$\mathrm{\Phi }\mathrm{\Psi }2E^{}=\left(4\frac{a^{}}{a}+2\frac{d^{}}{d}\right)E,$$
(91)
and the $`0i`$-component to:
$$\mathrm{\Psi }^{}=\left(\frac{a^{}}{a}+\frac{d^{}}{2d}\right)\mathrm{\Phi }.$$
(92)
So far these equations are equivalent to those one would have obtained in the low-energy limit. The difference comes from the $`00`$-component of the perturbed Einstein equations:
$`{\displaystyle \frac{d^{}}{d}}\mathrm{\Phi }^{}+{\displaystyle \frac{a^2}{Ld}}\left(2f\dot{d}f^{}+\dot{d}^3g^{}\right)\mathrm{\Phi }2k^2\mathrm{\Psi }6\mathrm{\Psi }^{}`$
$`4k^2E={\displaystyle \frac{\kappa ^2}{6}}a{\displaystyle \frac{d^{}}{d}}\mathrm{cot}\dot{d}\left(2\delta \rho +3\delta p\right)`$ (93)
and from the equation of motion for $`d`$:
$`\mathrm{\Phi }^{}+\left(8{\displaystyle \frac{d^{}a^2}{L}}f+{\displaystyle \frac{a}{L}}\left(g^{}\dot{d}^24f^{}\right)\right)\mathrm{\Phi }+3\mathrm{\Psi }^{}`$ (94)
$`+2k^2E={\displaystyle \frac{\kappa ^2}{6}}a\left(\mathrm{cot}\dot{d}\left(2\delta \rho +3\delta p\right)3\mathrm{coth}\dot{d}\delta \rho \right).`$
Note that one must, at this order, treat $`\mathrm{\Phi }`$,$`\mathrm{\Phi }^{}`$,$`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^{}`$ as four independent variables; differentiation with respect to conformal time will miss terms arising from higher order in $`d`$, since $`d`$ and $`d^{}`$ are of different order. We must then solve the five equations (B-94) simultaneously. Using (91) to eliminate $`\mathrm{\Phi }`$ from (B), we obtain
$`4k^2\left({\displaystyle \frac{d^{}}{d}}\right)E`$ $`=`$ $`6\mathrm{\Psi }^{\prime \prime }+2k^2\mathrm{\Psi }+6\left(\mathrm{\Phi }^{}+3\mathrm{\Psi }^{}\right)`$
$`{\displaystyle \frac{\kappa ^2}{L}}a^2\left(\delta \rho 3\delta p\right)`$
We may use a combination of (93) and (94) to find an expression for $`\mathrm{\Phi }^{}`$ in terms of $`\mathrm{\Phi },\mathrm{\Psi },\rho `$ and $`p`$ and hence write $`E`$ in terms of $`\mathrm{\Psi },\mathrm{\Psi }^{},\mathrm{\Psi }^{\prime \prime },\rho `$ and $`p`$. This can then be used in (93) to obtain a complicated expression for $`\mathrm{\Psi }`$ in terms of $`\mathrm{\Psi }^{}`$,$`\mathrm{\Psi }^{\prime \prime }`$,$`\delta \rho `$ and $`\delta p`$. We then only keep the leading order in $`d`$ for each coefficient, resulting in the much simplified equation
$$\mathrm{\Psi }^{\prime \prime }+\frac{d^{}}{d}\mathrm{\Psi }^{}+k^2\mathrm{\Psi }=\frac{\kappa ^2}{6}a^2\frac{\dot{d}}{d}\mathrm{coth}\dot{d}\delta \rho .$$
(96)
|
warning/0506/math0506396.html
|
ar5iv
|
text
|
# Singular surfaces, mod 2 homology, and hyperbolic volume, I
## 1. Introduction and general conventions
Let $`M`$ be any closed, orientable, hyperbolic $`3`$-manifold. The volume of $`M`$ is known to be an extremely powerful topological invariant, but its relationship to more classical topological invariants remains elusive. The main geometrical result of this paper, Theorem 9.6, asserts that if $`VolM3.08`$ then $`H_1(M,_2)`$ has rank at most $`6`$.
The Weeks-Hodgson census of closed hyperbolic $`3`$-manifolds contains two examples, m135(-1,3) and m135(1,3), for which the volume is $`<3.08`$ and the rank of the first homology with $`_2`$ coefficients is $`3`$. (They are both of volume $`2.666745\mathrm{}`$, and they have integer first homology isomorphic to $`_2_2_4`$ and $`_2_4_4`$ respectively.) There are no examples in that census for which the volume is $`<3.08`$ and the rank of the first homology with $`_2`$ coefficients is $`4`$. Thus there is still a substantial gap between our results and the known examples. However, the bound on the rank of $`H_1(M;_2)`$ given in this paper seems to be better by orders of magnitude than what could be readily deduced by previously available methods.
The proof of Theorem 9.6 relies on a purely topological result, Theorem 8.13, which states that if $`M`$ is a closed $`3`$-manifold which is simple (see 1.10), if $`\pi _1(M)`$ has a subgroup isomorphic to a genus-$`g`$ surface group for a given integer $`g`$, and if the rank of $`H_1(M;_2)`$ is at least $`4g1`$, then $`M`$ contains a connected incompressible closed surface of genus $`g`$. This may be regarded as a partial analogue of Dehn’s lemma for $`\pi _1`$-injective genus-$`g`$ surfaces.
Theorem 9.6 will be proved in Section 9 by combining Theorem 8.13 with a number of deep geometric results. These include the Marden tameness conjecture, recently established by Agol and by Calegari-Gabai ; a co-volume estimate for $`3`$-tame, $`3`$-free Kleinian groups due to Anderson, Canary, Culler and Shalen \[5, Proposition 8.1\]; and a volume estimate for hyperbolic Haken manifolds recently proved by Agol, Storm and Thurston .
The results of depend in turn on estimates developed by Perelman in his work on geometrization of $`3`$-manifolds.
By refining the methods of this paper one can obtain improvements of Theorems 8.13 and 9.6. In particular, in the case $`g=2`$, the lower bound of $`7`$ for the rank of $`H_1(M;_2)`$ in the hypothesis of Theorem 8.13 can be replaced by $`6`$, and the upper bound of $`6`$ in the conclusion of Theorem 9.6 can be replaced by $`5`$. The relevant refinements will be explored systematically in .
Our strategy for proving Theorem 8.13 is based on the method of two-sheeted coverings used by Shapiro and Whitehead in their proof of Dehn’s Lemma. (This method was inspired by Papakyriakopoulos’s tower construction , and was systematized by Stallings .) We consider a $`\pi _1`$-injective genus-$`g`$ singular surface in the $`3`$-manifold $`M`$, i.e. a map $`\varphi :KM`$, where $`K`$ is a closed orientable genus-$`g`$ surface, and $`\varphi _{\mathrm{}}`$ is injective. One can construct a “tower”
where the $`M_j`$ are simple (1.10) $`3`$-manifolds, $`N_j`$ is a simple $`3`$-dimensional submanifold of $`M_j`$ for $`j=0,\mathrm{},n`$, the $`p_j:M_jN_{j1}`$ are two-sheeted covering maps, $`\stackrel{~}{\varphi }_{}:H_1(K;_2)H_1(N_n;_2)`$ is surjective, and the diagram commutes up to homotopy. In general this diagram may contain both closed and bounded manifolds, but we use ideas from to construct the tower in such a way that if $`H_1(M,_2)`$ has rank $`4g1`$, then $`H_1(M_j,_2)`$ has rank $`4g2`$ whenever $`M_j`$ is closed. We also use ideas developed in based on Simon’s results on compactification of covering spaces, to construct the tower in such a way that the (possibly empty and possibly disconnected) surface $`N_j`$ is incompressible in $`M_j`$ for each $`jn`$.
The manifold $`N_n`$ always has non-empty boundary. This is obvious if $`M_n\mathrm{}`$. If $`M_n`$ is closed then $`H_1(M_n;_2)`$ has rank at least $`4g2`$, whereas the surjectivity of $`\stackrel{~}{\varphi }_{}:H_1(K;_2)H_1(N_n;_2)`$ implies that the rank of $`H_1(N_n;_2)`$ is at most $`2g`$. It follows that in this case $`N_n`$ is a proper submanifold of $`M_n`$, and hence $`N_n\mathrm{}`$.
We in fact show, using elementary arguments based on Poincaré-Lefschetz duality, that if the map $`\stackrel{~}{\varphi }_{}:H_2(K;)H_2(N_n;)`$ is trivial, then $`N_n`$ has a component $`F`$ of genus at most $`g`$. In the case where $`\stackrel{~}{\varphi }_{}:H_2(K;)H_2(N_n;)`$ is non-trivial, we use Gabai’s results from to show that $`N_n`$ contains a non-separating incompressible closed surface $`F`$ of genus at most $`g`$.
The rest of the proof consists of showing that if a given $`M_j`$, with $`0<jn`$, contains a closed incompressible surface of genus at most $`g`$, then $`N_{j1}`$ also contains such a surface. The surface in $`N_{j1}`$ will be incompressible in $`M_{j1}`$, as well as in $`N_{j1}`$, because $`N_{j1}`$ is incompressible in $`M_{j1}`$. It is at this step that we need to know that closed manifolds in the tower have first homology with $`_2`$-coefficients of rank at least $`4g2`$. Indeed, Proposition 4.4 implies that the existence of a closed incompressible surface of genus at most $`g`$ in a $`2`$-sheeted covering of a simple compact $`3`$-manifold $`N`$ implies the existence of such a surface in $`N`$ itself unless $`N`$ is closed and $`H_1(N;_2)`$ has rank at most $`4g3`$.
Proposition 4.4 involves the notion of a “book of $`I`$-bundles” which we define formally in 2.2. Books of $`I`$-bundles in PL $`3`$-manifolds arise naturally as neighborhoods of “books of surfaces” (2.6). We may think of a book of surfaces as being constructed from a $`2`$-manifold with boundary $`\widehat{\mathrm{\Pi }}`$, whose components have Euler characteristic $`0`$, and a closed $`1`$-manifold $`\mathrm{\Psi }`$, by attaching $`\widehat{\mathrm{\Pi }}`$ to $`\mathrm{\Psi }`$ by a covering map. The components of $`\mathrm{\Psi }`$ and $`\mathrm{\Pi }=int\widehat{\mathrm{\Pi }}`$ are respectively “bindings” and “pages.” A book of $`I`$-bundles comes equipped with a corresponding decomposition into “pages” which are $`I`$-bundles over surfaces, and “bindings” which are solid tori. (In the informal discussion that we give in this introduction, the extra structure defined by the decomposition will be suppressed from the notation.)
With these notions as background we shall now sketch the proof of Proposition 4.4. An incompressible surface $`F`$ in a two-sheeted covering space of $`N`$, if it is in general position, projects to $`N`$ via a map which has only double-curve singularities. After routine modifications one obtains a map from $`F`$ to $`N`$ with the additional property that its double curves are homotopically non-trivial. In particular, the image of such a map is a book of surfaces $`X`$. A regular neighborhood $`W`$ of $`X`$ in $`N`$ is then a book of $`I`$-bundles, which has Euler characteristic $`22g`$ if $`F`$ has genus at most $`g`$. Using the the simplicity of $`N`$ one can then produce a book of $`I`$-bundles $`V`$ with $`WVN`$ and $`\chi (W)22g`$, such that each page of $`W`$ has strictly negative Euler characteristic. (This step is handled by Lemma 2.5.)
We now distinguish two cases. In the case where some page $`P_0`$ of $`V`$ has the property that $`P_0V`$ is contained in a single component of $`V`$, we show that by splitting bindings of the book of surfaces $`X`$, one can produce an embedded (possibly disconnected) closed, orientable surface $`S`$ which is homologically non-trivial in $`N`$. Ambient surgery on $`S`$ in $`N`$ then produces a non-empty incompressible surface whose components have genus at most $`g`$. In the case where no such page $`P_0`$ exists, an Euler characteristic calculation shows that the boundary components of $`V`$ have genus at most $`g`$. In this case, ambient surgery on $`V`$ produces a non-empty incompressible surface whose components have genus at most $`g`$. We show that this surface is non-empty unless $`V`$ carries $`\pi _1(N)`$. But for a book of $`I`$-bundles $`V`$ whose Euler characteristic is at least $`22g`$, and whose pages are all of negative Euler characteristic, one can show that $`H_1(V;_2)`$ has rank at most $`4g3`$ (this is included in Lemma 2.11); so in the case where $`V`$ carries $`\pi _1(N)`$, the rank of $`H_1(N;_2)`$ is at most $`4g3`$.
The details and background needed for the proof of Proposition 4.4 occupy Sections 24. Section 5 provides the combinatorial background needed to construct the tower, while Sections 6 and 7 provide the homological background. The application of Gabai’s results mentioned above appears in Section 7. The material on towers proper, and the proof of the main topological theorem and its corollary, are given in Section 8, and the geometric applications are given in Section 9.
The rest of this introduction will be devoted to indicating some conventions that will be used in the rest of the paper.
###### 1.1.
In general, if $`X`$ and $`Y`$ are subsets of a set, we denote by $`XY`$ the set of elements of $`X`$ that do not belong to $`Y`$. In the case where we know that $`YX`$ and wish to emphasize this we will write $`XY`$ for $`XY`$.
###### 1.2.
A manifold may have a boundary. If $`M`$ is a manifold, we shall denote the boundary of $`M`$ by $`M`$ and its interior $`MM`$ by $`intM`$.
In many of our results about manifolds of dimension $`3`$ we do not specify a category. These results may be interpreted in the category in which manifolds are topological, PL or smooth, and submanifolds are respectively locally flat, PL or smooth; these three categories are equivalent in these low dimensions as far as classification is concerned. In much of the paper the proofs are done in the PL category, but the applications to hyperbolic manifolds in Section 9 are carried out in the smooth category.
###### 1.3.
A (possibly disconnected) codimension-$`1`$ submanifold $`S`$ of a manifold $`M`$ is said to be separating if $`M`$ can be written as the union of two $`3`$-dimensional submanifolds $`M_1`$ and $`M_2`$ such that $`M_1M_2=S`$.
###### 1.4.
We shall say that a map of topological spaces $`f:𝒳𝒴`$ is $`\pi _1`$-injective if for every path component $`X`$ of $`𝒳`$, the map $`f|X`$ induces an injection from $`\pi _1(X)`$ to $`\pi _1(Y)`$, where $`Y`$ is the path component of $`𝒴`$ containing $`f(X)`$. We shall say that a subset $`A`$ of a topological space $`X`$ is $`\pi _1`$-injective in $`X`$ if the inclusion map $`AX`$ is $`\pi _1`$-injective.
###### 1.5.
If $`X`$ is a space having the homotopy type of a finite CW complex, the Euler characteristic of $`X`$ will be denoted by $`\chi (X)`$. We have $`\chi (X)=_jdim_FH_j(X;F)`$ for any field $`F`$: the sum is independent of $`F`$ by virtue of the standard observation that it is equal to $`_j(1)^jc_j`$, where $`c_j`$ denotes the number of $`j`$-cells in a finite CW complex homotopy equivalent to $`X`$.
We shall often write $`\overline{\chi }(X)`$ as shorthand for $`\chi (X)`$.
###### 1.6.
If $`x`$ is a point of a compact PL space $`X`$, there exist a finite simplicial complex $`K`$ and a PL homeomorphism $`h:X|K|`$ such that $`v=h(x)`$ is a vertex of $`K`$. If $`L`$ denotes the link of $`v`$ in $`K`$ then the PL homeomorphism type of the space $`|L|`$ depends only on $`X`$ and $`x`$, not on the choice of $`K`$ and $`h`$. We shall refer to $`L`$ as the link of $`x`$ in $`X`$, with the understanding that it is defined only up to PL homeomorphism.
###### 1.7.
Suppose that $`x`$ is a point of a compact PL space $`X`$ and that $`n0`$ is an integer. The link of $`x`$ is PL homeomorphic to $`S^{n1}`$ if and only if $`x`$ is an $`n`$-manifold point of $`X`$, i.e. some neighborhood of $`x`$ is piecewise-linearly homeomorphic to $`𝐑^n`$.
If $`X`$ is a compact PL space of dimension at most $`2`$, we shall denote by $`\mathrm{\Pi }(X)`$ the set of all $`2`$-manifold points of $`X`$. Note that $`\mathrm{\Pi }(X)`$ is an open subset of $`X`$, and with its induced PL structure it is a PL $`2`$-manifold. Furthermore, $`X\mathrm{\Pi }(X)`$ is a compact PL subspace of $`X`$.
###### 1.8.
Let $`F`$ be a properly embedded orientable surface in an orientable $`3`$-manifold $`M`$. We define a compressing disk for $`F`$ in $`M`$ to be a disk $`DM`$ such that $`DF=D`$, and such that $`D`$ is not the boundary of a disk in $`F`$. It is a standard consequence of the loop theorem that $`F`$ is $`\pi _1`$-injective in $`M`$ if and only if there is no compressing disk for $`F`$ in $`M`$.
A closed orientable surface $`S`$ contained in the interior of an orientable $`3`$-manifold $`M`$ will be termed incompressible if $`S`$ is $`\pi _1`$-injective in $`M`$ and no component of $`S`$ is a $`2`$-sphere. (We have avoided using the term “incompressible” for surfaces that are not closed.)
###### 1.9.
An essential arc in a $`2`$-manifold $`F`$ is a properly embedded arc in $`F`$ which is not the frontier of a disk.
###### Definitions 1.10.
A $`3`$-manifold $`M`$ will be termed irreducible if every $`2`$-sphere in $`M`$ bounds a ball in $`M`$. We shall say that $`M`$ is boundary-irreducible if $`M`$ is $`\pi _1`$-injective in $`M`$, or equivalently if, for every properly embedded disk $`DM`$, the simple closed curve $`D`$ bounds a disk in $`M`$. We shall say that a $`3`$-manifold $`M`$ is simple if (i) $`M`$ is compact, connected, orientable, irreducible and boundary-irreducible; (ii) no subgroup of $`\pi _1(M)`$ is isomorphic to $`\times `$; and (iii) $`M`$ is not a closed manifold with finite fundamental group.
###### 1.11.
It is a theorem due to Meeks, Simon and Yau that a covering space of an irreducible orientable $`3`$-manifold is always irreducible. Given this result, it follows formally from our definition of simplicity that if a compact, orientable $`3`$-manifold $`M`$ is simple, then every connected finite-sheeted covering of $`M`$ is also simple.
###### 1.12.
The unit interval $`[0,1]`$ will often be denoted by $`I`$.
By an $`I`$-bundle we shall always mean a compact space equipped with a specific locally trivial fibration over some (often unnamed) base space, in which the fibers are homeomorphic to $`[0,1]`$. (The reader is referred to \[11, Chapter 10\] for a general discussion of $`3`$-dimensional $`I`$-bundles.)
By a Seifert fibered manifold we shall always mean a compact $`3`$-manifold equipped with a specific Seifert fibration.
In particular, the notion of a fiber of an $`I`$-bundle or a Seifert fibered manifold is well defined, although the fiber projection and base space will often not be explicitly named. A compact subset of an $`I`$-bundle or Seifert fibered space will be called horizontal if it meets each fiber in one point. A compact set will be called vertical if it is a union of fibers.
If $`𝒫`$ is an $`I`$-bundle, we define the horizontal boundary of $`𝒫`$ to be the subset of $`𝒫`$ consisting of all endpoints of fibers of $`𝒫`$. We shall denote the horizontal boundary of $`𝒫`$ by $`_h𝒫`$.
In the case where the base of the $`I`$-bundle $`𝒫`$ is a $`2`$-manifold $`F`$ (so that $`𝒫`$ is a $`3`$-manifold), we define the vertical boundary of $`𝒫`$ to be $`p^1(F)`$, where $`p:𝒫F`$ denotes the bundle map. Note that in this case we have $`𝒫=_v𝒫_h𝒫`$, and if $`𝒫`$ is orientable then $`_v𝒫`$ is always a finite disjoint union of annuli.
###### 1.13.
The rank of a finitely generated group $`\mathrm{\Gamma }`$ is the cardinality of a minimal generating set for $`\mathrm{\Gamma }`$. In particular, the trivial group has rank $`0`$.
A group $`\mathrm{\Gamma }`$ is said to be freely indecomposable if $`\mathrm{\Gamma }`$ is not the free product of two non-trivial subgroups.
###### 1.14.
If $`V`$ is a finite dimensional vector space over $`_2`$ then the dimension of $`V`$ will be denoted $`\mathrm{rk}_2V`$. If $`X`$ is a topological space, we will set $`\mathrm{rk}_2X=\mathrm{rk}_2H_1(X;_2)`$.
## 2. Books of $`I`$-bundles
###### Definition 2.1.
A generalized book of $`I`$-bundles is a triple $`𝒲=(W,,𝒫)`$, where $`W`$ is a (possibly empty) compact, orientable $`3`$-manifold, and $`,𝒫W`$ are submanifolds such that
* $``$ is a (possibly disconnected) Seifert fibered space,
* $`𝒫`$ is an $`I`$-bundle over a (possibly disconnected) $`2`$-manifold, and every component of $`𝒫`$ has Euler characteristic $`0`$,
* $`W=𝒫`$,
* $`𝒫`$ is the vertical boundary of $`𝒫`$, and
* $`𝒫`$ is vertical in the Seifert fibration of $``$.
We shall denote $`W`$, $``$ and $`𝒫`$ by $`|𝒲|`$, $`_𝒲`$ and $`𝒫_𝒲`$ respectively. The components of $`_𝒲`$ will be called bindings of $`𝒲`$, and the components of $`𝒫_𝒲`$ will be called its pages. The submanifold $`𝒫`$, whose components are properly embedded annuli in $`W`$, will be denoted $`𝒜_𝒲`$.
If $`B`$ is a binding of a generalized book of $`I`$-bundles $`𝒲`$, we define the valence of $`B`$ to be the number of components of $`𝒜_𝒲`$ that are contained in $`B`$.
A generalized book of $`I`$-bundles $`𝒲`$ will be termed connected if the manifold $`|𝒲|`$ is connected. Likewise, $`𝒲`$ will be termed boundary-irreducible if $`|𝒲|`$ is boundary-irreducible.
###### Definitions 2.2.
A book of $`I`$-bundles is a generalized book of $`I`$-bundles $`𝒲`$ such that
* $`|𝒲|\mathrm{}`$,
* each binding of $`𝒲`$ is a solid torus, and
* each binding of $`𝒲`$ meets at least one page of $`𝒲`$.
If $`B`$ is a binding of a book of $`I`$-bundles $`𝒲`$, there is a unique integer $`d>0`$ such that for every component $`A`$ of $`𝒜_𝒲`$ contained in $`B`$, the image of the inclusion homomorphism $`H_1(A;)H_1(B;)`$ has index $`d`$ in $`H_1(B;)`$. We shall call $`d`$ the degree of the binding $`B`$.
###### Lemma 2.3.
Suppose that $`𝒲`$ is a generalized book of $`I`$-bundles. Then there is a generalized book of $`I`$-bundles $`𝒲_0`$ such that
1. $`|𝒲_0|=|𝒲|`$,
2. every page of $`𝒲_0`$ has strictly negative Euler characteristic, and
3. every page of $`𝒲_0`$ is a page of $`𝒲`$.
###### Proof.
Set $`W=|𝒲|`$, $`=_𝒲`$ and $`𝒫=𝒫_𝒲`$. Let $`𝒬`$ denote the union of all components $`P`$ of $`𝒫`$ such that $`\chi (P)=0`$. Then $`𝒬`$ is an $`I`$-bundle over a compact surface $`A`$ whose components are annuli and Möbius bands, and $`𝒬`$ is the induced $`I`$-bundle over $`A`$. Hence every component $`Q`$ of $`𝒬`$ is a solid torus, and $`Q`$ consists of either a single annulus of degree $`2`$ in $`Q`$, or of two annuli of degree $`1`$ in $`Q`$. Since every such annulus is also vertical in the Seifert fibration of $``$, it follows that this Seifert fibration may be extended to a Seifert fibration of the manifold $`_0=𝒬`$, in such a way that each component of $`𝒬`$ contains either no singular fiber, or exactly one singular fiber of order $`2`$. Furthermore, since every component of $`𝒬`$ meets $``$, every component of $`_0`$ contains a component of $``$.
The manifold $`𝒫_0=𝒫𝒬`$ is a union of components of $`𝒫`$ and therefore inherits an $`I`$-bundle structure. It is now clear that $`𝒲_0=(W,_0,𝒫_0)`$ is a generalized book of $`I`$-bundles. It follows from the definition of $`𝒬`$ that $`𝒲_0`$ satisfies conclusions (2) and (3) of the lemma. ∎
###### Lemma 2.4.
Suppose that $`\widehat{B}`$ is a connected, Seifert-fibered submanifold of a simple, closed, orientable $`3`$-manifold $`M`$. Then either
1. $`\widehat{B}`$ is a solid torus, or
2. $`\widehat{B}`$ is contained in a ball in $`M`$, or
3. some component of $`Mint\widehat{B}`$ is a solid torus.
###### Proof.
Since $`M`$ is simple and $`\widehat{B}`$ is Seifert-fibered, we have $`\widehat{B}M`$, i.e. $`\widehat{B}\mathrm{}`$. Since the components of $`\widehat{B}`$ are tori and $`M`$ is simple, $`\widehat{B}`$ cannot be $`\pi _1`$-injective in $`M`$. Hence there is a compressing disk for $`\widehat{B}`$ in $`M`$. If $`D\widehat{B}`$ then $`\widehat{B}`$ is a boundary-reducible Seifert fibered space and hence (1) holds. The other possibility is that $`D\widehat{B}=D`$. In this case, let $`V`$ denote a regular neighborhood of $`D`$ relative to $`Mint\widehat{B}`$. The boundary of the manifold $`\widehat{B}V`$ has a unique sphere component $`S`$. Since $`M`$ is irreducible, $`S`$ bounds a ball $`\mathrm{\Delta }M`$. We must have either $`\mathrm{\Delta }\widehat{B}`$, which gives conclusion (2), or $`int\mathrm{\Delta }\widehat{B}=\mathrm{}`$; in the latter case, $`\mathrm{\Delta }V`$ is a solid torus component of $`Mint\widehat{B}`$, and so (3) holds. ∎
###### Lemma 2.5.
Suppose that $`M`$ is a simple, closed, orientable $`3`$-manifold, and that $`𝒲`$ is a connected generalized book of $`I`$-bundles such that $`W=|𝒲|M`$. Suppose that $`\chi (W)<0`$, and that $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`M`$. Then there is a connected book of $`I`$-bundles $`𝒱`$ with $`V=|𝒱|M`$, such that
1. $`VW`$,
2. $`\overline{\chi }(V)=\overline{\chi }(W)`$,
3. $`\chi (P)<0`$ for every page $`P`$ of $`𝒱`$,
4. $`V`$ is a union of components of $`W`$,
5. every component of $`\overline{VW}`$ is a solid torus,
6. every page of $`𝒱`$ is a page of $`𝒲`$, and
7. for each page $`P`$ of $`𝒱`$ we have $`PV=PW`$.
###### Proof.
Let $`𝒲_0=(W,,𝒫)`$ be a generalized book of $`I`$-bundles satisfying conditions (1)–(3) of Lemma 2.3. Since each page of $`𝒲_0`$ is also a page of $`𝒲`$, the hypothesis implies that each page of $`𝒲_0`$ is $`\pi _1`$-injective in $`M`$.
Let $`B`$ be any binding of $`𝒲_0`$. We will show that the Seifert fibers of $`B`$ are homotopically non-trivial in $`M`$. Since $`𝒲_0`$ is connected and $`\chi (|𝒲_0|)<0`$, the binding $`B`$ must meet some page $`P`$ of $`𝒲_0`$. Let $`A`$ be one of the annulus components of $`BP`$. Then $`A`$ is a component of the vertical boundary of $`P`$ and, since $`\chi (P)<0`$, it follows that $`A`$ is $`\pi _1`$-injective in $`P`$. Since $`P`$ is $`\pi _1`$-injective in $`M`$, it follows that $`A`$ is also $`\pi _1`$-injective in $`M`$. Recalling that the annulus $`A`$ is saturated in the Seifert fibration of $`B`$, we may conclude that each Seifert fiber of $`B`$ is homotopically non-trivial in $`M`$.
Now for any binding $`B`$ of $`𝒲_0`$ let us define $`\widehat{B}`$ to be the union of $`B`$ with all of the solid torus components of $`\overline{MB}`$. We will show that $`\widehat{B}`$ is a Seifert fibered submanifold of $`M`$ such that $`\widehat{B}𝒲_0=B`$. If $`J`$ is any solid torus component of $`\overline{MB}`$ then no page of $`𝒱_0`$ can be contained in $`J`$, since the pages are $`\pi _1`$-injective in $`M`$ and have negative Euler characteristic. Thus $`intJ`$ must be disjoint from all of the pages of $`𝒲_0`$. This implies that $`\widehat{B}𝒲_0=B`$. If $`FJ`$ is a fiber of the Seifert fibered space $`B`$ then, since $`F`$ is homotopically non-trivial in $`M`$, the simple closed curve $`FJ`$ cannot be a meridian curve for the solid torus $`J`$. It follows that the Seifert fibration of $`B`$ may be extended to a Seifert fibration of $`B=BJ`$, and hence that $`\widehat{B}`$ admits a Seifert fibration.
Next we will show that $`\widehat{B}`$ is, in fact, a solid torus. We know that $`\widehat{B}`$ must satisfy one of the conditions (1)—(3) of Lemma 2.4. Condition 2.4(3) is ruled out since, by construction, no component of $`Mint\widehat{B}`$ is a solid torus. The fact that the Seifert fibers of $`B`$ are homotopically non-trivial in $`M`$ implies that the inclusion homomorphism $`\pi _1(B)\pi _1(M)`$ has non-trivial image and thus $`B`$ cannot be contained in a ball in $`M`$. This rules out condition 2.4(2). Thus we conclude that condition 2.4(1) must hold, i.e that $`\widehat{B}`$ is a solid torus.
Since each binding of $`𝒲`$ must meet some page, and since no page can be contained in a solid torus, we have that if $`B_1`$ and $`B_2`$ are distinct bindings of $`𝒲_0`$, then $`\widehat{B}_1`$ is disjoint from $`\widehat{B}_2`$. We define $`^{}`$ to be the union of the solid tori $`\widehat{B}`$ as $`B`$ ranges over all bindings of $`𝒲_0`$, and we set $`V=^{}𝒫`$. We have $`^{}|𝒲_0|=`$. It follows that $`𝒱=(V,^{},𝒫)`$ is a book of $`I`$-bundles, and that every page of $`𝒱`$ has strictly negative Euler characteristic.
We shall now complete the proof by observing that $`V`$ satisfies Conclusions (1)—(7) of the present lemma. Conclusions (1), (4) and (5) are immediate from the construction of $`V`$, and they imply Conclusion (2). The pages of $`𝒱`$ are the same as the pages of $`𝒲_0`$, and each page of $`𝒲_0`$ is a page of $`𝒲`$ and has negative Euler characteristic. Hence Conclusions (3) and (6) hold. Since $`W`$ is the union of $`V`$ with a collection of tori that are disjoint from all pages, it follows that $`PV=PW`$ for every page $`P`$ of $`𝒱`$. This is Conclusion (7). ∎
Recall that in 1.7 we defined $`\mathrm{\Pi }(X)X`$ to be the set of $`2`$-manifold points in an arbitrary compact PL space $`X`$ of dimension at most $`2`$, and we observed that $`X\mathrm{\Pi }(X)`$ is a compact PL subset of $`X`$. It follows that $`\mathrm{\Pi }(X)`$ has the homotopy type of a compact PL space. In particular $`\chi (\pi )`$ is a well-defined integer for every component $`\pi `$ of $`\mathrm{\Pi }(X)`$.
###### Definition 2.6.
We define a book of surfaces to be a compact PL space $`X`$ such that
1. the link of every point of $`xX`$ is PL homeomorphic to the suspension of some non-empty finite set $`Z_x`$; and
2. for every component $`\pi `$ of $`\mathrm{\Pi }(X)`$ we have $`\chi (\pi )0`$.
The cardinality of the set $`Z_x`$ appearing in condition (1) is clearly uniquely determined by the point $`x`$. It will be called the order of $`x`$.
###### 2.7.
Note that a point $`x`$ in a book of surfaces $`X`$ has order $`2`$ if and only if $`x\mathrm{\Pi }(X)`$.
It also follows from the definition that if $`X`$ is a book of surfaces, the set $`X\mathrm{\Pi }(X)`$ is a compact PL $`1`$-manifold, which will be denoted by $`\mathrm{\Psi }(X)`$. The components of $`\mathrm{\Psi }(X)`$ and $`\mathrm{\Pi }(X)`$ may be respectively thought of as bindings and pages of $`X`$.
We also observe that if $`M`$ is a PL $`3`$-manifold and if $`S_1`$ and $`S_2`$ are closed surfaces in $`intM`$ which meet transversally, then $`S_1S_2`$ is a book of surfaces.
###### Lemma 2.8.
If $`X`$ is a book of surfaces, there exist a (possibly disconnected) compact PL surface $`F`$ and a PL map $`r:FX`$ such that
1. $`r|intF`$ is a homeomorphism of $`intF`$ onto $`\mathrm{\Pi }(X)`$, and
2. $`r|F`$ is a covering map from $`F`$ to $`\mathrm{\Psi }(X)`$.
###### Proof.
Let us identify $`X`$ with $`|K|`$, where $`|K|`$ is some finite simplicial complex. After subdividing $`K`$ if necessary we may assume that for every closed simplex $`\mathrm{\Delta }`$ of $`K`$ the set $`\mathrm{\Delta }\mathrm{\Psi }(X)`$ is a (possibly empty) closed face of $`\mathrm{\Delta }`$.
Let $`𝒟`$ denote the abstract disjoint union of all the closed $`2`$-simplices of $`X`$, and let $`i:𝒟X`$ denote the map which is the inclusion on each closed $`2`$-simplex. For each point $`z𝒟`$ let $`\mathrm{\Delta }_z`$ denote the closed $`2`$-simplex containing $`z`$. We define a relation $``$ on $`𝒟`$ by writing $`zw`$ if and only if (i) $`\mathrm{\Delta }_z\mathrm{\Delta }_w\mathrm{\Psi }(X)`$ and (ii) $`i(z)=i(w)`$. It is straightforward to show that $``$ is an equivalence relation. The quotient space $`F=\mathrm{\Delta }/`$ inherits a PL structure from $`𝒟`$. The definition of $``$ implies that there is a unique map $`r:FX`$ such that $`rq=i`$, where $`q:𝒟F`$ is the quotient map, and that $`r`$ maps $`E=r^1\mathrm{\Pi }(X)`$ homeomorphically onto $`\mathrm{\Pi }(X)`$.
If $`x`$ is a point of $`\mathrm{\Psi }(X)`$, then since $`X`$ is a book of surfaces, there exist a neighborhood $`A`$ of $`x`$ in $`\mathrm{\Psi }(X)`$, and a neighborhood $`V`$ of $`x`$ in $`X`$, such that $`A`$ is a PL arc, $`V`$ is a union of PL disks $`D_1\mathrm{}D_m`$, where $`m`$ is the order of $`x`$ in $`X`$, and $`D_iD_j=A`$ whenever $`ij`$. The definition of $``$ implies that $`r^1(V)`$ is a disjoint union of PL disks $`\stackrel{~}{D}_1,\mathrm{},\stackrel{~}{D}_m`$ such that $`r`$ maps $`\stackrel{~}{D}_i`$ homeomorphically onto $`D_i`$ for $`i=1,\mathrm{},m`$. Hence $`F`$ is a PL surface with interior $`E`$ and boundary $`r^1(\mathrm{\Psi }(X))`$, and $`r|F:F\mathrm{\Psi }(X)`$ is a covering map. ∎
###### Lemma 2.9.
Suppose that $`M`$ is an orientable PL $`3`$-manifold and that $`XintM`$ is a book of surfaces. Then there is a book of $`I`$-bundles $`𝒲`$ such that
1. $`|𝒲|=W`$ is a regular neighborhood of $`X`$;
2. $`|_𝒲|`$ is a regular neighborhood of $`\psi (X)`$;
3. for every page $`P`$ of $`𝒲`$, the set $`XP`$ is a section of the $`I`$-bundle $`P`$; and
4. $`𝒫_𝒲`$ is a regular neighborhood in $`M`$ of a deformation retract of $`\mathrm{\Pi }(X)`$.
###### Proof.
Let $``$ be a regular neighborhood of $`\mathrm{\Psi }(X)`$ in $`M`$ such that $`N=X`$ is a regular neigborhood of $`\mathrm{\Psi }(X)`$ in the PL space $`X`$. Every component of $``$ is a solid torus. Since $`\mathrm{\Pi }(X)`$ is an open $`2`$-manifold, $`Y=X\overline{M}`$ is a compact $`2`$-manifold and a deformation retract of $`\mathrm{\Pi }(X)`$. In particular, in view of condition (2) in the definition of a book of surfaces, every component of $`Y`$ has Euler characteristic $`0`$. Let $`𝒫`$ be a regular neighborhood of $`Y`$ in $`\overline{MN}`$. Then $`W=𝒫`$ is a regular neighborhood of $`X`$ in $`M`$. We may give $`𝒫`$ the structure of an $`I`$-bundle over $`Y`$ in such a way that $`Y`$ is identified with a section of the bundle. We have $`𝒫=_v𝒫`$, and $`\chi (P)0`$ for every component $`P`$ of $`𝒫`$.
Let $`F`$ be the surface, and $`r:FX`$ the map, given by Lemma 2.8. We have $`N=r(C)`$, where $`C`$ is a collar neighborhood of $`F`$ in $`F`$. Now if $`A`$ is any component of $`_v`$, then $`AY`$ is a component of $`Y`$ and therefore cobounds an annulus component of $`C`$ with some component $`\stackrel{~}{\psi }_A`$ of $`F`$. It follows from 2.8 that $`r|\stackrel{~}{\psi }_A`$ is a covering map of some degree $`d_A`$ to some component $`\psi _A`$ of $`\psi (X)`$. The annulus $`A`$ lies in the boundary of the component $`B_A`$ of $`𝒫`$ containing $`\psi _A`$, and the (unsigned) degree of $`A`$ in the solid torus $`B_A`$ is $`d_A`$. In particular, every component of $`_v𝒫`$ has non-zero degree in the component of $``$ containing it. This implies that $`𝒲=(W,,𝒫)`$ is a book of $`I`$-bundles.
Each page $`P`$ of $`𝒫`$ was constructed as an $`I`$-bundle over a component $`Y_0`$ of $`Y`$, where $`Y_0`$ is identified with a section of the bundle. Since $`Y_0=XP`$, Conclusion (3) of the lemma follows. Conclusions (1), (2) and (4) are immediate from the construction of $`𝒲`$. ∎
###### Lemma 2.10.
Suppose that $`𝒲`$ is a book of $`I`$-bundles, and let $`p`$ denote the number of pages of $`𝒲`$. Then
$$\mathrm{rk}_2H_2(|𝒲|;_2)p.$$
###### Proof.
It is most natural to prove a very mild generalization: if $`𝒲`$ is a generalized book of $`I`$-bundles whose bindings are all solid tori, and if $`p`$ denotes the number of pages of $`𝒲`$, then $`\mathrm{rk}_2H_2(|𝒲|;_2)p`$. We set $`W=|𝒲|`$ and use induction on $`p`$. If $`p=0`$ then the components of $`W`$ are solid tori and hence $`\mathrm{rk}_2H_2(W)=0`$. If $`p>0`$, choose a page $`P`$ of $`𝒲`$, and set $`W^{}=\overline{WP}`$ and $`𝒫^{}=𝒫_𝒲P`$. Then $`𝒫^{}`$ inherits an $`I`$-bundle structure from $`𝒫`$, and $`𝒲^{}=(W^{},,𝒫^{})`$ is a book of $`I`$-bundles with $`p1`$ pages. By the induction hypothesis we have $`\mathrm{rk}_2H_2(W^{})p1`$. On the other hand, if $`F`$ denotes the base surface of the $`I`$-bundle $`P`$, we have
$$H_2(W,W^{};_2)H_2(P,_vP;_2)H_2(F,F;_2)$$
and hence $`\mathrm{rk}_2H_2(W,W^{})=1`$. It follows that
$$\mathrm{rk}_2H_2(W)\mathrm{rk}_2H_2(W^{})+\mathrm{rk}_2H_2(W,W^{})p.$$
###### Lemma 2.11.
If $`𝒲`$ is a book of $`I`$-bundles, and if every page of $`𝒲`$ has strictly negative Euler characteristic, we have
$$\mathrm{rk}_2(|𝒲|)2\overline{\chi }(|𝒲|)+1.$$
###### Proof.
Set $`W=|𝒲|`$. By hypothesis we have $`\overline{\chi }(P)1`$ for every page $`P`$ of $`𝒲`$. Hence if $`P_1,\mathrm{},P_p`$ denote the pages of $`𝒲`$, we have
$$\overline{\chi }(W)=\underset{i=1}{\overset{p}{}}\overline{\chi }(P_i)p.$$
According to Lemma 2.10 we have
$$\mathrm{rk}_2H_2(W;_2)p\overline{\chi }(W).$$
Now $`W`$ is a connected $`3`$-manifold with non-empty boundary. Hence $`\mathrm{rk}_2H_0(W;_2)=1`$, and $`H_j(W;_2)=0`$ for each $`j>2`$. In view of 1.5, we have
$$\overline{\chi }(W)=\mathrm{rk}_2H_1(W;_2)\mathrm{rk}_2H_2(W;_2)1.$$
Hence
$$\mathrm{rk}_2(W)=\mathrm{rk}_2H_1(W;_2)=\overline{\chi }(W)+\mathrm{rk}_2H_2(W;_2)+12\overline{\chi }(W)+1.$$
## 3. Compressing submanifolds
###### Definition 3.1.
If $`𝒮`$ is a closed (possibly empty or disconnected) surface, we define a non-negative integer $`\kappa (𝒮)`$ by
$$\kappa (V)=\underset{S}{}(1+genus(S)^2),$$
where $`S`$ ranges over the components of $`𝒮`$.
###### Lemma 3.2.
Let $`𝒮`$ be a closed (possibly empty or disconnected) surface, let $`A𝒮`$ be a homotopically non-trivial annulus, and let $`𝒮^{}`$ be the surface obtained from the bounded surface $`\overline{𝒮A}`$ by attaching disks $`D_1`$ and $`D_2`$ to its two boundary components. Then $`\kappa (𝒮^{})<\kappa (𝒮)`$.
###### Proof.
Let us index the components of $`𝒮`$ as $`S_0,\mathrm{},S_n`$, where $`n0`$ and $`AS_0`$. If $`S_0A`$ is connected, the components of $`𝒮^{}`$ are $`S_0^{},S_1,\mathrm{},S_n`$, where $`S_0^{}=(S_0A)D_1D_2`$. We then have $`genusS_0^{}=(genusS_0)1`$, so that $`\kappa (𝒮)<\kappa (𝒮^{})`$.
If $`S_0A`$ is disconnected, then $`(S_0A)D_1D_2`$ has two components $`S_0^{}`$ and $`S_0^{\prime \prime }`$. If we denote the respective genera of $`S_0`$, $`S_0^{}`$ and $`S_0^{\prime \prime }`$ by $`g`$, $`g^{}`$ and $`g^{\prime \prime }`$, we have $`g=g^{}+g^{\prime \prime }`$; and since $`A`$ is homotopically non-trivial in $`S_0`$, both $`g^{}`$ and $`g^{\prime \prime }`$ are strictly positive. It follows that $`(1+(g^{})^2)+(1+(g^{\prime \prime })^2)<1+g^2`$, and we again deduce that $`\kappa (𝒮)<\kappa (𝒮^{})`$. ∎
###### 3.3.
Recall that a connected $`3`$-manifold $`H`$ is called a compression body if it can be constructed from a product $`T\times [1,1]`$, where $`T`$ is a connected, closed, orientable $`2`$-manifold, by attaching finitely many $`2`$\- and $`3`$-handles to $`T\times \{1\}`$. One defines $`_+H`$ to be the submanifold $`T\times \{1\}`$ of $`H`$, and one define $`_{}H`$ to be $`H_+H`$.
###### 3.4.
If $`H`$ is a connected compression body, it is clear that $`_+H`$ is connected, and that for each component $`F`$ of $`_{}H`$ we have $`genus(F)genus(_+H)`$.
###### 3.5.
It is a standard observation that a connected compression body $`H`$ with $`_{}H=\mathrm{}`$ is a handlebody.
###### 3.6.
Another standard observation is that any connected compression body $`H`$ with $`_H\mathrm{}`$ can be constructed from a product $`S\times [1,1]`$, where $`S`$ is a possibly disconnected, closed, orientable $`2`$-manifold, by attaching $`1`$\- and $`2`$-handles to $`S\times \{1\}`$. One then has $`_{}H=S\times \{1\}`$.
An immediate consequence of this observation is that if $`H`$ is a connected compression body then $`_{}H`$ is $`\pi _1`$-injective in $`H`$.
###### 3.7.
More generally, we shall define a compression body to be a compact, possibly disconnected $`3`$-manifold $``$ such that each component of $``$ is a compression body in the sense defined above. We define $`_+=_H_+H`$ and $`_{}=_H_{}H`$, where $`H`$ ranges over the components of $``$.
###### Proposition 3.8.
Let $`N`$ be a compact orientable, irreducible $`3`$-manifold, and let $`V`$ be a compact, connected, non-empty $`3`$-submanifold of $`intN`$. Suppose that $`\overline{NV}`$ is $`\pi _1`$-injective in $`N`$. Then at least one of the following conditions holds:
1. $`V`$ is contained in a ball in $`N`$; or
2. $`V\mathrm{}`$, and there exists a connected, incompressible closed surface in $`N`$ whose genus is at most the maximum of the genera of the components of $`V`$; or
3. $`N`$ is closed and every component of $`\overline{NV}`$ is a handlebody.
###### Proof.
First note that if $`V=N`$ then conclusion (3) holds. (The hypothesis $`VintN`$ implies that $`N`$ is closed, and the other assertion of (3) is vacuously true.) Hence we may assume that $`VN`$, so that $`V\mathrm{}`$.
Let $`𝒞`$ denote the set of all (possibly disconnected) compression bodies $`N`$ such that $`V=_+=V`$. Note that a regular neighborhood of $`V`$ relative to $`\overline{NV}`$ is an element of $`𝒞`$, and hence that $`𝒞\mathrm{}`$. Let us fix an element $``$ of $`𝒞`$ such that (in the notation of 3.1) we have $`\kappa (_{})\kappa (_{}^{})`$ for every $`^{}𝒞`$.
Note that $`V`$ is connected since $`𝒞`$.
Consider first the case in which $`_{}=\mathrm{}`$. In this case, it follows from 3.5 that every component of $``$ is a handlebody, and we have $`=_+=V`$. Since $`N`$ is connected and $`V\mathrm{}`$, we must have $`=\overline{NV}`$. In particular $`N`$ must be closed. Thus conclusion (3) of the proposition holds in this case.
Now consider the case in which some component $`S`$ of $`_{}`$ is a $`2`$-sphere. By irreducibility, $`S`$ bounds a ball $`BN`$. Since $`V`$ is connected, we have either $`VB`$ or $`B(V)=B`$. If $`VB`$, then in particular conclusion (1) of the proposition holds. If $`B(V)=B`$, then $`^{}B`$ is obtained from $``$ by attaching a $`3`$-handle to $`_{}`$, and hence $`^{}𝒞`$ (cf. 3.3). But we have $`_{}^{}=_{}S`$, and it follows from Definition 3.1 that $`\kappa (^{})=\kappa ()1`$. This contradicts the minimality of $`\kappa ()`$.
There remains the case in which $`_{}\mathrm{}`$, and every component of $`_{}`$ has positive genus. Let us fix a component $`Z`$ of $`\overline{NV}`$ which contains at least one component of $`_{}`$. Let us set $`F=Z_{}`$. Then $`F`$ is a non-empty (and possibly disconnected) closed surface in $`intZ`$, and each component of $`F`$ has positive genus. We claim that $`F`$ is incompressible in $`Z`$.
Suppose to the contrary that $`F`$ is compressible in $`Z`$. Then there is a disk $`DintZ`$ such that $`DF=D`$, and such that $`D`$ is a homotopically non-trivial simple closed curve in $`F`$. Since $`DintZNV`$, we have $`D_+=\mathrm{}`$. Furthermore, since $`DZ`$, we have $`D_{}=D(Z_{})=DF=D`$. Hence $`D=D`$. It follows that either $`D`$ or $`D=D`$.
If $`D`$, let $`H_0`$ denote the component of $``$ containing $`D`$, and let $`F_0_{}H_0`$ denote the component of $`F`$ containing $`D`$. Since $`D`$ is homotopically non-trivial it follows that the inclusion homomorphism $`\pi _1(F_0)\pi _1(H_0)`$ has non-trivial kernel. This contradicts 3.6.
If $`D=D`$, we fix a regular neighborhood $`E`$ of $`D`$ relative to $`\overline{Z}`$. Then $`^{}E`$ is obtained from $``$ by attaching a $`2`$-handle to $`_{}`$, and hence $`^{}𝒞`$ (cf. 3.3). The surface $`^{}`$ has the form $`(()A)D_1D_2`$, where $`A`$ is a homotopically non-trivial annulus, and $`D_1`$ and $`D_2`$ are disjoint disks in $`N`$ such that $`(D_1D_2)=A`$. It therefore follows from Lemma 3.2 that $`\kappa ()<\kappa (^{})`$. This contradicts the minimality of $`\kappa ()`$, and the incompressibility of $`F`$ in $`Z`$ is proved.
Since $`Z`$ is $`\pi _1`$-injective in $`N`$ by hypothesis, it now follows that $`F`$ is incompressible in $`N`$. Our choice of $`Z`$ guarantees that $`F\mathrm{}`$. Choose any component $`F_1`$ of $`F`$, and let $`H_1`$ denote the component of $``$ containing $`F_1`$. By 3.4, $`genus(F_1)`$ is at most the genus of the connected surface $`_+`$. But $`_+`$ is a component of $`V`$ since $`𝒞`$, and so $`genus(F_1)`$ is at most the maximum of the genera of the components of $`V`$. Hence conclusion (2) of the proposition holds in this case. ∎
## 4. Transporting surfaces downstairs
###### Lemma 4.1.
Let $`M`$ be a simple, compact, orientable $`3`$-manifold, let $`p:\stackrel{~}{M}M`$ be a $`2`$-sheeted covering, and let $`\tau :\stackrel{~}{M}\stackrel{~}{M}`$ denote the non-trivial deck transformation. Suppose that $`\stackrel{~}{M}`$ contains a closed, incompressible surface $`F_0`$ of positive genus. Then $`F_0`$ is ambiently isotopic to a surface $`F`$ such that $`F`$ and $`\tau (F)`$ meet transversally, and every component of $`F\tau (F)`$ is a homotopically non-trivial simple closed curve in $`M`$.
###### Proof.
Let $``$ denote the collection of all surfaces $`SM`$ such that $`S`$ is isotopic to $`F_0`$ and $`S`$ meets $`\tau (S)`$ transversely. Choose a surface $`F`$ so that the number of components of $`F\tau (F)`$ is minimal. We will show that every component of $`F\tau (F)`$ is a homotopically non-trivial simple closed curve.
Suppose there exists a homotopically trivial component $`\gamma `$ of $`F\tau (F)`$. Then, since $`F`$ is incompressible in $`M`$, the simple closed curve $`\gamma `$ must bound disks $`DF`$ and $`D^{}\tau (F)`$. We assume, without loss of generality, that the disk $`D^{}`$ is innermost on $`\tau (F)`$ in the sense that $`D^{}F=\gamma `$. This implies, in particular that $`DD^{}`$ is an embedded $`2`$-sphere in $`M`$.
Since $`M`$ is irreducible, the $`2`$-sphere $`DD^{}`$ bounds a ball $`B`$ in $`M`$. We may observe at this point that the curve $`\gamma `$ cannot be invariant under $`\tau `$. Otherwise, since $`D^{}`$ is the unique disk on $`\tau (F)`$ bounded by $`\gamma `$, it would follow that $`\tau (D)=D^{}`$, and hence that the sphere $`DD^{}`$ is invariant under $`\tau `$. Since $`\stackrel{~}{M}`$ contains an incompressible surface, it is not homeomorphic to $`S^3`$, and therefore $`B`$ is the unique $`3`$-ball bounded by $`DD^{}`$. Thus the assumption that $`\gamma `$ is invariant implies that the ball $`B`$ is invariant under the fixed point free map $`\tau `$, contradicting the Brouwer Fixed Point Theorem. This shows that $`\gamma `$ is not invariant under $`\tau `$. It follows, since $`D^{}`$ is innermost, that $`D^{}`$ is disjoint from its image under $`\tau `$.
Now let $`V`$ be a regular neighborhood of $`B`$, chosen so that $`VF`$ is a regular neighborhood of $`D`$ and $`VF^{}`$ is a regular neighborhood of $`D^{}`$. The disk $`F^{}V`$ divides $`V`$ into two balls, one of which, say $`U`$, is disjoint from the interior of $`D`$. Since $`D^{}\tau (D^{})=\mathrm{}`$, we may assume without loss of generality that $`V`$ has been chosen to be small enough so that $`U\tau (U)=\mathrm{}`$. Let $`E`$ denote the disk in $`U`$ which is bounded by $`FU`$ and which is disjoint from $`\tau (F)`$. We set $`A=\overline{FU}`$ and consider the surface $`F^{}=AE`$, which is clearly isotopic to $`F`$ by an isotopy supported in $`V`$. We will show that $`F^{}\tau (F^{})(F\tau (F))\gamma `$.
We write $`F^{}\tau (F^{})=(AE)(\tau (A)\tau (F))`$ as the union of the four sets $`A\tau (A)`$, $`A\tau (E)`$, $`E\tau (A)`$ and $`E\tau (E)`$. We have $`A\tau (A)F\tau (F)\gamma `$. Since $`EU`$ and $`U\tau (U)=\mathrm{}`$ we have $`E\tau (E)=\mathrm{}`$. The sets $`E`$ and $`\tau (F)\tau (A)`$ are disjoint by construction, and hence $`E\tau (A)=\mathrm{}`$. Finally, $`A\tau (E)=\tau (E\tau (A))=\mathrm{}`$.
We have shown that $`F^{}\tau (F^{})(F\tau (F))\gamma `$, and hence that $`F^{}\tau (F^{})`$ has fewer components that $`F\tau (F)`$. This contradicts the choice of $`F`$, and completes the proof of the lemma. ∎
###### Lemma 4.2.
Let $`N`$ be a simple, compact, orientable $`3`$-manifold, let $`p:\stackrel{~}{N}N`$ be a $`2`$-sheeted covering, and let $`\tau :\stackrel{~}{N}\stackrel{~}{N}`$ denote the non-trivial deck transformation. Suppose that $`F\stackrel{~}{N}`$ is a closed, incompressible surface such that $`F`$ and $`\tau (F)`$ meet transversally, and every component of $`F\tau (F)`$ is a homotopically non-trivial simple closed curve in $`N`$. Then $`Np(F)`$ is $`\pi _1`$-injective in $`N`$.
###### Proof.
Set $`F_1=\tau (F)`$, so that $`F_1`$ is incompressible in $`\stackrel{~}{N}`$. Set $`C=FF_1`$. Let $`\stackrel{~}{N}^{}`$ denote the $`3`$-manifold obtained by splitting $`\stackrel{~}{N}`$ along $`F`$, and let $`F_1^{}`$ denote the surface obtained by splitting $`F_1`$ along $`C`$. Then $`\stackrel{~}{N}`$ and $`F_1`$ may be regarded as quotient spaces of $`\stackrel{~}{N}^{}`$ and $`F_1^{}`$, and $`F_1^{}`$ is naturally identified with a properly embedded surface in $`\stackrel{~}{N}^{}`$. We have a commutative diagram
where the horizontal maps are quotient maps and the vertical maps are inclusions. The inclusion $`F_1\stackrel{~}{N}`$ is $`\pi _1`$-injective because $`F_1`$ is incompressible in $`\stackrel{~}{N}`$, and the quotient map $`F_1^{}F_1`$ is $`\pi _1`$-injective because the components of $`C`$ are homotopically non-trivial. By commutativity of the diagram it follows that the inclusion $`F_1^{}\stackrel{~}{N}^{}`$ is $`\pi _1`$-injective.
Now let $`\stackrel{~}{N}^{\prime \prime }`$ denote the $`3`$-manifold obtained by splitting $`\stackrel{~}{N}^{}`$ along $`F_1^{}`$. Since the inclusion $`F_1^{}\stackrel{~}{N}^{}`$ is $`\pi _1`$-injective, the quotient map $`\stackrel{~}{N}^{\prime \prime }\stackrel{~}{N}^{}`$ is also $`\pi _1`$-injective. On the other hand, the quotient map $`\stackrel{~}{N}^{}\stackrel{~}{N}`$ is $`\pi _1`$-injective because $`F`$ is incompressible in $`\stackrel{~}{N}`$. Hence the composite quotient map $`\stackrel{~}{N}^{\prime \prime }\stackrel{~}{N}`$ is $`\pi _1`$-injective. It follows that the inclusion map $`\stackrel{~}{N}(FF_1)\stackrel{~}{N}`$ is $`\pi _1`$-injective.
Now consider any component $`Z`$ of $`Np(F)`$. Choose a component $`\stackrel{~}{Z}`$ of $`p^1(Z)`$. Then $`\stackrel{~}{Z}`$ is a component of $`\stackrel{~}{N}(FF_1)`$, and hence the inclusion $`\stackrel{~}{Z}\stackrel{~}{N}`$ is $`\pi _1`$-injective. Thus in the commutative diagram
the inclusion homomorphism $`\pi _1(\stackrel{~}{Z})\pi _1(\stackrel{~}{N})`$ is injective, while the vertical homomorphisms are induced by covering maps and are therefore also injective. Since the image of $`\pi _1(\stackrel{~}{Z})`$ has index at most $`2`$ in $`\pi _1(Z)`$, the kernel of the inclusion homomorphism $`\pi _1(Z)\pi _1(N)`$ has order at most $`2`$. But $`\pi _1(Z)`$ is torsion-free because $`N`$ is simple. Hence $`\pi _1(Z)\pi _1(N)`$ is injective, as asserted by the Lemma. ∎
###### Lemma 4.3.
Suppose that $`N`$ is a simple, compact, orientable $`3`$-manifold, that $`p:\stackrel{~}{N}N`$ is a $`2`$-sheeted covering, that $`g2`$ is an integer, and that $`\stackrel{~}{N}`$ contains a closed, incompressible surface of genus $`g`$. Then there exist a connected book of $`I`$-bundles $`𝒱`$ with $`V=|𝒱|N`$, and a closed, orientable (possibly disconnected) surface $`SintV`$ such that
1. $`\overline{\chi }(V)=\overline{\chi }(S)=2g2`$;
2. every page of $`𝒱`$ has strictly negative Euler characteristic;
3. $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`N`$;
4. $`NV`$ is $`\pi _1`$-injective in $`N`$;
5. no component of $`S`$ is a sphere; and
6. for every page $`P`$ of $`𝒱`$, the set $`SP`$ is a section of the $`I`$-bundle $`P`$.
###### Proof.
According to Lemma 4.1, $`\stackrel{~}{N}`$ contains a closed, incompressible surface $`F`$ of genus $`g`$ such that $`F`$ and $`\tau (F)`$ meet transversally, and every component of $`F\tau (F)`$ is a homotopically non-trivial simple closed curve in $`N`$. It follows that $`q=p|F:FN`$ is an immersion with at most double-curve singularities.
The map $`q_{\mathrm{}}:\pi _1(F)\pi _1(N)`$ is injective because $`F`$ is incompressible in $`\stackrel{~}{N}`$ and $`p:\stackrel{~}{N}N`$ is a covering map.
Let us set $`X=q(F)`$, and let $`CX`$ denote the union of all double curves of $`q`$. Since the components of $`C`$ are homotopically non-trivial in $`N`$ and hence in $`X`$, the set $`\stackrel{~}{C}=q^1(C)`$ is a disjoint union of homotopically non-trivial simple closed curves in $`F`$. Hence $`F\stackrel{~}{C}`$ is $`\pi _1`$-injective in $`F`$, and each of its components has non-positive Euler characteristic. Since $`q_{\mathrm{}}:\pi _1(F)\pi _1(N)`$ is injective it follows that $`q|(F\stackrel{~}{C}):(F\stackrel{~}{C})N`$ is $`\pi _1`$-injective.
The set $`F\stackrel{~}{C}`$ is mapped homeomorphically onto $`XC`$ by $`q`$. In particular, each component of $`XC`$ has non-positive Euler characteristic. Furthermore, since $`q|(F\stackrel{~}{C}):(F\stackrel{~}{C})N`$ is $`\pi _1`$-injective, it now follows that $`XC`$ is $`\pi _1`$-injective in $`N`$.
In the notation of (1.7), we have $`\mathrm{\Pi }(X)=XC`$ , and the link in $`X`$ of every point of $`C`$ is homeomorphic to the suspension of a four-point set. Since every component of $`XC`$ has non-positive Euler characteristic, it follows from Definition 2.6 that $`X`$ is a book of surfaces. Since each component of $`C`$ is a simple closed curve, we have $`\overline{\chi }(X)=\overline{\chi }(F)=2g2`$.
Let $`W`$ denote a regular neighborhood of $`X`$ in $`N`$. According to Lemma 2.9, we may write $`W=|𝒲|`$ for some book of $`I`$-bundles $`𝒲`$ in such a way that Conclusions (2)–(4) of Lemma 2.9 hold. Since $`XC`$ is $`\pi _1`$-injective in $`N`$, it follows from Conclusion (4) of Lemma 2.9 that $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`N`$.
Since $`\chi (W)=\chi (X)=22g<0`$, and since $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`N`$, it follows from Lemma 2.5 that there is a connected book of $`I`$-bundles $`𝒱`$ with $`V=|𝒱|N`$, such that Conclusions (1)—(7) of Lemma 2.5 hold. Conclusion (2) of Lemma 2.5 gives $`\overline{\chi }(V)=\overline{\chi }(W)=\overline{\chi }(X)`$, so that
(4.3.1)
$$\overline{\chi }(V)=2g2.$$
It follows from Conclusions (1) and (6) of Lemma 2.5 that every binding of $`𝒲`$ is contained in a binding of $`𝒱`$. Since by Conclusion (2) of Lemma 2.9 we have $`Cint_𝒲`$, it follows that $`Cint_𝒱`$.
Let $`𝒰`$ denote a regular neighborhood of $`C`$ in $`int_𝒱`$. We may suppose $`𝒰`$ to be chosen so that $`𝒰`$ meets $`\mathrm{\Pi }(X)`$ transversally, and each component of $`𝒰X`$ is homeomorphic to $`\text{+}\times S^1`$, where + denotes a cone on a four-point set. Set $`X^{}=\overline{X(𝒰X)}`$ and $`F^{}=Fq^1(X^{})`$. Then $`F^{}`$ and $`X^{}`$ are (possibly disconnected) compact $`2`$-manifolds with boundary, and $`q^{}=q|F^{}`$ maps $`F^{}`$ homeomorphically onto $`X^{}`$. Let us fix an orientation of $`F`$, so that $`F^{}`$ inherits an orientation, and define an orientation of $`X^{}`$ by transporting the orientation of $`F^{}`$ via $`q`$.
Let $`U_1,\mathrm{},U_m`$ denote the components of $`𝒰`$. We set $`B_i=XU_i`$. Each component $`\beta `$ of $`B_i`$ is a boundary component of $`X^{}`$ and hence has an orientation induced from the orientation of $`X^{}`$, which determines a generator of $`H_1(U_i;)`$ via the inclusion isomorphism $`H_1(\beta ;)H_1(U_i;)`$. We shall say that two components of $`B_i`$ are similar if they determine the same generator of $`H_1(U_i;)`$ via this construction.
The set $`(U_i)B_i`$ has four components. Their closures are annuli, which we shall call complementary annuli. We shall say that two components of $`B_i`$ are adjacent if their union is the boundary of a complementary annulus, and opposite otherwise.
If $`\beta `$ and $`\beta ^{}`$ are opposite components of $`XU_i`$, then $`q^1(\beta )`$ and $`q^1(\beta ^{})`$ form the boundary of an annulus $`A`$ in $`F`$, which is mapped homeomorphically by $`q`$ to an embedded annulus in $`U_i`$. Since the orientation of $`F^{}`$ is the restriction of an orientation of $`F`$, the induced orientations of $`q^1(\beta )`$ and $`q^1(\beta ^{})`$ determine different generators of $`H_1(A;)`$. In view of our definitions it follows that opposite components of $`B_i`$ are dissimilar.
Let us call a complementary annulus bad if its boundary curves are similar, and good otherwise. If $`\beta `$ is any component of $`B_i`$, the two components of $`B_i`$ adjacent to $`\beta `$ are opposite each other; hence exactly one of them is similar to $`\beta `$. This shows that $`\beta `$ is contained in the boundary of exactly one bad annulus and one good annulus. We conclude that $`U_i`$ contains exactly two good annuli, say $`A_i`$ and $`A_i^{}`$, and that $`A_iA_i^{}=\mathrm{}`$.
The set
$$S=(X(X𝒰))(A_1\mathrm{}A_m)(A_1^{}\mathrm{}A_m^{})$$
is a (possibly disconnected) compact PL $`2`$-manifold embedded in $`V`$. Since $`A_i`$ and $`A_i^{}`$ are good annuli, the orientation of $`X^{}`$ extends to an orientation of $`S`$. In particular $`S`$ is orientable.
We shall show that Conclusions (1)–(6) of the present lemma hold when $`𝒱`$ and $`S`$ are defined as above.
According to (4.3.1) we have $`\overline{\chi }(V)=2g2`$. It follows from the construction of $`S`$ that $`\overline{\chi }(S)=\overline{\chi }(X)=2g2`$. Hence Conclusion (1) of the present lemma holds.
Conclusion (2) of the present lemma follows from Conclusion (3) of Lemma 2.5.
Since we have seen that $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`N`$, it follows from Conclusion (6) of Lemma 2.5 that $`𝒫_𝒱`$ is $`\pi _1`$-injective in $`N`$. This is Conclusion (3) of the present lemma.
It follows from Lemma 4.2 that $`NX=Nq(F)`$ is $`\pi _1`$-injective in $`N`$. It follows from conclusions (1) and (4) of Lemma 2.5 that every component of $`NV`$ is also a component of $`NW`$, and is therefore ambiently isotopic in $`N`$ to a component of $`NX`$. Hence $`NV`$ is $`\pi _1`$-injective in $`N`$. This is Conclusion (4) of the present lemma.
It follows from the construction of $`S`$ that $`S𝒫_𝒱=X𝒫_𝒱`$. If $`P`$ is any page of $`𝒱`$, then by Conclusion (6) of Lemma 2.5, $`P`$ is a page of $`𝒲`$, and hence $`SP=XP`$ is a section of the $`I`$-bundle $`P`$ according to Conclusion (3) of Lemma 2.9. This establishes Conclusion (6) of the present lemma.
In particular it follows that for every page $`P`$ of $`𝒱`$ the surface $`PS`$ is connected and has non-positive Euler characteristic. On the other hand, the construction of $`S`$ shows that every component of $`S_𝒱`$ is an annulus. Hence every component of $`S`$ has non-positive Euler characteristic, and Conclusion (5) of the present lemma follows. ∎
###### Proposition 4.4.
Suppose that $`N`$ is a simple, compact, orientable $`3`$-manifold, that $`p:\stackrel{~}{N}N`$ is a $`2`$-sheeted covering, that $`g2`$ is an integer, and that $`\stackrel{~}{N}`$ contains a closed, incompressible surface of genus $`g`$. Then either
1. $`N`$ contains a closed, connected, incompressible surface of genus at most $`g`$, or
2. $`N`$ is closed and there is a connected, book of $`I`$-bundles $`𝒱`$ with $`V=|𝒱|N`$ such that $`\overline{\chi }(V)=2g2`$, every page of $`𝒱`$ has strictly negative Euler characteristic, and every component of $`\overline{NV}`$ is a handlebody. In particular, the rank of $`H_1(N;_2)`$ is at most $`4g3`$. Furthermore, there is a closed, orientable (possibly disconnected) surface $`SintV`$ such that for every page $`P`$ of $`𝒱`$, the set $`SP`$ is a section of the $`I`$-bundle $`P`$.
The last sentence of alternative of (2) is not used in this paper, but will be needed in .
###### Proof of Proposition 4.4.
Let us fix a connected book of $`I`$-bundles $`𝒱`$ with $`V=|𝒱|N`$, and a closed, orientable surface $`SintV`$, such that Conclusions (1) to (6) of Lemma 4.3 hold. We distinguish two cases, depending on whether there does or does not exist a page of $`𝒱`$ whose horizontal boundary is contained in a single component of $`V`$.
Case I. There is a page $`P_0`$ of $`𝒱`$ such that $`_hP_0`$ is contained in a single component $`Y_0`$ of $`V`$.
According to conclusion (6) of Lemma 4.3, the set $`SP_0`$ is a section of the $`I`$-bundle $`P_0`$. Hence there is a properly embedded arc $`\alpha `$ in $`V`$, such that $`\alpha P_0`$, and such that $`\alpha `$ meets $`S`$ transversally in a single point. The endpoints of $`\alpha `$ lie in $`_hP_0Y_0`$. Since $`Y_0`$ is connected, there is an arc $`\beta Y_0`$ with $`\beta =\alpha `$.
Let $`\sigma `$ denote the class in $`H_2(N;_2)`$ represented by $`S`$. Since $`\alpha `$ is properly embedded in $`V`$ and meets $`X`$ transversally in a single point of $`\pi _0\mathrm{\Pi }(X)`$, the class $`\sigma `$ has intersection number $`1`$ with the class in $`H_1(N;_2)`$ represented by the simple closed curve $`\alpha \beta `$. In particular $`\sigma 0`$. Hence some component $`S_0`$ of $`S`$ represents a non-zero class in $`H_2(N;_2)`$. It follows from Conclusions (1) and (5) of Lemma 4.3 that $`\overline{\chi }(S_0)\overline{\chi }(S)=2g2`$, and hence that $`genus(S_0)g`$.
Among all closed, orientable surfaces in $`N`$ that represent non-trivial classes in $`H_2(N;_2)`$, let us choose one, say $`S_1`$, of minimal genus. Then $`genus(S_1)genus(S_0)g`$. If $`S_1`$ is compressible in $`N`$, a compression of $`S_1`$ produces a $`2`$-manifold $`S_2`$ with one or two components. Each component of $`S_2`$ has strictly smaller genus than $`S_1`$, and at least one of them represents a non-trivial class in $`H_2(N;_2)`$. This contradicts minimality. Hence $`S_1`$ is incompressible in $`N`$. Since $`genus(S_1)g`$, conclusion (1) of the present lemma holds in this case.
Case II. There is no page $`P_0`$ of $`𝒱`$ such that $`_hP_0`$ is contained in a single component of $`V`$.
In this case, every page of $`𝒱`$ is a trivial $`I`$-bundle. Furthermore, if $`T`$ is any component of $`V`$, then for every page $`P`$ of $`𝒱`$, at most one component of the horizontal boundary of $`P`$ is contained in $`T`$. Hence
$$\overline{\chi }(TP)\overline{\chi }(P)$$
for every page $`P`$ of $`𝒱`$. Letting $`P`$ range over the pages of $`𝒱`$, and using (4.3.1), we find that
$$\overline{\chi }(T)=\underset{P}{}\overline{\chi }(TP)\underset{P}{}\overline{\chi }(P)=\overline{\chi }(V)=2g2.$$
This shows that
(4.4.1)
$$genus(T)g$$
for every component $`T`$ of $`V`$.
According to Conclusion (4) of Lemma 4.3, $`NV`$ is $`\pi _1`$-injective in $`N`$. Thus $`VN`$ satisfies the hypotheses of Proposition 3.8. There are three subcases corresponding to the three alternatives (1)—(3) of Proposition 3.8.
First suppose that alternative (1) of Proposition 3.8 holds, i.e. that $`V`$ is contained in a ball. Then in particular for any page $`P`$ of $`𝒱`$, the inclusion homomorphism $`\pi _1(P)\pi _1(W)`$ is trivial. But according to Conclusions (2) and (3) of Lemma 4.3, we have $`\chi (P)<0`$ (so that $`\pi _1(P)`$ is non-trivial) and $`𝒫_𝒲`$ is $`\pi _1`$-injective in $`N`$. This contradiction shows that alternative (1) of Proposition 3.8 cannot hold in this situation.
Next suppose that alternative (2) of Proposition 3.8 holds, i.e. that there exists a connected, incompressible closed surface $`S_1`$ in $`N`$ whose genus is at most the maximum of the genera of the components of $`V`$. By (4.4.1) this maximum is at most $`g`$. Thus conclusion (1) of the present lemma holds in this subcase.
Finally, suppose that alternative (3) of Proposition 3.8 holds, i.e. that $`N`$ is closed and that every component of $`\overline{NV}`$ is a handlebody. We have that $`V=|𝒱|`$ where $`V`$ is a book of $`I`$-bundles whose pages all have negative Euler characteristic, and Conclusion (1) of Lemma 4.3 gives $`\overline{\chi }(V)=2g2`$. Since the components of $`\overline{NV}`$ are handlebodies, the inclusion of $`V`$ into $`N`$ induces a surjection from $`H_1(V;_2)`$ to $`H_1(N;_2)`$; hence the latter group has rank at most $`4g3`$ by Lemma 2.11. Furthermore, according to Conclusion (6) of Lemma 4.3, for every page $`P`$ of $`𝒱`$, the set $`SP`$ is a section of the $`I`$-bundle $`P`$. Thus conclusion (2) of the present proposition holds in this subcase. ∎
## 5. Singularity of PL maps
If $`K`$ is a finite simplicial complex, we shall denote the underlying space of $`K`$ by $`|K|`$. A simplicial map $`\varphi :K_1K_2`$ between finite simplicial complexes defines a map from $`|K_1|`$ to $`|K_2|`$ which we shall denote by $`|\varphi |`$.
Now suppose that $`X_1`$ and $`X_2`$ are compact topological spaces and that $`f:X_1X_2`$ is a continuous surjection. We define a triangulation of $`f`$ to be a quintuple $`(K_1,J_1,K_2,J_2,\varphi )`$, where each $`K_i`$ is a finite simpicial complex, $`J_i:|K_i|X_i`$ is a homeomorphism, and $`fJ_1=J_2\varphi `$. When it is unnecessary to specify the $`K_i`$ and $`J_i`$ we shall simply say that $`\varphi `$ is a triangulation of $`f`$.
Note that if $`f`$ is any PL map from a compact PL space $`X`$ to a PL space $`Y`$, then the surjection $`f:Xf(X)`$ admits a triangulation.
###### Definition 5.1.
Let $`K`$ and $`L`$ be finite simplicial complexes and let $`\varphi :KL`$ be a simplicial map. We define the degree of singularity of $`\varphi `$, denoted $`\mathrm{DS}(\varphi )`$, to be the number of ordered pairs $`(v,w)`$ of vertices of $`K`$ such that $`vw`$ but $`\varphi (v)=\varphi (w)`$.
If $`f`$ is any PL map from a compact PL space $`X`$ to a PL space $`Y`$, we define the absolute degree of singularity of $`f`$, denoted $`\mathrm{ADS}(f)`$, by
$$\mathrm{ADS}(f)=\underset{\varphi }{\mathrm{min}}\mathrm{DS}(\varphi ),$$
where $`\varphi `$ ranges over all triangulations of $`f:Xf(X)`$.
###### 5.2.
We emphasize that the definition of $`\mathrm{ADS}(f)`$ is based on regarding $`f`$ as a map from $`X`$ to $`f(X)`$. Hence if $`f`$ is any PL map from a compact PL space $`X`$ to a PL space $`Y`$, and $`Z`$ is a PL subspace of $`Y`$ containing $`f(X)`$, then the absolute degree of singularity of $`f`$ is unchanged when we regard $`f`$ as a PL map from $`X`$ to $`Z`$.
An almost equally trivial immediate consequence of the definition of absolute degree of singularity is expressed by the following result.
###### Lemma 5.3.
Suppose that $`X`$, $`Y`$ and $`Z`$ are PL spaces, that $`X`$ is compact, that $`f:XY`$ is a PL map, and that $`h`$ is a PL homeomorphism of $`f(X)`$ onto a PL subspace of $`Z`$. Then $`hf:XZ`$ has the same absolute degree of singularity as $`f`$.
###### Proof.
In view of 5.2 we may assume that $`f`$ is surjective and that $`h`$ is a PL homeomorphism of $`Y`$ onto $`Z`$. Now if $`(K_1,J_1,K_2,J_2,\varphi )`$ is a triangulation of $`f`$ then $`(K_1,J_1,K_2,hJ_2,h\varphi )`$ is a triangulation of $`hf`$, and $`\mathrm{DS}(h\varphi )=\mathrm{DS}(\varphi )`$. It follows that $`\mathrm{ADS}(hf)\mathrm{ADS}(f)`$. The same argument, with $`h^1`$ in place of $`h`$, shows that $`\mathrm{ADS}(f)\mathrm{ADS}(hf)`$. ∎
###### Proposition 5.4 (Stallings).
Suppose that $`Y`$ is a connected PL space and that $`p:\stackrel{~}{Y}Y`$ is a connected covering space, which is non-trivial in the sense that $`p`$ is not a homeomorphism. Suppose that $`f`$ is a PL map from a compact connected PL space $`X`$ to $`Y`$, such that the inclusion homomorphism $`\pi _1(f(X))\pi _1(Y)`$ is surjective. Suppose that $`\stackrel{~}{f}:X\stackrel{~}{Y}`$ is a lift of $`f`$. Then $`\mathrm{ADS}(\stackrel{~}{f})<\mathrm{ADS}(f)`$.
###### Proof.
We first prove the proposition in the special case where $`f:XY`$ is a surjection. In this case we set $`m=\mathrm{ADS}(f)`$, and we fix a triangulation $`(K_1,J_1,K_2,J_2,\varphi )`$ of the PL surjection $`f`$ such that $`\mathrm{DS}(\varphi )=m`$. Here, by definition, $`J_1:|K_1|X`$ and $`J_2:|K_2|Y`$ are homeomorphisms. Let us identify $`X`$ and $`Z`$ with $`|K_1|`$ and $`|K_2|`$ via these homeomorphisms. The covering space $`\stackrel{~}{Y}`$ of $`Y`$ may be identified with $`|\stackrel{~}{K}_2|`$ for some simplicial covering complex $`\stackrel{~}{K}_2`$ of $`K_2`$; thus $`p=|\sigma |`$ for some simplicial covering map $`\sigma :\stackrel{~}{K}_2K_2`$. The lift $`\stackrel{~}{f}`$ may be written as $`|\stackrel{~}{\varphi }|`$ for some simplicial lift $`\stackrel{~}{\varphi }:K_1\stackrel{~}{K}_2`$. We shall denote by $`W`$ the subcomplex $`\stackrel{~}{\varphi }(K_1)`$ of $`\stackrel{~}{K}_2`$.
Since $`\sigma \stackrel{~}{\varphi }=\varphi `$, the definition of degree of singularity implies that $`\mathrm{DS}(\stackrel{~}{\varphi })\mathrm{DS}(\varphi )=m`$. If equality holds here, then the restriction of $`\sigma `$ to the vertex set of $`W`$ is one-to-one. This implies that $`p`$ restricts to a one-to-one map from $`|W|`$ to $`Y`$. But we have $`W=\stackrel{~}{f}(X)`$, and the surjectivity of $`f`$ implies that $`p`$ maps $`|W|`$ onto $`Y`$; thus $`p`$ restricts to a homeomorphism from $`|W|`$ to $`Y`$. This is impossible since $`p:\stackrel{~}{Y}Y`$ is a non-trivial connected covering space. Hence we must have $`\mathrm{DS}(\stackrel{~}{\varphi })<m`$. Since by definition we have $`\mathrm{ADS}(\stackrel{~}{\varphi })DS(\stackrel{~}{\varphi })`$, the assertion of the proposition follows in the case where $`f`$ is surjective.
We now turn to the general case. Let us set $`Z=f(X)`$ and $`\stackrel{~}{Z}=p^1(Z)`$. Since $`\stackrel{~}{Y}`$ is a non-trivial connected covering space of $`Y`$, and since the inclusion homomorphism $`\pi _1(Z)\pi _1(Y)`$ is surjective, $`\stackrel{~}{Z}`$ is a non-trivial connected covering space of $`Z`$. According to 5.2, regarding $`\stackrel{~}{f}`$ and $`f`$ as maps into $`\stackrel{~}{Z}`$ and $`Z`$ does not affect their absolute degrees of singularity. Since $`f:XZ`$ is surjective, the inequality now follows from the special case that has already been proved. ∎
Following the terminology used by Simon in , we shall say that a $`3`$-manifold $`M`$ admits a manifold compactification if there is a homeomorphism $`h`$ of $`M`$ onto an open subset of a compact $`3`$-manifold $`Q`$ such that $`h(intM)=intQ`$.
###### Lemma 5.5.
Suppose that $`N`$ is a compact, orientable, connected, irreducible PL $`3`$-manifold and that $`D`$ is a separating, properly embedded disk in $`N`$. Let $`X`$ denote the closure of one of the connected components of $`ND`$. Let $`\nu D`$ be a base point, and let $`p:(\stackrel{~}{N},\stackrel{~}{\nu })(N,\nu )`$ denote the based covering space corresponding to the subgroup $`im(\pi _1(X,\nu )\pi _1(N,\nu ))`$ of $`\pi _1(N,\nu )`$. Then $`\stackrel{~}{N}`$ admits a manifold compactification.
###### Proof.
Let us set $`X_1=\overline{NX}`$. It will also be convenient to write $`X_0=X`$. Then the $`X_i`$ are compact submanifolds of $`N`$, and $`X_1X_2=D`$. We set $`H_i=\pi _1(X_i,\nu )`$ for $`i=0,1`$. We identify $`\pi _1(N,\nu )`$ with $`H_0H_1`$, so that the $`H_i`$ are in particular subgroups of $`\pi _1(N,\nu )`$. Thus $`(\stackrel{~}{N},\stackrel{~}{\nu })`$ is the based covering space corresponding to the subgroup $`H_0`$.
According to the general criterion given by Simon in \[20, Theorem 3.1\], $`\stackrel{~}{N}`$ will admit a manifold compactification provided that the following conditions hold:
1. $`X_0`$ and $`X_1`$ are irreducible,
2. $`D`$ is $`\pi _1`$-injective in $`X_0`$ and $`X_1`$,
3. each conjugate of $`H_0`$ in $`\pi _1(N,\nu )`$ intersects $`im(\pi _1(D,\nu )\pi _1(N,\nu ))`$ in a finitely generated subgroup, and
4. for each $`i\{0,1\}`$, and for each finitely generated subgroup $`Z`$ of $`H_i`$ which has the form $`H_ig^1H_0g`$ for some $`g\pi _1(N,\nu )`$, the based covering space of $`(X_i,\nu )`$ corresponding to $`Z`$ admits a manifold compactification.
Here conditions (ii) and (iii) hold trivially because $`\pi _1(D)`$ is trivial. Condition (i) follows from the irreducibility of $`N`$. (A ball bounded by a sphere in $`intX_i`$ must be contained in $`X_i`$ because the frontier of $`X_i`$ is the disk $`D`$, and $`D\mathrm{}`$.)
To prove (iv), we consider any $`i\{0,1\}`$ and any subgroup of $`H_i`$ having the form $`Z=H_ig^1H_0g`$ where $`g\pi _1(N,\nu )`$. Since $`\pi _1(N,\nu )=H_0H_1`$, we have either (a) $`Z=\{1\}`$ or (b) $`i=0`$ and $`gH_0`$. If (a) holds then the based covering of $`(X_i,\nu )`$ corresponding to $`Z`$ is equivalent to the universal cover of $`X_i`$. But since $`X_i`$ is irreducible and has a non-empty boundary, it is a Haken manifold. Hence by \[23, Theorem 8.1\], the universal cover of $`X_i`$ admits a manifold compactification. If (b) holds then the covering corresponding to $`Z`$ is homeomorphic to $`X_i`$ and is therefore a manifold compactification of itself. ∎
###### Lemma 5.6.
Suppose that $`N`$ is a compact, connected, orientable, irreducible PL $`3`$-manifold and that $`D`$ is a separating, properly embedded disk in $`N`$. Let $`X`$ denote the closure of one of the connected components of $`ND`$. Let $`\nu D`$ be a base point, and let $`p:(\stackrel{~}{N},\stackrel{~}{\nu })(N,\nu )`$ denote the based covering space corresponding to the subgroup $`im(\pi _1(X,\nu )\pi _1(N,\nu ))`$ of $`\pi _1(N,\nu )`$. Let $`\stackrel{~}{X}`$ denote the component of $`p^1(X)`$ containing $`\stackrel{~}{\nu }`$ (so that $`p`$ maps $`\stackrel{~}{X}`$ homeomorphically onto $`X`$). Then every compact PL subset of $`int\stackrel{~}{N}`$ is PL ambient-isotopic to a subset of $`\stackrel{~}{X}`$.
###### Proof.
Since $`N`$ is a compact, orientable, irreducible $`3`$-manifold with non-empty boundary, it is a Haken manifold. Hence by \[23, Theorem 8.1\], the universal cover of $`intN`$ is homeomorphic to $`𝐑^3`$. Thus $`int\stackrel{~}{N}`$ is covered by an irreducible manifold and is therefore irreducible.
According to Lemma 5.5, the manifold $`\stackrel{~}{N}`$ admits a manifold compactification. Thus there is a homeomorphism $`h`$ of $`\stackrel{~}{N}`$ onto an open subset of a compact $`3`$-manifold $`Q`$ such that $`h(int\stackrel{~}{N})=intQ`$. Since $`intQ`$ is homeomorphic to the irreducible manifold $`int\stackrel{~}{N}`$, the compact manifold $`Q`$ is itself irreducible.
The definition of $`\stackrel{~}{N}`$ implies that the inclusion map $`\iota :XN`$ admits a based lift $`\stackrel{~}{\iota }:(X,\nu )(\stackrel{~}{N},\stackrel{~}{\nu })`$, that $`\stackrel{~}{\iota }(X)=\stackrel{~}{X}`$, and that $`\stackrel{~}{\iota }_{\mathrm{}}:\pi _1(X,\nu )\pi _1(\stackrel{~}{N},\stackrel{~}{\nu })`$ is an isomorphism. Hence the inclusion $`\stackrel{~}{X}\stackrel{~}{N}`$ induces an isomorphism of fundamental groups, and if we set $`X^{}=h(\stackrel{~}{X})`$, the inclusion $`X^{}Q`$ induces an isomorphism of fundamental groups.
On the other hand, since the frontier of $`X`$ in $`N`$ is $`D`$, the frontier of $`X^{}`$ in $`Q`$ is $`D^{}=h(\stackrel{~}{\iota }(D))`$, a properly embedded disk in the compact $`3`$-manifold $`Q`$. Set $`Y=\overline{QX^{}}`$. Then in terms of a base point in $`D^{}`$ we have a canonical identification of $`\pi _1(Q)`$ with $`\pi _1(X^{})\pi _1(Y)`$. Since the inclusion $`X^{}Q`$ induces an isomorphism of fundamental groups, it follows that $`\pi _1(Y)`$ is trivial. We also know that $`Y`$ is irreducible because its frontier in the irreducible manifold $`Q`$ is a disk. Thus $`Y`$ is a compact, simply connected, irreducible $`3`$-manifold with non-empty boundary, and is therefore PL homeomorphic to a ball.
We have now exhibited $`Q`$ as the union of the compact $`3`$-dimensional submanifold $`X`$ and the PL $`3`$-ball $`Y`$, and their intersection is the disk $`D`$. It follows that any compact PL subset $`W`$ of $`intQ`$ is PL isotopic to a subset of $`intX`$. Since $`h`$ maps $`int\stackrel{~}{N}`$ homeomorphically onto $`intQ`$, and maps $`int\stackrel{~}{X}`$ homeomorphically onto $`intX^{}`$, the conclusion of the lemma follows. ∎
###### Lemma 5.7.
Suppose that $`K`$ is a compact, connected PL space such that $`\pi _1(K)`$ has rank $`2`$ and is freely indecomposable. Suppose that $`N`$ is a compact, connected, orientable PL $`3`$-manifold which is irreducible but boundary-reducible. Suppose that $`f:KintN`$ is a $`\pi _1`$-injective PL map, and that the inclusion homomorphism $`\pi _1(f(K))\pi _1(N)`$ is surjective. Then $`f`$ is homotopic to a map $`g`$ such that $`\mathrm{ADS}(g)<\mathrm{ADS}(f)`$.
###### Proof.
Since $`N`$ is boundary-reducible it contains an essential properly embedded disk. If $`N`$ contains a non-separating essential disk $`D_0`$, then there is a separating essential disk $`D_1`$ in $`ND_0`$ such that the closure of the component of $`ND_1`$ containing $`D_0`$ is a solid torus $`J`$. In this case $`\pi _1(\overline{NJ})`$ is non-trivial, since $`\pi _1(N)`$ has rank at least $`2`$; hence $`D_1`$ is an essential disk as well. Thus in all cases, $`N`$ contains a separating essential disk $`D`$. We may write $`N=X_0X_1`$ for some connected submanifolds $`X_0`$ and $`X_1`$ of $`N`$ with $`X_0X_1=D`$. We choose a base point in $`\nu D`$ and set $`A_i=im(\pi _1(X_i,\nu )\pi _1(N,\nu ))`$ for $`i=0,1`$. Then $`\pi _1(N,\nu )=A_0A_1`$.
If one of the $`A_i`$ were trivial, then one of the $`X_i`$ would be a ball since $`N`$ is irreducible, and $`D`$ would not be an essential disk. Hence the $`A_i`$ are non-trivial subgroups. It then follows from the free product structure of $`\pi _1(N,\nu )`$ that the $`A_i`$ are of infinite index in $`\pi _1(N,\nu )`$, and in particular that they are proper subgroups.
Since the subgroup $`H=f_{\mathrm{}}(\pi _1(K))`$, which is defined only up to conjugacy in $`\pi _1(N)`$, has rank at least $`2`$ and is freely indecomposable, it follows from the Kurosh subgroup theorem that $`H`$ is conjugate to a subgroup of one of the $`A_i`$. By symmetry we may assume that $`H`$ is conjugate to a subgroup of $`A_0`$. Hence after modifying $`f`$ by a homotopy we may assume that $`f`$ maps some base point $`\kappa `$ of $`K`$ to $`\nu `$ and that $`f_{\mathrm{}}(\pi _1(K,\kappa ))A_0`$. Hence if $`(\stackrel{~}{N},\stackrel{~}{\nu })`$ denotes the based covering space of $`(N,\nu )`$ corresponding to the subgroup $`A_0`$ of $`\pi _1(N)`$, then $`f`$ admits a lift $`\stackrel{~}{f}:(K,\kappa )(\stackrel{~}{N},\stackrel{~}{\nu })`$. Since $`A_0`$ is a proper subgroup of $`\pi _1(N,\nu )`$, the covering space $`\stackrel{~}{N}`$ is non-trivial. Hence, according to Proposition 5.4, we have $`\mathrm{ADS}(\stackrel{~}{f})<\mathrm{ADS}(f)`$.
Let $`\stackrel{~}{X}_0`$ denote the component of $`p^1(X_0)`$ containing $`\stackrel{~}{\nu }`$, so that $`p`$ maps $`\stackrel{~}{X}_0`$ homeomorphically onto $`X_0`$. According to Lemma 5.6, the compact PL subset $`\stackrel{~}{f}(K)`$ of $`int\stackrel{~}{N}`$ is PL ambient-isotopic to a subset of $`\stackrel{~}{X}_0`$. In particular, there is a PL homeomorphism $`j`$ of $`\stackrel{~}{f}(K)`$ onto a subset $`L`$ of $`\stackrel{~}{X}_0`$ such that $`j`$, regarded as a map of $`\stackrel{~}{f}(K)`$ into $`\stackrel{~}{N}`$, is homotopic to the inclusion $`\stackrel{~}{f}(K)\stackrel{~}{N}`$. It now follows that $`pj`$ maps $`\stackrel{~}{f}(K)`$ homeomorphically onto the subset $`p(L)`$ of $`X_0N`$. Hence by Lemma 5.3, if we set $`g=pj\stackrel{~}{f}:KN`$, we have $`\mathrm{ADS}(g)=\mathrm{ADS}(\stackrel{~}{f})<\mathrm{ADS}(f)`$. But since $`j:\stackrel{~}{f}(K)\stackrel{~}{N}`$ is homotopic to the inclusion $`\stackrel{~}{f}(K)\stackrel{~}{N}`$, the map $`g:KN`$ is homotopic to $`f`$. ∎
###### Proposition 5.8.
Suppose that $`K`$ is a compact, connected PL space such that $`\pi _1(K)`$ has rank at least $`2`$ and is freely indecomposable. Suppose that $`f`$ is a $`\pi _1`$-injective. PL map from $`K`$ to the interior of a compact, connected, orientable, irreducible PL $`3`$-manifold $`M`$. Then there exist a map $`g:KM`$ homotopic to $`f`$ with $`\mathrm{ADS}(g)\mathrm{ADS}(f)`$, and a compact, connected $`3`$-dimensional submanifold $`N`$ of $`intM`$ such that (i) $`intNg(K)`$, (ii) the inclusion homomorphism $`\pi _1(g(K))\pi _1(N)`$ is surjective, (iii) $`N`$ is incompressible in $`M`$, and (iv) $`N`$ is irreducible.
###### Proof.
Among all maps from $`K`$ to $`M`$ that are homotopic to $`f`$, we choose one, $`g`$, for which $`\mathrm{ADS}(g)`$ has the smallest possible value. In particular we then have $`\mathrm{ADS}(g)\mathrm{ADS}(f)`$. Note also that $`f_{\mathrm{}}:\pi _1(K)\pi _1(N)`$ is injective.
Now let $`N`$ be a compact, connected $`3`$-submanifold of $`M`$ satisfying conditions (i) and (ii) of the statement of the Proposition, and choose $`N`$ so as to minimize the quantity $`\kappa (N)`$ (see 3.1) among all compact, connected $`3`$-submanifolds satisfying (i) and (ii). We shall complete the proof by showing that $`N`$ satisfies (iii) and (iv).
We first show that (iv) holds, i.e. that $`N`$ is irreducible. If $`SintN`$ is a $`2`$-sphere, then $`S`$ bounds a ball $`BM`$. If we set $`N^{}=NB`$, then the pair $`N^{}`$ satisfies (i) and (ii). (It inherits property (ii) from $`N`$ because the inclusion homomorphism $`\pi _1(N)\pi _1(N^{})`$ is surjective.) But if $`BN`$, it is clear from Definition 3.1 that $`\kappa (N^{})<\kappa (N)`$, and the minimality of $`\kappa (N)`$ is contradicted. Hence we must have $`BN`$, and irreducibility is proved.
It remains to show that (iii) holds, i.e. that $`N`$ is incompressible. If this is false, then either $`N`$ has a sphere component, or there is a compressing disk $`D`$ for $`N`$. If $`N`$ has a sphere component $`S`$, then the irreducibility of $`N`$ implies that $`N`$ is a ball. But then the injectivity of $`g_{\mathrm{}}:\pi _1(K)\pi _1(N)`$ implies that $`\pi _1(K)`$ is trivial, a contradiction to the hypothesis that $`\pi _1(K)`$ has rank at least $`2`$.
If there is a compressing disk $`D`$ for $`N`$, then either $`DN=D`$ or $`DN`$. If $`DN=D`$, and if we set $`N^{}=NQ`$, where $`Q`$ is a regular neighborhood of $`D`$ relative to $`\overline{MN}`$, then the $`3`$-submanifold $`N`$ satisfies conditions (i) and (ii). (It inherits property (ii) because the inclusion homomorphism $`\pi _1(N)\pi _1(N^{})`$ is again surjective.) Now $`N^{}`$ has the form $`((N)A)D_1D_2`$, where $`AN`$ is a homotopically non-trivial annulus, and $`D_1`$ and $`D_2`$ are disjoint disks in $`M`$ such that $`(D_1D_2)N=A`$. It therefore follows from Lemma 3.2 that $`\kappa (N)<\kappa (N^{})`$. Again the minimality of $`\kappa (N)`$ is contradicted.
Finally, if $`DN`$, then $`N`$ is boundary-reducible. As we have already shown that $`N`$ is irreducible, it follows from Lemma 5.7 that $`f`$ is homotopic in $`N`$ to a map $`g^{}`$ such that $`\mathrm{ADS}(g^{})<\mathrm{ADS}(g)`$. In particular, $`g^{}`$ is homotopic to $`g`$ in $`M`$; and since, according to 5.2, the absolute degrees of singularity of $`g`$ and $`g^{}`$ do not depend on whether they are regarded as maps into $`N`$ or into $`M`$, we now have a contradiction to the minimality of $`\mathrm{ADS}(g)`$. ∎
## 6. Homology of covering spaces
In this short section we shall apply and extend some results from concerning homology of covering spaces of $`3`$-manifolds. In this section all homology groups are understood to be defined with coefficients in $`_2`$.
###### 6.1.
If $`N`$ is a normal subgroup of a group $`G`$, we shall denote by $`G\mathrm{\#}N`$ the subgroup of $`G`$ generated by all elements of the form $`gag^1a^1b^2`$ with $`gG`$ and $`a,bN`$. (This is a special case of the notation used in and . Here we are taking the prime $`p`$, which was arbitrary in and , to be $`2`$.
###### 6.2.
As in Section 1 of , for any group $`\mathrm{\Gamma }`$, we define subgroups $`\mathrm{\Gamma }_d`$ of $`\mathrm{\Gamma }`$ recursively for $`d0`$, by setting $`\mathrm{\Gamma }_0=\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_{d+1}=\mathrm{\Gamma }\mathrm{\#}\mathrm{\Gamma }_d`$. We regard $`\mathrm{\Gamma }_d/\mathrm{\Gamma }_{d+1}`$ as a $`_2`$-vector space.
###### Lemma 6.3.
Let $`M`$ be a closed $`3`$-manifold and set $`r=\mathrm{rk}_2M`$. Suppose that $`\stackrel{~}{M}`$ is a regular cover of $`M`$ whose group of deck transformations is isomorphic to $`(_2)^m`$ for some integer $`m0`$. Then
$$\mathrm{rk}_2(\stackrel{~}{M})mr\frac{m(m+1)}{2}.$$
###### Proof.
We set $`\mathrm{\Gamma }=\pi _1(M)`$ and define $`\mathrm{\Gamma }_d`$ for each $`d0`$ as in 6.1. We have $`\mathrm{rk}_2\mathrm{\Gamma }/\mathrm{\Gamma }_1=\mathrm{rk}_2M=r`$. It then follows from \[18, Lemma 1.3\] that $`\mathrm{rk}_2(\mathrm{\Gamma }_1/\mathrm{\Gamma }_2)r(r1)/2`$.
Let $`N`$ denote the normal subgroup of $`\mathrm{\Gamma }`$ corresponding to the regular covering space $`\stackrel{~}{M}`$. Since $`\mathrm{\Gamma }/N(_2)^m`$, we may write $`N=E\mathrm{\Gamma }_1`$ for some $`(rm)`$-generator subgroup $`E`$ of $`\mathrm{\Gamma }`$. It now follows from \[18, Lemma 1.4\] that
$`\mathrm{rk}_2\stackrel{~}{M}`$ $`=`$ $`\mathrm{rk}_2H_1(E\mathrm{\Gamma }_1)`$
$``$ $`\mathrm{rk}_2(\mathrm{\Gamma }_1/\mathrm{\Gamma }_2){\displaystyle \frac{(rm)(rm1)}{2}}`$
$``$ $`{\displaystyle \frac{r(r1)}{2}}{\displaystyle \frac{(rm)(rm1)}{2}}=mr{\displaystyle \frac{m(m+1)}{2}}.`$
The case $`m=2`$ of Lemma 6.3 will be applied in the proof of Lemma 8.5.
## 7. An application of a result of Gabai’s
This section contains the applications of Gabai’s results that were mentioned in the introduction. The main result of the section is Proposition 7.5.
###### Lemma 7.1.
Let $`X`$ be a PL space, let $`K`$ be a closed, connected, orientable surface of genus $`g>0`$, and let $`f:KX`$ be a PL map. Suppose that the homomorphism $`f_{}:H_2(K;_2)H_2(X;_2)`$ is trivial. Then the image of $`f_{}:H_1(K;_2)H_1(X;_2)`$ has dimension at most $`g`$.
###### Proof.
Since $`f_{}:H_2(K;_2)H_2(X;_2)`$ is trivial, it follows that the dual homomorphism $`f^{}:H^2(X;_2)H^2(K;_2)`$ is also trivial. Hence for any $`\alpha ,\beta H^1(X)`$ we have
$$f^{}(\alpha )f^{}(\beta )=f^{}(\alpha \beta )=0.$$
Thus if we set $`V=H^1(K;_2)`$ and let $`LV`$ denote the image of of $`f^{}:H^1(X;_2)H^1(K;_2)`$, we have $`LL=0`$, i.e.
$$LL^{}=\{vV:vL=0\}.$$
Hence if $`d`$ denotes the dimension of $`L`$, we have
$$ddimL^{}.$$
But by Poincaré duality, the cup product pairing on $`V`$ is non-singular, and so
$$dimL^{}=dimVdimL=2gd.$$
Hence $`dg`$. As the linear map $`f_{}:H_1(K;_2)H_1(X;_2)`$ is dual to $`f^{}:H^1(X;_2)H^1(K;_2)`$, its rank is the same as that of $`f^{}`$, namely $`d`$. The conclusion follows. ∎
###### Notation 7.2.
If $`F`$ is a closed, orientable $`2`$-manifold, we shall denote by $`\mathrm{TG}(F)`$ the total genus of $`F`$, i.e. the sum of the genera of its components.
###### Lemma 7.3.
For any compact, connected, orientable $`3`$-manifold $`N`$, we have
$$\mathrm{TG}(N)\mathrm{rk}_2N.$$
###### Proof.
In the exact sequence
$$H_2(N,N;_2)H_1(N;_2)H_1(N;_2),$$
Poincaré-Lefschetz duality implies that the vector spaces $`H_2(N,N;_2)`$ and $`H_1(N;_2)`$ are of the same dimension, $`\mathrm{rk}_2N`$. Hence we have
$$2\mathrm{TG}(N)=\mathrm{rk}_2N2\mathrm{rk}_2N$$
and the conclusion follows. ∎
###### Lemma 7.4.
If $`N`$ is a compact, connected, orientable $`3`$-manifold $`N`$ such that $`N`$ has at most one connected component, then $`H_2(N;)`$ is torsion-free.
###### Proof.
In the exact sequence
$$H_2(N;)H_2(N;)H_2(N,N;)$$
the inclusion map $`H_2(N;)H_2(N;)`$ is trivial since $`N`$ has at most one connected component. Hence the map $`H_2(N;)H_2(N,N;)`$ is injective, so that $`H_2(N;)`$ is isomorphic to a subgroup of $`H_2(N,N;)`$. But by Poincaré-Lefschetz duality, $`H_2(N,N;)`$ is isomorphic to $`H^1(N,)`$ and is therefore torsion-free. The conclusion follows. ∎
###### Proposition 7.5.
Suppose that $`N`$ is a compact (possibly closed) orientable $`3`$-manifold which is irreducible and boundary-irreducible. Suppose that $`K`$ is a closed, connected, orientable surface of genus $`g2`$, and that $`\varphi :KN`$ is a $`\pi _1`$-injective PL map. Then either
1. $`N`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$, or
2. the $`_2`$-vector subspace $`\varphi _{}(H_1(K;_2))`$ of $`H_1(N;_2)`$ has dimension at most $`g`$.
Furthermore, if $`\varphi _{}:H_1(K;_2)H_1(N;_2)`$ is surjective and $`N\mathrm{}`$, then (i) holds.
###### Proof.
We begin with the observation that $`N`$ is non-simply connected in view of the existence of the map $`\varphi `$. Since $`N`$ is also irreducible, it follows that no component of $`N`$ is a sphere. On the other hand, since $`N`$ is boundary-irreducible, every component of $`N`$ is $`\pi _1`$-injective in $`N`$. Thus every component of $`N`$ is parallel to an incompressible surface in $`N`$.
To prove the first assertion of the proposition we distinguish three cases, which are not mutually exclusive but cover all possibilities.
Case A. The homomorphism $`\varphi _{}:H_2(K;)H_2(N;)`$ is trivial.
Case B. The surface $`N`$ has at least two components.
Case C. The surface $`N`$ has at most one component and $`\varphi _{}:H_2(K;)H_2(N;)`$ is a non-trivial homomorphism.
To prove the assertion in Case A, we first consider the commutative diagram
in which the vertical maps are natural homomorphisms and the horizontal maps are induced by $`\varphi `$. The left-hand vertical arrow is surjective because the surface $`K`$ is orientable. Since the top horizontal map is trivial, it follows that the bottom horizontal map is trivial. Hence Lemma 7.1 asserts that the image of $`\varphi _{}:H_1(K;_2)H_1(N;_2)`$ has dimension at most $`g`$. Thus alternative (2) of the conclusion holds in Case A.
In Case B, using Lemma 7.3 and the surjectivity of $`\varphi _{}:H_1(K;_2)H_1(N;_2)`$, we find that
$$\mathrm{TG}(N)\mathrm{rk}_2N\mathrm{rk}_2K=2g.$$
Since $`N`$ has at least two components in this case, some component $`F`$ of $`N`$ must have genus at most $`g`$. By the observation at the beginning of the proof, $`F`$ is parallel to an incompressible surface in $`N`$. Thus alternative (1) of the conclusion holds in Case B.
To prove the assertion in Case C, we begin by considering the commutative diagram
in which the vertical maps are natural homomorphisms and the horizontal maps are induced by $`\varphi `$. Since $`M`$ has at most one component, Lemma 7.4 asserts that $`H_2(N;)`$ is torsion-free. Hence the right-hand vertical arrow in the diagram is injective. Since the top horizontal map is non-trivial, it follows that the bottom horizontal map is non-trivial. In other words, if $`[K]`$ denotes the fundamental class in $`H_2(K;)`$ then the class $`\alpha =f_{}([K])H_2(N;)`$ is non-zero.
We shall now apply a result from . For any $`2`$-manifold $``$ we shall denote by $`\chi _\mathrm{\_}()`$ the quantity
$$\underset{F}{}\mathrm{max}(\overline{\chi }(F),0),$$
where $`F`$ ranges over the components of $``$. As in , given a class $`z`$ in $`H_2(M;)`$, we denote by $`x_s(z)`$ and $`x(z)`$ respectively the “norm based on singular surfaces” and the Thurston norm of $`z`$. Since $`\alpha `$ is by definition realized by a map of the surface $`K`$ into $`N`$, and since $`\chi _\mathrm{\_}(K)=2g2`$, we have $`x_s(\alpha )2g2`$. But it follows from \[10, Corollary 6.18\] that $`x(\alpha )=x_s(\alpha )`$. Hence $`x(\alpha )2g2`$. By definition this means that if $``$ is a closed orientable embedded surface in $`intN`$ such that the fundamental class $`[]H_1(;)`$ is mapped to $`\alpha `$ under inclusion, and if $``$ is chosen among all such surfaces so as to minimize $`\chi _\mathrm{\_}()`$, then $`\chi _\mathrm{\_}()2g2`$. Since $`\alpha 0`$ we have $`\mathrm{}`$.
Since $`N`$ is irreducible, any sphere component of $``$ must be homologically trivial in $`N`$. We may assume that every torus component of $`F`$ is compressible, as otherwise alternative (1) of the conclusion holds. Under this assumption, if $`T`$ is a torus component of $``$, compressing $`T`$ yields a sphere which must be homologically trivial; hence $`T`$ is itself homologically trivial. Thus after discarding homologically trivial components of $``$ whose Euler characteristics are $`0`$, we may suppose that no component of $``$ is a sphere or torus. The minimality of $`\chi _\mathrm{\_}()`$ now implies that $``$ is incompressible.
Let $`F`$ be any component of $``$. Then $`F`$ is an incompressible closed surface in $`N`$, and we have
$$\chi _\mathrm{\_}(F)\chi _\mathrm{\_}()2g2.$$
Hence $`F`$ has genus at most $`g`$, and alternative (1) holds. This completes the proof of the first assertion of the proposition.
To prove the second assertion, suppose that $`\varphi _{}:H_1(K;_2)H_1(N;_2)`$ is surjective, that $`N\mathrm{}`$, and that alternative (2) holds. Then $`\mathrm{rk}_2Ng`$, and it follows from Lemma 7.3 that $`\mathrm{TG}(N)g`$. In particular, any component $`F`$ of the non-empty $`2`$-manifold $`N`$ has genus at most $`g`$. By the observation at the beginning of the proof, $`F`$ is parallel to an incompressible surface in $`N`$. Thus alternative (1) of the conclusion holds. ∎
## 8. Towers
In this section we prove a result, Proposition 8.10, which summarizes the tower construction described in the introduction. Our main topological result, Theorem 8.13, will then be proved by combining Proposition 8.10 with results from the earlier sections. We begin by introducing some machinery that will be needed for the statement and proof of Proposition 8.10.
###### Definition 8.1.
Suppose that $`n`$ is a non-negative integer. We define a height-$`n`$ tower of $`3`$-manifolds to be a $`(3n+2)`$-tuple
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n),$$
where $`M_0,\mathrm{},M_n`$ are compact, connected, orientable PL $`3`$-manifolds, $`N_j`$ is a compact, connected $`3`$-dimensional PL submanifold of $`M_j`$ for $`j=0,\mathrm{},n`$, and $`p_j:M_jN_{j1}`$ is a PL covering map for $`j=1,\mathrm{},n`$. We shall refer to $`M_0`$ as the base of the tower $`𝒯`$ and to $`N_n`$ as its top. We define the tower map associated to $`𝒯`$ to be the map
$$h=\iota _0p_1\iota _1p_2\mathrm{}p_n\iota _n:N_nM_0,$$
where $`\iota _j:N_jM_j`$ denotes the inclusion map for $`j=0,\mathrm{},n`$.
###### 8.2.
Consider any tower of $`3`$-manifolds
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n).$$
Note that for any given $`j`$ with $`0j<n`$, the manifold $`N_j`$ is closed if and only if its finite-sheeted covering space $`M_{j+1}`$ is closed. Note also that if, for a given $`j`$ with $`0jn`$, the submanifold $`N_j`$ of the (connected) manifold $`M_j`$ is closed, then we must have $`N_j=M_j`$, so that in particular $`M_j`$ is closed.
It follows that if $`M_j`$ is closed for a given $`j`$ with $`0jn`$, then $`M_i`$ is also closed for every $`i`$ with $`0ij`$. Thus either all the $`M_j`$ have non-empty boundaries, or there is an index $`j_0`$ with $`0j_0n`$ such that $`M_j`$ is closed when $`0jj_0`$ and $`M_j`$ has non-empty boundary when $`j_0<jn`$. Furthermore, in the latter case, for each $`j<j_0`$ we have $`N_j=M_j`$.
###### 8.3.
In particular, if in a tower of $`3`$-manifolds
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n)$$
the manifold $`M_j`$ is closed for a given $`jn`$, then for every $`i`$ with $`0i<j`$ the composition
$$p_{j1}\mathrm{}p_i:M_jM_i$$
is a well-defined covering map, whose degree is the product of the degrees of $`p_i,\mathrm{},p_{j1}`$.
###### Definition 8.4.
A tower of $`3`$-manifolds
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n)$$
will be termed good if it has the following properties:
1. $`M_j`$ and $`N_j`$ are irreducible for $`j=0,\mathrm{},n`$;
2. $`N_j`$ is a (possibly empty) incompressible surface in $`M_j`$ for $`j=0,\mathrm{},n`$;
3. the covering map $`p_j:M_jN_{j1}`$ has degree $`2`$ for $`j=1,\mathrm{},n`$; and
4. for each $`j`$ with $`2jn`$ such that $`M_j`$ is closed, the four-fold covering map (see 8.3)
$$p_jp_{j1}:M_jM_{j2}$$
is regular and has covering group isomorphic to $`_2\times _2`$.
###### Lemma 8.5.
Suppose that
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n)$$
is a good tower of $`3`$-manifolds and that $`j_0`$ is an index with $`0j_0n`$ such that $`M_{j_0}`$ is closed. Set $`r=\mathrm{rk}_2M_0`$ and assume that $`r3`$. For any index $`j`$ with $`0jj_0`$, we have
$$\mathrm{rk}_2M_j2^{j/2}(r3)+3$$
if $`j`$ is even, and
$$\mathrm{rk}_2M_j2^{(j1)/2}(r3)+2$$
if $`j`$ is odd.
In particular, we have $`\mathrm{rk}_2M_jr1`$ for each $`j`$ with $`0jn`$ such that $`M_j`$ is closed, and we have $`\mathrm{rk}_2M_j2r4`$ for each $`j`$ with $`2jn`$ such that $`M_j`$ is closed.
###### Proof.
According to 8.2, $`M_j`$ is closed for every index $`j`$ with $`0jj_0`$. We shall first show that for every even $`j`$ with $`0jj_0`$ we have $`\mathrm{rk}_2M_j2^{j/2}(r3)+3`$. For $`j=0`$ this is trivial since $`r=\mathrm{rk}_2M_0`$. Now, arguing inductively, suppose that $`j`$ is even, that $`0<jn`$, and that $`\mathrm{rk}_2M_{j2}2^{(j2)/2}(r3)+3`$. Since the definition of a good tower implies that $`M_j`$ is a regular $`(_2\times _2)`$-cover of $`M_{j2}`$, we apply Lemma 6.3 with $`m=2`$ to deduce that
$$\mathrm{rk}_2M_j2(\mathrm{rk}_2M_{j2})32(2^{(j2)/2}(r3)+3)3=2^{j/2}(r3)+3.$$
This completes the induction and shows that $`\mathrm{rk}_2M_j2^{j/2}(r3)+3`$ for every even index $`j`$ with $`2jj_0`$. Finally, if $`j`$ is an odd index with $`2<jj_0`$, then since $`j1`$ is even we have $`\mathrm{rk}_2M_{j1}2^{(j1)/2}(r3)+3`$; and since $`M_j`$ is a $`2`$-sheeted cover of $`M_{j1}`$, it is clear that $`\mathrm{rk}_2M_j\mathrm{rk}_2M_{j1}12^{(j1)/2}(r3)+2`$. ∎
###### Definition 8.6.
If
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n)$$
is a height-$`n`$ tower of $`3`$-manifolds, then for any $`m`$ with $`0mn`$, the $`(3m+2)`$-tuple
$$𝒯^{}=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_m,M_m,N_m)$$
is a height-$`n`$ tower. We shall refer to the tower $`𝒯^{}`$ as the height-$`m`$ truncation of $`𝒯`$. We shall say that a tower $`𝒯^+`$ is an extension of a tower $`𝒯`$, or that $`𝒯^+`$ extends $`𝒯`$, if $`𝒯`$ is a truncation of $`𝒯^+`$.
In particular, any tower may be regarded as an extension of itself. This will be called the degenerate extension.
###### Definition 8.7.
Let $`𝒯`$ be a tower of $`3`$-manifolds with base $`M`$ and top $`N`$, and let $`h:NM`$ denote the associated tower map. Let $`\varphi `$ be a PL map from a compact PL space $`K`$ to $`M`$. By a homotopy-lift of $`\varphi `$ through the tower $`𝒯`$ we mean a PL map $`\stackrel{~}{\varphi }:KN`$ such that $`h\stackrel{~}{\varphi }`$ is homotopic to $`\varphi `$. A homotopy-lift $`\stackrel{~}{\varphi }`$ of $`\varphi `$ will be termed good if the inclusion homomorphism $`\pi _1(\stackrel{~}{\varphi }(K))\pi _1(N)`$ is surjective.
###### Lemma 8.8.
Suppose that $`K`$ is a compact PL space with freely indecomposable fundamental group of rank $`k2`$. Suppose that $`𝒯=(M_0,N_0,p_1,\mathrm{},N_n)`$ is a good tower of $`3`$-manifolds of height $`n`$. Suppose that $`\varphi :KM_0`$ is a $`\pi _1`$-injective PL map, and that $`\stackrel{~}{\varphi }:KN_n`$ is a good homotopy-lift of $`\varphi `$ through the tower $`𝒯`$. Suppose that $`p_{n+1}:M_{n+1}N_n`$ is a two-sheeted covering space of $`N_n`$, and that the map $`\stackrel{~}{\varphi }:KN_n`$ admits a lift to the covering space $`M_{n+1}`$. Suppose that either
1. $`n1`$, the manifold $`N_n`$ is closed (so that $`M_{n+1}`$ is closed, cf. 8.2), and the covering map
$$p_np_{n+1}:M_{n+1}M_{n1}$$
is regular and has covering group isomorphic to $`_2\times _2`$; or
2. $`M_{n+1}\mathrm{}`$, or
3. $`n=0`$.
Then there exists a compact submanifold $`N_{n+1}`$ of $`M_{n+1}`$ with the following properties:
1. $`𝒯^+=(M_0,N_0,p_1,\mathrm{},N_n,p_{n+1},M_{n+1},N_{n+1})`$ is a good height-$`(n+1)`$ tower extending $`𝒯`$, and
2. there is a a good homotopy-lift $`\stackrel{~}{\varphi }^+`$ of $`\varphi `$ through the tower $`𝒯^+`$ such that
$$\mathrm{ADS}(\stackrel{~}{\varphi }^+)<\mathrm{ADS}(\stackrel{~}{\varphi }).$$
###### Proof.
Let $`h:N_nM_0`$ be the tower map associated to $`𝒯`$. We fix a lift $`f:KM_{n+1}`$ of the map $`\stackrel{~}{\varphi }:KN_n`$ to the covering space $`M_{n+1}`$. Since $`\stackrel{~}{\varphi }`$ is a homotopy lift of $`\varphi `$, the map $`hp_{n+1}f:KM_0`$ is homotopic to $`\varphi `$. Since $`\varphi _{\mathrm{}}:\pi _1(K)\pi _1(M_0)`$ is injective, it now follows that $`f_{\mathrm{}}:\pi _1(K)\pi _1(M_{n+1})`$ is also injective. We may therefore apply Proposition 5.8 to this map $`f`$, taking $`M=M_{n+1}`$ and $`N=N_{n+1}`$. We choose a map $`g:KM_{n+1}`$ homotopic to $`f`$, with $`\mathrm{ADS}(g)\mathrm{ADS}(f)`$, and a compact $`3`$-dimensional submanifold $`N=N_{n+1}`$ of $`intM_{n+1}`$, such that conditions (i)—(iv) of 5.8 hold with $`M=M_{n+1}`$.
It is clear from the definition that $`𝒯^+=(M_0,N_0,p_1,\mathrm{},N_n,p_{n+1},M_{n+1},N_{n+1})`$ is a tower extending $`𝒯`$. To show that the tower $`𝒯^+`$ is good, we first observe that conditions (1)—(4) of Definition 8.4 hold whenever $`jn`$ because $`𝒯`$ is a good tower. For $`j=n+1`$, Conditions (1) and (2) of Definition 8.4 follow from conditions (iv) and (iii) of 5.8, while condition (3) of Definition 8.4 follows from the hypothesis that $`p_{n+1}:M_{n+1}N_n`$ is a two-sheeted covering. The case $`j=n+1`$ of Condition (4) of Definition 8.4 is clear if alternative ($`\alpha `$) of the hypothesis holds, and is vacuously true if alternative ($`\beta `$) or ($`\gamma `$) holds. Hence $`𝒯^+`$ is a good tower.
Since by condition (i) of Proposition 5.8 we have $`intN_{n+1}g(K)`$, we may regard $`g:KM_{n+1}`$ as a composition $`\iota _{n+1}\stackrel{~}{\varphi }^+`$, where $`\iota _{n+1}:N_{n+1}M_{n+1}`$ is the inclusion map and $`\stackrel{~}{\varphi }^+`$ is a PL map from $`K`$ to $`N_{n+1}`$. Since $`g`$ is homotopic to $`f`$, the map $`hp_{n+1}\iota _{n+1}\stackrel{~}{\varphi }^+=hp_{n+1}g:KM_0`$ is homotopic to $`\varphi `$. It follows that $`\stackrel{~}{\varphi }^+`$ is a homotopy-lift of $`\varphi `$ through the tower $`𝒯^+`$. Condition (ii) of 5.8 asserts that the inclusion homomorphism $`\pi _1(\stackrel{~}{\varphi }^+(K))\pi _1(N_{n+1})`$ is surjective, which according to Definition 8.7 means that the homotopy-lift $`\stackrel{~}{\varphi }^+`$ of $`\varphi `$ is good.
Finally, since the homotopy-lift $`\stackrel{~}{\varphi }`$ of $`\varphi `$ is good by hypothesis, the inclusion homomorphism $`\pi _1(\stackrel{~}{\varphi }^+(K))\pi _1(N_n)`$ is surjective. As $`f`$ is a lift of $`\stackrel{~}{\varphi }`$ to the non-trivial covering space $`M_{n+1}`$ of $`N_n`$, it follows from Proposition 5.4 that $`\mathrm{ADS}(f)<\mathrm{ADS}(\stackrel{~}{\varphi })`$. But we chose $`g`$ in such a way that $`\mathrm{ADS}(g)\mathrm{ADS}(f)`$, and according to 5.2 we have $`\mathrm{ADS}(\stackrel{~}{\varphi }^+)=\mathrm{ADS}(g)`$. Hence $`\mathrm{ADS}(\stackrel{~}{\varphi }^+)<\mathrm{ADS}(\stackrel{~}{\varphi })`$. ∎
###### Lemma 8.9.
Suppose that $`K`$ is a closed orientable surface of genus $`g2`$. Suppose that $`𝒯`$ is a good tower of $`3`$-manifolds of height $`n`$. Let $`M`$ denote the base of $`𝒯`$, and assume that $`\mathrm{rk}_2Mg+3`$. Suppose that $`\varphi :KM`$ is a $`\pi _1`$-injective PL map, and that $`\stackrel{~}{\varphi }`$ is a good homotopy-lift of $`\varphi `$ through the tower $`𝒯`$. Then at least one of the following alternatives holds:
1. $`N_n`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$;
2. $`n1`$ and $`N_{n1}`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$; or
3. there exist a height-$`(n+1)`$ extension $`𝒯^+`$ of $`𝒯`$ which is a good tower, and a good homotopy-lift $`\stackrel{~}{\varphi }^+`$ of $`\varphi `$ through the tower $`𝒯^+`$, such that
$$\mathrm{ADS}(\stackrel{~}{\varphi }^+)<\mathrm{ADS}(\stackrel{~}{\varphi }).$$
###### Proof.
We write
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n),$$
so that $`M=M_0`$. We distinguish several cases.
Case A: $`N_n\mathrm{}`$ and the homomorphism $`\stackrel{~}{\varphi }_{}:H_1(K;_2)H_1(N_n;_2)`$ is surjective;
Case B: $`N_n\mathrm{}`$ and $`\stackrel{~}{\varphi }_{}:H_1(K;_2)H_1(N_n;_2)`$ is not surjective;
Case C: $`n=0`$;
Case D: $`n1`$ and $`N_n`$ is closed.
In Case A, all the hypotheses of the final assertion of Proposition 7.5 hold with $`\stackrel{~}{\varphi }`$ in place of $`\varphi `$. It therefore follows from the final assertion of Proposition 7.5 that conclusion (1) of the present lemma holds.
In Case B, the map $`\stackrel{~}{\varphi }:KN_n`$ admits a lift to some two-sheeted covering space $`p_{n+1}:M_{n+1}N_n`$ of $`N_n`$. Since $`N_n\mathrm{}`$, we have $`M_{n+1}\mathrm{}.`$ This is alternative ($`\beta `$) of the hypothesis of Lemma 8.8. It therefore follows from 8.8 that conclusion (3) of the present lemma holds.
In Case C the argument is identical to the one used in Case B, except that we have alternative $`(\gamma )`$ of Lemma 8.8 in place of alternative ($`\beta `$).
We now turn to Case D. In this case, as was observed in 8.2, we have $`N_n=M_n`$ and $`N_{n1}=M_{n1}`$, and $`p_n`$ is a two-sheeted covering map from $`M_n`$ to $`M_{n1}`$.
Let us set $`r=\mathrm{rk}_2Mg+3`$. According to Lemma 8.5, for any index $`j`$ such that $`1jn`$ and such that $`M_j`$ is closed, we have $`\mathrm{rk}_2M_jr1`$. In particular, if we set $`d=\mathrm{rk}_2M_{n1}`$, we have $`dr1g+2`$.
Now set $`\stackrel{~}{\varphi }^{}=p_n\stackrel{~}{\varphi }:KM_{n1}`$. Then $`X=\stackrel{~}{\varphi }_{}^{}(H_1(K;_2))`$ is a subspace of the $`d`$-dimensional $`_2`$-vector space $`V=H_1(M_{n1};_2)`$.
The hypotheses of Proposition 7.5 hold with $`N`$ and $`\varphi `$ replaced by $`M_{n1}`$ and $`\stackrel{~}{\varphi }^{}`$. Hence either $`M_{n1}`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$, or $`X`$ has dimension at most $`g`$. The first alternative gives conclusion (2) of the present lemma.
There remains the subcase in which $`X`$ has dimension at most $`g`$. Since $`dr1g+2`$, the dimension of $`X`$ is then at most $`gd2`$.
In this subcase we shall show that $`\stackrel{~}{\varphi }:KM_n`$ admits a lift to some two-sheeted covering space $`p_{n+1}:M_{n+1}M_n`$ of $`M_n=N_n`$ for which alternative ($`\alpha `$) of the hypothesis of Lemma 8.8 holds. It will will then follow from 8.8 that conclusion (3) of the present lemma holds.
Let $`q`$ denote the natural homomorphism from $`\pi _1(M_{n1})`$ to $`H_1(M_{n1};_2)`$. The two-sheeted cover $`M_n`$ of $`M_{n1}`$ corresponds to a normal subgroup of $`\pi _1(M_{n1})`$ having the form $`q^1(Z)`$, where $`Z`$ is some $`(d1)`$-dimensional vector subspace of $`V`$. Since $`\stackrel{~}{\varphi }^{}`$ admits the lift $`\stackrel{~}{\varphi }`$ to $`M_n`$, we have $`XZV`$. Since in addition we have $`\mathrm{rk}_2Xd2<d1=rankZ`$, there exists a $`(d2)`$-dimensional vector subspace $`Y`$ of $`V`$ with $`XYZ`$. The subgroup $`q^1(Y)`$ determines a regular covering space $`M_{n+1}`$ of $`M_{n1}`$ with covering group $`_2\times _2`$. Since $`q^1(Y)q^1(Z)`$, the degree-four covering map $`M_{n+1}M_{n1}`$ factors as the composition of a degree-two covering map $`p_{n+1}:M_{n+1}M_n`$ with $`p_n:M_nM_{n+1}`$. Thus the covering space $`p_{n+1}:M_{n+1}M_n`$ satisfies alternative ($`\alpha `$) of 8.8. It remains to show that $`\stackrel{~}{\varphi }`$ admits a lift to $`M_{n+1}`$.
Since $`\stackrel{~}{\varphi }_{\mathrm{}}^{}(\pi _1(K))q^1(X)q^1(Y)`$, the map $`\stackrel{~}{\varphi }^{}`$ admits a lift to the four-fold cover $`M_{n+1}`$ of $`M_{n1}`$. Since $`M_{n+1}`$ is a regular covering space of $`M_{n1}`$, there exist four different lifts of $`\varphi ^{}`$ to $`M_{n+1}`$. But $`\stackrel{~}{\varphi }^{}`$ can have at most two lifts to $`M_n`$, and each of these can have at most two lifts to $`M_{n+1}`$. Hence each lift of $`\stackrel{~}{\varphi }^{}`$ to $`M_n`$ admits a lift to $`M_{n+1}`$. In particular, $`\stackrel{~}{\varphi }`$ admits a lift to $`M_{n+1}`$. ∎
###### Lemma 8.10.
Suppose that $`K`$ is a closed, orientable surface of genus $`g2`$. Suppose that $`M`$ is a closed, orientable $`3`$-manifold such that $`\mathrm{rk}_2Mg+3`$, and that $`\varphi :KM`$ is a $`\pi _1`$-injective PL map. Suppose that
$$𝒯_0=(M_0,N_0,p_1,\mathrm{},N_{n_0})$$
is a good tower with base $`M`$ such that $`\varphi `$ admits a good homotopy-lift through $`𝒯`$. Then either
1. $`n_01`$, and $`N_{n_01}`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$, or
2. there exists a good tower $`𝒯_1`$ which is a (possibly degenerate) extension of $`𝒯_0`$, such that the top $`N`$ of $`𝒯_1`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$, and $`\varphi `$ admits a good homotopy-lift $`\stackrel{~}{\varphi }_1`$ through the tower $`𝒯_1`$.
###### Proof.
Let us fix a good homotopy-lift $`\stackrel{~}{\varphi }_0`$ of $`\varphi `$ through $`𝒯_0`$. Let $`𝒮`$ denote the set of all ordered pairs $`(𝒯,\stackrel{~}{\varphi })`$ such that $`𝒯`$ is a good tower which is an extension of $`𝒯_0`$ and $`\stackrel{~}{\varphi }`$ is a good homotopy-lift of $`\varphi `$ through $`𝒯`$. Then we have $`(𝒯_0,\varphi _0)𝒮`$, and so $`𝒮\mathrm{}`$. Hence there is an element $`(𝒯_1,\stackrel{~}{\varphi }_1)`$ of $`𝒮`$ such that $`\mathrm{ADS}(\stackrel{~}{\varphi }_1)\mathrm{ADS}(\stackrel{~}{\varphi })`$ for every element $`(𝒯,\stackrel{~}{\varphi })`$ of $`𝒮`$. Let us write
$$𝒯_1=(M_0,N_0,p_1,\mathrm{},N_{n_1}).$$
The hypotheses of Lemma 8.9 now hold with $`𝒯_1`$ and $`\stackrel{~}{\varphi }_1`$ in place of $`𝒯`$ and $`\stackrel{~}{\varphi }`$. Hence one of the following alternatives must hold:
1. $`N_{n_1}`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$;
2. $`n_11`$ and $`N_{n_11}`$ contains a connected (non-empty) closed incompressible surface of genus at most $`g`$; or
3. there exist a height-$`(n_1+1)`$ extension $`𝒯^+`$ of $`𝒯_1`$ which is a good tower, and a good homotopy-lift $`\stackrel{~}{\varphi }^+`$ of $`\varphi `$ through the tower $`𝒯^+`$, such that
$$\mathrm{ADS}(\stackrel{~}{\varphi }^+)<\mathrm{ADS}(\stackrel{~}{\varphi }_1).$$
If 8.9(1) holds, then the tower $`𝒯_1`$ has the property asserted in conclusion (1) of the present lemma. If 8.9(2) holds, and if $`n_1>n_0`$ (i.e. $`𝒯_1`$ is a non-degenerate extension of $`𝒯_0`$), then the height-$`(n_11)`$ truncation $`𝒯_1^{}`$ of $`𝒯_1`$ is an extension of $`𝒯_0`$, and conclusion (2) holds with $`𝒯_1^{}`$ in place of $`𝒯_1`$. If 8.9(2) holds and $`n_1=n_0`$ (i.e. $`𝒯_1`$ is a degenerate extension of $`𝒯_0`$), conclusion (2) of the present lemma holds. Finally, if 8.9(3) holds, then $`(𝒯^+,\stackrel{~}{\varphi }^+)𝒮`$, and we have a contradiction to the minimality of $`\mathrm{ADS}(\stackrel{~}{\varphi }_1)`$. ∎
###### Proposition 8.11.
Suppose that $`g`$ is an integer $`2`$, that $`M`$ is a closed, orientable $`3`$-manifold with $`\mathrm{rk}_2Mg+3`$, and that $`\pi _1(M)`$ has a subgroup isomorphic to a genus-$`g`$ surface group. Then there exists a good tower
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n),$$
with base $`M=M_0`$, such that $`N_n`$ contains a connected incompressible closed surface of genus $`g`$.
###### Proof.
Let $`K`$ denote a closed, orientable surface of genus $`g`$. The hypothesis implies that there is a $`\pi _1`$-injective PL map $`\varphi :KM`$. According to Proposition 5.8, there exist a PL map $`\stackrel{~}{\varphi }_0:KM`$ homotopic to $`\varphi `$, and a compact, connected $`3`$-submanifold $`N_0`$ of $`intM`$, such that (i) $`intN_0\stackrel{~}{\varphi }_0(K)`$, (ii) the inclusion homomorphism $`\pi _1(\stackrel{~}{\varphi }_0(K))\pi _1(N_0)`$ is surjective, (iii) $`N_0`$ is incompressible in $`M`$, and (iv) $`N_0`$ is irreducible. According to the definitions, this means that $`𝒯_0=(M,N_0)`$ is a good tower of height $`0`$ and that $`\stackrel{~}{\varphi }_0`$ is a good homotopy-lift of $`\varphi `$ through $`𝒯_0`$.
We apply Proposition 8.10 with these choices of $`𝒯_0`$ and $`\stackrel{~}{\varphi }_0`$. Conclusion (2) of 8.10 cannot hold since $`𝒯_0`$ has height $`0`$. Hence conclusion (1) must hold. The extension $`𝒯=𝒯_0`$ of $`𝒯_0`$ given by conclusion (1) is a good tower whose top contains a connected, closed incompressible surface of genus at most $`g`$. ∎
###### Lemma 8.12.
Suppose that
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n)$$
is a good tower of $`3`$-manifolds such that $`M_0`$ is simple. Then the manifolds $`M_j`$ and $`N_j`$ are all simple for $`j=0,\mathrm{},n`$.
###### Proof.
By hypothesis $`M_0`$ is simple. If $`M_j`$ is simple for a given $`jn`$, then since $`N_j`$ is a submanifold of $`M_j`$ bounded by a (possibly disconnected and possibly empty) incompressible surface, it is clear from Definition 1.10 that $`N_j`$ is simple. If $`j<n`$ it then follows from 1.11 that the two-sheeted covering space $`M_{j+1}`$ of $`N_j`$ is also simple. ∎
The following theorem is the main topological result of this paper.
###### Theorem 8.13.
Let $`g`$ be an integer $`2`$. Let $`M`$ be a closed simple $`3`$-manifold such that $`\mathrm{rk}_2M4g1`$ and $`\pi _1(M)`$ has a subgroup isomorphic to a genus-$`g`$ surface group. Then $`M`$ contains a closed, incompressible surface of genus at most $`g`$.
###### Proof.
Applying Proposition 8.11, we find a good tower
$$𝒯=(M_0,N_0,p_1,M_1,N_1,p_2,\mathrm{},p_n,M_n,N_n),$$
with base $`M_0`$ homeomorphic to $`M`$ and with some height $`n0`$, such that $`N_n`$ contains a connected incompressible closed surface $`F`$ of genus $`g`$. According to the definition of a good tower, $`N_n`$ is incompressible (and, a priori, possibly empty) in $`M_n`$. Hence $`N_n`$ is $`\pi _1`$-injective in $`M_n`$. Since $`F`$ is incompressible in $`N_n`$, it follows that it is also incompressible in $`M_n`$.
Since $`M`$ is simple it follows from Lemma 8.12 that all the $`M_j`$ and $`N_j`$ are simple.
Let $`m`$ denote the least integer in $`\{0,\mathrm{},n\}`$ for which $`M_m`$ contains a closed incompressible surface $`S_m`$ of genus at most $`g`$. To prove the theorem it suffices to show that $`m=0`$. Let $`h`$ denote the genus of $`S_m`$.
Suppose to the contrary that $`m1`$. Then the hypotheses of Proposition 4.4 hold with $`N_{m1}`$ and $`M_m`$ playing the respective roles of $`N`$ and $`\stackrel{~}{N}`$. Suppose that conclusion (1) of 4.4 holds, i.e. that $`N_{m1}`$ contains an incompressible closed surface $`S_{m1}`$ with $`genus(S_{m1})hg`$. According to the definition of a good tower, $`N_{m1}`$ is an incompressible (and possibly empty) surface in $`M_{m1}`$. Hence $`N_{m1}`$ is $`\pi _1`$-injective in $`M_{m1}`$. Since $`S_{m1}`$ is incompressible in $`N_{m1}`$, it follows that it is also incompressible in $`M_{m1}`$. We therefore have a contradiction to the minimality of $`m`$.
Hence conclusion (2) of 4.4 must hold; in particular, $`N_{m1}`$ is closed, so that $`N_{m1}=M_{m1}`$; and $`\mathrm{rk}_2M_{m1}=\mathrm{rk}_2N_{m1}4h34g3`$. On the other hand, since by hypothesis we have $`\mathrm{rk}_2M_04g1`$, it follows from Lemma 8.5 that for any index $`j`$ such that $`0jn`$ and such that $`M_j`$ is closed, we have $`\mathrm{rk}_2M_j4g2`$. This is a contradiction, and the proof is complete. ∎
## 9. Proof of the geometric theorem
As a preliminary to the proof of Theorem 9.6 we shall point out how the Marden tameness conjecture, recently established by Agol and by Calegari-Gabai , strengthens the results proved in .
We first recall some definitions from \[5, Section 8\]. Let $`\mathrm{\Gamma }`$ be a discrete torsion-free subgroup of $`Isom_+(^3)`$, and let $`k2`$ be an integer. We say that $`\lambda `$ is a $`k`$-Margulis number for $`\mathrm{\Gamma }`$, or for $`M=^3/\mathrm{\Gamma }`$, if for any $`k`$ elements $`\xi _1,\mathrm{},\xi _k\mathrm{\Gamma }`$, and for any $`z^3`$, we have that either
* the group $`\xi _1,\mathrm{},\xi _k`$ is generated by at most $`k1`$ abelian subgroups, or
* $`\mathrm{max}_{i=1}^kdist(\xi _iz,z)\lambda `$.
We say that $`\lambda `$ is a strong $`k`$-Margulis number for $`\mathrm{\Gamma }`$, or for $`M`$, if for any $`k`$ elements $`\xi _1,\mathrm{},\xi _k\mathrm{\Gamma }`$, and for any $`z^3`$, we have that either
* the group $`\xi _1,\mathrm{},\xi _k`$ is generated by at most $`k1`$ abelian subgroups, or
* $`{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{1+e^{dist(\xi _iz,z)}}}{\displaystyle \frac{k}{1+e^\lambda }}.`$
We note that if $`\lambda `$ is a strong $`k`$-Margulis number for $`\mathrm{\Gamma }`$, then $`\lambda `$ is also a $`k`$-Margulis number for $`\mathrm{\Gamma }`$.
A group $`\mathrm{\Gamma }`$ is termed $`k`$-free, where $`k`$ is a positive integer, if every subgroup of $`\mathrm{\Gamma }`$ whose rank is at most $`k`$ is free.
###### Theorem 9.1.
Let $`k2`$ be an integer and let $`\mathrm{\Gamma }`$ be a discrete subgroup of $`Isom_+(^3)`$. Suppose that $`\mathrm{\Gamma }`$ is $`k`$-free and purely loxodromic. Then $`\mathrm{log}(2k1)`$ is a strong $`k`$-Margulis number for $`\mathrm{\Gamma }`$.
###### Proof.
This is the same statement as \[5, Proposition 8.1\] except that the latter result contains the additional hypothesis that $`\mathrm{\Gamma }`$ is $`k`$-tame, in the sense that every subgroup of $`\mathrm{\Gamma }`$ whose rank is at most $`k`$ is topologically tame. (To say that a discrete torsion-free subgroup $`\mathrm{\Delta }`$ of $`Isom_+(^3)`$ is topologically tame means that $`^3/\mathrm{\Delta }`$ is diffeomorphic to the interior of a compact $`3`$-manifold.) But the main theorem of or asserts that any finitely generated discrete torsion-free subgroup $`\mathrm{\Delta }`$ of $`Isom_+(^3)`$ is topologically tame. ∎
By combining this with another result from , we shall prove:
###### Theorem 9.2.
Suppose that $`M`$ is an orientable hyperbolic $`3`$-manifold without cusps and that $`\pi _1(M)`$ is $`3`$-free. Then either $`M`$ contains a hyperbolic ball of radius $`(\mathrm{log}5)/2`$, or $`\pi _1(M)`$ is a free group of rank $`2`$.
###### Proof.
We may write $`M=^3/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }Isom(^3)`$ is discrete and purely loxodromic. Since $`\mathrm{\Gamma }\pi _1(M)`$ is $`3`$-free, it follows from Theorem 9.1 that $`\mathrm{log}5`$ is a strong $`3`$-Margulis number, and in particular a Margulis number, for $`\mathrm{\Gamma }`$ (or equivalently for $`\mathrm{\Gamma }`$). According to \[5, Theorem 9.1\], if $`M`$ is a hyperbolic $`3`$-manifold without cusps, if $`\pi _1(M)`$ is $`3`$-free and if $`\lambda `$ a $`3`$-Margulis number for $`M`$, then either $`M`$ contains a hyperbolic ball of radius $`\lambda /2`$, or $`\pi _1(M)`$ is a free group of rank $`2`$. The assertion follows. ∎
###### Corollary 9.3.
Suppose that $`M`$ is a closed orientable hyperbolic 3-manifold such that $`\pi _1(M)`$ is $`3`$-free. Then $`M`$ contains a hyperbolic ball of radius $`(\mathrm{log}5)/2`$. Hence the volume of $`M`$ is greater than $`3.08`$.
###### Proof.
It follows from Theorem 9.2 that either $`M`$ contains a hyperbolic ball of radius $`(\mathrm{log}5)/2`$ or $`\pi _1(M)`$ is a free group of rank $`2`$. The latter alternative is impossible, because $`\mathrm{\Gamma }`$, as the fundamental group of a closed hyperbolic $`3`$-manifold, must have cohomological dimension $`3`$, whereas a free group has cohomological dimension 1. Thus $`M`$ must contain a hyperbolic ball of radius $`(\mathrm{log}5)/2`$.
The lower bound on the volume now follows by applying Böröczky’s density estimate for hyperbolic sphere-packings as in . ∎
###### Theorem 9.4 (Agol-Storm-Thurston).
Suppose that $`M`$ is a closed orientable hyperbolic $`3`$-manifold containing a connected incompressible closed surface $`S`$. Then either $`Vol(M)>3.66`$, or the manifold $`X`$ obtained by splitting $`M`$ along $`S`$ has the form $`X=|𝒲|`$ for some (possibly disconnected) book of $`I`$-bundles $`𝒲`$.
###### Proof.
According to \[4, Corollary 2.2\], if $`S`$ is an incompressible closed surface in a closed orientable hyperbolic $`3`$-manifold $`M`$, if $`X`$ denotes the manifold obtained by splitting $`M`$ along $`S`$, and if $`K=\overline{X\mathrm{\Sigma }}`$ where $`\mathrm{\Sigma }`$ denotes the relative characteristic submanifold of the simple manifold $`X`$, then the volume of $`M`$ is greater than $`\overline{\chi }(K)3.66`$. Hence either $`Vol(M)>3.66`$ or $`\chi (K)=0`$. In the latter case, we shall show that $`X`$ is a book of $`I`$-bundles; this will complete the proof.
Note that each component of $`K`$ must have Euler characteristic $`0`$, because the components of the frontier of $`K`$ in $`X`$ are essential annuli in $`X`$. Since $`\chi (K)=0`$ it follows that each component of $`K`$ has Euler characteristic $`0`$. Hence if $`Y`$ denotes the union of all components of $`\mathrm{\Sigma }`$ with strictly negative Euler characteristic, and if we set $`Z=\overline{XY}`$, then each component of $`Z`$ has Euler characteristic $`0`$. But $`Z`$ is $`\pi _1`$-injective in $`X`$ because its frontier components are essential annuli. Since $`X`$ is simple, the components of $`Z`$ are solid tori. Since $`Y=\overline{XZ}`$ is an $`I`$-bundle with $`YZ=_vY`$, and the components of $`_vY`$ are $`\pi _1`$-injective in $`X`$ and hence in $`Z`$, it follows from the definition that $`X`$ is a book of $`I`$-bundles. ∎
###### Proposition 9.5.
Suppose that $`M`$ is a closed orientable hyperbolic 3-manifold such that $`\mathrm{rk}_2M7`$. Suppose that $`\pi _1(M)`$ has a subgroup isomorphic to a genus-$`2`$ surface group. Then $`VolM3.66`$.
###### Proof.
Since $`M`$ is simple and $`\mathrm{rk}_2M7`$, it follows from Theorem 8.13 that $`M`$ contains either a closed incompressible surface of genus $`2`$.
Suppose that $`VolM<3.66`$. Let $`X`$ denote the manifold obtained by splitting $`M`$ along $`S`$. According to Theorem 9.4, each component of $`MS`$ has the form $`|𝒲|`$ for some book of $`I`$-bundles $`𝒲`$.
Consider the subcase in which $`X`$ is connected. Since $`S`$ has genus $`2`$, we have $`\overline{\chi }(X)=2`$. By Lemma 2.11 it follows that
$$\mathrm{rk}_2(X)2\overline{\chi }(X)+15.$$
Hence
$$\mathrm{rk}_2M\mathrm{rk}_2X+16,$$
a contradiction to the hypothesis.
There remains the case in which $`X`$ has two components, say $`X_1`$ and $`X_2`$. Since $`S`$ has genus $`2`$, we have $`\overline{\chi }(X_i)=1`$ for $`i=1,2`$. By Lemma 2.11, it follows that
$$\mathrm{rk}_2(X_i)2\overline{\chi }(X_i)+1=3.$$
Hence
$$\mathrm{rk}_2M\mathrm{rk}_2X_1+\mathrm{rk}_2X_26,$$
and we have a contradiction. (The bound of $`6`$ in this last inequality could easily be improved to $`4`$, but this is obviously not needed.) ∎
We can now prove our main geometrical result.
###### Theorem 9.6.
Let $`M`$ be a closed orientable hyperbolic $`3`$-manifold such that $`VolM3.08`$. Then $`\mathrm{rk}_2M6`$.
###### Proof.
Assume that $`\mathrm{rk}_2M7`$. If $`\pi _1(M)`$ has a subgroup isomorphic to a genus-$`2`$ surface group, then it follows from Proposition 9.5 (with $`g=2`$) that $`VolM3.66>3.08`$, a contradiction to the hypothesis. There remains the possibility that $`\pi _1(M)`$ has no subgroup isomorphic to a genus-$`2`$ surface group. In this case, since $`\mathrm{rk}_2M5`$, it follows from \[5, Proposition 7.4 and Remark 7.5\] that $`\pi _1(M)`$ is $`3`$-free. Hence by Corollary 9.3 we have $`VolM>3.08`$, and again the hypothesis is contradicted. ∎
|
warning/0506/astro-ph0506390.html
|
ar5iv
|
text
|
# Measuring the primordial power spectrum: Principal component analysis of the cosmic microwave background
## 1 Introduction
Observations of the cosmic microwave background (CMB) anisotropies are presenting a fascinating opportunity for discerning between our models for the origin of structure in the universe in great detail. Indeed the most recent observations of the CMB from the Wilkinson Microwave Anisotropy Probe (*WMAP*) have vindicated a basic picture for the primordial perturbations which are nearly scale-invariant, Gaussian and adiabatic in nature, and which are dominated by a passive and growing-mode. This represents enormous progress by instrumentalists in the thirty years since Zel’dovich and Novikov lamented in their 1975 monologue over the observational prospects for measuring the CMB anisotropies: ‘*Given all the difficulties, it is not clear that we will ever successfully investigate the nature of the initial perturbations using the concept of \[Sakharov\] modulation \[of the acoustic peaks\]* ’ (Zel’dovich & Novikov 1975).
At this time, therefore, there is an overall consistency between observations (Peiris et al 2003; Barger, Lee & Marfatia 2003; Leach & Liddle 2003) and the inflationary paradigm which is well-known to contain a mechanism for generating large-scale perturbations of this type (see Liddle & Lyth 2000; Dodelson 2003). In the near future though, most progress in our understanding of the origin of structure is likely to come from empirical studies of the primordial perturbations where one of the known ingredients of the standard Gaussian adiabatic model is relaxed to a more general form. Indeed, this has been the spirit in which many authors have proceeded. In particular there has been a strong interest in measuring the shape of the primordial power spectrum, given the prospect of a factor twenty or so increase in the data to this sector of cosmology in the near future, coming from ground-based, balloon-borne and satellite experiments.
Model-independent methods for reconstructing the primordial power spectrum are being investigated where one only relies on the broad assumption that the overall picture of Gaussian adiabatic perturbations is correct. The available data are then confronted a more general primordial power spectrum sector, and the full parameter space is integrated out in a medium size computation. Many such power spectrum parametrisations exist and these include bandpowers (Wang, Spergel & Strauss 1999; Bridle et al 2003; Hannestad 2004), band-colours (Bridle et al 2003), wavelet bandpowers (Mukherjee & Wang 2003a,c), orthogonal wavelets (Mukherjee & Wang 2003b). The specific choices to be made such as the number and the location of the bandpowers will require a certain amount of optimisation. However, these promising methods are known to perform well on both real and simulated data without degrading too far the expected constraints on the remaining cosmological parameters (Bond et al 2004; Mukherjee & Wang 2005).
One can also apply inverse methods in order to reconstruct the primordial power spectrum, since the problem at hand is akin to deconvolution. Many methods have been investigated and these include semi-analytic iterative methods (Kogo et al 2005), the Richardon–Lucy deconvolution algorithm (Shafieloo & Souradeep 2004), regularised least-squares (Tegmark & Zaldarriaga 2002; Tocchini-Valentini, Douspis & Silk 2005). While these strategies may provide a reasonable glimpse of the form of the primordial power spectrum at a lower computational cost, they typically suffer a weakness that the cosmological parameters must be fixed to some representative values and are not integrated out. In addition there is usually a smoothing step involved either in the data or the deconvolved power spectrum requiring a careful treatment.
There is a data compression strategy which, although it is most similar in spirit to the model-independent methods described above, corresponds to asking a a slightly different question than whether we can reconstruct or deconvolve the primordial power spectrum. Although the question we refer to has been in the air and in the minds of many people for years, and is partially addressed by any CMB analysis that constrains the power-law slope of the primordial power spectrum, it is worth stating it here explicitly: *Are scale-invariant adiabatic perturbations an ingredient of our cosmology and how can we best measure any departures from scale-invariance?* This question is important because its eventual answer will represent the next step in our attempts to model and understand the underlying mechanism responsible generating the primordial perturbations. We will demonstrate in this paper that principal component analysis is very well suited for this purpose. Briefly summarised, the trick is to choose a complete orthonormal power spectrum basis which also reflects our expectation of where the departures from scale-invariance are likely to be best probed by the data, as has been repeatedly emphasised by Hu and collaborators (Hu & Okamoto 2004; Kadota et al 2005). The full cosmological plus power spectrum parameter space can be integrated out in a medium to large scale computation, and theoretical predictions for the power spectrum can be easily projected on onto the same power spectrum basis to make the comparison with observations.
The outline of this paper is to describe the principal component analysis formalism, providing a commentary of the relevant implementation details in §2; in §3 we test the method with simulated *Planck* data using three primordial power spectra which are respectively scale-invariant, scale-free, and broken scale-invariant; in §4 we apply the method to the *WMAP* data before concluding §5.
## 2 Principal component analysis formalism
In this paper we implement and investigate the principal component analysis (hereafter PCA) method detailed and described by Hu and Okamoto (2004) (hereafter HO04) which should be considered a companion paper. PCA has also been applied or discussed in countless other contexts in which data volumes have already or will soon be seeing sharp increases, for instance in galaxy-galaxy power spectrum estimation methods (Hamilton and Tegmark 2000), reionization history reconstruction (Hu and Holder 2003), dark energy reconstruction (Huterer & Starkman 2003) and most recently in the context of reconstructing the inflation potential (Kadota et al 2005). It can be thought of simply as a change of parameter basis, where the rotation is determined by properties of the observed or expected signal and noise. At the same time it is also a very useful lossless data compression technique.
The basic set-up in the context of the CMB is not at all unfamiliar to astrophysics, that of a deconvolution problem
$`C_{\mathrm{}}^{XX^{}}={\displaystyle \frac{2\pi }{\mathrm{}(\mathrm{}+1)}}{\displaystyle d\mathrm{ln}k𝒫\left(k\right)T_{\mathrm{}}^X(k;\{\omega _i\})T_{\mathrm{}}^X^{}(k;\{\omega _i\})},`$ (1)
where $`X=T,E`$ and the dependence of the CMB transfer functions $`T_{\mathrm{}}^X\left(k\right)`$ on the cosmological parameters $`\{\omega _i\}`$ has been written explicitly in order to show the added complication over and above an ordinary deconvolution problem of this type. Interestingly, there is a satisfactory solution to the problem of extracting the primordial power spectrum $`𝒫(k)`$, described in HO04, which involves exploiting what we know about the expected noise on $`C_{\mathrm{}}`$ and our precise and accurate knowledge of the CMB transfer function physics (Seljak et al 2003). Here we present the relevant equations from HO04.
The response of the $`C_{\mathrm{}}`$ with respect to some primordial power spectrum parameters $`\{p_i\}`$ can be investigated via a mode counting approach by constructing the Fisher information matrix
$`F_{ij}={\displaystyle \underset{\mathrm{}=2}{\overset{\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \frac{2\mathrm{}+1}{2}}\mathrm{Tr}[𝐃_\mathrm{}i𝐂_{\mathrm{}}^1𝐃_j\mathrm{}𝐂_{\mathrm{}}^1],`$ (2)
which has been written using a matrix notation, where
$`(𝐃_\mathrm{}i)_{XX^{}}=D_\mathrm{}i^{XX^{}}{\displaystyle \frac{C_{\mathrm{}}^{XX^{}}}{p_i}},`$ (3)
and where
$`D_\mathrm{}i^{XX^{}}`$ $`=`$ $`{\displaystyle \frac{C_{\mathrm{}}^{XX^{}}}{p_i}}|_{\mathrm{fid}}`$
$`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}(\mathrm{}+1)}}{\displaystyle d\mathrm{ln}k𝒫_0T_{\mathrm{}}^X(k)T_{\mathrm{}}^X^{}(k)W_i(\mathrm{ln}k)}.`$
We can take our power spectrum test function $`W_i`$ to be the triangle window
$`W_i(\mathrm{ln}k)=\mathrm{max}[1\left|{\displaystyle \frac{\mathrm{ln}k\mathrm{ln}k_i}{\mathrm{\Delta }\mathrm{ln}k}}\right|,0].`$ (5)
In this work we have used a discretisation $`\mathrm{\Delta }\mathrm{ln}k=0.00875`$ spanning a range of scales that traverses the acoustic peaks from $`0.004<k<0.2\mathrm{Mpc}^1`$. It is worth noting at this stage that this range need not include the largest scales responsible for the Sachs–Wolfe effect: the Fisher information on these scales tends to zero, and so it proves convenient to truncate these scales in order to later on invert the Fisher information matrix without numerical difficulties. The calculation of the power spectrum transfer functions $`D_\mathrm{}i^{XX^{}}`$ is achieved by making minor modifications to the CAMB CMB anisotropies code (Lewis, Challinor & Lasenby 2000) (based on CMBFAST (Seljak & Zaldarriaga 1996)), rather than using a full Boltzmann hierarchy code used in HO04. CAMB is run at slightly higher accuracy where we have increased by a factor four both the number of source and integration $`k`$ modes, and have calculated $`D_\mathrm{}i^{XX^{}}`$ at every $`\mathrm{}`$, rather than the usual splining method with $`\mathrm{\Delta }\mathrm{}50`$, in order to capture the high frequency information.
The choice of fiducial cosmological parameters is given by a baryon density $`\mathrm{\Omega }_\mathrm{B}h^2=0.024`$, cold dark matter density $`\mathrm{\Omega }_\mathrm{D}h^2=0.121`$, present Hubble rate $`H_0[\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1]=72`$, optical depth to last scattering $`\tau =0.17`$, and a curvature perturbation amplitude $`𝒫_0=23\times 10^{10}`$. We assume a spatially flat cosmology and ignore the effect of lensing. The latter will be important to take into account in a more thorough analysis in order avoid biasing of the recovered cosmological parameters (Hu & Okamoto 2004; Lewis 2005).
In Fig. 1 we illustrate the Fisher information matrix given by equation (2) which shows a band-diagonal structure with peaks of sensitivity to the primordial power spectrum on scales corresponding to the acoustic peaks; the sensitivity drops again on scales corresponding to the acoustic troughs, which can be remedied by information coming from the phase-shifted polarization peaks. Of course the sensitivity tends to zero on large scales due to a lack of modes to observe, and on small scales due to Silk damping and beam smoothing, since the $`C_{\mathrm{}}`$ of equation (2) is replaced by the total signal plus a Gaussian white noise model adjusted for a given experiment
$`C_{\mathrm{}}^{TT}|_{\mathrm{noise}}`$ $`=`$ $`\sigma _{\mathrm{noise}}^2e^{\mathrm{}(\mathrm{}+1)\theta ^2/8\mathrm{ln}2},`$
$`C_{\mathrm{}}^{EE}|_{\mathrm{noise}}`$ $`=`$ $`2\times \sigma _{\mathrm{noise}}^2e^{\mathrm{}(\mathrm{}+1)\theta ^2/8\mathrm{ln}2},`$ (6)
$`C_{\mathrm{}}^{TE}|_{\mathrm{noise}}`$ $`=`$ $`0,`$
where $`\sigma _{\mathrm{noise}}^2`$ is the noise variance in $`(\mu `$K-rad$`)^2`$ and $`\theta `$ is the FWHM of a Gaussian beam in radians. The noise model should be considered an important input to the analysis since it determines the range of scales that will be probed; it is an additional ingredient compared to the majority of analyses of the $`C_{\mathrm{}}`$ data. We use here a noise model for *Planck* with $`\sigma _{\mathrm{noise}}^2=3\times 10^4(\mu `$K-rad$`)^2`$ and $`\theta =7^{}`$, and a noise model for *WMAP* with $`\sigma _{\mathrm{noise}}^2=8.4\times 10^3(\mu `$K-rad$`)^2`$ and $`\theta =13^{}`$. In a realistic analysis the observed signal plus noise spectrum will be more appropriate.
As usual the Fisher information matrix can be inverted to obtain a covariance matrix $`C_{ij}`$ whose diagonal components provide a useful estimate, the Cramer–Rao bound, of the expected variance of the parameters $`p_i`$ with
$`\sigma ^2(p_i)=C_{ii}(F^1)_{ii}.`$ (7)
In Fig. 2 we plot this window of sensitivity to the primordial power spectrum (on a scale $`\delta \mathrm{ln}k0.05`$ set by the Fisher matrix bandwidth) for the *Planck* satellite, which can be seen to encompass the entire acoustic peak region. As noted in HO04, outside this range of scales, and in particular on large scales, we can only hope to recover wide-band ($`\delta \mathrm{ln}k0.05`$) averages of the primordial power spectrum at high accuracy.
The PCA basis $`\{m_i\}`$ is simply a linear combination of the power spectrum spike basis $`\{p_i\}`$
$`m_a=(\mathrm{\Delta }\mathrm{ln}k)^{1/2}{\displaystyle \underset{i}{}}S_{ia}p_i`$ (8)
where the $`S_{ia}`$ are the orthonormal eigenvectors of the covariance matrix. We can then work with a set of normalised principal components $`𝒮_{ia}=S_{ia}/\sqrt{\mathrm{\Delta }\mathrm{ln}k}`$ (hereafter the PCA modes) which will have unit variance when integrated over $`\mathrm{ln}k`$. In Figs. 3 and 4 we plot examples of the PCA modes with mode number from 1–4 and 17–20 respectively, generated using the *WMAP* noise model. The oscillations in the PCA modes become increasingly rapid at scales corresponding to the acoustic peaks where sensitivity to the primordial power spectrum is greatest, that is until we hit the numerical resolution. At this point the PCA modes branch into two wavepacket-like solutions travelling towards large and small scales, similar to the behaviour noted by Hamilton and Tegmark (2000), although this need not worry us. Note also that the PCA modes are invariant under changes in the discretisation scale $`\mathrm{\Delta }\mathrm{ln}k`$. However, we found that in order to obtain sensible estimates of the eigenvalues (projected errors) of the PCA modes themselves, the Fisher matrix should be discretised on a scale that renders it roughly diagonal, instead of band-diagonal.
The PCA modes can be straightforwardly integrated into the publicly available Markov Chain Monte Carlo (MCMC) package CosmoMC<sup>1</sup><sup>1</sup>1http://cosmologist.info/cosmomc/ (Lewis & Bridle 2002, February 2005 version) in order to explore the full cosmological plus power spectrum posterior parameter space. Specifically, we use the following power spectrum ansatz
$`{\displaystyle \frac{𝒫(k)}{𝒫_0}}=m_0+{\displaystyle \underset{a=1}{\overset{a_{\mathrm{max}}}{}}}m_a𝒮_a(k),`$ (9)
where we take $`𝒫_0=23\times 10^{10}`$, which should be calibrated from observations. Clearly if the underlying primordial power spectrum is close to scale-invariant then equation (9) admits a solution
$`m_a=0,a\text{Scale-invariance}.`$ (10)
More generally equation (9) is strongly suggestive of a general linear orthonormal model plus a noise term (see for instance Bretthorst 1988). In this way we are attempting to measure the spectrum of departures from scale-invariance which we call $`\mathrm{\Delta }𝒫/𝒫_0`$ and which is given by the second term in equation (9); in this context the dominant scale-invariant component $`m_0`$ is a Gaussian noise term.
Concerning the numerical implementation of the power spectrum equation (9), we simply perform a linear spline in $`\mathrm{ln}k`$ over the discrete PCA modes $`𝒮_{ia}`$, which are added together before the final convolution with CMB transfer functions to obtain the $`C_{\mathrm{}}`$; outside the PCA mode $`k`$-range the second term of equation (9) is dropped. We checked that the default $`k`$-source and $`k`$-integration settings for CAMB modified to calculate $`C_{\mathrm{}}`$ at $`\mathrm{\Delta }\mathrm{}=3`$ is accurate enough handle around the first forty modes of our current implementation; at this stage this is more than enough since we will only attempt to perform the MCMC with a maximum of sixteen PCA modes.
Having obtained measurements of the PCA mode amplitudes from the MCMC, it is then straightforward to project any power spectrum model, for instance a power-law spectrum, onto the PCA modes via
$`m_a`$ $`=`$ $`{\displaystyle d\mathrm{ln}k𝒮_a(k)\frac{\mathrm{\Delta }𝒫}{𝒫_0}(k)},`$ (11)
$`=`$ $`\mathrm{\Delta }\mathrm{ln}k{\displaystyle \underset{i}{}}𝒮_a(k_i)\left[\left({\displaystyle \frac{k_i}{k_0}}\right)^{n_\mathrm{S}1}1\right],`$
in order to make the comparison with observations.
We can easily make an empirical estimate of the total signal to noise of the measured departures from scale-invariance
$`{\displaystyle \frac{\mathrm{S}}{\mathrm{N}}}=\left[{\displaystyle \underset{a=1}{\overset{a_{\mathrm{max}}}{}}}{\displaystyle \frac{m_a^2}{\sigma _{m_a}^2}}\right]^{1/2},`$ (12)
where $`m_a`$ and $`\sigma _{m_a}^2`$ are the mean and variance of the individual mode amplitudes obtained from the MCMC. As noted by Kadota et al (2005), the PCA modes can be safely truncated as soon as $`\mathrm{S}/\mathrm{N}`$ saturates, assuming that the underlying primordial power spectrum is a reasonably smooth function. Incidentally, the total $`\mathrm{S}/\mathrm{N}`$ represents a useful figure of merit for optimising future CMB experiments to measure the primordial power spectrum sector. Other measures such as “risk” (Huterer & Starkman 2003) and Bayesian evidence (see for example MacKay 2003) could be used to provide a rationale for truncating the PCA mode amplitudes even further, given a power spectrum model of interest.
In the case that the recovered PCA mode amplitudes encode some complex information which can not be easily understood in the framework of power-law spectra, then it would be useful to obtain an estimate of $`\mathrm{\Delta }𝒫/P_0`$ in $`k`$-space in order to aid the process of modelling the power spectrum. Here we use an estimator
$`{\displaystyle \frac{\widehat{\mathrm{\Delta }𝒫}(k_i)}{𝒫_0}}={\displaystyle \underset{a=1}{\overset{a_{\mathrm{max}}}{}}}m_a𝒮_a(k_i),`$ (13)
and for the purposes of a comparison with the input spectrum, we estimate the noise variance via
$`\widehat{\sigma }_{\frac{\mathrm{\Delta }𝒫}{𝒫_0}}^2(k_i)=C_{ii}+{\displaystyle \underset{a=1}{\overset{a_{\mathrm{max}}}{}}}𝒮_a^2(k_i)\sigma _{m_a}^2,`$ (14)
where $`C_{ii}`$ is the covariance matrix, obtained from equation (7), accounting for the overall uncertainty in the narrow-band determination of $`\mathrm{\Delta }𝒫/𝒫_0`$ in regions of lower sensitivity on large scales, small scales, and in the temperature acoustic trough regions.
A bandpower representation of the primordial power spectrum could also obtained from the measured PCA mode amplitudes via a Monte Carlo procedure; in this case the Fisher information matrix could be used for guidance when choosing the location and widths of the bands. Obviously though, no further quantitative information about the primordial power spectrum can be gleaned in this way.
One final point worth making in this section concerns how one should deal with the inevitable degeneracies between the effect on the $`C_{\mathrm{}}`$ due to the cosmological parameters and the PCA power spectrum parameters, which will induce undesired off-diagonal components in the PCA covariance matrix. We sketch here the solution given in HO04: One must first form the joint Fisher information matrix, $`F_{\mu \nu }`$, for both power spectrum parameters and cosmological parameters
$`F_{\mu \nu }=\left[\begin{array}{ccc}\hfill F_{ij}& & \hfill B\\ \hfill B^\mathrm{T}& & \hfill F_{ab}\end{array}\right],`$ (17)
where $`F_{ab}`$ is the usual cosmological parameter Fisher information matrix (see for example Tegmark, Taylor and Heavens 1997) and $`B`$ are the cross terms. After inverting the full $`F_{\mu \nu }`$ to obtain a new covariance matrix $`C_{\mu \nu }`$, one simply retains the power spectrum parameter subblock $`C_{ij}`$, whose principal components will be “orthogonalized” to the effect of the cosmological parameters. In terms of implementation, one can use the matrix partitioning formulas (see for example Press et al 1992, Ó Ruanaidh & Fitzgerald 1996) to derive a “degraded” $`F_{ij}^{\mathrm{deg}}`$ subblock
$`F_{ij}^{\mathrm{deg}}=F_{ij}B^\mathrm{T}F_{\mathrm{ab}}B.`$ (18)
We will make use of this in the next section.
## 3 Tests with Simulated *Planck* data
As a means of gaining experience with the PCA method we generate simulated *Planck* data up to an $`\mathrm{}_{\mathrm{max}}=2250`$ using the Gaussian white noise model of equation (2) for a cosmological model with parameters $`\mathrm{\Omega }_\mathrm{B}h^2=0.024`$, $`\mathrm{\Omega }_\mathrm{D}h^2=0.121`$, $`H_0=72`$, $`\tau =0.17`$, and $`𝒫_0=2.3\times 10^9`$, which for simplicity are the same as those used to generate the PCA modes themselves. In a realistic data analysis scenario, the PCA modes would be generated with parameters close to the best-fit obtained from a traditional parameter determination approach. We consider three cases for the primordial power spectrum which is taken to described by a scale-invariant spectrum, a power-law spectrum with spectral index $`n_\mathrm{S}=0.97`$ and pivot scale $`k_0=0.05\mathrm{Mpc}^1`$, and then finally a broken scale-invariance model with a Gaussian bump in the acoustic peak region
$`{\displaystyle \frac{𝒫(k)}{𝒫_0}}=1+0.1\mathrm{exp}\left[\left({\displaystyle \frac{\mathrm{ln}\left[k/0.08\mathrm{Mpc}^1\right]}{0.3}}\right)^2\right].`$ (19)
We then perform MCMC over the full cosmological plus PCA mode parameter space using the simulated data up to an $`\mathrm{}_{\mathrm{max}}=2000`$. We have also varied the number of modes included in the analysis from zero to sixteen in steps of four in order to study the effect of truncating the PCA expansion on the recovery of the cosmological parameters.
The development of CosmoMC (Lewis & Bridle 2002) has reached a maturity that is very well suited to an analysis of this type where the number of power spectrum parameters begins to dominate over the number of cosmological parameters, but where we nonetheless expect by construction to obtain a stable multivariate Gaussian posterior solution. As a result we have taken full advantage of a conjugate gradients descent module which estimates the covariance and location of the posterior peak before the MCMC begins, thus alleviating the potential challenge working with so many parameters while also conserving some computing resources. On this note, the total number of $`C_{\mathrm{}}`$ likelihood evaluations required in our tests in the following section rises from around $`𝒩_{}=2\times 10^410^6`$ for zero and eight PCA modes respectively, and then tends to saturate at around this number. It seems reasonable that the number of likelihood evaluations ought not to exceed by much $`\mathrm{}^2`$, the total number of modes upon which the the $`C_{\mathrm{}}`$ spectrum depends. Moreover the ‘fast–slow’ split between power spectrum and cosmological parameter likelihood evaluation speeds, already implemented in CosmoMC, will be of increasing benefit as we attempt to measure up to perhaps thirty PCA mode amplitudes in the future (Kadota et al 2005).
### 3.1 The scale-invariant case
In Fig. 5 we illustrate the recovery of the first eight mode amplitudes for the $`n_\mathrm{S}=1`$ case and make comparison for the theoretical prediction for the mode amplitudes which are obtained by projecting some representative power-law spectra onto the PCA modes via equation (11); we find that the scale-invariant solution $`m_a=0`$ is very well recovered. Here it is worth mentioning that the Gaussian realisation for the simulated *Planck* data sets was taken to be the exact $`C_{\mathrm{}}`$ model, which explains why the recovery of the PCA mode amplitudes shows very little scatter around $`m_a=0`$. One can see that the first three PCA modes provide the bulk of constraining power for smooth power-law spectra leading to a constraint which will be roughly $`n_\mathrm{S}=1\pm 0.01`$, consistent with typical parameter forecasts in the literature.
We illustrate an estimate of the departures from scale-invariance $`\widehat{\mathrm{\Delta }𝒫}/𝒫_0`$ in Fig. 6, and the region with the most data weight can clearly be discerned showing consistency with a scale-invariant spectrum. In this case the recovery of the cosmological parameters is also excellent, and we recovered a stable Gaussian posterior (as a function of the number of PCA modes) with constraints given by $`\omega _\mathrm{B}h^2=0.0240\pm 0.0002`$, $`\omega _\mathrm{D}h^2=0.121\pm 0.02`$, $`H_0=71.9\pm 0.7`$, $`\tau =0.170\pm 0.005`$, $`𝒫/𝒫_0=1.00\pm 0.01`$ for the case of using eight PCA modes. Clearly the PCA method works well under these most idealised of circumstances.
### 3.2 The scale-free $`n_\mathrm{S}=0.97`$ case
The $`n_\mathrm{S}=0.97`$ case is more delicate since we know in advance that the power spectrum model equation (9) will not be able to accurately describe a tilted spectrum on large or small scales. We can therefore expect some biasing in the recovery of the cosmological parameters which will necessarily adjust to provide the overall excess of power on large scales relative to small scales; this is just the usual degeneracy between cosmological and power spectrum parameters.
In fact to get reasonable results at all, we found it necessary to apply equation (18) in order to orthogonalise the PCA modes to the effect of the primordial power spectrum amplitude $`𝒫_0`$. The qualitative effect on the PCA modes is the the positive definite mode 1 is removed. Having modified the PCA modes in this way, the cosmological parameters are recovered as $`\omega _\mathrm{B}h^2=0.0247\pm 0.0002`$, $`\omega _\mathrm{D}h^2=0.116\pm 0.001`$, $`H_0=74.6\pm 0.7`$, $`\tau =0.183\pm 0.006`$, $`𝒫/𝒫_0=1.02\pm 0.01`$ for the case of using eight PCA modes, showing biases at the 3$`\sigma `$ to 4$`\sigma `$ level. The fact that the recovered dark matter density shifts from $`\mathrm{\Omega }_\mathrm{D}h^2=0.113\pm 0.0010.116\pm 0.001`$ as the number of PCA modes is increased provides a useful indication that there are problems afoot with our power spectrum model equation (9).
Interestingly however, the PCA mode amplitudes are still very well recovered, and we illustrate in Fig. 7 that the first ten mode amplitudes, if somewhat attenuated in amplitude, provide strong evidence for a power-law primordial power spectrum, showing a distinctive pattern deviating from scale-invariance, $`m_\mathrm{a}=0`$. The corresponding departures from scale-invariance are shown in Fig. 8 where the recovered power spectrum shows strong evidence for a tilt, modulo some attenuation and oscillations in regions of lower sensitivity. In short there is enough signal to noise to overrule our assumption of scale-invariance, supplying us with strong evidence that model of equation (9) needs refining. It is likely that in a more refined analysis, one should orthogonalise the PCA modes to the effect of the spectral index and the other cosmological parameters in order to recover unbiased estimates of the cosmological parameters.
### 3.3 The Gaussian bump case
Although completely contrived, this is perhaps the most interesting and challenging case since the input primordial power spectrum now contains distinct feature within the acoustic peak region. We illustrate in Fig. 9 that the first sixteen PCA amplitudes are nonetheless rather well recovered and are consistent with the input Gaussian bump model. In this case we can see that, for instance, the second PCA mode strongly constrains the central position of the feature in $`k`$-space. In Fig. 10 we show that a bump like feature has indeed been recovered, again modulo some attenuation and oscillations in regions of lower sensitivity. The cosmological parameters are also very well recovered with $`\omega _\mathrm{B}h^2=0.0238\pm 0.0002`$, $`\omega _\mathrm{D}h^2=0.122\pm 0.002`$, $`H_0=71.6\pm 0.9`$, $`\tau =0.170\pm 0.005`$, $`𝒫/𝒫_0=1.00\pm 0.01`$. This represents an interesting success for the PCA method.
### 3.4 Summary and discussion
To summarise the tests so far, the PCA method has been demonstrated here to be very suitable and effective for measuring departures from scale-invariance, both scale-free and scale-dependent, in the most data-weighted regions of the $`C_{\mathrm{}}`$ spectrum. In a realistic data analysis setup the recovered PCA mode amplitudes, together with the PCA modes themselves will represent an extremely powerful compression of our information concerning the primordial power spectrum. At first sight this may represent an unnecessary data analysis stage compared the usual parameter determination methods where one fits to the $`C_{\mathrm{}}`$ data directly using the power spectrum model parameters on the same footing as the other cosmological parameters. However, the point here is to obtain first a detailed picture of the most important departures from scale-invariance in the primordial power spectrum while at the same time being able to weigh up the relative importance as well as locating both in $`k`$ and $`\mathrm{}`$ space any possible glitches or residual systematic effects in the $`C_{\mathrm{}}`$ data; then in the final data compression stage we can use the PCA mode amplitudes to rapidly test any wide class of specific power spectrum models with great ease and without recourse to any further $`C_{\mathrm{}}`$ likelihood evaluations, as was recently emphasised by Kadota et al (2005) for the case of inflation models.
## 4 Application to the *WMAP* data
In this Section we apply the PCA method to the currently available temperature and temperature-polarization cross correlation spectra from *WMAP* (Kogut et al 2003; Verde et al 2003; Hinshaw et al 2003) and bandpowers in the range $`600<\mathrm{}<2000`$ from the VSA (Grainge et al 2003; Dickinson et al 2004) ACBAR (Kuo et al 2004), CBI (Pearson et al 2003; Readhead et al 2004) and Boomerang B2K (Jones et al 2005, Piacentini et al 2005, Montroy et al 2005) instruments.
To emphasise once more, we are working within the framework of spatially flat $`\mathrm{\Lambda }`$CDM cosmologies, described by five basic cosmological parameters: the baryon density $`\mathrm{\Omega }_\mathrm{B}h^2`$, the cold dark matter density $`\mathrm{\Omega }_\mathrm{D}h^2`$, the optical depth to last scattering $`\tau `$, the ratio of the sound horizon to angular diameter distance at last scattering $`\theta =100r_\mathrm{s}^{}/D_a^{}`$ (instead of $`H_0`$) and the overall amplitude of scalar perturbations $`𝒫_0`$. In addition we throw into the mix the first four PCA modes generated with a noise model for *WMAP* given by $`\sigma _{\mathrm{noise}}^2=8.4\times 10^3(\mu `$K-rad$`)^2`$ and $`\theta =13^{}`$.
The measured amplitudes of the first four modes of Fig. 3 are displayed in Fig. 11 with the corresponding power spectrum in Fig. 12. The broad picture painted here is that we find no evidence for the breaking of scale-invariance: the mode amplitudes are very well fit my $`m_a=0`$. Only a single mode on scales corresponding to the second acoustic peak shows an S/N $`>1`$, which is barely worth mentioning aside from the fact that it can easily be accommodated by a slightly red primordial power spectrum: projecting power-law primordial power spectra onto the PCA basis and using a simple Gaussian likelihood function we find the constraint on the spectral index to be $`n_\mathrm{S}(k_0=0.04`$ Mpc$`{}_{}{}^{1})=0.94\pm 0.04`$, displayed in Fig. 13, and which is in accordance with conventional studies of the primordial power spectrum. It is also possible to make a detailed comparison with the primordial power spectrum bandpowers from fig. 4 of Bridle et al (2003), as well as with orthogonal wavelet expansion constraints in fig. 2 of Mukherjee & Wang (2003b). We all find the same very weak trend for a 20-30$`\%`$ drop in power between between the first acoustic peak at $`k=0.02`$ Mpc<sup>-1</sup> and the third acoustic peak scale at $`k=0.07`$ Mpc<sup>-1</sup>. Again, the trend is not so much interesting at this stage as the consistency between these complementary methods.
## 5 Conclusions
In this work we have implemented and investigated a principle component analysis (PCA) technique in order to study the possible departures from scale-invariance that may exist in the spectrum of primordial curvature perturbations, which are observable via the CMB anisotropies. The essence of this method is to decompose the primordial power spectrum into a scale-invariant component plus a series of orthonormal modes which reflect our expectation of where the departures from scale-invariance are likely to be best probed by the data. The information from the CMB is then be compressed into a series of mode amplitudes which can easily be compared with predictions from any wide class of primordial power spectra without recourse to any further $`C_{\mathrm{}}`$ likelihood evaluations.
The method was first tested on simulated *Planck* data using an input scale-invariant spectrum and we observed good performance in the simultaneous recovery of cosmological parameters and the principal component mode amplitudes via an MCMC exploration of the full parameter space. In the case of simulated data from an input power-law spectrum with spectral index $`n_\mathrm{S}=0.97`$, the recovery of the cosmological parameters was biased as they adjusted to provide an overall excess of large-scale to small-scale power. However, the biasing is evidenced by fluctuating cosmological parameter constraints as the number of power spectrum principal components is increased. Moreover, the PCA mode amplitudes were still very well recovered, showing strong evidence for a tilted primordial power spectrum and providing enough signal to noise to overrule our assumption of scale-invariance. Thus PCA can be used as a self-consistent means for justifying a more refined power spectrum model than the one considered here in equation (9). We also demonstrated that the PCA method is capable of measuring departures from scale-free spectra by considering simulated data from a primordial power spectrum containing a $`10\%`$ gaussian bump in the acoustic peak region, and observing good recovery of both the PCA mode amplitudes and the cosmological parameters.
Finally, as a proof of concept of the method we provided a first glimpse of the principal component mode amplitudes that can be obtained from the currently available CMB data from *WMAP*, VSA, ACBAR, CBI and Boomerang. We obtained measurements of the first four principal components corresponding to scales across the first and second acoustic peaks, finding no evidence for the breaking of scale-invariance with only a hint of a red primordial power spectrum with spectral index $`n_\mathrm{S}(k_0=0.04`$ Mpc$`{}_{}{}^{1})=0.94\pm 0.04`$, consistent with other studies in the literature, with a total signal to noise at not more than S/N $`2.5`$.
Assuming that the Gaussian adiabatic density perturbation scenario continues to hold as our observations of the CMB improve in the near future, then we will soon move into the regime where the information about the primordial power spectrum will completely outweigh the information about the cosmological parameters which become, from this perspective, well-understood nuisance parameters to be carefully integrated out. It seems very likely that principal component analysis, or else another very similar data compression technique, will be essential for fully exploiting the forthcoming temperature and polarization $`C_{\mathrm{}}`$ data.
## Acknowledgements
I thank Carlo Baccigalupi, Sergei Bashinsky, Uroš Seljak, Roberto Trotta and Ben Wandelt for their helpful discussions and comments, Antony Lewis for supplying his *Planck* simulations code via www.cosmocoffee.org, Christophe Ringeval for his help with the modifications to accuracy of the CAMB code, and the University of Geneva, where this work was started, and the Trieste Observatory for hosting the main computations in this work. I acknowledge the use of the Legacy Archive for Background Data Analysis (LAMBDA). Support for LAMBDA is provided by the Office of Space Science. I also acknowledge the use of the R statistical computing environment in which much of this analysis was performed, and the CosmoMC (Lewis & Bridle 2002) and CAMB (Lewis, Challinor & Lasenby 2000) packages.
|
warning/0506/math0506331.html
|
ar5iv
|
text
|
# On the origin of the BV operator on odd symplectic supermanifolds
## Two differentials on odd symplectic supermanifolds
In what follows, let $`M`$ be a supermanifold with an odd symplectic form $`\omega `$. On $`\mathrm{\Omega }(M)`$ we have two anticommuting differentials: one is de Rham’s $`d`$ and the other is the wedge product with $`\omega `$. We shall find that the Khudaverdian BV operator is (rougly speaking) the 3rd differential of the spectral sequence of this bicomplex.
###### Theorem.
Let $`(M,\omega )`$ be an odd symplectic supermanifold. In the spectral sequence of the bicomplex $`(\mathrm{\Omega }(M),\omega ,d)`$ we have:
1. the cohomology of the complex $`(\mathrm{\Omega }(M),\omega )`$ is naturally isomorphic to the semidensities on $`M`$
2. the next differential in the spectral sequence, de Rham’s $`d`$, vanishes on the cohomology of $`(\mathrm{\Omega }(M),\omega )`$
3. the next differential ($`d\left(\omega \right)^1d`$) coincides with the BV operator
4. all higher differentials are zero.
The proof of this theorem is completely straightforward; we shall do it leisurely in the rest of this note.
## Cohomology of $`\omega `$
It is fairly simple to describe the cohomology of the complex $`(\mathrm{\Omega }(M),\omega )`$ in local Darboux coordinates $`x^i`$, $`p_i`$ ($`i=1,\mathrm{},n`$, $`\omega =dp_idx^i`$), where $`x^i`$ are the even coordinates and $`p_i`$ the odd coordinates. Let $`UM`$ be the open subset covered by the coordinates. Then any cohomology class of $`(\mathrm{\Omega }(U),\omega )`$ has unique representative of the form
(1)
$$f(x,p)dx^1dx^2\mathrm{}dx^n.$$
In other words, using the coordinates, we can locally identify $`H(\mathrm{\Omega }(M),\omega )`$ with functions on $`M`$ (this identification *does* depend on the choice of coordinates).
To prove this claim we split $`(\mathrm{\Omega }(U),\omega )`$ into subcomplexes: we assign an auxiliary degree $`1`$ to each $`dx`$ and $`1`$ to each $`dp`$, and denote this degree $`\mathrm{auxdeg}`$ ($`x`$’s and $`p`$’s will have $`\mathrm{auxdeg}=0`$); the subspaces of $`\mathrm{\Omega }(U)`$ of fixed $`\mathrm{auxdeg}`$ are clearly subcomplexes. We shall see that each of them has zero cohomology, except for the one with degree $`n`$, where the differential vanishes. We shall prove it using an explicit homotopy. Let us consider the operator $`L:\mathrm{\Omega }(U)\mathrm{\Omega }(U)`$ given by
$$L:\alpha _{x^i}\mathrm{}_{p_i}\mathrm{}\alpha .$$
A direct computation shows that
$$L(\omega )+(\omega )L:\alpha (n\mathrm{auxdeg}\alpha )\alpha .$$
This concludes the proof.
Now we also see that $`d`$ is 0 on $`H(\mathrm{\Omega }(M),\omega )`$, since
$$d(f(x,p)dx^1dx^2\mathrm{}dx^n)=\frac{f}{p_k}dp_kdx^1dx^2\mathrm{}dx^n,$$
which is $`\omega `$-exact (having $`\mathrm{auxdeg}=n1`$).
## The third differential
Let us now compute the 3rd differential $`d\left(\omega \right)^1d`$ in local Darboux coordinates. We have
$$d(f(x,p)dx^1dx^2\mathrm{}dx^n)=\omega \alpha ,$$
where
$$\alpha =L(d(f(x,p)dx^1dx^2\mathrm{}dx^n))=\frac{f}{p_k}_{x^k}\mathrm{}dx^1dx^2\mathrm{}dx^n$$
and $`d\alpha `$ is (up to a $`\omega `$-exact term)
$$\frac{^2f}{x^kp_k}dx^1dx^2\mathrm{}dx^n.$$
The third differential in the spectral sequence is thus equal to the Batalin-Vilkoviski operator
$$\mathrm{\Delta }=\frac{^2}{x^kp_k}.$$
Now if $`M`$ is contractible and admits global Darboux coordinates, the cohomology of $`\mathrm{\Delta }`$ is isomorphic to $``$ (since $`\mathrm{\Delta }`$ can be identified with de Rham’s $`d`$ on a contractible subset of $`^n`$), and any cohomology class has a representative in $`\mathrm{\Omega }(M)`$ which is a constant multiple of $`dx^1dx^2\mathrm{}dx^n`$. This representative is $`d`$-closed and thus is annuled by all higher differentials in the spectral sequence. Since any $`M`$ can be covered by such patches, these higher differentials vanish for any $`M`$. This concludes the proof of the theorem, except for the part (1).
## Semidensities
Now we’ll prove the part (1) of the theorem. We locally identified the cohomology of $`H(\mathrm{\Omega }(M),\omega )`$ with functions on $`M`$ by choosing the representant (1); it is fairly easy to see that when we pass to another system of local Darboux coordinates, the function $`f`$ gets multiplied by the square root of the corresponding Berezinian. We shall however give a different proof, using Manin’s cohomological definition of Berezinian . The claim we are proving here, together with the proof, is taken from .
Let us recall Manin’s definition. Let $`V`$ be a vector supespace. Let us choose a vector superspace $`W`$ with an odd symplectic form $`\omega ^2W^{}`$, such that $`V`$ is its Lagrangian subspace (for example $`W=V\mathrm{\Pi }V^{}`$). Then $`\mathrm{Ber}(V^{})`$ (the 1-dimensional vector space of constant densities on $`V`$) is defined as the cohomology $`H(W^{},\omega )`$ (this definition is easily seen to be independent of the choice of $`W`$).
If now $`V^{}`$ is a Lagrangian complement of $`V`$ in $`W`$, then again $`\mathrm{Ber}(V^{})=H(W^{},\omega )`$; on the other hand,
$$\mathrm{Ber}(W^{})=\mathrm{Ber}(V^{})\mathrm{Ber}(V^{})=H(W^{},\omega )^2,$$
and thus $`H(W^{},\omega )=\mathrm{Ber}(W^{})^{1/2}`$ (we should write everywhere “naturally isomorphic” instead of “equal”, but hopefully it’s not a big crime). This identity is valid for any vector superspace with an odd symplectic form. We apply it to the bundle of symplectic vector spaces $`TM`$, which concludes the proof.
## Final remarks
1. We should say a few remarks about the spectral sequence of the bicomplex $`(\mathrm{\Omega }(M),\omega ,d)`$, since $`\mathrm{\Omega }(M)`$ is *not* bigraded. The spectral sequence is constructed in this way: we take $`\mathrm{\Omega }(M)[\mathrm{}]`$ (differential forms on $`M`$ depending polynomially on an indeterminate $`\mathrm{}`$; we could just as well take $`\mathrm{\Omega }(M)[[\mathrm{}]]`$), with the differential $`\mathrm{}d+\omega `$. Then multiplication by $`\mathrm{}`$ is an endomorphism of the complex $`(\mathrm{\Omega }(M)[\mathrm{}],\mathrm{}d+\omega )`$, and our spectral sequence is the Bochstein spectral sequence of this endomorphism. That is, we start with the short exact sequence of complexes
$$0(\mathrm{\Omega }(M)[\mathrm{}],\mathrm{}d+\omega )\stackrel{\mathrm{}}{}(\mathrm{\Omega }(M)[\mathrm{}],\mathrm{}d+\omega )(\mathrm{\Omega }(M),\omega )0,$$
out of which we get the exact couple
$$\begin{array}{ccccc}H(\mathrm{\Omega }(M)[\mathrm{}],\mathrm{}d+\omega )& & \stackrel{(\mathrm{})_{}}{}& & H(\mathrm{\Omega }(M)[\mathrm{}],\mathrm{}d+\omega )\\ & & & & \\ & & H(\mathrm{\Omega }(M),\omega )& & \end{array}$$
which generates the spectral sequence. If we denote $`E_{\mathrm{}}`$ its ultimate term, we have
$$H(\mathrm{\Omega }(M)[\mathrm{},\mathrm{}^1],\mathrm{}d+\omega )E_{\mathrm{}}[\mathrm{},\mathrm{}^1]$$
(and also $`H(\mathrm{\Omega }(M)[[\mathrm{}]][\mathrm{}^1],\mathrm{}d+\omega )E_{\mathrm{}}[[\mathrm{}]][\mathrm{}^1]`$).
2. The odd symplectic form $`\omega `$ on $`M`$ gives us an isomorphism between $`\mathrm{\Omega }(M)`$ and $`\mathrm{\Gamma }(STM)`$, i.e. the space of polynomial functions on $`T^{}M`$. This isomorphism transfers $`\omega `$ to multiplication by the odd Poisson structure $`\pi `$ corresponding to $`\omega `$ (recall that since $`\pi `$ is an odd Poisson structure, it is a function on $`T^{}M`$), and $`d`$ to $`\{\pi ,\}`$ (where $`\{,\}`$ is the Poisson bracket on $`T^{}M`$); the differential $`\mathrm{}d+\omega `$ becomes $`\pi +\mathrm{}\{\pi ,\}`$. This suggests some generalizations, e.g. we can take an odd Poisson structure which is not symplectic, or more generally, instead of $`T^{}M`$ we can take an aritrary supermanifold with an even symplectic form, and consider on it an odd function $`\pi `$ such that $`\{\pi ,\pi \}=0`$. The spectral sequences would still be defined, but it is not clear to me if they are good for anything.
3. The result in this note is extremely simple; it is written with the hope that it might be helpfull in situations where BV-like operators are much less trivial.
|
warning/0506/quant-ph0506050.html
|
ar5iv
|
text
|
# Simultaneous Classical-Quantum Capacities of Quantum Multiple Access Channels
##
*To my father*
Carl A. Yard Jr.
## Acknowledgments
It is hard to express the gratitude I have for my advisor, Tom Cover. I thank him for not only giving me the freedom to learn about and work on any topic of my choosing, but also for his unending encouragement and belief that I would eventually figure out how to do what I needed to do in order to graduate.
I am particularly indebted to my collaborators Igor Devetak and Patrick Hayden. The assistance and guidance that I received from them has helped me to learn an extraordinary amount about quantum information theory, a task which would have been insurmountable without their help. In fact, this dissertation is based on my collaboration with them. I would also like to give thanks to my reading committee, Yoshi Yamamoto, Abbas El Gamal and Alexander Fetter.
I thank the current members of the information theory group, Young-Han, Charles, Styrmir, George and Navid, as well as the past members Josh, David, Michael, Arak and Assaf, for years of stimulating conversations around the office and at our weekly group meetings.
Much thanks to my family for their love and support, especially Mom, who has been anticipating my graduation for quite some time. I also thank my many friends for being there for me over the years, most notably Glenn and Jeremy, who have helped me out of quite a few mathematical jams in which I have found myself over the years. I am gracious for the roommates I have had the pleasure of living with during graduate school, most notably Andy and Sophia, who lived with me during the writing of this dissertation. Thanks for dealing with me and my endless piles of paper.
I also thank Martin Morf, whose support and encouragement throughout my stay at Stanford has certainly help to keep me going. I am appreciative of Denise Murphy for rescuing me numerous times from various administrative messes in which I found myself over the years. I am also grateful to John Preskill for hosting me for a week at the Caltech Institute for Quantum Information where I worked on completing this dissertation. Finally, a million thanks to Meagen for her love and support.
## Chapter 1 Introduction
Information is embodied in physical objects. The paint on the ceiling of the Sistine chapel, the groove of a record, a single strand of DNA, and the spin of an individual electron each reflect a configuration of a particular physical system which governs the way in which it interacts with the rest of the world. Michaelangelo’s painting interacts with ambient light, emitting a spectrum of colors viewed by churchgoers and tourists alike. The spinning record induces vibrations in a needle which are converted to a flow of electrons which, when amplified, cause pressure waves to travel through the air. The structure of DNA encodes instructions for both self-replication and the construction of living beings. Interactions between individual electrons are mediated by photons and can be modeled with great precision using the tools of quantum electrodynamics. The physics of simplified models can be calculated using the rules of quantum mechanics. Quantum particles embody *quantum information*.
Claude Shannon, motivated by engineering problems in communication theory, initiated the study of information theory as an abstract discipline. On the first page of his seminal paper , he writes
> The fundamental problem of communication is that of reproducing at one point either exactly or approximately a message selected at another point. …the actual message is one *selected from* a set of possible messages.
He later states “we wish to consider certain general problems involving communication systems.” These general problems, the outgrowth from which is referred to nowadays as “Shannon theory,” concern characterizing the possibilities of reliably transmitting certain “information sources” over “information-bearing channels.” By making probabilistic assumptions regarding the behaviors of the sources and channels, a rich mathematical theory emerges which, in many cases, reasonably approximates the underlying physics. Specifically, Shannon showed that if a channel is modeled by a probability transition matrix $`p(y|x)`$, its capacity for the transmission of classical information is given by
$$C=\underset{p(x)}{\mathrm{max}}I(X;Y).$$
This formula will be discussed in Section 3.3, where we introduce the mutual information $`I(X;Y)`$, and in Section 4.1 where we discuss the proof of Shannon’s theorem. Shannon theory applies to network communication as well. A probability transition matrix $`p(z|x,y)`$ models a situation where two senders transmit to a single receiver, subject to noise and interference. The rates at which the senders can transmit independent information were determined by Ahlswede and Liao to admit a single-letter characterization, given by the convex hull of the closure of the set of pairs of nonnegative rates $`(R_X,R_Y)`$ satisfying
$`R_X`$ $`<`$ $`I(X;Z|Y)`$
$`R_Y`$ $`<`$ $`I(Y;Z|X)`$
$`R_X+R_Y`$ $`<`$ $`I(XY;Z)`$
for some $`p(x)p(y).`$ Further analysis by Cover, El Gamal and Salehi gives single-letter characterizations of a set of correlated sources which can be reliably transmitted over a multiple access channel, generalizing the above, as well as Slepian-Wolf source coding and cooperative multiple access channel capacity. They also give a multi-letter expression for the capacity region, showing that an i.i.d. source $`(U,V)`$ can be reliably transmitted if and only if
$`H(U|V)`$ $`<`$ $`{\displaystyle \frac{1}{n}}I(X^n;Z^n|U^nY^n)`$
$`H(V|U)`$ $`<`$ $`{\displaystyle \frac{1}{n}}I(Y^n;Z^n|V^nX^n)`$
$`H(U,V)`$ $`<`$ $`{\displaystyle \frac{1}{n}}I(X^nY^n;Z^n)`$
for some $`n`$ and $`p(x^n|u^n),p(y^n|v^n)`$, where by $`x^n`$ we mean the sequence of symbols $`(x_1,\mathrm{},x_n)`$. Here, $`H(U,V)`$ and $`H(U|V)`$ respectively denote the *entropy* and *conditional entropy* of the pair of random variables $`(U,V)`$ which model the source. Such a characterization is of limited practical use, however, as it does not apparently lead to a finite computation for deciding whether or not a source can be transmitted.
The concept of the entropy of a physical system initially arose out of attempts to characterize the optimal efficiency of physical machines such as steam engines, as well as to rule out the possibility of such constructions as perpetual motion machines. The extensivity of entropy demands that it be additive for independent physical systems. Boltzmann defined the entropy of a physical system to be proportional to the logarithm of the number of microstates, or indistinguishable configurations of its constituents, a definition he was likely led to because $`\mathrm{log}(W_1W_2)=\mathrm{log}(W_1)+\mathrm{log}(W_2)`$, making additivity of entropy evident. In order to circumvent mathematical subtleties which arise due to course grainings of the system’s configuration space, a probabilistic approach can be taken, allowing rigorous mathematical statements to be made about related systems which are essentially hidden Markov models . A crucial philosophical step was taken by Boltzmann in his work; he assumed that things were made of atoms. In his framework, heat was not a fluid that flowed from warm to cold bodies; rather, vibrational energy of the constituents of a physical system induces similar behavior in neighboring systems. Without direct physical evidence to support the existence of atoms, Boltzmann provided a mechanism for the flow of heat which assumed such ingredients did in fact exist. The existence of atoms was experimentally verified soon after Boltzmann’s untimely death.
In the years that followed, the structure of atoms was intensely investigated. The assumption that atoms obey the laws of Newtonian mechanics quickly resulted in various logical inconsistencies in the form of predictions which did not agree with experimental results. Quantum mechanics was reluctantly developed as a collection of fundamental assumptions about the nature of the physics of atoms and their constituents. A collection of mathematical rules was thus constructed which allowed theoretical calculations of certain aspects of experimental results. A caveat was that the new theory introduced randomness as a fundamental assumption of the theory, a feature which was quite unsettling to even the creators of quantum mechanics, most notably Einstein, who thought that “God does not play dice with the universe.”
Today, we live in a quantum world. Progress in applied physics and engineering has begun to make manipulation of matter on the quantum scale a reality. It is a strange world, at least when viewed with a classical mind. From the other side of the fence, however, classical physics can be seen to be part of a quantum world. The emergence of classicality due to phase transitions in systems of many particles is one way this occurs. Mathematically, as we will see in Sections 2.1.7 and 2.2.3, the tools and language of quantum theory enable the expression of concepts from classical probability theory.
This opens up the possibility of analyzing communication scenarios in which the senders and receivers process quantum information. In this case, the medium quite literally is the message, whereas rather than sending information by selecting a message from a set, physical systems are suitably prepared for transmission to a receiver. The possible types of quantum communication range from transmitting particles from sender to receiver to generating entanglement between the users of a channel. Quite remarkably, certain basic components from the classical theory find a place in the quantum extension. The techniques used to separate possible quantum information processing tasks from the impossible are directly analogous to those used Shannon’s in original program. Possibility questions of this nature have much in common with the original motivations of thermodynamics. The ways in which entropy arises in characterizing the answers further deepens this connection.
In this sense, one aspect of quantum information theory involves generalizing existing classical results to include quantum resources. While network Shannon theory is already a rich active area of research, its quantum extension has new aspects which do not fit into the former framework. This leads to a theory which, while including the old one as a special case, asks new questions leading to a deeper understanding of the physical nature of information. Apparently, quantum information is something new which cannot be properly analyzed with classical tools alone.
In this dissertation, we will analyze quantum channels with many senders and a single receiver, used in a variety of ways for the simultaneous transmission of classical and quantum information, representing an expanded version of the manuscript . At a high level, the results and approaches contained within mirror those of classical Shannon theory. Yet, the mathematical tools utilized are distinctly quantum. Let us end this introduction by giving a quote from Asher Peres and Daniel Terno’s paper on quantum information and relativity, where it is written that “the goals of quantum information theory are the intersection of those of quantum mechanics and information theory, while its tools are the union of these two theories.” Well said.
## Chapter 2 Background
### 2.1 The basics
#### 2.1.1 Quantum mechanics
Let us briefly review some elements of quantum mechanics. The physical state of an isolated system with $`d`$ quantum degrees of freedom is described by a complex unit vector $`|\psi ^d`$. The notation $`|\psi `$, known as a *ket* or *ket vector*, refers to a normalized column vector with $`d`$ complex components:
$$|\psi =\left(\begin{array}{c}\alpha _1\\ \alpha _2\\ \mathrm{}\\ \alpha _d\end{array}\right),\text{where}\underset{a}{}|\alpha _a|^2=1.$$
The conjugate transpose of $`|\psi `$ is a row vector
$$\psi |=\left(\begin{array}{cccc}\alpha _1^{}& \alpha _2^{}& \mathrm{}& \alpha _d^{}\end{array}\right).$$
$`\psi |`$ is called a *bra* or *bra vector*. This notation (and nomenclature) was introduced by Dirac, partly to emphasize the inner product structure of $`^d`$. Indeed, the inner product between two state vectors $`|\varphi `$ and $`|\psi `$ is written as a bra times a ket, or bra-ket
$$\varphi |\psi =\psi |\varphi ^{}.$$
It is often useful to write a basis for $`^d`$ by defining a collection of kets as
$$|1=\left(\begin{array}{c}1\\ 0\\ \mathrm{}\\ 0\end{array}\right),|2=\left(\begin{array}{c}0\\ 1\\ \mathrm{}\\ 0\end{array}\right),\mathrm{},|d=\left(\begin{array}{c}0\\ 0\\ \mathrm{}\\ 1\end{array}\right).$$
Then, a state such as $`|\psi `$ can be expanded in terms of this basis as
$$|\psi =\alpha _1|1+\mathrm{}+\alpha _d|d.$$
A measurement can be performed on the quantum system, obtaining classical information regarding the system’s current quantum state. Quantum mechanics is only able to predict the probabilities of occurrence for each outcome of the measurement. Further, the state of the system will generally be disturbed by the act of obtaining this classical inforamation. The simplest measurement to describe is a *pure state measurement*, which is completely described in terms of some orthogonal basis for $`^d`$. Such a basis will be called a *measurement basis*. Suppose that a pure state measurement in the measurement basis $`\{|1,\mathrm{},|d\}`$ is made on the state $`|\psi `$. Then,
* The measurement will return $`y`$ with probability
$$p(y)\mathrm{Pr}\{\text{measure }|y\}=|y|\psi |^2.$$
* If the measurement returns $`y`$, the post-measurement state is then $`|y`$.
In other words, the measurement result is modeled by a $`𝒴=\{1,\mathrm{},d\}`$-valued random variable $`Y`$, distributed as $`p(y)=|y|\psi |^2=|\alpha _y|^2`$, and the post-measurement state is a random vector $`|Y`$. If the same measurement is performed again, the same result $`Y`$ is obtained with certainty, leaving the system in the same state $`|Y`$ after the measurement.
#### 2.1.2 Pure state ensembles
Here, let us fix a basis $`\{|y\}_{y=1}^d`$ for $`^d`$. Imagine now a game with two parties, Alice and Bob. Assume that Alice has the ability to prepare any pure state from the finite collection of states $`\{|\psi _x\}_{x𝒳}`$. Then, the probability that Bob obtains measurement result $`y`$ given that Alice prepares state $`|\psi _x`$ is given by
$$p(y|x)=|y|\psi _x|^2.$$
(2.1)
Notice that this can be rewritten as
$`|y|\psi _x|^2`$ $`=`$ $`y|\psi _x(y|\psi _x)^{}`$ (2.2)
$`=`$ $`y|\psi _x\psi _x|y`$
$``$ $`y|\left(|\psi _x\psi _x|\right)|y.`$
We may interpret this as saying that if the 1-dimensional projection matrix $`|\psi _x\psi _x|`$ is written in the $`\{|y\}`$ basis, then $`p(y|x)`$ is equal to the diagonal matrix element corresponding to $`|y`$.
Now, suppose that Alice gives Bob a random state, choosing $`|\psi _x`$ with probability $`p(x)`$. In this case, we say that Alice is preparing an *ensemble* $`\{p(x),|\psi _x\}`$ of pure states. Together with elementary probability, (2.2) can be used to write the probability that Bob measures $`y`$ as
$`p(y)`$ $`=`$ $`{\displaystyle \underset{x}{}}p(x)p(y|x)`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)y||\psi _x\psi _x||y`$
$`=`$ $`y|\left({\displaystyle \underset{x}{}}p(x)|\psi _x\psi _x|\right)|y`$
$``$ $`y|\rho |y.`$
where the third line is by linearity. The fourth line defines the *density matrix*
$$\rho =\underset{x}{}p(x)|\psi _x\psi _x|$$
of the ensemble $`\{p(x),|\psi _x\}`$, which contains all the data required to compute all probabilities associated with any possible measurement on the ensemble, under the assumption that Bob doesn’t know the identities of the individual states. Note that $`\rho `$ is Hermitian
$$\rho ^{}=\underset{x}{}p(x)\left(|\psi _x\psi _x|\right)^{}=\underset{x}{}p(x)|\psi _x\psi _x|=\rho $$
and satisfies
$$\mathrm{Tr}\rho =\underset{x}{}p(x)\mathrm{Tr}|\psi _x\psi _x|=\underset{x}{}p(x)=1.$$
$`\rho `$, as we’ve constructed it, is also nonnegative definite. This is because for any $`|\varphi `$, we have
$$\varphi |\rho |\varphi =\underset{x}{}p(x)\varphi ||\psi _x\psi _x||\varphi =\underset{x}{}p(x)|\varphi |\psi _x|^20$$
where the last inequality is because each term in the sum is nonnegative.
#### 2.1.3 Density matrices
We have now seen that if a quantum system is prepared in a random pure state, one can write down its density matrix. This contains all of the data necessary to compute the probabilities of the outcomes of any measurement that can be made on that system, provided that the identities of the random pure states are unknown to the measurer. For a system in a pure state $`|\psi `$, we will use the abbreviation $`\psi |\psi \psi |`$ for the density matrix corresponding to that pure state (this is just the matrix which projects onto the subspace spanned by $`|\psi `$. Let us define here the collection of all density matrices of a $`d`$-level quantum system as
$$𝒟^d=\{\rho ^{d\times d}:\rho =\rho ^{},\rho 0,\mathrm{Tr}\rho =1\}.$$
In other words, a density matrix $`\rho 𝒟^d`$ is a Hermitian, nonnegative definite normalized matrix. We give the following facts about $`𝒟^d`$ without proof, as they are proven in detail in many texts on quantum mechanics :
###### Property.
$`𝒟^d`$ is convex.
###### Property.
The extremal points of $`𝒟^d`$ are the projections onto rank 1 subspaces of $`^d`$, corresponding to equivalence classes of pure states which are identified up to a global phase factor $`e^{i\theta }`$.
###### Property.
$`𝒟^d`$ is compact.
We may interpret the first fact as saying that if with probability $`p`$, one chooses to prepare a quantum system so that its density matrix is $`\rho `$, while with probability $`1p`$, it is instead prepared so that its density matrix is $`\sigma `$, someone who measures the resulting system (and is also ignorant about which preparation was made) computes measurement probabilities with the state $`p\rho +(1p)\sigma .`$
The second fact illustrates the fact that every density matrix can arise from some pure state ensemble. This can be seen more directly, since the Hermiticity of $`\rho `$ implies that it is diagonalizable as
$$\rho =\underset{i}{}\lambda _i|ii|$$
for some orthogonal basis $`\{|i\}`$ for $`^d`$. The positivity of $`\rho `$ implies that $`\lambda _i0`$, and the fact that $`\rho `$ is normalized implies that the $`\lambda _i`$ may be interpreted as probabilities, implying the existence of the required pure state ensemble. Note that there is in fact an uncountable number of ways in which a density matrix can arise by probabilistically preparing pure states.
More importantly, the fact that the extremal points of $`𝒟^d`$ are pure states implies that pure states are special, in that they cannot arise as nontrivial probabilistic preparations of other states. A quantum system in a pure state is in a *definite* state.
#### 2.1.4 Trace norm
For an arbitrary $`M^{d\times d},`$ its *trace norm* $`|M|_1`$ is defined as
$$|M|_1=\mathrm{Tr}\sqrt{MM^{}}.$$
This is easily seen to be equal to the sum of the singular values of $`M`$. Indeed, writing a singular value decomposition $`M=U\mathrm{\Lambda }V^{}`$, it follows that
$$|M|_1=\mathrm{Tr}\sqrt{U\mathrm{\Lambda }V^{}V\mathrm{\Lambda }U^{}}=\mathrm{Tr}U\sqrt{\mathrm{\Lambda }^2}U^{}=\underset{i}{\overset{d}{}}|\lambda _i|,$$
where $`\mathrm{\Lambda }=\text{diag}(\lambda _1,\mathrm{},\lambda _d)`$. As $`||_1`$ is a norm (or rather, a *unitarily invariant matrix norm*), it satisfies the following properties:
###### Property (Positivity).
$`|M|_10`$, while $`|M|_1=0`$ if and only if $`M=0`$.
###### Property (Homogeneity).
for any $`c`$, $`|cM|_1=|c||M|_1`$
###### Property (Unitary invariance).
$`|M|_1=|UMU^{}|_1`$ for any unitary $`U`$
###### Property (Triangle inequality).
$`|M+N|_1|M|_1+|N|_1`$
###### Property (Submultiplicativity).
$`|MN|_1|M|_1|N_1|`$
Positivity follows because the singular values of any matrix $`M`$ are always nonnegative, and are all equal to zero if and only if $`M=0`$. Homogeneity is true because the singular values of $`cM`$ equal $`|c|`$ times those of $`M`$, and unitary invariance holds because $`UMU^{}`$ and $`M`$ have the same singular values. For proofs of the triangle inequality and submultiplicativity, the reader is referred to .
The trace norm gives a natural metric space structure to $`^{d\times d}`$ which we will exploit considerably throughout this dissertation. Given two matrices $`M,N^{d\times d}`$ their *trace distance* is thus defined as the trace norm of their difference $`|MN|_1`$. For two density matrices $`\rho `$ and $`\sigma `$ of a $`d`$-level quantum system, their trace distance satisfies
$$0|\rho \sigma |_12,$$
where the lower bound is saturated if and only $`\rho =\sigma `$, while the upper bound is saturated if and only if $`\rho `$ and $`\sigma `$ are supported on orthogonal subspaces. Let us mention here the following alternative characterization of the trace distance between two density matrices
$$|\rho \sigma |_1=2\underset{0\mathrm{\Lambda }1}{\mathrm{max}}\mathrm{Tr}\mathrm{\Lambda }(\rho \sigma ).$$
The maximization above is over all nonnegative definite matrices $`\mathrm{\Lambda }`$ with spectrum bounded above by 1.
#### 2.1.5 Fidelity
Given two density matrices $`\rho `$ and $`\sigma `$ of a $`d`$-level system, their *fidelity* is defined <sup>1</sup><sup>1</sup>1Note that many authors (such as ) define this quantity as the square root of our definition. as
$$F(\rho ,\sigma )=\left(\mathrm{Tr}\sqrt{\sqrt{\rho }\sigma \sqrt{\rho }}\right)^2.$$
Fidelity can be expressed in terms of the trace norm as
$$F(\rho ,\sigma )=\left|\sqrt{\rho }\sqrt{\sigma }\right|_1^2,$$
a form which makes apparent the symmetry of fidelity in its two arguments. The following bounds are always satisfied whenever the arguments are density matrices
$$0F(\rho ,\sigma )1.$$
The lower bound is saturated if and only if $`\rho `$ and $`\sigma `$ have orthogonal support, while the upper bound is saturated if and only if $`\rho =\sigma `$. Contrary to the situation with the trace norm, a large value of the fidelity between two states signifies that they are close. Fidelity is not a norm, but it can be related to the trace norm in various ways which are summarized in Section 6.1.
If one of the arguments of the fidelity is a pure state, (say $`\rho =\varphi ),`$ then
$`F(|\varphi ,\sigma )`$ $`=`$ $`\left(\mathrm{Tr}\sqrt{|\varphi \varphi |\sigma |\varphi \varphi |}\right)^2`$
$`=`$ $`\left(\mathrm{Tr}|\varphi \varphi |\sqrt{\varphi |\sigma |\varphi }\right)^2`$
$`=`$ $`\varphi |\sigma |\varphi .`$
So $`F(|\varphi ,\sigma )`$ is just the diagonal matrix element of $`\sigma `$ corresponding to $`|\varphi `$, when $`\sigma `$ is written in a basis including $`|\varphi `$. Note that this is the success probability for a pure state measurement which tests a system prepared in the state $`\sigma `$ for the presence of the state $`|\varphi `$. When both arguments are pure states, we obtain
$$F(|\varphi ,|\psi )=|\varphi |\psi |^2.$$
Finally observe the following easily verifiable property.
###### Property (Linearity of fidelity).
Fidelity is linear in each argument, i.e.
$$F(c\rho ,\sigma )=cF(\rho ,\sigma )=F(\rho ,c\sigma ).$$
#### 2.1.6 POVMs
We describe here a certain general type of measurement which can be performed on a $`d`$-level quantum system, called a positive operator valued measurement (POVM). A POVM is specified in terms of a finite collection of matrices $`\{\mathrm{\Lambda }_x^{d\times d}\}_{x𝒳}`$ which are positive ($`\mathrm{\Lambda }_x0`$) and sum to the $`d\times d`$ identity matrix $`1_d`$
$$\underset{y}{}\mathrm{\Lambda }_y=1_d.$$
It is often said that the matrices $`\{\mathrm{\Lambda }_x\}_{x𝒳}`$ form a *partition of unity*. If the quantum system is in the state $`\rho `$, the probability of obtaining a measurement result $`y`$ is given by
$$p(y)=\mathrm{Pr}\{\text{measure }\mathrm{\Lambda }_y\}=\mathrm{Tr}\mathrm{\Lambda }_y\rho .$$
Conditioned on having received the measurement result $`y`$, the post-measurement state after such a measurement is computed as
$$\rho \rho _y=\frac{\sqrt{\mathrm{\Lambda }_y}\rho \sqrt{\mathrm{\Lambda }_y}}{p(y)}.$$
Here, $`\sqrt{\mathrm{\Lambda }}`$ is defined as the unique, positive operator which satisfies $`\sqrt{\mathrm{\Lambda }}\sqrt{\mathrm{\Lambda }}=\mathrm{\Lambda }`$. <sup>2</sup><sup>2</sup>2Note that some authors use a more general kind of measurement, described by matrices $`\{M_y\}`$ satisfying $`_yM_y^{}M_y=1_d`$. This amounts to choosing a different square root of each $`\{\mathrm{\Lambda }_y\}`$, giving post-measurement states which are unitarily equivalent to those of the convention above, conditioned on the measurement result. Such measurements can be modeled using the tools introduced in Section 2.3.8. The measurement results in an ensemble of density matrices $`\{p(y),\rho _y\}`$. A pure state measurement in the basis $`\{|x\}`$ can be expressed as the POVM $`\{|xx|\}`$ consisting of 1-dimensional projection matrices.
#### 2.1.7 Classical systems
Let $`𝒳`$ be a finite set and let $`X`$ be an $`𝒳`$-valued random variable, distributed according to $`p(x)`$. We can define a vector space $`^{|𝒳|}`$ with a fixed orthonormal basis $`\{|x^X\}_{x𝒳}`$, labeled by elements of the set $`𝒳`$. This sets up an identification $`|^X:𝒳^{|𝒳|}`$ between the elements of $`𝒳`$ and that particular basis. By this correspondence, the probability mass function $`p(x)`$ can be mapped to a density matrix
$$\rho =\underset{x𝒳}{}p(x)|xx|$$
(2.3)
which is diagonal in the basis $`\{|x\}_{x𝒳}.`$ Further, to every subset $`S𝒳`$ corresponds a projection matrix $`\mathrm{\Pi }_S=_{xS}|xx|`$ which commutes with $`\rho `$. In addition, the projections $`\mathrm{\Pi }_S`$ and $`\mathrm{\Pi }_T`$ corresponding to any two subsets $`S,T𝒳`$ commute. This way, we can express concepts from classical probability theory in the language of quantum probability. Consider the following translations from classical to quantum language:
$`\mathrm{Pr}\{XS\}`$ $`=`$ $`\mathrm{Tr}\rho \mathrm{\Pi }_S`$
$`\mathrm{Pr}\{XS\}`$ $`=`$ $`1\mathrm{Tr}\rho \mathrm{\Pi }_S=\mathrm{Tr}\rho (1^X\mathrm{\Pi }_S)`$
$``$ $`\mathrm{Tr}\rho \mathrm{\Pi }_{S^c}`$
$`\mathrm{Pr}\{XS\text{and}XT\}`$ $`=`$ $`\mathrm{Tr}\rho \mathrm{\Pi }_S\mathrm{\Pi }_T`$
$``$ $`\mathrm{Tr}\rho \mathrm{\Pi }_{ST}`$
$`\mathrm{Pr}\{XS\text{or}XT\}`$ $`=`$ $`1\mathrm{Pr}\{XS\text{and}XT\}`$
$`=`$ $`\mathrm{Tr}\rho \left(1^X(1^X\mathrm{\Pi }_S)(1^X\mathrm{\Pi }_T)\right)`$
$``$ $`\mathrm{Tr}\rho \mathrm{\Pi }_{ST}.`$
From the early development of quantum mechanics, noncommutativity has been seen to be the hallmark of quantum behavior. It is to be expected that classical probability, embedded in quantum theory’s framework, is described entirely with commuting matrices.
### 2.2 Composite quantum systems
Let us begin by introducing a number of conventions which will be used when dealing with multiple quantum systems. We will use capital letters from the beginning of the alphabet $`A,B,C,\mathrm{}`$ as labels for quantum systems. If $`A`$ is a quantum system, we will abbreviate its level as $`|A|`$ (which will *always* be finite), so that its pure states are unit vectors in $`^{|A|}`$. A generic pure state of $`A`$ will then be written as $`|\psi ^A`$, while a generic density matrix of $`A`$ will be written $`\rho ^A`$, to remind the reader to which system the state refers. Whenever we initially introduce a state, the superscript will identify the system it is describing, although later references to that state will not always include the superscript. This convention will not be cause for confusion, as different symbols will refer to different states. We will also write the $`|A|\times |A|`$ identity matrix on $`^{|A|}`$ as $`1^A`$.
If $`B`$ is another quantum systems, then $`A`$ and $`B`$ may be combined to form a *composite quantum system* $`AB`$. This new system has $`|A||B||AB|`$ levels. The pure states of the new system are instead unit vectors in the *tensor product* $`^{|A|}^{|B|}`$ vector space of the individual vector spaces. The simplest way to define $`^{|A|}^{|B|}`$ is as follows. First, fix arbitrary bases $`\{|a^A\}_{a=1}^{|A|}`$ and $`\{|b^B\}_{b=1}^{|B|}`$ for $`^{|A|}`$ and $`^{|B|}`$. Then, $`^{|A|}^{|B|}`$ can be formally defined as the linear span of the basis vectors formed by the product of the two individual bases
$$\left\{|a^A|b^B\right\}_{a,b=1}^{|A|,|B|}.$$
A convenient shorthand for the tensor product of pure states is to write
$$|a^A|b^B|a^A|b^B.$$
Then, any pure state of the quantum system can be written as
$$|\mathrm{\Psi }^{AB}=\underset{a=1}{\overset{|A|}{}}\underset{b=1}{\overset{|B|}{}}c_{ab}|a^A|b^B.$$
(2.4)
Observe that this new vector space we have constructed has dimension $`|A||B|`$. It is not difficult to show that this construction is universal, meaning that it is independent of the particular bases chosen for $`A`$ and for $`B`$.
It will be useful here to describe a certain convention which can be used to write down the tensor product of two column vectors as a single column vector. This will amount to fixing a way to enumerate the components of the tensor. Suppose that $`\stackrel{}{v}^{|A|}`$ and $`\stackrel{}{w}^{|B|}`$ are arbitrary column vectors
$$\stackrel{}{v}=\left(\begin{array}{c}v_1\\ v_2\\ \mathrm{}\\ v_{|A|}\end{array}\right)\text{ and }\stackrel{}{w}=\left(\begin{array}{c}w_1\\ w_2\\ \mathrm{}\\ v_{|B|}\end{array}\right).$$
As $`^{|A|}^{|B|}^{|A||B|}`$, we can “flatten” $`\stackrel{}{v}\stackrel{}{w}`$ into a single column vector, organizing its components according to the following convention
$$\text{flatten}:\stackrel{}{v}\stackrel{}{w}\left(\begin{array}{c}v_1\stackrel{}{w}\\ v_2\stackrel{}{w}\\ \mathrm{}\\ v_{|A|}\stackrel{}{w}\end{array}\right).$$
In this way, the earlier generic state (2.4) can be expressed as
$$\text{flatten}:|\mathrm{\Psi }^{AB}\left(\begin{array}{c}c_{11}\\ \mathrm{}\\ c_{1|B|}\\ c_{21}\\ \mathrm{}\\ c_{2|B|}\\ \mathrm{}\\ c_{|A||B|}\end{array}\right)$$
It is often the case that a pure state such as $`|\mathrm{\Psi }^{AB}`$ cannot be written as a tensor product of pure states of its constituent systems, i.e.
$$|\mathrm{\Psi }^{AB}|\psi ^A|\varphi ^B$$
for any pure states $`|\psi ^A`$ and $`|\psi ^B`$. If this is the case, then $`|\mathrm{\Psi }^{AB}`$ is said to be *entangled*. Nevertheless, for any pure state of the composite quantum system, there exists a pair of orthonormal bases $`\{|i^A\}`$ and $`\{|i^B\}`$ such that
$$|\mathrm{\Psi }^{AB}=\underset{i}{}d_i|i^A|i^B.$$
This form is called the *Schmidt decomposition* of $`|\mathrm{\Psi }^{AB}`$. Together, the combination of the orthonormal bases $`\{|i^A|i^B\}`$ is called the *Schmidt basis*, while the $`\{d_i\}`$ are called the *Schmidt coefficients*. These are easily calculated from the singular value decomposition of the matrix $`[c_{i,j}]`$ of coefficients in (2.4), where the Schmidt basis consists of the left and right eigenvectors, while the Schmidt coefficients are the singular values themselves.
Just as the tensor product builds larger vector spaces out of pairs of smaller ones, it also builds larger matrices from pairs of smaller ones. Fix two matrices $`M^{|C|\times |A|}`$ and $`N^{|D|\times |B|}`$. Recall that these are linear operators
$$M:^{|A|}^{|C|}\text{and}N:^{|B|}^{|D|}.$$
Their tensor product $`MN`$ is another linear operator
$$(MN):^{|A|}^{|B|}^{|C|}^{|D|}.$$
We will abbreviate this by writing
$$M:AC,N:BD\text{ and }(MN):ABCD.$$
This new object acts on the tensor product of vectors as
$$(MN)(|\psi ^A|\varphi ^B)=(M|\psi ^A)(N|\varphi ^B)$$
and linearity defines the action of $`MN`$ on all of $`^{|A|}^{|B|}`$. The tensor product is also *bilinear*, i.e. for any $`c`$,
$$c(MN)=(cM)N=M(cN).$$
In the same vein as the “flattened” representation $`^{|A|}^{|B|}^{|A||B|}`$ for the tensor product of vectors, there is more general mapping $`^{|A|\times |C|}^{|B||D|}^{|A||B|\times |C||D|}`$ given by
$$\text{flatten}:MN\left(\begin{array}{cccc}m_{11}N& m_{12}N& \mathrm{}& m_{1|C|}N\\ m_{21}N& m_{22}N& & \\ \mathrm{}& & \mathrm{}& \\ m_{|A|1}N& & & m_{|A||C|}N\end{array}\right).$$
Note our convention, where the blocks are labelled by elements of the left-most component of the tensor product. We will use that convention throughout this dissertation. It is easy to see that calculations can be made in this representation, namely that
$$\text{flatten}\{MN\}\text{flatten}\{|\psi |\varphi \}=\text{flatten}\{(MN)|\psi |\varphi \}.$$
As the composite system $`AB`$ is a quantum system itself, it includes a (strictly) larger collection of von Neumann measurements and unitary evolutions. Indeed, given any two bases $`\{|i^{AB}\}`$ and $`\{|i^{}^{AB}\}`$ for $`^{|A|}^{|B|}`$, they are related by a particular unitary matrix $`U`$, defined as
$$U=\underset{i^{}i}{}|i^{}i|.$$
It is not hard to see that any joint von Neumann measurement on the combined system $`AB`$ can be performed using separate product measurements on $`A`$ and $`B`$, provided that the unitary which takes intended measurement basis to the required product basis (and its inverse) are implementable.
Of particular interest is the subject of *local measurements* on a composite quantum system. Suppose that a measurement $`\{\mathrm{\Lambda }_x\}_{x𝒳}`$ is made on the $`A`$ part of the bipartite state $`\rho ^{AB}`$. New measurement operators $`\{\mathrm{\Lambda }_x1^B\}_{x𝒳}`$ can be constructed, so that
$$p(x)=\mathrm{Tr}\rho (\mathrm{\Lambda }_x1^B).$$
The post-measurement states are given as before
$$\rho \rho _x=\frac{\left(\sqrt{\mathrm{\Lambda }_x}1^B\right)\rho \left(\sqrt{\mathrm{\Lambda }_x}1^B\right)}{p(x)}.$$
It is instructive to see what happens if a local pure state measurement is made on part of a bipartite pure state $`|\mathrm{\Psi }^{AB}`$. Here, $`\mathrm{\Lambda }_x=|xx|`$, and we obtain
$$p(x)=\mathrm{Tr}\rho (|xx|1^B).$$
As a first step, express $`|\mathrm{\Psi }^{AB}`$ in terms of the new basis for $`A`$ as
$$|\mathrm{\Psi }^{AB}=\underset{xb}{}d_{xb}|x^A|b^b.$$
(2.5)
Note that the new coefficients $`d_{xb}`$ are related to the old ones via
$$\underset{a}{}U_{xa}c_{ab}=d_{xb},$$
where $`U:\{|a\}\{|x\}`$ is the unitary change of basis matrix. Then,
$`|\mathrm{\Psi }^{AB}`$ $`=`$ $`{\displaystyle \underset{xb}{}}d_{xb}|x^A|b^B`$
$`=`$ $`{\displaystyle \underset{x}{}}|x^A\left({\displaystyle \underset{b}{}}d_{xb}|b^B\right)`$
$``$ $`{\displaystyle \underset{x}{}}|x^A|\stackrel{~}{\psi }_x^B`$
$``$ $`{\displaystyle \underset{x}{}}\beta _x|x^A|\psi _x^B.`$
The third step above defines the unnormalized vector $`|\stackrel{~}{\psi }_x^B`$, where in the last, the normalization constant $`\beta _x\sqrt{\stackrel{~}{\psi }_x|\stackrel{~}{\psi }_x}=\sqrt{_b|d_{xb}|^2}`$ and normalized state $`|\psi _x\beta _x^1|\stackrel{~}{\psi }_x`$ are defined. Now, it is a simple task to compute
$`p(x)`$ $`=`$ $`\mathrm{Tr}\left(|xx|1^B\right)\mathrm{\Psi }^{AB}`$
$`=`$ $`\mathrm{\Psi }|^{AB}\left(|xx|1^B\right)|\mathrm{\Psi }^{AB}`$
$`=`$ $`\left({\displaystyle \underset{x^{\prime \prime }}{}}\beta _{x^{\prime \prime }}^{}x^{\prime \prime }|^A\psi _{x^{\prime \prime }}|^B\right)\left(|xx|1^B\right)\left({\displaystyle \underset{x^{}}{}}\beta _x^{}|x^{}^A|\psi _x^{}^B\right)`$
$`=`$ $`{\displaystyle \underset{x^{\prime \prime }x^{}}{}}\beta _{x^{\prime \prime }}^{}\beta _x^{}x^{\prime \prime }|xx|x^{}\psi _{x^{\prime \prime }}|\psi _x^{}`$
$`=`$ $`|\beta _x|^2.`$
Conditioned on having received the measurement result $`x`$, the post-measurement state is
$`\mathrm{\Psi }_x^{AB}`$ $`=`$ $`{\displaystyle \frac{\left(|xx|1^B\right)\mathrm{\Psi }^{AB}\left(|xx|1^B\right)}{p(x)}}`$
$`=`$ $`{\displaystyle \frac{\left(|xx|1^B\right)\left(_{x^{\prime \prime }x^{}}\beta _{x^{\prime \prime }}\beta _x^{}^{}|x^{\prime \prime }x^{}||\psi _{x^{\prime \prime }}\psi _x^{}|\right)\left(|xx|1^B\right)}{|\beta _x|^2}}`$
$`=`$ $`|xx||\psi _x\psi _x|.`$
Or rather,
$$|\mathrm{\Psi }_x^{AB}=|x^A|\psi _x^B.$$
So, we see that a measurement on $`A`$ causes the state of $`B`$ to “collapse” as well. Rather, we see that the measurement on $`A`$ creates a pure state ensemble $`\{p(x),|\psi _x^B\}`$ on $`B`$. If an arbitrary POVM is performed on $`A`$, an ensemble of density matrices on $`B`$ will generally result. To see this, we need to introduce the partial trace.
#### 2.2.1 Partial trace
If we are instead concerned only with the measurement probabilities, and not with the post-measurement states, it is convenient to work with a density matrix on $`A`$ to compute the measurement probabilities. This density matrix is defined in terms of the *partial trace* over $`B`$. Fixing a bipartite density matrix $`\mathrm{\Omega }^{AB}`$, the partial trace over $`B`$ of $`\mathrm{\Omega }^{AB}`$ can be defined as the unique density matrix $`\mathrm{Tr}_B\mathrm{\Omega }`$ on $`A`$ such that for every $`M^{|A|\times |A|}`$,
$$\mathrm{Tr}M(\mathrm{Tr}_B\mathrm{\Omega })\mathrm{Tr}(M1^B)\mathrm{\Omega }.$$
An equivalent way to define $`\mathrm{Tr}_B\mathrm{\Omega }`$ is as follows. If we write $`\mathrm{\Omega }_{a^{}ab^{}b}a^{}|b^{}|\mathrm{\Omega }|b|a`$ and $`(\mathrm{Tr}_B\mathrm{\Omega })_{a^{}a}a^{}|(\mathrm{Tr}_B\mathrm{\Omega })|a`$, then
$$(\mathrm{Tr}_B\mathrm{\Omega })_{a^{}a}=\underset{b}{}\mathrm{\Omega }_{aa^{}bb}.$$
With this in hand, we can express
$`\mathrm{Tr}\left(\mathrm{\Lambda }_x1^B\right)\mathrm{\Omega }=\mathrm{Tr}\mathrm{\Lambda }_x(\mathrm{Tr}_B\mathrm{\Omega }).`$
For any square matrix $`M`$ on $`AB`$, the partial traces over $`A`$ and $`B`$ satisfy the following easily verifyable properties:
$`\mathrm{Tr}M=\mathrm{Tr}_{AB}M=\mathrm{Tr}_A\mathrm{Tr}_BM=\mathrm{Tr}_B\mathrm{Tr}_AM.`$
A perhaps more concrete definition of the partial trace is obtained by writing a bipartite density matrix in the flattened representation
$$\mathrm{\Omega }^{AB}=\left(\begin{array}{ccc}\omega _{11}& \mathrm{}& \omega _{1|A|}\\ \mathrm{}& \mathrm{}& \\ \omega _{|A|1}& & \omega _{|A||A|}\end{array}\right)$$
where each $`\omega _{aa^{}}^{|B|\times |B|}`$. Then, $`\mathrm{Tr}_A\mathrm{\Omega }^{AB}`$ is obtained by summing the blocks on the diagonal
$$\mathrm{Tr}_A\mathrm{\Omega }^{AB}=\underset{a}{}\omega _{aa}$$
and $`\mathrm{Tr}_B\mathrm{\Omega }^{AB}`$ by taking the trace of each block separately
$$\mathrm{Tr}_B\mathrm{\Omega }^{AB}=\left(\begin{array}{ccc}\mathrm{Tr}\omega _{11}& \mathrm{}& \mathrm{Tr}\omega _{1|A|}\\ \mathrm{}& \mathrm{}& \\ \mathrm{Tr}\omega _{|A|1}& & \mathrm{Tr}\omega _{|A||A|}\end{array}\right).$$
In fact, this representation will allow us to define the following partial product
$$a^{}|\mathrm{\Omega }|a\omega _{a^{}a},$$
allowing the partial trace over $`A`$ to be expressed in the same was as with the usual trace
$$\mathrm{Tr}_A\mathrm{\Omega }=\underset{a}{}a|\mathrm{\Omega }|a=\underset{a}{}\omega _{aa}.$$
For the generic state $`\mathrm{\Psi }^{AB}`$ written in the form (2.5), let us compute
$`\mathrm{Tr}_A\mathrm{\Psi }^{AB}`$ $`=`$ $`{\displaystyle \underset{x}{}}x|^A\left({\displaystyle \underset{x^{\prime \prime }x^{}}{}}\beta _{x^{\prime \prime }}^{}\beta _x^{}|x^{\prime \prime }x^{}||\psi _{x^{\prime \prime }}\psi _x^{}|\right)|x^A`$
$`=`$ $`{\displaystyle \underset{x}{}}|\beta _x|^2\psi _x^B`$
#### 2.2.2 Purifications and extensions
Given an arbitrary density matrix $`\rho ^B`$, it is easy to construct a pure state $`|\mathrm{\Psi }^{AB}`$ such that $`\mathrm{Tr}_A\mathrm{\Psi }=\rho `$. The state $`|\mathrm{\Psi }^{AB}`$ is called a *purification* of $`\rho `$. The construction is as follows. First, choose any pure state ensemble $`\{p(x),|\psi _x^B\}`$ giving rise to $`\rho ^B`$, in the sense that
$$\underset{x}{}p(x)\psi _x=\rho .$$
Then, the state
$$|\mathrm{\Psi }^{AB}=\underset{x}{}\sqrt{p(x)}|x^A|\psi _x^B$$
is a purification of $`\rho ^B`$. This is easy to see by computing the partial trace over $`A`$, which was done for a pure state of the same form in the last subsection.
More generally we will speak of an *extension* $`\mathrm{\Omega }^{AB}`$ of a density matrix $`\rho ^A`$, which is just any density matrix (not necessarily a pure state) for which $`\mathrm{Tr}_B\mathrm{\Omega }=\rho `$. It is easy to see that any purification $`|\mathrm{\Psi }^{ABC}`$ of $`\mathrm{\Omega }^{AB}`$ is a purification of $`\rho ^A`$ as well, since
$$\mathrm{Tr}_{BC}\mathrm{\Psi }=\mathrm{Tr}_B(\mathrm{Tr}_C\mathrm{\Psi })=\mathrm{Tr}_B\mathrm{\Omega }=\rho .$$
to do: purifications! relate ensembles, purifications and measurements
#### 2.2.3 Classical-quantum (cq) systems
Consider now a collection of density matrices $`\left\{\sigma _x^A\right\}_{x𝒳},`$ indexed by the finite set $`𝒳`$. If those states occur according to the probability mass function $`p(x)`$, we may speak of an *ensemble* $`\{p(x),\sigma _x^A\}`$ of quantum states. In order to treat classical and quantum probabilities in the same framework, a joint density matrix can be constructed
$$\sigma ^{XA}=\underset{x𝒳}{}p(x)|xx|^X\sigma _x^A.$$
This is known as a *cq state*, and describes the classical and quantum aspects of the ensemble on the *extended Hilbert space* $`^{|𝒳|}^{|A|}`$ . The semiclassical nature of the ensemble is reflected in the embedding of a direct sum of Hilbert spaces $`_{x𝒳}^{|A|}`$ into $`^{|𝒳|}^{|A|}`$. This should be compared with what was done in Section 2.1.7, where a direct sum of one-dimensional vector spaces $`_{x𝒳}`$ was embedded into $`^{|𝒳|}`$. Just as the classical density matrix $`\rho `$ from (2.3) was diagonal in a basis corresponding to elements of $`𝒳`$, the cq density matrix $`\sigma `$ is *block-diagonal*, where the diagonal block corresponding to $`x`$ contains the non-normalized density matrix $`p(x)\sigma _x`$. The classical state is recoverable as $`\rho =\mathrm{Tr}_A\sigma ,`$ while the average quantum state is $`\mathrm{Tr}_X\sigma =_{x𝒳}\sigma _x`$. The classical-quantum formalism is not only of interest in its own right; information quantities evaluated on cq states play an important role in characterizing what is possible in quantum information theory.
### 2.3 Dynamics
We we have already seen an example of quantum dynamics; namely, the measurement process. In this section, we introduce the most general types of dynamical processes we will consider in this dissertation. The approach taken here will be to consider quantum channels whose inputs and/or outputs are classical-quantum systems. But first, let us review the notion of classical channels.
#### 2.3.1 Classical channels
A discrete classical channel with input symbols belonging to a finite alphabet $`𝒳`$ and output symbols from a finite alphabet $`𝒴`$ is modeled by a collection of *transition probabilities* $`p(y|x)`$. These probabilities comprise a *stochastic matrix* $`[p(y|x)]_{yx}`$, because the following two conditions are satisfied:
$$p(y|x)0\text{for each}(x,y)𝒳\times 𝒴$$
and
$$\underset{y}{}p(y|x)=1\text{for each}x𝒳$$
ensuring that to each input symbol $`x`$, there corresponds a conditional probability mass function on the output symbols $`𝒴`$.
Given an $`𝒳`$-valued random variable $`X`$ with probability mass function $`p(x)`$, the action of the channel then defines another random variable $`Y`$, jointly distributed with $`X`$ according to
$$p(x,y)=p(x)p(y|x).$$
Alternatively, we may view $`p(y|x)`$ a linear map from the simplex of probability mass functions on $`𝒳`$ to the simplex of probability mass functions on $`𝒴`$, via
$$p(x)p(y)=\underset{x}{}p(x)p(y|x).$$
In this sense, a classical channel is a model for a device which allows a sender to “prepare probability mass functions” at the output. This way of looking at classical channels leads to our first “partial” quantum generalization, described in the next section.
#### 2.3.2 Classical $``$ quantum (c $``$ q) channels
This generalization of classical channels consists of channels with a classical input and a quantum output. However, instead of preparing probability mass functions at the output, the sender prepares density matrices. A c $``$ q channel $`𝒳B`$ is specified by a collection of *conditional density matrices* $`\{\rho _x^B\}_{x𝒳}`$, labeled by the elements of a finite set $`𝒳`$. As with classical channels, such maps extend to mappings from the simplex of probability mass functions on the input alphabet $`𝒳`$ to the density matrices on the output quantum system $`B`$ via
$$p(x)\underset{x}{}\rho _x^B.$$
Such channels were implicitly considered in Section 2.2.3, where we saw that if the input is modeled by a random variable $`X`$ distributed according to $`p(x)`$, the combined input-output is a cq system with cq state
$$\rho ^{XB}=\underset{x}{}p(x)|xx|\rho _x^B.$$
The collection of c $``$ q channels with the same input set $`𝒳`$ and output quantum system $`B`$ has the structure of a compact convex set. Given two such channels with conditional density matrices $`\{\rho _x\}_{x𝒳}`$ and $`\{\sigma _x\}_{x𝒳}`$, if $`0\lambda 1`$, their corresponding convex combination has conditional density matrices $`\{\lambda \rho _x+(1\lambda )\sigma _x\}_{x𝒳}`$. The extremal points of this convex set have conditional density matrices which are extremal in the convex set of density matrices on $`B`$. In other words, the extremal points consist of channels which prepare pure states. This fact will be important when we discuss classical capacities of quantum channels in Section 4.2.
#### 2.3.3 Unitary quantum channels
The simplest quantum channel is a unitary transformation. For a closed quantum system $`A`$, this is the kind of evolution predicted by the Schrodinger equation
$$|\psi ^A|\psi ^{}^A=U|\psi ^A.$$
We will write
$$U:AA$$
to reflect the fact that $`U^{|A|\times |A|}`$ is a square matrix mapping
$$U:^{|A|}^{|A|}.$$
In this thesis, we will be exploring the consequences for processing quantum information which result from the ability to cause *any* unitary evolution to occur to a given quantum system. Ensuring that a quantum system undergoes a particular unitary evolution is generally a difficult engineering task, since it involves influencing the system in just the right way, from the outside, so as to inhibit its natural tendency to evolve in the way that it would have without any influence. To say that this will be of no concern to us here would be somewhat untrue. In fact, the central goal of this thesis is to show that, under that assumption that error-free processing of quantum information is possible, one can in fact *protect* and *correct* quantum information from this natural tendency to interfere with other quantum information and with the environment. Indeed, we will assume that it is possible to process quantum information *fault tolerantly*. If the state of $`A`$ is specified by a density matrix $`\rho ^A`$, the unitary channel acts as
$$U:\rho ^A\rho ^A=U\rho U^{}.$$
In other words, $`\rho `$ transforms according to the adjoint map associated to $`U`$. We will frequently abbreviate this map as
$$U(\rho )U\rho U^{}.$$
It will often be useful for us to speak of unitaries *between* quantum systems. For example, we may think of a quantum system $`A`$ at some time $`t`$, being turned into another quantum system $`B`$ at a later time $`t^{}`$, where $`|A|=|B|`$. If this process acts unitarily, we will write
$$U:AB$$
for the associated unitary channel. As an example, consider a physical scenario in which an electron placed at a position $`x`$ at time $`t`$ is transferred to some other position $`x^{}`$ by some later time $`t+T`$, after having been rotated by $`180^{}`$ about its $`z`$ axis. The quantum system $`A`$ thus represents the original preparation of the electron at $`x`$, while $`B`$ represents the evolved electron, $`T`$ seconds later, with its new state at position $`x^{}`$.
#### 2.3.4 Quantum channels
Quantum channels represent a physical process which transfers quantum states forward in time. The state at the output of the channel will be some noisy version of what was put in. Examples include an optical fiber over which the polarization of an input photon may become corrupted by noise, or a quantum dot which will hold a single electron for an uncertain amount of time.
Here, we will give a precise mathematical definition of quantum channels as functions from the density matrices of an *input* quantum system to the density matrices of an *output* quantum system, generalizing the notion of discrete memoryless classical channels described in Section 2.3.1, which map probability mass functions on the input alphabet to probability mass functions on the output. The mathematical properties which we require a channel to satisfy are from the standard literature on open quantum systems and quantum information theory, so much of the content here is presented without proof. Some standard references for this material include .
By a *quantum channel* $`𝒩:AB`$, we mean a mathematical object which maps density matrices on $`A`$ to density matrices on $`B`$, while satisfying the following three physically motivated properties described below.
###### Property (Linearity).
$$𝒩:^{|A|\times |A|}^{|B|\times |B|}$$
is a linear map, so that
$$𝒩\left(\underset{i}{}p_i\rho _i\right)=\underset{i}{}p_i𝒩(\rho _i).$$
###### Property (Trace preservation).
$`𝒩`$ preserves the trace of the input density operator
$$\mathrm{Tr}\rho =\mathrm{Tr}𝒩(\rho ).$$
This technical requirement will sometimes be relaxed to the requirement that $`𝒩`$ only be *trace-non-increasing*
$$\mathrm{Tr}\rho \mathrm{Tr}𝒩(\rho ).$$
With a slight loss in pedantry, we will generally refer to such maps as *trace-reducing*. In such a case, $`𝒩`$ can be interpreted as a channel which is executed with some probability less than one. To introduce the third property, let us show that there is a unique way in which $`𝒩`$ acts on the $`A`$ part of a composite quantum system $`AC`$. It is sufficient to see what happens when acting upon part of a pure state
$$|\mathrm{\Psi }^{AC}=\underset{ac}{}d_{ac}|a^A|c^C.$$
Here, we obtain
$`(𝒩1^C)(\mathrm{\Psi })`$ $`=`$ $`(𝒩1^C)\left({\displaystyle \underset{a^{}ac^{}c}{}}d_{a^{}c^{}}d_{ac}^{}|a^{}a||c^{}c|\right)`$
$`=`$ $`{\displaystyle \underset{a^{}ac^{}c}{}}d_{a^{}c^{}}d_{ac}^{}𝒩\left(|a^{}a|\right)|c^{}c|.`$
The action of $`𝒩1^C`$ is then uniquely defined on any density matrix $`\omega ^{AC}`$ by first writing any pure state decomposition
$$\omega ^{AC}=\underset{i}{}p_i\mathrm{\Psi }_i^{AC}.$$
Then by linearity,
$$(𝒩1^C)(\omega )=\underset{i}{}p_i(𝒩1^C)(\mathrm{\Psi }_i).$$
Now, we can mention the third characteristic property of a quantum channel.
###### Property (Complete positivity).
The channel must be *completely positive*, meaning that not only must $`𝒩:AB`$ take nonnegative definite matrices on $`A`$ to nonnegative definite matrices on $`B`$, but for any $`C`$ it must take nonnegative definite matrices on $`AC`$ to nonnegative definite matrices on $`BC`$.
A physically satisfying consequence of these three properties is that if a quantum channel acts on part of a convex combination of density matrices, the resulting operator will be a density matrix. Quantum channels also obey the following locality properties, which can be derived from the above three. We will later invoke these (quite frequently) without reference.
###### Property (Locality I).
Given a bipartite density matrix $`\rho ^{AB}`$ and two quantum channels
$$𝒩:AC\text{and}:BD,$$
the actions of $`𝒩`$ and $``$ commute with one another, i.e.
$$(𝒩1^B)(1^A)=(1^A)(𝒩1^B)=𝒩.$$
These equations are summarized by the leftmost commutative diagram below. The rightmost diagram is to remind the reader of the subsystems on which the corresponding states are defined.
###### Property (Locality II).
Given a bipartite density matrix $`\rho ^{AB}`$, a local operation on $`B`$ will not affect the reduced density matrix on $`A`$, i.e. given a quantum channel $`𝒩:BC`$, we have
$$\mathrm{Tr}_C(1^A𝒩)(\rho )=\mathrm{Tr}_B\rho .$$
This is summarized by the commutative diagram on the left below. On the right, we remind the reader of the subsystems involved.
This last property can be paraphrased as stating that $`\mathrm{Tr}_B\rho ^{AB}`$ is independent of any physical process which is carried out on $`B`$. Throughout this dissertation, we will often omit identity maps in expressions such as $`1^A𝒩`$, so that $`𝒩:BC`$ will be interpreted as the map $`𝒩:ABAC`$ whenever necessary. An advantage of this approach is that it allows long expressions to be simplified. This leaves no room for ambiguity, as the action of a channel on part of a larger system is always uniquely defined.
#### 2.3.5 Representing quantum channels
In this section, we review two useful representation theorems for quantum channels. The first, due to Stinespring, shows how unitary processes can give rise to quantum channels. The second, due to Kraus, shows how quantum channels can be viewed as measuring devices which “forget”, or “keep secret”, the measurement result.
Suppose that a quantum system $`A`$ is prepared and allowed to evolve unitarily with some extra system $`E`$ which is promised to be prepared in some known pure state $`|1^E`$ according to a unitary $`U:AEAE`$. Since the state of $`E`$ is guaranteed to be in the same state before the application of $`U`$, some of the elements of $`U`$ are irrelevant to the dynamics. For example, fixing bases $`\{|a^A\}`$ and $`\{|e^E\},`$ suppose that $`U`$ is given by
$$U=\underset{ae}{}|\varphi _{ae}^{AE}a|^Ae|^E$$
for some other orthogonal basis $`\{|\varphi _{ae}^{AE}\}`$ of the combined system $`AE`$. Then, an arbitrary pure state of $`A`$
$$|\varphi ^A=\underset{a}{}\alpha _a|a^A$$
will be mapped to
$`U|\varphi ^A|1^E`$ $`=`$ $`{\displaystyle \underset{ae}{}}|\varphi _{ae}^{AE}a|^Ae|^E\left({\displaystyle \underset{a^{}}{}}\alpha _a^{}|a^{}^A\right)|1^E`$
$`=`$ $`{\displaystyle \underset{aa^{}e}{}}\alpha _a^{}a|a^{}e|1|\varphi _{ae}^{AE}`$
$`=`$ $`{\displaystyle \underset{a}{}}\alpha _a|\varphi _{a1}^{AE}.`$
Thus, only the first $`|A|`$ columns of $`U`$ are relevant to this situation. Keeping only this “chunk” of the unitary $`U`$ defines an *isometry* $`𝒰:AAE`$. Mathematically, $`𝒱:AB`$ is an isometry if and only if it satisfies one (and thus both) of the following conditions:
$$𝒱^{}𝒱=1^A\text{and}𝒱𝒱^{}=\mathrm{\Pi }_A.$$
Above, $`\mathrm{\Pi }_A`$ is a projection matrix on $`B`$ satisfying $`\mathrm{Tr}\mathrm{\Pi }_A=|A|`$. In other words, an isometry is a length-preserving matrix whose range is a subspace of the target space, giving an image of the input space on the output space.
Returning to the isometry $`𝒰:AAE`$, consider what will happen if the extra system is disregarded. Given a density matrix $`\rho ^A`$, a mapping $`\mathrm{Tr}_E𝒰=𝒩:AA`$ results. This map $`𝒩`$ is a quantum channel, and the map $`𝒰`$ will be called an *isometric extension* of $`𝒩`$. We will generally use a subscript to identify the channel which is being extended, saying that $`𝒰_𝒩`$ *isometrically extends* $`𝒩`$. This way of representing a quantum channel is often referred to as the *Stinespring* representation, and we will use it almost exclusively throughout this dissertation. To be precise, we will often invoke the following proposition.
###### Proposition (Isometric extension representation).
A map $`𝒩:AB`$ is a quantum channel if and only if there exists an isometric extension $`𝒰_𝒩:ABE`$ of $`𝒩`$.
###### Remark.
In general, an isometric extension $`𝒰_𝒩`$ of $`𝒩`$ is *not unique*. This can be seen by defining $`𝒰_𝒩^{}=𝒱𝒰_𝒩`$, where $`𝒱:EE^{}`$ is any isometry into a (potentially different) environment $`E^{}`$. Since $`\mathrm{Tr}_E𝒰_𝒩=\mathrm{Tr}_E^{}𝒱𝒰_N`$, these extend the same channel $`𝒩`$.
Another way to represent a quantum channel is due to Kraus, and is called the operator sum representation (OSR). The following proposition was first proved in \[?\].
###### Proposition (Operator sum representation).
A map $`𝒩:AB`$ is a quantum channel if and only if it can be written as
$$𝒩(\rho )=\underset{i=1}{\overset{k}{}}N_i\rho N_i^{}$$
for matrices $`\{N_i^{|B|\times |A|}\}`$ which satisfy
$$\underset{i=1}{\overset{k}{}}N_i^{}N_i=1^A.$$
The matrices $`\{N_i\}`$ are called the operator sum matrices (OSR matrices) of the representation. Such a representation of $`𝒩`$ is generally not unique. It should be mentioned that this representation bears a strong resemblance to the measurement model of POVMs given in Section 2.1.6.
For a given channel, the two representations given above are intimately related, and having at hand one representation immediately gives the other as follows. If the action of $`𝒩`$ can be expressed in terms of OSR matrices $`\{N_i\}_{i=1}^k`$, an isometric extension $`𝒰_𝒩:ABE`$ into an environment of size $`|E|=k`$ can be constructed as
$$𝒰_𝒩=\underset{i=1}{\overset{k}{}}|i^EN_i.$$
This is perhaps easier expressed by writing $`𝒰_𝒩`$ as a block matrix (in the flattened representation), with blocks given by the OSR matrices as
$$𝒰_𝒩=\left(\begin{array}{c}N_1\\ N_2\\ \mathrm{}\\ N_k\end{array}\right).$$
Note that the dimensions match up; namely, $`𝒰_𝒩^{|E||B|\times |A|}`$. The reverse is also true, and the construction just involves identifying the OSR matrices with the corresponding blocks of a given isometric extension $`𝒰_𝒩`$.
###### Remark.
As the nonuniqueness of isometric extensions is due to the isometric freedom in describing the environment, the operator sum representation inherits this freedom as well.
#### 2.3.6 Complementary channels
Suppose that a channel $`𝒩:AB`$ is given. Fixing an isometric extension $`𝒰_𝒩:ABE`$ of $`𝒩`$, define the channel $`𝒩^c:AE`$ via $`𝒩^c=\mathrm{Tr}_B𝒰_𝒩`$. We will say that the channel $`𝒩^c`$ is *complementary* to $`𝒩`$. If the channel acts on a density matrix $`\rho ^A`$, the state $`𝒩^c(\rho )`$ on $`E`$ can be thought of as the disturbance induced into an initially pure environment by the action of the channel.
###### Remark.
While the choice of complementary channel is generally not unique, it is unique up to isometries on $`E`$, inheriting this freedom from the choice of isometric extension.
#### 2.3.7 Controlled quantum channels and cq $``$ q channels
Consider a collection of quantum channels $`\{_x:AB\}_{x𝒳}`$, labeled by a finite set $`𝒳`$. Introducing a controlling classical system $`X`$, available at the input and output, the collection of channels can be represented by a *controlled channel* $`:XAXB`$. This channel acts on a cq state
$$\sigma ^{XB}=\underset{x𝒳}{}p(x)|xx|^X\sigma _x^B$$
as
$$(\sigma )=\underset{x𝒳}{}p(x)|xx|^X_x(\sigma _x).$$
If the controlling system $`X`$ is not available at the output, the action of the channel is modified to
$$^{}(\sigma )=\mathrm{Tr}_X(\sigma )=\underset{x𝒳}{}p(x)_x(\sigma _x).$$
We will show next that for *any* quantum channel $`𝒩:XBC`$ which is only intended to act on cq states, less data is required to specify the action of the channel. In such a case, the channel can be represented in the same fashion as $`^{}`$, in the sense that the action of $`𝒩`$ on $`\sigma ^{XB}`$ decomposes as
$$𝒩(\sigma )=\underset{x𝒳}{}p(x)𝒩_x(\sigma _x)$$
for some channels $`\{𝒩_x:AB\}_{x𝒳}.`$ To see this, suppose that $`𝒩:XBC`$ has an operator sum decomposition
$$𝒩:\tau \underset{i=1}{\overset{d}{}}N_i\tau N_i^{},$$
where the $`|C|\times |𝒳||B|`$-dimensional matrices $`N_i`$ satisfy $`_{i=1}^dN_i^{}N_i=1^{XB}`$. Consider each $`N_i`$ to be composed of $`|𝒳|`$ blocks of size $`|C|\times |B|`$, as
$$N_i=\left(\begin{array}{cccc}N_{i1}& N_{i2}& \mathrm{}& N_{i|𝒳|}\end{array}\right).$$
The action of $`𝒩`$ on $`\sigma `$ then simplifies as
$`𝒩(\sigma )`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}N_i\sigma N_i^{}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \underset{x𝒳}{}}p(x)N_{ix}\sigma _xN_{ix}^{}`$
$`=`$ $`{\displaystyle \underset{x𝒳}{}}p(x){\displaystyle \underset{i=1}{\overset{d}{}}}N_{ix}\sigma _xN_{ix}^{}`$
$``$ $`{\displaystyle \underset{x𝒳}{}}p(x)𝒩_x(\sigma _x),`$
where in the last step we identify, for each $`x`$, the matrices $`\{N_{ix}\}_{x𝒳}`$ as the components of a trace preserving map $`𝒩_x`$.
#### 2.3.8 Quantum instruments (q $``$ cq channels)
A *quantum instrument* $`𝓝:ABX`$ is a quantum channel whose output is a cq system. Mathematically, it is specified by collection of completely positive, trace-reducing channels $`\{𝒩_x:AB\}_{x𝒳}`$, labeled by a finite set $`𝒳`$, such that the sum $`𝒩=_x𝒩_x`$, which acts on an arbitrary input state $`\rho ^A`$ as
$$𝒩(\rho )=\underset{x}{}𝒩_x(\rho ),$$
is trace preserving (and is thus a quantum channel). The action of the instrument on $`\rho ^A`$ is given by
$$𝓝(\rho )=\underset{x}{}|xx|𝒩_x(\rho ).$$
The measurement process can be modeled by a quantum instrument as follows. Given a POVM $`\{\mathrm{\Lambda }_x\}_{x𝒳},`$ consider the quantum instrument $`𝓝:AAX`$ with components acting as
$$𝒩_x(\rho )=\sqrt{\mathrm{\Lambda }_x}\rho \sqrt{\mathrm{\Lambda }_x}.$$
Then, the action of $`𝓝`$ is just
$`𝓝(\rho )`$ $`=`$ $`{\displaystyle \underset{x}{}}|xx|\sqrt{\mathrm{\Lambda }_x}\rho \sqrt{\mathrm{\Lambda }_x}`$
$`=`$ $`{\displaystyle \underset{x}{}}\mathrm{Tr}\mathrm{\Lambda }_x\rho |xx|{\displaystyle \frac{\sqrt{\mathrm{\Lambda }_x}\rho \sqrt{\mathrm{\Lambda }_x}}{\mathrm{Tr}\mathrm{\Lambda }_x\rho }}`$
$``$ $`{\displaystyle \underset{x}{}}p(x)|xx|\rho _x,`$
where the $`\{\rho _x\}`$ are the post-measurement states. In other words, $`𝒩_x(\rho )`$ is an unnormalized density matrix satisfying
$$\mathrm{Tr}𝒩_x(\rho )=p(x)$$
which is proportional to the post-measurement state. We will later utilize such a quantum instrument in order to simultaneously decode classical and quantum information which have been transmitted over a quantum multiple access channel.
It is also possible to use an instrument to model a measurement which ignores the post-measurement state. This is done with a *measuring instrument* $`𝓜:X`$, which can be defined in terms of the previous instrument as $`\mathrm{Tr}_A𝓝`$. This simpler instrument acts as
$`𝓜(\rho )`$ $`=`$ $`\mathrm{Tr}_A{\displaystyle \underset{x}{}}|xx|\sqrt{\mathrm{\Lambda }_x}\rho \sqrt{\mathrm{\Lambda }_x}`$
$`=`$ $`{\displaystyle \underset{x}{}}(\mathrm{Tr}_A\mathrm{\Lambda }_x)|xx|`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)|xx|,`$
which is exactly as one would expect a measuring device to act.
As an instrument is also a channel, it makes sense to speak of an isometric extension and complementary channel to an instrument. In the appendix (Section 11.1), we will demonstrate that any channel complementary to an instrument is another instrument with similar structure. Namely, the components of the complementary instrument are obtained as complements of the components of the original instrument.
## Chapter 3 Entropy and information quantities
In this chapter, we review the notion of quantum entropy, as well as some related information theoretical quantities which characterize the capacities to be introduced later.
### 3.1 Entropy
Let $`𝒳`$ be a finite set, and let $`X`$ be a $`𝒳`$-valued random variable, distributed according to $`p(x)`$. The *Shannon entropy* of $`X`$ is defined as
$$H(X)=\underset{x}{}p(x)\mathrm{log}p(x).$$
All logarithms in this dissertation will be to the base 2 ($`\mathrm{log}\mathrm{log}_2`$). Also, note that we will always take $`0\mathrm{log}0=0`$, as $`lim_{x0}x\mathrm{log}x=0`$ by continuity. Further note that $`H()`$ does not depend on the values taken by $`X`$. Rather, it is a functional of the probability mass function $`p(x)`$ of $`X`$. Indeed, $`𝒳`$ is merely abstract set whose elements are merely labels for events. For example, $`X`$ may be taken to represent the result of a fair coin flip, whereby $`𝒳=\{\text{heads},\text{tails}\}`$ and $`p(\text{heads})=p(\text{tails})=\frac{1}{2}`$. In this case, $`H(X)=1\text{ bit}`$. One interpretation to be gained from this example is that we obtain a bit of information by learning the result of a fair coin flip. In this sense, the coin flip example defines a “unit of information” equal to 1 bit.
$`H(X)`$ can also be interpreted as the number of bits, on average, required to represent the random variable $`X`$. Intuitively, entropy may be thought of as a measure of the amount of “information contained in” the random variable $`X`$. By definition, this is a statement concerning the asymptotic statistics of sequences of i.i.d. random variables $`X^n=(X_1,\mathrm{},X^n)`$. Such an operational definition has its roots in the source coding theorem, which dates to Shannon’s original paper , where the entropy was established as the fundamental limit on the compressibility of information. As this dissertation will focus on the closely related problem of *channel coding*, we will not pursue this interpretation further.
Suppose that a quantum system is prepared with density matrix $`\rho `$. We define the *von Neumann entropy* of $`\rho `$ as
$$H(\rho )\mathrm{Tr}\rho \mathrm{log}\rho .$$
Note that we overload the letter $`H`$ to mean both Shannon and von Neumann entropy. Writing an eigendecomposition of $`\rho `$ as
$$\rho =\underset{x}{}p(x)|xx|$$
we obtain an ensemble of orthogonal pure states $`\{p(x),|x\}`$ which also gives rise to the density matrix $`\rho `$. The von Neumann entropy of $`\rho `$ is then equal the Shannon entropy of the eigenvalues of $`\rho `$. Indeed,
$`H(\rho )`$ $`=`$ $`\mathrm{Tr}\left({\displaystyle \underset{x}{}}p(x)|xx|\right)\left({\displaystyle \underset{x}{}}\mathrm{log}p(x)|xx|\right)`$
$`=`$ $`\mathrm{Tr}\left({\displaystyle \underset{x}{}}p(x)\mathrm{log}p(x)|xx|\right)`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)\mathrm{log}p(x)=H(X)`$
where $`X`$ is a random variable with probability mass function $`p(x)`$.
If $`\rho ^A`$ is associated with system $`A`$, we will often write $`H(\rho )=H(A)_\rho ,`$ omitting the subscript when the state is apparent from the context. Given some multipartite state $`\mathrm{\Omega }^{AB}`$, the above notation gives a useful way to denote entropies of partial traces of $`\mathrm{\Omega }`$. For example, $`H(A)_\mathrm{\Omega }=H(\mathrm{Tr}_B\mathrm{\Omega })`$, while $`H(AB)_\mathrm{\Omega }=H(\mathrm{\Omega })`$. We now state the following elementary properties of entropy. These are proved in many introductory textbooks such as .
###### Property (Entropy is nonnegative).
$$H(A)0$$
This bound is saturated if and only if $`A`$ is in a pure state $`|\varphi ^A`$.
###### Property (Entropy is bounded).
$$H(A)\mathrm{log}|A|$$
This bound is saturated if and only if $`A`$ is prepared in a *maximally mixed state*
$$\pi ^A\frac{1}{|A|}1^A.$$
###### Property (Entropy is subadditive).
$$H(AB)H(A)+H(B)$$
This bound is saturated if and only if $`AB`$ is prepared in a *product state* $`\rho ^A\sigma ^B`$.
###### Property (Lieb’s inequality).
$$|H(A)H(B)|H(AB)$$
Let us compute the entropy of a generic cq state
$$\rho ^{XA}=\underset{x}{}p(x)|xx|\rho _x^A.$$
(3.1)
To do so, we first diagonalize each $`\rho _x`$ as
$$\rho _x=\underset{y}{}p_x(y)|y_xy_x|,$$
(3.2)
where for each $`x`$, the vectors $`\left\{|y_x\right\}_{y_x=1}^{|A|}`$ form (generally) different orthonormal bases for $`A`$. Then, we write
$`H(XA)`$ $`=`$ $`\mathrm{Tr}\left({\displaystyle \underset{x}{}}p(x)|xx|\rho _x\right)\mathrm{log}\left({\displaystyle \underset{x}{}}p(x)|xx|\rho _x\right)`$
$`=`$ $`\mathrm{Tr}{\displaystyle \underset{x}{}}p(x)|xx|\left(\rho _x\mathrm{log}\left(p(x)\rho _x\right)\right)`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)\mathrm{Tr}\left(\rho _x\mathrm{log}\left(p(x)\rho _x\right)\right)`$
$`=`$ $`{\displaystyle \underset{xy}{}}p(x)p_x(y)\mathrm{log}(p(x)p_x(y))`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)\left(\mathrm{log}p(x)+{\displaystyle \underset{y}{}}p_y(x)\mathrm{log}p_y(x)\right)`$
$`=`$ $`H(X)+{\displaystyle \underset{x}{}}p(x)H(\rho _x).`$
Together with subadditivity, the calculation of the joint entropy of a cq state allows a simple proof of the convexity of entropy .
###### Property (Convexity of entropy).
$$\underset{x}{}p_xH(\rho _x)H\left(\underset{x}{}p_x\rho _x\right).$$
###### Proof.
Consider the cq state $`\rho ^{XA}`$ from (3.1). Beginning with subadditivity, we have
$`H(X)+{\displaystyle \underset{x}{}}p(x)H(\rho _x)`$ $`=`$ $`H(XA)_\rho `$
$``$ $`H(X)+H(A)`$
$`=`$ $`H(X)+H\left({\displaystyle \underset{x}{}}p_x(\rho _x)\right).`$
Subtracting $`H(X)`$ from each side completes the argument. ∎
###### Property (Invariance of entropy).
For any density matrix $`\rho ^A`$ and any isometry $`𝒱:AB`$,
$$H(A)_\rho =H(B)_{𝒱(\rho )}.$$
###### Proof.
The eigenvalues of $`\rho `$ and of $`𝒱(\rho )`$ are the same. ∎
### 3.2 Conditional entropy
Let us begin by making the following formal definition for the *conditional entropy*
$$H(A|B)=H(AB)H(B).$$
By the calculation of $`H(XA)`$ for the cq state (3.1) from the previous section,
$$H(A|X)=\underset{x}{}p(x)H(\rho _x).$$
Observe that $`H(A|X)`$ is equal to the the average entropy of $`A`$, averaged over the classical part of the cq state. In classical information theory, conditional entropy is often *defined* as
$$H(Y|X)=\underset{xy}{}p(x,y)\mathrm{log}p(y|x).$$
It interesting to note that if we start with the cq state $`\rho ^{XA}`$ from (3.1), we may define a random variable $`Y`$ which is jointly distributed with $`X`$ in accordance with the conditional distribution $`p(y|x)p_x(y)`$, using the notation from (3.2). The equality $`H(A|X=x)=H(Y|X=x)`$ holds, and thus $`H(A|X)=H(Y|X)`$ holds as well.
However, for an arbitrary state on $`AB`$, this interpretation of $`H(A|B)`$ as an average entropy is not valid. In particular, suppose that $`|A|=|B|=2`$, and that $`AB`$ is in a pure state
$$|\mathrm{\Psi }^{AB}=\frac{1}{\sqrt{2}}\left(|0^A|0^B+|1^A|1^B\right).$$
Since $`\mathrm{Tr}_A\mathrm{\Psi }=\pi ^B`$, it follows that for this state,
$$H(A|B)=H(\mathrm{\Psi })H(\pi ^B)=01=1.$$
Defined in this formal way, conditional entropy can in fact be negative! As we will see in Sections 3.4 and 4.3, the negative of $`H(A|B)`$, referred to as the *coherent information*, plays a role in characterizing the quantum capacity of a quantum channel.
Let us conclude our discussion by noting the following property of conditional entropy. A proof can be found in .
###### Property.
$`H(A|B)_\rho `$ is concave as a function of $`\rho ^{AB}`$.
### 3.3 Mutual Information
Given two random variables $`X`$ and $`Y`$, jointly distributed according to $`p(x,y)`$, the *mutual information* $`I(X;Y)`$ measures the amount of correlation between the two random variables. $`I(X;Y)`$ is typically defined as an expected log likelihood ratio
$$I(X;Y)=\underset{xy}{}p(x,y)\mathrm{log}\frac{p(x,y)}{p(x)p(y)}.$$
Simple algebraic manipulations yield the following alternative formulas for $`I(X;Y)`$.
$`I(X;Y)`$ $`=`$ $`H(X)+H(Y)H(XY)`$
$`=`$ $`H(X)H(X|Y)`$
$`=`$ $`H(Y)H(Y|X).`$
Given a stochastic matrix $`p(y|x)`$ of conditional probabilities, further denotation of an input distribution $`p(x)`$ determines a joint distribution $`p(x,y)`$ for the random variables $`X`$ and $`Y`$. In Section 4.1, we will see that the capacity of a classical channel with transition matrix $`p(y|x)`$ is given by the expression
$$C=\underset{p(x)}{\mathrm{max}}I(X;Y).$$
A similar expression can be given for the capacity of a c $``$ q channel, in terms of the *quantum mutual information* evaluated on cq states.
Rather than define quantum mutual information in terms of a log-likelihood ratio, we opt here to give the following algebraic definition, valid for any composite quantum system $`AB`$.
$$I(A;B)=H(A)+H(B)H(AB).$$
Using the formal definition of conditional quantum entropy from the previous section, we could have equivalently defined $`I(A;B)`$ as
$`I(A;B)`$ $`=`$ $`H(A)H(A|B)`$
or as
$`I(A;B)`$ $`=`$ $`H(B)H(B|A).`$
Most relevant to this dissertation is the evaluation of mutual information on a cq state such as
$$\rho ^{XB}=\underset{x}{}p(x)|xx|^X\rho _x^B.$$
With respect to $`\rho ^{XB},`$ let us evaluate
$`I(X;B)`$ $`=`$ $`H(B)H(B|X)`$
$`=`$ $`H\left({\displaystyle \underset{x}{}}p(x)\rho _x\right){\displaystyle \underset{x}{}}p(x)H(\rho _x).`$
Together with the cq channel $`𝒳B`$ defined by the conditional density matrices $`\{\rho _x^B\}`$, the cq state $`\rho ^{XB}`$ represents the joint distribution on the input and output of the channel, serving the same purpose that $`p(x,y)=p(x)p(y|x)`$ did in the purely classical case. In fact, an analogous capacity formula is obtainable as well
$$C=\underset{p(x)}{\mathrm{max}}I(X;B).$$
This capacity is easily computable, as a consequence of the first of the following two convexity properties enjoyed by $`I(X;B)`$.
###### Property.
For a fixed cq channel $`𝒳B`$ defined by the conditional density matrices $`\{\rho _x^B\}`$, $`I(X;B)`$ is a concave function of $`p(x)`$.
###### Proof.
As the $`\rho ^B`$ is linear in $`p(x)`$, and $`H(B)`$ is concave in $`\rho ^B`$, $`H(B)`$ is concave in $`p(x)`$. But $`H(B|X)`$ is linear in $`p(x)`$, completing the argument. ∎
###### Property.
For a fixed input distribution $`p(x)`$, $`I(X;B)`$ is a convex function of the cq channel $`𝒳B`$.
###### Proof.
This follows because $`H(X|B)`$ is a concave function of $`\rho ^{XB}`$, which is itself linear in the conditional density matrices $`\{\rho _x^B\}`$. ∎
For an arbitrary quantum channel $`𝒩:AB`$, specification of a collection of input states $`\{\rho _x^A\}`$, or equivalently, of a cq channel with those conditional density matrices, yields a new cq channel $`𝒳B`$ with conditional density matrices $`\{𝒩(\rho _x)\}`$. This channel is mathematically equivalent to the composed actions of the cq and quantum channels. By the discussion above, optimization over input distributions $`p(x)`$ then gives the classical capacity of the newly constructed cq channel. However, the ultimate capacity of the quantum channel involves an optimization over collection of input states. Concavity of quantum mutual information in the input ensemble implies that extremal ensembles maximize capacity; such are ensembles of pure states. However, whether or not a single-letter converse can be obtained in this case remains a very important open problem in quantum information theory. As a result, the best known characterization of the capacity of a quantum channel for the transmission of classical information is
$$C(𝒩)=\underset{k\mathrm{}}{lim}\frac{1}{k}\underset{XA^k}{\mathrm{max}}I_c(X;B^k)_\sigma $$
where for each $`k`$, the maximization is over all pure state ensembles $`\{p(x),|\varphi _x^{A^k}\}`$ consisting of $`|𝒳|\mathrm{min}\{|A|,|B|\}^{2k}1`$ states. The mutual information is evaluated with respect to the corresponding cq states
$$\sigma ^{XB^k}=\underset{x}{}p(x)|xx|𝒩^k(\varphi _x^{A^k}).$$
A state such as $`\sigma `$ will be said to *arise from the channels* $`𝒩^k`$ in the above sense.
### 3.4 Coherent Information
Suppose a channel $`𝒩:A^{}B`$ is given. Fix an isometric extension $`𝒰_𝒩:A^{}BE`$, and let $`𝒩^c=\mathrm{Tr}_B𝒰_𝒩`$ the associated complementary channel. For a given input density operator $`\rho ^A^{}`$, the *coherent information* is defined as
$$I_c(\rho ,𝒩)=H(𝒩(\rho ))H(𝒩^c(\rho )).$$
Since any two complementary channels are equivalent up to an isometry on $`E`$, and since isometries preserve entropy, this quantity is independent of the particular complementary channel $`𝒩^c`$ chosen for the calculation. $`H(𝒩^c(\rho ))`$ is frequently referred to as the entropy exchange associated with sending a system with density matrix $`\rho `$ over the channel $`𝒩`$.
Coherent information can be used to characterize the capacity of a quantum channel for transmitting quantum information as
$$Q(𝒩)=\underset{k\mathrm{}}{lim}\frac{1}{k}\underset{\rho ^{A^k}}{\mathrm{max}}I_c(\rho ,𝒩^k).$$
In Section 4.3 we will give an operational definition of quantum capacity, as well a discussion of the proof of this capacity formula. It should be noted that this multi-letter characterization is the most general expression known for an arbitrary quantum channel. However, as we illustrate in Sections 9.2 and 9.3, there are classes of channels for which a single-letter expression suffices.
Let us now explore other ways of writing $`I_c(\rho ,𝒩)`$. With respect to the joint output-environment state $`𝒰_𝒩(\rho )`$ on $`BE`$, observe that
$$H(𝒩(\rho ))=H(B)\text{ and }H(𝒩^c(\rho ))=H(E).$$
Then,
$$I_c(\rho ,𝒩)=H(B)H(E).$$
It is possible to write this quantity without making explicit mention of the environment. To do this, first fix any purification $`|\mathrm{\Psi }^{AA^{}}`$ of $`\rho ^A^{}`$. Then, use this to write a global pure state
$$|\mathrm{\Omega }^{ABE}=𝒰_𝒩|\mathrm{\Psi }^{AA^{}}.$$
Since $`|\mathrm{\Omega }^{ABE}`$ is pure, it follows that $`H(E)=H(AB)`$. This allows us to rewrite
$$H(B)H(E)=H(B)H(AB)=H(A|B).$$
###### Remark.
Written this way, it is clear that coherent information can be positive or negative. However, $`Q(𝒩)0`$ for every channel $`𝒩`$, as $`I_c(|\varphi ^A^{},𝒩)=0`$ for every pure state $`|\varphi ^A^{}`$.
Observe that since any two purifications of $`\rho ^A^{}`$ are the same up to local unitaries on $`A`$, and such unitaries preserve $`H(AB)`$, this last expression is independent of the particular purification $`|\mathrm{\Psi }^{AA^{}}`$ chosen for $`\rho ^A^{}`$. Further note that to compute $`H(A|B)`$, it suffices to consider the joint state
$$\omega ^{AB}=\mathrm{Tr}_E\mathrm{\Omega }^{ABE}=𝒩(\mathrm{\Psi }^{AA^{}}).$$
It is common to write
$$I_c(AB)_\omega H(A|B)_\omega $$
acknowledging the directionality of coherent information *from* $`A`$ *to* $`B`$. While we will freely interchange the two notations for coherent information throughout this dissertation, we will generally write $`I_c(AB)`$ when characterizing capacity regions and proving the converses for the main theorems, while the notation $`I_c(\rho ,𝒩)`$ will be utilized more frequently in the coding theorems.
Let us review a few facts concerning coherent information.
###### Property.
For a maximally entangled state
$$|\mathrm{\Phi }^{AB}=\frac{1}{\sqrt{k}}\underset{i=1}{\overset{k}{}}|i^A|i^B$$
we have
$$I_c(AB)_\mathrm{\Phi }=\mathrm{log}k.$$
###### Proof.
$$H(B)_\mathrm{\Phi }H(AB)_\mathrm{\Phi }=H(\pi _k)H(\mathrm{\Phi })=\mathrm{log}k0.$$
###### Property.
For any state on $`AB`$,
$$I_c(AB)\mathrm{min}\{H(A),H(B)\}.$$
###### Proof.
We begin by observing that
$`I_c(AB)`$ $`=`$ $`H(B)H(AB)H(B)`$
$``$ $`H(B).`$
To see that $`I_c(AB)H(A)`$, we start with Lieb’s inequality
$$|H(B)H(A)|H(AB).$$
Getting rid of the absolute value and subtracting $`H(B)`$ from each side yields
$$H(A)H(A|B).$$
Multiplying the sides by $`1`$ completes the argument. ∎
###### Property.
For any channel $`𝒩:A^{}B`$ and any $`\rho ^A^{}`$,
$$I_c(\rho ,𝒩)\mathrm{log}|A^{}|.$$
###### Proof.
Fix a purification $`|\mathrm{\Phi }^{AA^{}}`$ of $`\rho ^A^{}`$. Then
$`I_c(\rho ,𝒩)`$ $`=`$ $`I_c(AB)_{𝒩(\mathrm{\Phi })}`$
$``$ $`H(A)_\mathrm{\Phi }`$
$`=`$ $`H(A^{})_\mathrm{\Phi }`$
$``$ $`\mathrm{log}|A^{}|.`$
###### Property.
For fixed $`\rho ^A^{}`$, $`I_c(\rho ,𝒩)`$ is a convex function of $`𝒩`$.
###### Proof.
Fixing a purification $`|\mathrm{\Psi }^{AA^{}}`$ of $`\rho ^A^{}`$, observe that the state $`\omega ^{AB}=𝒩(\mathrm{\Psi })`$ is linear function of $`𝒩`$. But $`H(A|B)`$ is concave in $`\omega ^{AB}`$, and thus in $`𝒩`$, so $`I_c(\rho ,𝒩)=H(A|B)`$ is convex in $`𝒩`$. ∎
###### Remark.
This property is in close agreement to the corresponding statement that $`I(X;Y)`$ is convex in $`p(y|x)`$. However, $`I(\rho ,𝒩)`$ is not generally concave or convex in $`\rho `$, for a given fixed $`𝒩`$.
### 3.5 Conditional coherent information
In the appendix (Section 11.1), we show that if $`𝓝:A^{}BX`$ is an instrument with components $`\{p(x)𝒩_x\}`$, then any isometric extension $`𝒰:A^{}BEX`$ of $`𝓝`$ can be expressed as
$$𝒰=\underset{x}{}\sqrt{p(x)}|x^X|x^X^{}𝒰_x,$$
where the $`𝒰_x:A^{}BE^{}`$ are isometric extensions of the $`𝒩_x`$, and $`EE^{}X^{}`$. We also verify there that $`\mathrm{Tr}_E𝒰=𝓝`$, while $`\mathrm{Tr}_{BX}𝒰=𝓝^𝒄`$, where
$$𝓝^𝒄=\underset{x}{}p(x)|xx|^X^{}𝒩_x^c.$$
Above, each component $`p(x)𝒩_x^c`$ is formed from a complement $`𝒩_x^c:A^{}E^{}`$ of the corresponding normalized component $`𝒩_x`$ of $`𝓝`$. The main observation here is that the environment $`E=E^{}X^{}`$ of the instrument includes the common environment $`E^{}`$ to the component channels $`𝒩_x^c`$ *as well as* a part $`X^{}`$ which purifies the classical component $`X`$ of $`𝓝.`$
For any $`\rho ^A^{}`$, the coherent information over $`𝓝`$ can thus be expressed as
$`I_c(\rho ,𝓝)`$ $`=`$ $`H\left(𝓝(\rho )\right)H\left(𝓝^c(\rho )\right)`$
$`=`$ $`H\left({\displaystyle \underset{x}{}}p(x)|xx|^X𝒩_x(\rho )\right)H\left({\displaystyle \underset{x}{}}p(x)|xx|^X^{}𝒩_x^c(\rho )\right)`$
$`=`$ $`H(X)+{\displaystyle \underset{x}{}}p(x)H\left(𝒩_x(\rho )\right)H(X){\displaystyle \underset{x}{}}p(x)H\left(𝒩_x^c(\rho )\right)`$
$`=`$ $`{\displaystyle \underset{x}{}}p(x)I_c(\rho ,𝒩_x).`$
In the third line, we mirror the calculation of the entropy of a cq state performed in Section 3.2. The coherent information over $`𝓝`$ is thus just the average of the coherent information over each $`𝒩_x`$. Another way to see this is to note that
$`I_c(\rho ,𝓝)`$ $`=`$ $`H\left(𝓝(\rho )\right)H\left(𝓝^c(\rho )\right)`$
$`=`$ $`H(BX)H(E^{}X^{})`$
$`=`$ $`H(B|X)H(E^{}|X).`$
A third derivation fixes a purification $`|\mathrm{\Psi }^{AA^{}}`$ of $`\rho ^A^{}`$ and defines the state
$`|\mathrm{\Omega }^{ABEX}`$ $``$ $`|\mathrm{\Omega }^{ABE^{}X^{}X}`$
$`=`$ $`𝒰|\mathrm{\Psi }^{AA^{}},`$
noting that
$`H(BX)H(E)`$ $`=`$ $`H(BX)H(ABX)`$
$`=`$ $`H(A|BX)`$
$`=`$ $`I_c(ABX).`$
## Chapter 4 Capacity theorems for single-user channels
In this chapter we recall various existing capacity theorems from the literature. After reviewing the proof of the capacity theorem for a classical channel, we will see that the main ingredients of that proof have counterparts for quantum channels, both for the transmission of classical and of quantum information. The common element to all of the situations is as follows. Each assumes that the sender and receiver are able to transmit an unlimited number of times over a collection of identical channels. It is useful to think of these channels as acting in parallel, as sequential transmissions can be thought of as parallel transmissions “in time”. After giving an operational definition of a set of rates at which the sender can communicate to the receiver arbitrarily well, the capacity is then *defined* to be the supremum, or least upper bound, of those achievable rates, representing the ultimate rate at which arbitrarily reliable communication can occur, provided that the channel can be used any number of times. The capacity is then described, or *characterized*, in terms of some optimization of entropic quantities over a well-defined collection of classical probabilities or quantum states. Later, when we characterize various capacity regions for quantum multiple access channels, we will invoke the single-user coding theorems for quantum channels introduced in this chapter.
### 4.1 Classical capacities of classical channels
Suppose that two parties, Alice and Bob, are connected by a large number of identical classical channels with probability transition matrix $`p(y|x)`$. This is to be interpreted as follows. At any given time, Alice can choose to send a symbol $`x𝒳`$ to Bob. Because of noise, Bob “hears” a corrupted version of the symbol $`x`$. Specifically, he receives the symbol $`y𝒴`$ with the conditional probability $`p(y|x)`$. Fixing a probability distribution $`p(x)`$ on Alice’s input symbols defines a random variable $`X`$. Together with the conditional probabilities $`p(y|x)`$, this yields a joint distribution $`p(x,y)`$ of a pair of correlated random variables $`X`$ and $`Y`$. The classical capacity of the channel $`p(y|x)`$ is the logarithm of the number of distinguishable inputs, whereby Alice uses the channel many times to send Bob a message which he can ascertain arbitrarily well. Shannon gave the following formula for the capacity:
$$C=\underset{p(x)}{\mathrm{max}}I(X;Y).$$
(4.1)
Mathematically, he proved that this expression equals a certain operationally defined capacity which we now review. Suppose Alice tries use the channel $`n`$ times to send information to Bob at a rate of $`R`$ bits per channel use. To this end, she selects a collection of *codewords*, consisting of $`2^{nR}`$ sequences of input symbols $`x^n(m)`$, one sequence for each message she would like to send, and reveals them to Bob. This can be modeled by an *encoding function*
$$f:2^{nR}𝒳^n.$$
Since the channel is noisy, Bob will receive a noisy version of Alice’s message, denoted $`Y^n(m)`$. Let the decoding function
$$g:𝒴^n2^{nR}$$
describe some scheme by which Bob attempts to decide which message Alice had intended for him to receive. Using this scheme, Alice and Bob have effectively created a new channel
$$Q(\widehat{m}|m)=\underset{y^ng^1(\widehat{m})}{}p(y^n|f(m)),$$
whereby each message $`m2^{nR}`$ Alice may choose to send induces a distribution on the possible messages Bob may decode. We might allow Alice to use a stochastic encoder $`p(x^n|m),`$ in which case the effective channel would be
$$Q(\widehat{m}|m)=\underset{y^ng^1(\widehat{m})}{}\underset{x^n}{}p(y^n|x^n)p(x^n|m).$$
If Alice sends the message $`m2^{nR}`$, the probability Bob decodes the message incorrectly can be expressed in a number of ways:
$`P_e(m)`$ $``$ $`\mathrm{Pr}\{\widehat{M}m|M=m\}`$
$`=`$ $`\mathrm{Pr}\{g(Y^n(m))m\}`$
$`=`$ $`1Q(m|m)`$
$`=`$ $`{\displaystyle \underset{\stackrel{\widehat{m}2^{nR}}{\widehat{m}m}}{}}Q(\widehat{m}|m).`$
Associated to the coded channel $`Q(\widehat{m}|m)`$ is its *maximal probability of error*
$$P_{\text{max}}=\underset{m2^{nR}}{\mathrm{max}}P_e(m)$$
and its *average probability of error*
$$P_{\text{ave}}=2^{nR}\underset{m2^{nR}}{}P_e(m).$$
One may phrase the goal of successful communication as that of *simulating* a fictitious identity channel id$`:2^{nR}2^{nR}`$ from Alice to Bob, where id$`(\widehat{m}|m)=\delta _{\widehat{m},m}`$. Perfect simulation would amount to using a zero-error code. Approximate simulation can be gauged in a number of ways. For example, one could require that either $`P_{\text{ave}}`$ or $`P_{\text{max}}`$ is small. Clearly, the former will imply the latter.
Suppose that Alice chooses her message $`M`$ randomly according to the distribution $`P(m)`$. If she sends her message through the identity channel to Bob, the two will hold a perfectly correlated pair of random variables $`(M,M),`$ distributed as
$$\text{dist}(M,M)_P(m,\widehat{m})=P(m)\delta _{m,\widehat{m}}.$$
However, Alice will actually be sending through the coded channel $`Q(\widehat{m}|m)`$, generating a pair of noisy correlated random variables $`(M,\widehat{M})`$ distributed as
$$\text{dist}(M,\widehat{M})_P(m,\widehat{m})=P(m)Q(\widehat{m}|m).$$
One way to judge the success of the simulation is to consider the the $`\mathrm{}_1`$ norm $`\mathrm{\Delta }(P)`$ between the two distributions dist$`(M,\widehat{M})_P`$ and dist$`(M,M)_P`$. This is calculated as
$`\mathrm{\Delta }(P)`$ $`=`$ $`\left|\text{dist}(M,M)_P\text{dist}(M,\widehat{M})_P\right|_1`$
$`=`$ $`{\displaystyle \underset{m,\widehat{m}=1}{\overset{\mu }{}}}|P(m)\delta _{m,\widehat{m}}P(m)Q(\widehat{m}|m)|`$
$`=`$ $`{\displaystyle \underset{m,\widehat{m}=1}{\overset{\mu }{}}}P(m)|\delta _{m,\widehat{m}}Q(\widehat{m}|m)|`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{\mu }{}}}P(m)\left((1Q(\widehat{m}|m))+{\displaystyle \underset{\stackrel{\widehat{m}=1}{\widehat{m}m}}{\overset{\mu }{}}}Q(\widehat{m}|m)\right)`$
$`=`$ $`2{\displaystyle \underset{m=1}{\overset{\mu }{}}}P(m)P_e(m)`$
$`=`$ $`2𝔼_PP_e(M).`$
In other words, the $`\mathrm{}_1`$ distance between the ideal and the actual joint distributions is precisely equal to twice the expected error probability. Observe that
$$\mathrm{\Delta }(\text{unif}(2^{nR}))=2P_{\text{ave}}\text{ and }\mathrm{\Delta }(\delta _m)=2P_e(m).$$
Further note that requiring that the maximal error probability be less than $`ϵ`$ is equivalent to demanding that $`\mathrm{\Delta }(\delta _m)2ϵ`$ for each $`m`$, where $`\delta _m`$ is a point distribution at $`\{M=m\}`$. It is worth noting that the latter requirement is also equivalent to requiring that $`\mathrm{\Delta }(P)2ϵ`$ for all distributions $`P(m)`$.
So, communication can be viewed in the light of generating near perfect common randomness over noisy quantum channels. We have phrased things in this way as it makes the road to quantum communication a bit easier. Rather than asking the sender and receiver to end up with classical correlations, we will see later in Section 4.3 that they attempt to build *quantum correlations*.
Any code $`(f,g)`$ which encodes $`2^{nR}`$ messages using $`n`$ instances of a channel $`p(y|x)`$ such that $`P_e(m)ϵ`$ for all $`m2^{nR}`$ will be called an $`(R,n,ϵ)`$ maximal error code for the channel $`p(y|x)`$. A rate $`R`$ is said to be *achievable* if there exists a sequence of $`(R,n,ϵ_n)`$ maximal error codes with $`ϵ_n0`$. The (operational) capacity of the channel $`p(y|x)`$ is then defined to be the supremum of the set of achievable rates. Shannon’s capacity theorem states that this operationally defined capacity is equal to the number $`C`$, defined in (4.1).
The channel capacity theorem is proved in two main parts. First, it is proven that for any rate $`R<C`$, $`R`$ is achievable. This is provided by a *coding theorem*, which is generally structured as follows. Given $`ϵ>0`$ and some rate $`R<C`$, it is shown that there is a long enough blocklength $`n`$ so that there exists an $`(R,n,ϵ)`$ code. As $`ϵ`$ was arbitrary, this immediately implies the existence of a sequence of such codes which achieves the rate $`R`$, corresponding to any sequence of error probabilities which go to zero. The second component is called the *converse*. In this part, it is shown that every achievable rate $`R`$ satisfies $`R<C`$. These components are summarized in Figure 4.1.
One route to proving the coding theorem involves first showing that codes with a weaker error constraint exist. Rather than requiring that every message have a low error probability, it is sufficient to show that the error probability, averaged over all codewords $`m2^{nR}`$ is small. A code satisfying this weaker constraint will be called an *average error code.* A way to prove such a coding theorem is through the technique of random coding. For an arbitrary distribution $`p(x)`$, define the product distribution $`p(x^n)=_{i=1}^np(x_i)`$. A rate $`R`$ random encoder is then defined by randomly selecting $`2^{nR}`$ codewords
$$𝒞=\{X^n(1),\mathrm{},X^n(2^{nR})\}$$
i.i.d. according to $`p(x^n)`$. The following coding proposition, or some variant thereof, is proved in many textbooks on information theory, such as in .
###### Proposition 0 (Classical channel coding theorem).
Given is a channel $`p(y|x)`$, an input distribution $`p(x)`$, and a number $`0R<I(X;Y)`$, where $`I(X;Y)`$ is computed with respect to $`p(x,y)=p(x)p(y|x)`$. For every $`ϵ>0`$, there is $`n`$ sufficiently large so that if $`2^{nR}`$ codewords $`𝒞=\{X^n(1),\mathrm{},X^n(2^{nR})\}`$ are chosen i.i.d. according to the product distribution $`p(x^n)=_ip(x_i)`$, there exists a decoding function $`g:𝒴^n2^{nR}`$ which depends on the random choice of codebook $`𝒞`$ and correctly identifies the input message with expected average probability of error less than $`ϵ`$, in the sense that
$$𝔼_𝒞2^{nR}\underset{m2^{nR}}{}\mathrm{Pr}\{g(Y^n(m))=m\}1ϵ.$$
Observe that, because of the symmetry in the code construction, the expectation of each term in the above summation is the same. It is thus possible to reexpress that error condition as
$$𝔼_𝒞\mathrm{Pr}\{g(Y^n(1))=1\}1ϵ,$$
showing that at the level of random codes, one may assume that the message $`m=1`$ has been sent without losing any generality.
It is a simple task to “derandomize” any code which is guaranteed to exist by Proposition 0. Suppose that Alice chooses a message uniformly distributed on the set $`\{1,\mathrm{},2^{nR}\}`$, represented by the random variable $`M`$, to send to Bob. Then
$`𝔼_𝒞\mathrm{Pr}\{g(Y^n(M)=M)\}`$ $`=`$ $`2^{nR}{\displaystyle \underset{m2^{nR}}{}}𝔼_𝒞\mathrm{Pr}\{g(Y^n(m)=m)\}`$
$``$ $`1ϵ.`$
It is then immediate that there must exist a particular deterministic code yielding an average probability of success at least as large as $`1ϵ`$.
So far, this is enough to conclude that every input distribution $`p(x)`$ yields a lower bound to the average error capacity of $`p(y|x)`$. This is because each $`p(x)`$ corresponds to a set of achievable rates $`\{R:0R<I(X;Y)\},`$ and the largest such set is given by optimizing over all $`p(x)`$.
Recall that we have defined the operational capacity $`C`$ in terms of the maximal probability of error constraint. However, we have only outlined how to show that codes with low average error exist. By Markov’s inequality from probability theory, if the average error probability is less than $`ϵ`$, then at least half of the codewords have an error probability less than $`\sqrt{ϵ}`$. By only using these codewords, a rate $`R\frac{1}{n}`$ code with maximal error probability $`\sqrt{ϵ}`$ is obtained, and thus every rate less than $`R`$ is achievable with maximal error, showing that the maximal and average error capacities are the same.
While the coding proposition implies the existence of sequences of codes achieving any rate less than capacity, it remains to prove that no such sequences exist for rates above capacity. Rather than reproduce the entire converse theorem, we outline the basic structure of the theorem. First, one assumes that $`R`$ is an achievable rate. This means that there should exist a sequence of $`(2^{nR},n,ϵ_n)`$ codes with $`ϵ_n0`$. For any $`n`$, let $`p(x^n,y^n)=p(x^n)_ip(y_i|x_i)`$ be the joint distribution on $`X^n`$ and $`Y^n`$ induced by selecting codewords uniformly at random from the corresponding code in the sequence. An initial step in the proof shows that
$$R<\frac{1}{n}I(X^n;Y^n)+ϵ_n^{}$$
where $`ϵ_n^{}0`$ as $`ϵ_n0`$, and $`I(X^n;Y^n)`$ is evaluated with respect to the induced distribution $`p(x^n,y^n)`$. For *any* joint distribution on $`X^n`$ and $`Y^n`$, the following can be easily proved:
$$\frac{1}{n}I(X^n;Y^n)\frac{1}{n}\underset{i=1}{\overset{n}{}}I(X_i;Y_i)\underset{i}{\mathrm{max}}I(X_i;Y_i).$$
If $`i^{}`$ achieves the maximum on the right hand side, the marginal distribution $`p(x_i^{})`$ provides a “witness” to the fact that the rate $`R`$ is in fact achievable ($`R`$ is thus less than the maximum mutual information over all input distributions). This proves that the capacity formula is *additive*, and thus that every achievable rate is upper bounded by the solution of a “single-letter” optimization problem. For this reason, this second conceptual step in the converse is known as single-letterization. Without it, one would only be able to write the capacity as
$$C=\underset{k\mathrm{}}{lim}\frac{1}{k}\underset{p(x^k)}{\mathrm{max}}I(X^k;Y^k)$$
a result which follows by applying Proposition 0 to extensions of the channel
$$p(y^k|x^k)=\underset{i=1}{\overset{k}{}}p(y_i|x_i).$$
Such an expression has become known as a “regularized” expression for the capacity. Actually, this is a persistent problem in quantum information theory. The best known expressions characterizing the capacities of an arbitrary quantum channel to transmit classical or quantum information are regularized maximizations of information quantities over appropriate sets of input states.
### 4.2 Classical capacities of quantum channels
Suppose that Alice and Bob are connected via some large number $`n`$ of instances of a quantum channel $`𝒩`$, and that Alice wishes to transmit classical messages to Bob. The overall maximal rate at which is this is possible is the *classical capacity* $`C(𝒩)`$ of the channel $`𝒩`$, which is the logarithm of the number of physical input preparations Alice can make, per channel use, so that Bob can distinguish them arbitrarily well by measuring the induced states at the outputs of the channels. The best known expression for the classical capacity of a quantum channel, due to Holevo Schumacher and Westmoreland , is the following regularized formula, known as the HSW Theorem:
$$C(𝒩)=\underset{k\mathrm{}}{lim}\frac{1}{k}\underset{XA^k}{\mathrm{max}}I(X;B^k)_\omega .$$
Here, the maximization is over all pure state input ensembles $`\{p(x),|\varphi _x^{A^k}\}`$ of states for Alice to prepare at the inputs to $`k`$ parallel instances of the channel $`𝒩`$. For a given ensemble, the mutual information is computed relative to the corresponding cq state
$$\omega ^{XB^k}=\underset{m}{}|xx|^X𝒩^k(\varphi _x).$$
Operationally, the classical capacity of $`𝒩`$ is defined in analogy to that of a classical channel. A $`(2^{nR},n)`$ code consists of $`2^{nR}`$ message states $`\{|\varphi _1^{A^k},\mathrm{},|\varphi _{2^{nR}}^{A^k}\}`$ for Alice and a corresponding measurement for Bob, mathematically modeled as POVM with $`2^{nR}`$ outcomes $`\{\mathrm{\Lambda }_m\}_{m2^{nR}}`$. We call this code an $`(2^{nR},n,ϵ)`$ code if the following constraint on success probability, averaged over all messages, is satisfied:
$$2^{nR}\underset{m2^{nR}}{}\mathrm{Tr}\mathrm{\Lambda }_m𝒩^n(\varphi _m)1ϵ.$$
A rate $`R`$ is achievable if there exists a sequence of $`(2^{nR},n,ϵ_n)`$ codes with $`ϵ_n0`$, and the capacity $`C(𝒩)`$ is the supremum of all achievable rates.
As with the capacity of a classical channel, the proof that $`C(𝒩)`$ can be expressed in such a regularized form has two parts, a coding theorem and a converse. The following coding theorem is attributed to Holevo , Schumacher and Westmoreland .
###### Proposition 1 (HSW Theorem).
Given is a cq state $`\sigma ^{XB}=_xp(x)|xx|^X\rho _x^B`$ and a number $`0R<I(X;B)_\sigma .`$ For every $`ϵ>0`$, there is $`n`$ sufficiently large so that if $`2^{nR}`$ codewords $`𝒞=\{X^n(m)\}`$ are chosen i.i.d. according to the product distribution $`p(x^n)=_{i=1}^np(x_i)`$, corresponding to input preparations
$$\rho _{x^n}=\rho _{x_1}\mathrm{}\rho _{x_n},$$
there exists a decoding POVM $`\{\mathrm{\Lambda }_m\}`$ on $`B^n`$ which depends on the random choice of codebook $`𝒞`$ and correctly identifies the index $`m`$ with average probability of error less than $`ϵ,`$ in the sense that
$`𝔼_𝒞2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}\mathrm{Tr}\rho _{X^n(m)}\mathrm{\Lambda }_m1ϵ.`$ (4.2)
Due to the symmetry of the distribution of $`𝒞`$ under codeword permutations, it is clear that the expectations of each term in the above sum are equal. In other words,
$`𝔼_𝒞2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}\mathrm{Tr}\rho _{X^n(m)}\mathrm{\Lambda }_m=𝔼_𝒞\mathrm{Tr}\rho _{X^n(1)}\mathrm{\Lambda }_1,`$ (4.3)
The arguments for derandomization and for obtaining a good maximal error code are identical to those used for classical channels in the previous section.
A proof of the converse begins, as before, by assuming that $`R`$ is an achievable rate. Taking a cq state $`\omega ^{XB^n}`$ induced by an $`(R,n,ϵ_n)`$ code in the achieving sequence, Fano’s inequality (Lemma 5) and the Holevo Bound (Lemma 7) are used <sup>1</sup><sup>1</sup>1These details are given more explicitly in the converse proofs of the main theorems (Section 7.2). to show that
$$R<\frac{1}{n}I(X;B^n)_\omega ,$$
where again $`ϵ_n^{}0`$ as $`ϵ_n0`$. However, it is an important open problem as to whether a single-letterization step can be proved. No counterexample to additivity is known, and it is widely believed that none exists.
### 4.3 Quantum capacities of quantum channels
The quantum capacity $`Q(𝒩)`$ of a quantum channel $`𝒩:A^{}B`$ is the answer to a number of physical questions regarding the possibilities of performing various operational information processing tasks over many parallel instances of the channel $`𝒩`$. $`Q(𝒩)`$ is the logarithm of various quantities:
* the amount of entanglement that can be created (entanglement generation)
* the amount of entanglement that can be sent (entanglement transmission)
* the size of a Hilbert space all of whose states can be reliably transmitted
(subspace transmission)
* the size of a Hilbert space all of whose entangled states can be reliably transmitted (strong subspace transmission).
All of these quantities have units of *qubits per channel use*, and as the rates at which these tasks are possible all coincide, it is justifiable to say that they all represent “sending quantum information,” and hence to speak of a single quantum capacity $`Q(𝒩).`$ The best known characterization of the quantum capacity is a regularized maximization of the coherent information
$$Q(𝒩)=\underset{k\mathrm{}}{lim}\frac{1}{k}\underset{XA^{}}{\mathrm{max}}I_c(AB^k)_\omega ,$$
where for each $`k`$, the maximization is over all states of the form
$$\omega ^{AB^k}=𝒩^k(\mathrm{\Psi }^{AA^k}).$$
Such a state $`\omega `$ will be said to *arise from* $`𝒩^k`$ or rather, to arise from the action of $`𝒩^k`$ on the bipartite pure state $`|\mathrm{\Phi }^{AA^k}`$. Here, the regularization is known to be necessary for a general quantum channel, as opposed to the case with the classical capacity $`C(𝒩)`$, where the existence of a single-letterization step in the converse is an open problem. The existence of a counterexample to additivity is known .
Of the different operational definitions of $`Q(𝒩)`$, the simplest to describe is entanglement generation, since it can defined without explicit mention of encodings. Suppose that a large number $`n`$ of channels $`𝒩:A^{}B`$ are available from Alice to Bob. Alice and Bob will use the channels to build a large maximally entangled state between degrees of freedom of some physical systems located in their respective laboratories. To this end, Alice prepares some bipartite pure state $`|\mathrm{{\rm Y}}^{AA^n}`$, entangled between some system $`A`$ of dimension $`|A|=2^{nQ}`$ in her laboratory, and the inputs $`A^n`$ of the channels. After the actions of the channels, Alice’s system $`A`$ is correlated with the outputs $`B^n`$ of the channels quantum mechanically. Bob then performs some post-processing procedure, modeled by a quantum operation $`𝒟:B^n\widehat{A}`$, to transfer the quantum correlations from the outputs $`B^n`$ of the channels to an “output” physical system $`\widehat{A}`$, also of dimension $`|\widehat{A}|=2^{nQ}`$ in his laboratory. Their goal is to produce a state which is close to some target maximally entangled state $`|\mathrm{\Phi }^{A\widehat{A}}`$. More specifically, we say that they generate entanglement at rate $`Q`$ if they produce a maximally entangled state of the form
$$|\mathrm{\Phi }^{A\widehat{A}}=\frac{1}{\sqrt{2^{nQ}}}\underset{a2^{nQ}}{}|a^A|a^{\widehat{A}}.$$
We will call such a state a *rate $`Q`$ maximally entangled state*. The blocklength $`n`$ will always be apparent from the context.
$`(|\mathrm{{\rm Y}}^{AA^n},𝒟)`$ will be called a $`(Q,n,ϵ)`$ *entanglement generation code* for the channel $`𝒩`$ if, for the rate $`Q`$ maximally entangled state $`|\mathrm{\Phi }^{A\widehat{A}}`$, we have
$$F(|\mathrm{\Phi }^{A\widehat{A}},𝒟𝒩^n(\mathrm{{\rm Y}}^{AA^n}))1ϵ.$$
A rate $`Q`$ is an *achievable rate for entanglement generation* over the channel $`𝒩`$ if there exists a sequence of $`(Q,n,ϵ_n)`$ entanglement generation codes with $`ϵ_n0`$. The *entanglement generating capacity* $`Q^{\text{eg}}(𝒩)`$ of $`𝒩`$ is then defined operationally as the supremum of all such achievable rates.
We will now introduce a number of coding propositions from , each a more refined version of the previous one. While the first is sufficient to prove achievability for single-user channels, the others have additional properties which we will need later when we characterize various capacity regions of quantum multiple access channels.
###### Proposition (Entanglement generation coding theorem).
Given is a channel $`𝒩:A^{}B`$, a density matrix $`\rho ^A^{}`$, and a number $`0Q<I_c(\rho ,𝒩).`$ For every $`ϵ>0`$, there is $`n`$ sufficiently large so that there is a $`(Q,n,ϵ)`$ entanglement generation code $`(|\mathrm{{\rm Y}}^{AA^n},𝒟)`$ for $`𝒩`$.
Recall the discussion in Section 3.4 regarding the two different ways of expressing coherent information. Given an input density operator $`\rho ^A^{}`$, if $`|\mathrm{\Psi }^{AA^{}}`$ is any purification of $`\rho `$, then the identity
$$I_c(\rho ,𝒩)=I_c(AB)_{𝒩(\mathrm{\Psi })}$$
holds. This proposition then guarantees that for every state $`\omega ^{AB}=𝒩(\mathrm{\Psi })`$ arising from the action of $`𝒩`$ on a state $`|\mathrm{\Psi }^{AA^{}}`$, every rate $`0Q<I_c(AB)_\omega `$ is an achievable rate. This works by applying the coding theorem to the input state $`\rho ^A^{}=\mathrm{Tr}_A\mathrm{\Psi }`$.
As with the classical capacity, it is also true that for each integer $`k>0`$, if $`\omega ^{}`$ arises from $`𝒩^k`$, then every rate $`0Q<\frac{1}{k}I_c(AB^k)_\omega ^{}`$ is achievable as well. We then conclude that
$$Q(𝒩)\underset{k\mathrm{}}{lim}\frac{1}{k}\mathrm{max}I_c(AB^k).$$
The usual Shannon-theoretic prescription for converse theorems applies here as well, although as mentioned above, it known that a single-letterization step cannot be proved for arbitrary $`𝒩`$. Suppose that $`Q`$ is achievable, and fix a $`(Q,n,ϵ_n)`$ entanglement generation code $`(|\mathrm{{\rm Y}}^{AA^n},𝒟)`$ in the achieving sequence of codes. The encoding $`|\mathrm{{\rm Y}}`$ gives rise to the state $`\omega ^{AB^n}=𝒩^n(\mathrm{{\rm Y}})`$. It is a simple consequence of the quantum data processing inequality (Lemma 6) and continuity of coherent information in the input density operator (Lemma 3) that <sup>2</sup><sup>2</sup>2these details are given more explicitly in the converse proofs of the main theorems (Section 7.2).
$$Q\frac{1}{n}I_c(AB^n)_\omega +ϵ^{}_n$$
where $`ϵ_n^{}0`$. By standard arguments we then conclude that
$$Q(𝒩)\underset{k\mathrm{}}{lim}\frac{1}{k}\mathrm{max}I_c(AB^k).$$
The state $`\mathrm{Tr}_A\mathrm{{\rm Y}}`$ which is induced by Alice’s encoding at the inputs $`A^n`$ of $`𝒩^n`$ is called the *code density operator* of the entanglement generation code. With randomization, it is possible to make this operator arbitrarily close to the product state $`\rho ^n`$, where $`\rho ^A^{}`$ is the input density matrix used when invoking the proposition. If Alice and Bob have access to a shared source of randomness, they may utilize an ensemble of codes to this end. This is very useful for our multiple access coding theorems, as it guarantees that if one sender codes randomly, the induced channel seen by the other sender is close to a product channel, allowing coding theorems for product channels to be invoked.
A $`(Q,n,ϵ)`$ *random entanglement generation code* consists of a collection of deterministic $`(Q,n,ϵ)`$ entanglement transmission codes $`(|\mathrm{{\rm Y}}^\beta ^{AA^n},𝒟^\beta )`$ and a probability distribution $`P_\beta `$, corresponding to a source of shared common randomness available to both sender and receiver. We will often omit the subscript, once the randomness of the code has been clarified, and it will be understood that $`|\mathrm{{\rm Y}}`$ and $`𝒟`$ constitute a pair of classically correlated random objects. Associated to a random code is its expected, or average code density operator
$$\varrho ^{A^n}=𝔼_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}=\underset{\beta }{}P_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}^\beta $$
which is the expectation, over the shared randomness, of the state at the channel inputs $`A^n`$. The following extension of the previous coding proposition pertains to these random codes and is also proved in .
###### Proposition (Random entanglement generation coding theorem).
Given is a channel $`𝒩:A^{}B`$, a density matrix $`\rho ^A^{}`$, and a number $`0Q<I_c(\rho ,𝒩).`$ For every $`ϵ>0`$, there is $`n`$ sufficiently large so that there is a $`(Q,n,ϵ)`$ random entanglement generation code $`(P_\beta ,|\mathrm{{\rm Y}}^\beta ^{AA^n},𝒟^\beta )`$ for $`𝒩`$ with average code density operator
$$\varrho ^{A^n}=𝔼_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}=\underset{\beta }{}P_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}^\beta $$
satisfying
$$|\varrho \rho ^n|_1ϵ.$$
Finally, there are certain features of the decoder structure of random entanglement generation codes that are necessary for proofs which utilize quantum side information at the decoder. This final form of the coding proposition is the most powerful, utilizing features which are implicit from the proof of the coding theorem of . This will be the proposition which is invoked later in the dissertation.
###### Proposition 2.
Given is a channel $`𝒩:A^{}B`$, a density matrix $`\rho ^A^{}`$, and a number $`0Q<I_c(\rho ,𝒩).`$ For every $`ϵ>0`$, there is $`n`$ sufficiently large so that there is a random $`(Q,n,ϵ)`$ entanglement generation code $`(P_\beta ,|\mathrm{{\rm Y}}^\beta ^{AA^n},𝒟^\beta )`$ for $`𝒩`$ with average code density operator
$$\varrho ^{A^n}=𝔼_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}=\underset{\beta }{}P_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}^\beta $$
satisfying
$$|\varrho \rho ^n|_1ϵ.$$
Furthermore, given any particular isometric extension $`𝒰_𝒩:A^{}BE`$ of $`𝒩`$, it is possible to choose isometric extensions $`𝒰_𝒟^\beta :B^n\widehat{A}F`$ of the deterministic decoders so that
$$F(|\mathrm{\Phi }^{A\widehat{A}}|\lambda ^{E^nF},𝒰_𝒟^\beta 𝒰_𝒩^n|\mathrm{{\rm Y}}^\beta ^{AA^n})1ϵ$$
for every $`\mathrm{}`$ and the same fixed pure state $`|\lambda ^{E^nF}`$.
## Chapter 5 Main results
### 5.1 Quantum multiple access channels
For this dissertation, a quantum multiple access channel will have two senders and a single receiver. While many-sender generalizations of the theorems which appear here are readily obtainable, we focus on the case with two senders for simplicity. Such a channel $`𝒩:AB^{}C`$ will generally be one in which Alice and Bob simultaneously transmit to Charlie. We will assume throughout that no other resources are available to the three parties. Namely, none of the parties share any prior classical or quantum correlations between themselves, nor do they have access to any other auxiliary channels. If Alice inputs a physical system with density matrix $`\rho _1^A^{}`$, while Bob’s input has density matrix $`\rho _2^B^{}`$, Charlie will receive the state $`𝒩(\rho _1\rho _2)`$.
In the next section, we give an operational definition of the four-dimensional region $`𝒮(𝒩)`$, which consists of the rates at which each sender can simultaneously send classical and quantum information to Charlie. Sections 5.3 and 5.4 state the main results of this dissertation. These results characterize the two-dimensional shadows of $`𝒮(𝒩)`$ corresponding to the situation where Alice sends classically while Bob sends quantum information (Theorem 1), and that where each sends quantum information (Theorem 2).
These theorems will be proved by first showing in Chapter 7 that the characterizations given in Sections 5.3 and 5.4 describe other sets of operationally defined rates, corresponding to weaker constraints on good codes than those to be introduced in this chapter. In Chapter 8, it will ultimately be shown that the other sets of operationally defined rates equal those introduced in this chapter.
### 5.2 $`𝒮(𝒩)`$ \- the general problem
Assume that Alice and Bob are connected to Charlie by $`n`$ instances of a multiple access channel $`𝒩:A^{}B^{}C`$, where Alice and Bob respectively have control over the $`A^n`$ and $`B^n`$ inputs. We will describe a scenario in which Alice wishes to transmit classical information at a rate of $`R_a`$ bits per channel use, while simultaneously transmitting quantum information at a rate of $`Q_a`$ qubits per channel use. At the same time, Bob will be transmitting classical and quantum information at rates of $`R_b`$ and $`Q_b`$ respectively. Alice attempts to convey any one of $`2^{nR_a}`$ messages to Charlie, while Bob tries to send him one of $`2^{nR_b}`$ such messages. We will also assume that the senders are presented with systems $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$, where $`|\stackrel{~}{A}|=2^{nQ_a}`$ and $`|\stackrel{~}{B}|=2^{nQ_b}`$. Each will be required to complete the following two-fold task. Firstly, they must individually transfer the quantum information embodied in $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ to their respective inputs $`A^n`$ and $`B^n`$ of the channels, in such a way that it is recoverable by Charlie at the receiver. Second, they must simultaneously make Charlie aware of their independent messages $`M_a`$ and $`M_b`$. Alice and Bob will encode with maps from the cq systems holding their classical and quantum messages to their respective inputs of $`𝒩^n`$, which we denote
$$_1:M_a\stackrel{~}{A}A^n\text{ and }_2:M_b\stackrel{~}{B}B^n.$$
Charlie decodes with a quantum instrument
$$𝓓:C^n\widehat{M}_a\widehat{M}_b\widehat{A}\widehat{B}.$$
The output systems are assumed to be of the same sizes and dimensions as their respective input systems. For the quantum systems, we assume that there are pre-agreed upon unitary correspondences id$`{}_{a}{}^{}:\stackrel{~}{A}\widehat{A}`$ and id$`{}_{b}{}^{}:\stackrel{~}{B}\widehat{B}`$ between the degrees of freedom in the quantum systems presented to Alice and Bob which embody the quantum information they are presented with and the target systems in Charlie’s laboratory to which that information should be transferred. The goal for quantum communication will be to, in the strongest sense, simulate the actions of these corresponding identity channels. We similarly demand low error probability for each pair of classical messages. Formally, $`(_1,_2,𝓓)`$ will be said to comprise an $`(R_a,R_b,Q_a,Q_b,n,ϵ)`$ *strong subspace transmission code* for the channel $`𝒩`$ if for all $`m_a2^{nR_a}`$, $`m_b2^{nR_b}`$, $`|\mathrm{\Psi }_1^{A\stackrel{~}{A}}`$, $`|\mathrm{\Psi }_2^{B\stackrel{~}{B}}`$, where $`A`$ and $`B`$ are purifying systems of arbitrary dimensions,
$$F(|m_a^{\widehat{M}_a}|m_b^{\widehat{M}_b}|\mathrm{\Psi }_1^{A\widehat{A}}|\mathrm{\Psi }_2^{B\widehat{B}},\mathrm{\Omega }_{m_am_b})1ϵ$$
where
$`\mathrm{\Omega }_{m_am_b}^{\widehat{M}_a\widehat{M}_bA\widehat{A}B\widehat{B}}=𝓓𝒩^n\left(_1\left(|m_am_a|^{M_a}\mathrm{\Psi }_1^{A\stackrel{~}{A}}\right)_2\left(|m_bm_b|^{M_b}\mathrm{\Psi }_2^{B\stackrel{~}{B}}\right)\right).`$
We will say that a rate vector $`(R_a,R_b,Q_a,Q_b)`$ is *achievable* if there exists a sequence of $`(R_a,R_b,Q_a,Q_b,n,ϵ_n)`$ strong subspace transmission codes with $`ϵ_n0`$. The simultaneous capacity region $`𝒮(𝒩)`$ is then defined as the closure of the collection of achievable rates. Setting various rate pairs equal to zero uncovers six two-dimensional rate regions. The next section contains our first theorem, which gives a multi-letter characterization of the two shadows relevant to the situation where one user only sends classical information, while the other only sends quantum information. The following section contains a theorem which describes the rates at which each sender can send quantum information via a multi-letter formula.
### 5.3 $`𝒞𝒬(𝒩)`$ \- classical-quantum capacity region
Suppose that Alice only wishes to send classical information at a rate of $`R`$ bits per channel use, while Bob will only send quantum mechanically at $`Q`$ qubits per use of the channel. The rate pairs $`(R,Q)`$ at which this is possible comprise a classical-quantum (cq) region $`𝒞𝒬(𝒩)`$ consisting of rate vectors in $`𝒮(𝒩)`$ of the form $`(R,0,0,Q)`$. Our first theorem gives a characterization of $`𝒞𝒬(𝒩)`$ as a regularized union of rectangles.
###### Theorem 1.
$`𝒞𝒬(𝒩)`$ = the closure of the union of pairs of nonnegative rates $`(R,Q)`$ satisfying
$`R`$ $``$ $`{\displaystyle \frac{1}{k}}I(X;C^k)_\omega `$
$`Q`$ $``$ $`{\displaystyle \frac{1}{k}}I_c(BC^kX)_\omega `$
for some $`k`$, some pure state ensemble $`\{p(x),|\varphi _x^{A^k}\}`$ and some bipartite pure state $`|\mathrm{\Psi }^{BB^k}`$ giving rise to the state
$`\omega ^{XBC^k}={\displaystyle \underset{x}{}}p(x)|xx|^X𝒩^k(\varphi _x\mathrm{\Psi })).`$ (5.1)
Further, it is sufficient to consider ensembles for which
$$|𝒳|\mathrm{max}\{|A^{}|,|C|\}^{2k}.$$
It should also be noted that this characterization does not apparently lead to a finite computation for determining the capacity regions, as it does not admit a single-letter characterization in general. However, as an application, the following example contains a channel for which this region is additive.
###### Example.
Consider an erasure channel into which Alice inputs a classical bit (or rather, a qubit that will be dephased into the $`|0^A^{},|1^A^{}`$ basis), while Bob inputs a qubit. If Alice inputs $`|0^A^{}`$, Charlie receives Bob’s qubit without error. If Alice inputs $`|1^B^{}`$, Charlie receives a pure erasure state $`|e^C`$ which is orthogonal to the degrees of freedom of Bob’s input state. The cq capacity region of this channel is equal to the collection of pairs of nonnegative cq rates $`(R,Q)`$ which satisfy
$`R`$ $``$ $`H(p)`$
$`Q`$ $``$ $`12p`$
for some $`0p\frac{1}{2}`$. This region is pictured in Figure 5.1.
###### Proof.
In Section 9.1, we prove this for the more general case where Bob inputs a $`d`$-level quantum system. ∎
###### Remark.
It is also possible to characterize $`𝒞𝒬(𝒩)`$ as a regularized union of pentagons, a form which is analogous to the result of for classical multiple access channels. As we do not yet know an example of a channel for which this characterization is single-letter (and not equivalent to the rectangle region above), we defer further consideration of this characterization until Chapter 10.
###### Remark.
The proof of the bound on $`|𝒳|`$ is found in the appendix (Section 11.3).
### 5.4 $`𝒬(𝒩)`$ \- quantum-quantum capacity region
The situation in which each sender only attempts to convey quantum infomation to Charlie is described by the quantum-quantum (qq) rate region $`𝒬(𝒩)`$ which consists of rate vectors in $`𝒮(𝒩)`$ of the form $`(0,0,Q_a,Q_b)`$. Our second theorem gives a characterization of $`𝒬(𝒩)`$ as a regularized union of pentagons.
###### Theorem 2.
$`𝒬(𝒩)`$ = the closure of the union of pairs of nonnegative rates $`(Q_a,Q_b)`$ satisfying
$`Q_a`$ $``$ $`{\displaystyle \frac{1}{k}}I_c(ABC^k)_\omega `$
$`Q_b`$ $``$ $`{\displaystyle \frac{1}{k}}I_c(BAC^k)_\omega `$
$`Q_a+Q_b`$ $``$ $`{\displaystyle \frac{1}{k}}I_c(ABC^k)_\omega `$
for some $`k`$ and some bipartite pure states $`|\mathrm{\Psi }_1^{AA^k}`$, $`|\mathrm{\Psi }_2^{BB^k}`$ giving rise to
$`\omega ^{ABC^k}=𝒩^k(\mathrm{\Psi }_1\mathrm{\Psi }_2).`$ (5.2)
###### Example.
An example of a channel for which this region is single-letter is a channel into which Alice and Bob each input a qubit. With probability $`p`$, each of their qubits undergoes a phase flip, or $`180^{}`$ rotation about the $`z`$-axis, before being received by Charlie. Otherwise, Charlie receives both qubits without error. The qq capacity region of this channel is given by a single pentagon, consisting of the pairs of nonnegative qq rates $`(Q_a,Q_b)`$ which satisfy
$`Q_a`$ $``$ $`1`$
$`Q_b`$ $``$ $`1`$
$`Q_a+Q_b`$ $``$ $`2H(p).`$
###### Proof.
See Section 9.4. ∎
###### Remark.
There does not appear to be any obstacle preventing application of the methods used in this paper to prove many-sender generalizations of Theorems 1 and 2. For simplicity, we have focused on the situations with two senders.
###### Remark.
Contrary to the corresponding result for classical multiple access channels, the regions of Theorems 1 and 2 do not require convexification. That this follows from the multi-letter nature of the regions will be demonstrated in the appendix (Section 11.2).
## Chapter 6 Supplementary results
In this chapter, we collect a number of auxiliary results which will be used to prove the main theorems. The first section contains some relationships satisfied by the distance measures of trace distance and fidelity which will comprise the machinery used to prove the coding theorems. The main novel contribution of that section is the statement and proof of Lemma 2. The next section contains other lemmas, proved elsewhere, which will needed later. In the third section we review strong subadditivity of quantum entropy, and explore a number of its consequences. These include quantum versions of the classical data processing inequality, as well as the fact that conditioning decreases conditional quantum entropy or equivalently, increases coherent information. We also obtain a particularly elegant proof of the Holevo bound on the accessible information of an ensemble of quantum states.
### 6.1 Further properties of distance measures
We first collect some relevant results which will be used in what follows, starting with some relationships between our distance measures. If $`\rho `$ and $`\sigma `$ are density matrices defined on the same (or isomorphic) Hilbert spaces, set
$$F=F(\rho ,\sigma )\text{ and }T=|\rho \sigma |_1.$$
Then, the following inequalities hold (see e.g. )
$`1\sqrt{F}`$ $`T/2`$ $`\sqrt{1F},`$ (6.1)
$`1T`$ $`F`$ $`1T^2/4.`$ (6.2)
From these inequalities, we can derive the following more useful relationships
$`F>1ϵ`$ $``$ $`T2\sqrt{ϵ}`$ (6.3)
$`Tϵ`$ $``$ $`F>1ϵ,`$ (6.4)
which are valid for $`0ϵ1.`$ Uhlmann has given the following characterization of fidelity
$$F(\rho ,\sigma )=\underset{|\mathrm{\Psi }_\rho ,|\mathrm{\Phi }_\sigma }{\mathrm{max}}|\mathrm{\Psi }_\rho |\mathrm{\Phi }_\sigma |^2=\underset{|\mathrm{\Psi }_\rho }{\mathrm{max}}|\mathrm{\Psi }_\rho |\mathrm{\Phi }_\sigma |^2$$
where the first maximization is over all purifications of each state, and the second maximization holds for any fixed purification $`|\mathrm{\Phi }_\sigma `$ of $`\sigma .`$ This characterization is useful in two different ways. First, for any two states, it guarantees the existence of purifications of those states whose squared inner product equals the fidelity. Second, one can derive from that characterization the following monotonicity property associated with an arbitrary trace-preserving channel $`𝒩`$,
$`F(\rho ,\sigma )`$ $``$ $`F(𝒩(\rho ),𝒩(\sigma ))`$ (6.5)
An analogous property is shared by the trace distance ,
$`|\rho \sigma |_1`$ $``$ $`\left|𝒩(\rho )𝒩(\sigma )\right|_1,`$ (6.6)
which holds even if $`𝒩`$ is trace-reducing. A simple proof for the trace-preserving case can be found in . These inequalities reflect the fact that completely-positive maps are *contractive* and cannot improve the distinguishability of quantum states; the closer states are to each other, the harder it is to tell them apart. Another useful property will be the multiplicativity of the fidelity under tensor products
$`F(\rho _1\rho _2,\sigma _1\sigma _2)=F(\rho _1,\sigma _1)F(\rho _2,\sigma _2).`$ (6.7)
Since the trace distance comes from a norm, it satisfies the triangle inequality. The fidelity does not come from a norm, but it is possible to derive the following analog by applying (6.1) and (6.2) to the triangle inequality for the trace distance
$`F(\rho _1,\rho _3)12\sqrt{1F(\rho _1,\rho _2)}2\sqrt{1F(\rho _2,\rho _3)}.`$ (6.8)
It will be possible to obtain a sharper triangle-like inequality as a consequence of the following lemma, which states that if a measurement succeeds with high probability on a state, it will also do so on a state which is close to that state in trace distance.
###### Lemma 1.
Suppose that $`\rho ,\sigma ,\mathrm{\Lambda }^{d\times d},`$ where $`\rho `$ and $`\sigma `$ are density matrices, and $`0\mathrm{\Lambda }1.`$ Then, $`\mathrm{Tr}\mathrm{\Lambda }\sigma \mathrm{Tr}\mathrm{\Lambda }\rho |\rho \sigma |_1.`$
###### Proof.
$`\mathrm{Tr}\mathrm{\Lambda }\sigma `$ $`=`$ $`\mathrm{Tr}\mathrm{\Lambda }\rho \mathrm{Tr}\mathrm{\Lambda }(\rho \sigma )`$
$``$ $`\mathrm{Tr}\mathrm{\Lambda }\rho \underset{0\mathrm{\Lambda }1}{\mathrm{max}}2\mathrm{Tr}\mathrm{\Lambda }(\rho \sigma )`$
$`=`$ $`\mathrm{Tr}\mathrm{\Lambda }\rho |\rho \sigma |_1,`$
where the last equality invokes a characterization of the trace distance between density matrices given in Section 2.1.4. ∎
Since $`F(\varphi ,\rho )=\mathrm{Tr}\varphi \rho `$ when $`\varphi `$ is a pure state, a corollary of Lemma 1 is a fact we will refer to as the “special triangle inequality.”
###### Corollary (Special triangle inequality).
$`F(\varphi ,\sigma )F(\varphi ,\rho )|\rho \sigma |_1,`$
The following lemma can be thought of either as a type of transitivity property inherent to any bipartite state with a component near a pure state, or as a partial converse to the monotonicity of fidelity.
###### Lemma 2.
For arbitrary quantum systems $`A`$ and $`B`$, let $`|\varphi ^A`$ be a pure state, $`\rho ^B`$ a density matrix, and $`\mathrm{\Omega }^{AB}`$ a density matrix of the composite system $`AB`$ with partial traces $`\mathrm{\Omega }^A=\mathrm{Tr}_B\mathrm{\Omega }`$ and $`\mathrm{\Omega }^B=\mathrm{Tr}_A\mathrm{\Omega }.`$ Then
$$F(\varphi \rho ,\mathrm{\Omega })1|\rho \mathrm{\Omega }^B|_13\left(1F(\varphi ,\mathrm{\Omega }^A)\right).$$
###### Proof.
We begin by defining the subnormalized density matrix $`\stackrel{~}{\omega }`$ via the equation
$`(\varphi 1)\mathrm{\Omega }(\varphi 1)=\varphi \stackrel{~}{\omega },`$ (6.9)
which we interpret as the upper-left block of $`\mathrm{\Omega }`$, when the basis for $`^{|A|}`$ is chosen in such a way that $`|\varphi =(1,0,\mathrm{},0)^T.`$ Notice that $`F(\varphi ,\mathrm{Tr}_B\mathrm{\Omega })=\mathrm{Tr}\stackrel{~}{\omega }(1ϵ).`$ Writing the normalized state $`\omega =\stackrel{~}{\omega }/(1ϵ),`$ we see that it is close to $`\stackrel{~}{\omega }`$ in the sense that
$`|\omega \stackrel{~}{\omega }|_1`$ $``$ $`ϵ|\stackrel{~}{\omega }|_1`$ (6.10)
$``$ $`ϵ.`$
Now we write
$`\sqrt{F(\varphi \rho ,\mathrm{\Omega })}`$ $`=`$ $`\mathrm{Tr}\sqrt{\sqrt{(\varphi \rho )}\mathrm{\Omega }\sqrt{(\varphi \rho )}}`$ (6.11)
$`=`$ $`\mathrm{Tr}\sqrt{(1\sqrt{\rho })(\varphi 1)\mathrm{\Omega }(\varphi 1)(1\sqrt{\rho })}`$
$`=`$ $`\mathrm{Tr}\sqrt{(1\sqrt{\rho })(\varphi \stackrel{~}{\omega })(1\sqrt{\rho })}`$
$`=`$ $`\mathrm{Tr}\sqrt{\varphi (\sqrt{\rho }\stackrel{~}{\omega }\sqrt{\rho })}`$
$`=`$ $`\mathrm{Tr}\sqrt{\sqrt{\rho }\stackrel{~}{\omega }\sqrt{\rho }}`$
$`=`$ $`\sqrt{F(\stackrel{~}{\omega },\rho )}`$
$`=`$ $`\sqrt{(1ϵ)F(\omega ,\rho )}`$
$``$ $`\sqrt{(1ϵ)(1|\omega \rho |_1)}.`$
The first line is the definition of fidelity and the third follows from (6.9). The last equality relies on the fact that the fidelity, as we’ve defined it, is linear in either of its two inputs, while the inequality follows from (6.2).
Noting that $`\mathrm{\Omega }^B\stackrel{~}{\omega }`$, we define another positive operator $`\omega ^{}=\mathrm{\Omega }^B\stackrel{~}{\omega },`$ which satisfies $`\mathrm{Tr}\omega ^{}ϵ`$ and can be interpreted as the sum of the rest of the diagonal blocks of $`\mathrm{\Omega }.`$ The trace distance in the last line above can be bounded via double application of the triangle inequality as
$`|\rho \omega |_1`$ $``$ $`|\rho (\rho \omega ^{})|_1+|(\rho \omega ^{})\stackrel{~}{\omega }|_1+|\stackrel{~}{\omega }\omega |_1`$ (6.12)
$``$ $`\mathrm{Tr}\omega ^{}+\left|\rho \mathrm{\Omega }^B\right|_1+ϵ`$
$``$ $`\left|\rho \mathrm{\Omega }^B\right|_1+2ϵ,`$
where the second line follows from (6.10). Combining (6.11) with (6.12), we obtain
$`F(\varphi \rho ,\mathrm{\Omega })`$ $``$ $`(1ϵ)(1|\rho \mathrm{\Omega }^B|_12ϵ)`$
$``$ $`1\left|\rho \mathrm{\Omega }^B\right|_13ϵ.`$
### 6.2 Other useful lemmas
This continuity lemma from shows that if two bipartite states are close to each other, the difference between their associated coherent informations is small.
###### Lemma 3 (Continuity of coherent information).
Let $`\rho ^{AB}`$ and $`\sigma ^{AB}`$ be two states of a finite-dimensional bipartite system $`AB`$ satisfying $`|\rho \sigma |_1ϵ`$. Then
$$|I_c(AB)_\rho I_c(AB)_\sigma |2H(ϵ)+4\mathrm{log}|A|ϵ,$$
where $`H(ϵ)`$ is the binary entropy function.
Next is Winter’s “gentle measurement” lemma , which implies that a measurement which is likely to be successful in identifying a state tends not to significantly disturb that state.
###### Lemma 4 (Gentle measurement).
Let a density matrix $`\rho ^A`$ be given, where $`|A|`$ is finite. If $`\mathrm{\Lambda }^{|A|\times |A|}`$ is nonnegative with spectrum bounded above by 1, then
$$\mathrm{Tr}\rho \mathrm{\Lambda }1ϵ$$
implies
$$\left|\sqrt{\mathrm{\Lambda }}\rho \sqrt{\mathrm{\Lambda }}\rho \right|_1\sqrt{8ϵ}.$$
We will also need a lemma from classical information theory which bounds the conditional entropy of two random variables with the same support in terms of the probability they are different.
###### Lemma 5 (Fano’s inequality).
Let $`M`$,$`\widehat{M}`$ be $``$-valued random variables, and write $`P_e=\mathrm{Pr}\{M\widehat{M}\}`$. Then
$$H(M|\widehat{M})H(P_e)+P_e\mathrm{log}||.$$
###### Proof.
See . ∎
### 6.3 Strong subadditivity and its consequences
In this section, we recall an inequality which holds for *any* tripartite quantum system $`ABC`$. This inequality goes by the name *strong subadditivity*, and was originally proved in , stating that
$$H(AB)+H(BC)H(B)+H(ABC).$$
(6.13)
As much has been written about the proof of strong subadditivity of quantum entropy (see e.g. ), we will not discuss the proof of the theorem here. Rather, we will endeavor to show here how strong subadditivity can be used as a mathematical “hammer of Thor,” enabling short and elegant proofs of many known entropy inequalities in quantum information theory. In fact, many of these results will turn out to be *equivalent* to strong subadditivity, in the sense that the latter is easily derivable from many of them.
Begin by subtracting $`H(B)+H(BC)`$ from either side of (6.13) to yield
$$H(A|B)H(A|BC).$$
(6.14)
This inequality can be interpreted as a demonstration that *conditioning reduces entropy*. Collecting the terms on a single side yields the compact formula
$$I(A;B|C)0,$$
showing that the quantum conditional mutual information is always positive. Note that in the classical case, these inequalities are quite simple to prove . For instance, positivity of $`I(X;Z|Y)`$ follows from positivity of mutual information which, in turn, is a consequence of positivity of the Kullback-Leibler distance $`D(P||Q)`$.
Classically, it is simple to show that $`I(X;Z|Y)=0`$ if and only if $`XYZ`$ forms a Markov chain in that order. Necessary and sufficient conditions for saturation of quantum strong subadditivity were recently determined in , who showed that $`I(A;C|B)=0`$ if and only if
$$\rho ^{ABC}=\underset{x}{}p_x\rho _x^{AB_x^A}\rho _x^{B_x^CC}$$
where $`B=_xB_x^AB_x^C`$. In other words, if and only if there is a local measurement that can be performed on $`B`$ which determines $`x`$ without disturbing the global state. Conditioned on knowing $`x`$, the global system is in a product state. Such a measurement is commonly known as a “which path measurement.”
Recall the definition of coherent information as
$$I_c(AB)=H(A|B).$$
Simply reexpressing (6.14) in terms of coherent information yields the inequality
$$I_c(ABC)I_c(AC),$$
(6.15)
which can be interpreted in light of the quantum capacity theorem as saying that losing access to part of the output of a quantum channel can only decrease capacity. Observe that a similar property is obeyed by the classical mutual information, namely that
$$I(X;YZ)I(X;Y).$$
More generally, coherent information can be shown to obey an analog of the classical data processing inequality (see e.g. ), which says that if $`XYZ`$ is a Markov chain, then
$$I(X;Y)I(X;Z).$$
A quantum version of the data processing inequality can be proved easily from strong subadditivity.
###### Lemma 6 (Quantum data processing inequality).
Let a bipartite density matrix $`\rho ^{AB}`$ and a channel $`𝒩:BC`$ be given. Then
$$I_c(AB)_\rho I_c(AC)_{𝒩(\rho )}.$$
###### Proof.
Choose any isometric extension $`𝒰_𝒩:BCE`$ of $`𝒩`$. Then
$`I_c(AB)_\rho `$ $`=`$ $`I_c(ACE)_{𝒰_𝒩(\rho )}`$
$``$ $`I_c(AC)_{𝒩(\rho )},`$
where the first step is because isometries preserve entropy, while the second is by (6.15). ∎
It is thus apparent that post-processing of $`B`$ can never increase coherence with $`A`$. It is also possible to derive strong subadditivity from data processing, by taking $`𝒩=\mathrm{Tr}_C`$, so the data processing inequality is another equivalent way to express strong subadditivity.
The quantum data processing inequality can be used to derive a more direct analog of the classical data processing inequality, dealing with quantum mutual information rather than coherent information. A simple corollary of Lemma 6 is
###### Corollary (Quantum mutual information data processing inequality).
With the same conditions as in Lemma 6,
$$I(A;B)_\rho I(A;C)_{𝒩(\rho )}.$$
###### Proof.
The conclusion of Lemma 6 can be rewritten in terms of conditional entropies as
$$H(A|B)_\rho H(A|C)_{𝒩(\rho )}.$$
Adding $`H(A)`$ to each side yields the required inequality. ∎
As a simple consequence of this corollary, we obtain a completely elementary proof of the Holevo bound , an essential step in the converse part of the HSW capacity theorem.
###### Lemma 7 (Holevo bound).
Let a cq state
$$\rho ^{XB}=\underset{x}{}p(x)|xx|^X\rho _x^B$$
be given. For any measurement on $`B`$ with POVM $`\{\mathrm{\Lambda }_y\}_{y𝒴}`$, the following inequality holds.
$$I(X;B)_\rho I(X;Y).$$
###### Proof.
Construct a measuring instrument $`𝓝:BY`$ (as in Section 2.3.8), acting as
$$𝓝(\tau ^B)=\underset{y}{}\left(\mathrm{Tr}\mathrm{\Lambda }_x\tau \right)|yy|^Y.$$
Application of the previous version of the data processing inequality proves the result. ∎
The following inequality from will be useful in Section 9.2, where we give a proof that *degradable* quantum channels have single-letter quantum capacities.
###### Lemma 8 (Joint subadditivity of conditional entropy).
For any quadripartite state on $`ABCD`$, the following entropy inequality applies
$$H(AB|CD)H(A|C)+H(B|D).$$
(6.16)
###### Proof.
Using the original formulation of strong subadditivity (6.13), we may write the following two inequalities:
$`H(ABCD)+H(C)`$ $``$ $`H(AC)+H(BCD)`$
$`H(BCD)`$ $``$ $`H(BD)+H(CD)H(D).`$
Combining these gives
$$H(ABCD)+H(C)H(AC)+H(BD)+H(CD)H(D).$$
Rearranging terms gives the required result. ∎
Observe that this lemma can equivalently be expressed in terms of coherent information as
$$I_c(AC)+I_c(BD)I_c(ABCD).$$
(6.17)
Note that if (6.17) is computed on a state of the form
$$\mathrm{\Omega }^{ABCD}=|\varphi \varphi |^B\omega ^{ACD},$$
it follows that
$$I_c(BD)=H(D)H(BD)=H(D)H(D)=0$$
and
$$I_c(ABCD)=H(CD)H(ABCD)=H(CD)H(ACD)=I_c(ACD),$$
implying that
$$I_c(AC)I_c(ACD),$$
which is just the original strong subadditivity inequality we started with. So we see that strong subadditivity is equivalent to Lemma 8, as well as to the fact that coherent information is superadditive.
## Chapter 7 Entanglement generation capacities
As a first step towards proving the theorems stated in Chapter 5, we introduce a less restrictive communication scenario, *entanglement generation*. While the criterion of strong subspace transmission is analogous to a classical requirement that the maximal error probability be small, the entanglement generation criterion will rather be related to an average error constraint on good codes.
##### classical-quantum scenario
Alice sends classical information to Charlie at rate $`R`$, while Bob sends quantum information at rate $`Q`$. Rather than being required to transmit half of *any* quantum state Bob is presented with, Bob will only need to *create* near maximal quantum correlations with Charlie at rate $`Q`$. To this end, Bob begins by preparing a bipartite pure state $`|\mathrm{{\rm Y}}^{BB^n},`$ entangled between a physical system $`B`$ located in his laboratory, and the $`B^n`$ part of the inputs of $`𝒩^n`$.
At the same time, Charlie will only need to identify Alice’s classical message with a low average error probability, averaged over all of Alice’s classical messages. As with strong subspace transmission, Charlie’s post-processing procedure will be modeled by a quantum instrument. While the outer bound provided by our converse theorem will apply to any decoding modeled by an instrument, the achievability proof will require a less general approach, consisting of the following steps.
In order to ascertain Alice’s message $`M`$, Charlie first performs some measurement on $`C^n`$, whose statistics are given by a POVM $`\{\mathrm{\Lambda }_m\}_{m2^{nR}}`$. We let the result of that measurement be denoted $`\widehat{M},`$ his declaration of the message sent by Alice. Based on the result of that measurement, he will perform one of $`2^{nR}`$ decoding operations $`𝒟_m^{}:C^n\widehat{B}.`$ These two steps can be mathematically combined to define a quantum instrument $`𝓓:C^n\widehat{M}\widehat{B}`$ with (trace-reducing) components
$$𝒟_m:\tau 𝒟_m^{}(\sqrt{\mathrm{\Lambda }_m}\tau \sqrt{\mathrm{\Lambda }_m}).$$
The instrument acts as
$$𝓓:\tau \underset{m=1}{\overset{2^{nR}}{}}|mm|^{\widehat{M}}𝒟_m(\tau ),$$
and induces the trace preserving map $`𝒟:C^n\widehat{B}`$, acting according to
$$𝒟:\tau \mathrm{Tr}_{\widehat{M}}𝓓(\tau )=\underset{m=1}{\overset{2^{nR}}{}}𝒟_m(\tau ).$$
We again remark that this is the most general decoding procedure required of Charlie. Any situation in which he were to iterate the above steps by measuring, manipulating, measuring again, and so on, is asymptotically just as good as a single instance of the above mentioned protocol. This is because the inner and outer bounds provided by the coding theorem and converse coincide. $`(\{\varphi _m\}_{m2^{nR}},\mathrm{{\rm Y}}^{BB^n},𝓓)`$ will be called an $`(R,Q,n,ϵ)`$ *cq entanglement generation code* for the channel $`𝒩`$ if
$`2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}P_s^{\text{eg}}(m,\mathrm{{\rm Y}})1ϵ,`$ (7.1)
where
$`P_s^{\text{eg}}(m,\mathrm{{\rm Y}})=F(|m|\mathrm{\Phi }^{B\widehat{B}},𝓓𝒩^n(\varphi _m^{A^n}\mathrm{{\rm Y}}^{BB^n})).`$ (7.2)
We will say that $`(R,Q)`$ is an *achievable cq rate pair for entanglement generation* if there exists a sequence of $`(R,Q,n,ϵ_n)`$ cq entanglement generation codes with $`ϵ_n0`$. The capacity region $`𝒞𝒬_{\text{eg}}(𝒩)`$ is defined to be the closure of the collection of all achievable cq rate pairs for entanglement generation.
##### quantum-quantum scenario
As above, Alice and Bob are no longer required to transmit arbitrary quantum correlations with which they are presented. Rather, each has the goal of creating near-maximal entanglement with Charlie. For encoding, Alice and Bob respectively prepare the states $`|\mathrm{{\rm Y}}_1^{AA^n}`$ and $`|\mathrm{{\rm Y}}_2^{BB^n},`$ entangled with the $`A^n`$ and $`B^n`$ parts of the inputs of $`𝒩^n`$. Their goal is to do this in such a way so that Charlie, after applying a suitable decoding operation $`𝒟:C^n\widehat{A}\widehat{B}`$, can hold the $`\widehat{A}\widehat{B}`$ part of a state which is close to $`|\mathrm{\Phi }_1^{A\widehat{A}}|\mathrm{\Phi }_2^{B\widehat{B}}`$. Formally, $`(\mathrm{{\rm Y}}_1^{AA^n},\mathrm{{\rm Y}}_2^{BB^n},𝒟)`$ is an $`(R,Q,n,ϵ)`$ *qq entanglement generation code* for the channel $`𝒩`$ if
$`F(\mathrm{\Phi }_1\mathrm{\Phi }_2,𝒟𝒩^n(\mathrm{{\rm Y}}_1\mathrm{{\rm Y}}_2))1ϵ.`$ (7.3)
$`(R,Q)`$ is an *achievable qq rate pair for entanglement generation* if there is a sequence of $`(R,Q,n,ϵ_n)`$ qq entanglement generation codes with $`ϵ_n0`$. The capacity region $`𝒬_{\text{eg}}(𝒩)`$ is the closure of the collection of all such achievable rates.
### 7.1 The coding theorems
For any quantum multiple access channel $`𝒩:A^{}B^{}C`$, we first prove that the single-letter regions $`𝒞𝒬^{(1)}(𝒩)`$ and $`𝒬^{(1)}(𝒩)`$, defined as the restrictions to $`k=1`$ of the respective characterizations from Sections 5.3 and 5.4, are respectively contained in $`𝒞𝒬_{\text{eg}}(𝒩)`$ and in $`𝒬_{\text{eg}}(𝒩)`$. It will then follow that
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k}𝒞𝒬^{(1)}(𝒩^k)𝒞𝒬_{\text{eg}}(𝒩)\text{and}\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k}𝒬^{(1)}(𝒩^k)𝒬_{\text{eg}}(𝒩)$$
by applying the coding theorems to extensions $`𝒩^k`$ of $`𝒩`$.
###### Proof of Theorem 1 (coding theorem).
Our method of proof for the coding theorem will work as follows. We will employ random HSW codes and random entanglement generation codes to ensure that the average state at the input of $`𝒩^n`$ is close to a product state. Each sender will utilize a code designed for the product channel induced by the other’s random input, whereby existing coding theorems for product channels will be invoked. The quantum code used will be one which achieves the capacity of a modified channel, in which the classical input is copied, without error, to the output of the channel. As the random HSW codes will exactly induce a product state input, the existence of these quantum codes will follow directly from Proposition 2.
The random HSW codes will be those which exist for product channels. As random entanglement generation codes exist with average code density matrix arbitrarily close to a product state, this will ensure that the resulting output states are distinguishable with high probability. Furthermore, obtaining the classical information will be shown to cause but a small disturbance in the overall joint quantum state of the system. As we will show, it is possible to mimic the channel for which the quantum code is designed by placing the identities of the estimated classical message states into registers appended to the outputs of each channel in the product.
The decoder for the modified channel will then be shown to define a quantum instrument which satisfies the success condition for a cq entanglement transmission code, on average. This feature will then be used to infer the existence of a particular, deterministic code which meets the same requirement.
Fix a pure state ensemble $`\{p(x),|\varphi _x^A^{}\}`$ and a bipartite pure state $`|\mathrm{\Psi }^{BB^{}}`$ which give rise to the cq state
$$\omega ^{XBC}=\underset{x}{}p(x)|xx|^X(1^B𝒩)(\varphi _x^A^{}\mathrm{\Psi }^{BB^{}}),$$
which has the form of (5.1). Define $`\rho _1^A^{}=_xp(x)\varphi _x`$ and $`\rho _2^B^{}=\mathrm{Tr}_B\mathrm{\Psi }`$. We will demonstrate the achievability of the corner point $`(I(X;C),I_c(BCX))_\omega `$ by showing that for every $`ϵ,\delta >0,`$ if $`R=I(X;C)_\omega \delta `$ and $`Q=I_c(BCX)_\omega \delta `$, there exists an $`(R,Q,n,ϵ)`$ cq entanglement generation code for the channel $`𝒩`$, provided that $`n`$ is sufficiently large. The rest of the region will follow by timesharing.
For encoding, Alice will choose $`2^{nR}`$ sequences $`𝒞=\{X^n(m)\}_{m2^{nR}}`$, i.i.d. according to the product distribution $`p(x^n)=_{i=1}^np(x_i)`$. As each sequence corresponds to a preparation of channel inputs $`|\varphi _m^{A^n}=|\varphi _{X_1(m)}\mathrm{}|\varphi _{X_n(m)},`$ the expected average density operator associated with Alice’s input to the channel is precisely
$$𝔼_𝒞2^{nR}\underset{m=1}{\overset{2^{nR}}{}}|\varphi _m\varphi _m|=\underset{x^n}{}p(x^n)|\varphi _{x^n}\varphi _{x^n}|=\rho _1^n.$$
Define a new channel $`𝓝^{}:B^{}C\widehat{X}`$ (which is also an instrument) by
$$𝓝^{}:\rho \underset{x}{}p(x)𝒩(\varphi _x\rho )|xx|^{\widehat{X}},$$
This can be interpreted as a channel which reveals the identity of Alice’s input state to Charlie, with the added assumption that Alice chooses her inputs at random. Alternatively, one can view this as a channel with state information available to the receiver, where nature is randomly choosing the “state” $`x`$ at Alice’s input. By Proposition 2, there exists a $`(Q,n,ϵ)`$ random entanglement generation code $`\{q_\beta ,|\mathrm{{\rm Y}}^\beta ^{AA^n},𝒟^\beta \}`$ for the channel $`𝓝^{}`$ with average code density operator $`\varrho ^{B^n}=_\beta q_\beta \mathrm{Tr}_A\mathrm{{\rm Y}}^\beta `$ satisfying
$$|\varrho \rho _2^n|_1ϵ.$$
In what follows, we will use the shorthand $`|\mathrm{{\rm Y}}`$ for the random vector which takes the value $`|\mathrm{{\rm Y}}^\beta `$ with probability $`q_\beta `$. We further abbreviate
$$𝔼_\beta (\mathrm{{\rm Y}})\underset{\beta }{}q_\beta (\mathrm{{\rm Y}}^\beta ),$$
where $``$ is any function of the random vector $`\mathrm{{\rm Y}}`$.
Now, by Proposition 1, for the channel $`𝒩_1:\rho 𝒩(\rho \rho _2)`$ which would result if Bob’s average code density operator were *exactly* equal to $`\rho _2^n,`$ there exists a decoding POVM $`\{\mathrm{\Lambda }_m\}_{m2^{nR}}`$ which would identify Alice’s index $`m`$ with expected average probability of error less than $`ϵ`$, in the sense that
$$𝔼_𝒞2^{nR}\underset{m=1}{\overset{2^{nR}}{}}\mathrm{Tr}\mathrm{\Lambda }_m\tau _m^{}1ϵ,$$
where
$$\tau _m^{}=𝒩^n(\varphi _m\rho _2^n).$$
By the symmetry of the random code construction, we utilize (4.3) to write this as
$$𝔼_𝒞\mathrm{Tr}\mathrm{\Lambda }_1\tau _1^{}1ϵ.$$
Define the *actual* output of the channel corresponding to $`M=m`$ as
$$\tau _m=𝒩^n(\varphi _m\mathrm{Tr}_A\mathrm{{\rm Y}}),$$
as well as its extension
$$\xi _m^{BC^n}=𝒩^n(\varphi _m\mathrm{{\rm Y}}),$$
where $`|\mathrm{\Phi }^{B\stackrel{~}{B}}`$ is the maximally entangled state which Bob is required to transmit. Note that
$$𝔼_\beta \tau _m=𝔼_\beta \mathrm{Tr}_B\xi _m=𝒩^n(\varphi _m\varrho ).$$
It follows from monotonicity of trace distance that
$$\left|𝔼_\beta \tau _1\tau _1^{}\right|_1ϵ,$$
which, together with Lemma 1, implies that
$`𝔼_𝒞2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}\mathrm{Tr}\mathrm{\Lambda }_m𝔼_\beta \tau _m=𝔼_{𝒞\beta }\mathrm{Tr}\mathrm{\Lambda }_1\tau _112ϵ.`$
This allows us to bound the expected probability of correctly decoding Alice’s message as
$`𝔼_{𝒞\beta }\mathrm{Tr}(1\mathrm{\Lambda }_1)\xi _112ϵ.`$ (7.4)
In order to decode, Charlie begins by performing the measurement $`\{\mathrm{\Lambda }_m\}_{m2^{nR}}.`$ He declares Alice’s message to be $`\widehat{M}=m`$ if measurement result $`m`$ is obtained. Charlie will then attempt to simulate the channel $`𝓝^n`$, by associating a separate classical register $`\widehat{X}_i`$ to each channel $`𝒩:A_i^{}C_i`$ in the product, preparing the states $`|X_i(m)^{\widehat{X}_i}`$, for each $`1in`$. Additionally, he stores the result of the measurement in the system $`\widehat{M}`$, his declaration of the message intended by Alice. This procedure results in the global state
$$\mathrm{\Gamma }^{BC^n\widehat{X}^n\widehat{M}}=\underset{m=1}{\overset{2^{nR}}{}}\left(1\sqrt{\mathrm{\Lambda }_m}\right)\xi _1\left(1\sqrt{\mathrm{\Lambda }_m}\right)|X^n(m)X^n(m)|^{\widehat{X}^n}|mm|^{\widehat{M}}.$$
Let $`\mathrm{\Theta }^{BC^n\widehat{X}^n}=\mathrm{Tr}_{\widehat{M}}\mathrm{\Gamma }`$. If Charlie was able to perfectly reconstruct Alice’s classical message, $`\mathrm{\Gamma }`$ would instead be
$$\mathrm{\Gamma }^{}=\xi _1|X^n(1)X^n(1)|^{\widehat{X}^n}|11|^{\widehat{M}},$$
with $`\mathrm{\Theta }^{}=\mathrm{Tr}_{\widehat{M}}\mathrm{\Gamma }^{}`$. When averaged over Alice’s random choice of HSW code, $`\mathrm{\Theta }^{}`$ is precisely equal to the state which would arise via the action of the modified channel $`𝓝^{}.`$ This is because
$`𝔼_𝒞\mathrm{\Theta }^{}`$ $`=`$ $`{\displaystyle \underset{x^n}{}}p(x^n)\xi _{x^n}|x^nx^n|^{\widehat{X}^n}`$ (7.5)
$`=`$ $`𝓝^n(\mathrm{{\rm Y}}),`$
where we have written the joint state which results when Alice prepares $`\varphi _{x^n}`$ as
$$\xi _{x^n}^{BC^n}=𝒩^n(\varphi _{x^n}\mathrm{{\rm Y}}).$$
However, our choice of a good HSW code ensures that he can almost perfectly reconstruct Alice’s message. A consequence of this will be that the two states $`\mathrm{\Theta }`$ and $`\mathrm{\Theta }^{}`$ are almost the same, as we will now demonstrate.
In what follows, we will need to explicitly keep track of the randomness in our codes, by means of superscripts which are to be interpreted as indexing the deterministic codes which occur with the probabilities $`p_𝒞`$ and $`q_\beta `$. Rewriting (7.4) as
$$\underset{𝒞\beta }{}p_𝒞q_\beta \mathrm{Tr}\left(1\mathrm{\Lambda }_1^𝒞\right)\xi _1^{𝒞\beta }12ϵ,$$
it is clear that we may write
$$\mathrm{Tr}\left(1\mathrm{\Lambda }_1^𝒞\right)\xi _1^{𝒞\beta }1ϵ_{𝒞\beta },$$
for positive numbers $`\{ϵ_{𝒞\beta }\}`$ chosen to satisfy
$$\underset{𝒞\beta }{}p_𝒞q_\beta ϵ_{𝒞\beta }=2ϵ.$$
By the gentle measurement lemma,
$$\left|\left(1\sqrt{\mathrm{\Lambda }_1^𝒞}\right)\xi _1^{𝒞\beta }\left(1\sqrt{\mathrm{\Lambda }_1^𝒞}\right)\xi _1^{𝒞\beta }\right|_1\sqrt{8ϵ_{𝒞\beta }},$$
and thus, by the concavity of the square root function,
$`𝔼_{𝒞\beta }\left|\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\right|_1`$
$`=`$ $`{\displaystyle \underset{𝒞\beta }{}}p_𝒞q_\beta \left|\left(1\sqrt{\mathrm{\Lambda }_1^𝒞}\right)\xi _1^{𝒞\beta }\left(1\sqrt{\mathrm{\Lambda }_1^𝒞}\right)\xi _1^{𝒞\beta }\right|_1`$
$``$ $`4\sqrt{ϵ}.`$
Along with (7.4) and monotonicity with respect to $`\mathrm{Tr}_{\widehat{M}}`$, this estimate lets us write
$`𝔼_{𝒞\beta }|\mathrm{\Theta }\mathrm{\Theta }^{}|_1`$ $``$ $`𝔼_{𝒞\beta }|\mathrm{\Gamma }\mathrm{\Gamma }^{}|_1`$
$`=`$ $`𝔼_{𝒞\beta }\left|\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\right|_1`$
$`+𝔼_{𝒞\beta }{\displaystyle \underset{m=2}{\overset{2^{nR}}{}}}\left|\left(1\sqrt{\mathrm{\Lambda }_m}\right)\xi _1\left(1\sqrt{\mathrm{\Lambda }_m}\right)\right|_1`$
$`=`$ $`𝔼_{𝒞\beta }\left|\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\left(1\sqrt{\mathrm{\Lambda }_1}\right)\xi _1\right|_1`$
$`+𝔼_{𝒞\beta }{\displaystyle \underset{m=2}{\overset{2^{nR}}{}}}\mathrm{Tr}(1\mathrm{\Lambda }_m)\xi _1`$
$``$ $`4\sqrt{ϵ}+2ϵ`$
$``$ $`5\sqrt{ϵ},`$ (7.7)
provided that $`ϵ\frac{1}{2}.`$ Since the the entanglement fidelity is linear in $`𝒟(\mathrm{\Theta }),`$ which is itself linear in $`\mathrm{\Theta },`$ we can also use the special triangle inequality to write
$`F(|\mathrm{\Phi },𝒟(𝔼_{𝒞\beta }\mathrm{\Theta }))`$ $`=`$ $`F(|\mathrm{\Phi },𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta }))`$
$``$ $`F(|\mathrm{\Phi },𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta }^{}))\left|𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta }^{})𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta })\right|_1.`$
Using our earlier observation from (7.5) and the definition of a $`(Q,n,ϵ)`$ entanglement transmission code, we can bound the first term as
$`F(|\mathrm{\Phi },𝒟(𝔼_𝒞\mathrm{\Theta }^{}))`$ $`=`$ $`F(|\mathrm{\Phi },𝒟𝓝^n\mathrm{{\rm Y}})`$
$``$ $`1ϵ.`$
An estimate on the second term is obtained via
$`\left|𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta })𝔼_\beta 𝒟(𝔼_𝒞\mathrm{\Theta }^{})\right|_1`$ $``$ $`𝔼_\beta \left|𝒟(𝔼_𝒞\mathrm{\Theta })𝒟(𝔼_𝒞\mathrm{\Theta }^{})\right|_1`$
$``$ $`𝔼_\beta \left|𝔼_𝒞\mathrm{\Theta }𝔼_𝒞\mathrm{\Theta }^{}\right|_1`$
$``$ $`𝔼_{𝒞\beta }\left|\mathrm{\Theta }\mathrm{\Theta }^{}\right|_1`$
$``$ $`5\sqrt{ϵ},`$
where first three lines are by convexity, monotonicity, and convexity once again of the trace norm. The last inequality follows from (7.7). Putting these together gives
$`𝔼_{𝒞\beta }F(|\mathrm{\Phi },𝒟(\mathrm{\Theta }))`$ $``$ $`1ϵ5\sqrt{ϵ}`$ (7.8)
$``$ $`16\sqrt{ϵ}.`$
At last, observe that the final decoded state $`\mathrm{\Omega }`$ (which still depends on both sources of randomness $`𝒞`$ and $`\beta `$) is equal to
$$\mathrm{\Omega }^{B\widehat{B}\widehat{M}}=𝒟(\mathrm{\Gamma }^{BC^n\widehat{X}^n\widehat{M}})𝓓(\xi _1^{BC^n}),$$
implicitly defining the desired decoding instrument $`𝓓:C^n\widehat{B}\widehat{M}`$. The expectation of (7.1) can now be bounded as
$`𝔼_{𝒞\beta }2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}P_s^{\text{eg}}(m)`$ $`=`$ $`𝔼_{𝒞\beta }P_s^{\text{eg}}(1)`$
$`=`$ $`F(|1|\mathrm{\Phi },𝔼_{𝒞\beta }\mathrm{\Omega })`$
$``$ $`1\left|\mathrm{Tr}_{B\widehat{B}}𝔼_{𝒞\beta }\mathrm{\Gamma }|11|\right|_13\left(1F(|\mathrm{\Phi },𝒟(\mathrm{\Theta }))\right)`$
$``$ $`12\sqrt{2ϵ}18\sqrt{ϵ}`$
$``$ $`121\sqrt{ϵ}.`$
The third line above is by Lemma 2. The first estimate in the fourth line follows from (7.4), while the second estimate is by (7.8), together with (6.3). We may now conclude that there are particular values of the randomness indices $`\beta `$ and $`𝒞`$ such that the same bound is satisfied for a deterministic code. We have thus proven that $`(\{\varphi _m\}_{m2^{nR}},,𝓓)`$ comprises a $`(R,Q,n,21\sqrt{ϵ})`$ entanglement generation code. This concludes the coding theorem. ∎
###### Proof of Theorem 2 (coding theorem).
Begin by fixing bipartite pure states $`|\mathrm{\Psi }_1^{A^{\prime \prime }A^{}}`$ and $`|\mathrm{\Psi }_2^{B^{\prime \prime }B^{}}`$ which give rise to the state
$$\omega ^{A^{\prime \prime }B^{\prime \prime }C}=(1^{A^{\prime \prime }B^{\prime \prime }}𝒩^n)(\mathrm{\Psi }_1\mathrm{\Psi }_2),$$
and defining $`\rho _1^A^{}=\mathrm{Tr}_A\mathrm{\Psi }_1`$, $`\rho _2^B^{}=\mathrm{Tr}_B\mathrm{\Psi }_2.`$ Letting $`ϵ,\delta >0`$ be arbitrary, we will show that there exists a $`(Q_a,Q_b,n,ϵ)`$ qq entanglement transmission code where
$$Q_a=I_c(A^{\prime \prime }C)_\omega \delta \text{ and }Q_b=I_c(B^{\prime \prime }A^{\prime \prime }C)_\omega \delta $$
provided that $`Q_a,Q_b0`$. Note that the rates in Theorem 2 will be implied by taking the channel to be $`𝒩^k,`$ with $`\omega ^{ABC^k}`$ defined similarly.
Let us begin by choosing an isometric extension $`𝒰_𝒩:A^{}B^{}CE`$ of $`𝒩`$. Define the ideal channel $`𝒩_1:A^{}C`$ which would effectively be seen by Alice were Bob’s average code density operator exactly equal to $`\rho _2^n`$ as
$$𝒩_1:\tau 𝒩(\tau \rho _2).$$
We now use $`𝒰_𝒩`$ to define a particular isometric extension $`𝒰_{𝒩_1}:A^{}CE^{}`$ of $`𝒩_1`$, where $`E^{}=B^{\prime \prime }E`$, as
$$𝒰_{𝒩_1}:\tau 𝒰_𝒩(\tau \mathrm{\Psi }_2).$$
Observe that Bob’s fake input $`B^{\prime \prime }`$ is treated as part of the environment of Alice’s ideal induced channel. We then further define the channel $`𝒩_2:B^{}A^{\prime \prime }C`$ by
$$𝒩_2:\tau 𝒩(\mathrm{\Psi }_1\tau ).$$
In contrast to the interpretation of $`𝒩_1`$, this may be viewed as the channel which would be seen by Bob if Alice were to input the $`A^{}`$ part of the purification $`|\mathrm{\Psi }_2^{A^{\prime \prime }A^{}}`$ of $`\rho _2^A^{}`$ to her input of the channel and then send the $`A^{\prime \prime }`$ system to Charlie via a noiseless quantum channel. As in the proof of Theorem 1, Charlie will first decode Alice’s information, after which he will attempt to simulate the channel $`𝒩_2`$, allowing a higher transmission rate for Bob than if Alice’s information was treated as noise. Since quantum information cannot be copied, showing that this is indeed possible will require different techniques than were utilized in the previous coding theorem. Although ensembles of random codes will be used in this proof, we introduce the technique of *coherent coding*, in which we pretend that the common randomness is purified. The main advantage of this approach will be that working with states in the enlarged Hilbert space allows monotonicity to be easily exploited in order to provide the estimates we require. Additionally, before we derandomize at the end of the proof, it will ultimately be only Bob who is using a random code. Alice will be able to use any deterministic code from her random ensemble, as Charlie will implement a decoding procedure which produces a global state which is close to that which would have been created had Alice coded with the coherent randomness. To show this, we will first analyze the state which would result if both senders used their full ensembles of codes. Then we show that if Alice uses any code from her ensemble, Charlie can create the proper global state himself, allowing him to effectively simulate $`𝒩_2`$ and ultimately decode both states at the desired rates.
By Proposition 2, for large enough $`n`$, there exists a $`(Q_a,n,ϵ)`$ random entanglement generation code $`(p_{\mathrm{}},|\mathrm{{\rm Y}}_1^{\mathrm{}}^{AA^n},𝒟_1^{\mathrm{}})`$ for the channel $`𝒩_1,`$ where $`Q_a=I_c(\rho _1,𝒩_1)\delta =I_c(A^{\prime \prime }C)\delta .`$ There similarly exists a $`(Q_b,n,ϵ)`$ random entanglement generation code $`(q_m,|\mathrm{{\rm Y}}_2^m^{BB^n},𝒟_2^m)`$ for $`𝒩_2`$, with $`Q_b=I_c(\rho _2,𝒩_2)\delta =I_c(B^{\prime \prime }A^{\prime \prime }C)\delta `$. Proposition 2 further guarantees that these codes can be chosen so that their respective average code density operators
$$\varrho _1^{A^n}=\underset{\mathrm{}}{}p_{\mathrm{}}\mathrm{Tr}_A\mathrm{{\rm Y}}_1^{\mathrm{}}\text{ and }\varrho _2^{B^n}=\underset{m}{}q_m\mathrm{Tr}_B\mathrm{{\rm Y}}_2^m$$
satisfy
$`|\varrho _i\rho _i^n|_1ϵ.`$ (7.9)
Recall that by Proposition 2 we may choose isometric extensions $`𝒰_{𝒟_1}^{\mathrm{}}:C^n\widehat{A}F`$ implementing the $`𝒟_1^{\mathrm{}}`$ from Alice’s random code which satisfy
$`F(|\mathrm{\Phi }_1^{A\widehat{A}}|\lambda ^{FE^n},𝒰_{𝒟_1}^{\mathrm{}}𝒰_{𝒩_1}^n\left(\mathrm{{\rm Y}}_1^{\mathrm{}}\right))1ϵ`$ (7.10)
for every random code index $`\mathrm{}`$ and the same fixed state $`|\lambda ^{FE^n}`$.
Let the code common randomness between Alice and Charlie be held between the systems $`L_A`$ and $`L_C`$, represented by the state
$$\gamma _1^{L_AL_C}=\underset{\mathrm{}}{}p_{\mathrm{}}|\mathrm{}\mathrm{}|^{L_A}|\mathrm{}\mathrm{}|^{L_C},$$
defining a similar state $`\gamma _2^{M_BM_C}`$ for the Bob-Charlie common randomness. For convenience, let us further pretend that $`\gamma _1`$ is part of a pure state
$$|\mathrm{\Gamma }_1^{L_EL_AL_B}=\underset{\mathrm{}}{}\sqrt{p_{\mathrm{}}}|\mathrm{}^{L_E}|\mathrm{}^{L_A}|\mathrm{}^{L_C}.$$
Similarly, let $`\gamma _2`$ by purified by $`|\mathrm{\Gamma }_2^{M_EM_BM_C}`$. Write controlled encoding isometries $`_1:L_AL_AA^n`$ and $`_2:M_BM_BB^n`$ as
$$_1=\underset{\mathrm{}}{}|\mathrm{}|\mathrm{{\rm Y}}_1^{\mathrm{}}\mathrm{}|\text{ and }_2=\underset{m}{}|m|\mathrm{{\rm Y}}_2^mm|.$$
The states which would arise if Alice and Bob each encoded *coherently* are
$`|\mathrm{{\rm Y}}_1^{LAA^n}`$ $``$ $`_1|\mathrm{\Gamma }_1={\displaystyle \underset{\mathrm{}}{}}\sqrt{p_{\mathrm{}}}|\mathrm{}^L|\mathrm{{\rm Y}}_1^{\mathrm{}}`$
$`|\mathrm{{\rm Y}}_2^{MBB^n}`$ $``$ $`_2|\mathrm{\Gamma }_2={\displaystyle \underset{m}{}}\sqrt{q_m}|m^M|\mathrm{{\rm Y}}_2^m.`$
Note that we have abbreviated $`L=L_EL_AL_C`$ and $`M=M_EM_BM_C`$. As each $`|\mathrm{{\rm Y}}_i`$ is a purification of $`\varrho _i`$, together with (7.9), Uhlmann’s theorem tells us that there exist unitaries $`V_1:LAA^{\prime \prime n}`$ and $`V_2:MBB^{\prime \prime n}`$ such that
$`F(V_i|\mathrm{{\rm Y}}_i,|\mathrm{\Psi }_i^n)1ϵ.`$ (7.11)
Further define a corresponding controlled isometric decoder $`𝒰_{𝒟_1}:L_CC^nL_C\widehat{A}F`$ for Alice’s code as
$$𝒰_{𝒟_1}=\underset{\mathrm{}}{}|\mathrm{}\mathrm{}|^{L_C}𝒰_{𝒟_1}^{\mathrm{}}.$$
Let us now imagine that each of Alice and Bob encodes using the coherent common randomness, resulting in a joint pure state $`𝒰_𝒩^n|\mathrm{{\rm Y}}_1|\mathrm{{\rm Y}}_2`$ on $`LAMBC^nE^n`$. If Charlie then applies the full controlled decoder from Alice’s code, the resulting global pure state would be
$$|\mathrm{\Theta }^{LA\widehat{A}MBFE^n}=𝒰_{𝒟_1}𝒰_𝒩^n|\mathrm{{\rm Y}}_1|\mathrm{{\rm Y}}_2.$$
For each $`\mathrm{}`$, let us define an isometry $`𝒪^{\mathrm{}}:B^nA\widehat{A}FE^n`$ as
$$𝒪^{\mathrm{}}=𝒰_{𝒟_1}^{\mathrm{}}𝒰_𝒩^n(\mathrm{{\rm Y}}_1)$$
which we use to define the pure states
$$|\theta _{\mathrm{}}^{A\widehat{A}MFBE^n}=𝒪^{\mathrm{}}|\mathrm{{\rm Y}}_2.$$
These definitions allow us to express
$$|\mathrm{\Theta }=\underset{\mathrm{}}{}\sqrt{p_{\mathrm{}}}|\mathrm{}^L|\theta _{\mathrm{}}.$$
Further writing $`|\lambda ^{}^{FMBE^n}V_2^1|\lambda ^{FB^{\prime \prime n}E^n},`$ the following bound applies
$`F(|\mathrm{\Phi }_1^{A\widehat{A}}|\lambda ^{}^{FMBE^n},|\theta _{\mathrm{}})`$ $`=`$ $`F(|\mathrm{\Phi }_1|\lambda ^{}^{FMBE^n},𝒪^{\mathrm{}}|\mathrm{{\rm Y}}_2)`$
$`=`$ $`F(|\mathrm{\Phi }_1|\lambda ^{FB^{\prime \prime n}E^n},𝒪^{\mathrm{}}V_2|\mathrm{{\rm Y}}_2)`$
$``$ $`12\sqrt{1F(|\mathrm{\Phi }_1|\lambda ^{FB^{\prime \prime n}E^n},𝒪^{\mathrm{}}|\mathrm{\Psi }_2^n))}`$
$`2\sqrt{1F(V_2|\mathrm{{\rm Y}}_2,|\mathrm{\Psi }_2^n)}`$
$``$ $`12\sqrt{1F(|\mathrm{\Phi }_1|\lambda ^{FE^n},𝒰_{𝒟_1}^{\mathrm{}}𝒰_{𝒩_1}^n|\mathrm{{\rm Y}}_1^{\mathrm{}})}2\sqrt{ϵ}`$
$``$ $`14\sqrt{ϵ}.`$
Above, the second equality is because the actions of $`𝒪^{\mathrm{}}`$ and $`V_2`$ commute, the first inequality is by the triangle inequality and monotonicity with respect to $`𝒪^{\mathrm{}},`$ while for the second inequality, we have just rewritten the first term and used (7.11) for the second. The last bound is from (7.10). Observe that we are still free to specify the global phases of the outputs of the $`𝒰_{𝒟_1}^{\mathrm{}}`$ so that the above bound further implies $`\theta _{\mathrm{}}||\mathrm{\Phi }_1|\lambda ^{}(14\sqrt{ϵ})^{1/2}`$ for each $`\mathrm{}`$. Consequently,
$`F(|\mathrm{\Theta },|\mathrm{\Gamma }_1|\mathrm{\Phi }_1|\lambda ^{})`$ $`=`$ $`\left|{\displaystyle \underset{\mathrm{}\mathrm{}^{}}{}}\sqrt{p_{\mathrm{}}p_{\mathrm{}^{}}}\mathrm{}||\mathrm{}^{}\theta _{\mathrm{}}||\mathrm{\Phi }_1|\lambda ^{}\right|^2`$
$`=`$ $`\left|{\displaystyle \underset{\mathrm{}}{}}p_{\mathrm{}}\theta _{\mathrm{}}||\mathrm{\Phi }_1|\lambda ^{}\right|^2`$
$``$ $`14\sqrt{ϵ}.`$
Essentially, the subsystems $`L`$, $`A\widehat{A}`$ and $`MBFE^n`$ of $`|\mathrm{\Theta }`$ are mutually decoupled.
As mentioned earlier, it will be sufficient for Alice to use *any* deterministic code from the random ensemble to encode. Without loss of generality, we assume that Alice chooses to use the first code $`(\mathrm{}=1)`$ in her ensemble. Bob, on the other hand, will need to use randomness to ensure that Alice’s effective channel is close to a product channel. The state on $`AMBC^nE^n`$ which results from these encodings is $`𝒰_𝒩^n|\mathrm{{\rm Y}}_1^1|\mathrm{{\rm Y}}_2`$.
We will now describe a procedure by which Charlie first decodes Alice’s information, then produces a global state which is close to $`|\mathrm{\Theta }`$, making it look like Alice had in fact utilized the coherent coding procedure. This will allow Charlie to apply local unitaries to effectively simulate the channel $`𝒩_2`$ for which Bob’s random code was designed, enabling him to decode Bob’s information as well. These steps will constitute Charlie’s decoding $`𝒟:M_CC^nM_C\widehat{A}\widehat{B}`$, which depends on the Bob-Charlie common randomness. The existence of a deterministic decoder will then be inferred.
Charlie first applies the isometric decoder $`𝒰_{𝒟_1}^1`$, placing all systems into the state $`|\theta _1`$. He then removes his local system $`\widehat{A}`$ (it is important that he keep $`\widehat{A}`$ in a safe place, as it represents the decoder output for Alice’s quantum information) and replaces it with the corresponding parts of the locally prepared pure state $`|\mathrm{\Phi }_1^{A^{}\widehat{A}^{}}.`$ Charlie also locally prepares the state $`|\mathrm{\Gamma }_1^L`$. The resulting state
$$\mathrm{\Theta }^{}=\mathrm{\Gamma }_1^L\mathrm{\Phi }_1^{A^{}\widehat{A}^{}}\mathrm{Tr}_{A\widehat{A}}\theta _1,$$
satisfies
$`F(\mathrm{\Theta }^{},\mathrm{\Theta })`$ $``$ $`1\left|\mathrm{Tr}_{A\widehat{A}}\theta _1\lambda ^{}\right|_1\left|\lambda ^{}\mathrm{Tr}_{LA\widehat{A}}\mathrm{\Theta }\right|_1`$ (7.12)
$`3\left(1F(|\mathrm{\Gamma }|\mathrm{\Phi }_1,\mathrm{Tr}_{MBFE^n}\mathrm{\Theta })\right)`$
$``$ $`12\sqrt{4\sqrt{ϵ}}2\sqrt{4\sqrt{ϵ}}12\sqrt{ϵ}`$
$``$ $`19ϵ^{1/4}`$
whenever $`ϵ12^4`$. The first line combines Lemma 2 and the triangle inequality. The first two estimates in the second line are from applying (6.3) and monotonicity with respect to $`\mathrm{Tr}_{A\widehat{A}}`$ and $`\mathrm{Tr}_{LA\widehat{A}}`$ to the previous two estimates. The last estimate in that line is from monotonicity with respect to the map $`\mathrm{Tr}_{MBFE^n}`$ applied to the previous estimate. Next, Charlie will apply $`V_1𝒰_{D_1}^1`$ to $`\mathrm{\Theta }^{}`$ <sup>1</sup><sup>1</sup>1This operation only acts on Charlie’s local systems, i.e. $`V_1𝒰_{𝒟_1}^1:LA^{}\widehat{A}^{}FA^{\prime \prime n}C^n`$. in order to simulate the channel $`𝒩_2`$. To see that this will work, define $`:LA\widehat{A}FE^nA^{\prime \prime n}C^n`$ as $`\mathrm{Tr}_{E^n}V_1𝒰_{𝒟_1}^1`$ and observe that by monotonicity with respect to $`𝒩^n(\mathrm{{\rm Y}}_2)`$ and (7.11), the states on $`MBA^{\prime \prime n}C^n`$ satisfy
$`F((\mathrm{\Theta }),𝒩_2^n(\mathrm{{\rm Y}}_2))`$ $`=`$ $`F(V_1𝒩^n(\mathrm{{\rm Y}}_1\mathrm{{\rm Y}}_2),𝒩^n(\mathrm{\Psi }_1^n\mathrm{{\rm Y}}_2))`$
$``$ $`F(V_1|\mathrm{{\rm Y}}_1,|\mathrm{\Psi }_1^n)`$
$``$ $`1ϵ.`$
We may now use the triangle inequality and monotonicity with respect to $``$ to combine our last two estimates, yielding
$`F((\mathrm{\Theta }^{}),𝒩_2^n(\mathrm{{\rm Y}}_2))`$ $``$ $`12\sqrt{1F((\mathrm{\Theta }^{}),(\mathrm{\Theta }))}`$ (7.13)
$`2\sqrt{1F((\mathrm{\Theta }),𝒩_2^n(\mathrm{{\rm Y}}_2))}`$
$``$ $`12\sqrt{9ϵ^{1/4}}2\sqrt{ϵ}`$
$``$ $`17ϵ^{1/8}`$
whenever $`ϵ2^{8/3}`$. We have thus far shown that Charlie’s decoding procedure succeeds in simulating the channel $`𝒩_2^n`$, while simultaneously recovering Alice’s quantum information. Charlie now uses the controlled decoder $`𝒟_2:M_CA^{\prime \prime n}C^nM_C\widehat{B}`$ defined as
$$𝒟_2=\underset{m}{}|mm|^{M_C}𝒟_2^m$$
to decode Bob’s quantum information. This entire procedure has defined our decoder $`𝒟:M_CC^nM_C\widehat{A}\widehat{B}`$ which gives rise to a global state $`\mathrm{\Omega }^{A\widehat{A}B\widehat{B}}`$ representing the final output state of the protocol, averaged over Bob’s common randomness. This state satisfies
$`F(|\mathrm{\Phi }_1,\mathrm{Tr}_{B\widehat{B}}\mathrm{\Omega })`$ $``$ $`F(\mathrm{\Theta },\mathrm{\Theta }^{})`$
$``$ $`19ϵ^{1/4},`$
because of monotonicity with respect to $`\mathrm{Tr}_{LMBFE^n}`$ applied to the bound (7.12). By using the triangle inequality, the fact that Bob’s codes are $`ϵ`$-good for each $`m`$, and monotonicity of the estimate (7.13) with respect to $`\mathrm{Tr}_M𝒟_2`$, the global state can further be seen to obey
$`F(|\mathrm{\Phi }_2,\mathrm{Tr}_{A\widehat{A}}\mathrm{\Omega })`$ $`=`$ $`F(|\mathrm{\Phi }_2,\mathrm{Tr}_M𝒟_2(\mathrm{\Theta }^{}))`$
$``$ $`12\sqrt{1F(|\mathrm{\Phi }_2,\mathrm{Tr}_M𝒟_2𝒩_2^n(\mathrm{{\rm Y}}_2))}`$
$`2\sqrt{1F(\mathrm{Tr}_M𝒟_2𝒩_2^n(\mathrm{{\rm Y}}_2),\mathrm{Tr}_M𝒟_2(\mathrm{\Theta }^{}))}`$
$``$ $`12\sqrt{ϵ}2\sqrt{7ϵ^{1/8}}`$
$``$ $`17ϵ^{1/16}`$
as long as $`ϵ2^{16/7}`$. Along with (6.3), a final application of Lemma 2 combines the above two bounds to give
$`F(|\mathrm{\Phi }_1|\mathrm{\Phi }_2,\mathrm{\Omega })`$ $``$ $`1\left|\mathrm{\Phi }_1\mathrm{Tr}_{B\widehat{B}}\mathrm{\Omega }\right|_13\left(1F(|\mathrm{\Phi }_2,\mathrm{Tr}_{A\widehat{A}}\mathrm{\Omega })\right)`$
$``$ $`12\sqrt{9ϵ^{1/4}}21ϵ^{1/16}`$
$``$ $`122ϵ^{1/16},`$
provided that $`ϵ6^{16}`$. Since this estimate represents an average over Bob’s common randomness, there must exist a particular value $`m^{}`$ of the common randomness so that the corresponding deterministic code is at least as good as the random one, thus concluding the coding theorem. ∎
### 7.2 The converse theorems
We will now demonstrate that
$$𝒞𝒬_{\text{eg}}(𝒩)\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k}𝒞𝒬^{(1)}(𝒩^k)\text{and}𝒬_{\text{eg}}(𝒩)\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k}𝒬^{(1)}(𝒩^k),$$
where the single-letter regions $`𝒞𝒬^{(1)}(𝒩)`$ and $`𝒬^{(1)}(𝒩)`$ are those defined at the beginning of the last section.
###### Proof of Theorem 1 (converse).
Suppose there exists a sequence of $`(R,Q,n,ϵ_n)`$ entanglement generation codes with $`ϵ_n0`$. Fixing a blocklength $`n`$, let $`\{\varphi _m\},\mathrm{{\rm Y}}^{BB^n},𝓓`$ comprise the corresponding cq entanglement generation code. The state induced by the encoding is
$$\omega ^{MBC^n}=2^{nR}\underset{m=1}{\overset{2^{nR}}{}}|mm|^M(1^B𝒩^n)(\varphi _m\mathrm{{\rm Y}}).$$
After application of the decoding instrument $`𝓓:C^n\widehat{B}\widehat{M}`$, this state becomes
$$\mathrm{\Omega }^{M\widehat{M}B\widehat{B}}=(1^{MB}𝓓)(\omega ).$$
An upper bound on the classical rate of the code can be obtained as follows:
$`nR`$ $`=`$ $`H(M)_\mathrm{\Omega }`$
$`=`$ $`I(M;\widehat{M})_\mathrm{\Omega }+H(M|\widehat{M})_\mathrm{\Omega }`$
$``$ $`I(M;\widehat{M})_\mathrm{\Omega }+H(ϵ_n)+nRϵ_n`$
$``$ $`I(M;C^n)_\omega +nϵ_n^{}.`$
The first inequality follows from Fano’s inequality (Lemma 5) while in the second we use the Holevo bound (Lemma 7) and define $`ϵ_n^{}=\frac{1}{n}+Rϵ_n`$. The quantum rate of the code is upper bounded as
$`I_c(BC^nM)_\omega `$ $``$ $`I_c(B\widehat{B}M)_\mathrm{\Omega }`$
$``$ $`I_c(B\widehat{B})_\mathrm{\Omega }`$
$``$ $`I_c(B\widehat{B})_\mathrm{\Phi }2H(ϵ_n)8nQ\sqrt{ϵ_n}`$
$`=`$ $`nQnϵ_n^{\prime \prime }.`$
Above, the first two inequalities are consequences of the data processing inequality (Lemma 6), while the last inequality applies a combination of Lemma 3 and (6.3), along with the definition $`ϵ_n^{\prime \prime }=\frac{2}{n}+nQ\sqrt{ϵ_n}`$. Setting $`X=M`$, we have thus proven that
$$R\frac{1}{n}I(X;C^n)+ϵ_n^{},Q\frac{1}{n}I_c(BC^nX)+ϵ^{\prime \prime }_n$$
whenever $`(R,Q)`$ is an achievable cq rate pair for entanglement generation, where $`ϵ_n^{},ϵ_n^{\prime \prime }0`$. It follows that for any achievable rate pair $`(R,Q)`$ and any $`\delta >0`$, we have
$$(R\delta ,Q\delta )\frac{1}{n}𝒞𝒬^{(1)}(𝒩^n)𝒞𝒬(𝒩).$$
Since $`𝒞𝒬(𝒩)`$ is closed by definition, this completes the proof. ∎
###### Proof of Theorem 2 (converse).
Suppose that $`(Q_a,Q_b)`$ is an achievable qq rate pair for entanglement generation. By definition, this means that there must exist a sequence of $`(Q_a,Q_b,n,ϵ_n)`$ entanglement generation codes with $`ϵ_n0`$. Fixing a blocklength $`n`$, let $`|\mathrm{{\rm Y}}_1^{AA^n},|\mathrm{{\rm Y}}_2^{BB^n}`$ and $`𝒟:C^n\widehat{A}\widehat{B}`$ comprise the corresponding encodings and decodings. Define
$$\omega ^{ABC^n}=(1^{AB}𝒩^n)(\mathrm{{\rm Y}}_1\mathrm{{\rm Y}}_2)$$
to be the result of sending the respective $`A^n`$ and $`B^n`$ parts of $`\mathrm{{\rm Y}}_1`$ and $`\mathrm{{\rm Y}}_2`$ through the channel $`𝒩^n`$. Further defining
$$\mathrm{\Omega }^{AB\widehat{A}\widehat{B}}=(1^{AB}𝒟)(\omega )$$
as the corresponding state after decoding, the entanglement fidelity of the code is given by
$`F_{AB}=F(|\mathrm{\Phi }_1|\mathrm{\Phi }_2,\mathrm{\Omega })1ϵ_n.`$ (7.14)
where $`|\mathrm{\Phi }_1^{A\widehat{A}}`$ and $`|\mathrm{\Phi }_2^{B\widehat{B}}`$ are the maximally entangled target states. The sum rate can be bounded as
$`I_c(ABC^n)_\omega `$ $``$ $`I_c(AB\widehat{A}\widehat{B})_\mathrm{\Omega }`$
$``$ $`I_c(AB\widehat{A}\widehat{B})_{\mathrm{\Phi }_1\mathrm{\Phi }_2}2H(ϵ_n)8n(Q_a+Q_b)\sqrt{ϵ_n}`$
$``$ $`n(Q_a+Q_b)nϵ_n^{}.`$
The first step is by the data processing inequality (Lemma 6). The second step uses Lemma 3 and (6.3), along with monotonicity applied to (7.14). The last step has defined $`ϵ_n^{}=\frac{2}{n}8(Q_a+Q_b)\sqrt{ϵ_n}`$ and holds because the binary entropy $`H()`$ is upper bounded by 1. We can bound Alice’s rate $`Q_a`$ by writing
$`I_c(ABC^n)_\omega `$ $``$ $`I_c(AC^n)_\omega `$
$``$ $`I_c(A\widehat{A}\widehat{B})_\mathrm{\Omega }`$
$``$ $`I_c(A\widehat{A})_\mathrm{\Omega }`$
$``$ $`I_c(A\widehat{A})_{\mathrm{\Phi }_1}2H(ϵ_n)8nQ_a\sqrt{ϵ_n}`$
$``$ $`nQ_anϵ_n^{}.`$
The first three steps above are by data processing (Lemma 6). The remaining steps hold for the same reasons as in the previous chain of inequalities. Similarly, Bob’s rate also must satisfy
$$nQ_bI_c(BAC^n)_\omega +nϵ^{}_n.$$
Since $`ϵ_n0`$ implies $`ϵ_n^{}0,`$ this means that for every $`\delta >0`$, any achievable qq rate pair $`(Q_a,Q_b)`$ must satisfy
$$(Q_a\delta ,Q_b\delta )\frac{1}{n}𝒬^{(1)}(𝒩^n)𝒬(𝒩).$$
Since $`𝒬(𝒩)`$ is closed by definition, this completes the proof. ∎
## Chapter 8 Transmission of quantum information
In the previous chapter, we have proven the main theorems for the restricted case in which all quantum communication has been in the sense of *generating* quantum correlations between senders and receiver. The results of this chapter will complete the proofs of the main theorems, by extending the weaker error criteria of entanglement generation (which incidentally, are analogous to a classical requirement on the *average* probability of error) to the stronger requirements of strong subspace transmission in the main theorem statements. As a first step, we demonstrate how the results of the last chapter immediately imply the ability to perform an intermediate task, *entanglement transmission*, where the senders are required to transmit preexisting maximal entanglement, while still adhering to an average error criterion on the classical error. We then show how to use a given entanglement transmission code to construct a strong subspace transmission codes achieving any rates less then those of the original code, while paying a negligible price in fidelity.
### 8.1 Entanglement transmission
##### Classical-quantum scenario
In this scenario, rather than generating entanglement with Charlie, Bob will act to transmit *preexisting* entanglement to him. We assume that Bob is presented with the $`\stackrel{~}{B}`$ part of the maximally entangled state $`|\mathrm{\Phi }^{B\stackrel{~}{B}}.`$ It is assumed that he has complete control over $`\stackrel{~}{B}`$, while he has no access to $`B`$. He will perform a physical operation in order to transfer the quantum information embodied in his system $`\stackrel{~}{B}`$ to the inputs $`B^n`$ of the channel, modeled by an encoding operation $`:\stackrel{~}{B}B^n`$. The goal of this encoding will be to make it possible for Charlie, via post-processing of the information embodied in the system $`C^n`$, to hold the $`\widehat{B}`$ part of a state which is close to that which would have resulted if Bob had sent his system through a perfect quantum channel $`\text{id}:\stackrel{~}{B}\widehat{B}`$. Here, we imagine that $`\stackrel{~}{B}`$ and $`\widehat{B}`$ denote two distinct physical systems with the same number of quantum degrees of freedom. The role of the identity channel is to set up a unitary correspondence, or isomorphism, between the degrees of freedom of $`\stackrel{~}{B}`$ in Bob’s laboratory and those of $`\widehat{B}`$ in Charlie’s. We will often tacitly assume that such an identity map has been specified ahead of time in order to judge how successful an imperfect quantum transmission has been. This convention will be taken for granted many times throughout the paper, wherein specification of an arbitrary state $`|\mathrm{\Psi }^{B\stackrel{~}{B}}`$ will immediately imply specification of the state $`|\mathrm{\Psi }^{B\widehat{B}}=(1^B\text{id})|\mathrm{\Psi }^{B\stackrel{~}{B}}.`$ Decoding is the same as it is for entanglement generation.
$`(\{\varphi _m\}_{m2^{nR}},,𝓓)`$ will be called an $`(R,Q,n,ϵ)`$ *cq entanglement transmission code* for the channel $`𝒩`$ if
$`2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}P_s^{\text{et}}(m)1ϵ,`$ (8.1)
where
$`P_s^{\text{et}}(m)=F(|m|\mathrm{\Phi }^{B\widehat{B}},𝓓𝒩^n(\varphi _m^{A^n}(\mathrm{\Phi }^{B\stackrel{~}{B}})).`$ (8.2)
Achievable rate pairs and the capacity region $`𝒞𝒬_{\text{et}}(𝒩)`$ are defined analogous to those for entanglement generation.
##### Quantum-quantum scenario
Alice and Bob each respectively have control over the $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ parts of the separate maximally entangled states $`|\mathrm{\Phi }_1^{A\stackrel{~}{A}},|\mathrm{\Phi }_2^{B\stackrel{~}{B}}`$, while neither has access to $`A`$ or $`B`$. Alice transfers the correlations in her system to the $`A^n`$ parts of the inputs of $`𝒩^n`$ with an encoding operation $`_1:\stackrel{~}{A}A^n`$. Bob acts similarly with $`_2:\stackrel{~}{B}B^n`$. Their goal is to preserve the respective correlations, so that Charlie can apply a decoding operation $`𝒟:C^n\widehat{A}\widehat{B}`$, in order to end up holding the $`\widehat{A}\widehat{B}`$ part of a state which is close to $`|\mathrm{\Phi }_1^{A\widehat{A}}|\mathrm{\Phi }_2^{B\widehat{B}}`$. Formally, $`(_1,_2,𝒟)`$ is a $`(Q_a,Q_b,n,ϵ)`$ *qq entanglement transmission code* for the channel $`𝒩`$ if
$`F(|\mathrm{\Phi }_1|\mathrm{\Phi }_2,𝒟𝒩^n(_1_2)(\mathrm{\Phi }_1\mathrm{\Phi }_2))1ϵ.`$ (8.3)
Achievable qq rate pairs for entanglement generation and the capacity region $`𝒬_{\text{et}}(𝒩)`$ are defined as in the previous scenario.
### 8.2 Equivalence of entanglement transmission and entanglement generation
#### 8.2.1 $`𝒞𝒬_{\text{eg}}𝒞𝒬_{\text{et}}`$ and $`𝒬_{\text{eg}}𝒬_{\text{et}}`$
###### Proof.
This essentially follows as an artifact of the entanglement generation coding theorem from . There, the input preparation $`|\mathrm{{\rm Y}}^{AA^n}`$ for a $`(Q,n)`$ entanglement generation code is constructed with the particular form
$$|\mathrm{{\rm Y}}^{AA^n}=\frac{1}{\sqrt{2^{nQ}}}\underset{a2^{nQ}}{}|a^A|\varphi _a^{A^n},$$
where the $`\{|\varphi _a\}`$ are orthogonal. Observe that the if the encoder acts on the $`\stackrel{~}{A}`$ part of the maximally entangled state
$$|\mathrm{\Phi }^{A\widehat{A}}=\frac{1}{\sqrt{2^{nQ}}}\underset{a2^{nQ}}{}|a^A|a^{\stackrel{~}{A}}$$
with an encoding isometry $`:\stackrel{~}{A}A^n`$ defined via
$$=\underset{a2^{nQ}}{}|\varphi _a^{A^n}a|^{\stackrel{~}{A}},$$
the identity $`|\mathrm{\Phi }^{A\stackrel{~}{A}}=|\mathrm{{\rm Y}}^{AA^n}`$ holds trivially. It is thus a simple task to modify the proofs of Chapter 7 to instead prove the existence of the entanglement transmission codes described in the previous section. Indeed, if $`(|\mathrm{{\rm Y}},\{\varphi _m\},𝓓)`$ is a $`(R,Q,n,ϵ)`$ cq entanglement generation code, there then exists an encoder $``$ so that $`(,\{\varphi _m\},𝓓)`$ is a $`(R,Q,n,ϵ)`$ cq entanglement transmission code. Identical reasoning shows that to every qq entanglement generation code, there a qq entanglement transmission code with the same parameters. ∎
#### 8.2.2 $`𝒞𝒬_{\text{et}}𝒞𝒬_{\text{eg}}`$
###### Proof.
Suppose there exists an $`(R,Q,n,ϵ)`$ cq entanglement transmission code, consisting of classical message states $`\{|\varphi _m^{A^n}\}_{m2^{nR}},`$ a quantum encoding map $`:\stackrel{~}{B}\widehat{B}`$, and a decoding instrument $`𝓓:C^n\widehat{M}\widehat{B}.`$ Write any pure state decomposition of the encoded state
$$(1^B)(\mathrm{\Phi })=\underset{i}{}p_i|\mathrm{{\rm Y}}_i\mathrm{{\rm Y}}_i|.$$
Then, the success condition (8.1) for a cq entanglement transmission code can be rewritten as
$`1ϵ`$ $``$ $`2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}P_s^{\text{et}}(m)`$ (8.4)
$`=`$ $`2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}F(|\mathrm{\Phi }^{B\widehat{B}},𝒟_m𝒩^n\left(\varphi _m^{A^n}\left({\displaystyle \underset{i}{}}p_i\mathrm{{\rm Y}}_i\right)\right))`$ (8.5)
$`=`$ $`{\displaystyle \underset{i}{}}p_i\left(2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}F(|\mathrm{\Phi }^{B\widehat{B}},𝒟_m𝒩^n(\varphi _m^{A^n}\mathrm{{\rm Y}}_i))\right)`$ (8.6)
$`=`$ $`{\displaystyle \underset{i}{}}p_i\left(2^{nR}{\displaystyle \underset{m=1}{\overset{2^{nR}}{}}}P_s^{\text{eg}}(m,\mathrm{{\rm Y}}_i)\right),`$ (8.7)
so that there is a particular value $`i^{}`$ of $`i`$ for which
$$2^{nR}\underset{m=1}{\overset{2^{nR}}{}}P_s^{\text{eg}}(m,\mathrm{{\rm Y}}_i^{}))1ϵ.$$
Hence, $`(\{|\varphi _m\}_{m2^{nR}},|\mathrm{{\rm Y}}_i^{},𝓓)`$ comprises an $`(R,Q,n,ϵ)`$ cq entanglement generation code. ∎
#### 8.2.3 $`𝒬_{\text{et}}𝒬_{\text{eg}}`$
###### Proof.
Suppose there exists a $`(Q_a,Q_b,n,ϵ)`$ entanglement transmission code $`(_1,_2,𝒟)`$ which transmits the maximally entangled states $`|\mathrm{\Phi }_1,|\mathrm{\Phi }_2`$. As in the cq case, the encoded states can be decomposed as
$$(1^A_1)(\mathrm{\Phi }_1)=\underset{i}{}p_i\mathrm{{\rm Y}}_{1i}$$
and
$$(1^B_2)(\mathrm{\Phi }_2)=\underset{j}{}q_j\mathrm{{\rm Y}}_{2i}.$$
The reliability condition (8.3) can then be rewritten as
$$\underset{ij}{}p_iq_jF(|\mathrm{\Phi }_1|\mathrm{\Phi }_2,𝒟𝒩^n(\mathrm{{\rm Y}}_{1i}\mathrm{{\rm Y}}_{2j}))1ϵ,$$
which implies the existence of a particular pair $`(i^{},j^{})`$ of values of $`(i,j)`$ such that
$$F(|\mathrm{\Phi }_1|\mathrm{\Phi }_2,𝒟𝒩^n(\mathrm{{\rm Y}}_{1i^{}}\mathrm{{\rm Y}}_{2j^{}}))1ϵ.$$
Hence, $`(|\mathrm{{\rm Y}}_{1i^{}},|\mathrm{{\rm Y}}_{2j^{}},𝒟)`$ comprises a $`(Q_a,Q_b,n,ϵ)`$ qq entanglement generation code. ∎
### 8.3 Strong subspace transmission revisited
The criteria of entanglement generation and transmission, both in the cq and qq cases, are directly analogous to the requirement in classical information theory that the average probability of error, averaged over all codewords, be small. However, the requirements imposed in Section 5.2 are analogous to the stronger classical condition that the *maximal* probability of error be small, or that the probability of error for *each* pair of codewords be small. There are examples of classical multiple access channels for which, when each encoder is a deterministic function from the set of the messages to the set of input symbols, the maximal error capacity region is *strictly* smaller than the average error region . However, it is known that if stochastic encoders are allowed (see Problem 3.2.4 in ), the maximal and average error capacity regions are equal.
It is well-known that randomization is not necessary for such an equivalence to hold for single-user channels, as Markov’s inequality implies that a fraction of the codewords with the worst probability of error can be purged, while incurring a negligible loss of rate. The obstacle to utilizing such an approach for classical multiple access channels, and hence for quantum ones as well, is that there is no guarantee that a large enough subset of bad pairs of codewords decomposes as the product of subsets of each sender’s codewords.
A particularly attractive feature of the requirements of Section 5.2 is that they ensure *composability*; when combined with other protocols satisfying analogous criteria, the joint protocol will satisfy similar properties. As an example, recent work on organizing and classifying quantum Shannon-theoretic protocols by means of *resource inequalities* , makes heavy use of such concatenation of quantum information processing protocols.
In the next two subsections, we cast the requirements outlined in Section 5.2 into somewhat simpler forms which are specific to each of the cq and qq cases. We will use these forms in order to prove the equivalences of entanglement transmission and strong subspace transmission in both the cq and qq cases.
#### 8.3.1 classical-quantum scenario
Strong subspace transmission can be considered a more ambitious version of entanglement transmission, whereby rather than requiring Bob to transmit half of a maximally entangled state $`|\mathrm{\Phi }^{B\stackrel{~}{B}},`$ it is instead required that he faithfully transmit the $`\stackrel{~}{B}`$ part, presented to him, of *any* bipartite pure state $`|\mathrm{\Psi }^{B\stackrel{~}{B}},`$ where $`|B|`$ can be any finite number. The reader should note that this constitutes a generalization of the usual subspace transmission , as whenever $`|\mathrm{\Psi }^{B\stackrel{~}{B}}=|\psi ^B|\phi ^{\stackrel{~}{B}}`$, this amounts to requiring that $`|\phi `$ be transmitted faithfully. We further demand that the maximal error probability for the classical messages be small.
As with entanglement transmission, Alice will send classical information at rate $`R`$ by preparing one of $`2^{nR}`$ pure states $`\{|\varphi _m^{A^n}\}_{m2^{nR}}`$. As previously discussed, our more restrictive information transmission constraints can only be met by allowing Alice to employ a stochastic encoding. We assume that Alice begins by generating some randomness, modeled by the random variable $`X`$. To send message $`M=m`$, she prepares a state $`\varphi _{f(m)}`$, where $`f(m)f_X(m)`$ is a random encoding function, depending on the randomness in $`X`$. In the language of Section 5.2, this amounts to the definition of a c $``$ q encoding function $`_1:MA^n`$. Observe that our definition there already allows for randomness to be part of the encoding process.
Bob will apply an encoding $`:\stackrel{~}{B}B^n`$ (this is just his encoding $`_2`$ from Section 5.2 without a classical input), and Charlie will employ a decoding instrument $`𝓓:C^n\widehat{M}\widehat{B}`$. These maps require a more complicated structure than was required for entanglement generation and transmission. Indeed, these will be constructed by means of a protocol, to be described below, out of the entanglement transmission codes which were proved to exist in Section 8.2.1. The success probability for the protocol, conditioned on $`m`$ being sent and $`|\mathrm{\Psi }^{B\stackrel{~}{B}}`$ being presented, can be expressed as
$`P_s(m,\mathrm{\Psi })`$ $`=`$ $`F(|f(m)^{\widehat{M}}|\mathrm{\Psi }^{B\widehat{B}},𝓓𝒩^n\left(\varphi _{f(m)}^{A^n}(\mathrm{\Psi }^{B\stackrel{~}{B}})\right)).`$
We will say that $`(f,X,\{|\varphi _m\}_{m2^{nR}},,𝓓)`$ is an $`(R,Q,n,ϵ)`$ *cq strong subspace transmission code* for the channel $`𝒩`$ if, for every $`m2^{nR}`$ and every $`|\mathrm{\Psi }^{B\stackrel{~}{B}}`$,
$`𝔼_XP_s(m,\mathrm{\Psi })1ϵ.`$ (8.8)
The rate pair $`(R,Q)`$ is an *achievable cq rate pair for strong subspace transmission* if there is a sequence of $`(R,Q,n,ϵ_n)`$ cq random strong subspace transmission codes with $`ϵ_n0`$, and the capacity region $`𝒞𝒬(𝒩)`$ is closure of the collection of all such achievable rates.
#### 8.3.2 quantum-quantum scenario
This scenario is the obvious combination of the relevant concepts from the previous scenario and the qq entanglement transmission scenario. Alice and Bob are respectively presented with the $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ parts of some pure bipartite states $`|\mathrm{\Psi }_1^{A\stackrel{~}{A}}`$ and $`|\mathrm{\Psi }_2^{B\stackrel{~}{B}}`$. As before, we place no restriction on $`|A|`$ and $`|B|`$, other than that they are finite. They employ their respective encodings $`_1`$ and $`_2`$ (which are just the encodings from Section 5.2 without classical inputs), while Charlie decodes with $`𝒟`$. As in the above cq case, the structure of these maps will be more complicated than in the previous two scenarios. $`(_1,_2,𝒟)`$ is then a $`(Q_a,Q_b,n,ϵ)`$ *qq strong subspace transmission code* if
$`F(|\mathrm{\Psi }_1^{A\widehat{A}}|\mathrm{\Psi }_2^{B\widehat{B}},𝒟𝒩^n(_1_2)(\mathrm{\Psi }_1^{A\stackrel{~}{A}}\mathrm{\Psi }_2^{B\stackrel{~}{B}}))1ϵ,`$ (8.9)
for every pair of pure bipartite states $`|\mathrm{\Psi }_1^{A\stackrel{~}{A}}`$ and $`|\mathrm{\Psi }_2^{B\stackrel{~}{B}}`$. Achievable rates and the capacity region $`𝒬(𝒩)`$ are defined as in the cq case.
### 8.4 Equivalence of entanglement transmission and strong subspace transmission
Let us first prove the easy directions. To see that $`𝒞𝒬𝒞𝒬_{\text{et}}`$, note that given a strong subspace transmission code, if Alice uses any deterministic value $`x`$ for her locally generated randomness $`X`$, the average classical error will be equal to the expected maximal classical error of the randomized code. Since the ability to transmit any state includes the maximally entangled case, this completes the claim. The inclusion $`𝒬𝒬_{\text{et}}`$ follows trivially. As any states can be transmitted, this certainly includes the case of a pair of maximally entangled states.
#### 8.4.1 $`𝒞𝒬_{\text{et}}𝒞𝒬`$
###### Proof.
Suppose there exists an $`(R,Q,n,ϵ^2/2)`$ entanglement transmission codes with classical message states $`\{|\varphi _m^{A^n}\}_{m2^{nR}},`$ quantum encoding $`:\stackrel{~}{B}\widehat{B},`$ and decoding instrument $`𝓓:C^n\widehat{M}\widehat{B}`$ with trace-reducing components $`\{𝒟_m:C^n\widehat{B}\}`$, which transmits a maximally entangled state $`|\mathrm{\Phi }^{A\stackrel{~}{A}}`$, where $`|A|<\mathrm{}`$ (although $`|\stackrel{~}{A}|=2^{nQ}`$).
We will initially prove the equivalence by constructing a code which requires two independent sources of shared common randomness $`X`$ and $`Y`$. $`X`$ is assumed to be available to Alice and to Charlie, while $`Y`$ is available to Bob and to Charlie. Then, we will argue that it is possible to eliminate the dependence on the shared randomness, by using the channel to send a negligibly small “random seed”, which can be recycled to construct a code which asymptotically achieves the same performance as the randomized one.
We begin by demonstrating how shared common randomness between Alice and Charlie allows Alice to send any message with low probability of error. Setting $`\mu =2^{nR}`$, let the random variable $`X`$ be uniformly distributed on the set $`\{1,\mathrm{},\mu \}`$. To send message $`M=m`$, Alice computes $`m^{}=m+X`$ modulo $`\mu `$. She then prepares the state $`|\varphi _m^{}`$ for transmission through the channel. Bob encodes the $`\stackrel{~}{B}`$ part of $`|\mathrm{\Phi }^{B\stackrel{~}{B}}`$ with $``$, and each sends appropriately through the channel. Charlie decodes as usual with the instrument $`𝓓`$. Denoting the classical output as $`\widehat{M}^{}`$, his declaration of Alice’s message is then $`\widehat{M}=\widehat{M}^{}X`$ modulo $`\mu `$. Defining the trace-reducing maps $`_m:\stackrel{~}{B}\widehat{B}`$ by
$$_m:\tau 𝒟_m𝒩^n(\varphi _m(\tau )),$$
and the trace-reducing average map as
$$:\tau \frac{1}{\mu }\underset{m=1}{\overset{\mu }{}}_m(\tau ),$$
we can rewrite the success criterion (8.1) for entanglement transmission as
$`F(|\mathrm{\Phi },(\mathrm{\Phi }))1ϵ^2/2,`$
which, together with (6.3), implies that for the identity map $`\text{id}:\stackrel{~}{B}\widehat{B}`$,
$`\left|(\text{id})(\mathrm{\Phi })\right|_1ϵ.`$ (8.10)
The above randomization of the classical part of the protocol can be mathematically expressed by replacing the $`_m`$ with $`_{m+X}`$. As tracing over the common randomness $`X`$ is equivalent to computing the expectation with respect to $`X`$, we see that $`𝔼_X_{m+X}=`$, or rather
$$𝔼_XF(|\mathrm{\Phi },_{m+X}(\mathrm{\Phi }))=F(|\mathrm{\Phi },(\mathrm{\Phi })).$$
It is thus clear that the maximal error criterion for the randomized protocol is equal to the average criterion for the original one.
We continue by randomizing the quantum part of the classically randomized protocol. Setting $`d=2^{nQ}=|\stackrel{~}{B}|,`$ let $`\{U_y\}_{yd^2}`$ be the collection of Weyl unitaries, or generalized Pauli operators, on the $`d`$-dimensional input space. Observe that for any $`\rho `$, acting with a uniformly random choice of Weyl unitary has a completely randomizing effect, in the sense that
$$\frac{1}{d^2}\underset{y=1}{\overset{d^2}{}}U_y\rho U_y^1=\pi _d.$$
Let the random variable $`Y`$ be uniformly distributed on $`\{1,\mathrm{},d^2\}`$. It will be convenient to define the common randomness state
$$\mathrm{{\rm Y}}^{Y_BY_C}=\frac{1}{d^2}\underset{y=1}{\overset{d^2}{}}|yy|^{Y_B}|yy|^{Y_C},$$
where the system $`Y_B`$ is in the possession of Bob, while $`Y_C`$ is possessed by Charlie. Define now the controlled unitaries $`𝒰_B:Y_B\stackrel{~}{B}Y_B\stackrel{~}{B}`$ and $`𝒰_C:Y_C\widehat{B}Y_C\widehat{B}`$ by
$$𝒰_B=\underset{y=1}{\overset{d^2}{}}|yy|^{Y_B}U_y$$
and
$$𝒰_C=\underset{y=1}{\overset{d^2}{}}|yy|^{Y_C}U_y^1.$$
Suppose Bob is given the $`\stackrel{~}{B}`$ part of an arbitrary pure state $`|\mathrm{\Psi }^{B\stackrel{~}{B}}`$, and Alice sends the classical message $`M=m`$. For encoding, Bob will apply $`𝒰_B`$ to the combined system $`\mathrm{{\rm Y}}\mathrm{\Psi }`$. Charlie decodes with $`𝒰_C𝒟`$. If $``$ were equal to the perfect quantum channel $`\text{id}:\stackrel{~}{B}\widehat{B}`$, this procedure would result in the state
$$\frac{1}{d^2}\underset{y=1}{\overset{d^2}{}}|yy|^{Y_B}|yy|^{Y_C}\mathrm{\Psi }.$$
Note that the common randomness is still available for reuse. Abbreviating $`|yy|^Y=|yy|^{Y_B}|yy|^{Y_C}`$, and $`|\mathrm{\Psi }_y^{B\stackrel{~}{B}}=(1^BU_y)|\mathrm{\Psi }`$, we write
$`\sigma ^{YB\stackrel{~}{B}}`$ $`=`$ $`𝒰_B(\mathrm{{\rm Y}}\mathrm{\Psi })`$ (8.11)
$`=`$ $`{\displaystyle \frac{1}{d^2}}{\displaystyle \underset{y=1}{\overset{d^2}{}}}|yy|^Y\mathrm{\Psi }_y.`$ (8.12)
Observe that $`\sigma `$ is an extension of the maximally mixed state $`\pi ^{\stackrel{~}{B}}`$, and can be seen to arise by storing in $`Y`$ the result of a von Neumann measurement along the basis $`\{|y^F\}_{yd^2}`$ on the $`F`$ part of the pure state
$$|\mathrm{\Gamma }^{FB\stackrel{~}{B}}=\frac{1}{d}\underset{y=1}{\overset{d^2}{}}|y^F|\mathrm{\Psi }_y^{B\stackrel{~}{B}}.$$
Since $`\mathrm{Tr}_{R^{}R}\mathrm{\Gamma }=\mathrm{Tr}_{YR}\sigma =\pi ^{\stackrel{~}{B}}`$, $`|\mathrm{\Gamma }`$ is maximally entangled between $`FB`$ and $`\stackrel{~}{B}`$. So, there exists an isometry $`V:BFB`$ such that $`(V1^{\stackrel{~}{B}})|\mathrm{\Phi }^{B\stackrel{~}{B}}=|\mathrm{\Gamma }.`$ This implies that there is a quantum operation $`𝒪:BYB`$ such that $`(𝒪1^{\stackrel{~}{B}})(\mathrm{\Phi })=\sigma `$. Define the trace-reducing map $`𝒯:\stackrel{~}{B}\widehat{B},`$ which represents the coded channel with common randomness accounted for, by
$$𝒯:\tau \mathrm{Tr}_Y𝒰_C𝒰_B(\mathrm{{\rm Y}}\tau ).$$
Recalling our denotation of the noiseless quantum channel $`\text{id}:\stackrel{~}{B}\widehat{B}`$, as well as our convention that id acts as the identity on any system which is not $`\stackrel{~}{B}`$, we now bound
$`1F(|\mathrm{\Psi },𝒯(\mathrm{\Psi }))`$ $``$ $`\left|(𝒯\text{id})(\mathrm{\Psi })\right|_1`$
$``$ $`\left|(𝒰_C𝒰_B\text{id})(\mathrm{{\rm Y}}\mathrm{\Psi })\right|_1`$
$`=`$ $`\left|(\text{id})𝒰_B(\mathrm{{\rm Y}}\mathrm{\Psi })\right|_1`$
$`=`$ $`\left|(\text{id})(\sigma )\right|_1`$
$``$ $`\left|(\text{id})(\mathrm{\Phi })\right|_1`$
$``$ $`ϵ,`$
where the first line is by (6.1) and the second by monotonicity with respect to $`\mathrm{Tr}_Y`$. The third follows from unitary invariance of the trace. The second to last inequality is a consequence of monotonicity with respect to $`𝒪`$, while the last is by (8.10). Note that by monotonicity, this implies that any density matrix $`\mathrm{\Omega }^{B\stackrel{~}{B}}`$ satisfies
$`|𝒯(\mathrm{\Omega })\mathrm{\Omega }|_1ϵ.`$ (8.13)
We have thus shown that if Alice and Charlie have access to a common randomness source of rate $`R`$, while Bob and Charlie can access one of rate $`2Q`$, the conditions for strong subspace transmission can be satisfied. Next, we will illustrate that, by modifying our protocol, it is possible to reduce the amount of shared randomness required. Using the previous blocklength-$`n`$ construction, we will concatenate $`N`$ such codes, where each utilizes the *same* shared randomness, to construct a new code with blocklength $`nN`$. For an arbitrary $`|\mathrm{\Psi }^{(N)}^{B\stackrel{~}{B}^N}`$, further define the commuting operations $`\{𝒯_i\}_{iN},`$ where $`𝒯_i:\stackrel{~}{B}_i\widehat{B}_i`$ is $`𝒯`$ acting on the $`i`$’th tensor factor of $`\mathrm{\Psi }^{(N)}.`$ Setting $`\xi _0\mathrm{\Psi }^{(N)}`$, we then recursively define the density operators $`\xi _i=𝒯_i(\xi _{i1}),`$ noting that $`\xi _N=𝒯_N\mathrm{}𝒯_1(\xi _0)=𝒯^N(\mathrm{\Psi }^{(N)})`$. Because of (8.13), $`|\xi _{i+1}\xi _i|_1=|𝒯_{i+1}(\xi _i)\xi _i|_1ϵ`$, and we can use the triangle inequality to estimate
$`\left|𝒯^N(\mathrm{\Psi }^{(N)})\mathrm{\Psi }^{(N)}\right|_1`$ $`=`$ $`\left|\xi _N\xi _0\right|_1`$
$``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\left|\xi _i\xi _{i1}\right|_1`$
$``$ $`Nϵ.`$
By choosing $`N=\frac{1}{\sqrt{ϵ}}`$, it is clear that we have reduced Alice’s and Bob’s shared randomness rates respectively to $`\sqrt{ϵ}R`$ and $`2\sqrt{ϵ}Q`$, while the error on the $`N`$-blocked protocol is now $`\sqrt{ϵ}`$. Next, we argue that by using two more blocks of length $`n`$, it is possible to simulate the shared randomness by having Alice send $`nR`$ random bits $`X`$ using the first block, while Bob locally prepares two copies of $`\mathrm{\Phi }`$, $`\mathrm{\Phi }^{B_1\stackrel{~}{B}_1}\mathrm{\Phi }^{B_2\stackrel{~}{B}_2}`$, and transmits the $`\stackrel{~}{B}_1\stackrel{~}{B}_2`$ parts over the channel using both blocks. Charlie decodes each block separately, obtaining a random variable $`\widehat{X}`$ and the $`\widehat{B}_1`$ and $`\widehat{B}_2`$ parts of the post-decoded states $`\mathrm{\Omega }_1^{B_1\widehat{B_1}}`$ and $`\mathrm{\Omega }_2^{B_2\widehat{B_2}}.`$ Bob and Charlie then measure their respective parts of $`\mathrm{\Omega }_1\mathrm{\Omega }_2`$ in some previously agreed upon orthogonal bases to obtain a simulation $`\widehat{\mathrm{{\rm Y}}}`$ of the perfect shared randomness state which, by monotonicity and telescoping, satisfies
$`|\mathrm{{\rm Y}}\widehat{\mathrm{{\rm Y}}}|_1`$ $``$ $`|\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Omega }_1\mathrm{\Omega }_2|_1`$
$``$ $`ϵ^2.`$
Further, the noisy shared randomness for the classical messages can be shown to satisfy
$`\left|\mathrm{dist}(X,X)\mathrm{dist}(X,\widehat{X})\right|_1`$ $`=`$ $`2\mathrm{Pr}\{X=\widehat{X}\}`$
$``$ $`ϵ^2.`$
By monotonicity of trace distance and the triangle inequality, using the noisy common randomness state $`\widehat{\mathrm{{\rm Y}}}`$ increases the estimate for each block by $`2ϵ^2`$. For identical reasons, the same increase is incurred by using the noisy common randomness $`(X,\widehat{X})`$. Thus, accounting for both sources of noisy common randomness, the estimate (8.13) is changed to $`2ϵ`$, provided that $`ϵ\frac{1}{4}`$. The noisy common randomness thus increases the bound on the error of the $`N`$-blocked protocol to $`2\sqrt{ϵ}`$, while costing each of Alice and Bob a negligible rate overhead of $`\frac{2}{N+2}`$ in order to seed the protocol.
The above protocol can be considered as defining an encoding map $`^{}:\stackrel{~}{B}^NB^{(N+2)n}`$ and decoding instrument $`𝓓:C^{(N+2)n}\widehat{B}^N\widehat{M}^N`$. Thus, the protocol takes an $`(R,Q,n,ϵ_n)`$ cq entanglement transmission code and constructs an $`(R^{},Q^{},n^{},ϵ_n^{}^{})`$ strong subspace transmission code with cq rate pair $`(R^{},Q^{})=(\frac{R}{1+ϵ_n^{}^{}},\frac{Q}{1+ϵ_n^{}^{}}),`$ where $`n^{}=\left(2+\frac{1}{\sqrt{ϵ_n}}\right)n`$, and $`ϵ_n^{}^{}=2\sqrt{ϵ_n}`$. Now, if the rates $`(R,Q)`$ are achievable cq rates for entanglement transmission, there must exist a sequence of $`(R,Q,n,2ϵ_n^2)`$ entanglement transmission codes with $`ϵ_n0`$. Since this means that $`\frac{1}{1+2\sqrt{ϵ_n}}`$ increases to unity, we have shown that for any $`\delta >0`$, every rate pair $`(R\delta ,Q\delta )`$ is an achievable cq rate pair for strong subspace transmission. Since the capacity regions for each scenario are defined as the closure of the achievable rates, this completes the proof. ∎
#### 8.4.2 $`𝒬_{\text{et}}𝒬`$
###### Proof.
We will employ similar techniques as were used in the previous proof to obtain this implication. Suppose there exists a $`(Q_a,Q_b,n,\frac{1}{2}ϵ^2)`$ qq entanglement transmission code $`(_1,_2,𝒟)`$, with $`_1:\stackrel{~}{A}A^n`$, $`_2:\stackrel{~}{B}B^n`$, and $`𝒟:C^n\widehat{A}\widehat{B}.`$ Setting $`a=|\stackrel{~}{A}|=2^{nQ_a}`$ and $`b=|\stackrel{~}{B}|=2^{nQ_b}`$, define the common randomness states
$$\mathrm{{\rm Y}}_X^{X_AX_C}=\frac{1}{a^2}\underset{x=1}{\overset{a^2}{}}|xx|^{X_A}|xx|^{X_C}$$
and
$$\mathrm{{\rm Y}}_Y^{Y_BY_C}=\frac{1}{b^2}\underset{x=1}{\overset{b^2}{}}|yy|^{Y_B}|yy|^{Y_C}$$
These states will be used as partial inputs to the controlled unitaries
$`𝒰_A`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{a^2}{}}}|xx|^{X_A}U_x,`$
$`𝒰_C`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{a^2}{}}}|xx|^{X_C}U_x^1,`$
$`𝒱_B`$ $`=`$ $`{\displaystyle \underset{y=1}{\overset{b^2}{}}}|yy|^{Y_B}V_x,`$
$`𝒱_C`$ $`=`$ $`{\displaystyle \underset{y=1}{\overset{b^2}{}}}|yy|^{Y_C}V_x^1`$
where, as before, we have utilized the Weyl unitaries $`\{U_x\}_{xa^2}`$ and $`\{V_y\}_{yb^2}`$, which respectively completely randomize any states on $`a`$-dimensional and $`b`$-dimensional spaces. Suppose Alice and Bob are respectively presented with the $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ parts of the arbitrary pure states $`|\mathrm{\Psi }_1^{A\stackrel{~}{A}}`$ and $`|\mathrm{\Psi }_2^{B\stackrel{~}{B}}.`$ Writing $`=𝒟𝒩^n(_1_2)`$, and defining the map $`𝒯:\stackrel{~}{A}\stackrel{~}{B}\widehat{A}\widehat{B}`$ by
$$𝒯:\tau (𝒰_C𝒱_C)(𝒰_A𝒱_B)(\tau \mathrm{{\rm Y}}_1\mathrm{{\rm Y}}_2),$$
the overall joint state of the randomized protocol is given by $`𝒯(\mathrm{\Psi }_1\mathrm{\Psi }_2)`$. Abbreviating
$$|xyxy|^{XY}=|xx|^{X_A}|xx|^{X_C}|yy|^{Y_B}|yy|^{Y_C}$$
and defining $`|\mathrm{\Psi }_x^{A\stackrel{~}{A}}=(1^AU_x)|\mathrm{\Psi }_1`$, $`|\mathrm{\Psi }_y^{B\stackrel{~}{B}}=(1^BV_y)|\mathrm{\Psi }_2,`$ we write
$$\sigma ^{XYAB\stackrel{~}{A}\stackrel{~}{B}}=\frac{1}{a^2b^2}\underset{xy}{}|xyxy|\mathrm{\Psi }_x\mathrm{\Psi }_y.$$
By similar arguments as in the cq case, there exists a map $`𝒪:ABABQR`$ so that
$$(𝒪1^{\stackrel{~}{A}\stackrel{~}{B}})(\mathrm{\Phi }_1\mathrm{\Phi }_2)=\sigma .$$
Again, for the same reasons as in the cq case, we have
$`|(𝒯\text{id})(\mathrm{\Psi }_1\mathrm{\Psi }_2)|_1`$ $``$ $`|(\text{id})(\sigma )|_1`$
$``$ $`|(\text{id})(\mathrm{\Phi }_1\mathrm{\Phi }_2)|_1`$
$``$ $`ϵ.`$
The rest of the proof is nearly identical to that from the previous section, so we omit these details, so as not to have to repeat our previous arguments here. ∎
## Chapter 9 Single-letter examples
Due to the regularized form of our Theorems 1 and 2, the possibility of actually computing the capacity regions seems generally out of reach. Here we give some examples of channels whose capacity region does in fact admit a single-letter characterization, in the sense that no regularization is necessary. In the first section below, we show that a certain erasure quantum erasure multiple access channel has an additive cq capacity region. The next two sections describe classes of channels which have additive single-user capacities. The contents of these two sections are essentially an elaboration of results which appear elsewhere in . The last section demonstrates that the qq capacity region of a certain collective phase-flip channel has an additive capacity region.
### 9.1 Proof of additivity of $`𝒞𝒬`$ for quantum erasure multiple access channel
Our first example is a multiple access erasure channel $`𝒩:A^{}B^{}C`$, where $`|A^{}|=2,|B^{}|=d`$ and $`|C|=d+1.`$ Alice will send classical information while Bob will send quantum. Fixing bases $`\{|0^A^{},|1^A^{}\},\{|1^B^{},\mathrm{}|d^B^{}\},\{|0^C,\mathrm{},|d^C\},`$ the channel has $`d+1`$ operation elements
$`N_0`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{d}{}}}|0^C0|^A^{}j|^B^{}`$
$`N_i`$ $`=`$ $`|i^C1|^A^{}i|^B^{},i=1,\mathrm{}d.`$
The action of the channel can be interpreted as follows. First, a projective measurement of Alice’s input along $`\{|0,|1\}`$ is performed. If the result is $`0`$, Charlie’s output is prepared in a pure state $`|0`$. Otherwise, Bob’s input is transferred perfectly to the remaining degrees of freedom in Charlie’s output. Bob’s input is “erased”, or otherwise ejected into the environment, whenever Alice sends $`|0`$, and is perfectly preserved when she sends $`|1`$. Indeed, the action of $`𝒩`$ on $`\tau ^A^{}\rho ^B^{}`$ is given by
$$𝒩(\tau \rho )=\tau _{00}|00|+\tau _{11}\rho .$$
We will show that the cq capacity region of this channel, $`𝒞𝒬(𝒩_{\text{erasure}})`$, has a single-letter characterization given by the collection of pairs of nonnegative classical-quantum rates $`(R,Q)`$ such that
$`R`$ $``$ $`H(p)`$
$`Q`$ $``$ $`(12p)\mathrm{log}d`$
for some $`0p\frac{1}{2}`$, constituting a generalization of results in on single-user erasure channels to a multiuser setting. Figure 5.1 contains a plot of this region for the case where $`d=2`$.
In the sense of (5.1), any state $`\mathrm{\Omega }^{XBC^k}`$ which arises from $`𝒩^k`$ can be specified by fixing some pure state ensemble $`\{p(x),|\varphi _x^{A^k}\}`$ and a pure bipartite state $`|\mathrm{\Psi }^{BB^k}`$. We thus write
$$\mathrm{\Omega }=\underset{x}{}p(x)|xx|^X(1^B𝒩^k)(\varphi _x\mathrm{\Psi }).$$
For a binary string $`y^k`$, let $`|y^k^{A^k}=|y_1^A^{}\mathrm{}|y_k^A^{}`$ be the associated computational basis state. Writing $`p(y^k|x)=|y^k|\varphi _x|^2`$ defines the random variable $`Y^k`$, which is correlated with $`X`$, and can be interpreted as the erasure pattern associated with the state $`\mathrm{\Omega }`$. We next define another state of the form (5.1),
$$\mathrm{\Omega }^{XY^kBC^k}=\underset{x,y^k}{}p(x)p(y^k|x)|xx|^X|y^ky^k|^{Y^k}𝒩^k(|y^ky^k|\mathrm{\Phi }),$$
for
$$|\mathrm{\Phi }^{BB^k}=\underset{j^k}{}|j^k^B|j_1^{B_1^{}}\mathrm{}|j_k^{B_n^{}},$$
where the summation is over $`d`$-ary strings of length $`k`$, $`j^k=(j_1,\mathrm{},j_k).`$ Finally, for
$`q_i`$ $`=`$ $`\mathrm{Pr}\{Y_i=0\},`$
$`q`$ $`=`$ $`{\displaystyle \frac{1}{k}}{\displaystyle \underset{i=1}{\overset{k}{}}}q_i,`$
$`|\phi ^{BC}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{d}}}{\displaystyle \underset{j=1}{\overset{d}{}}}|j^B|j^C,`$
define a third state
$$\omega ^{UBC}=q|00|^U\pi _d^B|00|^C+(1q)|11|^U\phi ^{BC}.$$
The above states can easily be seen to satisfy the following chain of inequalities
$`I(X;C^k)_\mathrm{\Omega }`$ $`=`$ $`I(X;C^k)_\mathrm{\Omega }^{}`$
$`=`$ $`I(X;Y^k)_\mathrm{\Omega }^{}`$
$``$ $`H(Y^k)_\mathrm{\Omega }^{}`$
$``$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}H(Y_i)_\mathrm{\Omega }^{}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}H(q_i)`$
$``$ $`kH(q)`$
$`=`$ $`kH(U)_\omega `$
$`=`$ $`kI(U;C)_\omega .`$
The only nontrivial step above is that we have used the concavity of the binary entropy function in the last inequality. Furthermore, it is not hard to see that
$`I_c(BC^kX)_\mathrm{\Omega }`$ $``$ $`I_c(BC^kXY^k)_\mathrm{\Omega }^{}`$
$`=`$ $`I_c(BC^kY^k)_\mathrm{\Omega }^{}`$
$`=`$ $`kI_c(BCU)_\omega .`$
Thus, we have shown that for any state $`\mathrm{\Omega }^{XBC^k}`$ arising from $`𝒩^k`$ in the sense of (5.1), there is a state $`\omega ^{UBC}`$ arising from $`𝒩`$ in the same sense, allowing the multi-letter information quantities to be bounded by single-letter information quantities; i.e. $`𝒞𝒬(𝒩)=𝒞𝒬^{(1)}(𝒩)`$. ∎
As it is clear that $`I(U;C)_\omega =H(q)`$, we focus on calculating
$`I_c(BCU)_\omega `$ $`=`$ $`q\left(H(|00|^C)H(\pi _d^B|00|^C)\right)+(1q)\left(H(\pi _d^C)H(\phi ^{BC})\right)`$
$`=`$ $`q(0\mathrm{log}d)+(1q)(\mathrm{log}d0)`$
$`=`$ $`(12q)\mathrm{log}d.`$
Note that the above quantity is a weighted average of a positive and a negative coherent information. It is perhaps tempting to interpret these terms as follows. The positive term can be considered as resulting from a preservation of quantum information, while the negative term can be seen as signifying a complete loss of quantum information to the environment. The overall coherent information is positive only when $`q<\frac{1}{2}`$, a result which is in agreement with the result of Bennett et al. on the quantum capacity of a binary erasure channel. Varying $`0q\frac{1}{2},`$ the rate pairs
$`(R,Q)`$ $`=`$ $`(I(U;C),I_c(BCU))_\omega `$
$`=`$ $`(H(q),(12q)\mathrm{log}d)`$
can be seen to parameterize the outer boundary of $`𝒞𝒬(𝒩)`$, as is pictured in figure 5.1 for the case $`d=2.`$
As an aside, we remark that this calculation, together with the quantum channel capacity theorem from , gives a direct derivation of the quantum capacity of a quantum erasure channel, without relying on the no-cloning and hashing arguments used in .
### 9.2 Degradable channels
While for the single-user capacity $`Q(𝒩)`$ of an arbitrary quantum channel $`𝒩:A^{}B`$ is known not to be additive in general, there is a certain class of channels for which additivity follows relatively easily. This is the class of so-called *degradable channels* . A channel $`𝒩`$ is degradable if its complement $`𝒩^c:A^{}E`$ is a stochastically degraded version of $`𝒩`$, i.e. if there exists a *degrading channel* $`𝒩^d:BE`$ such that
$$𝒩^c=𝒩^d𝒩.$$
Below, we will give a version of the proof from of the additivity of the quantum capacity of an arbitrary degradable channel. Then, we argue that the maximum sum rate bound of the qq capacity region is additive for such channels.
Assume that $`𝒩_1:A_1^{}C_1`$ and $`𝒩_2:A_2^{}C_2`$ are degradable, with isometric extensions $`𝒰_i:A_i^{}C_iE_i.`$ Fix an input state $`|\mathrm{\Psi }^{AA_1^{}A_2^{}}`$ which gives rise to the global state $`|\mathrm{\Omega }^{AC^2E^2}=𝒰_1𝒰_2(\mathrm{\Psi }^{AA^2})`$, where the $`𝒰_i`$ are isometric extensions of the $`𝒩_i`$. By degradability, there exist $`𝒩_i^d`$’s so that $`𝒩_i^c=𝒩_i^d𝒩_i`$, where $`𝒩_i^c=\mathrm{Tr}_{C_i}𝒰_i`$. Letting $`𝒱_i:C_iE_iF_i`$ isometrically extend each $`𝒩_i^d`$, define $`\mathrm{\Theta }^{E^2F^2}=𝒱_1𝒱_2(\mathrm{Tr}_{AE^2}\mathrm{\Omega })`$. Then
$`I_c(AC^2)_\mathrm{\Omega }`$ $`=`$ $`H(C^2)_\mathrm{\Omega }H(E^2)_\mathrm{\Omega }`$
$`=`$ $`H(F^2E^2)_\mathrm{\Theta }H(E^2)_\mathrm{\Theta }`$
$`=`$ $`H(F^2|E^2)_\mathrm{\Theta }`$
$``$ $`H(F_1|E_1)_\mathrm{\Theta }+H(F_2|E_2)_\mathrm{\Theta }`$
$`=`$ $`H(F_1E_1)_\mathrm{\Theta }H(E_1)_\mathrm{\Theta }+H(F_2E_2)_\mathrm{\Theta }H(E_2)_\mathrm{\Theta }`$
$`=`$ $`H(C_1)_\mathrm{\Omega }H(E_1)_\mathrm{\Omega }+H(C_2)_\mathrm{\Omega }H(E_2)_\mathrm{\Omega }`$
$`=`$ $`H(C_1)_\mathrm{\Omega }H(AC^2E_2)_\mathrm{\Omega }+H(C_2)_\mathrm{\Omega }H(AC^2E_1)_\mathrm{\Omega }`$
$`=`$ $`I_c(AC_2E_2C_1)_\mathrm{\Omega }+I_c(AC_1E_1C_2)_\mathrm{\Omega }`$
$`=`$ $`I_c(A_1C_1)_{\omega _1}+I_c(A_2C_2)_{\omega _2}`$
where the inequality is by Lemma 8. In the last line, we set $`\omega _i^{A_iC_i}=𝒩_i(\mathrm{\Psi })`$, identifying $`A_1AA_2^{}`$ and $`A_2AA_1^{}`$. All other steps are either by the fact that isometries preserve entropy or by other trivial rewritings.
Now, if we are given $`k`$ identical channels $`𝒩:A_i^{}C_i`$ and we fix an input state $`|\mathrm{\Psi }^{AA^k}`$ giving rise to $`|\mathrm{\Omega }^{AC^nE^n}=𝒰^k(\mathrm{\Psi }^{AA^k})`$, recursive application of the above yields
$$I_c(AC^k)_\mathrm{\Omega }\underset{i}{}I_c(A_1C_i)_{\omega _i}$$
where $`A_i=AA_1^{}\mathrm{}A_{i1}^{}A_{i+1}^{}\mathrm{}A_k^{}`$, $`\omega _i^{A_iC_i}=𝒩_i(\mathrm{\Psi })`$, and $`𝒩_i`$ is $`𝒩`$ acting on the $`i`$th tensor factor. Choosing
$$i^{}=\mathrm{arg}\underset{i}{\mathrm{max}}\{I_c(A_iC_i)_{\omega _i}\}$$
yields
$$\frac{1}{k}I_c(AC^k)_\mathrm{\Omega }I_c(A_i^{}C_i^{})_{\omega _i^{}}\underset{\omega ^{AC}}{\mathrm{max}}I_c(AC)_\omega =Q^{(1)}(𝒩),$$
where the maximization is as over all $`\omega =𝒩(\varphi ^{AA^{}})`$.
Let us phrase this conclusion using different notation. Let $`\tau ^{A^kB^k}`$ be arbitrary, and define $`\tau _i^{A_iB_i}=\mathrm{Tr}_{/A_iB_i}\tau ,`$ where $`\mathrm{Tr}_{/A_iB_i}`$ denotes the partial trace over all systems which are *not* $`A_iB_i`$. Then
$$I_c(\tau ,𝒩^k)kI_c(\tau _i^{},𝒩),$$
where
$$i^{}=\mathrm{arg}\underset{i}{\mathrm{max}}I_c(\tau _i,𝒩).$$
Now, if $`\rho ^{A^k}`$ and $`\sigma ^{B^k}`$ are arbitrary, and we define $`\rho _i=\mathrm{Tr}_{/A_i}\rho `$ and $`\sigma _i=\mathrm{Tr}_{/B_i}\sigma `$, observe that if $`\tau =\rho \sigma `$, then $`\tau _i=\rho _i\sigma _i`$. This immediately implies that
$$I_c(\rho \sigma ,𝒩^k)kI_c(\rho _i^{}\sigma _i^{},𝒩),$$
where
$$i^{}=\mathrm{arg}\underset{i}{\mathrm{max}}I_c(\rho _i\sigma _i,𝒩),$$
proving that the maximum sum rate of any degradable channel is additive, even when the inputs are restricted to be product states. This fact will be useful in Section 9.4, where we give a channel whose qq capacity region is single-letter.
### 9.3 Generalized dephasing channels
In this section we describe a certain subclass of the class of degradable channels. These are channels $`𝒩:A^{}B`$ with $`|A|=|B|=d`$ for which there is a particular orthogonal basis $`\{|x^A^{}\}`$ which can be transmitted through the channel without error
$$𝒩\left(|xx|\right)=|xx|$$
although superpositions of these basis vectors are potentially subject to noise. Here, $`\{|x^B\}`$ is a corresponding orthogonal basis for $`B`$. Such a channel has an isometric extension $`𝒰:A^{}BE`$ given by
$$𝒰=\underset{x}{}|x^B|\varphi _x^Ex|^A^{},$$
where the states $`|\varphi _x^E`$ are not necessarily orthogonal. To see that these channels are degradable, observe that for any input state $`\rho ^A^{}`$,
$`𝒩^c(\rho )`$ $`=`$ $`\mathrm{Tr}_B𝒰(\rho )`$
$`=`$ $`{\displaystyle \underset{x}{}}x|^B\left({\displaystyle \underset{x^{\prime \prime }x^{}}{}}|x^{\prime \prime }^B|\varphi _{x^{\prime \prime }}^Ex^{\prime \prime }|^A^{}\rho |x^{}^A^{}x^{}|^B\varphi _x^{}|^E\right)|x^B`$
$`=`$ $`{\displaystyle \underset{x}{}}x|\rho |x\varphi _x^E.`$
Note that $`𝒩^c(\rho )`$, depends only on the diagonal matrix elements of $`\rho `$ (when it is expressed in the dephasing basis. However, these are exactly the matrix elements which are unaffected by the action of $`𝒩`$, making degradability evident. In fact, the degrading channel is precisely $`𝒩^c`$, i.e.
$$𝒩^c=𝒩^c𝒩.$$
It is interesting to relate the isometric extension $`𝒰_𝒩`$ to the operator sum representation for $`𝒩`$. To do this, first express $`𝒰_𝒩`$ in the flattened representation
$$𝒰_𝒩=\left(\begin{array}{cccc}|\varphi _1& & & \\ & |\varphi _2& & \\ & & \mathrm{}& \\ & & & |\varphi _d\end{array}\right).$$
Supposing that $`|E|=k`$, note that the matrix is “block diagonal”, with $`d`$ $`k\times 1`$ blocks, where this is expressed as a map to the system $`EB`$. Regrouping the rows into to $`k`$ groups of size $`d`$ we rewrite
$$𝒰_𝒩=\left(\begin{array}{cccc}1|\varphi _1& & & \\ & 1|\varphi _2& & \\ & & \mathrm{}& \\ & & & 1|\varphi _d\\ 2|\varphi _1& & & \\ & 2|\varphi _2& & \\ & & \mathrm{}& \\ & & & 2|\varphi _d\\ & & \mathrm{}& \\ k|\varphi _1& & & \\ & k|\varphi _2& & \\ & & \mathrm{}& \\ & & & k|\varphi _d\end{array}\right)=\left(\begin{array}{c}N_1\\ N_2\\ \mathrm{}\\ N_k\end{array}\right).$$
This is just the flattened representation for the map to the system $`BE`$ (the order of $`E`$ and $`B`$ have been reversed). Note that we have identified the $`|E|`$ blocks with the matrices of the operator sum representation
$$𝒩(\rho )=\underset{e=1}{\overset{k}{}}N_e\rho N_e^{}.$$
So we see that the operator sum matrices are all diagonal in the $`\{|i\}`$ basis and are given explicitly as
$$N_e=\underset{x}{}e||\varphi _x|x^Bx|^A^{}.$$
Reversing the above steps, it is clear that $`𝒩`$ is a generalized dephasing channel if and only if it has an operator sum representation consisting of matrices which commute.
Let us mention that in the special case where the $`\{\varphi _x\}`$ are mutually orthogonal, the channel is *completely dephasing*. We denote this channel as $`\mathrm{\Delta }`$, and note that it corresponds to a channel which performs a pure state measurement in the dephasing basis while ignoring the result. This has the effect of setting all of the off-diagonal matrix elements of $`\rho `$ equal to zero. $`\mathrm{\Delta }`$ obeys the following equations:
$`𝒩^c`$ $`=`$ $`𝒩^c\mathrm{\Delta }`$
$`H(\mathrm{\Delta }(\rho ))`$ $``$ $`H(\rho ).`$
The first is because $`𝒩^c`$ only depends on the diagonal components of $`\rho `$, while the second is proved in . Observe that the inequality is saturated for diagonal $`\rho `$. Because of this, we may write
$`Q(𝒩)`$ $`=`$ $`\underset{\rho }{\mathrm{max}}I_c(\rho ,𝒩)`$
$`=`$ $`\underset{\rho }{\mathrm{max}}\left\{H\left(𝒩(\rho )\right)H\left(𝒩^c(\rho )\right)\right\}`$
$`=`$ $`\underset{\rho }{\mathrm{max}}\left\{H\left(𝒩\mathrm{\Delta }(\rho )\right)H\left(𝒩^c\mathrm{\Delta }(\rho )\right)\right\}`$
$`=`$ $`\underset{p(x)}{\mathrm{max}}\left\{H(X)H\left({\displaystyle \underset{x}{}}p(x)\varphi _x\right)\right\}.`$
### 9.4 Proof of additivity of $`𝒬`$ for collective phase-flip channel
While the description of the capacity region $`𝒬`$ in Theorem 2 generally requires taking a many-letter limit, we give here an example of a quantum multiple access channel $`𝒩_p:A^{}B^{}C`$ for which that description can be single-letterized. The channel $`𝒩_p`$ takes as input two qubits, one from Alice and the other from Bob. With probability $`p`$, the channel causes each qubit to undergo a phase flip, by rotating each by 180 about its z-axis before it is received by the receiver Charlie. The action of $`𝒩_p`$ on an input density operator $`\rho ^{A^{}B^{}}`$ is described in terms of the operator sum representation as
$$𝒩_p(\rho )=(1p)\rho +p(\sigma _z\sigma _z)\rho (\sigma _z\sigma _z),$$
where
$$\sigma _z=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
is the Pauli phase flip matrix. We will demonstrate that $`𝒬(𝒩_p)`$ is equal to the collection of all pairs of nonnegative rates $`(Q_a,Q_b)`$ which satisfy
$`Q_a`$ $``$ $`1`$
$`Q_b`$ $``$ $`1`$
$`Q_a+Q_b`$ $``$ $`2H(p).`$
###### Proof.
In order to prove this, we first recall that the maximum of the sum rate bound $`I_c(ABC)`$ over all inputs of the form (5.2) is additive. Next, we calculate $`Q(𝒩_p)`$, the single-user capacity of the channel, and observe that it is achieved for inputs of the form (5.2), implying that the maximum sum rate bound equals the capacity. Then, we show that for the same inputs, the bounds $`I_c(ABC)`$ and $`I_c(BAC)`$ on the individual rates are as large as is possible. The characterization in terms of a single pentagon will then follow.
We first note that the the operator sum matrices $`\sqrt{p}\sigma _x\sigma _x`$ and $`\sqrt{1p}1_4`$ commute. By results in the previous section, we conclude that $`𝒩_p`$ is an example of a *generalized dephasing channel* and thus, the following two conditions are satisfied:
* for any state $`\mathrm{\Omega }^{ABC^k}=𝒩^k(\mathrm{\Psi }^{ABA^kB^k})`$ arising from $`𝒩^k`$ (where Alice and Bob can jointly prepare any state at the inputs), there is a state $`\omega ^{ABC}=𝒩(\psi ^{ABA^{}B^{}})`$ for which
$$I_c(ABC)_\omega \frac{1}{k}I_c(ABC^k)_\mathrm{\Omega }.$$
Furthermore, the input density operator $`\rho ^{A^{}B^{}}=\mathrm{Tr}_{AB}\psi ^{ABA^{}B^{}}`$ is diagonal in the dephasing basis of $`𝒩_p`$.
* for any state $`\mathrm{\Omega }^{ABC^k}=𝒩^k(\mathrm{\Psi }_1^{AA^k}\mathrm{\Psi }_2^{BB^k})`$ arising from $`𝒩^k`$ in the sense of (5.2), there is a state $`\omega ^{ABC}=𝒩(\varphi _1^{AA^{}}\varphi _2^{BB^{}})`$ arising from $`𝒩`$ in the same sense for which
$$I_c(ABC)_\omega ^{}\frac{1}{k}I_c(ABC^k)_\mathrm{\Omega }^{}.$$
The first condition above says that the single-user capacity $`Q(𝒩_p)`$ is additive. It also guarantees that the relevant maximization is achieved by an input density operator $`\rho ^{A^{}B^{}}`$ which is diagonal in the dephasing basis. The second condition guarantees that the constrained single-user capacity of $`𝒩_p`$, when the users are constrained to preparing product input states, is additive.
In order to compute $`Q(𝒩_p)`$, let us first write an isometric extension $`𝒰:ABCE`$ of $`𝒩_p`$ as
$$𝒰|i^A|i^B=|ij^C|\varphi _{ij}^E,$$
where
$$|\varphi _{00}^E=|\varphi _{11}^E=\sqrt{1p}|0^E+\sqrt{p}|1^E|\varphi _+^E$$
and
$$|\varphi _{01}^E=|\varphi _{10}^E=\sqrt{1p}|0^E\sqrt{p}|1^E|\varphi _{}^E.$$
A complementary channel $`𝒩_p^c`$ is then defined as
$`𝒩_p^c(\rho )`$ $`=`$ $`\mathrm{Tr}_C𝒰(\rho )`$
$`=`$ $`{\displaystyle \underset{ij}{}}|\varphi _{ij}i|j|\rho |i|j\varphi _{ij}|`$
$`=`$ $`{\displaystyle \underset{ij}{}}\rho _{ij}\varphi _{ij}`$
$`=`$ $`(\rho _{00}+\rho _{11})\varphi _++(\rho _{01}+\rho _{10})\varphi _{}.`$
Observe that the output of the $`𝒩_p^c`$ depends only on the diagonal elements of $`\rho `$, when $`\rho `$ is written in the dephasing basis $`\{|00,|01,|10,|11\}.`$ Define $`\alpha =\rho _{00}+\rho _{11}`$. As $`Q(𝒩_p)`$ is achieved when $`\rho `$ is diagonal in this basis, let us calculate
$`H(C)`$ $`=`$ $`H(A^{}B^{})`$
$`=`$ $`H\left(\{\rho _{00},\rho _{01},\rho _{10},\rho _{11}\}\right)`$
$`=`$ $`H(\alpha )+\alpha H\left({\displaystyle \frac{\rho _{00}}{\alpha }}\right)+(1\alpha )H\left({\displaystyle \frac{\rho _{01}}{1\alpha }}\right)`$
$``$ $`H(\alpha )+1,`$
where the inequality is saturated when $`\rho _{00}=\rho _{11}=\frac{\alpha }{2}`$ and $`\rho _{01}=\rho _{10}=\frac{1\alpha }{2}`$. It thus suffices to optimize over the class of states
$$\rho ^{A^{}B^{}}=\left(\begin{array}{cccc}\frac{\alpha }{2}& 0& 0& 0\\ 0& \frac{1\alpha }{2}& 0& 0\\ 0& 0& \frac{1\alpha }{2}& 0\\ 0& 0& 0& \frac{\alpha }{2}\end{array}\right)$$
for which
$$H(C)=H(\rho )=1+H(\alpha ),$$
Note that we may express
$`\varphi _\pm ^E={\displaystyle \frac{1}{2}}\left(1\pm \sqrt{p(1p)}\sigma _x(12p)\sigma _z\right),`$
allowing us to write
$`𝒩_p^c(\rho )=\alpha \varphi _++(1\alpha )\varphi _{}={\displaystyle \frac{1}{2}}\left(1+(2\alpha 1)\sqrt{p(1p)}\sigma _x(12p)\sigma _z\right),`$
so that $`H(E)=H\left(\frac{1}{2}(1+\sqrt{p(1p)(2\alpha 1)^2+(12p)^2})\right).`$ Thus,
$`I_c(\rho ,𝒩)`$ $`=`$ $`H\left(𝒩_p(\rho )\right)H\left(𝒩_p^c(\rho )\right)`$
$`=`$ $`1+H(\alpha )H\left({\displaystyle \frac{1}{2}}(1+\sqrt{p(1p)(2\alpha 1)^2+(12p)^2})\right)`$
$``$ $`h(\alpha ).`$
For fixed $`p`$, $`h(\alpha )`$ is symmetric about $`\alpha =\frac{1}{2}`$, and has a first derivative which is positive for $`0\alpha <\frac{1}{2}`$ (and is thus negative for $`\frac{1}{2}<\alpha 1`$). Because $`h(\alpha )`$ is continuous on $`0\alpha 1`$, its maximum is attained when $`\alpha =\frac{1}{2}`$, so that
$$\underset{\rho }{\mathrm{max}}I_c(\rho ,𝒩)=I_c(\pi ^{A^{}B^{}},𝒩)=1+H\left(\frac{1}{2}\right)H(p)=2H(p).$$
So we see that the maximum is already achieved for a product state $`\pi ^{A^{}B^{}}=\pi ^A^{}\pi ^B^{}`$. Define the Bell states
$`|\psi _\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|00\pm |11\right).`$
As $`|\psi _+`$ purifies the maximally mixed state $`\pi _2`$, let us define the global state
$$\omega ^{ABC}=𝒩(\psi _+^{AA^{}}\psi _+^{BB^{}}).$$
Identifying $`C=\widehat{A}\widehat{B}`$ in the obvious way, let us reexpress
$$\omega ^{A\widehat{A}B\widehat{B}}=(1p)\psi _+^{A\widehat{A}}\psi _+^{B\widehat{B}}+p\psi _{}^{A\widehat{A}}\psi _{}^{B\widehat{B}}.$$
It is now a simple task to calculate
$`H(ABC)`$ $`=`$ $`H(\omega )=H(p)`$
$`H(C)`$ $`=`$ $`H(\pi ^C)=2`$
$`H(AC)`$ $`=`$ $`H(A\widehat{A})+H(\widehat{B})=H(p)+1=H(BC).`$
Combining these gives the relevant coherent informations
$`I_c(ABC)`$ $`=`$ $`H(C)H(ABC)=2H(p)`$
$`I_c(ABC)`$ $`=`$ $`H(BC)H(ABC)=1+H(p)H(p)=1`$
$`I_c(BAC)`$ $`=`$ $`H(AC)H(ABC)=1.`$
As we saw in Section 3.4, $`I_c(ABC)\mathrm{log}|A^{}|=1`$ and $`I_c(BAC)\mathrm{log}|B^{}|=1`$ for any state arising from $`𝒩`$. The individual rate bounds are thus saturated and the claim follows. ∎
## Chapter 10 Discussion
There have been a number of results analyzing multiterminal coding problems in quantum Shannon theory. For an i.i.d. classical-quantum source $`XB`$, Devetak and Winter have proved a Slepian-Wolf-like coding theorem achieving the cq rate pair $`(H(X|B),H(B))`$ for classical data compression with quantum side information. Such codes extract classical side information from $`B^n`$ to aid in compressing $`X^n.`$ The extraction of side information is done in such a way as to cause a negligible disturbance to $`B^n`$. Our Theorem 1 is somewhat of this flavor. There, the quantum state of $`C^n`$ is measured to extract Alice’s classical message which, in turn, is used as side information for decoding Bob’s quantum information. Analogous results to ours were obtained by Winter in his analysis of a multiple access channel with classical inputs and a quantum output, whereby the classical decoded message of one sender can be used as side information to increase the classical capacity of another sender.
We further mention the obvious connection between our coding theorems and the subject of channel codes with side information available to the receiver. The more difficult problem of classical and quantum capacities when side information is available at the *encoder* is analyzed by Devetak and Yard in , constituting quantum generalizations of results obtained by Gelfand and Pinsker for classical channels with side information.
In an earlier draft of , we characterized $`𝒬(𝒩)`$ as the closure of a regularized union of rectangles
$`0`$ $`R`$ $`{\displaystyle \frac{1}{k}}I_c(AC^k)`$
$`0`$ $`S`$ $`{\displaystyle \frac{1}{k}}I_c(BC^k).`$
This solution had been conjectured on the basis of a duality between classical Slepian-Wolf distributed source coding and classical multiple-access channels , as well as on a purported no-go theorem for distributed data compression of so-called irreducible pure state ensembles that appeared in an early version of . After the earlier preprint was made available, Andreas Winter announced recent progress with Jonathan Oppenheim and Michal Horodecki on the quantum Slepian-Wolf problem, offering a characterization identical in functional form to the classical one, while also supplying an interpretation of negative rates and apparently evading the no-go theorem. Motivated by the earlier mentioned duality, he informed us that the qq capacity region could also be characterized in direct analogy to the classical case. Subsequently, we found that we could modify our previous coding theorem to achieve the new region, provided that the rates are nonnegative. After those events unfolded, the authors of found an error in the proof of their no-go theorem, leading to a revised version consistent with the newer developments. Our earlier characterization of $`𝒬(𝒩)`$, while correct, is contained in the rate region of Theorem 2 for any finite $`k`$, frequently strictly so. The newer theorem, therefore, gives a more accurate approximation to the rate region for finite $`k`$. In fact, for any state arising from the channel which does not saturate the strong subadditivity inequality , the corresponding pentagon and rectangle regions are distinct. As seen in Section 9.2, another beneficial feature of the new characterization is that for any channel which is *degradable*, the maximum sum rate bound $`R+S\mathrm{max}I_c(ABC)`$ is additive, where the maximization is over all states of the form (5.2). Furthermore, recall that in Section 9.4, the pentagon characterization was single-letterized for the collective phase flip channel. On the other hand, computer calculations have revealed that the rectangle region does *not* lead to a single-letter characterization of that channel. This seems to indicate that the newer characterization is the “correct” one, at least for that particular channel.
More recently, we discovered that the same technique used to prove the new characterization of $`𝒬(𝒩)`$ implies a new cq coding theorem, and thus a new characterization of $`𝒞𝒬(𝒩)`$. By techniques nearly identical to those employed in the coding theorem for Theorem 2, it is possible to achieve the cq rate pair
$$(R,Q)=(I(X;BC),I_c(BC))$$
corresponding to Bob’s quantum information being used as side information for decoding Alice’s classical message. This is accomplished by having Charlie isometrically decode Bob’s quantum information, then coherently decode to produce an effective channel $`𝒩_1:A^{}BC`$ so that Alice can transmit classically at a higher rate. The new characterization is then a regularized union of pentagons, consisting of pairs of nonnegative rates $`(R,Q)`$ satisfying
$`r`$ $``$ $`I(X;BC)`$
$`S`$ $``$ $`I_c(BCX)`$
$`r+S`$ $``$ $`I(X;C)+I_c(BCX)=I(X;BC)+I_c(BC).`$
Surprisingly, it is thus possible to characterize each of $`𝒞𝒬(𝒩)`$ and $`𝒬(𝒩)`$ in terms of pentagons, in analogy to the original classical result. This situation makes apparent the dangers of being satisfied with regularized expressions for capacity regions. Without being able to prove single-letterization steps in the converses, it is hard to differentiate which characterization is the “right” one. While it is intuitively satisfying to see analogous formulae appear in both the classical and quantum theories, the regularized nature of the quantum results blurs the similarity. Indeed, the problems with single-letterization for single-user channels appear to be amplified when analyzing quantum networks (see e.g. ). While $`𝒬`$ is additive for the collective phase flip channel of Section 9.4, this behavior does not appear to be generic for the classes of degradable or generalized dephasing channels, as the saturation of the individual rate bounds for that example seem to be the source of additivity. Perhaps this indicates that the necessity of understanding the capacities of single-user channels at a level beyond regularized optimizations is even more pressing than previously thought. It should be mentioned that for the erasure channel analyzed in Section 9.1, the newer description of $`𝒞𝒬(𝒩)`$ is not an issue, as the new corner point is contained in the old rectangle for any state arising from any number of parallel instances of the erasure channel.
Consider the full simultaneous classical-quantum region $`𝒮(𝒩)`$ defined in Section 5.2. This region can be characterized in a way that generalizes Theorems 1 and 2 as the regularization of the region $`𝒮^{(1)}(𝒩)`$, defined as the vectors of nonnegative rates $`(R_a,R_b,Q_a,Q_b)`$ satisfying
$`R_a`$ $``$ $`I(X;C|Y)`$
$`R_b`$ $``$ $`I(Y;C|X)`$
$`R_a+R_b`$ $``$ $`I(XY;C)`$
$`Q_a`$ $``$ $`I_c(ABCXY)`$
$`Q_b`$ $``$ $`I_c(BACXY)`$
$`Q_a+Q_b`$ $``$ $`I_c(ABCXY)`$
for some state of the form
$$\sigma ^{XYABC}=\underset{x,y}{}p(x)p(y)|xx|^X|yy|^Y𝒩(\psi _x^{AA^{}}\varphi _y^{BB^{}}),$$
arising from the action of $`𝒩`$ on the $`A^{}`$ and $`B^{}`$ parts of some pure state ensembles $`\{p(x),|\psi _x^{AA^{}}\}`$, $`\{p(y),|\varphi _y^{BB^{}}\}`$. Briefly, achievability of this region is obtained as follows. Using techniques introduced in , each sender “shapes” their quantum information into HSW codewords. Decoding is accomplished by first decoding all of the classical information, then using that information as side information for a quantum decoder. A formal proof of the achievability of this region is found in . The main result of , the regularized optimization of the cq result from over pairs of input ensembles, and our Theorems 1 and 2 follow as corollaries of the corresponding capacity theorem. Indeed, the six two-dimensional “shadows” of the above region, obtained by setting pairs of rates equal to zero, reproduce those aforementioned results. This characterization, however, only utilizes the rectangle description of $`𝒞𝒬(𝒩)`$. It is indeed possible to write a more accurate regularized description of $`𝒮(𝒩)`$ which generalizes the pentagon characterizations of $`𝒞𝒬(𝒩)`$ and $`𝒬(𝒩)`$, although we will not pursue that at this time.
## Chapter 11 Appendix
### 11.1 Quantum instruments and coherent information
For some finite set $`𝒮`$, consider a labelled collection of channels $`\{𝒩_s\}_{s𝒮}`$, where $`𝒩_s:A^{}B`$. Define an instrument $`𝓝:A^{}SB`$ to act as
$$𝓝:\tau \underset{s}{}p(s)|ss|^S𝒩_s(\tau ).$$
An instrument channel such as $`𝓝`$ may be interpreted as one with classical state information made available to the receiver. We will show that every channel $`𝒩^c:A^{}E`$ which is complementary to $`𝓝`$ is an instrument as well, as the environment $`E`$ contains a copy of $`S`$. In other words, the classical state information is also available to an eavesdropper with full control of the environment.
An isometric extension $`𝒰`$ of $`𝓝`$ may be constructed as follows. First, fix isometric extensions $`𝒰_s:A^{}E^{}B`$ for the individual $`𝒩_s`$’s. Then, define $`𝒰:A^{}SEB`$ via
$$𝒰=\underset{s}{}\sqrt{p(s)}|s^S|s^{E^{\prime \prime }}𝒰_s,$$
taking $`E=E^{}E^{\prime \prime }`$. That this is indeed an isometry is evident, because $`𝒰^{}𝒰=_sp(s)𝒰_s^{}𝒰_s=_sp(s)1^A^{}=1^A^{}`$. We may further check that $`𝒰`$ is in fact an extension of $`𝓝`$, by calculating
$`\mathrm{Tr}_E𝒰(\tau )`$ $`=`$ $`\mathrm{Tr}_E^{}\mathrm{Tr}_{E^{\prime \prime }}𝒰(\tau )`$
$`=`$ $`\mathrm{Tr}_E^{}{\displaystyle \underset{s}{}}p(s)|ss|^S𝒰_s(\tau )`$
$`=`$ $`{\displaystyle \underset{s}{}}p(s)|ss|^S𝒩_s(\tau )`$
$`=`$ $`𝓝(\tau ).`$
Thus, the action of the complementary channel $`𝒩^c`$ can be defined via $`𝒰`$ as
$`𝒩^c(\tau )`$ $`=`$ $`\mathrm{Tr}_{BS}𝒰(\tau )`$
$`=`$ $`\mathrm{Tr}_B{\displaystyle \underset{s}{}}p(s)|ss|^{E^{\prime \prime }}𝒰_s(\tau )`$
$`=`$ $`{\displaystyle \underset{s}{}}p(s)|ss|^{E^{\prime \prime }}𝒩_s^c(\tau ),`$
where the $`𝒩_s^c=\mathrm{Tr}_B𝒰_s`$ are complementary channels to the $`𝒩_s`$’s.
### 11.2 Proof of convexity of $`𝒞𝒬`$ and $`𝒬`$
Let $`𝒩:A^{}B^{}C`$ be a quantum multiple access channel. We will prove that $`𝒬(𝒩)`$ is convex, as the proof for $`𝒞𝒬`$ is identical. Let $`k_0`$ and $`k_1`$ be positive integers, and fix any two states of the form (5.2), $`\sigma _0^{A_0B_0C^{k_0}}`$ and $`\sigma _1^{A_1B_1C^{k_1}}.`$ Then $`(R_0,S_0),(R_1,S_1)𝒬(𝒩)`$, where for $`i\{0,1\}`$,
$`R_i`$ $`=`$ $`{\displaystyle \frac{1}{k_i}}I_c(A_iC^{k_i})_{\sigma _i}`$
$`S_i`$ $`=`$ $`{\displaystyle \frac{1}{k_i}}I_c(B_iC^{k_i})_{\sigma _i}.`$
We will now show that for any rational $`0\lambda 1`$, $`\lambda (R_0,S_0)+(1\lambda )(R_1,S_1)𝒬(𝒩).`$ We first write $`\lambda =\frac{\alpha }{\beta },`$ for integers satisfying $`\beta >0,`$ $`\beta \alpha 0`$. Setting $`p_0=\alpha k_1,`$ $`p_1=(\beta \alpha )k_0,`$ and $`k=p_0k_0+p_1k_1`$, define the composite systems $`A=A_0^{p_0}A_1^{p_1}`$ and $`B=B_0^{p_0}B_1^{p_1}`$, as well as the density matrix $`\sigma ^{ABC^k}=\sigma _0^{p_0}\sigma _1^{p_1},`$ which is also of the form (5.2). Additivity of coherent information across product states and some simple algebra gives
$`{\displaystyle \frac{1}{k}}I_c(AC^k)_\sigma `$ $`=`$ $`{\displaystyle \frac{p_0}{k}}I_c(A_0C^{k_0})_{\sigma _0}+{\displaystyle \frac{p_1}{k}}I_c(A_1C^{k_1})_{\sigma _1}`$
$`=`$ $`{\displaystyle \frac{p_0k_0R_0+p_1k_1R_1}{p_0k_0+p_1k_1}}`$
$`=`$ $`\lambda R_0+(1\lambda )R_1.`$
An identical calculation shows that $`\frac{1}{k}I_c(BC^k)_\sigma =\lambda S_0+(1\lambda )S_1.`$ As $`𝒬(𝒩)`$ was defined as the topological closure of rate pairs corresponding to states which appropriately arise from the channel, the result follows because the set of previously considered $`\lambda `$’s comprises a dense subset of the unit interval.
### 11.3 Proof of cardinality bound on $`𝒳`$.
Begin by fixing a finite set $`𝒳`$, a labelled collection of pure states $`\{|\varphi _x^A^{}\}_{x𝒳}`$, and a pure bipartite state $`|\mathrm{\Psi }^{BB^{}}.`$ For each $`x`$, these define the states $`\sigma _x^{BC}=𝒩(\varphi _x\mathrm{\Psi })`$ and $`\omega _x^C=\mathrm{Tr}_B\sigma _x`$. Assume for now that $`|A^{}||C|`$. Define a mapping $`f:𝒳^{|C|^2+1}`$, via
$$f:xf_x(\omega _x,H(\omega _x),I_c(BC)_{\sigma _x}),$$
where we are considering $`\omega _x`$ to be synonymous with its $`|C|^21`$ dimensional parameterization. By linearity, this extends to a map from probability mass functions on $`𝒳`$ to $`^{|C|^2+1},`$ where
$$f:p(x)\underset{x}{}p(x)f_x(\omega _p,H(C|X)_p,I_c(BCX)_p),$$
Our use of the subscript $`p`$ should be clear from the context. The use of Caratheodory’s theorem for bounding the support sizes of auxiliary random variables in information theory (see ) is well-known. Perhaps less familiar is the observation that a better bound can often be obtained by use of a related theorem by Fenchel and Eggleston , which states that if $`S^n`$ is the union of at most $`n`$ connected subsets, and if $`y`$ is contained in the convex hull of $`S`$, then $`y`$ is also contained in the convex hull of at most $`n`$ points in $`S`$. As the map $`f`$ is linear, it maps the simplex of distributions on $`𝒳`$ into a single connected subset of $`^{|C|^2+1}`$. Thus, for any distribution $`p(x)`$, there is another distribution $`p^{}(x)`$ which puts positive probability on at most $`|C|^2+1`$ states, while satisfying $`f(p)=f(p^{}).`$ If it is instead the case that $`|A^{}|<|C|,`$ this bound can be reduced to $`|A|^2+1`$ by replacing the first components of the map $`f`$ with a parameterization of $`\varphi _x^A^{}`$, as specification of a density matrix on $`A^{}`$ is enough to completely describe the resulting state on $`C`$. It is therefore sufficient to consider $`|X|\mathrm{min}\{|A^{}|,|C|\}^2+1`$ when computing $`𝒞𝒬^{(1)}(𝒩)`$.
|
warning/0506/hep-ex0506038.html
|
ar5iv
|
text
|
# Latest Jets Results from the Tevatron at √𝒔=1.96𝐓𝐞𝐕
Latest jet results from the Tevatron are presented in this conference note. These are namely: new results on central inclusive jet production using both cone and $`k_T`$ algorithms, measurement of decorrelation in azimuthal angle between the two jets with the highest transverse momenta, and study of jet shapes. Results are based on data collected in $`\mathrm{p}\overline{\mathrm{p}}`$ collisions at $`\sqrt{s}=1.96\mathrm{TeV}`$ in the years 2001-2004. Depending on the analysis, integrated luminosity of the sample was up to $`378\mathrm{pb}^1`$.
Producing events with high transverse momenta ($`p_\mathrm{T}`$) jets we probe properties of matter and space at very short distances. At the Fermilab Tevatron, jets are being produced with $`p_\mathrm{T}`$ up to about $`600\mathrm{GeV}`$ which corresponds to distances about thousand times smaller than the proton size. Up to this scale, the Tevatron high $`p_\mathrm{T}`$ jet data provide tests of our understanding of proton structure and strong force that acts between proton constituents: quarks and gluons. Another aspect of the Tevatron jet physics program is the detailed study of jet properties and understanding these properties within the framework of perturbative Quantum Chromodynamics (pQCD). These studies have direct impact on many non-QCD analyzes that work with jets. In addition, strong interaction induced processes provide unavoidable background for these analyzes in hadron-hadron colliders. Better understanding of QCD processes thus improve our chances for discovery of potential signals of new physics.
To the first category of high $`p_\mathrm{T}`$ jet program belong updated results on central jet inclusive $`p_\mathrm{T}`$ spectra both from the CDF and the DØ Collaborations. Second category is represented by the DØ measurement of dijet azimuthal decorrelations, study of jet shapes by CDF, and the CDF measurement of inclusive jet $`p_\mathrm{T}`$ spectra using $`k_\mathrm{T}`$ algorithm with different sizes of jets.
The analyzes presented in this talk are based on data collected with CDF $`^\mathrm{?}`$ and DØ $`^\mathrm{?}`$ detectors in Run II of the Tevatron. This period started in year 2001 after upgrade of collider and both detectors. The Tevatron delivers now proton antiproton collisions at $`\sqrt{s}=1.96\mathrm{TeV}`$ while in previous run (Run I) it was $`\sqrt{s}=1.8\mathrm{TeV}`$. Also Tevatron luminosity was increased significantly. Both aspects of collider upgrade affected high $`p_\mathrm{T}`$ jet physics. Due to them, both experiments already collected about one order of magnitude more events with high $`p_\mathrm{T}`$ jets than in Run I.
The CDF and DØ Collaborations measured inclusive production of central jets. DØ performed the measurement in two rapidity bins: $`|y|<0.4`$ and $`0.4<|y|<0.8`$ (rapidity $`y`$ is defined as $`y=\frac{1}{2}\mathrm{ln}\frac{E+p_z}{Ep_z}`$, where $`E`$ is jet energy and $`p_z`$ is jet longitudinal momentum along the beam axis). The measurement was based on $`378\mathrm{pb}^1`$ of data. Jets were reconstructed using iterative seed-based cone algorithm (including midpoints) with radius $`R=0.7`$ $`^\mathrm{?}`$. Jet energies where calibrated using $`\gamma +\mathrm{jet}`$ sample. Uncertainty on jet energy was found to be between 4-5%, for jet energies between $`30`$-$`400\mathrm{GeV}`$, and about 6% at $`500\mathrm{G}\mathrm{e}\mathrm{V}`$. It translates to about 25% (30%, 60%) error on jet production cross section for $`p_\mathrm{T}=100\mathrm{GeV}`$ ($`300\mathrm{GeV}`$, $`500\mathrm{GeV}`$). Results of the measurement are shown in Fig. 2. The error on jet cross section was dominated by the uncertainty on jet energy calibration.
CDF performed the measurement for jets with $`0.1<|\eta |<0.7`$ where $`\eta `$ is jet pseudorapidity ($`\eta =\mathrm{ln}\mathrm{tan}\vartheta /2`$, where $`\vartheta `$ is jet polar angle measured from the beam axis). In this case, jets where reconstructed with CDF Run I cone algorithm with radius $`R=0.7`$ $`^\mathrm{?}`$. Measurement was based on $`177\mathrm{pb}^1`$ of data. Final result is shown in Fig. 2. Systematic error is again dominated by the uncertainty of jet energy calibration which was about 3%. CDF calibrated their jet energies using calorimeter electron and hadronic responses measured during testbeam.
CDF and DØ results were compared with next-to-leading (NLO) QCD predictions. In both cases, a good agreement was observed over the entire region of jet $`p_\mathrm{T}`$ (from $`50\mathrm{GeV}`$ up to about $`600\mathrm{GeV}`$) in which the cross section is rapidly falling down by 8 orders of magnitude. More detailed comparison with NLO QCD prediction is given in Fig. 4 (DØ) and Fig. 4 (CDF). CTEQ6.1M $`^\mathrm{?}`$ and MRST2004 $`^\mathrm{?}`$ parton distribution functions (PDF) where used in the NLO QCD calculations. Both sets of PDF lead to similar predictions. Data are sensitive to running of strong coupling $`\alpha _\mathrm{S}`$ and also to proton structure functions. At high $`p_\mathrm{T}`$, theoretical uncertainty is dominated by the uncertainty on gluon distribution function at high $`x`$ (where $`x`$ is fraction of proton momentum carried by gluon).
The CDF Collaboration also studied central jet production for jets reconstructed with $`k_\mathrm{T}`$-algorithm. CDF used Ellis-Sopper version of $`k_\mathrm{T}`$-algorithm $`^\mathrm{?}`$ adapted for hadron-hadron colliders. In this case, the et size is controlled by the parameter $`D`$. Obtained jet $`p_T`$ cross section for $`D=0.7`$ and its comparison with NLO QCD predictions lead to the same conclusions as in the case of cone jets. CDF performed the measurement for three different sizes of $`k_\mathrm{T}`$ jets ($`D=0.5`$, 0.7, and 1.0). For $`p_\mathrm{T}>150\mathrm{GeV}`$, there was, with respect to NLO QCD, no difference between them. At low $`p_T`$ end, observed jet production is higher above the NLO QCD prediction for higher values of $`D`$ (see Fig. 5). The results suggest that the larger jets are more sensitive to hadronization effects and/or to the soft underlying event physics.
DØ studied radiative processes in QCD by examining their impact on angular distributions. DØ measured the distribution of azimuthal angle between two jets with highest $`p_\mathrm{T}`$, $`\mathrm{\Delta }\varphi _{\mathrm{dijet}}`$ $`^\mathrm{?}`$. Second leading jet was required to have $`p_\mathrm{T}>40\mathrm{GeV}`$, and both leading jets were required to have rapidity $`|y|<0.5`$. Measurement was performed in four bins of leading jet transverse momentum $`p_T^{\mathrm{max}}`$. Fully corrected distribution of dijet azimuthal angle is presented in Fig. 7. As the data show, decorrelations increase with decrease of $`p_T^{\mathrm{max}}`$.
The first non-trivial description of $`\mathrm{\Delta }\varphi _{\mathrm{dijet}}`$ distribution is in pQCD given by the tree level $`23`$ parton matrix element. The limitations of this leading order (LO) prediction are apparent (see dashed line in Fig. 7). With three partons, it is imposible to produce final state with $`\mathrm{\Delta }\varphi {}_{\mathrm{dijet}}{}^{}<2\pi /3`$. NLO calculations, obtained with NLOJET++ $`^\mathrm{?}`$, provide much better agreement with data in much wider range of $`\mathrm{\Delta }\varphi _{\mathrm{dijet}}`$. However, in the region where $`\mathrm{\Delta }\varphi {}_{\mathrm{dijet}}{}^{}\pi `$, any fixed order pQCD calculations become unreliable. Resummations of soft parton emissions in all orders of pQCD are needed in order to describe this region properly.
General purpose Monte Carlo (MC) generators, like Herwig $`^\mathrm{?}`$ or pythia $`^\mathrm{?}`$, provide such resummations in the so called leading logarithm approximation through the developement of parton showers. Dijet azimuthal decorrelations are then sensitive to the details of the parton shower mechanism. Comparison between data and MC generator predictions is given in Fig. 7. herwig provides good desciption of the data in the whole range of $`\mathrm{\Delta }\varphi _{\mathrm{dijet}}`$, while pythia gives much smaller decorrelations than observed in the data. We found that pythia predictions were sensitive to the parameters of initial-state parton shower (ISR). Shaded region in Fig. 7 indicates the changes in azimuthal decorrelations as the maximal allowed virtuality of partons in the shower is increased from its default value by factor of four.
Another way how to study the effects of multiparton radiation is to examine an energy deposition within a jet, so called jet shapes. CDF measured averaged deposition of transverse momentum as a function of distance $`r`$ from jet exis:
$$\mathrm{\Psi }(r)=\frac{1}{N_{jet}}\underset{jets}{}\frac{p_\mathrm{T}(0,r)}{p_\mathrm{T}(0,R)},r=\sqrt{\mathrm{\Delta }^2\varphi +\mathrm{\Delta }^2\eta },$$
(1)
where $`p_\mathrm{T}(0,r)`$ is jet transverse momentum deposited in cone with radius $`r`$. Measurement was performed for the central rapidity jets in wide range of $`p_\mathrm{T}`$. Jets were reconstructed with midpoint cone algorithm $`^\mathrm{?}`$ with cone size $`R=0.7`$. An example of jet shape measurement in one low-$`p_\mathrm{T}`$ bin is given in Fig. 9, overall dependence on jet $`p_\mathrm{T}`$ is then summarized in Fig. 9. With increasing $`p_\mathrm{T}`$, jets are getting more narrow. This is due to two reasons: running of $`\alpha _\mathrm{S}`$, and change of proportion of gluon and quark jets. At low $`p_\mathrm{T}`$, the sample is dominated by gluon induced jets which are wider than quark jets, while at high $`p_\mathrm{T}`$, the sample is dominated by quark jets.
Herwig’s and pythia’s predictions are compared with the data as well in the two figures. pythia with default setting produces too narrow jets which indicates that there is not enough radiation in there. This is independent on pythia’s model of soft physics, multi-parton interactions (MPI), being switched on or off. In Run I, CDF tuned pythia to their data from studies of soft underlying event. So called Tune A gives very good description of jet shapes. Herwig describes data quite well, except the low-$`p_\mathrm{T}`$ region where the jets are narrower than in the data.
To summarize, new results on jet $`p_\mathrm{T}`$ spectra from CDF and DØ experiments has been presented. They are consistent with NLO QCD predictions. Aspects of multi-parton radiation were studied by the two experiments in the measurement of decorrelation in azimuthal angle between the two leading jets and in the measurement of jet shapes. NLO QCD provides good description of $`\mathrm{\Delta }\varphi _{\mathrm{dijet}}`$-distributions except the region where $`\mathrm{\Delta }\varphi {}_{\mathrm{dijet}}{}^{}\pi `$. The results are also useful for tuning parton shower models in MC generators.
|
warning/0506/cond-mat0506751.html
|
ar5iv
|
text
|
# 17O NMR Measurements on Superconducting Na0.35CoO₂⋅𝑦H2O
## Abstract
An <sup>17</sup>O NMR measurement was performed on nonoriented polycrystalline Na<sub>0.35</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O with superconducting (SC) transition temperature $`T_\text{c}=4.6`$ K. A weak temperature dependence was observed in the Knight shift at the O site ($`{}_{}{}^{17}K`$). The spin part of $`{}_{}{}^{17}K`$ ($`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$) is estimated from the plot of $`{}_{}{}^{17}K`$ against bulk susceptibility $`\chi `$. The $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ decreases in the SC state, indicative of the decrease in the in-plane component of the spin susceptibility. The nuclear spin-lattice relaxation rate $`1/T_1`$ at the O site $`{}_{}{}^{17}(1/T_1)`$ shows a good scaling with $`1/T_1`$ at the Co site $`{}_{}{}^{59}(1/T_1)`$. This indicates that the spin fluctuations at the O site originate from the Co spin dynamics. The relationships between $`{}_{}{}^{17}(1/T_1T)`$ and $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ and between $`{}_{}{}^{17}(1/T_1T)`$ and $`{}_{}{}^{59}(1/T_1T)`$ show the development of incommensurate fluctuations at q $``$ 0 other than q $`=`$ 0 below 30 K. A clear indication of ferromagnetic correlations at q $`=`$ 0 was not observed from the present <sup>17</sup>O-NMR studies.
superconductivity, hydrate sodium cobalt oxide, NMR, spin fluctuations
The recently discovered hydrate cobaltate superconductor Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O has attracted much attention, because its two-dimensional CoO<sub>2</sub> layers where superconductivity occurs form a triangular structure, in contrast to the tetragonal structure in cuprate superconductors. An interplay between superconductivity and geometrical frustrations is invoked due to the crystal structure. Various experiments as well as the spin-lattice relaxation rate $`1/T_1`$ in the superconducting (SC) state revealed that this superconductivity is of an unconventional type, where electron correlations play an important role in superconductivity. So far, we have shown the relationship between magnetic fluctuations and SC transition temperature $`T_\text{c}`$ from the Co nuclear-quadrupole-resonance (NQR) measurements on various samples with different values of $`T_\text{c}`$ : $`T_\text{c}`$ increases with increasing spin fluctuations at $`T_\text{c}`$, and the highest $`T_\text{c}`$ in the system is observed in the vicinity of the magnetic phase. In addition, it has been shown that the spin fluctuations correlate with NQR frequency $`\nu _\text{Q}`$, which is related to the distortion of the CoO<sub>6</sub> octahedron along the $`c`$-axis. Therefore, we suggest from an experimental point of view that the crystal-field splitting between the $`a_{1g}`$ and $`e_g^{}`$ states of Co-$`3d`$ $`t_{2g}`$ orbitals is an important parameter for determining superconductivity. Quite recently, a scenario based on experimental results has been proposed from a theoretical point of view . To verify this scenario, it is quite important to identify the properties of the magnetic fluctuations in the normal state and the symmetry of the SC pairs. Spin-triplet superconductivity induced by ferromagnetic fluctuations is expected in several theoretical models .
To date, there have been several reports about the <sup>59</sup>Co Knight shift ($`{}_{}{}^{59}K`$) in the SC state. Although a decrease in $`{}_{}{}^{59}K`$ was observed in the SC state, quantitative discussion was not carried out due to an ambiguity in the spin part of $`{}_{}{}^{59}K`$. This ambiguity originates mainly from the large electric field gradient (EFG) with a temperature dependence and the large orbital part of $`{}_{}{}^{59}K`$. Therefore, we change the NMR nucleus to <sup>17</sup>O, since <sup>17</sup>O has a small EFG frequency ($`\nu _Q`$) and a small orbital shift in general.
In this letter, we report on the spin part of the Knight shift at the O site ($`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$) estimated from a plot of $`{}_{}{}^{17}K`$ against $`\chi `$, and the behavior of $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ in the SC state. A decrease in $`{}_{}{}^{17}K`$ is observed in the field up to 8 T, which is consistent with the results in $`{}_{}{}^{59}K`$ . In addition, we found the development of AFM fluctuations below 30 K from the measurement of $`1/T_1`$ at the O site. Plausible spin-fluctuation character is discussed on the basis of our NMR experiments.
The sample we used in the measurement shows superconductivity at $`T_\text{c}4.6`$ K, which was determined by a dc susceptibility measurement. The resonance frequency of the $`\pm 5/2\pm 7/2`$ transition of Co NQR is 12.45 MHz, slightly larger than that of the highest-$`T_c`$ sample (12.40 MHz). According to the phase diagram we developed, this sample is situated near the border between superconductivity and magnetism. Strong magnetic fluctuations with critical character are expected in this sample. The detailed sample preparation and the <sup>17</sup>O exchange procedures will be reported elsewhere. <sup>16</sup>O nuclei in the CoO<sub>2</sub> layers are partially replaced by <sup>17</sup>O isotopes with the nuclear spin $`I=5/2`$, but O of the H<sub>2</sub>O are not. This is because the exchange annealing was carried out in anhydrate Na<sub>0.7</sub>CoO<sub>2</sub>, and then water was added to the compound. An <sup>17</sup>O NMR measurement was carried out using the nonoriented powder sample because mixing some material to fix the powdered sample orientation might degrade sample quality.
Figures 1 (a) and 1 (b) show NMR spectra taken at the frequencies of 24.7 and 47.8 MHz. The spectra show the typical powder-pattern structure with a non-zero asymmetric parameter $`\eta =(\nu _{xx}\nu _{yy})/\nu _{zz}`$ ($`\nu _{\alpha \alpha }`$ : EFG along the $`\alpha `$ direction). In general, when nuclei with a nuclear quadrupole moment are in a magnetic field, the total interaction is a sum of the Zeeman and electric-quadrupole (eqQ) interactions and is written as
$``$ $`=_{\mathrm{Zeeman}}+_{\mathrm{eqQ}}`$
$`=(1+K)\gamma _n\mathrm{}IH+{\displaystyle \frac{\mathrm{}\nu _\text{Q}}{6}}\left\{(3I_z^2I^2)+{\displaystyle \frac{1}{2}}\eta (I_+^2+I_{}^2)\right\}.`$
Here, $`K`$ and $`\gamma _n`$ are the Knight shift and nuclear gyromagnetic ratio, respectively. When the Zeeman interaction is dominant, $`_{\mathrm{eqQ}}`$ is treated as a perturbation. The eigenvalue of the total Hamiltonian depends on the angle between the applied field and the principal axis of EFG. In a nonoriented powder sample, the principal axis of the EFG is randomly distributed in all directions with respect to the external field. The powder-pattern spectrum observed for the nonoriented sample is reproduced using fitting parameters of center field $`H_0`$, $`\nu _\text{Q}`$, $`\eta `$ and $`K`$. Here, $`\nu _\mathrm{Q}`$ and $`\eta `$ are evaluated to be 0.168 MHz and 0.21, respectively, both of which are independent of applied field. We note that the $`\eta `$ at the O site is close to that at the Co site (0.208 $`\pm `$ 0.007) .
The powder-pattern spectrum of the <sup>17</sup>O-NMR becomes broader with increasing field, resulting in a structureless spectrum as seen in Fig. 1 (b). This is due to the large susceptibility at low temperatures. Although the <sup>17</sup>O-NMR spectrum becomes broader in higher magnetic fields, the Knight-shift measurement in high fields is effective for detecting the temperature dependence up to higher temperatures because the difference between the resonance field and the reference one becomes larger with increasing applied field. Fig. 1 (c) shows the central peak in the NMR spectrum recorded at various temperatures at the frequency of 68.9 MHz.
The temperature dependence of the Knight shift of <sup>17</sup>O ($`{}_{}{}^{17}K`$) was measured in a magnetic field of 13 T, and is shown in Fig. 2. Although the typical experimental error is approximately 0.02% due to the broad NMR peak, the Knight shift shows a temperature dependence beyond the error bars as shown in Fig. 2. The $`{}_{}{}^{17}K`$ gradually increases with decreasing temperature, which is scaled with $`\chi `$ measured in a magnetic field of 7 T. The inset of Fig. 2 shows the plot of $`{}_{}{}^{17}K`$ against $`\chi `$ with temperature as an implicit parameter. A good linear relationship was observed between $`\chi `$ and $`K`$, showing that the weak temperature dependence of bulk susceptibility is an intrinsic effect. From the linear relationship, the hyperfine coupling constant $`{}_{}{}^{17}A_{\mathrm{hf}}^{}`$ and temperature-independent Knight shift $`{}_{}{}^{17}K_{0}^{}`$ are evaluated to be 0.59 T/$`\mu _\mathrm{B}`$ and 0.021$`\pm 0.006`$%, respectively. Such a weak temperature dependence of $`{}_{}{}^{17}K`$ would be difficult to detect in small magnetic fields. $`{}_{}{}^{17}K_{0}^{}`$ is ascribed to the orbital shift and the shift as being due to the eqQ interaction. The spin part of $`{}_{}{}^{17}K`$ ($`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$) is estimated by subtracting $`{}_{}{}^{17}K_{0}^{}`$ from the total observed Knight shift. It is shown that the orbital susceptibility $`\chi _{\mathrm{orb}}`$ is negligibly small and that the spin susceptibility $`\chi _{\mathrm{spin}}`$ is dominant in the total susceptibility. Using the estimated $`\chi _{\mathrm{spin}}`$ and the Sommerfeld term ($`\gamma _{\mathrm{el}}=`$ 15 mJ/mol K<sup>2</sup>) obtained in the specific-heat measurement reported previously, the Wilson ratio $`R_\text{W}`$ is estimated to be $``$ 2.6. The Wilson ratio greater than two suggests the enhancement of the spin susceptibility compared with that estimated from the specific-heat measurement.
Next, we discuss the variation of $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ in the SC state. The resonance field of the central peak is shifted abruptly below $`T_\mathrm{c}`$. Figure 3(a) shows the relative shift of the central peak $`\delta H`$ below $`T_\mathrm{c}`$ with respect to the resonance field at $`T_\mathrm{c}`$. In general, $`\delta H`$ in the SC state is ascribed to the SC diamagnetic effect $`\mathrm{\Delta }H_{\mathrm{dia}}`$, and the decrease in spin susceptibility is given by $`\mathrm{\Delta }H_{\mathrm{spin}}=\mathrm{\Delta }K_{\mathrm{spin}}\times H_0`$, where $`H_0`$ is the applied field. These two effects exhibit different field dependences. $`\mathrm{\Delta }H_{\mathrm{dia}}`$ becomes smaller, whereas $`\mathrm{\Delta }H_{\mathrm{spin}}`$ becomes larger with increasing $`H_0`$. The increase in $`|\delta H|`$ when the field increases from 4 to 8 T suggests that the effect of $`\mathrm{\Delta }H_{\mathrm{spin}}`$ is dominant in this field range. Taking $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ 0.06% into consideration, we discuss quantitatively the temperature variation of $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ in the SC state. Figure 3(b) shows the temperature variation of $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$. Since $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ is determined from the central peak of the powder-pattern spectra, $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ corresponds to the isotropic term of $`{}_{}{}^{17}K`$, ($`{}_{}{}^{17}K_{\mathrm{iso}}^{}=(^{17}K_a+^{17}K_b+^{17}K_c)/3`$). According to the $`H_{c2}`$ versus $`T`$ phase diagram , superconductivity is almost destroyed when $`H`$ = 2 T is applied along the $`c`$ axis. Thus, we consider that $`{}_{}{}^{17}K_{c}^{}`$ is unchanged, and that $`{}_{}{}^{17}K_{a}^{}`$ and $`{}_{}{}^{17}K_{b}^{}`$ decrease in the SC state. The decrease in $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ shows that the in-plane component of the spin susceptibility decreases in the SC state. This suggests that the superconductivity is in a spin-singlet state ($`d`$-wave in this case), or in a spin-triplet state with the SC d-vector in the CoO<sub>2</sub> plane, since the spin susceptibility along the SC d-vector decreases when magnetic fields are applied parallel to the d-vector. The possibility of the spin-triplet state with the d-vector along the $`c`$-axis is excluded since the spin susceptibility should be unchanged in the fields perpendicular to the d-vector. The experimentally observed decrease in $`{}_{}{}^{17}K_{\mathrm{iso}}^{}`$ is reproduced by the two-dimensional line-node ($`\mathrm{\Delta }(\varphi )=\mathrm{cos}(2\varphi )`$) model with the singlet pairing, which incorporates the residual density of states DOS ($`N_{\mathrm{res}}`$) originating from the imperfection and/or inhomogeneity of the compound. The fitting parameters are $`2\mathrm{\Delta }/k_\mathrm{B}T_\mathrm{c}=3.5`$ and $`N_{\mathrm{res}}/N_0=0.32`$ ($`N_0`$: DOS at the Fermi level), which were determined to reproduce the temperature dependence of $`1/T_1`$ in the SC state . It should be noted that the onset temperature of the decrease in $`{}_{}{}^{17}K`$ is unchanged with respect to the applied field, which is in good agreement with the dc susceptibility measurement. It is likely that the insensitivity of $`T_\mathrm{c}`$ with respect to the strength of applied fields is related to the two-dimensional character of this superconductivity. From the theoretical point of view, the possibility of the spin-triplet state with the d-vector in the CoO<sub>2</sub> plane is proposed. It is quite important to measure $`{}_{}{}^{17}K_{c}^{}`$ in the SC state using an aligned-powder sample in order to identify the SC-pairing symmetry thoroughly.
The spin-lattice relaxation rate $`1/T_1`$ of <sup>17</sup>O was measured at the central peak in the powder-pattern spectrum. The recovery of the nuclear magnetization after saturation pulses can be fitted by the expected theoretical curve over the entire temperature and field range. $`1/T_1`$ was measured in various fields, but did not show any appreciable field dependence. Figure 4 shows the temperature dependence of $`1/T_1T`$ of <sup>17</sup>O ($`{}_{}{}^{17}(1/T_1T)`$) measured in magnetic field of 4 and 13 T, along with $`1/T_1T`$ of <sup>59</sup>Co measured by Co NQR. It was found that both results of $`1/T_1T`$ show the same temperature dependence. A similar result was recently reported by Ning and Imai. We plot $`{}_{}{}^{59}(1/T_1T)`$ against $`{}_{}{}^{17}(1/T_1T)`$ in the inset of Fig. 4. A good linear relationship is observed between the two quantities, indicating that $`{}_{}{}^{17}(1/T_1T)`$ arises from the spin dynamics at the Co site. The intercept of $`{}_{}{}^{59}(1/T_1T)`$ in the inset gives the relaxation rate induced by the orbital effect of the Co-$`3d`$ state, which is $`1.8\pm 2.3`$ s<sup>-1</sup> K<sup>-1</sup>. It was found that the orbital contribution in $`{}_{}{}^{59}(1/T_1T)`$ is negligibly smaller than the spin contribution. The slope in the inset gives the hyperfine coupling constant at the Co site, which is estimated to be $`{}_{}{}^{59}A_{\mathrm{hf}}^{}`$ 5.9 T/$`\mu _B`$ using $`{}_{}{}^{17}A_{\mathrm{hf}}^{}`$. This value is 1.2 times larger than the value estimated from the $`{}_{}{}^{59}K\chi `$ plot. Both results of $`1/T_1T`$ become enhanced below 100 K, indicative of the development of the spin fluctuations. If the $`q`$ dependence of the hyperfine coupling constant at the O site is taken into account, the AFM fluctuations at the q-vectors far from q $`=`$ 0 should be excluded because such AFM fluctuations are filtered out at the O site.
Now, we discuss the character of magnetic fluctuations in this compound on the basis of $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ and $`{}_{}{}^{17}(1/T_1T)`$. When $`K_{\mathrm{spin}}`$ and $`1/T_1T`$ arise from noninteracting electrons, the Korringa relation, which is described as ($`1/T_1TK^2)_0S_0=(\gamma _\text{e}/\gamma _\text{n})^2(\mathrm{}/4\pi k_\text{B})`$, holds. Here, $`\gamma _\mathrm{e}`$ is the gyromagnetic ratio of an electron. The character of spin fluctuations is discussed using the value of $`R`$, which is estimated from a comparison between the experimental value and $`S_0`$ as $`R^{17}(1/T_1TK_{\mathrm{spin}}^2)_{\mathrm{exp}}/S_0`$. $`R`$, which signifies the effect of electronic correlations, exceeds unity when antiferromagnetic (AFM) fluctuations are dominant, and it is much less than unity when ferromagnetic (FM) fluctuations become significant. The temperature dependence of $`R`$ is shown in the inset of Fig. 5. $`R`$ is temperature independent below 100 K to 20 K, and its value is $`1.2`$. The value being slightly larger than unity suggests that the moderate AFM fluctuations continue down to 20 K. $`R`$ increases and becomes $`2`$ at $`T_\mathrm{c}`$, suggesting that AFM fluctuations are enhanced at low temperatures. However, note that the magnetic fluctuations are considered to possess a critical character at low temperatures, which is suggested from the Co-NQR frequency and large $`{}_{}{}^{59}(1/T_1T)`$ value at $`T_\mathrm{c}`$. $`R2`$ at $`T_\mathrm{c}`$ and $`1.2`$ above 20 K are not so large, although the compound is situated close to the magnetic instability. Rather, if the larger Wilson ratio $`R_\mathrm{W}2.6`$ is taken into consideration, the presence of q $`=`$ 0 fluctuations is also suggested. The critical fluctuations with a moderate enhancement of $`R`$ may occur when the fluctuation spectrum of the AFM correlations has a peak close to q $``$ 0 and has significant contributions at q $`=`$ 0. This is consistent with the increase in $`{}_{}{}^{17}K\chi (0)`$ with decreasing temperature. AFM fluctuations at q far from q $`=`$ 0 would be excluded from the identical temperature dependences of $`{}_{}{}^{17}(1/T_1T)`$ and $`{}_{}{}^{59}(1/T_1T)`$. In order to reveal the entire $`q`$ structure in the spin-fluctuation spectrum, inelastic neutron studies of the SC compounds are highly desired.
Finally, we compare the present experimental results with the previous ones. The main panel of Fig. 5 shows the bulk susceptibilities of the present and previous samples measured at 1 T, along with $`{}_{}{}^{17}(1/T_1T)`$. In the previous paper, we suggested the presence of the FM fluctuations from the linear relationship between $`1/T_1T`$ and $`\chi `$ . However, the marked increase in $`\chi `$ below 30 K, which was observed for the previous sample, is not found for the present sample, although both samples show a gradual increase in $`\chi `$ in the temperature range between 150 and 30 K. Since $`\chi `$ in the present sample shows a linear relationship with $`{}_{}{}^{17}K`$, the moderate temperature dependence is evidenced as an intrinsic behavior, but the origin of the difference in $`\chi `$ below 30 K is not clear at the moment. In order to determine the intrinsic behavior in $`\chi `$ below 30 K, a systematic Knight-shift measurement of various samples is quite important.
In conclusion, from the Knight-shift measurement at the O site in the SC state, we found a decrease in $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ below $`T_\mathrm{c}`$ in the nonoriented powdered sample, which is consistent with the <sup>59</sup>Co-NMR results . The result suggests that the superconductivity is in the spin-singlet pairing state or spin-triplet state with the SC d-vector in the CoO<sub>2</sub> plane. $`{}_{}{}^{17}K_{\mathrm{spin}}^{}`$ is evaluated from the $`{}_{}{}^{17}K\chi `$ plot, and is compared with $`{}_{}{}^{17}(1/T_1T)`$. The development of AFM fluctuations is suggested below 30 K, and a clear indication of FM fluctuations was not observed in the <sup>17</sup>O NMR measurements. However, from the larger Wilson Ratio $`R_\mathrm{W}2.6`$ and the identical temperature dependences of $`{}_{}{}^{17}(1/T_1T)`$ and $`{}_{}{}^{59}(1/T_1T)`$, incommensurate fluctuations with small wave vectors are suggested to be enhanced at low temperatures, which have a significant contribution at q $`=`$ 0.
We thank C. Michioka and Y. Maeno for experimental support and valuable discussions. We also thank H. Ikeda, S. Fujimoto, K. Yamada, Y. Yanase, M. Mochizuki, and M. Ogata for valuable discussions. This work was partially supported by CREST of the Japan Science and Technology Corporation (JST) and the 21 COE program on “Center for Diversity and Universality in Physics” from MEXT of Japan, and by Grants-in-Aid for Scientific Research from the Japan Society for the Promotion of Science (JSPS) and MEXT.
|
warning/0506/math0506213.html
|
ar5iv
|
text
|
# Evaluation of Sylvester type determinants using block-triangularization
## 1 Introduction
Richard Askey in shows two ways, one matrix-theoretic and another based on orthogonal polynomials, to evaluate determinants
$$D_{N+1}(x):=\left|\begin{array}{cccccccc}x& 1& & & & & & \\ N& x& 2& & & 0& & \\ & N1& x& 3& & & & \\ & & \mathrm{}& & & & & \\ & 0& & & 2& x& N& \\ & & & & 0& 1& x& \end{array}\right|,$$
which were first considered by Sylvester . In addition, he obtains several generalizations of Sylvester’s determinants and explores their connection to orthogonal polynomials.
The purpose of this note is to show how the determinants from can be evaluated in yet another way, based on partial information about left or right eigenvectors of the corresponding matrices coupled with a simple similarity trick. In all cases except one, only one (the most obvious) eigenvector is used. The exceptional case is the Sylvester determinant itself, where two eigenvectors are readily available and hence used to derive the result.
## 2 Sylvester’s determinant and two close variants
Let us start with the Sylvester determinant. We want to prove that
$$D_{N+1}(x)=\underset{j=0}{\overset{N}{}}(x+N2j)$$
(1)
(which is formulas (2.3), (2.4) from ). Since the values of $`D_1`$ and $`D_2`$ agree with (1), it is enough to show that
$$D_{N+1}(x)=(xN)(x+N)D_{N1}(x)$$
(2)
for $`N1`$. The Sylvester determinant $`D_{N+1}`$ is the characteristic polynomial of the matrix $`𝒟_{N+1}`$, i.e., $`D_{N+1}(x)=_{\lambda \sigma (D_{N+1})}(x+\lambda )`$, where
$$𝒟_{N+1}:=\left[\begin{array}{ccccccc}0& 1& 0& 0& \mathrm{}& 0& 0\\ N& 0& 2& 0& \mathrm{}& 0& 0\\ 0& N1& 0& 3& \mathrm{}& 0& 0\\ 0& 0& N2& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& 0& N\\ 0& 0& 0& 0& \mathrm{}& 1& 0\end{array}\right].$$
Note that $`(1,1,1,\mathrm{},1)`$ is a left eigenvector of $`𝒟_{N+1}`$ corresponding to eigenvalue $`N`$ and $`(1,1,1,\mathrm{},(1)^N)`$ is its left eigenvector corresponding to eigenvalue $`N`$. The similarity transformation $`𝒟_{N+1}𝒯_{N+1}𝒟_{N+1}𝒯_{N+1}^1`$ where
$$𝒯_{N+1}:=\left[\begin{array}{cccccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1& \hfill \mathrm{}& 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1& \hfill \mathrm{}& (1)^N\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill \mathrm{}& 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill \mathrm{}& 0\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \mathrm{}\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill \mathrm{}& 1\end{array}\right]$$
therefore reduces $`𝒟_{N+1}`$ to a block lower-triangular form
$$\left[\begin{array}{cc}diag(N,N)& 0\\ & _{N1}\end{array}\right],$$
where
$$_{N1}:=\left[\begin{array}{ccccccc}0& 3+(N+1)& 0& N+1& 0& N+1& \mathrm{}\\ N2& 0& 4& 0& 0& 0& \mathrm{}\\ 0& N3& 0& 5& 0& 0& \mathrm{}\\ 0& 0& N4& 0& 6& 0& \mathrm{}\\ 0& 0& 0& N5& 0& 7& \mathrm{}\\ 0& 0& 0& 0& N6& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
Now it remains to show that $`_{N1}`$ is similar to $`𝒟_{N1}`$. But this is indeed so, since $`𝒮_{N1}^1_{N1}𝒮_{N1}=𝒟_{N1}`$ where
$$𝒮_{N1}:=\left[\begin{array}{cccccc}\hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill \mathrm{}\\ \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill \mathrm{}\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill \mathrm{}\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill \mathrm{}\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill \mathrm{}\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}\end{array}\right].$$
This proves (2) and therefore (1).
A related determinant
$$B_{N+1}(x):=\left|\begin{array}{cccccccc}x& a& 0& 0& \mathrm{}& 0& 0& 0\\ N(a1)& x1& 2a& 0& \mathrm{}& 0& 0& 0\\ 0& (N1)(a1)& x2& 3a& \mathrm{}& 0& 0& 0\\ & & \mathrm{}& & & & & \\ 0& 0& 0& 0& \mathrm{}& 2(a1)& x(N1)& Na\\ 0& 0& 0& 0& \mathrm{}& 0& a1& xN\end{array}\right|$$
can be evaluated analogously and even more simply. We use the notation
$$_{N+1}:=\left[\begin{array}{cccccccc}0& a& 0& 0& \mathrm{}& 0& 0& 0\\ N(a1)& 1& 2a& 0& \mathrm{}& 0& 0& 0\\ 0& (N1)(a1)& 2& 3a& \mathrm{}& 0& 0& 0\\ & & \mathrm{}& & & & & \\ 0& 0& 0& 0& \mathrm{}& 2(a1)& (N1)& Na\\ 0& 0& 0& 0& \mathrm{}& 0& a1& N\end{array}\right]$$
for the matrix that satisfies $`B_{N+1}(x)=_{\lambda \sigma (_{N+1})}(x+\lambda )`$.
To prove that
$$B_{N+1}(x)=\underset{j=0}{\overset{N}{}}[x+(N2j)aN+j]$$
(3)
(formula (2.8) from ), we will show that
$$B_{N+1}(x)=(x+NaN)B_N(xa).$$
(4)
The vector $`(1,1,1,\mathrm{},1)`$ is a left eigenvector of $`_{N+1}`$ corresponding to eigenvalue $`NaN`$. The similarity transformation $`_{N+1}𝒯_{N+1}_{N+1}𝒯_{N+1}^1`$ where
$$𝒯_{N+1}:=\left[\begin{array}{cccccc}1& 1& 1& 1& \mathrm{}& 1\\ 0& 1& 0& 0& \mathrm{}& 0\\ 0& 0& 1& 0& \mathrm{}& 0\\ 0& 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& 1\end{array}\right]$$
reduces $`_{N+1}`$ to a block lower triangular form
$$\left[\begin{array}{cc}NaN& 0\\ & _N\end{array}\right],$$
where
$$_N:=\left[\begin{array}{cccccc}1N(a1)& 2aN(a1)& N(a1)& N(a1)& N(a1)& \mathrm{}\\ (N1)(a1)& 2& 3a& 0& 0& \mathrm{}\\ 0& (N2)(a1)& 3& 4a& 0& \mathrm{}\\ 0& 0& (N3)(a1)& 4& 5a& \mathrm{}\\ 0& 0& 0& (N4)(a1)& 5& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
Finally, $`_N`$ is similar to $`_NaI`$:
$$_N𝒮_N=𝒮_N(_NaI),$$
where
$$𝒮_N:=\left[\begin{array}{cccccc}\hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill \mathrm{}& \hfill 0\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0& \hfill \mathrm{}& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1& \hfill \mathrm{}& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill \mathrm{}& \hfill 0\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill \mathrm{}& \hfill 1\end{array}\right].$$
This proves (3).
As is shown in , the formula (3) for $`B_{N+1}`$ implies the formula (1) for the Sylvester determinant $`D_{N+1}`$, since
$$D_{N+1}(x)=\underset{a\mathrm{}}{lim}\frac{B_{N+1}(ax)}{a^{N+1}}.$$
Also, (3) gives a formula for another related determinant,
$$A_{N+1}(x):=\left|\begin{array}{cccccccc}x& 1& 0& 0& \mathrm{}& 0& 0& 0\\ N& x2& 2& 0& \mathrm{}& 0& 0& 0\\ 0& (N1)& x4& 3& & & & \\ & \mathrm{}& & & & & & \\ 0& 0& 0& 0& & 2& x2(N1)& N\\ 0& 0& 0& 0& \mathrm{}& & 1& x2N\end{array}\right|$$
via the relation
$$A_{N+1}(x)=2^{N+1}B_{N+1}\left(\frac{x}{2}\right)|_{a=\frac{1}{2}}.$$
## 3 Determinants for Krawtchouk and dual Hahn polynomials
The determinant for the Krawtchouk polynomial $`K_{N+1}(x;p,N)`$ is given by
$`K_{N+1}(x;p,N):=`$
$`\left|\begin{array}{ccccccc}x+pN& pN& 0& 0& \mathrm{}& 0& 0\\ & & & & & & \\ \left(1p\right)& x+pN+\left(12p\right)& p\left(N1\right)& 0& \mathrm{}& 0& 0\\ & & & & & & \\ 0& 2\left(1p\right)& x+pN+2\left(12p\right)& p\left(N2\right)& \mathrm{}& 0& 0\\ & & & & & & \\ & & \mathrm{}& & & & \\ & & & & & & \\ 0& 0& 0& 0& \mathrm{}& N\left(1p\right)& x+pN+N\left(12p\right)\end{array}\right|.`$
With the usual notation $`(a)_k`$ for the shifted factorial (see \[1, formula (3.10)\]), we need to prove \[1, (3.25)\]:
$$K_{N+1}(x;p,N)=(x)_{N+1}.$$
(6)
For that, it is enough to show that
$$K_{N+1}(x;p,N)=(x)K_N(x;p,N1).$$
(7)
As usual, we need to find an ‘obvious’ eigenvector of the matrix
$$𝒦_{N+1}:=\left[\begin{array}{ccccccc}pN& pN& 0& 0& \mathrm{}& 0& 0\\ & & & & & & \\ (1p)& pN+(12p)& p(N1)& 0& \mathrm{}& 0& 0\\ & & & & & & \\ 0& 2(1p)& pN+2(12p)& p(N2)& \mathrm{}& 0& 0\\ & & & & & & \\ & & \mathrm{}& & & & \\ & & & & & & \\ 0& 0& 0& 0& \mathrm{}& N(1p)& pN+N(12p)\end{array}\right]$$
corresponding to eigenvalue $`0`$. In this case, it happens to be the right eigenvector $`(1,1,1,1,\mathrm{},(1)^N)^T`$. We then use the transformation
$$𝒯_{N+1}:=\left[\begin{array}{cccccc}\hfill 1& 0& 0& 0& \mathrm{}& 0\\ \hfill 1& 1& 0& 0& \mathrm{}& 0\\ \hfill 1& 0& 1& 0& \mathrm{}& 0\\ \hfill 1& 0& 0& 1& \mathrm{}& 0\\ \hfill \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \hfill (1)^N& 0& 0& 0& \mathrm{}& 1\end{array}\right]$$
(8)
to transform $`𝒦_{N+1}`$ into a block upper-triangular form
$$𝒯_{N+1}^1𝒦_{N+1}𝒯_{N+1}=\left[\begin{array}{cc}0& \\ 0& _N\end{array}\right],$$
where
$$_N:=\left[\begin{array}{cccccc}pN+PN+(12p)& p(N1)& 0& \mathrm{}& 0& 0\\ pN+2(1p)& pN+2(12p)& p(N2)& \mathrm{}& 0& 0\\ pN& 0& 3(1p)& \mathrm{}& 0& 0\\ pN& 0& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ (1)^NpN& 0& 0& \mathrm{}& N(1p)& pN+N(12p)\end{array}\right].$$
Now, $`_N`$ is similar to the matrix $`𝒦_N+I`$, viz.
$$𝒮_N_N=(𝒦_N+I)𝒮_N$$
via the transformation
$$𝒮_N:=\left[\begin{array}{cccccc}1& 0& 0& 0& \mathrm{}& 0\\ 1& 1& 0& 0& \mathrm{}& 0\\ 0& 1& 1& 0& \mathrm{}& 0\\ 0& 0& 1& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& 1\end{array}\right].$$
(9)
This proves (7) and hence (6).
The determinant for dual Hahn polynomials,
$`R_{N+1}(\lambda \left(x\right);\gamma ,\delta ,N):=`$
$`\left|\begin{array}{ccccc}\lambda \left(x\right)+N\left(\gamma +1\right)& N\left(\gamma +1\right)& 0& 0& \\ & & & & \\ \left(N+\delta \right)& \begin{array}{c}\hfill \lambda \left(x\right)+N\left(\gamma +3\right)\\ \hfill \left(\gamma \delta +2\right)\end{array}& \left(N1\right)\left(\gamma +2\right)& 0& \mathrm{}\\ & & & & \\ 0& 2\left(N+\delta 1\right)& \begin{array}{c}\hfill \lambda \left(x\right)+N\left(\gamma +5\right)\\ \hfill 2\left(\gamma \delta +4\right)\end{array}& \left(N2\right)\left(\gamma +3\right)& \\ & & & & \\ & & \mathrm{}& & \\ & & & & \\ 0& 0& \mathrm{}& N\left(\delta +1\right)& \left\{\begin{array}{c}\hfill \lambda \left(x\right)+N\left(2N+\gamma +1\right)\\ \hfill N\left(2N+\gamma \delta \right)\end{array}\right\}\end{array}\right|,`$
where
$$\lambda (x):=x(x+\gamma +\delta +1),$$
appears in as the first example beyond Krawtchouk polynomials in terms of simplicity of recurrence coefficients. This determinant seems much harder to evaluate, but this in fact requires just an additional shift of parameters. The formula for $`R_{N+1}`$ is \[1, (4.5)\]:
$$R_{N+1}(\lambda (x);\gamma ,\delta ,N)=(x)_N(x+\gamma +\delta +1)_{N+1}.$$
(14)
It can be proved from the relation
$$R_{N+1}(\lambda (x);\gamma ,\delta ,N)=\lambda (x)R_N(\lambda (x)+(\gamma +\delta +2);\gamma +1,\delta +1,N1).$$
(15)
The latter can be checked exactly as above, i.e., by applying the transformation $`𝒯_{N+1}`$ of the form (8) to the matrix
$`_{N+1}(\gamma ,\delta ):=`$
$`\left[\begin{array}{ccccc}N(\gamma +1)& N(\gamma +1)& 0& 0& \\ & & & & \\ (N+\delta )& \begin{array}{c}\hfill N(\gamma +3)\\ \hfill (\gamma \delta +2)\end{array}& (N1)(\gamma +2)& 0& \mathrm{}\\ & & & & \\ 0& 2(N+\delta 1)& \begin{array}{c}\hfill N(\gamma +5)\\ \hfill 2(\gamma \delta +4)\end{array}& (N2)(\gamma +3)& \\ & & & & \\ & & \mathrm{}& & \\ & & & & \\ 0& 0& \mathrm{}& N(\delta +1)& \left\{\begin{array}{c}\hfill N(2N+\gamma +1)\\ \hfill N(2N+\gamma \delta )\end{array}\right\}\end{array}\right],`$
obtaining a block upper-triangular matrix
$$𝒯_{N+1}^1_{N+1}(\gamma ,\delta )𝒯_{N+1}=\left[\begin{array}{cc}0& \\ 0& _N\end{array}\right],$$
where
$`_N(\gamma ,\delta ):=`$
$`\left[\begin{array}{ccccc}\begin{array}{c}\hfill N(\gamma +1)+N(\gamma +3)\\ \hfill (\gamma \delta 2)\end{array}& (N1)(\gamma +2)& 0& 0& \mathrm{}\\ & & & & \\ \begin{array}{c}\hfill N(\gamma +1)\\ \hfill +2(N+\delta 1)\end{array}& \begin{array}{c}\hfill N(\gamma +5)\\ \hfill 2(\gamma \delta +4)\end{array}& (N2)(\gamma +3)& 0& \mathrm{}\\ & & & & \\ N(\gamma +1)& 3(N+\delta 2)& \begin{array}{c}\hfill N(\gamma +7)\\ \hfill 3(\gamma \delta +6)\end{array}& (N3)(\gamma +4)& \mathrm{}\\ & & & & \\ N(\gamma +1)& 0& 4(N+\delta +3)& \begin{array}{c}\hfill N(\gamma +9)\\ \hfill 4(\gamma \delta +8)\end{array}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right],`$
and then transforming $`_N(\gamma ,\delta )`$ into $`_N(\gamma +1,\delta +1)+(\gamma +\delta +2)I`$ by
$$_N(\gamma ,\delta )𝒮_N_N(\gamma ,\delta )𝒮_N^1,$$
where $`𝒮_N`$ has the form (9). This proves (15) and (14).
## 4 Determinants for Hahn, Racah, and $`q`$-Racah polynomials
The last examples appearing in are three types of orthogonal polynomials, whose recurrence coefficients are fractional rather than polynomial. Recall that the recurrence coefficients appear as entries on the main diagonal and the first sub- and super- diagonals of the tridiagonal determinant that gives the value of the corresponding polynomial. The alternating sums of these coefficients are still zero, just like in the two examples we just considered in section 3. In other words, the determinants we now consider all have the form
$$\left|\begin{array}{ccccccc}\lambda (x)b_0& a_0& 0& \mathrm{}& 0& 0& 0\\ c_1& \lambda (x)b_1& a_1& \mathrm{}& 0& 0& 0\\ 0& c_2& \lambda (x)b_2& \mathrm{}& 0& 0& 0\\ & \mathrm{}& & & & & \\ 0& 0& 0& \mathrm{}& c_{N1}& \lambda (x)b_{N1}& a_{N1}\\ 0& 0& 0& \mathrm{}& 0& c_N& \lambda (x)b_N\end{array}\right|\mathit{=:}det(\lambda (x)+𝒢_{N+1}),$$
(25)
where
$$b_n=a_nc_n\text{for all}n.$$
(26)
The coefficients $`a_n`$ and $`c_n`$ for Hahn polynomials are
$`a_n`$ $`=`$ $`{\displaystyle \frac{(n+\alpha +\beta +1)(n+\alpha +1)(Nn)}{(2n+\alpha +\beta +1)(2n+\alpha +\beta +2)}}`$
$`c_n`$ $`=`$ $`{\displaystyle \frac{n(n+\alpha +\beta +N+1)(n+\beta )}{(2n+\alpha +\beta )(2n+\alpha +\beta +1)}},`$
with $`\lambda (x)=x`$, while the coefficients for Racah polynomials are
$`a_n`$ $`=`$ $`{\displaystyle \frac{(n+\alpha +1)(n+\alpha +\beta +1)(n+\gamma +1)(Nn)}{(2n+\alpha +\beta +1)(2n+\alpha +\beta +2)}}`$
$`c_n`$ $`=`$ $`{\displaystyle \frac{n(n+\alpha +\beta +N+1)(n+\alpha +\beta \gamma )(n+\beta )}{(2n+\alpha +\beta )(2n+\alpha +\beta +1)}},`$
with $`\beta +\gamma +1=N`$ and $`\lambda (x)=x(x+\gamma +\delta )`$, and the coefficients for $`q`$-Racah polynomials are
$`a_n`$ $`=`$ $`{\displaystyle \frac{(1abq^{n+1})(1q^{n+1})(1q^{nN})(1cq^{n+1})}{(1abq^{2n+1})(1abq^{2n+2})}}`$ (27)
$`c_n`$ $`=`$ $`{\displaystyle \frac{cq^N}{b}}{\displaystyle \frac{(1q^n)(1bq^n)(1abc^1q^n)(1abq^{n+N+1})}{(1abq^{2n})(1abq^{2n+1})}},`$ (28)
with $`bdq=q^N`$ and $`\lambda (x)=(q^x)(1q^{x+1}cd)`$ (see \[1, Sec. 4\]).
Recall that Hahn and Racah polynomials are just limiting cases of $`q`$-Racah polynomials \[1, Sec. 4\]. Precisely, if we denote Hahn polynomials by $`H_{N+1}(\lambda (x);\alpha ,\beta ,N)`$, Racah polynomials by $`RA_{N+1}(\lambda (x);\alpha ,\beta ,\gamma ,N)`$, and $`q`$-Racah polynomials by $`QRA_{N+1}(\lambda (x);q,a,b,c,N)`$, then we see that
$`RA_{N+1}(\lambda (x);\alpha ,\beta ,\gamma ,N)`$ $`=`$ $`\underset{q1}{lim}(1)^{N+1}{\displaystyle \frac{QRA_{N+1}(\lambda (x);q,q^\alpha ,q^\beta ,q^\gamma ,N)}{(1q)^{2N+2}}},`$
$`H_{N+1}(\lambda (x);\alpha ,\beta ,N)`$ $`=`$ $`\underset{\gamma \mathrm{}}{lim}{\displaystyle \frac{RA_{N+1}(\lambda (x);\alpha ,\beta ,\gamma ,N)}{\gamma ^{N+1}}}.`$
So, it is enough to evaluate $`q`$-Racah polynomials. The formula is given in in the form
$$QRA_{N+1}(\lambda (x);q,a,b,c,N)=(1)^{N+1}(q^x;q)_{N+1}(q^{x+1}cd;q)_{N+1},$$
where $`(;)_k`$ is the $`q`$-analogue of the shifted factorial \[1, formula (4.13)\]:
$$(a;q)_k:=\underset{j=0}{\overset{k1}{}}(1aq^j).$$
Its equivalent form that is more suitable for an inductive proof is
$$QRA_{N+1}(\lambda (x);q,a,b,c,N)=\underset{n=0}{\overset{N}{}}(\lambda (x)\lambda (n)).$$
(29)
This formula will be proved once we show that
$$QRA_{N+1}(\lambda (x);q,a,b,c,N)=\lambda (x)q^NQRA_N(\stackrel{~}{\lambda }(x1);q,aq,b,cq,N1),$$
(30)
where $`\stackrel{~}{\lambda }(x)`$ corresponds to the parameters $`(q,aq,b,cq,N1)`$ so that
$$\stackrel{~}{\lambda }(x1)=(1q^{x+1})(1\frac{cq}{bq^N}q^x)=q\left(\lambda (x)+1+\frac{c}{bq^N}\frac{1}{q}\frac{cq}{bq^N}\right).$$
(31)
To prove (30), let us start with an observation about our ansatz matrices $`𝒢_{N+1}`$ satisfying (25) and (26). Suppose that such a $`𝒢_{N+1}`$ is transformed using the matrix $`𝒯_{N+1}`$ given in (8). Then, as we already saw, $`𝒯_{N+1}^1𝒢_{N+1}𝒯_{N+1}`$ is block upper triangular:
$$\left[\begin{array}{cc}0& \\ 0& _N\end{array}\right].$$
Next, the transformation $`𝒮_N`$ given by (9) reduces $`_N`$ to the tridiagonal form
$$𝒮_N_N𝒮_N^1=\left[\begin{array}{cccccccc}a_0+c_1& a_1& 0& 0& \mathrm{}& 0& 0& 0\\ c_1& a_1+c_2& a_2& 0& \mathrm{}& 0& 0& 0\\ 0& c_2& a_2+c_3& a_3& \mathrm{}& 0& 0& 0\\ & & \mathrm{}& & & & & \\ 0& 0& 0& 0& \mathrm{}& c_{N2}& a_{N2}+c_{N1}& a_{N1}\\ 0& 0& 0& 0& \mathrm{}& 0& c_{N1}& a_{N1}+c_N\end{array}\right].$$
(32)
Proving (30) therefore amounts to showing that the matrix (32), with the entries $`a_n`$ and $`c_n`$ coming from the determinant $`QRA_{N+1}(\lambda (x);q,a,b,c,N)`$, is similar to the matrix $`\frac{1}{q}(𝒢_N+\stackrel{~}{\lambda }(x1)I_N)\lambda (x)I_N`$, where $`𝒢_N`$ is an $`N\times N`$-matrix of the form (25)–(26), with the entries $`a_n`$ and $`c_n`$ coming from the determinant $`QRA_N(\lambda (x1);q,aq,b,cq,N1)`$. (Note that, due to (31), the difference $`\lambda (x)\frac{1}{q}\stackrel{~}{\lambda }(x1)`$ does not depend on $`x`$.) The needed similarity is realized by the diagonal matrix
$$\mathrm{\Lambda }_N:=diag((1abq^2),\frac{1}{q}(1abq^4),\frac{1}{q^2}(1abq^6),\mathrm{},\frac{1}{q^{N1}}(1abq^{2N})).$$
so that
$$\mathrm{\Lambda }_N^1𝒮_N_N𝒮_N^1\mathrm{\Lambda }_N=\frac{1}{q}(𝒢_N+\stackrel{~}{\lambda }(x1)I_N)\lambda (x)I_N.$$
(33)
Verification of (33) is straightforward for off-diagonal entries. For diagonal entries, it reduces to verification of the identity
$$\left(a_n+c_{N+1}1\frac{c}{bq^N}+\frac{1}{q}+\frac{cq}{bq^N}\right)q=a_{n+1}\frac{1abq^{2n+4}}{1abq^{2n+2}}+q^2c_n\frac{1abq^{2n}}{1abq^{2n+2}},$$
where $`a_n`$ are given by (27) and $`c_n`$ by (28). This last identity can be checked using MATLAB Symbolic Math Toolbox. This finishes the proof of (30) and (29).
The determinant for Racah polynomials is therefore
$$RA_{N+1}(\lambda (x);\alpha ,\beta ,\gamma ,N)=(x)_{N+1}(x+\gamma +\delta +1)_{N+1}$$
and the determinant for Hahn polynomials is
$$H_{N+1}(\lambda (x);\alpha ,\beta ,N)=(x)_{N+1}.$$
|
warning/0506/math0506458.html
|
ar5iv
|
text
|
# Linearization coefficients of Bessel polynomials
## 1 Introduction
In this paper we consider the Bessel polynomials $`q_n`$ of degree $`n`$
$$q_n\left(u\right)=\underset{k=0}{\overset{n}{}}\alpha _k^{(n)}u^k,$$
(1)
where
$$\alpha _k^{(n)}=\frac{\left(\genfrac{}{}{0pt}{}{n}{k}\right)}{\left(\genfrac{}{}{0pt}{}{2n}{k}\right)}\frac{2^k}{k!}=\frac{n!(2nk)!\mathrm{\hspace{0.17em}2}^k}{(2n)!(nk)!k!}.$$
(2)
The first examples of these polynomials are
$$q_0\left(u\right)=1,q_1\left(u\right)=1+u,q_2\left(u\right)=1+u+\frac{u^2}{3}.$$
They are normalized according to
$$q_n\left(0\right)=1,$$
and thus differ from the polynomials $`\theta _n\left(u\right)`$ in by the constant factor $`\frac{\left(2n\right)!}{n!2^n},`$ i.e.
$$\theta _n(u)=\frac{(2n)!}{n!2^n}q_n(u).$$
For $`\nu >0`$ we recall that the probability density on $``$
$$f_\nu (x)=\frac{A_\nu }{(1+x^2)^{\nu +{\scriptscriptstyle \frac{1}{2}}}},A_\nu =\frac{\mathrm{\Gamma }(\nu +\frac{1}{2})}{\mathrm{\Gamma }(\frac{1}{2})\mathrm{\Gamma }(\nu )}$$
(3)
has the characteristic function
$$_{\mathrm{}}^{\mathrm{}}e^{ixy}f_\nu (x)𝑑x=k_\nu (|y|),y,$$
(4)
where
$$k_\nu (u)=\frac{2^{1\nu }}{\mathrm{\Gamma }\left(\nu \right)}u^\nu K_\nu \left(u\right),u0,$$
(5)
and $`K_\nu `$ is the modified Bessel function of the third kind. If $`\nu =n+\frac{1}{2}`$ with $`n=0,1,2,\mathrm{}`$ then
$$k_\nu (u)=e^uq_n(u),u0,$$
(6)
and $`f_\nu `$ is called a Student-t density with $`2\nu =2n+1`$ degrees of freedom. For $`\nu =\frac{1}{2}`$ then $`f_\nu `$ is density of a Cauchy distribution. Note that for simplicity we have avoided the usual scaling of the Student-t distribution.
In this paper, we provide the solutions of the three following problems: assuming in the rest of the paper that $`a`$ is a constant with $`0a1`$,
1. positivity and explicit values of the connection coefficients $`c_k^{\left(n\right)}(a)`$ in the expansion
$$q_n\left(au\right)=\underset{k=0}{\overset{n}{}}c_k^{\left(n\right)}\left(a\right)q_k\left(u\right),$$
(7)
2. positivity and explicit value of the linearization coefficients $`\beta _i^{\left(n\right)}(a)`$ in the expansion
$$q_n(au)q_n((1a)u)=\underset{i=0}{\overset{n}{}}\beta _i^{\left(n\right)}(a)q_{n+i}\left(u\right),$$
(8)
3. positivity of the linearization coefficients $`\beta _k^{(n,m)}(a)`$ in the expansion
$$q_n(au)q_m((1a)u)=\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}(a)q_k\left(u\right).$$
(9)
Note that $`\beta _i^{(n)}(a)=\beta _{n+i}^{(n,n)}(a)`$ and that (7) is a special case of (9) corresponding to $`m=0`$ with $`c_k^{(n)}(a)=\beta _k^{(n,0)}(a)`$. Note also that $`u=0`$ in (9) yields
$$\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}(a)=1,$$
so (9) is a convex combination. As polynomial identities, (7)-(9) of course hold for all complex $`a,u`$, but as we will see later, the positivity of the coefficients holds only for $`0a1.`$
Because of (4) and (6) formula (9) is equivalent with the following identity between Student-t densities
$$\frac{1}{a}f_{n+{\scriptscriptstyle \frac{1}{2}}}\left(\frac{x}{a}\right)\frac{1}{1a}f_{m+{\scriptscriptstyle \frac{1}{2}}}\left(\frac{x}{1a}\right)=\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}(a)f_{k+{\scriptscriptstyle \frac{1}{2}}}(x)$$
(10)
for $`0<a<1`$ and $``$ is the ordinary convolution of densities.
Although (9) is more general than (7),(8), we stress that we give explicit formulas below for $`c_k^{(n)}(a)`$ and $`\beta _i^{(n)}(a)`$ from which the positivity is clear. The positivity of $`\beta _k^{(n,m)}(a)`$ for the general case can be deduced from the special cases via a recursion formula, see Lemma 3.4 below.
These problems have an important application in statistics: the Behrens-Fisher problem consists in testing the equality of the means of two normal populations. Fisher <sup>1</sup><sup>1</sup>1the collected papers of R.A. Fisher are available at the following address http://www.library.adelaide.edu.au/digitised/fisher/ has shown that this test can be performed using the $`d`$statistics defined as
$$d_{f_1,f_2,\theta }=t_1\mathrm{sin}\theta t_2\mathrm{cos}\theta ,$$
where $`t_1`$ and $`t_2`$ are two independent Student-t random variables with respective degrees of freedom $`f_1`$ and $`f_2`$ and $`\theta [0,\frac{\pi }{2}]`$. Many different results have been obtained on the behaviour of the $`d`$statistics. Tables of the distribution of $`d_{f_1,f_2,\theta }`$ have been provided in 1938 by Sukhatme at Fisher’s suggestion. In 1956, Fisher and Healy explicited the distribution of $`d_{f_1,f_2,\theta }`$ as a mixture of Student-t distributions (Student-t distribution with a random, discrete number of degrees of freedom) for small, odd values of $`f_1`$ and $`f_2`$. This work was extended by Walker and Saw who provided, still in the case of odd numbers of degrees of freedom, an explicit way of computing the coefficients of the Student-t mixture as solutions of a linear system; however, they did not prove the positivity of these coefficients, claiming only
> “Extensive numerical investigation indicates also that $`\eta _i0`$ for all i; however, an analytic proof has not been found.”
This conjecture is proved in Theorem 2 and 3 below. Section 2 of this paper gives the explicit solutions to problems 1, 2 and 3, whereas section 3 is dedicated to their proofs. The last section gives an extension of Theorem 2 in terms of inverse Gamma distributions.
Using the fact that the Student-t distribution is a scale mixture of normal distributions by an inverse gamma distribution our positivity result is equivalent to an analogous positivity result for inverse gamma distributions. This result has been observed for small values of the degrees of freedom in . In the coefficients are claimed to be non-negative but the paper does not contain any arguments to prove it.
## 2 Results
### 2.1 Solution of problem 1 and a stochastic interpretation
###### Theorem 2.1
The coefficients $`c_k^{\left(n\right)}\left(a\right)`$ in (7) write
$$c_k^{\left(n\right)}\left(a\right)=a^k\frac{\left(\genfrac{}{}{0pt}{}{n}{k}\right)}{\left(\genfrac{}{}{0pt}{}{2n}{2k}\right)}\underset{r=1}{\overset{\left(nk\right)\left(k+1\right)}{}}\left(\genfrac{}{}{0pt}{}{n+1}{k+1r}\right)\left(\genfrac{}{}{0pt}{}{nk1}{r1}\right)\left(1a\right)^r$$
for $`0kn1`$ while $`c_n^{(n)}(a)=a^n`$ and, as $`0a1,`$ they are positive.
A stochastic interpretation of Theorem 2.1 writes as follows: replacing u by $`|u|`$ and multiplying equation (7) by $`\mathrm{exp}\left(|u|\right)`$, we obtain
$$e^{\left(1a\right)|u|}e^{a|u|}q_n\left(a|u|\right)=\underset{k=0}{\overset{n}{}}c_k^{\left(n\right)}\left(a\right)q_k(|u|)e^{|u|}.$$
(11)
Equality (11) can be interpreted as follows: the convex combination of an independent Cauchy variable $`C`$ and a Student-t variable $`X_n`$ with $`2n+1`$ degrees of freedom follows a Student-t distribution with random number $`2K\left(\omega \right)+1`$ of degrees of freedom:
$$\left(1a\right)C+aX_n\stackrel{𝑑}{=}X_{K\left(\omega \right)},$$
where $`K\left(w\right)[0,n]`$ is a discrete random variable such that
$$\mathrm{Pr}\{K\left(w\right)=k\}=c_k^{\left(n\right)}\left(a\right),0kn.$$
### 2.2 Solution of problem 2 and a probabilistic interpretation
###### Theorem 2.2
The coefficients $`\beta _i^{\left(n\right)}\left(a\right)`$ in (8) write
$`\beta _i^{(n)}(a)`$ $`=`$ $`(4a(1a))^i\left({\displaystyle \frac{n!}{(2n)!}}\right)^22^{2n}{\displaystyle \frac{(2n2i)!(2n+2i)!}{(ni)!(n+i)!}}`$
$`\times `$ $`{\displaystyle \underset{j=0}{\overset{ni}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{2j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{i}}\right)(2a1)^{2j}`$
and, as $`0a1,`$ they are positive.
A probabilistic interpretation of this result writes as follows.
###### Corollary 2.3
With $`a=\mathrm{sin}\theta `$, $`f_1=f_2=2n+1,`$ statistic $`d_{f_1,f_2,\theta }`$ follows a Student-t distribution with a random number of degrees of freedom $`F\left(\omega \right)`$ distributed according to
$$\mathrm{Pr}\left\{F\left(\omega \right)=2n+2i+1\right\}=\beta _i^{\left(n\right)}\left(a\right),0in.$$
### 2.3 Problem 3
###### Theorem 2.4
The coefficients $`\beta _k^{(n,m)}(a)`$ in (9) are positive for $`0a1.`$
We were unable to derive the explicit values of the coefficients $`\beta _k^{(n,m)}(a)`$: however, their positivity allows us to claim the following Corollary.
###### Corollary 2.5
With $`a=\mathrm{sin}\theta `$, $`f_1=2n+1,f_2=2m+1,`$ statistic $`d_{f_1,f_2,\theta }`$ follows a Student-t distribution with a random number of degrees of freedom $`F\left(\omega \right)`$ distributed according to
$$\mathrm{Pr}\left\{F\left(\omega \right)=2k+1\right\}=\beta _k^{(n,m)}\left(a\right),nmkn+m.$$
###### Theorem 2.6
For $`k2`$ let $`n_1,\mathrm{},n_k`$ be nonnegative integers and let $`a_1,\mathrm{},a_k`$ be positive real numbers with sum 1. Then
$$q_{n_1}(a_1u)q_{n_2}(a_2u)\mathrm{}q_{n_k}(a_ku)=\underset{j=l}{\overset{L}{}}\beta _jq_j(u),u$$
(12)
with nonnegative coefficients $`\beta _j`$ with sum 1 and $`l=\mathrm{min}(n_1,\mathrm{},n_k),L=n_1+\mathrm{}+n_k`$.
## 3 Proofs
### 3.1 Generalities about Bessel polynomials
As a preparation to the proofs we give some recursion formulas for $`q_n`$. They follow from corresponding formulas for $`\theta _n`$ from , but they can also be proved directly from the definitions (1) and (2). The formulas are
$$q_{n+1}(u)=q_n(u)+\frac{u^2}{4n^21}q_{n1}(u),n1,$$
(13)
$$q_n^{}(u)=q_n(u)\frac{u}{2n1}q_{n1}(u),n1.$$
(14)
We can write
$$u^n=\underset{i=0}{\overset{n}{}}\delta _i^{(n)}q_i(u),n=0,1,\mathrm{}$$
(15)
and $`\delta _i^{(n)}`$ is given by a formula due to Carlitz , see \[9, p. 73\] or :
$$\delta _i^{\left(n\right)}=\{\begin{array}{cc}\frac{\left(n+1\right)!}{2^n}\frac{\left(1\right)^{ni}\left(2i\right)!}{\left(ni\right)!i!\left(2i+1n\right)!}& \text{for }\frac{n1}{2}in\\ 0& \text{for }0i<\frac{n1}{2}\end{array}.$$
(16)
Later we need the following extension of (13) which we formulate using the Pochhammer symbol $`(z)_n:=z(z+1)\mathrm{}(z+n1)`$ for $`z,n=0,1,\mathrm{}.`$
###### Lemma 3.1
For $`0kn`$ we have
$$u^{2k}q_{nk}(u)=\underset{i=0}{\overset{k}{}}\gamma _i^{(n,k)}q_{n+i}(u)$$
where
$$\gamma _i^{(n,k)}=2^{2k}\left(\genfrac{}{}{0pt}{}{k}{i}\right)(nk+\frac{1}{2})_{k+i}(n\frac{1}{2})_{ki}.$$
(17)
Proof: The Lemma is trivial for $`k=0`$ and reduces to the recursion (13) for $`k=1`$ written as
$$u^2q_{n1}(u)=2^2(n\frac{1}{2})_2\left(q_{n+1}(u)q_n(u)\right).$$
(18)
We will prove the formula (17) by induction in $`n`$, so assume it holds for some $`n`$ and all $`0kn`$. Multiplying the formula of the lemma by $`u^2`$ we get
$$u^{2k+2}q_{nk}(u)=\underset{i=0}{\overset{k}{}}\gamma _i^{(n,k)}u^2q_{n+i}(u),$$
hence by (18)
$`u^{2(k+1)}q_{n+1(k+1)}(u)={\displaystyle \underset{i=0}{\overset{k}{}}}\gamma _i^{(n,k)}2^2(n+i+\frac{1}{2})_2\left[q_{n+i+2}(u)q_{n+i+1}(u)\right]`$
$`=`$ $`\gamma _k^{(n,k)}2^2(n+k+\frac{1}{2})_2q_{n+k+2}(u)`$
$`+`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}2^2(n+i+\frac{1}{2})\left[\gamma _{i1}^{(n,k)}(n+i\frac{1}{2})\gamma _i^{(n,k)}(n+i+\frac{3}{2})\right]q_{n+1+i}(u)`$
$``$ $`\gamma _0^{(n,k)}2^2(n+\frac{1}{2})_2q_{n+1}(u).`$
Using the induction hypothesis we easily get
$$\gamma _k^{(n,k)}2^2(n+k+\frac{1}{2})_2=2^{2k+2}(nk+\frac{1}{2})_{2k+2}=\gamma _{k+1}^{(n+1,k+1)},$$
and
$$\gamma _0^{(n,k)}2^2(n+\frac{3}{2})(n+\frac{1}{2})=2^{2k+2}(nk+\frac{1}{2})_{k+1}(n\frac{3}{2})_{k+1}=\gamma _0^{(n+1,k+1)}.$$
Concerning the coefficient $`C`$ to $`q_{n+1+i}(u)`$ above we have
$$C=2^{2k+2}(n+i+\frac{1}{2})[\left(\genfrac{}{}{0pt}{}{k}{i1}\right)(nk+\frac{1}{2})_{k+i1}(n\frac{1}{2})_{ki+1}(n+i\frac{1}{2})$$
$$\left(\genfrac{}{}{0pt}{}{k}{i}\right)(nk+\frac{1}{2})_{k+i}(n\frac{1}{2})_{ki}(n+i+\frac{3}{2})]$$
$`=`$ $`2^{2k+2}(nk+\frac{1}{2})_{k+1+i}(n\frac{1}{2})_{ki}`$
$`\times `$ $`\left[\left({\displaystyle \genfrac{}{}{0pt}{}{k}{i1}}\right)(n\frac{1}{2}+ki)\left({\displaystyle \genfrac{}{}{0pt}{}{k}{i}}\right)(n+i+\frac{3}{2})\right]`$
$`=`$ $`2^{2k+2}(nk+\frac{1}{2})_{k+1+i}(n\frac{1}{2})_{ki}\left[\left({\displaystyle \genfrac{}{}{0pt}{}{k+1}{i}}\right)(n\frac{3}{2})\right]`$
$`=`$ $`2^{2k+2}\left({\displaystyle \genfrac{}{}{0pt}{}{k+1}{i}}\right)(nk+\frac{1}{2})_{k+1+i}(n\frac{3}{2})_{k+1i}=\gamma _i^{(n+1,k+1)}.`$
$`\mathrm{}`$
We stress that Lemma 3.1 is the special case $`\nu =n+\frac{1}{2}`$ of the following recursion for modified Bessel functions of the third kind.
###### Lemma 3.2
For all $`\nu >0`$ and all nonnegative integers $`j<\nu `$ we have for $`u>0`$
$$u^{\nu +j}K_{\nu j}\left(u\right)=\underset{i=0}{\overset{j}{}}(2)^{ji}\left(\genfrac{}{}{0pt}{}{j}{i}\right)\frac{\mathrm{\Gamma }(\nu +1)}{\mathrm{\Gamma }(\nu +1(ji))}u^{\nu +i}K_{\nu +i}\left(u\right)$$
and
$$u^{2j}k_{\nu j}\left(u\right)=\underset{i=0}{\overset{j}{}}\left(1\right)^{ji}2^{2j}\left(\genfrac{}{}{0pt}{}{j}{i}\right)\frac{\mathrm{\Gamma }\left(\nu +1\right)\mathrm{\Gamma }\left(\nu +i\right)}{\mathrm{\Gamma }\left(\nu +1\left(ji\right)\right)\mathrm{\Gamma }\left(\nu j\right)}k_{\nu +i}\left(u\right).$$
Proof: The second formula follows from the first using formula (5), and the first can be proved by induction using the following recursion formula for modified Bessel functions of the third kind, cf. \[12, p. 79\]
$$K_{\nu 1}(u)=K_{\nu +1}(u)\frac{2\nu }{u}K_\nu (u).$$
We skip the details. $`\mathrm{}`$
### 3.2 Proof of Theorem 2.1
From (1) and (15) we get
$$q_n\left(au\right)=\underset{j=0}{\overset{n}{}}\alpha _j^{\left(n\right)}a^j\underset{i=0}{\overset{j}{}}\delta _i^{\left(j\right)}q_i\left(u\right)=\underset{k=0}{\overset{n}{}}c_k^{\left(n\right)}(a)q_k\left(u\right)$$
with
$`c_k^{\left(n\right)}\left(a\right)`$ $`=`$ $`{\displaystyle \underset{j=k}{\overset{n}{}}}a^j\alpha _j^{\left(n\right)}\delta _k^{\left(j\right)}`$
$`=`$ $`a^k{\displaystyle \frac{n!}{\left(2n\right)!}}{\displaystyle \frac{\left(2k\right)!}{k!}}{\displaystyle \underset{j=k,j2k+1}{\overset{n}{}}}\left(a\right)^{jk}{\displaystyle \frac{\left(2nj\right)!\left(j+1\right)}{\left(nj\right)!\left(jk\right)!\left(2k+1j\right)!}}`$
In particular $`c_n^{\left(n\right)}\left(a\right)=a^n`$ and for $`0kn1`$
$$c_k^{\left(n\right)}\left(a\right)=a^k\frac{n!}{\left(2n\right)!}\frac{\left(2k\right)!}{k!}p(a),$$
(19)
where
$$p(a)=\underset{i=0}{\overset{(nk)(k+1)}{}}(a)^i\frac{(2nki)!(k+i+1)}{(nki)!i!(k+1i)!}.$$
We clearly have
$$p(a)=\underset{r=0}{\overset{(nk)(k+1)}{}}(1)^r\frac{p^{(r)}(1)}{r!}(1a)^r$$
with
$$p^{(r)}(1)=\underset{i=r}{\overset{(nk)(k+1)}{}}(1)^i\frac{(2nki)!(k+i+1)}{(nki)!(ir)!(k+1i)!}$$
and we only consider $`0r(nk)(k+1)`$. To sum this we shift the summation by $`r`$. For simplicity we define $`T:=(nkr)(k+1r)`$ and get
$$(1)^rp^{(r)}(1)=\underset{i=0}{\overset{T}{}}(1)^i\frac{(2nkri)!(k+r+i+1)}{(nkri)!i!(k+1ri)!}.$$
We write $`k+r+1+i=(2k+2)(k+1ri)`$ and split the above sum accordingly
$`(1)^rp^{(r)}(1)`$ $`=`$ $`(2k+2){\displaystyle \underset{i=0}{\overset{T}{}}}(1)^i{\displaystyle \frac{(2nkri)!}{(nkri)!i!(k+1ri)!}}`$
$``$ $`{\displaystyle \underset{i=0}{\overset{T}{}}}(1)^i{\displaystyle \frac{(2nkri)!}{(nkri)!i!(kri)!}}.`$
Note that for nonnegative integers $`a,b,c`$ with $`b,ca`$ we have
$$\underset{i=0}{\overset{bc}{}}(1)^i\frac{(ai)!}{(bi)!(ci)!i!}=\frac{a!}{b!c!}\underset{i=0}{\overset{bc}{}}\frac{(b)_i(c)_i}{(a)_ii!}=\frac{a!}{b!c!}{}_{2}{}^{}F_{1}^{}(b,c;a;1),$$
where we use that the sum is an $`{}_{2}{}^{}F_{1}^{}`$ evaluated at 1. Its value is given by the Chu-Vandermonde formula, see , hence
$$\underset{i=0}{\overset{bc}{}}(1)^i\frac{(ai)!}{(bi)!(ci)!i!}=\frac{a!(ca)_b}{(a)_bb!c!}.$$
The two sums above are of this form and we get
$$(1)^rp^{(r)}(1)=\frac{(2nkr)!}{(nkr)!(k+1r)!}Q,$$
where
$$Q=(2k+2)\frac{(2k2n+1)_{nkr}}{(k+r2n)_{nkr}}(k+1r)\frac{(2k2n)_{nkr}}{(k+r2n)_{nkr}}$$
$$=\frac{(2k2n+1)_{nkr1}}{(k+r2n)_{nkr}}[(2k+2)(krn)(k+1r)(2k2n)]$$
$$=2r(n+1)\frac{(n+r+1k)_{nkr1}}{(n+1)_{nkr}},$$
where we used $`(a)_n=(1)^n(1an)_n`$ twice. This gives
$$(1)^rp^{(r)}(1)=2r\left(\genfrac{}{}{0pt}{}{n+1}{k+1r}\right)\frac{(2n2k1)!}{(nkr)!}$$
and finally
$$p(a)=\underset{r=0}{\overset{(nk)(k+1)}{}}(1a)^r\frac{2r}{r!}\left(\genfrac{}{}{0pt}{}{n+1}{k+1r}\right)\frac{(2n2k1)!}{(nkr)!}.$$
Note that the term corresponding to $`r=0`$ is zero. If we insert this expression for $`p(a)`$ in (19), we get the formula of Theorem 2.1. $`\mathrm{}`$
###### Remark 3.3
The evaluation above of $`(1)^rp^{(r)}(1)`$ can be done using generating functions like in . The authors want to thank Mogens Esrom Larsen for the idea to use the Chu-Vandermonde identity twice.
### 3.3 Proof of Theorem 2.2
The starting point is the following formula of Macdonald, see
$$K_\nu (z)K_\nu (X)=\frac{1}{2}_0^{\mathrm{}}\mathrm{exp}[\frac{s}{2}\frac{z^2+X^2}{2s}]K_\nu (\frac{zX}{s})\frac{ds}{s},$$
(20)
which we will use for $`\nu =n+\frac{1}{2},z=au,X=(1a)u`$. Multiplying (20) by
$$\left(\frac{2^{1\nu }}{\mathrm{\Gamma }(\nu )}\right)^2(a(1a)u^2)^\nu $$
and using (5) we find
$$k_\nu (au)k_\nu ((1a)u)=$$
$$\frac{1}{2^\nu \mathrm{\Gamma }(\nu )}_0^{\mathrm{}}\mathrm{exp}\left[\frac{s}{2}u^2\frac{a^2+(1a)^2}{2s}\right]s^{\nu 1}k_\nu \left(\frac{a(1a)u^2}{s}\right)𝑑s.$$
We now insert that with $`\nu =n+\frac{1}{2}`$ we have $`k_\nu (|u|)=e^{|u|}q_n(|u|)`$ and hence after some simplification
$$e^{|u|}q_n(a|u|)q_n((1a)|u|)=$$
$$\frac{1}{2^{n+\frac{1}{2}}\mathrm{\Gamma }(n+\frac{1}{2})}_0^{\mathrm{}}\mathrm{exp}\left[\frac{s}{2}\frac{u^2}{2s}\right]s^{n\frac{1}{2}}q_n\left(\frac{a(1a)u^2}{s}\right)𝑑s.$$
We next insert the expression (1) for $`q_n`$ under the integral sign. This gives
$$e^{|u|}q_n(a|u|)q_n((1a)|u|)=$$
$$\underset{k=0}{\overset{n}{}}\alpha _k^{(n)}(a(1a))^ku^{2k}\frac{1}{2^{n+\frac{1}{2}}\mathrm{\Gamma }(n+\frac{1}{2})}_0^{\mathrm{}}\mathrm{exp}\left[\frac{s}{2}\frac{u^2}{2s}\right]s^{nk\frac{1}{2}}𝑑s.$$
Using the following formula from \[7, 3.471(9)\]
$$_0^{\mathrm{}}x^{\nu 1}\mathrm{exp}(\frac{\beta }{x}\gamma x)𝑑x=2\left(\frac{\beta }{\gamma }\right)^{\nu /2}K_\nu \left(2\sqrt{\beta \gamma }\right)$$
(21)
and again (5) the above is equal to
$$=\underset{k=0}{\overset{n}{}}\alpha _k^{(n)}(a(1a))^k\frac{2^{nk+\frac{1}{2}}\mathrm{\Gamma }(nk+\frac{1}{2})}{2^{n+\frac{1}{2}}\mathrm{\Gamma }(n+\frac{1}{2})}e^{|u|}u^{2k}q_{nk}(|u|).$$
Finally, using Pochhammer symbols and skipping absolute values since we are now dealing with a polynomial identity, we get
$$q_n(au)q_n((1a)u)=\underset{k=0}{\overset{n}{}}\alpha _k^{(n)}(a(1a))^k\frac{(\frac{1}{2})_{nk}}{2^k(\frac{1}{2})_n}u^{2k}q_{nk}(u).$$
(22)
Using the expression for $`u^{2k}q_{nk}(u)`$ from Lemma 3.1 and the expression for $`\alpha _k^{(n)}`$ in (22) we then get
$`q_n(au)q_n((1a)u)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{k}\right)(\frac{1}{2})_{nk}}{\left(\genfrac{}{}{0pt}{}{2n}{k}\right)(\frac{1}{2})_nk!}}(a(1a))^k{\displaystyle \underset{i=0}{\overset{k}{}}}\gamma _i^{(n,k)}q_{n+i}(u)`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}q_{n+i}(u){\displaystyle \underset{k=i}{\overset{n}{}}}(a(1a))^k{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{k}\right)(\frac{1}{2})_{nk}}{\left(\genfrac{}{}{0pt}{}{2n}{k}\right)(\frac{1}{2})_nk!}}\gamma _i^{(n,k)}`$
hence
$$q_n(au)q_n((1a)u)=\underset{i=0}{\overset{n}{}}\beta _i^{(n)}(a)q_{n+i}(u)$$
with
$`\beta _i^{(n)}(a)`$ $`=`$ $`{\displaystyle \underset{k=i}{\overset{n}{}}}(a(1a))^k{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{k}\right)(\frac{1}{2})_{nk}}{\left(\genfrac{}{}{0pt}{}{2n}{k}\right)(\frac{1}{2})_nk!}}2^{2k}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{i}}\right)(nk+{\displaystyle \frac{1}{2}})_{k+i}(n{\displaystyle \frac{1}{2}})_{ki}`$
$`=`$ $`(a(1a))^i{\displaystyle \underset{l=0}{\overset{ni}{}}}(a(1a))^l{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{i+l}\right)(\frac{1}{2})_{nil}}{\left(\genfrac{}{}{0pt}{}{2n}{i+l}\right)(\frac{1}{2})_n(i+l)!}}2^{2i+2l}\left({\displaystyle \genfrac{}{}{0pt}{}{i+l}{i}}\right)`$
$`\times `$ $`(nil+{\displaystyle \frac{1}{2}})_{l+2i}(n{\displaystyle \frac{1}{2}})_l`$
Collecting
$$(\frac{1}{2})_{nil}(nil+\frac{1}{2})_{l+2i}=(\frac{1}{2})_{n+i}$$
we get
$$\beta _i^{(n)}(a)=(a(1a))^i\underset{l=0}{\overset{ni}{}}(a(1a))^l\frac{\left(\genfrac{}{}{0pt}{}{n}{i+l}\right)(\frac{1}{2})_{n+i}}{\left(\genfrac{}{}{0pt}{}{2n}{i+l}\right)(\frac{1}{2})_n}\frac{2^{2i+2l}}{i!l!}(n\frac{1}{2})_l$$
$$=(a(1a))^i\frac{n!(\frac{1}{2})_{n+i}2^{2i}}{(2n)!(\frac{1}{2})_ni!}\underset{l=0}{\overset{ni}{}}(4a(1a))^l\frac{(2nil)!(n\frac{1}{2})_l}{(nil)!l!}$$
$$=(a(1a))^i\left(\frac{n!}{(2n)!}\right)^2\frac{(2n+2i)!}{(n+i)!i!}\underset{l=0}{\overset{ni}{}}\left(1(2a1)^2\right)^l\frac{(2nil)!(n\frac{1}{2})_l}{(nil)!l!}.$$
Expanding $`(1(2a1)^2)^l`$ using the binomial formula and interchanging the sums we get
$$\underset{l=0}{\overset{ni}{}}(1(2a1)^2)^l\frac{(2nil)!(n\frac{1}{2})_l}{(nil)!l!}$$
$$=\underset{j=0}{\overset{ni}{}}(1)^j\frac{(2a1)^{2j}}{j!}\underset{l=j}{\overset{ni}{}}\frac{(2nil)!(n\frac{1}{2})_l}{(nil)!(lj)!}.$$
We claim that
$`S:`$ $`=`$ $`{\displaystyle \underset{l=j}{\overset{ni}{}}}{\displaystyle \frac{(2nil)!(n\frac{1}{2})_l}{(nil)!(lj)!}}`$ (23)
$`=`$ $`{\displaystyle \frac{(2n2i)!}{(ni)!}}2^{2i2n}(1)^jj!i!\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{2j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{i}}\right).`$
and we have obtained the final formula
$`\beta _i^{(n)}(a)=(4a(1a))^i\left({\displaystyle \frac{n!}{(2n)!}}\right)^22^{2n}{\displaystyle \frac{(2n2i)!(2n+2i)!}{(ni)!(n+i)!}}`$
$`\times {\displaystyle \underset{j=0}{\overset{ni}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{2j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{i}}\right)(2a1)^{2j}.`$
We will see that (23) is a Chu-Vandermonde formula. In fact, shifting the summation index putting $`l=j+m`$ we get
$$S=(n\frac{1}{2})_j\underset{m=0}{\overset{nij}{}}\frac{(2nijm)!(n+j\frac{1}{2})_m}{m!(nijm)!},$$
and calling the general term in this sum $`c_m`$ we get
$$\frac{c_{m+1}}{c_m}=\frac{(m+i+jn)(m+jn\frac{1}{2})}{(m+1)(m+i+j2n)},$$
which shows that the sum is a $`{}_{2}{}^{}F_{1}^{}`$. We have
$$S=(n\frac{1}{2})_jc_0\times {}_{2}{}^{}F_{1}^{}((nij),jn\frac{1}{2};i+j2n;1),$$
and using the Chu-Vandermonde identity, cf. :
$${}_{2}{}^{}F_{1}^{}(n,a;c;1)=\frac{(ca)_n}{(c)_n},$$
we get
$`S`$ $`=`$ $`(n{\displaystyle \frac{1}{2}})_j{\displaystyle \frac{(2nij)!}{(nij)!}}{\displaystyle \frac{(n+i+\frac{1}{2})_{nij}}{(i+j2n)_{nij}}}`$
$`=`$ $`(n{\displaystyle \frac{1}{2}})_j{\displaystyle \frac{(2nij)!}{(nij)!}}{\displaystyle \frac{(j+\frac{1}{2})_{nij}}{(n+1)_{nij}}},`$
where we used $`(a)_n=(1)^n(1an)_n`$ twice. We can now simplify to get
$$S=(n\frac{1}{2})_j\frac{n!}{(nij)!}\frac{(\frac{1}{2})_{ni}}{(\frac{1}{2})_j},$$
and applying the formula
$$(\frac{1}{2})_p=\frac{(2p)!}{p!2^{2p}}$$
twice we get
$$S=\frac{(2n2i)!}{(ni)!}2^{2i2n}(n\frac{1}{2})_j\frac{n!}{(nij)!}\frac{j!2^{2j}}{(2j)!}.$$
Now we can write
$$(n\frac{1}{2})_j2^{2j}=(1)^j\frac{(2n+1)!(nj)!}{n!(2n2j+1)!},$$
and hence
$$S=\frac{(2n2i)!}{(ni)!}2^{2i2n}(1)^jj!i!\left(\genfrac{}{}{0pt}{}{2n+1}{2j}\right)\left(\genfrac{}{}{0pt}{}{nj}{i}\right).$$
$`\mathrm{}`$
### 3.4 Proof of Theorem 2.4
For $`n,m0`$ and $`a`$, we can write
$$q_n\left(au\right)q_m\left(\left(1a\right)u\right)=\underset{k=0}{\overset{m+n}{}}\beta _k^{(n,m)}\left(a\right)q_k\left(u\right)$$
(24)
for some uniquely determined coefficients since the left-hand side is a polynomial in $`u`$ of degree $`n+m`$. Clearly $`\beta _k^{(n,m)}(a)`$ is a polynomial in $`a`$ satisfying
$$\beta _k^{(n,m)}\left(a\right)=\beta _k^{(m,n)}\left(1a\right).$$
(25)
We shall prove that $`\beta _k^{(n,m)}(a)0`$ for $`0a1`$ and that $`\beta _k^{(n,m)}\left(a\right)=0`$ if $`k<nm`$, which will be a consequence of the following recursion formula.
###### Lemma 3.4
For $`n,m1`$, we have
$$\frac{1}{2k+1}\beta _{k+1}^{(n,m)}\left(a\right)=\frac{a^2}{2n1}\beta _k^{(n1,m)}\left(a\right)+\frac{\left(1a\right)^2}{2m1}\beta _k^{(n,m1)}\left(a\right),$$
(26)
where $`k=0,1,\mathrm{},m+n1.`$
Furthermore $`\beta _0^{(n,m)}(a)=0.`$
Proof: Differentiating (24) with respect to $`u`$ gives
$$aq_n^{}\left(au\right)q_m\left(\left(1a\right)u\right)+\left(1a\right)q_n\left(au\right)q_m^{}\left(\left(1a\right)u\right)=\underset{k=1}{\overset{m+n}{}}\beta _k^{(n,m)}\left(a\right)q_k^{}\left(u\right)$$
and using the formula (14) we find
$`a\left(q_n\left(au\right){\displaystyle \frac{au}{2n1}}q_{n1}\left(au\right)\right)q_m\left(\left(1a\right)u\right)`$
$`+`$ $`\left(1a\right)q_n\left(au\right)\left(q_m\left(\left(1a\right)u\right){\displaystyle \frac{\left(1a\right)u}{2m1}}q_{m1}\left(\left(1a\right)u\right)\right)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{m+n}{}}}\beta _k^{(n,m)}\left(a\right)\left(q_k\left(u\right){\displaystyle \frac{u}{2k1}}q_{k1}\left(u\right)\right)`$
and using (24) once more we get
$`{\displaystyle \frac{a^2u}{2n1}}q_{n1}\left(au\right)q_m\left(\left(1a\right)u\right){\displaystyle \frac{\left(1a\right)^2u}{2m1}}q_n\left(au\right)q_{m1}\left(\left(1a\right)u\right)`$
$`=`$ $`\beta _0^{(n,m)}(a)u{\displaystyle \underset{k=0}{\overset{n+m1}{}}}\beta _{k+1}^{(n,m)}\left(a\right)(2k+1)^1q_k(u).`$
For $`u=0`$ this gives $`\beta _0^{(n,m)}(a)=0`$ and dividing by $`u`$ and equating the coefficients of $`q_k\left(u\right)`$, we get the desired formula. $`\mathrm{}`$
Now the proof of Theorem 2.4 is easy by induction in $`k`$ and by the symmetry formula (25), we can assume $`nm.`$ Let $`0a1`$. We prove that $`\beta _k^{(n,m)}(a)0`$ for $`kn+m`$ and that it is zero for $`k<m`$ (under the assumption $`nm`$). This is true for $`k=0`$ by Lemma 3.4 when $`m1`$, and for $`m=0`$ it follows by Theorem 2.1. Assume now that the results hold for $`k=k_0`$ and assume $`k_0+1n+m`$. The nonnegativity for $`k=k_0+1`$ now follows by Lemma 3.4, and likewise if $`k_0+1<mn`$ the coefficient is 0 since $`k_0<(n1)(m1)`$. $`\mathrm{}`$
### 3.5 Proof of Theorem 2.6
By Theorem 2.4 the result holds for $`k=2`$. Assuming it holds for $`k12`$ we have
$$q_{n_1}(a_1u)\mathrm{}q_{n_{k1}}(a_{k1}u)=\underset{j=l^{}}{\overset{L^{}}{}}\gamma _jq_j((1a_k)u),u$$
(27)
with $`l^{}=\mathrm{min}(n_1,\mathrm{},n_{k1}),L^{}=n_1+\mathrm{}+n_{k1}`$ and $`\gamma _j0`$ because we can write
$$a_ju=\frac{a_j}{1a_k}(1a_k)u,j=1,\mathrm{},k1.$$
If we multiply (18) with $`q_{n_k}(a_ku)`$ we get
$$\underset{j=l^{}}{\overset{L^{}}{}}\gamma _jq_{n_k}(a_ku)q_j((1a_k)u)=\underset{j=l^{}}{\overset{L^{}}{}}\gamma _j\underset{i=n_kj}{\overset{n_k+j}{}}\beta _i^{(n_k,j)}(a_k)q_i(u),$$
and the assertion follows. $`\mathrm{}`$
## 4 Inverse Gamma distribution
Grosswald proved that the Student-t distribution is infinitely divisible. This is a consequence of the infinite divisibility of the inverse Gamma distribution because of subordination. It was proved later that the inverse Gamma distribution is a generalized Gamma convolution in the sense of Thorin, which is stronger than self-decomposability and in particular stronger than infinite divisibility.
The following density on the half-line is an inverse Gamma density with scale parameter $`\frac{1}{4}`$ and shape parameter $`\nu >0`$:
$$C_\nu \mathrm{exp}(\frac{1}{4t})t^{\nu 1},C_\nu =\frac{1}{2^{2\nu }\mathrm{\Gamma }(\nu )}.$$
(28)
Let the corresponding probability measure be denoted $`\stackrel{~}{\gamma }_\nu `$ and let
$$g_t(x)=\frac{1}{\sqrt{4\pi t}}\mathrm{exp}(\frac{x^2}{4t}),t>0,x$$
denote the Gaussian semigroup of normal densities (in the normalization of ). Then
$$f_\nu (x)=_0^{\mathrm{}}g_t(x)𝑑\stackrel{~}{\gamma }_\nu (t)$$
(29)
is the Student-t density (3) with $`2\nu `$ degrees of freedom. The corresponding probability measure is denoted $`\sigma _\nu `$. This formula says that $`\sigma _\nu `$ is subordinated to the Gaussian semigroup by an inverse Gamma distribution, and it implies the infinite divisibility of Student-t from the infinite divisibility of inverse Gamma. Since the Laplace transformation is one-to-one, it is clear that if two probabilities $`\gamma _1,\gamma _2`$ on $`]0,\mathrm{}[`$ lead to the same subordinated density
$$_0^{\mathrm{}}g_t(x)𝑑\gamma _1(t)=_0^{\mathrm{}}g_t(x)𝑑\gamma _2(t),x,$$
then $`\gamma _1=\gamma _2`$.
If we denote $`\tau _a(x)=ax`$, the distribution $`\tau _a(\sigma _{n+{\scriptscriptstyle \frac{1}{2}}})\tau _{1a}(\sigma _{m+{\scriptscriptstyle \frac{1}{2}}})`$ is given in (10). However note that $`\tau _a(g_t(x)dx)=g_{ta^2}(x)dx`$ so
$$\tau _a(\sigma _\nu )=_0^{\mathrm{}}g_{ta^2}(x)𝑑\stackrel{~}{\gamma }_\nu (t)𝑑x,$$
(30)
hence
$$\tau _a(\sigma _{\nu _1})\tau _{1a}(\sigma _{\nu _2})=_0^{\mathrm{}}_0^{\mathrm{}}(g_{ta^2}dx)(g_{s(1a)^2}dx)𝑑\stackrel{~}{\gamma }_{\nu _1}(t)𝑑\stackrel{~}{\gamma }_{\nu _2}(s)$$
$$=_0^{\mathrm{}}_0^{\mathrm{}}(g_{ta^2+s(1a)^2}dx)𝑑\stackrel{~}{\gamma }_{\nu _1}(t)𝑑\stackrel{~}{\gamma }_{\nu _2}(s)$$
$$=_0^{\mathrm{}}g_u(x)𝑑\tau _{a^2}(\stackrel{~}{\gamma }_{\nu _1})\tau _{(1a)^2}(\stackrel{~}{\gamma }_{\nu _2})(u)𝑑x.$$
Therefore, using (30) we see that for $`\nu _1=n+\frac{1}{2},\nu _2=m+\frac{1}{2}`$ with $`n,m=0,1,\mathrm{}`$ the formula (10) rewritten as
$$\tau _a(\sigma _{n+{\scriptscriptstyle \frac{1}{2}}})\tau _{1a}(\sigma _{m+{\scriptscriptstyle \frac{1}{2}}})=\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}(a)\sigma _{k+{\scriptscriptstyle \frac{1}{2}}}$$
is equivalent to
$$\tau _{a^2}(\stackrel{~}{\gamma }_{n+{\scriptscriptstyle \frac{1}{2}}})\tau _{(1a)^2}(\stackrel{~}{\gamma }_{m+{\scriptscriptstyle \frac{1}{2}}})=\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}(a)\stackrel{~}{\gamma }_{k+{\scriptscriptstyle \frac{1}{2}}}.$$
(31)
This shows that Theorem 2.4 is equivalent to the following result about inverse Gamma distributions:
The distribution of $`a^2Z_n+(1a)^2Z_m`$, where $`Z_n,Z_m`$ are independent inverse Gamma random variables with distribution (28) for $`\nu =n+\frac{1}{2},m+\frac{1}{2}`$ respectively, has a density which is a convex combination of inverse Gamma densities.
This result can be extended to the multivariate Student-t distributions as follows. A rotation invariant $`N`$variate Student-t probability density is given for $`𝕩=(x_1,\mathrm{},x_N)^N`$ by
$$f_{N,\nu }\left(𝕩\right)=A_{N,\nu }\left(1+|𝕩|^2\right)^{\nu \frac{N}{2}},A_{N,\nu }=\frac{\mathrm{\Gamma }\left(\nu +\frac{N}{2}\right)}{\mathrm{\Gamma }\left(\nu \right)(\mathrm{\Gamma }(\frac{1}{2}))^N},$$
where
$$𝕩,𝕪=\underset{i=1}{\overset{N}{}}x_iy_i,|𝕩|=\left(𝕩,𝕩\right)^{{\scriptscriptstyle \frac{1}{2}}},𝕩,𝕪^N.$$
It is easy to verify that $`f_{N,\nu }(𝕩)`$ is subordinated to the N-variate Gaussian semigroup
$$g_{N,t}(𝕩)=(4\pi t)^{{\scriptscriptstyle \frac{N}{2}}}\mathrm{exp}\left(\frac{|𝕩|^2}{4t}\right),t>0,𝕩^N$$
by the inverse Gamma density (28), i.e.
$$f_{N,\nu }(𝕩)=_0^{\mathrm{}}g_{N,t}(𝕩)𝑑\stackrel{~}{\gamma }_\nu (t).$$
(32)
Therefore the characteristic function is given by
$$_^Ne^{i𝕩,𝕪}f_{N,\nu }\left(𝕩\right)𝑑𝕩=k_\nu (|𝕪|)$$
(33)
generalizing (4). In fact
$$_^Ne^{i𝕩,𝕪}f_{N,\nu }\left(𝕩\right)𝑑𝕩=_0^{\mathrm{}}\left(_^Ne^{i𝕩,𝕪}g_{N,t}\left(𝕩\right)𝑑𝕩\right)𝑑\stackrel{~}{\gamma }_\nu (t)$$
$$=_0^{\mathrm{}}e^{t|𝕪|^2}𝑑\stackrel{~}{\gamma }_\nu (t)$$
and the result follows by (21).
As a conclusion, the Theorems 2.1, 2.2 and 2.4 apply in the multivariate case. For example, an equivalent form of (10) writes as follows: with $`0<a<1`$,
$$\frac{1}{a^N}f_{N,n+\frac{1}{2}}\left(a^1𝕩\right)\frac{1}{\left(1a\right)^N}f_{N,m+\frac{1}{2}}\left(\left(1a\right)^1𝕩\right)=\underset{k=nm}{\overset{n+m}{}}\beta _k^{(n,m)}\left(a\right)f_{N,k+\frac{1}{2}}\left(𝕩\right).$$
Department of Mathematics, University of Copenhagen, Universitetsparken 5, DK-2100, Copenhagen, Denmark.
Email: berg@math.ku.dk
Laboratoire d’Informatique IGM, UMR 8049, Université de Marne-la-Vallée, 5 Bd. Descartes, F-77454 Marne-la-Vallée Cedex, France.
Email: vignat@univ-mlv.fr
|
warning/0506/gr-qc0506067.html
|
ar5iv
|
text
|
# A group field theory for 3d quantum gravity coupled to a scalar field
## I Introduction
Spin foam models daniele ; ale represent a purely combinatorial and algebraic implementation of the sum-over-histories approach to quantum gravity, in any signature and spacetime dimension, with an abstract 2-complex playing the role of a discrete spacetime, and algebraic data from the representation theory of the Lorentz group playing the role of geometric data assigned to it. Indeed, the first model of quantum gravity to be ever proposed, the Ponzano-Regge model, was a spin foam model for Euclidean quantum gravity without cosmological constant PR0 . This approach has recently been developed to a great extent in the 3-dimensional case. It is now clear that it provides a full quantisation of pure gravity PR1 , whose relation with the one obtained by other approaches is well understood PR2 ; lec . Moreover, matter can be consistently included in the picture PR1 ; barrettfeynman , providing a link between spin foam models and effective field theory PR3 living on a non-commutative geometry. This picture allows us to naturally address the semi-classical limit of spin foam models and shows that quantum gravity in dimension $`3`$ effectively follows the principle of the so-called deformed (or doubly) special relativity dsr .
The group field theory formalism laurentgft represents a generalisation of matrix models of 2-dimensional quantum gravity mmreview . It is a universal structure lying behind any spin foam model for quantum gravity DP-F-K-R ; carlomike , providing a third quantisation point of view on gravity laurentgft and allowing us to sum over pure quantum gravity amplitudes associated with different topologies NPsum . In this picture, spin foams, and thus spacetime itself, appear as (higher-dimensional analogues of) Feynman diagrams of a field theory defined on a group manifold and spin foam amplitudes are simply the Feynman amplitudes weighting the different graphs in the perturbative expansion of the quantum field theory. On the other hand, we can construct a non commutative field theory whose Feynman diagram amplitudes reproduce the coupling of matter fields to 3d quantum gravity for a trivial topology of spacetime PR3 . Remarkably, the momenta of the fields are labelled also by group elements. Moreover, in three dimensions there is a duality between matter and geometry, and the insertion of matter can be understood as the insertion of a topological defect charged under the Poincaré group PR1 . This suggests that one should be able to treat the third quantisation of gravity and the second quantisation of matter fields in one stroke (see mikovic for an early attempt). The purpose of this paper is to study how the coupling of matter to quantum gravity is realised in the group field theory, and whether it is possible to write down a group field theory for gravity and particles that reproduces the amplitudes derived in PR1 coupling quantum matter to quantum geometry.
This is what we achieve in the present work. The way the correct amplitudes are generated as Feynman amplitudes of the group field theory is highly non-trivial. It requires an extension of the usual group field theory (gft) formalism to a higher number of field variables, and produces an interesting intertwining of gravity and matter degrees of freedom, as we are going to discuss in the following. The formalism we present is still based on the classical $`SU(2)`$, it bears strong similarity with a recent work of Krasnov Kirill who considered gft based on the quantum group $`DSU(2)`$. This should not be a surprise since it is well understood that the particle spin foam amplitudes have a quantum group structure hidden in them PR1 ; PR2 . The gft model we propose here is, however, very different from the ones considered in Krasnov’s work since our Feynman graphs reproduce explicitly the insertion of particles coupled to gravity.
## II The Ponzano-Regge spin foam model coupled to point particles
The general form of Feynman graph amplitudes for spinning particles coupled to 3 dimensional quantum gravity -eg the Ponzano-Regge model- has been written in PR1 . In this paper we focus on the case of spinless particles and we recall in this section the definition of these amplitudes before deriving them from a gft. We start from a triangulation $`\mathrm{\Delta }`$ of our spacetime $`M`$ and consider also the dual $`\mathrm{\Delta }^{}`$: dual vertices, edges and faces correspond respectively to tetrahedra, faces and edges of $`\mathrm{\Delta }`$. We choose our Feynman graph, $`\mathrm{}`$, to be embedded in the triangulation $`\mathrm{\Delta }`$ such that edges of $`\mathrm{}`$ are edges of the triangulation. Each edge of $`\mathrm{}`$ is labelled by an angle $`\theta [0,\pi ]`$
$$\theta =\kappa m,\kappa =4\pi G_N,$$
where $`G_N`$ is Newton’s constant, $`\kappa `$ is the inverse Planck mass and $`m`$ is the mass of the particle. To each angle $`\theta `$ we associate an element of the Cartan subgroup $`H`$ of $`\mathrm{SU}(2)`$
$$h_\theta =\left(\begin{array}{cc}e^{i\theta }& 0\\ 0& e^{i\theta }\end{array}\right),$$
which corresponds to a rotation of angle $`2\theta `$ around a given axis. Given a group $`G`$, here $`\mathrm{SU}(2)`$, we assign group elements $`g_e^{}`$ to all dual edges $`e^{}`$ of the triangulation. We constrain the holonomies around dual faces $`f^{}e`$ to be flat if there is no particle and we constrain it to be in the conjugacy class $`\theta `$ if $`e`$ is an edge of $`\mathrm{}`$. More precisely, let us denote by $`G_e(=G_f^{})`$ the product of the group elements around a dual face (or plaquette) $`f^{}e`$:
$$G_e=G_f^{}=\underset{e^{}f^{}}{}g_e^{}^{ϵ_f^{}(e^{})},$$
where $`ϵ_f^{}(e^{})=\pm 1`$ records the orientation of the edge $`e^{}`$ in the boundary of the (dual) face $`f^{}`$. The amplitude is given by
$$𝒵_M(\mathrm{}_\theta )=\mathrm{\Delta }(\theta )^{|E_{\mathrm{}}|}\underset{e^{}}{}dg_e^{}\underset{e\mathrm{}}{}du_e\underset{e\mathrm{}}{}\delta (G_e)\underset{e\mathrm{}}{}\delta (G_eu_eh_\theta u_e^1),$$
(1)
where $`dg`$ is the normalised Haar measure and $`\delta (g)`$ the corresponding delta function on $`G`$, $`\mathrm{\Delta }(\theta )\mathrm{sin}(\theta )`$ and $`|E_{\mathrm{}}|`$ is the number of edges in the particle graph $`\mathrm{}`$. We see that two types of group elements arise in the construction of this amplitude, the $`g_e^{}`$ describe pure gravity excitations and the $`u_e`$ variables are associated with the particle degrees of freedom. They arise because the insertion of a particle locally breaks the Lorentz and translational symmetries of the gravity model and the former gauge transformation becomes dynamical at the location of the particle PR1 . The $`u_e`$ are then interpreted as giving the direction of the particle momenta propagating along the edge $`e`$. The fact that we are talking about spinless particles implies that $`u_e`$ should not be thought of as an element of $`G`$ but as an element of $`G/H`$ (with $`H=U(1)`$), that is momentum space. The insertion of spinning particles can be achieved by taking into account a non trivial dependence under the $`H`$ part of $`u_e`$. The main lesson which we learn from this amplitude, and which gives the key idea leading to a gft construction of such an amplitude, is the fact that we need both $`G`$ variables describing the gravity excitation and $`G/H`$ variables describing the propagation of particles. We can expand the $`\delta `$ functions in terms of characters
$$\delta (g)=\underset{j}{}d_j\chi _j(g),$$
$`d_j=2j+1`$ being the dimension of the spin $`j`$ representation and perform the integration over $`g_e^{}`$, $`u_e`$ in order to obtain a state sum model
$$𝒵_M(\mathrm{}_\theta )=\mathrm{\Delta }(\theta )^{|E_{\mathrm{}}|}\underset{\{j_e\}}{}\underset{e\mathrm{}}{}d_{j_e}\underset{e\mathrm{}}{}\chi _{j_e}(h_{\theta _e})\underset{t}{}\left\{\begin{array}{ccc}j_{e_{t_1}}& j_{e_{t_2}}& j_{e_{t_3}}\\ j_{e_{t_4}}& j_{e_{t_5}}& j_{e_{t_6}}\end{array}\right\},$$
(2)
where the summation is over all edges of $`\mathrm{\Delta }`$ and the product of normalised 6j symbols is over all tetrahedra $`t`$. For each tetrahedron, the admissible triples of edges, e.g. $`(j_{e_{t_1}},j_{e_{t_2}},j_{e_{t_3}})`$, corresponds to faces of this tetrahedra. Boulatov Boul was the first to show that the amplitude (1) can be obtained as a Feynman graph evaluation of a group field theory. It is important to note however that this amplitude is generically divergent. It is now understood diffeo that this divergence is due to a translational gauge symmetry (equivalent on-shell to diffeomorphism symmetry) acting on the Ponzano-Regge model. This symmetry should be gauge-fixed in order to obtain well defined and triangulation independent amplitudes. The gauge-fixing is easily implemented by choosing a maximal tree<sup>4</sup><sup>4</sup>4A connected set of edges of $`\mathrm{\Delta }\backslash \mathrm{}`$ passing through every vertex of $`\mathrm{\Delta }\backslash \mathrm{}`$ $`T`$ of $`\mathrm{\Delta }\backslash \mathrm{}`$. The gauge-fixed amplitude can be obtained from (1) by replacing the product $`_e\mathrm{}\delta (G_e)`$ by a product over delta functions with $`eT\mathrm{}`$. In terms of the state sum model (1) the gauge-fixing inserts in the summation a factor $`_{eT}\delta _{j_e,0}`$ which eliminates the sum over $`j_e`$, $`eT`$. The overall amplitude does not depend on the choice of $`T`$.
## III A gft model for 3d quantum gravity coupled to scalar matter
### III.1 Action and Feynman rules
We shall now define a field theory on a group manifold, whose Feynman expansion gives the above modified Ponzano-Regge model. We consider a generic real field $`\varphi `$, over the Cartesian product of six copies of $`\mathrm{SU}(2)`$
$$\varphi (g_1,g_2,g_3;u_1,u_2,u_3):\underset{6}{\underset{}{\mathrm{SU}(2)\times \mathrm{}\times \mathrm{SU}(2)}}.$$
(3)
This is the basic object of the theory and, just as in the other group field formulations of spin foam models, it has the interpretation of a ‘3rd quantised’ chunk of quantum geometry DP-F-K-R ; laurentgft . However, in this extended formulation based on a 6-argument field, this chunk of quantum geometry carries also additional degrees of freedom, labelled by the extra $`u`$ variables, that acquire the physical meaning of particle degrees of freedom (more precisely particle momenta) when a mass parameter is inserted in a suitable way, as we are going to show in the following. Let us now list the symmetries that this field is required to satisfy.
* We require that $`\varphi `$ is invariant under (even) elements $`\sigma `$ of the permutation group of three elements $`S_3`$, acting on pairs of field variables $`(g_i,u_i)`$:
$$\varphi (g_1,g_2,g_3;u_1,u_2,u_3)=\varphi (g_{\sigma (1)},g_{\sigma (2)},g_{\sigma (3)};u_{\sigma (1)},u_{\sigma (2)},u_{\sigma (3)}).$$
(4)
If we require the field to be invariant under even permutations of the three pairs of arguments, then this is equivalent to dealing with a complex field instead, with the odd permutations mapping the field to its complex conjugate DP-P , and the Feynman amplitudes produced by the corresponding group field theory are in one-to-one correspondence with orientable 2-complexes, as explained in DP-F-K-R ; DP-P . We can more generally require the field to transform under an arbitrary representation (not necessarily reducible) of $`S_3`$. This will affect the type of 2-complexes generated by the perturbative expansion of the theory. We stress, however, that this would not imply any change for the amplitudes of the Feynman diagrams.
* Furthermore, we pick a $`\mathrm{U}(1)`$ subgroup $`H`$, of $`\mathrm{SU}(2)`$, with the interpretation of the invariance subgroup for the particle momenta, and project three of the arguments into $`\mathrm{SU}(2)/\mathrm{U}(1)`$ equivalence classes
$$P_b\varphi (g_1,g_2,g_3;u_1,u_2,u_3)_{H^3}\underset{i=1}{\overset{3}{}}db_i\varphi (g_1,g_2,g_3;u_1b_1,u_2b_2,u_3b_3),$$
(5)
so that the field becomes in fact a function over three copies each of $`\mathrm{SU}(2)`$ and $`\mathrm{SU}(2)/\mathrm{U}(1)`$.
* Finally, we project the first half of the field, i.e. the part dependent on the first three arguments, into its $`SU(2)`$ invariant part, by imposing invariance under simultaneous right action of $`\mathrm{SU}(2)`$ on the first three arguments:
$$P_\alpha \varphi (g_1,g_2,g_3;u_1,u_2,u_3)_{\mathrm{SU}(2)}𝑑\alpha \varphi (g_1\alpha ,g_2\alpha ,g_3\alpha ;u_1,u_2,u_3).$$
(6)
This last symmetry has a geometric interpretation, as in the usual Boulatov model. It imposes the closure of the triangle of which the field $`\varphi `$ represents the 2nd quantisation, by constraining the spin variables dual to the $`g_i`$ variables associated to its three edges.
Given such a field, we can write down a Boulatov-like action, with the extra $`u`$ variables simply mimicking the relations among the gravity degrees of freedom $`g`$:
$$\begin{array}{cc}\hfill S[\varphi ]=& \frac{1}{2}\underset{i=1}{\overset{3}{}}dg_idu_i[P_\alpha P_b\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\overline{\alpha }}P_{\overline{b}}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)]\hfill \\ & +\frac{\lambda }{4!}\underset{i=1}{\overset{6}{}}dg_idu_i[P_{\alpha _1}P_{b_1}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\alpha _2}P_{b_2}\varphi (g_4,g_5,g_3;u_4,u_5,u_3)]\hfill \\ & \times [P_{\alpha _3}P_{b_3}\varphi (g_4,g_2,g_6;u_4,u_2,u_6)][P_{\alpha _4}P_{b_4}\varphi (g_1,g_5,g_6;u_1,u_5,u_6)].\hfill \end{array}$$
(7)
As we will see the Feynman amplitudes obtained from this model are proportional to those obtained by the Boulatov model, i.e. the usual Ponzano-Regge spin foam amplitudes describing pure 3d Riemannian quantum gravity. This shows that the $`u`$ variables are completely redundant at this stage and do not have any real physical meaning. They are going to acquire it soon, however. Now we introduce a mass parameter in the theory, turning this redundant description of pure quantum gravity into a model for gravity coupled to scalar matter. We define a mass insertion operator $`P_\theta `$, acting on the field $`\varphi `$ as follows:
$$P_\theta \varphi (g_1,g_2,g_3;u_1,u_2,u_3)\varphi (u_1h_\theta u_1^1g_1,g_2,g_3;u_1ϵ,u_2,u_3),$$
(8)
where $`h_\theta =exp(\theta J_0)H`$; $`\theta `$ is half the deficit angle created by the presence of a mass $`m`$, $`\theta =4\pi Gm`$, $`\frac{1}{4\pi G}`$ being the Planck mass; $`J_0`$ is the generator of the $`\mathrm{U}(1)`$ subgroup $`H`$, the same subgroup under which the $`u_i`$ variables of the field are invariant and $`ϵ`$ is the non-trivial Weyl group element, given in the fundamental representation by:
$$ϵ=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
We define then a new group field theory model, representing 3d Riemannian quantum gravity coupled to scalar matter, whose dynamics are given by the action:
$$\begin{array}{cc}\hfill S[\varphi ]=& \frac{1}{2}\underset{i=1}{\overset{3}{}}dg_idu_i([P_\alpha P_b\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\overline{\alpha }}P_{\overline{b}}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)]\hfill \\ & 𝔞[P_\alpha P_b\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_\theta P_{\overline{\alpha }}P_{\overline{b}}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)])\hfill \\ & +\frac{\lambda }{4!}\underset{i=1}{\overset{6}{}}dg_idu_i[P_{\alpha _1}P_{b_1}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\alpha _2}P_{b_2}\varphi (g_4,g_5,g_3;u_4,u_5,u_3)]\hfill \\ & \times [P_{\alpha _3}P_{b_3}\varphi (g_4,g_2,g_6;u_4,u_2,u_6)][P_{\alpha _4}P_{b_4}\varphi (g_1,g_5,g_6;u_1,u_5,u_6)],\hfill \end{array}$$
(9)
where $`𝔞`$ is a free parameter and all integrals are with respect to the normalised Haar measure. Let us write the Feynman rules of the theory in coordinate space. We have to identify kinetic and vertex operators. For this purpose, we write the action as
$$\begin{array}{cc}\hfill S[\varphi ]=& \frac{1}{2}\underset{i=1}{\overset{3}{}}dg_idu_i\underset{j=1}{\overset{3}{}}d\overline{g}_jd\overline{u}_j\varphi (g_i,u_i)𝒦(g_i,\overline{g}_j,u_i,\overline{u}_j)\varphi (\overline{g}_j,\overline{u}_j)\hfill \\ & +\underset{i,j}{}dg_{ij}du_{ij}𝒱(g_{ij},u_{ij})\varphi (g_{1j},u_{2j})\varphi (g_{2j},u_{2j})\varphi (g_{3j},u_{3j})\varphi (g_{4j},u_{4j}).\hfill \end{array}$$
(10)
where in this integral, $`ij`$, and $`\varphi (g_{1j},u_{1j})=\varphi (g_{11},g_{12},g_{13};u_{11},u_{12},u_{13})`$, and so forth. The kinetic and vertex operators are
$$𝒦(g_i,\overline{g}_j,u_i,\overline{u}_j)=\underset{i=1}{\overset{3}{}}db_i\delta (g_i\overline{g}_i^1)\delta (u_ib_i\overline{u}_i^1)𝔞\underset{i=1}{\overset{3}{}}db_i\delta (u_1h_\theta u_1^1g_1\overline{g}_1^1)\delta (u_1b_1ϵ\overline{u}_1^1)\underset{j=2}{\overset{3}{}}\delta (g_j\overline{g}_j^1)\delta (u_jb_j\overline{u}_j^1),$$
(11)
$$𝒱(g_{ij},u_{ij})=\frac{\lambda }{4!}\underset{i=1}{\overset{4}{}}d\alpha _i\underset{j>i}{}db_{ij}\delta (\alpha _jg_{ji}^1g_{ij}\alpha _i^1)\delta (b_{ji}u_{ji}^1u_{ij}b_{ij}^1).$$
(12)
We have purposefully discarded the $`\alpha `$ variables in the $`\delta `$-functions of the kinetic term. This does not change, up to an overall factor, the computation of the amplitudes since $`P_\theta `$ commutes with $`P_\alpha `$. We have also sidelined the sum over permutations, for ease of notation. Care should be taken in inverting the kinetic term, however, on the subspace of symmetric fields only. On this subspace, the kinetic term is indeed diagonal and we can proceed as follows. We now define the operators $`I`$ and $`K_\theta `$:
$$𝒦(g_i,\overline{g}_j,u_i,\overline{u}_j)I𝔞K_\theta .$$
(13)
The propagator is the inverse of the kinetic operator. Furthermore, the operator $`K_\theta `$ satisfies $`(K_\theta )^2=I`$, as laid out below:
$$\begin{array}{cc}\hfill (K_\theta )^2& =d^3\stackrel{~}{g}d^3\stackrel{~}{u}K_\theta (g,\stackrel{~}{g},u,\stackrel{~}{u})K_\theta (\stackrel{~}{g},\overline{g},\stackrel{~}{u},\overline{u})\hfill \\ & =d^3\stackrel{~}{g}d^3\stackrel{~}{u}d^3bd^3\stackrel{~}{b}\delta (u_1h_\theta u_1^1g_1\stackrel{~}{g}_1^1)\delta (g_2\stackrel{~}{g}_2^1)\delta (g_3\stackrel{~}{g}_3^1)\delta (u_1b_1ϵ\stackrel{~}{u}_1^1)\delta (u_2b_2\stackrel{~}{u}_2^1)\delta (u_3b_3\stackrel{~}{u}_3^1)\hfill \\ & \times \delta (\stackrel{~}{u}_1h_\theta \stackrel{~}{u}_1^1\stackrel{~}{g}_1\overline{g}_1^1)\delta (\stackrel{~}{g}_2\overline{g}_2^1)\delta (\stackrel{~}{g}_3\overline{g}_3^1)\delta (\stackrel{~}{u}_1\stackrel{~}{b}_1ϵ\overline{u}_1^1)\delta (\stackrel{~}{u}_2\stackrel{~}{b}_2\overline{u}_2^1)\delta (\stackrel{~}{u}_3\stackrel{~}{b}_3\overline{u}_3^1).\hfill \end{array}$$
(14)
We integrate with respect to the $`\stackrel{~}{g}`$ variables
$$\begin{array}{cc}\hfill (K_\theta )^2=d^3\stackrel{~}{u}d^3bd^3\stackrel{~}{b}\delta (\stackrel{~}{u}_1h_\theta & \stackrel{~}{u}_1^1u_1h_\theta u_1^1g_1\overline{g}_1^1)\delta (g_2\overline{g}_2^1)\delta (g_3\overline{g}_3^1)\hfill \\ & \times \delta (u_1b_1ϵ\stackrel{~}{u}_1^1)\delta (u_2b_2\stackrel{~}{u}_2^1)\delta (u_3b_3\stackrel{~}{u}_3^1)\delta (\stackrel{~}{u}_1\stackrel{~}{b}_1ϵ\overline{u}_1^1)\delta (\stackrel{~}{u}_2\stackrel{~}{b}_2\overline{u}_2^1)\delta (\stackrel{~}{u}_3\stackrel{~}{b}_3\overline{u}_3^1),\hfill \end{array}$$
(15)
and then the $`\stackrel{~}{u}`$ variables
$$\begin{array}{cc}\hfill (K_\theta )^2=d^3bd^3\stackrel{~}{b}\delta (u_1b_1ϵh_\theta ϵ^1& b_1^1u_1^1u_1h_\theta u_1^1g_1\overline{g}_1^1)\delta (g_2\overline{g}_2^1)\delta (g_3\overline{g}_3^1)\hfill \\ & \times \delta (u_1b_1ϵ\stackrel{~}{b}_1ϵ\overline{u}_1^1)\delta (u_2b_2\stackrel{~}{b}_2\overline{u}_2^1)\delta (u_3b_3\stackrel{~}{b}_3\overline{u}_3^1).\hfill \end{array}$$
(16)
But $`b_1ϵh_\theta ϵ^1b_1^1=b_1h_\theta ^1b_1^1=h_\theta ^1`$ since $`ϵh_\theta ϵ^1=h_\theta ^1`$ and since $`b_1`$ and $`h_\theta ^1`$ are in the same commutative $`U(1)`$ subgroup. Furthermore $`ϵ\stackrel{~}{b}_1ϵ=\stackrel{~}{b}_1^1`$. Finally, redefining $`(b_1\stackrel{~}{b}_1^1)b_1`$, $`b_2\stackrel{~}{b}_2b_2`$ and $`b_3\stackrel{~}{b}_3b_3`$ gives us
$$(K_\theta )^2=d^3b\delta (g_1\overline{g}_1^1)\delta (g_2\overline{g}_2^1)\delta (g_3\overline{g}_3^1)\delta (u_1b_1\overline{u}_1^1)\delta (u_2b_2\overline{u}_2^1)\delta (u_3b_3\overline{u}_3^1)=I.$$
(17)
This leads to a nice closed form for the propagator
$$𝒫(g_i,\overline{g}_j,u_i,\overline{u}_j)=\frac{I+𝔞K_\theta }{1𝔞^2}.$$
(18)
### III.2 Feynman amplitudes and spin foam formulation
To construct a generic Feynman amplitude, we will analyse the structure and gluing properties of the propagator and vertex operator:
#### III.2.1 Vertex Operator
We scrutinise (12) in two parts:
$$𝒱(g_{ij},u_{ij})=\frac{\lambda }{4!}\underset{i=1}{\overset{4}{}}d\alpha _i\underset{j>i}{}db_{ij}\underset{\text{g variable part}}{\underset{}{\delta (\alpha _jg_{ji}^1g_{ij}\alpha _i^1)}}\underset{\text{u variable part}}{\underset{}{\delta (b_{ji}u_{ji}^1u_{ij}b_{ij}^1)}}.$$
The $`\delta `$-functions over the $`g`$ variables are the usual holonomies around the six wedges dual to the edges $`e`$, of a tetrahedron. So the model already has the structure of a 2-complex dual to a triangulation. The $`u`$ variables represent the momenta of the particles and as such are identified with the edges of the tetrahedron. Each edge of the tetrahedron is shared by two triangles. The $`\delta `$-functions above ensure that the momentum associated to an edge is the same when viewed from either of these triangles. This extra structure is not present in the Boulatov model.
#### III.2.2 Propagator
The operator (18) glues two tetrahedra at a triangular interface. From the analysis of the vertex operator above, we know that each triangle of a tetrahedron has three wedges (dual to each of its three edges) and three momenta associated to it. The propagator has two terms with different gluing properties.
$$𝒫(g_i,\overline{g}_j,u_i,\overline{u}_j)=\underset{𝒫_{massless}}{\underset{}{\frac{1}{1𝔞^2}I}}+\underset{𝒫_{massive}}{\underset{}{\frac{𝔞}{1𝔞^2}K_\theta }}.$$
For $`𝒫_{massless}`$, the $`\delta `$-functions over the $`g`$ variables glue three wedges of one tetrahedron to three wedges of another tetrahedron pairwise, giving the holonomies around three composite wedges. The $`\delta `$-functions over the $`u`$ variables ensure that the momenta on the edges of the two triangles match. For $`𝒫_{massive}`$, two of the $`\delta `$-functions over the $`g`$-variables act as for $`𝒫_{massless}`$, similarly for the $`u`$ variables. One of the $`\delta `$-functions, however, couples the $`g_1`$ and $`u_1`$ variables, effectively placing a massive particle at the point where the edge of the tetrahedron intersects the composite wedge. The final $`\delta `$-function to consider, $`\delta (u_1b_1ϵ\overline{u}_1^1)`$, places a tag on the edge $`e`$, of the tetrahedron with the mass insertion. The tag ensures that if another propagator further around the sequence of Feynman graph edges $`e^{}`$, forming the boundary of a (dual) face $`f^{}`$, inserts a mass along the edge e, it will cancel with the first according to the property $`(K_\theta )^2=I`$. We find in general that we can only have two possibilities when a face is fully assembled: no particle on the edge $`e`$, or one particle on the edge, according to whether there have been an even or an odd number of insertions of $`K_\theta `$ on the boundary of the face. To be clearer, we will calculate more explicitly the amplitude for a generic dual face in the next subsection. In the end the partition function for the field theory, when expanded in Feynman graphs (i.e. in a power series for $`\lambda `$) takes the form:
$$𝒵=𝒟\varphi e^{S[\varphi ]}=\underset{\mathrm{\Gamma }}{}\frac{\lambda ^{v[\mathrm{\Gamma }]}}{sym[\mathrm{\Gamma }]}𝒵[\mathrm{\Gamma }],$$
(19)
where $`v[\mathrm{\Gamma }]`$ is the number of vertices in the Feynman graph; $`sym[\mathrm{\Gamma }]`$ is its symmetry factor; and $`Z[\mathrm{\Gamma }]`$ is the amplitude for each Feynman graph, being given explicitly by:
$$𝒵[\mathrm{\Gamma }]=𝒩[\mathrm{\Gamma }]\underset{e^{}}{}d\alpha _e^{}\underset{e}{}du_e\underset{e\mathrm{}}{}\delta (G_e)\underset{e\mathrm{}}{}\delta (G_eu_eh_\theta u_e^1),$$
(20)
where $`𝒩[\mathrm{\Gamma }]`$ is a normalisation factor, discussed later, arising from the (partial) redundancy of the $`u`$ variables in each diagram and the $`𝔞`$ dependence; $`\alpha _e^{}`$ is the holonomy along an edge of the Feynman graph; $`G_e=_{e^{}f^{}}\alpha _e^{}^{\pm 1}`$ is the holonomy around a face $`f^{}`$ of the Feynman graph, which is dual to the edge of $`e`$ of the triangulation; and $`\mathrm{}`$ is the set of edges of the triangulation that have a particle present. We recognise in (20) the Ponzano-Regge amplitude (1). We see that the group field theory we have defined gives, in addition to a sum over all possible quantum gravity spin foams arising as usual as Feynman graphs of the theory, a sum over all possible massive spinless particle insertions in the spin foam, interpreted as a sum over all possible Feynman diagrams for a scalar field theory. Each gravity + particle configuration is weighted exactly by the amplitude of the Ponzano-Regge model coupled with massive spinless particles given in PR1 , and provides us also with a definite normalisation factor for each of these amplitudes.
### III.3 Amplitude for a generic face of the Feynman graph
Consider a face $`f^{}`$ of a Feynman graph, the boundary of which is formed by $`N`$ contiguous edges, $`e^{}`$, labelled $`e_1^{},\mathrm{},e_N^{}`$. The amplitude for this face is
$$A(f^{})=(d\mathrm{})𝒫_1^f^{}𝒱_{12}^f^{}\mathrm{}𝒫_N^f^{}𝒱_{N1}^f^{},$$
(21)
where $`𝒫_i^f^{}`$ are the $`\delta `$-functions from the propagator along $`e_i^{}`$ relating to the face $`f^{}`$; $`𝒱_{ij}^f^{}`$ are the $`\delta `$-functions from the vertex where $`e_i^{}`$ and $`e_j^{}`$ meet, pertaining to the face $`f^{}`$. If we contract all the $`g`$ variable $`\delta `$-functions we get a final $`\delta `$-function of the form
$$A(f^{})=\underset{n=0}{\overset{N}{}}(d\mathrm{})\delta (G_f^{}u_{a_1}h_\theta u_{a_1}^1\times \mathrm{}\times u_{a_n}h_\theta u_{a_n}^1)\times \left(\delta \text{-functions over the }u\text{ variables}\right),$$
(22)
the sum is over different combinations of mass insertions; $`nN`$ is the number of particle insertions in that specific term and $`G_f^{}`$ is the holonomy around the face built up from a product of $`\alpha _e^{}`$. If there are $`n`$ particle insertions in the $`g`$ variable part then there will be $`n`$ $`\delta `$-functions in the $`u`$ variable part with $`ϵ`$ inserted. Once we contract the $`u`$ variables we get that the masses cancel just as in the calculation of $`(K_\theta )^2=I`$. In the end, for each term in the sum we have two possibilities: For $`n`$ even we get
$$\delta (G_f^{})\delta (b_f^{}),$$
(23)
where $`b_f^{}`$ is a product of the $`b`$ variables around the face. Therefore we get no particle insertion and a pure gravity face modulo an extra factor which we take into the normalisation. For $`n`$ odd we get
$$\delta (G_f^{}uh_\theta u^1)\delta (b_f^{}ϵ),$$
(24)
thus a single particle insertion and another factor which we take into the normalisation.
### III.4 Overall normalisation of the Feynman graph
For a kinetic term with the structure $`I𝔞K_\theta `$, the propagator takes the form:
$$\frac{1}{1𝔞^2}\left(I+𝔞K_\theta \right),$$
as we have shown. This produces $`𝔞`$-dependent amplitudes when the expansion in Feynman graphs is performed. We have not specified what this $`𝔞`$ is and how it depends on the physical parameters of gravity or matter; indeed there is quite some freedom involved in choosing a specific expression for $`𝔞`$ and only further analysis of the model we proposed can narrow the range of possibilities down to a restricted one<sup>5</sup><sup>5</sup>5a natural possibility in view of (1) is to consider $`𝔞=\mathrm{\Delta }(\theta )`$ . This parameter will enter the normalisation coefficients controlling the relative strength of Feynman diagrams with and without particles. The normalisation factor will clearly contain an overall factor $`\left(1𝔞^2\right)^{|e^{}|}`$ where $`|e^{}|`$ denotes the number of dual edges of the two complex. There will be an additional factor $`𝔞`$ each time $`K_\theta `$ is inserted along a dual edge. If an even number of $`K_\theta `$ are inserted along a face no particle circulates along that face. A useful formula in order to get the right normalisation factor in a given example is
$$(I+𝔞K_\theta )^n=\frac{1}{2}\left((1+𝔞)^n+(1𝔞)^n\right)I+\frac{1}{2}\left((1+𝔞)^n(1𝔞)^n\right)K_\theta .$$
(25)
Along with this numerical factor, there is a singular factor coming from a redundant $`\delta `$-function for each face as shown in (23,24). For each face not carrying a particle we have a factor
$$_{U(1)}𝑑b\delta (b)=\underset{j}{}(2j+1),$$
(26)
and for each face carrying a particle we have a factor
$$_{U(1)}𝑑b\delta (bϵ)=\underset{j}{}(2j+1)(1)^j,$$
(27)
the sum, being over integer $`j`$, is obtained from the character expansion of the delta function. These expressions are unfortunately ill defined and they need a regularisation. A proper regularisation that can preserve all the symmetries of the theory is to replace the usual $`\mathrm{SU}(2)`$ group by a quantum group $`U_q(\mathrm{SU}(2))`$. It is expected that with this choice the key features of the model can be preserved and that the corresponding normalisation coefficients are given by
$$\underset{j=0}{\overset{N}{}}[2j+1]_qt^j=\frac{1+t+(q+q^1)t^{N+1}}{(1q^2t)(1q^2t)},$$
(28)
with $`q=\mathrm{exp}(i\frac{\pi }{2N+1})`$, $`[n]_q=(q^nq^n)/(qq^1)`$ and $`t=1`$ for a face with particle and $`t=+1`$ otherwise. Let us recall that if we consider the original Boulatov model such factors do not arise since they come from a redundant summation over spins dual to the $`u`$ variables. The Feynman graph amplitudes of the Boulatov model are, however, not equal to the physical quantum gravity amplitudes. In order to get the quantum gravity amplitudes one has to divide out the infinite volume of a gauge symmetry. This is conveniently done by restricting the summation over spins as presented in the first section. This gauge-fixing procedure which is well defined at the level of the quantum gravity amplitudes is, however, not fully understood at the level of the gft. When we extend the gft to include the momenta variables $`u`$, we also extend in a trivial way the gauge symmetry of the quantum gravity amplitude. This is where the additional factors (26, 27) come from. The gauge-fixing is trivially realised by fixing to $`0`$ the spin dual to the variables $`u_e`$. This is, we feel, one of the main open challenges in this domain: to understand if such a gauge fixing can be understood already at the gft level both for the Boulatov model and for our extension including particles; that is, whether we can already for the Boulatov model identify at the level of the gft the translational (or diffeomorphism) symmetry responsible for the divergences of the naïve gravity amplitude. A similar identification should also be implemented for our particle model. We do not resolve this issue in the present work.
## IV Discussion
### IV.1 Features of the model
We have seen that the model correctly generates spin foam configurations with some dual faces carrying particle data, i.e. a mass label, indicating that a particle of the given mass is propagating along the edge of the triangulation dual to that face. This means that the group field theory produces all possible Feynman graphs for a scalar field embedded in the triangulation dual to the quantum gravity 2-complex, specifying the field propagator on each line of the Feynman graph, and this only. Interestingly, this is enough, in this 3-dimensional setting, for specifying fully the dynamics of the particles, i.e. their interaction. In fact, this is dictated by the Bianchi identity constraining the sum of curvatures in the boundary of any 3-cell of the dual complex around any given vertex of the triangulation. When one or more particles are meeting at that vertex, thus interacting there, this implies momentum conservation for their interaction, which is the only content of any $`\varphi ^n`$ theory. Because any number of particles can be incident to any given vertex of the triangulation in the model we propose, this means that this corresponds to a scalar field theory with a potential given by a sum over any power of the field operators: $`\varphi ^3(x)+\varphi ^4(x)+\mathrm{}.`$.
We have seen that the crucial property of the modified kinetic term we propose for producing mass insertions in the spin foam amplitudes is, besides the extension of the field to 6 arguments, the property $`(K_\theta )^2=I`$ for the operator $`K_\theta `$ inserting the mass of the particles in the group field theory action. It is interesting to note that the presence of this operator which inserts particles also breaks a symmetry that the pure gravity model possesses. This bears some similarity with the fact that particles in 3d can be understood as defects breaking the translational symmetry of the theory without matter PR1 . The symmetry is the following: Let’s consider the transformation
$$\varphi (g_1,g_2,g_3;u_1,u_2,u_3)\varphi (𝐯_1g_1,g_2,g_3;𝐰_1u_1,u_2,u_3),$$
(29)
where $`𝐯_1,𝐰_1`$ are arbitrary fixed group elements. This transformation is clearly a symmetry of the pure gravity action (7). This symmetry is, however, broken by the insertion of a mass term and only the transformation
$$\varphi (g_1,g_2,g_3;u_1,u_2,u_3)\varphi (𝐯_1g_1,g_2,g_3;𝐯_1u_1,u_2,u_3),$$
(30)
preserves the action (9).
### IV.2 A direct generalisation
The model we presented in section III produces mass insertions in different faces of the spin foam 2-complex by application of the operator $`K_\theta `$ in the 1st argument of the field, carrying the gravity variable $`g_1`$. Because of permutation symmetry, the fact that one has chosen the 1st argument of the field for inserting a mass parameter and not, say, the 2nd is irrelevant, as one can easily convince oneself. Still, one may find the fact that a mass parameter is inserted in only -one- of the arguments of the field a bit unsatisfactory, for symmetry reasons. Here we want to discuss briefly what happens when one relaxes this condition. The result is that one can write down a generalised version of the model presented above, that is, however, basically equivalent to it, and leads to the same type of graphs being generated. One can consider defining a generalised kinetic term with matter insertions, defined using a sum of operators $`K_\theta (1)`$, $`K_\theta (2)`$, $`K_\theta (3)`$, each $`K_\theta (i)`$ inserting a mass parameter in the i-th argument of field, thus having as kinetic term an operator with the structure $`I𝔞(K_\theta (1)+K_\theta (2)+K_\theta (3))`$. It is obvious that a model like this would generate exactly the same type of graphs and amplitudes as the one we have defined above. It is also easy to realise that, once such an operator is introduced, there will be graphs in which the insertion of a $`K_\theta (1)`$ and a $`K_\theta (2)`$, say, in different propagators would produce the same amplitude that would have been generated by the use in the kinetic term of an operator of the form $`K_\theta (1,2)`$, i.e. an operator inserting a mass parameter in both the 1st and 2nd arguments of the field at once; and indeed one could generalise further the kinetic term to an operator of the form: $`I𝔞(K_\theta (1)+K_\theta (2)+K_\theta (3))𝔟(K_\theta (1,2)+K_\theta (2,3)+K_\theta (1,3))`$. Carrying this line of reasoning even further, one is led to the kinetic term:
$$𝒦=I𝔞(K_\theta (1)+K_\theta (2)+K_\theta (3))𝔟(K_\theta (1,2)+K_\theta (2,3)+K_\theta (1,3))𝔠K_\theta (1,2,3)IK_{total},$$
(31)
where the $`K_\theta (1,2,3)`$ is defined as the operator inserting a mass parameter in all the first three arguments of the field. Again, this generalised kinetic term leads to the same kind of Feynman graphs and amplitudes, as it is easy to verify. In fact the structure of the amplitudes is determined by the property $`K_\theta (i)^2=1`$, as we have explained, so it is enough for each mass insertion produced by the operators $`K_\theta (i,j)`$ and $`K_\theta (i,j,k)`$ to satisfy that property in order for the resulting amplitude to be of the form we have described. Of course, the normalisation factors for the amplitudes are going to be different from those of the simpler model presented in section III and it will depend in general on three different coupling constants $`𝔞`$, $`𝔟`$ and $`𝔠`$. This gives additional freedom that may well turn out to be useful in some situation.
It is interesting to note that there are several choices of coupling constant that lead to further simplifications. For example, one can show that in order for the added terms to satisfy $`(K_{total})^2I`$, then one needs to choose
$$𝔞=0,𝔟=0,\mathrm{or}𝔟=0,𝔞+𝔠=0,$$
(32)
that is
$$𝒦=I𝔠K_\theta (1,2,3),\mathrm{or}𝒦=I𝔞\left(K_\theta (1)+K_\theta (2)+K_\theta (3)K_\theta (1,2,3)\right).$$
(33)
Finally, a simple and highly symmetrical choice is
$$\begin{array}{cc}\hfill 𝒦& =(I𝔞K_\theta (1))(I𝔞K_\theta (2))(I𝔞K_\theta (3))\hfill \\ & =I𝔞(K_\theta (1)+K_\theta (2)+K_\theta (3))+𝔞^2(K_\theta (1,2)+K_\theta (2,3)+K_\theta (1,3))𝔞^3K_\theta (1,2,3).\hfill \end{array}$$
(34)
Since the $`(I𝔞K_\theta (i))`$ commute, we can easily compute the propagator
$$𝒫=\frac{(I+𝔞K_\theta (1))(I+𝔞K_\theta (2))(I+𝔞K_\theta (3))}{(1𝔞^2)^3}.$$
(35)
It is not clear at the present stage, however, which specific properties one should ask the gft propagator to fulfill.
### IV.3 Alternatives
After having discussed some possible generalisations of the proposed model leading to very similar structures, we would like to discuss two genuine alternatives to it. One based on a much simpler action constructed inserting a mass parameter in the simplest way in the usual 3-argument field, leading however to a problematic structure for the Feynman amplitudes, and thus showing the need for the 6-argument extension on which we have based our model. The other keeps the same structure of the model presented in section III, but inserts the mass parameter by means of a modification of the interaction term in the gft action, instead of the kinetic one. As we will show, this modification is completely harmless. Consider first a three argument field, the same on which the Boulatov model is based, and insert the mass parameter $`h_\theta `$ and the velocities for the relevant particles $`u_i`$ in the 1st argument of the field. The model we obtain is therefore realised by forgetting about the presence of the $`u_i`$ variables in the extra slots of the generalised field used in the model presented in section III. The action is then:
$$\begin{array}{cc}\hfill S[\varphi ]=\frac{1}{2}\underset{i=1}{\overset{3}{}}dg_idu_1(\varphi (g_1,g_2,g_3)& \varphi (g_1,g_2,g_3)𝔞\varphi (g_1,g_2,g_3)\varphi (u_1h_\theta u_1^1g_1,g_2,g_3))\hfill \\ & +\frac{\lambda }{4!}\underset{i=1}{\overset{6}{}}dg_i\varphi (g_1,g_2,g_3)\varphi (g_4,g_5,g_3)\varphi (g_4,g_2,g_6)\varphi (g_1,g_5,g_6).\hfill \end{array}$$
(36)
The Feynman rules for this action can be read out easily, and the amplitudes constructed in the usual way. The kinetic operator is:
$$𝒦(g_i,\overline{g}_j)=\underset{i=1}{\overset{3}{}}\delta (g_i\overline{g}_i^1)𝔞𝑑u_1\delta (u_1h_\theta u_1^1g_1\overline{g}_1^1)\underset{i=2}{\overset{3}{}}\delta (g_i\overline{g}_i^1)I𝔞K_\theta ,$$
(37)
while the interaction operator is the usual Boulatov one. It is easy to see from the expression for the kinetic operator, that the mass-inserting operator does not satisfy any property like $`(K_\theta )^2=I`$ anymore. Instead, its square gives:
$$(K_\theta )^2=𝑑u_1𝑑\stackrel{~}{u}_1\delta (\stackrel{~}{u}_1h_\theta \stackrel{~}{u}_1^1u_1h_\theta u_1^1g_1\overline{g}_1^1)\underset{i=2}{\overset{3}{}}\delta (g_i\overline{g}_i^1).$$
(38)
This makes the propagator much more complicated, being given by:
$$𝒫(g_i,\overline{g}_j)=I+\underset{n>0}{}(𝔞K_\theta )^n,$$
(39)
resulting in considerably more arduous Feynman diagrammatics. An example of an amplitude of this model, for a typical dual face is:
$$A(f^{})=\delta (G_f^{}u_{a_1}h_\theta u_{a_1}^1\times \mathrm{}\times u_{a_n}h_\theta u_{a_n}^1),$$
(40)
for $`n`$ less than the number of edges bounding the dual face $`f^{}`$. We see that there is no ‘multiple mass cancellation’ anymore, and we end up having multiple mass insertions on each face, in the typical Feynman diagram, each of which is associated to a different $`SU(2)/U(1)`$ velocity element. There is a possible physical interpretation of the resulting configuration, which is in terms of multi-particle states. In other words, we would have more than one particle with a given mass associated to a dual face and so to an edge of the triangulation. We are not in the position of being able to exclude this interpretation, or definitely reject the model that generates these configurations altogether; however, we feel that such an interpretation is problematic for at least two reasons: 1) it would imply that more than one particle is propagating along the same link of the triangulation, so they would be located at the same point in the manifold and interact together with other (multi-)particles at the vertices of the triangulation; most important, 2) interpreting the particle configurations in the triangulation as Feynman graphs of some effective field theory would become much less straightforward than for the model proposed in section III, and applying any procedure to extract this effective field theory, like that used in PR3 , to the amplitudes generated by this ‘simplified’ model would be quite cumbersome, if at all possible. We note that the same problem of ‘multiple mass insertions’ is generated by most other modifications of the structure of the model presented in section III, affecting the 3 extra arguments of the field, i.e. the $`u_i`$ variables. Such modifications lead to losing the property $`(K_\theta )^2=I`$ which is responsible for ‘mass cancellation’ in the dual faces when the operator is inserted more than once, and that leads to the presence of only one mass in each dual face, and furthermore, to consistency with the interpretation of the mass-labelled graphs in the triangulation as Feynman graphs of a field theory.
We conclude by mentioning instead a harmless modification of the model, that may even turn out to be useful in future studies. One can keep the structure of the field to be a function of 6 arguments, and one can keep the same form for the pure gravity action based on this 6 argument field, but choose to insert a mass parameter not in the kinetic term but in the interaction term of the group field theory action. The group field theory action would then be:
$$\begin{array}{cc}\hfill S[\varphi ]=& \frac{1}{2}\underset{i=1}{\overset{3}{}}dg_idu_i[P_\alpha P_b\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\overline{\alpha }}P_{\overline{b}}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)]\hfill \\ & +\frac{\lambda }{4!}\underset{i=1}{\overset{6}{}}dg_idu_i[P_{\alpha _1}P_{b_1}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\alpha _2}P_{b_2}\varphi (g_4,g_5,g_3;u_4,u_5,u_3)]\hfill \\ & \times [P_{\alpha _3}P_{b_3}\varphi (g_4,g_2,g_6;u_4,u_2,u_6)][P_{\alpha _4}P_{b_4}\varphi (g_1,g_5,g_6;u_1,u_5,u_6)]\hfill \\ & +\frac{\lambda }{4!}\underset{i=1}{\overset{6}{}}dg_idu_i𝔞[P_\theta P_{\alpha _1}P_{b_1}\varphi (g_1,g_2,g_3;u_1,u_2,u_3)][P_{\alpha _2}P_{b_2}\varphi (g_4,g_5,g_3;u_4,u_5,u_3)]\hfill \\ & \times [P_{\alpha _3}P_{b_3}\varphi (g_4,g_2,g_6;u_4,u_2,u_6)][P_{\alpha _4}P_{b_4}\varphi (g_1,g_5,g_6;u_1,u_5,u_6)].\hfill \end{array}$$
(41)
As we have anticipated, however, it is straightforward to verify that this action leads to the same type of Feynman graphs and amplitudes that one gets instead by modifying the kinetic term, the only difference being that one would get a different normalisation factor, in front of each amplitude.
## V Conclusions
We have defined a new group field theory for 3-dimensional Riemannian quantum gravity, and constructed its perturbative expansion in Feynman diagrams. These diagrams correspond to spin foam 2-complexes describing quantum gravitational degrees of freedom dual to 3d triangulations, as in the Boulatov model, but they also carry additional labels which describe massive spinless particles propagating in the spacetime one reconstructs from the gravity degrees of freedom. They have the interpretation of Feynman graphs for a scalar field theory embedded in the triangulation representing spacetime. The amplitudes for these diagrams have exactly the form obtained in PR1 from classical considerations, i.e. are given by the Ponzano-Regge model coupled to massive spinless particles, and are shown in PR3 to admit an effective (non-commutative) scalar field theory description, confirming the above physical interpretation.
The model presented possesses some quite non-trivial and interesting features, that we highlighted in the paper, and that deserve further analysis. This represents a first step in a program of analysing the coupling of matter and gauge fields to quantum gravity in the group field theory approach. The next steps would be first of all the inclusion of spin degrees of freedom and the construction of a group field theory reproducing the amplitudes given in PR1 for massive spinning particles. Second, one should develop the study of the gft observables in the presence of particles and its relation with spin networks with open ends. Next, one would like to show how the effective field theory picture for the particle degrees of freedom can be obtained directly, and possibly in a simpler way, from the group field theory formulation. Finally, the problem of the coupling of gauge fields and the description of their interaction with both gravity and matter fields, again in the group field theory formalism, should be tackled.
To conclude, as we mention in the text, one of the most pressing issues for our model and the Boulatov model is to understand whether there are symmetries at the gft level that justify the gauge fixing needed at the level of the Feynman graphs to reproduce physical 3d gravity amplitudes. This is necessary in order to promote this type of gft to a fundamental model of three dimensional gravity coupled to matter.
Acknowledgments: We would like to thank K. Krasnov for keeping us informed of his progress and PI for an invitation which initiated this collaboration. J.R. would like to thank the occupants of room B0.10 for correcting innumerable typos.
|
warning/0506/cond-mat0506506.html
|
ar5iv
|
text
|
# Implementation of the histogram method for equilibrium statistical models using moments of a distribution
## I Introduction
The calculation of the expectation value of some operator $`\mathrm{\Phi }`$ in a system in contact with a heath-bath of temperature $`T`$
$$\mathrm{\Phi }_T=\frac{_i\mathrm{\Phi }_i\mathrm{exp}(E_i/k_BT)}{_i\mathrm{exp}(E_i/k_BT)},$$
(1)
is the usual task of equilibrium statistical mechanics. Here $`E_i`$ and $`\mathrm{\Phi }_i`$ are, respectively, the energy and the value of $`\mathrm{\Phi }`$ associated to the $`i`$-th configuration, and the sum runs over all possible configurations. These expectation values are usually impossible to calculate analytically, but there are many techniques to approximate expectation values in the canonical ensemble. Some of them are quite powerful and general, like high-temperature expansions and renormalization group techniques.
A different approach, quite successful in many applications, is the numerical simulation (“Monte Carlo Simulation”) of the ensemble metropolis1 ; newman-barkema ; landau-binder . This paper is concerned with the Histogram Method, introduced years ago to extrapolate the results of a Monte Carlo simulation conducted at a point in parameter space —for instance, some given temperature and magnetic field—, to a range of those parameters. It was formulated by by Ferrenberg and Swendsen ferrenberg1 ; ferrenberg2 , although there are some earlier proposals valleau1 . For completeness, a short description of this technique follows: Assume for instance that one has a model where the only control parameter is the temperature, and wants to compute the behavior of some function $`f`$ of the energy $`E`$. At some given $`T`$ one performs a Monte Carlo simulation using any given algorithm that generates $`Q`$ configurations with the correct probability, given by their Boltzmann weight, measures $`f(E)`$ for each configuration, and calculates finally the average value $`f(E)`$. This average may be taken directly, or going through the preliminary construction of a normalized histogram $`W_T(E)`$ for the energies found in the simulation. One can write then
$$f(E)_T\frac{1}{Q}\underset{q=1}{\overset{Q}{}}f(E_q)=\underset{E}{}W_T(E)f(E),$$
(2)
which is approximation to Eq. (1) with $`\mathrm{\Phi }=f(E)`$. This equation can also written in terms of the density of states $`g(E)`$
$$f(E)_T=\frac{_Eg(E)f(E)\mathrm{exp}(E/(k_BT))}{_Eg(E)\mathrm{exp}(E/(k_BT))}.$$
(3)
Comparing Eqs. (2) and (3) it is clear that $`W_T(E)`$ is proportional to $`g(E)\mathrm{exp}(E/(k_BT))`$, and so an approximation for the density of states —ignoring normalization for the moment— is given by $`g(E)W_T(E)\mathrm{exp}(E/(k_BT))`$. Notice that $`g(E)`$ is a property of the system itself, independent of temperature.
With an approximated density of states at hand the results of the simulation are extended to other temperatures. Eq. (3) written for a different temperature $`T^{}`$ gives
$`f(E)_T^{}`$ $`=`$ $`{\displaystyle \frac{_Eg(E)f(E)\mathrm{exp}(E/(k_BT^{}))}{_Eg(E)\mathrm{exp}(E/(k_BT^{}))}}`$ (4)
$``$ $`{\displaystyle \frac{_EW_T(E)f(E)\mathrm{exp}(E/(k_BT)E/(k_BT^{}))}{_EW_T(E)\mathrm{exp}(E/(k_BT)E/(k_BT^{}))}},`$
and one says that the histogram $`W_T(E)`$ has been reweighted.
## II Expansion in Moments
In general, with the Histogram Method results obtained at a given point in the space of control parameters can be extrapolated to a neighborhood of that point. However, the implementation of the method can be cumbersome for models with more than one parameter, because the size of the histograms generated becomes too large. For instance, an Ising model in a simple cubic lattice, with nearest-neighbor, next-nearest-neighbor and magnetic couplings has a 3D parameter space. When applying the histogram method to this system one faces the problem of handling a histogram with a total number of bins that rapidly becomes impossible to accommodate in the computer’s physical memory. For instance, for a very modest size $`L=16`$, one needs to store some $`16^34\times 10^3`$ possible values for each of the 3 couplings, and this gives a histogram with around $`64\times 10^9`$ bins. Working with 4 Byte integers (which may be insufficient for very large runs) gives a memory requirement around 250 GBytes, realizable in large workstations and supercomputers, but usually beyond the reach of the commodity machines used by many scientist today not-dense .
Among possible solutions one could coarse-grain the binning of the histogram (this is unavoidable if one is dealing with continuous variables, like in the XY or Heisenberg Models); one could also store a record of the operators of interest for all configurations generated by the Monte Carlo algorithm (this has been recommended for Hamiltonians with continuous variables newman-barkema ; ferrenberg3 ). Here I propose a different option that arises from two simple facts: one, most of the quantities one is interested in calculating can be expressed in terms of moments of the same operators that constitute the Hamiltonian. Two, the reweighting function can be easily expanded in a power series. Then, the physical quantities one wants to estimate at different values of the couplings can be obtained as power series in the increments of said couplings, using the moments of the operators obtained at some fixed point in parameter space as coefficients. A similar approach was proposed in Ref. rickman1 , using cumulants instead of moments; This approach via moments is bit simpler. Here I also show some details of implementation needed to insure larger ranges of numerical convergence.
To fix ideas, I shall use a simple example, namely, the Ising Model with zero magnetic field. The Hamiltonian is
$$=J\underset{<a,b>}{}S_aS_b.$$
(5)
Assume now that one calculates the expectation value of some power of the magnetization, say $`M^2`$, at some given temperature $`T`$. Define for convenience the adimensional energy operator $`\mathrm{\Theta }_{<a,b>}S_aS_b`$ and the couplings $`KJ/k_BT`$. From Eq. (1) one gets
$$M^2_K=\frac{_iM_i^2\mathrm{exp}(K\mathrm{\Theta }_i)}{_i\mathrm{exp}(K\mathrm{\Theta }_i)}.$$
(6)
Now, if one wants to calculate this same expectation value at some other temperature, one could store the histogram for $`\mathrm{\Theta }`$ and then reweight it. But one can also write the desired expectation value
$$M^2_K^{}=\frac{_iM_i^2\mathrm{exp}(K^{}\mathrm{\Theta }_i)}{_i\mathrm{exp}(K^{}\mathrm{\Theta }_i)}=\frac{_iM_i^2\mathrm{exp}(K\mathrm{\Theta }_i)\mathrm{exp}((K^{}K)\mathrm{\Theta }_i)}{_i\mathrm{exp}(K\mathrm{\Theta }_i)\mathrm{exp}((K^{}K)\mathrm{\Theta }_i)},$$
(7)
and expand the exponential containing $`K^{}K`$ in a Taylor series
$$M^2_K^{}=\frac{_iM_i^2\mathrm{exp}(K\mathrm{\Theta }_i)_{n=0}^{\mathrm{}}(K^{}K)^n\mathrm{\Theta }_i^n/n!}{_i\mathrm{exp}(K\mathrm{\Theta }_i)_{n=0}^{\mathrm{}}(K^{}K)^n\mathrm{\Theta }_i^n/n!}.$$
(8)
For finite systems the order of the sums can be safely interchanged, giving
$$M^2_K^{}=\frac{_{n=0}^{\mathrm{}}[(KK^{})^n/n!]_iM_i^2\mathrm{\Theta }_i^n\mathrm{exp}(K\mathrm{\Theta }_i)}{_{n=0}^{\mathrm{}}[(KK^{})^n/n!]_i\mathrm{\Theta }_i^n\mathrm{exp}(K\mathrm{\Theta }_i)}.$$
(9)
One can now divide numerator and denominator by the partition function $`Z(K)`$ and find that the terms at the end of both expressions are the expectation values (moments, in short), of $`M^2\mathrm{\Theta }^n`$ and $`\mathrm{\Theta }^n`$, respectively, taken at $`K`$. Denoting $`\delta KKK^{}`$ one gets
$$M^2_K^{}=\frac{_{n=0}^{\mathrm{}}(\delta K)^nM^2\mathrm{\Theta }^n_K/n!}{_{n=0}^{\mathrm{}}(\delta K)^n\mathrm{\Theta }^n_K/n!}.$$
(10)
Since the expansion of the exponential function is convergent everywhere, there is no a priori limitation for this approach. But it is clear that one can never actually calculate all the required moments, and a truncation needs to be applied. This immediately introduces very strong bounds in the applicability of the method, since for large $`\delta K`$ the two sums in the previous expression require the inclusion of higher and higher moments if reasonable results are going to be obtained. A discussion about the number of moments needed to insure certain range of convergence is given later on.
Once the series are truncated, one faces a more serious difficulty: consider again the $`M^2`$ example, and assume that a simulation has been conducted in a lattice with $`N`$ spins. Since $`M`$ is an extensive quantity, $`M^2`$ is of order $`N^2`$, and the range of $`K`$ where some degree of numerical convergence can be achieved is quite small. To find a way around this problem —at least partially—, begin by working with densities instead of with the original operators. Going back to our example, define $`mM/N`$ and $`\theta \mathrm{\Theta }/N`$. The previous expression is rewritten as
$$m^2_K^{}=\frac{_{n=0}^{\mathrm{}}(N\delta K)^nm^2\theta ^n_K/n!}{_{n=0}^{\mathrm{}}(N\delta K)^n\theta ^n_K/n!}$$
(11)
Now it is clear that, since $`\theta `$ is of order one, one may expect numerical convergence of the strongly truncated series up to
$$\mathrm{\Delta }K\frac{1}{N}\mathrm{\Delta }T\frac{T^2}{N},$$
(12)
which is a very narrow range; here ‘strongly truncated’ means a series with a few terms, say, less than 20. This is actually a very conservative bound, since the $`n!`$ in the denominators improve numerical convergence (and guarantees it for the infinite series).
The convergence of the truncated series can be improved if one notices that the desired moments are taken from a sharply peaked distribution. The combination of the fast increasing density of states $`g(\mathrm{\Theta })`$ and the fast decreasing Boltzmann weight implies the existence of a narrow distribution $`W_K(\theta )`$ —that is, the histogram in $`\theta `$—, centered in some value $`\theta _c`$. Going back now to Eq. (6), one may see that the zero values of the operators appearing in the energy (and therefore in the energy density) can be easily shifted. For instance, if one shifts the density $`\theta `$ by some value $`\theta _R`$, the expectation value for $`m^2`$ becomes
$$m^2_K=\frac{_im_i^2\mathrm{exp}(NK\theta _i)}{_i\mathrm{exp}(NK\theta _i)}=\frac{_im_i^2\mathrm{exp}(NK(\theta _i\theta _R))}{_i\mathrm{exp}(NK(\theta _i\theta _R))},$$
(13)
since the factors of $`\mathrm{exp}(NK\theta _R)`$ cancel in the fraction. Now, if the value one chooses for $`\theta _R`$ is close to the center of the distribution, the moments one gets for $`\delta \theta \theta \theta _R`$ are going to be numerically small —and become smaller the higher the moment and the larger the lattice—, simply because of the narrowness of the originating distribution. Retracing the steps taken going from Eqs. (7) to (10) one gets the final expression for $`m^2_K^{}`$ as
$$m^2_K^{}=\frac{_{n=0}^n^{}(N\delta K)^nm^2(\delta \theta )^n_K/n!}{_{n=0}^n^{}(N\delta K)^n(\delta \theta )^n_K/n!},$$
(14)
where both numerator and denominator series have been truncated to $`n^{}`$ terms. Now the fast growth in the $`N\delta K`$ coefficients is partially balanced by a fast decrease in the values of $`(\delta \theta )^n_K`$ and $`m^2(\delta \theta )^n_K`$, and in this way the applicability of the method is extended to a much wider range in couplings. It is not too difficult to calculate the range of expected convergence of the extrapolation. In fact, for temperatures away from the critical point the width of $`W_K(\theta )`$ is order $`1/\sqrt{N}`$, from where one gets a range of convergence for the extrapolation of
$$\mathrm{\Delta }K\frac{1}{\sqrt{N}}\mathrm{\Delta }T\frac{T^2}{\sqrt{N}},$$
(15)
which comes from estimating $`N\mathrm{\Delta }K\mathrm{\Delta }\theta 1`$. Close to the critical point the width of $`W_T(\theta )`$ increases, and scales as $`\mathrm{\Delta }\theta N^{\alpha /(2\nu d)}/\sqrt{N}`$. From here one estimates then a range of convergence
$$\mathrm{\Delta }T\frac{T^2}{\sqrt{N^{1+\alpha /(\nu d)}}}.$$
(16)
Here $`\alpha `$ and $`\nu `$ are the exponents for specific heath and correlation distance close to the critical point, and $`d`$ is the dimensionality. For a small $`\alpha `$ this is still a much wider range than the one gotten when no shifts are used. Notice that for the 2D Ising Model one replaces $`N^{\alpha /(\nu d)}`$ by $`\mathrm{ln}N`$ and gets an scaling
$$\mathrm{\Delta }T\frac{T^2}{\sqrt{N\mathrm{ln}N}}.$$
(17)
Notice finally that the example given here, that extrapolates for $`m^2`$, can be easily extended to any other lattice operator.
### II.1 Example: the 2-D Ising Model: extrapolations at $`h=0`$
A numerical test of the method was done for the two-dimensional Ising Model in square lattices of sizes 16, 32 and 64, with the Hamiltonian given in Eq. (5). Two Monte Carlo simulations were carried out using the Wolff algorithm wolf1 , at a temperature of $`k_BT/J=2.27`$, close to the critical temperature for the model. All moments of the form $`m^l(\delta \theta )^n`$ and $`(\delta |m|)^l(\delta \theta )^n`$ for $`0l,n16`$ were recorded. First a short run with no shifts in either $`\theta `$ or $`|m|`$ was done, using $`2\times 10^6`$, $`3\times 10^6`$ and $`4.4\times 10^6`$ Wolff iterations for $`L=16`$, $`32`$ and $`64`$, respectively, after discarding transients of $`120`$, $`170`$ and $`270`$ iterations. The running time for this set of simulations was about $`1,800`$ seconds in a 2.4 GHz Pentium 4 processor. The exact solutions for the energy and the specific heath were obtained from the analytic free energy for finite lattices found by Kaufman kaufman1 . Fig. (1) shows the results of the extrapolation for the adimensional energy density $`\theta `$ using either 14 or 15 moments numberofmoments , when no shift in $`\theta `$ has been implemented, compared with the exact results. It is clear that the range of applicability of the extrapolation is extremely narrow, as expected from Eq. (12). That estimate gives here $`\mathrm{\Delta }T_{32}=0.0050`$ and $`\mathrm{\Delta }T_{64}=0.0013`$. This estimate is quite conservative, and the figure gives ranges of numerical convergence which are 2 to 3 times longer.
The second simulation was a larger run where the results for $`\theta `$ and $`|m|`$ from the first were used as reference values $`\theta _R`$ and $`|m|_R`$ for the shifts. The total numbers of Wolff iterations were $`4\times 10^7`$, $`6\times 10^7`$ and $`8.8\times 10^7`$ for $`L=16`$, $`32`$ and $`64`$, respectively, after transients of $`120`$, $`170`$ and $`250`$ iterations were discarded. The total time for this set of simulations was about $`37,000`$ seconds with the same processor. Data were divided in 20 blocks in order to generate error estimates. Extrapolation for energy and specific heath were calculated. The results for the energies are given in Fig. (2), those for specific heath in Fig. (3), and in both figures a comparison with the exact results is given. It is clear that the range of applicability of the extrapolation is larger, and actually a bit larger than the estimation made in Eq. (17), which gives $`\mathrm{\Delta }T_{32}=0.061`$ and $`\mathrm{\Delta }T_{64}=0.028`$. Two important points should be remarked: First, for $`\theta `$ the direction in which the extrapolated curves deviate from the correct results, for low $`T`$, depends on the number of moments taken into account; they deviate upwards when one uses 14 moments, and downwards when using 15 moments. A similar behavior is found for $`c`$, except that now the two extrapolations deviate in different directions at both ends of the range of convergence. This gives a very simple way of bounding the range of convergence of the algorithm. Second, the statistical errors in the extrapolations grow as one moves away from the simulated temperature, but the effect of this growth is smaller than the effect of the change in the number of included moments. The behavior of both types of errors are given in Fig. (4), which shows the difference between extrapolated and exact energies for $`L=32`$. The behavior for errors is similar for the specific heath. Notice that the statistical errors in the extrapolated results actually become smaller for temperatures a bit below the point where the simulation was carried out; this curious phenomenon has been studied in full histogram extrapolations ferrenberg3 .
The moments of $`|m|`$ were used to generate an extrapolated estimation of the susceptibility $`\chi ^{}`$, defined by $`\chi ^{}=NK^{}(|m|^2|m|^2)=NK^{}(m^2|m|^2)`$. This expression is used instead of the true susceptibility $`\chi `$ which in a numerical experiment does not give the expected peak close to the critical temperature ferrenberg4 . Fig. (5) shows a comparison between the extrapolated $`\chi ^{}`$ and values obtained in other individual simulations. These were obtained at their nominal temperatures using the Wolff algorithm. It is remarkable how the extrapolated results manage to reproduce reasonable well the peak in the susceptibility. Here one can also notice that the order of the approximation again decides in which direction the extrapolated results deviate form the actual values, and so a simple comparison between the 15- and 16-moments expansions gives bounds for the region of convergence of the method.
### II.2 Example: the 2-D Ising Model with magnetic field
As mentioned before, this algorithm becomes attractive especially in cases where one has to deal with Hamiltonians composed of several operators, where the sizes of the histograms needed for reweighting may overflow the available memory. As a simple example consider again the 2-dimensional nearest-neighbor Ising Model, but now with a magnetic field. The Hamiltonian is
$$=J\underset{<a,b>}{}S_aS_bH\mu \underset{a}{}S_a,$$
(18)
giving a Boltzmann weight
$$\mathrm{exp}(/k_BT))=\mathrm{exp}(K\mathrm{\Theta }+hM),$$
(19)
where the additional definitions $`M_aS_a`$ and $`hH\mu /k_BT`$ have been introduced. Consider now a simulation carried out at some temperature and magnetic field. Assuming that one wants to extrapolate the expectation value of some operator $`\mathrm{\Phi }`$, one gets
$$\mathrm{\Phi }_{K^{},h^{}}=\frac{_i\mathrm{\Phi }_i\mathrm{exp}(K\mathrm{\Theta }_i+hM_i)\mathrm{exp}((K^{}K)\mathrm{\Theta }_i+(h^{}h)M_i)}{_i\mathrm{exp}(K\mathrm{\Theta }_i+hM_i)\mathrm{exp}((K^{}K)\mathrm{\Theta }_i+(h^{}h)M_i)}.$$
(20)
As before, it is better to change all operators into their densities, defining $`mM/N`$ and $`\varphi \mathrm{\Phi }/N`$. Also, the density $`\theta `$ should be shifted so that its distribution is centered close to zero (this is unnecessary for $`m`$, since its distribution is symmetric around $`m=0`$). Performing these operations and the Taylor expansions for the exponential containing $`\delta K`$ and $`\delta h`$, one gets, repeating the same steps that gave Eq. (14), the following expression
$$\varphi _{K^{}h^{}}=\frac{_{l=0}^{\mathrm{}}_{n=0}^{\mathrm{}}(N\delta K)^l(N\delta h)^n\varphi (\delta \theta )^lm^n_{K,h}/(l!n!)}{_{l=0}^{\mathrm{}}_{n=0}^{\mathrm{}}(N\delta K)^l(N\delta h)^n(\delta \theta )^lm^n_{K,h}/(l!n!)}.$$
(21)
One should pay attention here to the fact that for low temperatures the histogram $`W_{K,h}(\theta ,m)`$ becomes bimodal in $`m`$, and the assumption of a distribution with a single narrow peak is no longer valid. However, close to the critical point the width in $`m`$ of such histogram remains small, and one can still get by with the first few moments. Occasionally, it may be convenient to work in terms of $`|m|`$, whose distribution remains unimodal.
The data obtained from the previous simulations were now used to generate the behavior of the magnetization and the susceptibility at non-zero values of $`H`$. For the magnetization one gets, after truncation
$$m_{K^{}h^{}}=\frac{_{l=0}^l^{}_{n=0}^n^{}(N\delta K)^l(Nh)^n(\delta \theta )^lm^{n+1}_{K,h=0}/(l!n!)}{_{l=0}^l^{}_{n=0}^n^{}(N\delta K)^l(Nh)^n(\delta \theta )^lm^n_{K,h=0}/(l!n!)};$$
(22)
an analogous expression is obtained for $`m^2_{K^{},h^{}}`$, and from here the true susceptibility can be computed. The results are shown in Fig. (6), which shows: (a) magnetization vs. $`H`$ for $`L=16`$, $`32`$ and $`64`$ at the critical temperature $`T_c=2/\mathrm{ln}(1+\sqrt{2})=2.269185\mathrm{}`$, and (b) magnetization vs. $`H`$ for $`L=32`$ at $`T=2.20`$, $`2.27`$ and $`2.34`$. In all cases the magnetization has been extrapolated from the simulations that were carried out at $`H=0`$ and $`T=2.27`$. In a slight departure from what was done before, here the denominator in Eq. (22) was calculated using 16 moments while the sums in the numerator were truncated to $`15`$ or $`14`$ moments. The individual points have been calculated using a modified version of the Wolff algorithm, were the acceptance ratio depends on the change of energy due to a cluster flipping in the presence a magnetic field $`H`$. It is clear from the figure that the expansion in moments reproduces quite well the behavior of the magnetization for each temperature and lattice size, and in particular it manages to show the large growth of $`m`$ with $`H`$ as $`T`$ is reduced. The splits in the extrapolation curves correspond to the separation of the $`15`$\- and $`14`$-moment extrapolations, and mark the end of the ranges of convergence. For small lattices these splits do not appear in the $`H`$ range tested here.
Finally, Figs. (7) and (8) show the true susceptibility $`\chi `$ as a function of $`H`$, for $`T=T_c`$ and $`L=16`$, $`32`$ and $`64`$, and for $`L=32`$ and $`T=2.20`$, $`2.27`$ and $`2.34`$. It should be noticed that the extrapolation manages to cover quite well the whole peak in $`\chi `$. The results for $`L=32`$ and different temperatures also show an excellent agreement between extrapolations and individual simulations, and show how the method can really extrapolate in more than one parameter. Notice that, as expected, $`\chi `$ grows as $`T`$ is reduced, even below $`T_c`$; for $`H=0`$ the true susceptibility is just proportional to $`m^2`$. Otherwise, the behavior the extrapolated quantities vis a vis the individual simulations at $`H0`$ is in all respects analogous the results found before: a very good correspondence for small $`H`$ —that is, close to where the simulation was carried out—, followed by a large deviation, which depends on the number of moments included in the extrapolation.
## III Conclusions
This paper shows how to implement the histogram method for extrapolation of results of a Monte Carlo simulation using the moments of the histogram. This approach has several advantages over the direct method —histogram construction and posterior reweighting—, and over the method of storing configurations for their reweighting (named “histogram on the fly” in Ref. ferrenberg3 ). To start with, the resulting expressions for the extrapolated quantities are given by very simple and conceptually appealing formulas. Second, the ranges of applicability of the method become evident simply by changing the number of moments included in the extrapolation. Third, the amount of computer memory and physical storage needed are so small that one may without any problem generate several repetitions of the simulations so as to generate in a simple way the error estimates for the extrapolated quantities. And finally, this approach eliminates completely the need to choose a binning size in cases of continuous variables.
One needs however to balance these benefits against the cost of the extra approximation involved in the method. After all, replacing a full histogram for its first few moments necessarily reduces precision. How many moments one really needs to keep in any given simulation so that not too much information is lost is an issue that has to be considered with some care. On the one hand, it is clear that increasing too much the number of stored moments not only defeats one of the motivations for this approach, which is to work with a limited memory, but also necessarily runs into the limits of reliability imposed by the statistical errors of the simulation. On the other, not keeping enough moments implies a waste of simulation time. As a first approach to the answer to this question one can consider the following estimation, done here for a one-coupling Hamiltonian: consider the preliminary run needed for the estimation of $`\theta _c`$. This preliminary run can also be used to obtain a rough estimate of the width $`\sigma _\theta `$ of the distribution. Now, the order of magnitude of the moments $`(\delta \theta )^n`$ will be around $`\sigma _\theta ^n`$, and so the terms needed in an estimation with a shifted coupling look like (see Eq. (14))
$$\frac{(\mathrm{\Delta }K)^n(\delta \theta )^n}{n!}\left(\frac{e\mathrm{\Delta }K\sigma _\theta }{n}\right)^n.$$
(23)
It is not difficult to show, using a steepest descent approximation, that as a function of $`n`$ this expression behaves as a Gaussian centered in $`n_{\mathrm{max}}=\mathrm{\Delta }K\sigma _\theta `$, with a width given by $`\sqrt{n_{\mathrm{max}}}`$ and height $`\mathrm{exp}(n_{\mathrm{max}})`$. Therefore one gets
$$\frac{(\mathrm{\Delta }K)^n(\delta \theta )^n}{n!}\mathrm{exp}(\mathrm{\Delta }K\sigma _\theta )\mathrm{exp}\left(\frac{(n\mathrm{\Delta }K\sigma _\theta )^2}{2\mathrm{\Delta }K\sigma _\theta }\right).$$
(24)
The conclusion is then the following: given an initial estimation of the width $`\sigma _\theta `$ of the distribution, and assuming that a certain maximum extrapolation range $`\mathrm{\Delta }K`$ is going to be used, the main contribution to the extrapolation comes from the moments with $`n`$ around $`n_{\mathrm{max}}=\mathrm{\Delta }K\sigma _\theta `$. Besides, one finds that the moments with $`n`$ such that $`nn_{\mathrm{max}}\sqrt{n_{\mathrm{max}}}`$ are basically irrelevant. The numerical results shown here display a much larger dependence on the number of moments used than on the statistical spread of the data, suggesting that several more moments may have been used in the extrapolation before reaching the limits given by statistical spread.
###### Acknowledgements.
I want to thank F. Sastre for a careful reading of the manuscript and for many helpful comments and suggestions. This work has been supported by CONACyT through grant No. 40726-F.
|
warning/0506/math0506221.html
|
ar5iv
|
text
|
# On seminormal monoid rings
## 1. Introduction
Let $`M`$ be an affine monoid, i.e. $`M`$ is a finitely generated commutative monoid which can be embedded into $`^m`$ for some $`m`$. Let $`K`$ be a field and $`K[M]`$ be the affine monoid ring associated to $`M`$. Sometimes $`M`$ is also called an affine semigroup and $`K[M]`$ a semigroup ring. The study of affine monoids and affine monoid rings has applications in many areas of mathematics. It establishes the combinatorial background for the theory of toric varieties, which is the strongest connection to algebraic geometry. In the last decades many authors have studied the relationship between ring properties of $`K[M]`$ and monoid properties of $`M`$. See Bruns and Herzog for a detailed discussion and Bruns, Gubeladze and Trung for a survey about open problems.
A remarkable result of Hochster states that if $`M`$ is a normal, then $`K[M]`$ is Cohen–Macaulay. The converse is not true. It is a natural question to characterize the Cohen–Macaulay property of $`K[M]`$ for arbitrary affine monoids $`M`$ in terms of combinatorial and topological information related to $`M`$. Goto, Suzuki and Watanabe could answer this question for simplicial affine monoids. Later Trung and Hoa generalized their result to arbitrary affine monoids. But the characterization is technical and not easy to check. Thus it is interesting to consider classes of monoids which are not necessarily simplicial, but nevertheless admit simple criteria for the Cohen-Macaulay property.
One of the main topics in the thesis of the second author were seminormal affine monoids and their monoid rings. Recall that an affine monoid $`M`$ is called seminormal if $`z\mathrm{G}(M)`$, $`z^2M`$ and $`z^3M`$ imply that $`zM`$. Here $`\mathrm{G}(M)`$ denotes the group generated by $`M`$. Hochster and Roberts \[10, Proposition 5.32\] noted that $`M`$ is seminormal if and only if $`K[M]`$ is a seminormal ring. See Traverso or Swan for more details on the latter subject. In general there exist Cohen–Macaulay affine monoid rings which are not seminormal, and there exist seminormal affine monoid rings which are not Cohen–Macaulay. One of the main goals of this paper is to understand the problem in which cases $`K[M]`$ is Cohen–Macaulay for a seminormal affine monoid $`M`$. Another question is to characterize the seminormality property of affine monoids.
Let us go into more detail. Let $`R`$ be a Noetherian ring and let $`N`$ be a finitely generated $`R`$-module. The module $`N`$ satisfies Serre’s condition ($`S_k`$) if
$$\mathrm{depth}N_𝔭\mathrm{min}\{k,dimN_𝔭\}$$
for all $`𝔭\mathrm{Spec}R`$. For trivial reasons Cohen–Macaulay rings satisfy Serre’s condition ($`S_k`$) for all $`k1`$. The main result of and a result of Schäfer and Schenzel show that for a simplicial affine monoid $`M`$ the ring $`K[M]`$ is Cohen–Macaulay if and only if $`K[M]`$ satisfies ($`S_2`$). After some prerequisites we study in Section 3 the question to characterize the ($`S_2`$) property for $`K[M]`$ if $`M`$ is a seminormal monoid. In the following let $`\mathrm{C}(M)`$ be the cone generated by $`M^m`$. The main result in this section already appeared in the thesis of the second author and states:
###### Theorem.
Let $`M^m`$ be a seminormal monoid and let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$. Then the following statements are equivalent:
1. $`K[M]`$ satisfies $`(S_2)`$;
2. For all proper faces $`F`$ of $`\mathrm{C}(M)`$ one has
$$M\mathrm{int}F=\underset{FF_j}{}\mathrm{G}(MF_j)\mathrm{int}F;$$
3. $`\mathrm{G}(MF)=_{FF_j}\mathrm{G}(MF_j)`$.
Here $`\mathrm{int}F`$ denotes the relative interior of $`F`$ with respect to the subspace topology on the affine hull of $`F`$. Let us assume for a moment that $`M`$ is positive, i.e. $`0`$ is the only invertible element in $`M`$. In order to decide whether $`K[M]`$ is a Cohen–Macaulay ring, one must understand the local cohomology modules $`H_𝔪^i(K[M])`$ where $`𝔪`$ denotes the maximal ideal of $`K[M]`$ generated by all monomials $`X^a`$ for $`aM\{0\}`$, because the vanishing and non-vanishing of these modules control the Cohen–Macaulayness of $`K[M]`$. Already Hochster and Roberts noticed that certain components of $`H_𝔪^i(K[M])`$ vanish for a seminormal monoid. Our result 4.3 in Section 4 generalizes their observation, and we can prove the following
###### Theorem.
Let $`M^m`$ be a positive affine seminormal monoid such that $`H_𝔪^i(K[M])_a0`$ for some $`a\mathrm{G}(M)`$. Then $`a\mathrm{G}(MF)F`$ for a face $`F`$ of $`\mathrm{C}(M)`$ of dimension $`i`$. In particular,
$$H_𝔪^i(K[M])_a=0\text{if }a\mathrm{C}(M).$$
As a consequence of this theorem and a careful analysis of the groups $`H_𝔪^i(K[M])`$ we obtain in 4.7 that under the hypothesis of the previous theorem $`M`$ is seminormal if and only if $`H_𝔪^i(K[M])_a=0`$ for all $`i`$ and all $`a\mathrm{G}(M)`$ such that $`a\mathrm{C}(M)`$. Note that this result has a variant for the normalization, discussed in Remark 4.8.
Using further methods from commutative algebra we prove in Theorem 4.9:
###### Theorem.
Let $`M^m`$ be a positive affine monoid of rank $`d`$ such that $`K[M]`$ satisfies $`(S_2)`$ and $`H_𝔪^d(K[M])_a=0`$ for all $`a\mathrm{G}(M)\left(\mathrm{C}(M)\right)`$. Then $`M`$ is seminormal.
Since there are Cohen–Macaulay affine monoid rings which are not seminormal (like $`K[t^2,t^3]`$), one can not omit the assumption about the vanishing of the graded components of $`H_𝔪^d(K[M])`$ outside $`\mathrm{C}(M)`$.
In Section 5 we define the numbers
$`c_K(M)`$ $`=sup\{i:K[MF]\text{ is Cohen–Macaulay for all faces }F,dimFi\},`$
$`n(M)`$ $`=sup\{i:MF\text{ is normal for all faces }F,dimFi\}.`$
By Hochster’s theorem on normal monoids we have that $`c_K(M)n(M)`$. The main result 5.3 in Section 5 are the inequalities
$$\mathrm{depth}Rc_K(M)\mathrm{min}\{n(M)+1,d\}$$
for an affine seminormal monoid $`M^m`$ of rank $`d`$. Since seminormal monoids of rank $`1`$ are normal, we immediately get that $`\mathrm{depth}K[M]2`$ if $`d2`$. In particular, $`K[M]`$ is Cohen–Macaulay if $`\mathrm{rank}M=2`$.
One obtains a satisfactory result also for $`\mathrm{rank}M=3`$, which was already shown in by different methods. In fact, in Corollary 5.6 we prove that $`K[M]`$ is Cohen–Macaulay for a positive affine seminormal monoid $`M^m`$ with $`\mathrm{rank}M3`$ if $`K[M]`$ satisfies $`(S_2)`$. One could hope that $`K[M]`$ is always Cohen–Macaulay if $`M`$ is seminormal and $`K[M]`$ satisfies ($`S_2`$). But this is not the case and we present a counterexample in 7.1. The best possible result is given in 5.6: ($`S_2`$) is sufficient for $`K[M]`$ to be Cohen-Macaulay in the seminormal case if the cross-section of $`\mathrm{C}(M)`$ is a simple polytope.
In Section 6 we study the seminormality of affine monoid rings with characteristic $`p`$ methods. The main observation is that a positive affine monoid $`M^m`$ is seminormal if there exists a field $`K`$ of characteristic $`p`$ such that $`K[M]`$ is $`F`$-injective. This fact is a consequence of our analysis of the local cohomology groups of positive affine monoid rings. In 6.2 we give a precise description for which prime numbers $`p`$ and fields of characteristic $`p`$, we have that $`K[M]`$ is $`F`$-injective. Implicitly, this result was already observed by Hochster and Roberts in \[10, Theorem 5.33\]. In fact, if $`M^m`$ is a positive affine seminormal monoid and $`K`$ is a field of characteristic $`p>0`$, then the following statements are equivalent:
1. The prime ideal $`(p)`$ is not associated to the $``$-module $`\mathrm{G}(M)F/\mathrm{G}(MF)`$ for any face $`F`$ of $`\mathrm{C}(M)`$;
2. $`R`$ is $`F`$-split;
3. $`R`$ is $`F`$-pure;
4. $`R`$ is $`F`$-injective.
As a direct consequence we obtain that $`M`$ is normal if the equivalent statements hold for every field $`K`$ of characteristic $`p>0`$.
In the last section we present examples and counterexamples related to the results of this paper. In particular, we will show that for every simplicial complex $`\mathrm{\Delta }`$ there exists a seminormal affine monoid $`M`$ such that the only obstruction to the Cohen-Macaulay property of $`K[M]`$ is exactly the simplicial homology of $`\mathrm{\Delta }`$. Choosing $`\mathrm{\Delta }`$ as a triangulation of the real projective plane we obtain an example whose Cohen-Macaulay property depends on $`K`$. A similar result was proved by Trung and Hoa . Our construction has the advantage of yielding a seminormal monoid $`M`$, and is geometrically very transparent.
We are grateful to Aldo Conca for directing our attention toward the results of Hochster and Roberts in .
## 2. Prerequisites
We recall some facts from convex geometry. Let $`X`$ be a subset of $`^m`$. The convex hull $`\mathrm{conv}(X)`$ of $`X`$ is the set of convex combinations of elements of $`X`$. Similarly, the set $`\mathrm{C}(X)`$ of positive linear combinations of elements of $`X`$ is called the cone generated by $`X`$. By convention $`\mathrm{C}(\mathrm{})=\{0\}`$ and $`\mathrm{conv}(\mathrm{})=\mathrm{}`$. A cone $`C`$ is called positive (or pointed) if $`0`$ is the only invertible element in $`C`$. To an affine form $`\alpha `$ on $`^m`$ (i. e. a polynomial of degree $`1`$) we associate the affine hyperplane $`H_\alpha =\alpha ^1(0)`$, the closed half-space $`H_\alpha ^+=\alpha ^1([0,\mathrm{}))`$, and the open half-space $`H_\alpha ^>=\alpha ^1(0,\mathrm{})`$. An intersection $`P=_{i=1}^nH_{\alpha _i}^+`$ of finitely many closed half-spaces is called a polyhedron. A (proper) face of a polyhedron $`P`$ is the (proper) intersection of $`P`$ with a hyperplane $`H_\beta `$ such that $`PH_\beta ^+`$. Also $`P`$ is considered as a face of itself. A facet is a maximal proper face. Recall that there are only a finite number of faces. A polytope is a bounded polyhedron. The set $`\mathrm{conv}(F)`$ is a polytope for every finite subset $`F`$ of $`^m`$, and every polytope is of this form. A cone is a finite intersection of half-spaces of the form $`H_{\alpha _i}^+`$ where the $`\alpha _i`$ are linear forms (i. e. homogeneous polynomials of degree $`1`$). The set $`\mathrm{C}(F)`$ is a cone for every finite subset $`F`$ of $`^m`$, and every cone is of this form. Let $`P`$ be a polyhedron and $`F`$ a face of $`P`$. Then we denote the relative interior of $`F`$ with respect to the subspace topology on the affine hull of $`F`$ by $`\mathrm{int}F`$. Note that $`P`$ decomposes into the disjoint union $`\mathrm{int}F`$ of the (relative) interiors of its faces. For more details on convex geometry we refer to the books of Bruns and Gubeladze , Schrijver and Ziegler .
An affine monoid $`M`$ is a finitely generated commutative monoid which can be embedded into $`^m`$ for some $`m`$. We always use $`+`$ for the monoid operation. In the literature $`M`$ is also called an affine semigroup in this situation. We call $`M`$ positive if $`0`$ is the only invertible element in $`M`$. Observe that $`M`$ is positive if and only if $`\mathrm{C}(M)`$ is pointed.
Let $`K`$ be a field and $`K[M]`$ be the $`K`$-vector space with $`K`$-basis $`X^a`$, $`aM`$. The multiplication $`X^aX^b=X^{a+b}`$ for $`a,bM`$ induces a ring structure on $`K[M]`$ and this $`K`$-algebra is called the affine monoid ring (or algebra) associated to $`M`$. The embedding of $`M`$ into $`^m`$ induces an embedding of $`K[M]`$ into the Laurent polynomial ring $`K[^m]=K[X_i^{\pm 1}:i=1,\mathrm{},m]`$ where $`X_i`$ corresponds to the $`i`$th element of the canonical basis of $`^m`$. Note that $`K[M]`$ is a $`^m`$-graded $`K`$-algebra with the property that $`dim_KK[M]_a1`$ for all $`a^m`$. It is easy to determine the $`^m`$-graded prime ideals of $`K[M]`$. In fact every $`^m`$-graded prime ideal is of the form $`𝔭_F=(X^a:aM,aF)`$ for a unique face $`F`$ of $`\mathrm{C}(M)`$ (see \[4, Theorem 6.1.7\] for a proof). In particular, the prime ideals of height 1 correspond to the facets of $`\mathrm{C}(M)`$.
Recall that a Noetherian domain $`R`$ is normal if it is integrally closed in its field of fractions. The normalization $`\overline{R}`$ of $`R`$ is the set of elements in the quotient field of $`R`$ which are integral over $`R`$. An affine monoid $`M`$ is called normal, if $`z\mathrm{G}(M)`$ and $`mzM`$ for some $`m`$ imply $`zM`$. Here $`\mathrm{G}(M)`$ is the group generated by $`M`$. It is easy to see that $`M`$ is normal if and only if $`M=\mathrm{G}(M)\mathrm{C}(M)`$. If $`M`$ is an arbitrary submonoid of $`^m`$, then its normalization is the monoid $`\overline{M}=\{z\mathrm{G}(M):mzM\text{ for some }m\}`$. By Gordan’s lemma $`\overline{M}`$ is affine for an affine monoid $`M`$.
Hochster proved that $`K[M]`$ is normal if and only if $`M`$ is a normal monoid, in fact we have that $`\overline{K[M]}=K[\overline{M}]`$. In particular, Hochster showed that if $`M`$ is normal, then $`K[M]`$ is a Cohen–Macaulay ring. One can characterize the Cohen–Macaulay property of $`K[M]`$ in terms of combinatorial and topological information associated to $`M`$. This amounts to an analysis of the $`^m`$-graded structure of the local cohomology of $`K[M]`$; see Trung and Hoa for a criterion of this type.
A Noetherian domain $`R`$ is called seminormal if for an element $`x`$ in the quotient field $`Q(R)`$ of $`R`$ such that $`x^2,x^3R`$ we have $`xR`$. The seminormalization $`^+R`$ of $`R`$ is the intersection of all seminormal subrings $`S`$ such that $`RSQ(R)`$. An affine monoid $`M`$ is called seminormal, if $`z\mathrm{G}(M)`$, $`z^2M`$ and $`z^3M`$ imply that $`zM`$. The seminormalization $`^+M`$ of $`M`$ is the intersection of all seminormal monoids $`N`$ such that $`MN\mathrm{G}(M)`$. It can be shown that $`^+M`$ is again an affine monoid. Hochster and Roberts \[10, Proposition 5.32\] proved that an affine monoid $`M`$ is seminormal if and only if $`K[M]`$ is a seminormal ring. We frequently use the following characterization of seminormal monoids. See Reid and Roberts \[12, Theorem 4.3\] for a proof for positive monoids, but this proof works also for arbitrary affine monoids.
###### Theorem 2.1.
Let $`M`$ be an affine monoid $`M^m`$. Then
$$^+M=\underset{F\text{ }\text{face of }\mathrm{C}(M)}{}\mathrm{G}(MF)\mathrm{int}F.$$
In particular, $`M`$ is seminormal if and only if it equals the right hand side of the equality.
## 3. Seminormality and Serre’s condition $`(S_2)`$
Let $`R`$ be a Noetherian ring and let $`N`$ be a finitely generated $`R`$-module. Recall that $`N`$ satisfies Serre’s condition ($`S_k`$) if
$$\mathrm{depth}N_𝔭\mathrm{min}\{k,dimN_𝔭\}$$
for all $`𝔭\mathrm{Spec}R`$. Affine monoid rings trivially satisfy ($`S_1`$), since they are integral domains. We are interested in characterizing $`(S_2)`$ for affine monoid rings.
While the validity of $`(S_k)`$ in $`K[M]`$ may depend on the field $`K`$ for $`k>2`$, $`(S_2)`$ can be characterized solely in terms of $`M`$, as was shown in .
Let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$ and let
$$M_i=\{a\mathrm{G}(M):a+bM\text{ for some }bMF_i\}$$
for $`i=1,\mathrm{},t`$. Note that the elements of $`M_i`$ correspond to the monomials in the homogeneous localization $`K[M]_{(𝔭_{F_i})}`$. We set
$$M^{}=\underset{i=1}{\overset{t}{}}M_i.$$
###### Proposition 3.1.
Let $`M`$ be an affine monoid $`M^m`$, and $`K`$ a field. Then the following statements are equivalent:
1. $`K[M]`$ satisfies $`(S_2)`$;
2. $`M=M^{}`$.
Observe that $`(M^{})^{}=M^{}`$. Thus $`K[M^{}]`$ always satisfies $`(S_2)`$. For seminormal monoids $`M`$ the equality $`M=M^{}`$ can be expressed in terms of the lattices $`\mathrm{G}(MF)`$ as we will see in Corollary 3.4. First we describe the monoids $`M_i`$ under a slightly weaker condition.
###### Lemma 3.2 (Proposition 4.2.6 in ).
Let $`M^m`$ be an affine monoid and let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$ with defining linear forms $`\alpha _1,\mathrm{},\alpha _t`$. Then:
1. $`M_iH_{\alpha _i}=\mathrm{G}(MF_i).`$
2. If $`\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M`$, then
$$M_i=\left(\mathrm{G}(M)H_{\alpha _i}^>\right)\mathrm{G}(MF_i).$$
###### Proof.
Every element in $`M_i`$ is of the form $`c=ab`$ for some $`aM`$ and $`bMF_i`$. Hence $`\alpha _i(c)0`$ with equality if and only if $`aMF_i`$. It follows that
$$M_i\left(\mathrm{G}(M)H_{\alpha _i}^>\right)\mathrm{G}(MF_i)\text{and}M_iH_{\alpha _i}\mathrm{G}(MF_i).$$
If $`c\mathrm{G}(MF_i)`$ and $`c=ab`$ for some $`a,bMF_i`$, then clearly by the definition of $`M_i`$ we have that $`cM_i`$. Thus we see that $`M_iH_{\alpha _i}=\mathrm{G}(MF_i)`$.
For (ii) it remains to show that if $`c\mathrm{G}(M)H_{\alpha _i}^>`$, then $`cM_i`$. Pick $`dM\mathrm{int}F_i`$ such that $`\alpha _j(c+d)>0`$ for all $`ji`$. Hence $`c+d\mathrm{int}\mathrm{C}(M)`$. But
$$c+d\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M,$$
by the additional assumption in (ii). Thus $`cM_i`$. ∎
In the following proposition we consider $`\mathrm{C}(M)`$ as a face of itself.
###### Proposition 3.3 (Proposition 4.2.7 in ).
Let $`M`$ be an affine monoid $`M^m`$ and let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$. If $`\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M`$, then
$$M^{}=\underset{F\text{ }\text{face of }\mathrm{C}(M)}{}[\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}F]$$
with the convention that $`_{FF_i}\mathrm{G}(MF_i)=\mathrm{G}(M)`$ if $`F=\mathrm{C}(M)`$.
###### Proof.
We apply 3.2 several times. By assumption we have
$$\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M^{}.$$
Let $`F`$ be a proper face of $`\mathrm{C}(M)`$. Choose a facet $`F_j`$ with defining linear form $`\alpha _j`$. Either $`\mathrm{int}FF_j`$ and thus
$$\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}F\mathrm{G}(MF_j)\mathrm{int}FM_j,$$
or $`\mathrm{int}F`$ is contained in $`\mathrm{C}(M)H_{\alpha _j}^>`$ and
$$\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}F\mathrm{G}(M)H_{\alpha _j}^>M_j.$$
Hence
$$\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}FM^{}.$$
Note that
$$M^{}\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)\mathrm{int}\mathrm{C}(M).$$
For a proper face $`F`$ of $`\mathrm{C}(M)`$ it follows from 3.2 that
$$M^{}\mathrm{int}F\underset{FF_i}{}M_i\mathrm{int}F=\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}F.$$
All in all we see that
$$M^{}=\underset{F\text{ face of }\mathrm{C}(M)}{}[\underset{FF_i}{}\mathrm{G}(MF_i)\mathrm{int}F].\mathit{}$$
The equivalence of part (ii) and (iii) in the following corollary was shown in Theorem 4.2.14 in .
###### Corollary 3.4.
Let $`M^m`$ be an affine monoid, $`K`$ a field, and let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$. If $`\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M`$, then the following statements are equivalent:
1. $`K[M]`$ satisfies $`(S_2)`$;
2. $`M=M^{}`$;
3. For all proper faces $`F`$ of $`\mathrm{C}(M)`$ one has
$$M\mathrm{int}F=\underset{FF_j}{}\mathrm{G}(MF_j)\mathrm{int}F.$$
If $`M`$ is seminormal, then the following is equivalent to (i) – (iii):
1. $`\mathrm{G}(MF)=_{FF_j}\mathrm{G}(MF_j)`$.
###### Proof.
The equivalence of (i) and (ii) was already stated in 3.1. The equivalence of (ii) and (iii) is an immediate consequence of 3.3.
In the seminormal case one has $`\mathrm{G}(FM)\mathrm{int}FM`$ for all faces, so that (iv) implies (iii). For the converse implication one uses the fact that $`\mathrm{G}(MF)`$ is generated by its elements in $`\mathrm{int}(F)`$ (see Bruns and Gubeladze ). ∎
###### Remark 3.5.
If $`M`$ is seminormal, then we know from 2.1, that $`\mathrm{G}(M)\mathrm{int}\mathrm{C}(M)M`$. Hence we can apply 3.2, 3.3 and 3.4 in this situation.
The corollary shows that a seminormal monoid satisfies $`(S_2)`$ if and only if the restriction of the groups $`\mathrm{G}(MF)`$ happens only in the passage from $`\mathrm{C}(M)`$ to its facets.
## 4. Local cohomology of monoid rings
For the rest of the paper $`K`$ always denotes a field, and $`M^m`$ is an affine positive monoid of rank $`d`$. Recall that the seminormalization of $`M`$ is
$$^+M=\underset{F\text{ face of }\mathrm{C}(M)}{}\mathrm{G}(MF)\mathrm{int}F.$$
and the normalization of $`M`$ is
$$\overline{M}=\mathrm{G}(M)\mathrm{C}(M).$$
In this section we want to compute the local cohomology of $`K[M]`$ and compare it with the local cohomology of $`K[^+M]`$ and $`K[\overline{M}]`$.
If $`M`$ is a positive affine monoid, then $`K[M]`$ is a $`^m`$-graded $`K`$-algebra with a unique graded maximal ideal $`𝔪`$ generated by all homogeneous elements of nonzero degree. By the local cohomology of $`K[M]`$ we always mean the local cohomology groups $`H_𝔪^i(K[M])`$. Observe that $`^+M`$ and $`\overline{M}`$ are also positive affine monoids. Since the $`K`$-algebras $`K[^+M]`$ and $`K[\overline{M}]`$ are finitely generated modules over $`K[M]`$ and the extensions of $`𝔪`$ are primary to their maximal ideals, the local cohomology groups of $`K[^+M]`$ and $`K[\overline{M}]`$ coincide with $`H_𝔪^i(K[^+M])`$ and $`H_𝔪^i(K[\overline{M}])`$ respectively. Because of this fact and to avoid cumbersome notation we always write $`H_𝔪^i(K[^+M])`$ and $`H_𝔪^i(K[\overline{M}])`$ for the local cohomology of $`K[^+M]`$ and $`K[\overline{M}]`$. The same applies to $`^m`$-graded residue class rings of $`K[M]`$.
In the following $`R`$ will always denote the ring $`K[M]`$, and thus $`^+R`$ and $`\overline{R}`$ will stand for $`K[^+M]`$ and $`K[\overline{M}]`$, respectively.
Let $`F`$ be a proper face of $`\mathrm{C}(M)`$. Then $`𝔭_F=(X^a:aM,aF)`$ is a monomial prime ideal of $`R`$, and conversely, if $`𝔭`$ is a monomial prime ideal, then $`F(𝔭)=_+\{aM:X^a𝔭\}`$ is a proper face of $`\mathrm{C}(M)`$. These two assignments set up a bijective correspondence between the monomial prime ideals of $`R`$ and the proper faces of $`\mathrm{C}(M)`$.
Note that the natural embedding $`K[MF]K[M]`$ is split by the face projection $`K[M]K[MF]`$ that sends all monomials in $`F`$ to themselves and all other monomials to $`0`$. Its kernel is $`𝔭_F`$. Therefore we have a natural isomorphism $`K[M]/𝔭_FK[MF]`$.
The next lemma states a crucial fact for the analysis of the local cohomology of $`R`$. For this lemma and its proof we need the following notation. For $`W^m`$ we define
$$W=\{a:aW\}.$$
For a $`^m`$-graded local Noetherian $`K`$-algebra $`R`$ with $`R_0=K`$ (like the monoid ring $`K[M]`$ for a positive affine monoid $`M`$) and a $`^m`$-graded $`R`$-module $`N`$ we set
$$N^{}=\mathrm{Hom}_K(N,K).$$
(Here we mean by $`\mathrm{Hom}_K(N,K)`$ the homogenous homomorphisms from $`N`$ to $`K`$.) Note that $`N^{}`$ is again a $`^m`$-graded $`R`$-module by setting
$$(N^{})_a=\mathrm{Hom}_K(N_a,K)\text{ for }a^m.$$
###### Lemma 4.1.
Let $`M^m`$ be a positive affine monoid of rank $`d`$. The $`R`$-module $`\overline{\omega }`$ of $`R`$ generated by the monomials $`X^b`$ with $`b\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)`$ is the canonical module of the normalization $`\overline{R}`$. If $`a\mathrm{G}(M)`$, then
$$H_𝔪^i(\overline{\omega })_a\{\begin{array}{cc}0\hfill & \text{if }i<d\text{ or }a\mathrm{C}(M),\hfill \\ K\hfill & \text{if }i=d\text{ and }a\mathrm{C}(M).\hfill \end{array}$$
###### Proof.
Danilov and Stanley showed that $`\overline{\omega }`$ is the canonical module of $`\overline{R}`$ (see Bruns and Herzog \[4, Theorem 6.3.5\] or Stanley ). Thus $`\overline{\omega }`$ is Cohen–Macaulay of dimension $`d`$. This implies $`H_𝔪^i(\overline{\omega })=0`$ for $`i<d`$. Furthermore, by graded local duality we have that
$$H_𝔪^d(\overline{\omega })^{}\mathrm{Hom}_{\overline{R}}(\overline{\omega },\overline{\omega })\overline{R}$$
as $`^m`$-graded modules. This concludes the proof. ∎
For the central proofs in this paper it is useful to extend the correspondence between the faces of $`\mathrm{C}(M)`$ and the monomial prime ideals of $`R`$ to a bijection between the unions of faces of $`\mathrm{C}(M)`$ and the monomial radical ideals. If $`𝔮`$ is a monomial radical ideal, then we let $`F(𝔮)`$ denote the union of the faces $`F(𝔭)`$ such that $`𝔭𝔮`$, and if $`F`$ is the union of faces, then the corresponding radical ideal $`𝔮_F`$ of $`R`$ is just the intersection of all monomial prime ideals $`𝔭_G`$ such that $`GF`$. We need the following lemma about monomial prime ideals of $`R`$.
###### Lemma 4.2.
Let $`M^m`$ be an affine monoid and $`F_1,\mathrm{},F_t,G`$ be faces of $`\mathrm{C}(M)`$. Then:
1. $`𝔭_{F_1\mathrm{}F_t}=𝔭_{F_1}+\mathrm{}+𝔭_{F_t}`$;
2. $`𝔭_G+_{i=1}^t𝔭_{F_i}=_{i=1}^t(𝔭_G+𝔭_{F_i})`$.
###### Proof.
Observe that the $`𝔭_{F_i}`$ are $`^m`$-graded ideals of $`R`$, i. e. they are monomial ideals in this ring. In other words, their bases as $`K`$-vector spaces are subsets of the set of monomials $`X^a`$, $`aM`$. Using this fact it is easy to check the equalities claimed. ∎
We are ready to prove a vanishing result for the local cohomology of seminormal monoid rings. Hochster and Roberts \[10, Remark 5.34\] already noticed that certain “positive” graded components of $`H_𝔪^i(R)`$ vanish for a seminormal monoid. We can prove a much more precise statement.
###### Theorem 4.3.
Let $`M^m`$ be a positive affine seminormal monoid and $`R=K[M]`$. If $`H_𝔪^i(R)_a0`$ for $`a\mathrm{G}(M)`$, then $`a\overline{MF}`$ for a face $`F`$ of $`\mathrm{C}(M)`$ of dimension $`i`$. In particular,
$$H_𝔪^i(R)_a=0\text{if }a\mathrm{C}(M).$$
###### Proof.
The assertion is trivial for $`\mathrm{rank}M=0`$. Thus assume that $`d=\mathrm{rank}M>0`$. Since $`M`$ is seminormal, $`\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)`$ is contained in $`M`$. Thus $`\overline{\omega }`$, which as a $`K`$-vector space is generated by the monomials $`X^a`$ with $`a\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)`$, is an ideal of $`R`$. Now consider the exact sequence
$$0\overline{\omega }RR/\overline{\omega }0.$$
By Lemma 4.1 the long exact local cohomology sequences splits into isomorphisms
(1) $`H_𝔪^i(R)H_𝔪^i(R/\overline{\omega })\text{for }i<d1`$
and the exact sequence
(2) $`0H_𝔪^{d1}(R)H_𝔪^{d1}(R/\overline{\omega })H_𝔪^d(\overline{\omega })H_𝔪^d(R)0.`$
The local cohomology of $`\overline{\omega }`$ has been determined in 4.1. Thus
$$H_𝔪^d(R)_a=0\text{for }a\mathrm{C}(M).$$
This takes care of the top local cohomology. For the lower cohomologies we note that
$$R/\overline{\omega }R/\underset{i=1}{\overset{t}{}}𝔭_{F_i}$$
where $`F_1,\mathrm{},F_t`$ are the facets of $`\mathrm{C}(M)`$.
Therefore it is enough to prove the following statement which generalizes the theorem: *let $`𝔮`$ be a monomial radical ideal of $`R`$; if $`H_𝔪^i(R/𝔮)_a0`$, then $`a\overline{MG}`$ for a face $`GF(𝔮)`$ of dimension $`i`$*.
The case $`𝔮=(0)`$ has already been reduced to the case $`𝔮=\overline{\omega }`$. So we can assume that $`𝔮(0)`$ and use induction on $`\mathrm{rank}M`$ and on the number $`t`$ of minimal monomial prime ideals $`𝔭_1,\mathrm{},𝔭_t`$ of $`𝔮`$.
If $`t=1`$, then $`R/𝔭_1\mathrm{}K[MF(𝔭_1)]`$. Now we can apply induction on $`\mathrm{rank}M`$. Let $`t>1`$. We set $`𝔮^{}=_{j=1}^{t1}𝔭_j`$. Then we have the standard exact sequence
$$0R/𝔮R/𝔮^{}R/𝔭_tR/(𝔮^{}+𝔭_t)0.$$
The local cohomologies of $`R/𝔮^{}`$ and $`R^{}=R/𝔭_t`$ are under control by induction. But this applies to $`R/(𝔮^{}+𝔭_t)`$, too. In fact, by Lemma 4.2 one has
$$R/(𝔮^{}+𝔭_t)R/\underset{j=1}{\overset{t1}{}}(𝔭_t+𝔭_j)\mathrm{}R^{}/𝔮^{\prime \prime }$$
where $`𝔮^{\prime \prime }=_{j=1}^{t1}\left((𝔭_t+𝔭_j)/𝔭_t\right)`$. Thus $`𝔮^{\prime \prime }`$ is a monomial radical ideal of $`R^{}K[MF(𝔭_t)]`$. Now it is enough to apply the long exact cohomology sequence
$$\mathrm{}H_𝔪^{i1}(R^{}/𝔮^{\prime \prime })H_𝔪^i(R/𝔮)H_𝔪^i(R/𝔮^{})H_𝔪^i(R^{})\mathrm{}\mathit{}$$
We describe a complex which computes the local cohomology of $`R`$. Writing $`R_F`$ for the homogeneous localization $`R_{(𝔭_F)}`$, let
$$L^\text{.}(M):0L^0(M)\mathrm{}L^t(M)\mathrm{}L^d(M)0$$
be the complex with
$$L^t(M)=\underset{F\text{ face of }\mathrm{C}(M),dimF=t}{}R_F$$
and the differential $`:L^{t1}(M)L^t(M)`$ induced by
$$_{G,F}:R_GR_F\text{ to be }\{\begin{array}{cc}0\hfill & \text{if }GF,\hfill \\ ϵ(G,F)\mathrm{nat}\hfill & \text{if }GF,\hfill \end{array}$$
where $`ϵ`$ is a fixed incidence function on the face lattice of $`\mathrm{C}(M)`$ in the sense of \[4, Section 6.2\]. In \[4, Theorem 6.2.5\] it was shown that $`L^\text{.}(M)`$ is indeed a complex and that for an $`R`$-module $`N`$ we have that
$$H_𝔪^i(N)=H^i(L^\text{.}(M)\mathrm{}_RN)\text{ for all }i0.$$
Next we construct another, “smaller” complex which will be especially useful for the computation of the local cohomology of $`R`$ if $`M`$ is seminormal. Let
$$^+L^\text{.}(M):0^+L^0(M)\mathrm{}^+L^t(M)\mathrm{}^+L^d(M)0$$
be the complex with
$$^+L^t(M)=\underset{F\text{ face of }\mathrm{C}(M),dimF=t}{}K[\overline{MF}]$$
and the differential $`^+:^+L^{t1}(M)^+L^t(M)`$ is induced by the same rule as $``$ above.
###### Proposition 4.4.
Let $`M^m`$ be a positive affine monoid, $`a\mathrm{C}(M)\mathrm{G}(M)`$. Then
$$H_𝔪^i(R)_a=H^i(^+L(M))_a$$
###### Proof.
We know that $`H_𝔪^i(R)_a`$ is the cohomology of the complex $`L(M)_a`$. To prove the claim it suffices to determine $`(R_F)_a`$ for a face $`F`$ of $`\mathrm{C}(M)`$ if $`a\mathrm{C}(M)\mathrm{G}(M)`$. It is an easy exercise to show that
$$(R_F)_a=K[\overline{MF}]_a$$
where one has to use the fact that $`\overline{MF}=\mathrm{G}(MF)F`$. ∎
It follows from 4.3 and 4.4 that the local cohomology of $`^+R`$ is a direct summand of the local cohomology of $`R`$ as a $`K`$-vector space.
###### Corollary 4.5.
Let $`M^m`$ be a positive affine monoid. Then
$$\underset{a\mathrm{C}(M)\mathrm{G}(M)}{}H_𝔪^i(R)_a\underset{a\mathrm{C}(M)\mathrm{G}(M)}{}H_𝔪^i(^+R)_a=H_𝔪^i(^+R).$$
###### Corollary 4.6.
Let $`M^m`$ be a positive affine monoid of rank $`d`$. Then:
1. If $`R`$ is Cohen–Macaulay, then $`^+R`$ is Cohen–Macaulay.
2. If $`\mathrm{depth}Rk`$, then $`\mathrm{depth}^+Rk`$.
3. If $`R`$ satisfies ($`S_k`$), then $`^+R`$ satisfies ($`S_k`$).
###### Proof.
It is well-known that the Cohen–Macaulay property and depth can be read off the local cohomology groups. This is also true for Serre’s property ($`S_k`$) since we have that $`R`$ satisfies ($`S_k`$) if and only if $`dimH_𝔪^j(R)^{}jk`$ for $`j=0,\mathrm{},dimR1`$ and an analogous characterization of Serre’s property ($`S_k`$) for $`^+R`$. (See Schenzel for a proof of the latter fact.) ∎
The results of this section allow us to give a cohomological characterization of seminormality for positive monoid rings.
###### Theorem 4.7.
Let $`M^m`$ be a positive affine monoid. Then the following statements are equivalent:
1. $`M`$ is seminormal;
2. $`H_𝔪^i(R)_a=0`$ for all $`i`$ and all $`a\mathrm{G}(M)`$ such that $`a\mathrm{C}(M)`$.
###### Proof.
Consider the sequence
$$0R^+R^+R/R0$$
of finitely generated $`^m`$-graded $`R`$-modules. Observe that $`H_𝔪^i(R)_aH_𝔪^i(^+R)_a`$ for $`a\mathrm{C}(M)`$. Thus it follows from the long exact cohomology sequence
$$\mathrm{}H_𝔪^i(R)H_𝔪^i(^+R)H_𝔪^i(^+R/R)\mathrm{}$$
that $`H_𝔪^i(R)_a=0`$ for $`a\mathrm{C}(M)`$ and all $`i`$ if and only if $`H_𝔪^i(^+R/R)_a=0`$ for all $`a`$ and all $`i`$. This is equivalent to $`^+R/R=0`$. Hence $`R`$ and, thus, $`M`$ are seminormal. ∎
###### Remark 4.8.
The previous results have variants for the normalization. If we restrict the direct sum in Corollary 4.5 to those $`a`$ that belong to $`\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)`$ then the local cohomology of $`K[^+M]`$ must be replaced by that of $`\overline{R}`$. Moreover, the local cohomology of $`R`$ vanishes in all degrees $`a`$ outside $`\mathrm{int}\mathrm{C}(M)`$ if and only if $`M`$ is normal.
This follows by completely analogous arguments since we have that $`H_𝔪^i(\overline{R})=0`$ for $`i<d`$ and $`H_𝔪^d(\overline{R})_a0`$ if and only if $`a\mathrm{int}\mathrm{C}(M)\mathrm{G}(M)`$.
With different methods than those used so far we can prove another seminormality criterion. It involves only the top local cohomology group, but needs a stronger hypothesis on $`M`$.
###### Theorem 4.9.
Let $`M^m`$ be a positive affine monoid of rank $`d`$. If $`R`$ satisfies $`(S_2)`$ and $`H_𝔪^d(R)_a=0`$ for all $`a\mathrm{G}(M)\left(\mathrm{C}(M)\right)`$, then $`M`$ is seminormal.
###### Proof.
The assumption and 4.5 imply that the $`d`$th local cohomology of $`R`$ and $`^+R`$ coincide as $`R`$-modules. Since $`M`$ and therefore $`^+M`$ are positive, there exists a $``$-grading on $`R`$ and $`^+R`$ such that both $`K`$-algebras are generated in positive degrees. We choose a common Noether normalization $`S`$ of $`R`$ and $`^+R`$ with respect to this $``$-grading.
Since $`R`$ satisfies $`(S_2)`$ it is a reflexive $`S`$-module. By 4.6 also $`^+R`$ satisfies $`(S_2)`$ and is a reflexive $`S`$-module. In the following let $`\omega _S`$ be the canonical module of $`S`$ which is in our situation just a shifted copy of $`S`$ with respect to the $``$-grading. By graded local duality and reflexivity we get the following chain of isomorphisms of graded $`S`$-modules:
$`R`$ $`\mathrm{Hom}_S(\mathrm{Hom}_S(R,S),S)\mathrm{Hom}_S(\mathrm{Hom}_S(R,\omega _S),\omega _S)`$
$`\mathrm{Hom}_S(H_𝔪^d(R)^{},\omega _S)\mathrm{Hom}_S(H_𝔪^d(^+R)^{},\omega _S)`$
$`\mathrm{Hom}_S(\mathrm{Hom}_S(^+R,\omega _S),\omega _S)\mathrm{Hom}_S(\mathrm{Hom}_S(^+R,S),S)`$
$`^+R.`$
Hence $`M=^+M`$ is seminormal. ∎
Again we can obtain a similar normality criterion if we replace $`\mathrm{C}(M)`$ by $`\mathrm{int}\mathrm{C}(M)`$ in Theorem 4.9. In the rest of this section we further analyze the local cohomologies of $`R`$.
###### Proposition 4.10.
Let $`M`$ be seminormal and $`R=K[M]`$. Then:
1. $`H_𝔪^d(R)_a0`$ (and so of $`K`$-dimension $`1`$) $``$ $`a\overline{M}_F\mathrm{G}(MF)`$ where $`F`$ runs through the facets of $`\mathrm{C}(M)`$.
2. $`H_𝔪^{d1}(R)_a0`$ $``$ $`a\mathrm{C}(M)_F\mathrm{G}(MF)`$ and $`dim_KH_𝔪^{d1}(R/\overline{\omega })_a2`$ where $`F`$ again runs through the facets of $`\mathrm{C}(M)`$.
3. $`R`$ is Cohen–Macaulay $``$ $`R/\overline{\omega }`$ is Cohen–Macaulay and $`dim_KH_𝔪^{d1}(R/\overline{\omega })_a1`$ for all $`a\mathrm{G}(M)`$.
###### Proof.
$`H_𝔪^d(R)`$ is the cokernel of the map
$$\underset{F\text{ facet of }\mathrm{C}(M)}{}K[\overline{MF}]K[\mathrm{C}(M)\mathrm{G}(M)],$$
which implies (i).
(ii) We have the exact sequence (2)
$$0H_𝔪^{d1}(R)H_𝔪^{d1}(R/\overline{\omega })H_𝔪^d(\overline{\omega })\stackrel{𝜋}{}H_𝔪^d(R)0.$$
Lemma 4.1 and (i) yield the $`^m`$-graded structure of $`\mathrm{Ker}\pi `$: its nonzero graded components have dimension $`1`$ and live in exactly the degrees $`a`$ with $`a\mathrm{C}(M)_F\mathrm{G}(MF)`$.
On the other hand, $`H_𝔪^{d1}(R/\overline{\omega })`$ can have non-zero components only in these degrees, as follows from the generalization of Theorem 4.3 stated in its proof. Thus $`H_𝔪^{d1}(R)`$ is limited to these degrees and is non-zero at $`a`$ if and only if $`dim_KH_𝔪^{d1}(R/\overline{\omega })_a2`$.
(iii) The isomorphisms $`H_𝔪^i(R)H_𝔪^i(R/\overline{\omega })`$ for $`i<d1`$ reduce the claim immediately to (ii). ∎
Having computed the $`d`$-th local cohomology of $`K[M]`$, we can easily describe the Gorenstein property of $`K[M]`$ in combinatorial terms:
###### Corollary 4.11.
Let $`M`$ be seminormal and $`R=K[M]`$ Cohen–Macaulay. For each facet $`F`$ of $`\mathrm{C}(M)`$ let $`\gamma _F`$ denote the index of the group extension $`\mathrm{G}(MF)\mathrm{G}(M)F`$, and $`\sigma _F`$ the unique $``$-linear form on $`\mathrm{G}(M)`$ such that $`\sigma _F(\overline{M})=_+`$ and $`\sigma _F(x)=0`$ for all $`xF`$. Then the following are equivalent:
1. $`R`$ is Gorenstein;
2. 1. $`\gamma _F2`$ for all facets $`F`$ of $`\mathrm{C}(M)`$;
2. there exists $`b\overline{M}`$ such that
$`bF\mathrm{G}(MF),`$ $`\text{if }\gamma _F=2,`$
$`\sigma _F(b)=1,`$ $`\text{else}.`$
###### Proof.
The multigraded support of the $`K`$-dual $`\omega _R`$ of $`H_𝔪^d(R)`$ is $`N=\overline{M}_F\mathrm{G}(MF)`$, and $`R`$ is Gorenstein if and only if there exists $`b\overline{M}`$ such that $`N=b+M`$, or, equivalently, $`\omega _R=RX^b`$. It remains to be shown that such $`b`$ exists if and only if the conditions in (ii) are satisfied. We leave the exact verification to the reader. (Note that for each facet $`F`$ there exists $`c\overline{M}\mathrm{int}\mathrm{C}(M)M`$ such that $`\sigma _F(x)=1`$.) ∎
Finally, we give an interpretation of $`H_𝔪^i(R)_a`$ for $`a\mathrm{C}(M)`$ which will be useful in later sections.
###### Remark 4.12.
Let $`M^m`$ be a positive affine monoid of rank $`d`$, $`i\{0,\mathrm{},d\}`$, $`a\mathrm{C}(M)\mathrm{G}(M)`$ and $`(M,a)=\{F\text{ face of }\mathrm{C}(M):a\overline{MF}\}`$. Then $`H_𝔪^i(R)_a`$ is the $`i`$th cohomology of the complex
$$𝒞^\text{.}(M,a):0𝒞^0(M,a)\mathrm{}𝒞^t(M,a)\mathrm{}𝒞^d(M,a)0,$$
where
$$𝒞^t(M,a)=\underset{G(M,a),dimG=t}{}Ke_G$$
and the differential is given by $`(e_G)=_{F\text{ face of }\mathrm{C}(M),GF}ϵ(G,F)e_F`$. (Here $`ϵ`$ is the incidence function on the face lattice of $`\mathrm{C}(M)`$ which we fixed above to define the complex $`L^\text{.}(M)`$.)
###### Theorem 4.13.
Let $`M^m`$ be a positive affine monoid of rank $`d`$ and $`a\mathrm{C}(M)\mathrm{G}(M)`$. If the set $`(M,a)=\{F\text{ face of }\mathrm{C}(M):a\overline{MF}\}`$ has a unique minimal element $`G`$, then
$$H_𝔪^i(R)_a=0\text{for all }i=0,\mathrm{},d1.$$
###### Proof.
It follows from 4.12 that $`H_𝔪^i(R)_a`$ is the $`i`$th cohomology of the complex
$$𝒞^\text{.}(M,a):0𝒞^0(M,a)\mathrm{}𝒞^t(M,a)\mathrm{}𝒞^d(M,a)0,$$
where
$$𝒞^t(M,a)=\underset{G(M,a),dimG=t}{}Ke_G.$$
By taking a cross section of $`\mathrm{C}(M)`$ we see that the face lattice of $`\mathrm{C}(M)`$ is the face lattice of a polytope (see \[4, Proposition 6.1.8\]). If $`(M,a)`$ has a unique minimal element, then this set is again the face lattice of a polytope $`P`$, as can be seen from Ziegler \[20, Theorem 2.7\]. Note that if $`(M,a)`$ has only one element, then $`P`$ is the empty set. But this can only happen if $`(M,a)=\{\mathrm{C}(M)\}`$ and then we have homology only in cohomological degree $`d`$. If $`(M,a)`$ has more than one element, then $`𝒞^\text{.}(M,a)`$ is the $`K`$-dual of a cellular resolution which computes the singular cohomology of $`P`$. A nonempty polytope is homeomorphic to a ball and thus the complex $`𝒞^d(M,a)`$ is acyclic. Hence in this case $`H_𝔪^i(R)_a=0`$ for $`i=0,\mathrm{},d`$, and this concludes the proof. ∎
The following corollary collects two immediate consequences of Theorem 4.13.
###### Corollary 4.14.
Let $`M`$ be a positive affine monoid.
1. Let $`F`$ be the unique face of $`\mathrm{C}(M)`$ such that $`a\mathrm{int}F`$. If $`a\mathrm{G}(MF)`$, then
$$H_𝔪^i(R)_a=0\text{ for all }i=0,\mathrm{},d1.$$
2. Suppose that $`M`$ is seminormal. Then $`R`$ is Cohen–Macaulay if $`(M,a)`$ has a unique minimal element for all $`a\overline{M}`$.
Finally we note that nonzero lower local cohomologies must be large in the seminormal case.
###### Proposition 4.15.
Let $`M^m`$ be a positive affine seminormal monoid. If $`H_𝔪^i(R)_a0`$ for some $`a\mathrm{C}(M)`$, then $`dim_KH_𝔪^i(R)=\mathrm{}`$. In particular, if a seminormal monoid is Buchsbaum, then it must be Cohen–Macaulay.
###### Proof.
If $`H_𝔪^i(R)_a0`$, then the complex $`𝒞^\text{.}(M,a)`$ has nontrivial cohomology in degree $`i`$. Consider the multiples $`ka`$ for $`k`$. If $`a\overline{MF}=\mathrm{G}(MF)F`$, then $`ka\overline{MF}`$ for all $`k`$. If $`a\overline{MF}`$, then there exist infinitely many $`k`$ such that $`ka\overline{MF}`$. Since the face lattice of $`\mathrm{C}(M)`$ is finite we can choose a sequence $`(k_n)_n`$ such that $`k_n<k_{n+1}`$ and $`a\overline{MF}`$ if and only if $`k_na\overline{MF}`$. Thus $`𝒞^\text{.}(M,a)=𝒞^t(M,k_na)`$ for all $`n0`$ which implies $`H_𝔪^i(R)_{k_na}0`$. Hence $`dim_KH_𝔪^i(R)=\mathrm{}`$.
If $`R`$ is Buchsbaum, then $`dim_KH_𝔪^i(R)<\mathrm{}`$ for all $`i<d`$. Thus the local cohomology must vanish in this case for $`i<d`$ which implies that $`R`$ is already Cohen–Macaulay. ∎
## 5. The Cohen–Macaulay property and depth
If a seminormal monoid $`M`$ fails to be normal by the smallest possible margin, then $`K[M]`$ is Cohen–Macaulay as the following result shows:
###### Proposition 5.1.
Let $`M`$ be seminormal such that $`MF`$ is normal for each facet $`F`$ of $`\mathrm{C}(M)`$. Then $`R`$ is Cohen–Macaulay.
###### Proof.
It is enough to show that $`(M,a)`$ has a unique minimal element for all $`a\overline{M}`$. Let $`a\mathrm{G}(M)\mathrm{C}(M)`$. If $`a\overline{MG}`$ for all facets $`G`$ of $`\mathrm{C}(M)`$, then $`\mathrm{C}(M)`$ is the unique minimal element of $`(M,a)`$.
Otherwise we have $`a\overline{MG}=MGM`$ for some facet $`G`$ of $`\mathrm{C}(M)`$. We choose the unique face $`F^{}`$ of $`\mathrm{C}(M)`$ with $`a\mathrm{int}F^{}`$. It follows that $`aF^{}M`$, and $`F^{}`$ is the unique minimal element of $`(M,a)`$. ∎
###### Remark 5.2.
Another, albeit more complicated proof of the proposition can be given as follows. The main result of Brun, Bruns and Römer implies for $`R=K[M]`$ that
1. $`R/\overline{\omega }`$ is Cohen–Macaulay,
2. $`H_𝔪^d(R/\overline{\omega })=_FH_𝔪^{dimF}(K[MF])`$ where $`F`$ runs through the proper faces of $`\mathrm{C}(M)`$,
provided all the rings $`K[MF]`$ are Cohen–Macaulay. If they are even normal, then the local cohomology modules in (ii) do not “overlap” because the $`^m`$-graded support of $`H_𝔪^{dimF}(K[MF])`$ is restricted to $`\mathrm{int}F`$, and the relative interiors of faces are pairwise disjoint. Now we can conclude from Proposition 4.10 that $`R`$ is Cohen–Macaulay.
In general, without normality of the facets the local cohomology modules in (ii) will overlap (see Example 7.1). This limits all attempts to prove stronger assertions about the Cohen–Macaulay property in the seminormal case.
Using the results and techniques of Section 4, we can give lower bounds for the depth of seminormal monoid rings. Let $`M^m`$ be an affine seminormal monoid. We define
$`c_K(M)`$ $`=sup\{i:K[MF]\text{ is Cohen–Macaulay for all faces }F,dimFi\},`$
$`n(M)`$ $`=sup\{i:MF\text{ is normal for all faces }F,dimFi\}.`$
Observe that if $`MF`$ is normal for a face $`F`$ of $`\mathrm{C}(M)`$, then also $`MG`$ is normal for all faces $`GF`$ of $`\mathrm{C}(M)`$. Hence it would be enough to consider all $`i`$-dimensional faces of $`\mathrm{C}(M)`$ in the definition of $`n(M)`$. However, as we will see in Section 7, this is not true for the Cohen–Macaulay property.
###### Theorem 5.3.
Let $`M^m`$ be an affine seminormal monoid of rank $`d`$, and $`R=K[M]`$. Then
$$\mathrm{depth}Rc_K(M)\mathrm{min}\{n(M)+1,d\}.$$
###### Proof.
The proof of the first inequality follows essentially the same idea as that of Theorem 4.3.
The assertion is trivial for $`\mathrm{rank}M=0`$. Thus assume that $`d=\mathrm{rank}M>0`$. There is nothing to prove if $`c_K(M)=d`$. So we can assume that $`c_K(M)<d`$. Since $`M`$ is seminormal, we can again use the exact sequence
$$0\overline{\omega }RR/\overline{\omega }0.$$
Since $`\mathrm{depth}\overline{\omega }=d`$ according to Lemma 4.1, it is enough to show that $`\mathrm{depth}R/\overline{\omega }c_K(M)`$.
Again we write $`\overline{\omega }=_{j=1}^t𝔭_{F_j}`$ where $`F_1,\mathrm{},F_t`$ are the facets of $`\mathrm{C}(M)`$. However, contrary to Theorem 4.3, the bound does not hold for arbitrary residue class rings with respect to monomial radical ideals $`𝔮`$, since the combinatorial structure of the set $`F(𝔮)`$ may contain obstructions.
Therefore we order the facets $`F_1,\mathrm{},F_t`$ in such a way that they form a shelling sequence for the face lattice of $`\mathrm{C}(M)`$. Such a sequence exists by the Brugesser-Mani theorem (applied to a cross section polytope of $`\mathrm{C}(M)`$). See \[20, Lecture 8\]. The generalization of the first inequality of the theorem to be proved is the following: *let $`F_1,\mathrm{},F_t`$ be a shelling sequence for $`\mathrm{C}(M)`$ and let $`u\{1,\mathrm{},t\}`$, then $`\mathrm{depth}R/𝔮\mathrm{min}\{d1,c_K(M)\}`$ for $`𝔮=_{j=1}^u𝔭_{F_j}`$.*
If $`u=1`$, then $`R/𝔭_1\mathrm{}K[MF(𝔭_1)]`$. Now we can apply induction on $`\mathrm{rank}M`$. Let $`u>1`$. We set $`𝔮^{}=_{j=1}^{u1}𝔭_{F_j}`$. Again we have the standard exact sequence
$$0R/𝔮R/𝔮^{}R/𝔭_{F_u}R/(𝔮^{}+𝔭_{F_u})0.$$
Therefore
$$\mathrm{depth}R/𝔮\mathrm{min}\{1+\mathrm{depth}(R/(𝔮^{}+𝔭_{F_u}),\mathrm{depth}R/𝔭_{F_u},\mathrm{depth}R/𝔮^{}\}.$$
By induction on $`u`$ we have $`\mathrm{depth}R/𝔭_{F_u},\mathrm{depth}R/𝔮^{}\mathrm{min}\{d1,c_K(M)\}`$.
Now the crucial point is that $`F_u_{j=1}^{u1}F_j=_{j=1}^{u1}F_uF_j`$ is the union of certain facets $`G_1,\mathrm{},G_v`$ of $`F_u`$ that form the starting segment of a shelling sequence for $`F_u`$ (by the very definition of a shelling). As in the proof of Theorem 4.3 we have
$$R/(𝔮^{}+𝔭_{F_u})R/\underset{j=1}{\overset{t1}{}}(𝔭_{F_u}+𝔭_{F_j})\mathrm{}R^{}/𝔮^{\prime \prime }$$
where $`𝔮^{\prime \prime }=_{j=1}^{t1}\left((𝔭_{F_u}+𝔭_{F_j})/𝔭_{F_u}\right)`$. Therefore $`𝔮^{\prime \prime }`$ is the radical ideal of $`R^{}=K[MF_1]`$ corresponding to the union $`G_1,\mathrm{},G_v`$. By induction we have
$$\mathrm{depth}R^{}/𝔮^{\prime \prime }\mathrm{min}\{d2,c(MF_u)\}\mathrm{min}\{d2,c_K(M)\},$$
and this completes the proof for the inequality $`\mathrm{depth}Rc_K(M)`$.
By Hochster’s theorem the second inequality holds if $`M`$ itself is normal. Suppose that $`n(M)<d`$ and let $`F`$ be a face of dimension $`n(M)+1`$. Then we must show that $`K[MF]`$ is Cohen–Macaulay. Thus the second inequality reduces to the claim that $`R`$ is Cohen–Macaulay if the intersections $`FM`$ are normal for all facets $`F`$ of $`\mathrm{C}(M)`$ (and $`M`$ is seminormal). This has been shown in Proposition 5.1. ∎
There is a general lower bound for the depth of seminormal monoid rings of rank $`2`$. It follows from the proposition since seminormal monoids of rank $`1`$ are normal.
###### Corollary 5.4.
Let $`M^m`$ be an affine seminormal monoid of rank $`d2`$. Then
$$\mathrm{depth}R2.$$
In particular, $`R`$ is Cohen–Macaulay if $`d=2`$.
One could hope that seminormality plus some additional assumptions on $`M`$ already imply the Cohen–Macaulay property of $`R`$. But most time this is not the case as will be discussed in Example 7.1. However, we will now show that Serre’s condition $`(S_2)`$ implies the Cohen–Macaulay property of $`R`$ if $`\mathrm{C}(M)`$ is a simple cone (to be explained below). More generally, we want to show that simple faces of $`\mathrm{C}(M)`$ cannot contain an obstruction to the Cohen–Macaulay property in the presence of $`(S_2)`$.
Let $`F`$ be a proper face of $`\mathrm{C}(M)`$. We call the face $`F`$ simple if the partially ordered set $`\{G\text{ face of }\mathrm{C}(M):FG\}`$ is the face lattice of a simplex. Observe that by \[20, Theorem 2.7\] the latter set is always the face lattice of a polytope, because the face lattice of $`\mathrm{C}(M)`$ is the face poset of a cross section of $`\mathrm{C}(M)`$. Let $`F`$ be a simple face of $`\mathrm{C}(M)`$. It is easy to see that every face $`G`$ of $`\mathrm{C}(M)`$ containing the simple face $`F`$ is also simple.
We say that $`\mathrm{C}(M)`$ is simple if a cross section polytope of $`\mathrm{C}(M)`$ is a simple polytope. This amounts to the simplicity of all the edges of $`\mathrm{C}(M)`$. (Note that the apex $`\{0\}`$ is a simple face if and only if $`\mathrm{C}(M)`$ is a simplicial cone.)
###### Proposition 5.5.
Let $`M^m`$ be a positive affine seminormal monoid such that $`R`$ satisfies $`(S_2)`$. Let $`a\mathrm{G}(M)\mathrm{C}(M)`$ and $`a\mathrm{int}F`$ for a proper face $`F`$ of $`\mathrm{C}(M)`$. If $`H_𝔪^i(R)_a0`$ for some $`i`$, $`0id1`$, then $`F`$ is not a simple face of $`\mathrm{C}(M)`$.
###### Proof.
Assume that $`F`$ is a simple face. Consider the intersection
$$H=\underset{G\text{ face of }\mathrm{C}(M),FG,a\overline{MG}}{}G$$
which is a simple face containing $`F`$ because $`F`$ is simple. Let $`F_1,\mathrm{},F_t`$ be the facets of $`\mathrm{C}(M)`$. For each facet $`F_j`$ such that $`HF_j`$ there exists a face $`G`$ of $`\mathrm{C}(M)`$ with $`FG`$, $`a\overline{MG}`$ such that $`GF_j`$ because $`H`$ is simple. This follows from the fact that the partially ordered set $`\{L:L\text{ is a face of }\mathrm{C}(M)`$, $`HL\}`$ is the face poset of a simplex, and for a simplex the claim is trivially true. We observe that $`a\mathrm{G}(MG)\mathrm{G}(MF_j)`$ for those facets $`F_j`$ with $`HF_j`$. By Corollary 3.4 we have
$$\mathrm{G}(MH)=\underset{HF_j}{}\mathrm{G}(MF_j).$$
Therefore $`a\mathrm{G}(MH)H=\overline{MH}`$.
All in all we get that the set $`(M,a)=\{L\text{ face of }\mathrm{C}(M):a\overline{ML}\}`$ has the unique minimal element $`H`$, and 4.13 implies that
$$H_𝔪^i(R)_a=0$$
which is a contradiction to our assumption. Thus $`F`$ is not a simple face of $`\mathrm{C}(M)`$. ∎
The latter result gives a nice Cohen–Macaulay criterion in terms of $`\mathrm{C}(M)`$ for a seminormal monoid. It implies Theorem 4.4.7 in , and can be viewed as a variant of the theorem by Goto, Suzuki and Watanabe by which $`(S_2)`$ implies the Cohen–Macaulay property of $`R`$ if $`\mathrm{C}(M)`$ is simplicial.
###### Corollary 5.6.
Let $`M^m`$ be a positive affine seminormal monoid such that $`R=K[M]`$ satisfies $`(S_2)`$ and such that $`\mathrm{C}(M)`$ is a simple cone. Then $`R`$ is Cohen–Macaulay for every field $`K`$. In particular, if $`\mathrm{rank}M3`$, then $`R`$ is Cohen–Macaulay.
###### Proof.
Every proper face of $`\mathrm{C}(M)`$, with the potential exception of $`\{0\}`$, is simple. Thus it follows from 4.3 and 5.5 that $`H_𝔪^i(R)_a=0`$ for $`a0`$ and $`i=0,\mathrm{},d1`$. For $`a=0`$ this results from Corollary 4.14. Hence $`R`$ is Cohen–Macaulay.
If $`\mathrm{rank}M3`$, then the cross section of $`\mathrm{C}(M)`$ is a polytope of dimension $`2`$, which is always simple. Thus we can apply the corollary. ∎
We will point out in Remark 7.2 that the corollary is the best possible result if one wants to conclude the Cohen–Macaulay property of $`R`$ only from the seminormality of $`M`$ and the validity of ($`S_2`$).
## 6. Seminormality in characteristic $`p`$
In this section we study local cohomology properties of seminormal rings in characteristic $`p>0`$. Let $`K`$ be a field with $`\mathrm{char}K=p>0`$. In this situation we have the Frobenius homomorphism $`F:RR,ff^p`$. Through this homomorphism $`R`$ is an $`F(R)`$-module. Now $`R`$ is called $`F`$-injective if the induced map on the local cohomology $`H_𝔪^i(R)`$ is injective for all $`i`$. It is called $`F`$-pure if the extension $`F(R)R`$ is pure, and $`F`$-split if $`F(R)`$ is a direct $`F(R)`$-summand of $`R`$. In general we have the implications
$$F\text{-split }\text{ }F\text{-pure }\text{ }F\text{-injective}.$$
For example see for general properties of these notions.
###### Proposition 6.1.
Let $`M^m`$ be a positive affine monoid. If there exists a field $`K`$ of characteristic $`p`$ such that $`R`$ is $`F`$-injective, then $`M`$ is seminormal.
###### Proof.
Assume that there exists an $`a\mathrm{G}(M)`$, $`a\mathrm{C}(M)`$, and an $`i\{0,\mathrm{},\mathrm{rank}M\}`$ such that $`H_𝔪^i(R)_a0`$. Since $`R`$ is $`F`$-injective, it follows that $`H_𝔪^i(R)_{p^ma}0`$ for all $`m`$. Write $`R=S/I_M`$ as a $`^m`$-graded quotient of a polynomial ring $`S`$. Then by graded local duality $`H_𝔪^i(R)^{}\mathrm{Ext}_S^{ni}(R,\omega _S)`$ is a finitely generated $`^m`$-graded $`R`$-module. This implies that $`H_𝔪^i(R)_{p^ma}=0`$ for $`m0`$, which is a contradiction. Thus $`H_𝔪^i(R)_a=0`$ for all $`a\mathrm{C}(M)`$. It follows from 4.7 that $`M`$ is seminormal. ∎
If $`M`$ is seminormal there exist only finitely many prime numbers such that $`R`$ is not $`F`$-injective. Moreover, we can characterize this prime numbers precisely.
###### Proposition 6.2.
Let $`M^m`$ be a positive affine seminormal monoid and let $`K`$ be a field of characteristic $`p>0`$. Then the following statements are equivalent:
1. The prime ideal $`(p)`$ is not associated to the $``$-module $`\mathrm{G}(M)F/\mathrm{G}(MF)`$ for any face $`F`$ of $`\mathrm{C}(M)`$;
2. $`R`$ is $`F`$-split;
3. $`R`$ is $`F`$-pure;
4. $`R`$ is $`F`$-injective.
###### Proof.
(i) $``$ (ii) We show by a direct computation that $`F(R)=K^p[pM]`$ is a direct $`K^p[pM]`$-summand of $`R`$. Since $`K^p[pM]`$ is a $`K^p[pM]`$-summand of $`K[pM]`$, it is enough to show that $`K[pM]`$ is a direct $`K[pM]`$-summand of $`R`$. The monoid $`M`$ is the disjoint union of the residue classes modulo $`p\mathrm{G}(M)`$. This induces a direct sum decomposition of $`R`$ as a $`K[pM]`$-module. We claim that $`pM`$ is the intersection of $`M`$ and $`p\mathrm{G}(M)`$. This will show the remaining assertion.
To prove the claim we have only to show that an element $`w`$ of the intersection of $`M`$ and $`p\mathrm{G}(M)`$ is an element of $`pM`$. Write $`w=pz`$ for some $`z\mathrm{G}(M)`$. We have that $`w\mathrm{int}F`$ for a face $`F`$ of $`\mathrm{C}(M)`$. Then $`p`$ annihilates the element $`z\mathrm{G}(M)F`$ modulo $`\mathrm{G}(MF)`$. By assumption $`p`$ is a nonzero-divisor on that module. Thus $`z\mathrm{G}(MF)`$. Since $`z\mathrm{G}(MF)\mathrm{int}F`$ and $`M`$ is seminormal we have that $`zM`$ by 2.1. Hence $`wpM`$ as desired.
(ii) $``$ (iii) and (iii) $``$ (iv) : This holds in general as was remarked above.
(iv) $``$ (i) Assume that $`(p)`$ is associated to some of the $``$-modules $`\mathrm{G}(M)F/\mathrm{G}(MF)`$ for the faces $`F`$ of $`\mathrm{C}(M)`$. Choose a maximal face $`F`$ with this property. Choose an element $`\overline{0}\overline{a}\mathrm{G}(M)F/\mathrm{G}(MF)`$ which is annihilated by $`p`$. We may assume that $`a\mathrm{int}F`$.
If follows from 4.12 that $`H_𝔪^{dimF+1}(R)_a0`$, since the poset $`(M,a)`$ consists all faces which contain $`F`$ but not $`F`$ itself. But $`H_𝔪^{dimF+1}(R)_{pa}=0`$, because $`(pa)`$ consists all faces which contain $`F`$ including $`F`$ itself. This poset is the face poset of a polytope and thus acyclic. We have derived a contradiction to the assumption that $`R`$ is $`F`$-injective, because the map $`H_𝔪^{dimF+1}(R)_aH_𝔪^{dimF+1}(R)_{pa}`$ is not injective. ∎
###### Corollary 6.3.
Let $`M^m`$ be a positive affine monoid. Then the following statements are equivalent:
1. $`M`$ is normal;
2. $`R`$ is $`F`$-split for every field $`K`$ of characteristic $`p>0`$;
3. $`R`$ is $`F`$-pure for every field $`K`$ of characteristic $`p>0`$;
4. $`R`$ is $`F`$-injective for every field $`K`$ of characteristic $`p>0`$.
###### Proof.
It is easy to see that $`M`$ is normal if and only if $`\mathrm{G}(M)F=\mathrm{G}(MF)`$ for all faces of $`\mathrm{C}(M)`$. Thus (i) is equivalent to (ii), (iii) and (iv) by 6.1 and 6.2. ∎
In a sense, it is inessential for the results of this section that $`K`$ has characteristic $`p`$. In order to obtain variants that are valid for every field $`K`$, one has to replace the Frobenius endomorphism by the natural inclusion $`K[pM]K[M]`$.
## 7. Examples and Counterexamples
In this section we present various examples and counterexamples related to the results of this paper. We choose a field $`K`$.
We saw in 5.6 that $`R=K[M]`$ is Cohen–Macaulay for a positive seminormal monoid $`M`$ of rank $`d3`$. Since a Cohen–Macaulay ring always satisfies $`(S_2)`$ one could hope that seminormality and $`(S_2)`$ already imply the Cohen–Macaulay property of $`R`$. This is not the case as the following example shows.
###### Example 7.1.
For this example and the following one we fix some notation. Let $`P`$ be a $`3`$-dimensional pyramid with a square base embedded into $`^4`$ in degree 1. For example let $`P`$ be the convex hull of the vertices given by
$$m_0=(0,0,1,1),m_1=(1,1,0,1),m_2=(1,1,0,1),$$
$$m_3=(1,1,0,1),m_4=(1,1,0,1).$$
Figure 1 shows projections of the pyramid onto its base.
Let $`C`$ be the cone generated by $`P`$, so that $`P`$ is a cross section of $`C`$. The facets of $`C`$ are the cones
$$F_0=\mathrm{C}(m_1,m_2,m_3,m_4),F_1=\mathrm{C}(m_0,m_1,m_2),F_2=\mathrm{C}(m_0,m_2,m_3),$$
$$F_3=\mathrm{C}(m_0,m_3,m_4),F_4=\mathrm{C}(m_0,m_1,m_4).$$
Let $`M`$ be the monoid generated by all integer points of even degree in the facets $`F_1`$ and $`F_3`$ and all their faces, and all remaining integer points in the interior of all other faces of $`C`$ including $`C`$. (The facets $`F_1`$ and $`F_3`$ have been shaded in the left diagram in Figure 1.) Thus $`M`$ is positive, $`\mathrm{C}(M)=C`$ and $`M`$ is not normal. It follows from 2.1 that $`M`$ is seminormal and from 3.4 that $`R`$ satisfies $`(S_2)`$.
We claim that $`R`$ is not Cohen–Macaulay. Observe that all faces of $`C`$ are simple except the face $`\mathrm{C}(m_0)`$. Thus 4.3 and 5.5 imply that $`H_𝔪^i(R)_a`$ can be nonzero only for some $`a\mathrm{C}(m_0)\mathrm{G}(M)`$. Choose an arbitrary $`a`$ in the relative interior of $`\mathrm{C}(m_0)\mathrm{G}(M)`$ of odd degree. The set $`(M,a)`$ introduced in 4.12 is
$$\{F_2,F_4,C\}$$
and we see that the complex $`𝒞^\text{.}(M,a)`$ has cohomology in cohomological degree $`3`$. Hence $`H_𝔪^3(R)_a0`$ and therefore $`R`$ is not Cohen–Macaulay.
The reader may check that $`R/\overline{\omega }`$ is Cohen–Macaulay, but both $`K[F_2M]`$ and $`K[F_4M]`$ have nonzero third local cohomology in degree $`a`$.
###### Remark 7.2.
Example 7.1 can be generalized in the following way: if $`C`$ is not a simple cone, then there exists a seminormal affine monoid $`M`$ with $`C=\mathrm{C}(M)`$ such that $`K[M]`$ satisfies ($`S_2`$), but is not Cohen–Macaulay for any field $`K`$.
Next we consider the question whether the Cohen–Macaulay property or the $`(S_2)`$ property are inherited by face projections. A counterexample to this claim is already given in \[9, Example 2.2\]. We can modify 7.1 a little bit to get the same result for seminormal monoids.
###### Example 7.3.
With the same notation as in 7.1 let $`C`$ be the cone over the pyramid $`P`$ with facets $`F_0,\mathrm{},F_4`$. Now let $`M`$ be the monoid generated by all integer points of even degree in the facet $`F_1`$ and all its faces (as indicated in the right diagram in Figure 1), and all remaining integer points in the interior of all other faces of $`C`$. Thus $`M`$ is positive, $`\mathrm{C}(M)=C`$ and $`M`$ is not normal. It still follows from 2.1 that $`M`$ is seminormal and by 3.4 that $`R`$ satisfies $`(S_2)`$. Since all proper faces of $`C`$ except $`\mathrm{C}(m_0)`$ are simple, we only have to check the vanishing of the local cohomology for points in the in $`\mathrm{int}\mathrm{C}(m_0)`$. Let $`a\mathrm{int}\mathrm{C}(m_0)`$. If $`a`$ has even degree, then
$$(M,a)=\{F\text{ face of }C:\mathrm{C}(m_0)F\}.$$
If $`a`$ has odd degree, then
$$(M,a)=\{F_2F_3,F_3F_4,F_2,F_3,F_4,C\}.$$
In any case, we can check that the complex $`𝒞^\text{.}(M,a)`$ is acyclic and therefore $`H_𝔪^i(R)=0`$ for $`i<\mathrm{rank}M`$. Thus $`R`$ is Cohen–Macaulay and must satisfy $`(S_2)`$.
But $`K[MF_3]`$ has only depth $`1`$, as can be seen from a similar discussion as for $`R`$. So it does not satisfy $`(S_2)`$. Hence neither the Cohen–Macaulay property, nor $`(S_2)`$ are inherited by face projections of seminormal monoid rings.
Let $`\mathrm{\Delta }`$ be a simplicial complex contained in the simplex $`\mathrm{\Sigma }`$ with vertex set $`V`$. We consider the dual simplex $`\mathrm{\Sigma }^{}`$ whose facets correspond bijectively to the vertices $`vV`$ of $`\mathrm{\Delta }`$. Next we erect the pyramid $`\mathrm{\Pi }`$ over $`\mathrm{\Sigma }^{}`$ with apex $`t`$. Then the faces of $`\mathrm{\Pi }`$ that contain $`t`$ are in bijective correspondence with the faces $`G`$ of $`\mathrm{\Sigma }`$:
$$F\mathrm{\Sigma }F^{}\mathrm{\Sigma }^{}\stackrel{~}{F}=F^{}\{t\}$$
where $``$ denotes the join. Observe that this correspondence reverses the partial order by inclusion. Choose a realization of $`\mathrm{\Pi }`$ as a rational polytope, also denoted by $`\mathrm{\Pi }`$.
Next we plane off those faces of $`\mathrm{\Pi }`$ that correspond to the minimal non-faces of $`\mathrm{\Delta }`$. For such a non-face $`G`$ we choose a support hyperplane $`H`$ of $`\mathrm{\Pi }`$ with $`\mathrm{\Pi }H=\stackrel{~}{G}`$. Moving this hyperplane by a sufficiently small rational displacement towards the interior of $`\mathrm{\Pi }`$, and intersecting $`\mathrm{\Pi }`$ with the positive halfspace of the displaced parallel hyperplane $`H^{}`$ we obtain a polytope $`\mathrm{\Pi }^{}`$ such that exactly the faces $`F`$ of $`\mathrm{\Sigma }`$ that do not contain $`G`$ are preserved in $`\mathrm{\Pi }^{}`$: $`\stackrel{~}{F}\mathrm{\Pi }^{}\mathrm{}FG`$.
Repeating this construction for each minimal non-face of $`\mathrm{\Delta }`$ we finally reach a polytope $`\mathrm{\Pi }^{\prime \prime }`$ in which exactly the faces $`\stackrel{~}{F}`$, $`F\mathrm{\Delta }`$, have survived in the sense that $`F^{}=\stackrel{~}{F}\mathrm{\Pi }^{\prime \prime }`$ is a non-empty face of $`\mathrm{\Pi }^{\prime \prime }`$.
Moreover, the only facets of $`\mathrm{\Pi }^{\prime \prime }`$ containing $`F^{}`$ are the facets $`\{v\}^{}`$ corresponding to the vertices $`vF`$. On the other hand, every face $`E`$ of $`\mathrm{\Pi }^{\prime \prime }`$ that is not of the form $`F^{}`$ is contained in at least one “new” facet of $`\mathrm{\Pi }^{\prime \prime }`$ created by the planing of $`\mathrm{\Pi }`$.
Note that $`\mathrm{\Pi }`$ is a simplex and therefore a simple polytope. The process by which we have created $`\mathrm{\Pi }^{\prime \prime }`$ does not destroy simplicity if the displacements of the hyperplanes are sufficiently small and “generic”.
Set $`d=dim\mathrm{\Pi }^{\prime \prime }+2`$ and embed $`\mathrm{\Pi }^{\prime \prime }`$ into $`^{d2}\times \{0\}^{d1}`$. Then let $`\mathrm{\Gamma }`$ be the pyramid over $`\mathrm{\Pi }^{\prime \prime }`$ with apex $`v=(0,\mathrm{},0,1)`$. The construction of $`\mathrm{\Gamma }`$ that leads to the pyramid of Example 7.1 is illustrated in Figure 2.
Note that all faces of $`\mathrm{\Gamma }`$, except $`\{v\}`$, are simple. ($`\{v\}`$ is simple only if $`\mathrm{\Delta }=\mathrm{\Sigma }`$, or equivalently, $`\mathrm{\Pi }^{\prime \prime }=\mathrm{\Pi }`$.)
In the last step we embed $`\mathrm{\Gamma }`$ into $`^{d1}\times \{1\}^d`$ by the assignment $`x(x,1)`$ and choose the cone $`C=_+\mathrm{\Gamma }`$. It has dimension $`d`$. The point $`v`$ (in $`^d`$) has the coordinates $`(0,\mathrm{},0,1,1)`$. Therefore it has value $`1`$ under the linear form $`\mathrm{deg}:^d`$, $`\mathrm{deg}(y)=y_d`$.
Set $`L=\{a^d:\mathrm{deg}(a)0(2)\}`$. To each facet $`F`$ of $`C`$ we assign the lattice
$$L_F=\{\begin{array}{cc}F^d,\hfill & F=_+\{v\}^{}\text{ for some }vV,\hfill \\ FL,\hfill & \text{else}.\hfill \end{array}$$
Finally, we let $`M`$ be the monoid formed by all $`aC^d`$ such that $`aL_F`$ for all facets $`F`$ of $`C`$ containing $`a`$. Clearly $`M`$ is seminormal. Moreover, with the notation of Corollary 3.4, we have $`M=M^{}`$, since we have restricted the lattice facet-wise, and thus $`K[M]`$ satisfies $`(S_2)`$ for all fields $`K`$.
Let $`a\overline{M}=\mathrm{G}(M)C`$. (By construction we have $`\mathrm{G}(M)=^d`$.) If the face $`F`$ of $`C`$ with $`a\mathrm{int}(F)`$ is different from $`_+v`$, then it is a simple face of $`C`$, and $`H_𝔪^i(K[M])_a=0`$ for all $`i<n`$ by Proposition 5.5. If $`F=_+v`$ and $`\mathrm{deg}(a)0(2)`$, then we arrive at the same conclusion by Corollary 4.14. However, if $`F=_+v`$ and $`\mathrm{deg}(a)1(2)`$, then the poset $`(M,a)`$ is isomorphic to the dual of $`\mathrm{\Delta }`$ (as a poset). Hence the cochain complex $`𝒞^{}(M,a)`$ is isomorphic to the chain complex of $`\mathrm{\Delta }`$ (up to shift). We have
(3)
$$H_𝔪^i(K[M])_a=\stackrel{~}{H}_{di1}(\mathrm{\Delta };K),i=0,\mathrm{},d.$$
###### Theorem 7.4.
Let $`\mathrm{\Delta }`$ be a simplicial complex and $`K`$ a field. Then there exists a seminormal monoid $`M`$ of rank $`d`$ with $`M=M^{}`$ and such that $`R=K[M]`$ has the following properties:
1. For every $`a\overline{M}`$
1. $`H_𝔪^i(R)_a=0`$ for all $`i<d`$, or
2. $`H_𝔪^{}(R)_a`$ is given by (3).
2. Moreover, case (b) holds true for at least one $`a\overline{M}`$.
3. The following are equivalent:
1. $`\mathrm{\Delta }`$ is acyclic over $`K`$;
2. $`R`$ is Cohen-Macaulay.
If we choose $`\mathrm{\Delta }`$ as a triangulation of the real projective plane, then we obtain a monoid algebra $`K[M]`$ which is Cohen-Macaulay if and only if $`\mathrm{char}K2`$.
|
warning/0506/math0506465.html
|
ar5iv
|
text
|
# Sample Paths in Wavelet Theory
## 1 Introduction. Wavelets
We will be primarily interested in the infinite products that typically occur in wavelet analysis \[12, Ch. 5\]; see also . As is well known, the traditional approach to wavelet-scaling functions is based on analysis of wavelet filters which are periodic functions in one or more variables, and which have some degree of regularity (say, Lipschitz at low frequencies). In addition, the filter functions must satisfy a variant of the following two features: (i) some quadrature condition, and (ii) some low-pass condition. But conditions (i)–(ii) are known to be incompatible with several new classes of examples. For these wider classes of wavelets, we show that a new *uncertainty principle* for wavelet filters is at play. Reflecting this uncertainty principle for wavelet filters, we note that recent papers on wavelet sets, and on other frequency-localized wavelets (see, e.g., ), necessitate the consideration of wavelet filters which are defined only a.e., and for which conditions (i)–(ii) are more subtle. As a result, the standard deterministic methods then do not immediately apply.
In this paper, we use a random walk model (see ) in a new analysis of wavelets, but for a class of wavelets where the filters are more singular than is the case in the standard references . Our results apply also to iterative algorithms for wavelet packets . In these problems solutions are constructed from an algorithm which relates a certain function $`\phi (t)`$, $`t`$, called a scaling function, to its scaled version $`\phi (Nt)`$ where $`N`$ is a fixed integer, $`2`$ or more, or an expansive integral $`d`$ by $`d`$ matrix $`A`$ if $`t`$ is in $`^d`$.
For $`d=1`$, the formula which determines $`\phi `$ is
$$\phi \left(t\right)=N\underset{k}{}a_k\phi \left(Ntk\right),t.$$
(1.1)
It is called the *scaling identity*. Consider the filter function $`m\left(x\right)=\underset{k}{}a_ke^{i2\pi kx}`$ for $`z=\mathrm{exp}(i2\pi x)`$ in the circle $`𝕋`$, or the $`d`$-torus $`𝕋^d`$. It is traditional to impose some degree of regularity, e.g., Lipschitz, or low-pass for $`x=0`$; but we shall not do this. Instead we shall work with a probabilistic notion which serves as a minimal restriction on the filter $`m\left(x\right)`$ guaranteeing the existence of $`L^2`$ solutions to the scaling identity (1.1). It is not at all clear *a priori* that this identity should have $`L^2`$ solutions, or even solutions that are given by some kind of convergent algorithm. The coefficients in (1.1) are called *masking coefficients*, *filters*, or filter coefficients. For special values of the coefficients, it turns out that there is a normalized $`L^2`$ solution $`\phi `$, and that pointwise convergence of a good approximation can be established in a meaningful way. Moreover, we show that the convergence behavior is dictated by properties of an associated transfer operator $`R`$, or Ruelle operator.
A main point in this paper is our suggestion of a different and more versatile approach to how we impose low-pass conditions on the basic wavelet filters. Our approach is probabilistic, and it is based directly on the Ruelle operator $`R=R_W`$, and on a certain Perron–Frobenius eigenfunction $`h`$ for $`R`$ (see Theorem 5.2). Our work is motivated in part by recent considerations of frequency-localized wavelets, as pioneered in for example in the papers . These papers suggest an interesting context for wavelets that are necessarily frequency-localized. But at the same time, it is evident from this work that the corresponding wavelet filters will then typically not satisfy the low-pass properties that are otherwise known to hold for more traditional time- (or space-) localized wavelets.
The expression on the right-hand side in equation (1.1) is also called a *subdivision* because of the function values $`\phi (Ntk)`$. An iteration of the operations on the right hand side in (1.1) on some initial function is called a *cascade approximation*. This approximation is one approach to the function $`\phi `$, and the other is the infinite product formula (3.5) for the Fourier transform of $`\phi `$.
The relation between $`\phi `$ and its scaled refinement is well understood when we pass from the time domain to the frequency domain, via the Fourier transform. In that case the relation is multiplicative, and involves a certain periodic matrix function $`m`$, called the *low-pass filter*. For further details, see formulas (3.2) and (3.4) below.
The study of the filters $`m`$ is part of signal processing (see, e.g., ). But by a ‘mathematical miracle’, they have become one of the most useful tools in wavelet constructions; and at the same time, they have pointed to a host of exciting applications of wavelet mathematics; see, e.g., . To get a path space measure for some of the wavelet problems, we use the quadrature mirror properties (see (2.3) below), or their generalizations, which are assumed for the *filter function* $`m`$, also called a frequency response function. It is periodic, in one or several variables. In $`^d`$, there is a variety of choices of a period lattice for the problems at hand.
For a number of applications (see especially ), we must consider vector versions of the scaling identity (1.1); and in that case, the solution $`\phi `$, the *scaling function*, will be viewed as a vector-valued function, i.e., a function from $`^d`$ into some Hilbert space, typically finite-dimensional. In that case, the coefficients $`a_k`$ in (1.1) will be matrix-valued; and the product on the right-hand side in (1.1) will be matrix acting on vector. Following , we then say that the initial resolution subspace $`V_0`$ in $`L^2(^d)`$ has multiplicity. In the more traditional multiresolution analysis (MRA) approach to wavelets, $`\phi `$ is a scalar function, and $`V_0`$ is the closed span of the $``$-translates of $`\phi `$. The case of multiplicity is called the generalized MRA (GMRA) approach. The present considerations deal with the issue of passing from the filter function (possibly matrix-valued) to the scaling function $`\phi `$. However, the formulation of our ideas in the scalar case may easily be modified to the matrix/vector case; and for the sake of simplicity, our technical discussion below will be presented in the scalar case. We leave to the reader the spelling out of the generalization to GMRAs.
The fact that there are solutions in $`L^2(^d)`$ is not at all obvious; see . In application to images, the subspace $`V_0`$ (where $`V_0:=\mathrm{cl}\mathrm{span}\{\phi (k)k^d\}`$) may represent a certain resolution, and hence there is a choice involved, but we know by standard theory, see, e.g., , that under apropriate conditions, such choices are possible. As a result there are extremely useful, and computationally efficient, wavelet bases in $`L^2(^d)`$. A resolution subspace $`V_0`$ within $`L^2(^d)`$ can be chosen to be arbitrarily fine: Finer resolutions correspond to larger subspaces.
As noted for example in , a variant of the scaling equation is also used in computer graphics: there data is successively subdivided and the refined level of data is related to the previous level by prescribed masking coefficients. The latter coefficients in turn induce generating functions which are direct analogues of wavelet filters.
One reason for the computational efficiency of wavelets lies in the fact that *wavelet coefficients* in *wavelet expansions* for functions in $`V_0`$ may be computed using matrix iteration, rather than by a direct computation of inner products: the latter would involve integration over $`^d`$, and hence be computationally inefficient, if feasible at all. The deeper reason for why we can compute wavelet coefficients using matrix iteration is an important connection to the subband filtering method from signal/image processing involving digital filters, down-sampling and up-sampling. In this setting filters may be realized as functions $`m_0`$ on a $`d`$-torus, e.g., quadrature mirror filters.
As emphasized for example in and , because of down-sampling, the matrix iteration involved in the computation of wavelet coefficients involves so-called slanted Toeplitz matrices $`F`$ from signal processing. The slanted matrices $`F`$ are immediately available; they simply record the numbers (masking coefficients) from the $`\phi `$-scaling equation. These matrices further have the computationally attractive property that the iterated powers $`F^k`$ become sucessively more sparse as $`k`$ increases, i.e., the matrix representation of $`F^k`$ has mostly zeros, and the non-zero terms have an especially attractive geometric configuration. In fact subband signal processing yields a finite family, $`F`$, $`G`$, etc., of such slanted matrices; for example (with $`L`$ for “low frequency” and $`H`$ for “high frequency”):
$$\phi =\underset{k}{}P_k\phi (2k),\psi =\underset{k}{}Q_k\phi (2k),$$
$$\{\begin{array}{cc}\hfill F:& y_n^L=\frac{1}{\sqrt{2}}\underset{k}{}P_{k2n}x_k,\hfill \\ \hfill G:& y_n^H=\frac{1}{\sqrt{2}}\underset{k}{}Q_{k2n}x_k.\hfill \end{array}$$
The wavelet coefficients at scaling level $`k`$ of a numerical signal $`s`$ from $`V_0`$ are then simply the coordinates of $`GF^ks`$. By this we mean that a signal in $`V_0`$ is represented by a vector $`s`$ via a fixed choice of scaling function; see .
Given some choice of scaling function $`\phi `$, then under suitable conditions (e.g., orthogonality or *a priori* frame estimates) we may define the operator $`W:V_0\mathrm{}^2`$ which links the resolution subspace $`V_0`$ with the sequence space $`\mathrm{}^2`$ of signals $`s=(s_k)`$ as follows:
$$W\left(\underset{k^d}{}s_k\phi (k)\right)=\left(s_k\right).$$
Here $`W`$ becomes a well defined linear operator, $`W:V_0\mathrm{}^2`$
Then the matrix product $`GF^k`$ is applied to $`s`$; and the matrices $`GF^k`$ get more slanted as $`k`$ increases.
Our approach begins with the observation that the computational feature of this engineering device can be said to begin with an endomorphism $`r_A`$ of the $`d`$-torus $`𝕋^d=^d/^d`$, an endomorphism which results from simply passing matrix multiplication by $`A`$ on $`^d`$ to the quotient by $`^d`$. It is then immediate that the inverse images $`r_A^1(x)`$ are finite for all $`x`$ in $`𝕋^d`$, in fact $`\mathrm{\#}r_A^1(x)=|detA|`$. From this we recover the scaling identity, and we note that the wavelet scaling equation is a special case of a more general identity known in computational fractal theory and in symbolic dynamics . We show that wavelet algorithms and harmonic analysis naturally generalize to affine iterated function systems. Moreover, in this general context, we are able to build the ambient Hilbert spaces for a variety of dynamical systems which arise from the iterated dynamics of endomorphisms of compact spaces .
As a consequence, the fact that the ambient Hilbert space in the traditional wavelet setting is the more familiar $`L^2(^d)`$ is merely an artifact of the choice of filters $`m_0`$. As we further show, by enlarging the class of admissible filters, there is a variety of other ambient Hilbert spaces possible with corresponding wavelet expansions: the most notable are those which arise from iterated function systems (IFS) of fractal type, for example for the middle-third Cantor set, and scaling by $`3`$.
With examples, with theorems, and with graphics, we hope to bring these threads to light in this little book: The journey from wavelets to fractals via signal processing.
More generally, there is a variety of other natural dynamical settings (affine IFSs) that invite the same computational approach.
The two most striking examples which admit such a harmonic analysis are perhaps complex dynamics and subshifts. Both will be worked out in detail. In the first case, consider a given rational function $`r(z)`$ of one complex variable. We then get an endomorphism $`r`$ acting on an associated Julia set $`X`$ in the complex plane $``$ as follows: This endomorphism $`r:XX`$ results by restriction to $`X`$ . (Details: Recall that $`X`$ is by definition the complement of the points in $``$ where the sequence of iterations $`r^n`$ is a normal family. Specifically, the Fatou set $`F`$ of $`r(z)`$ is the largest open set in $``$ where $`r^n`$ is a normal sequence of functions, and we let $`X`$ be the complement of $`F`$. Here $`r^n`$ denotes the $`n`$’th iteration of the rational function $`r(z)`$.) The induced endomorphism $`r`$ of $`X`$ is then simply the restriction to $`X`$ of $`r(z)`$. If $`r`$ then denotes the resulting endomorphism, $`r:XX`$, it is known that $`\mathrm{\#}r^1(x)=`$degree of $`r`$, for every $`x`$ in $`X`$ (except for a finite set of singular points).
In the second case, for a particular one-sided subshift, we may take $`X`$ as the corresponding state space, and again we have a naturally induced finite-to-one endomorphism of $`X`$ of geometric and computational significance.
But in the general framework, there is not a natural candidate for the ambient Hilbert space. That is good in one sense, as it means that the subband filters $`m_0`$ which are feasible will constitute a richer family of functions on $`X`$.
In all cases, the analysis is governed by a random-walk model with successive iterations where probabilities are assigned on the finite sets $`\mathrm{\#}r^1(x)`$ and are given by the function $`W:=|m_0|^2`$. This leads to a *transfer operator* $`R_W`$ which has features in common with the classical operator considered first by Perron and Frobenius for positive matrices, in particular it has a Perron–Frobenius eigenvalue, and positive Perron–Frobenius eigenvectors, one on the right, a function, and one on the left, a measure; see . This Perron–Frobenius measure, also sometimes called the Ruelle measure, is an essential ingredient for our construction of an ambient Hilbert space. All of this, we show, applies to a variety of examples, and as we show, has the more traditional wavelet setup as a special case, in fact the special case when the Ruelle measure on $`𝕋^d`$ is the Dirac mass corresponding to the point $`0`$ in $`𝕋^d`$ (additive notation) representing zero frequency in the signal processing setup.
There are two more ingredients entering in our construction of the ambient Hilbert space: a path-space measure governed by the $`W`$-probabilities, and certain finite cycles for the endomorphism $`r`$. For each $`x`$ in $`X`$, we consider paths by infinite iterated tracing back with $`r^1`$ and recursively assigning probabilities with $`W`$. Hence we get a measure $`P_x`$ on a space of paths for each $`x`$. These measures are in turn integrated in $`x`$ using the Ruelle measure on $`X`$. The resulting measure will now define the inner product in the ambient Hilbert space.
Since the first question is to decide when $`\phi `$ is in $`L^2`$, we iterate and get an infinite matrix product involving the matrix $`W:=m^{}m`$. Since $`W`$ is positive semidefinite, we may create a positive path measure of a random walk starting at $`x`$ in some period interval. In several dimensions, $`x`$ starts in a fundamental domain $`D`$ for some fixed lattice, for example $`^d`$. The paths starting at $`x`$ arise by iteration of the inverse branches of $`xAxmod^d`$. There are $`N=|detA|`$ distinct branches. These $`N`$ branches may be viewed as endomorphisms of $`D`$.
In this paper, we will consider pointwise convergence of the infinite product (3.5) below for the Fourier transform $`\widehat{\phi }`$ of the scaling function $`\phi `$. But the traditional low-pass/regularity considerations for the filter $`m(x)`$ are more general. Once pointwise convergence (Theorem 5.2) is established, then the quadrature property for $`m(x)`$ will imply that the scaling function $`\phi `$ is automatically in $`L^2()`$.
We construct our random walks in a general framework which includes both wavelets, wavelet packets, and some of the other more classical problems. We further show how some of the classical questions may be phrased and solved with the use of path space measures.
It is interesting to contrast our proposed approach with the more traditional one used in wavelet analysis: see, e.g., . Traditionally, some kind of Lipschitz or Dini regularity condition must be assumed for the filter. Then the corresponding infinite product may be made precise, and we can turn to the question of when the wavelet generators are in $`L^2(^d)`$; see, e.g., and . As it turns out, both of these issues have natural formulations, and solutions, in terms of the path space measures. And the results allow a wider generality. In a variety of wavelet questions for band-limited wavelets, the regularity conditions just aren’t satisfied for the filters $`m`$ that are dictated by the setting and the applications; see, e.g., .
## 2 Path space
A well tested tool in analysis, and in mathematical physics, centers around the application of path-space measures. This tool is used in attacking a variety of singular convergence, or approximation, problems. We will adopt this viewpoint in our study of wavelet approximations. Traditionally, the setting for wavelet questions has included assumptions concerning continuity, or some kind of differentiability. In contrast, we shall work almost entirely in the measurable category. One advantage of our present approach is that we stay in the measurable category when addressing problems from multiresolution analysis (MRA). Earlier work on the use of probability in wavelets includes that of R.F. Gundy et al.; see .
Our present viewpoint is more general than , and it starts with the random walks naturally associated with a measure space $`(X,)`$ and a given measurable onto map $`\sigma :XX`$ such that $`\mathrm{\#}\sigma ^1\left(\left\{x\right\}\right)=N`$ for all $`xX`$, where $`N`$, $`2N<\mathrm{}`$, is fixed. Iteration of the branches
$$\sigma ^1\left(\left\{x\right\}\right):=\{yX\sigma \left(y\right)=x\}$$
(2.1)
then yields a combinatorial tree. If
$$\omega =(\omega _1,\omega _2,\mathrm{})\mathrm{\Omega }:=\{0,1,\mathrm{},N1\}^{},$$
an associated path may be thought of as an infinite extension of the finite walks
$$\tau _{\omega _n}\mathrm{}\tau _{\omega _2}\tau _{\omega _1}x$$
starting at $`x`$, where $`\left(\tau _i\right)`$, $`i=0,1,\mathrm{},N1`$, is a system of inverses, i.e., where
$$\sigma \tau _i=\underset{X}{id},0i<N.$$
(2.2)
If $`W:X[\mathrm{\hspace{0.17em}0},1]`$ is a given measurable function such that
$$\underset{y:\sigma \left(y\right)=x}{}W\left(y\right)=1,$$
(2.3)
then an associated measure $`P_x`$ on $`\mathrm{\Omega }`$ may be defined as follows. Suppose some function $`fC\left(\mathrm{\Omega }\right)`$ depends only on a finite number of coordinates, say $`\omega _1,\mathrm{},\omega _n`$, then set
$$_\mathrm{\Omega }fdP_x=\underset{(\omega _1,\mathrm{},\omega _n)}{}f(\omega _1,\mathrm{},\omega _n)W\left(\tau _{\omega _1}x\right)W\left(\tau _{\omega _2}\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right).$$
(2.4)
Extensions of this formula to $`\mathrm{\Omega }`$ can be done in a number of ways: see the cited references.
A special feature of this construction, which will be explored in the present monograph, is that of attractive convergence properties for infinite products of the form
$$\underset{n,\omega _1,\mathrm{},\omega _n}{}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)$$
(2.5)
over certain subsets of $`\mathrm{\Omega }`$. As it turns out, these infinite products are determined by the measures $`\left(P_x\right)_{xX}`$, and by the *Ruelle transition operator*
$$\left(R_Wg\right)\left(x\right):=\underset{y:\sigma \left(y\right)=x}{}W\left(y\right)g\left(y\right),gL^{\mathrm{}}\left(X\right).$$
(2.6)
The operator $`R`$ in (2.6) is called the transition operator, the Ruelle operator, or the Ruelle–Perron–Frobenius operator, and it will play a major role in what follows.
Many problems in dynamics are governed by transition probabilities $`W`$, and $`P_x`$, and by an associated transition operator $`R_W`$ as in (2.6). Wavelet theory is a case in point, and we show that fundamental convergence questions for wavelets, and properties of the solutions to (1.1), depend on the positive solutions $`h`$ to the eigenvalue problem $`R_W(h)=h`$. In Theorem 5.2, we show that there is a particular solution $`h`$, see (3.11), which determines the issue of pointwise convergence of the infinite product. We refer to the discussion around (3.11)–(3.12) below. The solutions $`h`$ are called *harmonic*, and the function $`h`$ in (3.11) is a special harmonic function which will play a central role in our main result, Theorem 5.2 below.
## 3 Multiresolutions
The multiresolution approach to wavelets involves functions on $``$. It begins with the fixed-point problem
$$\phi \left(t\right)=N\underset{k}{}a_k\phi \left(Ntk\right),t,$$
(3.1)
where a given sequence $`\left(a_k\right)_k`$ is chosen with special filtering properties, e.g., quadrature mirror filters; see . The equation (3.1) is called the scaling identity, and the $`a_k`$’s the response coefficients, or the masking coefficients. Introducing the Fourier series
$$m\left(x\right)=\underset{k}{}a_ke^{i2\pi kx}$$
(3.2)
and Fourier transform
$$\widehat{\phi }\left(x\right)=_{}e^{i2\pi tx}\phi \left(t\right)𝑑t,$$
(3.3)
we get the relation
$$\widehat{\phi }\left(x\right)=m\left(\frac{x}{N}\right)\widehat{\phi }\left(\frac{x}{N}\right),x,$$
(3.4)
which suggests a closer inspection of the infinite products
$$\underset{n=1}{\overset{\mathrm{}}{}}m\left(\frac{x}{N^n}\right).$$
(3.5)
Since we shall want solutions $`\phi `$ to (3.1) which are in $`L^2\left(\right)`$, (3.4)–(3.5) suggest the corresponding convergence questions for the function $`W:=\left|m\right|^2`$.
When $`\left(P_x\right)_{x[\mathrm{\hspace{0.17em}0},1]}`$ is the family of measures on $`\mathrm{\Omega }`$ corresponding to $`W=\left|m\right|^2`$, then the formal infinite product
$$\left|\widehat{\phi }\left(x\right)\right|^2=\underset{n=1}{\overset{\mathrm{}}{}}W\left(\frac{x}{N^n}\right)$$
(3.6)
is $`P_x\left(\left\{(0,0,0,\mathrm{})\right\}\right)`$, i.e., the measure of the singleton $`(0,0,\mathrm{})`$ (an infinite string of zeroes) in $`\{0,\mathrm{},N1\}^{}`$. Our main theorem (Theorem 5.2) gives a necessary and sufficient condition for the pointwise convergence of (3.6) in a rather general context. We note that because of assumption (2.3), $`\phi `$ will automatically be in $`L^2()`$, once pointwise convergence is established. There is a natural way (based on Euclid’s algorithm) of embedding $``$ into
$$\mathrm{\Omega }=\{0,\mathrm{},N1\}^{}\times \{0,\mathrm{},N1\}^{}$$
(3.7)
such that
$$P_x\left(\right)=\underset{k}{}\left|\widehat{\phi }\left(x+k\right)\right|^2.$$
(3.8)
Even though the measures $`P_x`$ (Section 6) are defined *a priori* on the over-countable probability space, the surprise is that they are in fact supported only on a fixed thin (countable) subset of $`\mathrm{\Omega }`$.
###### Remark 3.1.
Note that, in general, it is not at all clear that the measures $`\left(P_x\right)`$, $`xX`$, on $`\mathrm{\Omega }`$ should even have atoms. Typically, they don’t! But if atoms exist, i.e., when there are points $`\omega \mathrm{\Omega }`$ such that $`P_x\left(\left\{\omega \right\}\right)>0`$, we note that this yields convergence of an associated infinite product. Let $`_0:=\{0,1,2,\mathrm{}\}=\left\{0\right\}`$. Using Euclid, and the $`N`$-adic expansion
$$k=i_1+i_2N+\mathrm{}+i_nN^{n1}\text{ for }k_0,$$
(3.9)
we see that the points
$$\omega \left(k\right)=(i_1,\mathrm{},i_n,\underset{\mathrm{}\text{ string of zeroes}}{\underset{}{0,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},\mathrm{}}})$$
represent a copy of $`_0`$ sitting in $`\mathrm{\Omega }`$. With the identification $`k\omega \left(k\right)`$, we set
$$P_x\left(_0\right)=\underset{k=0}{\overset{\mathrm{}}{}}P_x\left(\left\{\omega \left(k\right)\right\}\right).$$
(3.10)
But in general, this function $`P_x\left(_0\right)`$ might be zero. Our first observation about
$$h\left(x\right):=P_x\left(_0\right),xX,$$
(3.11)
is that it solves the eigenvalue problem
$$R_Wh=h.$$
(3.12)
We say that $`h`$ is a minimal harmonic function relative to $`R_W`$. Note that in general $`h=0`$ may happen!
###### Remark 3.2.
Let $`X`$, $``$, $`\sigma `$, $`\tau _0`$, $`\mathrm{}`$, $`\tau _{N1}`$, and $`W`$ be as described above, and let $`\left(P_x\right)_{xX}`$ be the corresponding transition probabilities. Let
$$\mathrm{𝟎}=\underset{\mathrm{}\text{ string of zeroes}}{(\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},\mathrm{})}\mathrm{\Omega }=\{0,1,\mathrm{},N1\}^{}.$$
While in general, often $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)=0`$, the case when $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)>0`$ is important. The condition $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)>0`$ is a way of making precise sense of the infinite product
$$P_x\left(\left\{\mathrm{𝟎}\right\}\right)=\underset{n=1}{\overset{\mathrm{}}{}}W\left(\tau _0^nx\right).$$
(3.13)
If, for example, $`lim_n\mathrm{}W\left(\tau _0^nx\right)<1`$, then it is immediate from (3.13) that $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)=0`$.
Suppose $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)>0`$. Then it follows that
$$P_x\left(\left\{(i_1,\mathrm{},i_n,\mathrm{𝟎})\right\}\right)=W\left(\tau _{i_1}x\right)\mathrm{}W\left(\tau _{i_n}\mathrm{}\tau _{i_1}x\right)P_{\tau _{i_n}\mathrm{}\tau _{i_1}x}\left(\left\{\mathrm{𝟎}\right\}\right).$$
(3.14)
Using (3.9), we shall identify $`k_0`$ with the point $`\omega \left(k\right)\mathrm{\Omega }`$, and write $`P_x\left(\left\{k\right\}\right)`$ for the expression in (3.14).
An important question for dyadic wavelets in $`L^2\left(\right)`$ is the issue of when these wavelets form *orthonormal bases* (ONB’s). A dyadic wavelet function $`\psi L^2\left(\right)`$ generates an ONB if the double-indexed family
$$\left\{\mathrm{\hspace{0.17em}2}^{n/2}\psi \left(2^ntk\right)\right|n,k\}$$
(3.15)
satisfies (i) and (ii) below:
1. $`{\displaystyle _{}}\overline{\psi _{n,k}\left(t\right)}\psi _{m,l}\left(t\right)𝑑t=\delta _{n,m}\delta _{k,l}`$, with
$$\psi _{n,k}\left(t\right):=2^{n/2}\psi \left(2^ntk\right),$$
(3.16)
and
1. the closed linear span of $`\{\psi _{n,k}n,k\}`$ is $`L^2\left(\right)`$.
In our analysis of the scaling identity
$$\phi \left(t\right)=2\underset{k}{}a_k\phi \left(2tk\right)$$
(3.17)
(a special case of (3.1)), we will be looking at two functions $`\phi `$ and $`\psi `$; the second one may be taken to be
$$\psi \left(t\right)=2\underset{k}{}\left(1\right)^{k+1}\overline{a}_{1k}\phi \left(2tk\right).$$
(3.18)
This analysis is the approach to wavelets which goes under the name of *multiresolution analysis*. The function $`\psi `$ which is used in (3.15)–(3.16) is the solution to (3.18). The two standing conditions which are placed on the numbers $`\left(a_k\right)_k`$, called *masking coefficients*, are
$$\underset{k}{}\overline{a}_ka_{k+2n}=\frac{1}{2}\delta _{0,n},n,$$
(3.19)
and
$$\underset{k}{}a_k=1.$$
(3.20)
These conditions in themselves do not imply orthonormality in (3.16), but only the following much weaker property:
$$\underset{n,k}{}\left|\psi _{n,k}f\right|^2=f^2=_{}\left|f\left(t\right)\right|^2𝑑t,fL^2\left(\right).$$
(3.21)
A system of functions $`\left(\psi _{n,k}\right)`$ satisfying (3.21) is called a *Parseval frame*, or a *normalized tight frame*.
Given (3.19)–(3.20), it turns out that the ONB property for the wavelet is equivalent to either one of the following two conditions for the normalized scaling function $`\phi `$:
$$_{}\overline{\phi \left(t\right)}\phi \left(tk\right)𝑑t=\delta _{0,k},$$
(3.22)
or
$$\phi _{L^2\left(\right)}=1.$$
(3.23)
## 4 Sampling
In this section and the next, we study an intriguing relationship between the following three problems:
1. When does the scaling identity (3.1) have $`L^2`$ solutions?
2. How may the transition probabilities $`P_x`$ be used in *sampling* certain functions at the points $`x+k`$ as $`k`$ runs over a set of integers?
3. When is the infinite product (3.5) pointwise convergent?
In the wavelet applications, $`X=[\mathrm{\hspace{0.17em}0},1]`$, and the system $`\sigma ,\tau _0,\mathrm{},\tau _{N1}`$ is as follows:
$$\{\begin{array}{cc}& \sigma \left(x\right)=Nxmod1,\hfill \\ & \tau _j\left(x\right)=\frac{x+j}{N},j=0,1,\mathrm{},N1.\hfill \end{array}$$
(4.1)
###### Lemma 4.1.
Setting $`F\left(x\right):=P_x\left(\left\{\mathrm{𝟎}\right\}\right)`$, and
$$k=i_1+i_2N+\mathrm{}+i_nN^{n1}(_0),$$
(4.2)
we get the formula
$$P_x\left(\left\{k\right\}\right)=F\left(x+k\right),$$
(4.3)
where we have identified $`k`$ with the point $`\omega \left(k\right):=(i_1,\mathrm{},i_n,\mathrm{𝟎})`$ in $`\mathrm{\Omega }`$.
###### Proof..
To see this, identify functions on $`[\mathrm{\hspace{0.17em}0},1]`$ with $`1`$-periodic functions on $``$, and note that the second formula in (4.1) yields $`\tau _{i_n}\mathrm{}\tau _{i_1}\left(x\right)=\left(x+k\right)/N^n`$ where $`k`$ is given by (4.2). Hence, if $`1sn`$, then
$$W\left(\tau _{i_s}\mathrm{}\tau _{i_1}\left(x\right)\right)=W\left(\frac{x+i_1N+\mathrm{}+i_sN^{s1}}{N^s}\right)=W\left(\frac{x+k}{N^s}\right).$$
It follows that (4.3) is really just a rewrite of (3.14). The right-hand side of (3.14) yields
$$W\left(\frac{x+k}{N}\right)\mathrm{}W\left(\frac{x+k}{N^n}\right)F\left(\frac{x+k}{N^n}\right)=F\left(x+k\right),$$
(4.4)
which is the desired conclusion. ∎
###### Remark 4.2.
To extend $`P_x()`$ from $`_0`$ to $``$, recall that $`k_0`$ is identified with the singleton $`\omega \left(k\right)=(i_1,\mathrm{},i_n,\mathrm{𝟎})`$ via (4.2). Now, if $`N^nk<0`$, then set
$$P_x\left(\left\{k\right\}\right):=P_x\left(\left\{\omega \left(N^{n+1}+k\right)\right\}\right).$$
(4.5)
To help the reader gain some intuitive feeling for the conclusion in Lemma 4.1, observe that the right-hand side of (4.3) represents a *sampling* of the function $`F`$ at the integral translates on $``$, starting at $`x`$, i.e., $`x+k`$. Obviously, different subsets of $`\mathrm{\Omega }`$ would yield different sets of sampling points for $`F`$, including nonuniformly distributed sampling points; see .
Starting with Shannon , the theory of sampling has emerged as a significant tool in signal processing; see, e.g., the beautifully written survey as well as the references cited therein. Thus, in a general context, our formula (4.3) offers a probabilistic prescription for sampling of functions on the real line, and at the same time it stresses the ‘random’ feature of sampling.
## 5 A convergence theorem for infinite products
We now show how this viewpoint from sampling theory is closely related to some fundamental properties of the measures $`P_x`$. In particular, our Theorem 5.2 gives a necessary and sufficient condition for pointwise convergence of the infinite product (3.5), or more generally (2.5), with the condition for convergence stated in terms of the harmonic function $`h`$ of (3.11). The relationship between $`h`$ and the measure family $`P_x`$ is studied more systematically below.
Let $`A`$. Returning to (4.3), we set
$$P_x\left(A\right):=\underset{kA}{}P_x\left(\left\{k\right\}\right)=\underset{kA}{}F\left(x+k\right).$$
(5.1)
As in Lemma 4.1, the number $`N`$, $`N2`$, and the function $`W`$ are given. The measures $`P_x`$ are constructed from these data using (2.4), and we have the two functions
$$F\left(x\right):=P_x(\{(\underset{\mathrm{}\text{ string of zeroes}}{\underset{}{0,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},\mathrm{}}})\})$$
(5.2)
and
$$h\left(x\right):=P_x\left(\right),x.$$
(5.3)
Finally, for $`k`$, set
$$N^k:=\{N^kjj\}.$$
(5.4)
Using (4.5), we see that $`N^k`$ is represented in $`\mathrm{\Omega }=\{0,1,\mathrm{},N1\}^{}`$ as
$$(\underset{k\text{ zeroes}}{\underset{}{0,\mathrm{},\mathrm{\hspace{0.33em}0}}},\underset{\begin{array}{c}\text{a finite string}\\ \text{of symbols}\\ \omega _i\{0,\mathrm{},N1\}\end{array}}{\underset{}{\omega _1,\omega _2,\mathrm{},.}},\underset{\mathrm{}\text{ string of zeroes}}{\underset{}{0,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},\mathrm{}}}).$$
(5.5)
An infinite string of zeroes will be denoted $`\mathrm{𝟎}`$.
###### Lemma 5.1.
Let $`N`$, $`W`$, $`F`$, $`P_x`$, and $`h`$ be as described above; see (5.2)–(5.3). Then $`h`$ satisfies the following cocycle identity:
$$h\left(x\right)W\left(x\right)=P_{Nx}\left(N\right),x.$$
(5.6)
In particular, if $`h\left(x\right)1`$ $`\mathrm{a}.\mathrm{e}.x`$, then we recover the function $`W`$ from the transition probabilities $`P_x`$, $`\mathrm{a}.\mathrm{e}.x`$.
###### Proof..
We calculate the left-hand side in (5.6), using the earlier equations:
$`h\left(x\right)W\left(x\right)`$ $`\underset{\begin{array}{c}\text{by (}\text{4.3}\text{)}\\ \text{and (}\text{5.3}\text{)}\end{array}}{=}W\left(x\right)_jF\left(x+j\right)`$
$`\underset{\text{by periodicity}}{=}{\displaystyle \underset{j}{}}W\left(x+j\right)F\left(x+j\right)`$
$`\underset{\text{by (}\text{4.4}\text{)}}{=}{\displaystyle \underset{j}{}}F\left(N\left(x+j\right)\right)`$
$`\underset{\text{ }}{=}{\displaystyle \underset{j}{}}F\left(Nx+Nj\right)`$
$`\underset{\text{by (}\text{4.3}\text{)}}{=}{\displaystyle \underset{j}{}}P_{Nx}\left(\left\{Nj\right\}\right)`$
$`\underset{\text{by (}\text{5.1}\text{)}}{=}P_{Nx}\left(N\right).`$
This is the desired identity (5.6), and the proof is completed. ∎
###### Theorem 5.2.
Let $`N`$, $`W`$, $`F`$, $`P_x`$, and $`h`$ be as described above. Let $`x`$, and suppose that $`P_x\left(\left\{\mathrm{𝟎}\right\}\right)>0`$. Then the following two conditions are equivalent.
1. The limit on the right-hand side below exists, and
$$F\left(x\right)=\underset{n\mathrm{}}{lim}\underset{k=1}{\overset{n}{}}W\left(\frac{x}{N^k}\right).$$
2. The limit on the left-hand side below exists, and
$$\underset{n\mathrm{}}{lim}h\left(\frac{x}{N^n}\right)=1.$$
###### Proof..
(a)$``$(b). An iteration of the identity (5.6) in Lemma 5.1 above yields
$$P_x\left(N^k\right)=\left(\underset{j=1}{\overset{k}{}}W\left(\frac{x}{N^j}\right)\right)h\left(\frac{x}{N^k}\right).$$
(5.7)
Using (5.5), and working in $`\mathrm{\Omega }`$, we find
$$\underset{k}{}N^k=\left\{\mathrm{𝟎}\right\}.$$
(5.8)
An application of a standard result in measure theory \[41, Theorem 1.19(e), p. 16\] now yields existence of the following limit:
$$\underset{k\mathrm{}}{lim}P_x\left(N^k\right)=P_x\left(\left\{\mathrm{𝟎}\right\}\right)\underset{\text{by (}\text{5.2}\text{)}}{=}F\left(x\right).$$
(5.9)
Since (a) is assumed, and $`F\left(x\right)>0`$, we conclude that the limit $`h\left(x/N^k\right)`$, for $`k\mathrm{}`$, must exist as well, and further that
$$F\left(x\right)=F\left(x\right)\underset{k\mathrm{}}{lim}h\left(\frac{x}{N^k}\right).$$
Using again $`F\left(x\right)>0`$, we finally conclude that (b) holds.
(b)$``$(a). Recall that formula (5.7) holds in general. If (b) is assumed, we then conclude that the limit $`_{j=1}^kW\left(x/N^j\right)`$ exists as $`k\mathrm{}`$. The limit on the left-hand side in (5.7) exists and is $`F\left(x\right)`$. As a result, we get
$$F\left(x\right)=\underset{k\mathrm{}}{lim}\underset{j=1}{\overset{k}{}}W\left(\frac{x}{N^j}\right)\underset{k\mathrm{}}{lim}h\left(\frac{x}{N^k}\right)\underset{\text{using (}\text{b}\text{)}}{=}\underset{k\mathrm{}}{lim}\underset{j=1}{\overset{k}{}}W\left(\frac{x}{N^j}\right),$$
which is (a). ∎
## 6 Ruelle’s wavelet transition operator
The next definitions (Definitions 6.1) and the lemma (Lemma 6.2) give us the precise details behind the two critical notions from Theorem 5.2 above; i.e., the cocycles, and the random walk measures $`P_x`$. An existence and uniqueness theorem for $`P_x`$ is then presented in Section 8.
###### Definitions 6.1.
1. The *Ruelle operator* $`R=R_W`$ is defined by
$$\left(Rf\right)\left(x\right)=\underset{yX,\sigma \left(y\right)=x}{}W\left(y\right)f\left(y\right),xX,fL^{\mathrm{}}\left(X\right),$$
(6.1)
and maps $`L^{\mathrm{}}\left(X\right)`$ into itself.
2. Let $`\mathrm{\Omega }`$ be the compact Cartesian product
$$\mathrm{\Omega }=_N^{}=\{0,\mathrm{},N1\}^{}=\underset{1}{\overset{\mathrm{}}{}}\{0,\mathrm{},N1\}.$$
(6.2)
3. A bounded measurable function $`V:X\times \mathrm{\Omega }`$ is said to be a *cocycle* if
$$V(x,(\omega _1,\omega _2,\mathrm{}))=V(\tau _{\omega _1}\left(x\right),(\omega _2,\omega _3,\mathrm{}))$$
(6.3)
for all $`\omega =(\omega _1,\omega _2,\mathrm{})\mathrm{\Omega }`$.
4. A function $`h:X`$ is said to be *harmonic*, or $`R_W`$*-harmonic*, if
$$R_Wh=h.$$
(6.4)
5. Let $`n`$, and let $`i_1,\mathrm{},i_n_N`$. Then the subset
$$A(i_1,\mathrm{},i_n):=\{w\mathrm{\Omega }\omega _1=i_1,\mathrm{},\omega _n=i_n\}$$
(6.5)
is called a *cylinder set*.
We shall use the following correspondence (see and for details and proofs) between the cocycles $`V`$ from (c) and the harmonic functions $`h`$ from (d): If $`V`$ is given as in (c), then the function $`h(x):=P_x[V(x,)]`$ has the properties in (d). Conversely, for every $`h`$ satisfying (d), including (6.4), there is a martingale limit which lets us recover a cocycle $`V`$, $`P_x`$ a.e., such that $`h(x):=P_x[V(x,)]`$. In the present discussion, we are interested in the cocycles $`V`$ which arise as indicator functions $`\chi _S`$ for certain cyclic and invariant subsets $`S`$ of $`\mathrm{\Omega }`$. We refer to Section 8 below.
### Existence of the measures $`P_x`$
The cylinder sets generate the topology of $`\mathrm{\Omega }`$, and its Borel sigma-algebra. In determining Radon measures on $`\mathrm{\Omega }`$, it is therefore convenient to first specify them on cylinder sets. This approach was initiated by Kolmogorov ; see also Nelson . Recall that $`\mathrm{\Omega }`$ is compact in the Tychonoff topology, and that we may use the Stone–Weierstraß theorem on $`C\left(\mathrm{\Omega }\right)=`$the algebra of all continuous functions on $`\mathrm{\Omega }`$.
###### Lemma 6.2.
Let $`X`$, $``$, $`\mu `$, $`\sigma `$, $`\tau _0`$, $``$, $`\tau _{N1}`$, and $`W`$ be given as described above. We make the following more restrictive assumption on $`W`$:
$$\underset{yX,\sigma \left(y\right)=x}{}W\left(y\right)=1\mathrm{a}.\mathrm{e}.xX.$$
(6.6)
Then for every $`xX`$ there is a unique positive Radon probability measure $`P_x`$ on $`\mathrm{\Omega }`$ such that
$$P_x\left(A(i_1,\mathrm{},i_n)\right)=W\left(\tau _{i_1}x\right)W\left(\tau _{i_2}\tau _{i_1}x\right)\mathrm{}W\left(\tau _{i_n}\mathrm{}\tau _{i_1}x\right).$$
(6.7)
The main fact about the Ruelle operator
$$R_Wf\left(x\right)=\underset{i}{}W\left(\tau _ix\right)f\left(\tau _ix\right)$$
is this : Under suitable conditions on $`W`$, there is a unique probability measure (Ruelle measure $`\nu `$) satisfying
$$\nu R_W=\nu ,$$
and a unique continuous minimal eigenfunction $`h_{\mathrm{min}}`$,
$$R_Wh_{\mathrm{min}}=h_{\mathrm{min}}.$$
This function $`h_{\mathrm{min}}`$ is minimal in the ordered convex set
$$\left\{hC\left(X\right)\right|0h1,R_Wh=h,_Xh𝑑\nu =1\},$$
and moreover
$$h_{\mathrm{min}}\left(x\right)=P_x\left(\right).$$
###### Remark 6.3.
It turns out that the general case $`_y\mathrm{}1`$ may be reduced to (6.6). So (6.6) is not really a restriction.
###### Proof of Lemma 6.2.
If $`P`$ is a Radon measure on $`\mathrm{\Omega }`$, we set
$$P\left[f\right]:=_\mathrm{\Omega }f\left(\omega \right)𝑑P\left(\omega \right)\text{ for all }fC\left(\mathrm{\Omega }\right).$$
(6.8)
Set
$$\begin{array}{c}C_{fin}\left(\mathrm{\Omega }\right)=\{fC\left(\mathrm{\Omega }\right)n\text{ such that }f\left(\omega \right)=f(\omega _1,\mathrm{},\omega _n)\text{}\hfill \\ \hfill \text{i.e., }f\text{ depends only on the first }n\text{ coordinates in }\mathrm{\Omega }\}.\end{array}$$
(6.9)
Increasing $`n`$ in definition (6.9), we get an ascending nest of subalgebras of $`C\left(\mathrm{\Omega }\right)`$,
$$𝔄_1𝔄_2\mathrm{}𝔄_n𝔄_{n+1}\mathrm{},$$
(6.10)
with
$$\overline{\underset{n=1}{\overset{\mathrm{}}{}}𝔄_n}=C\left(\mathrm{\Omega }\right),$$
where $`\overline{}`$ stands for norm-closure. We set
$$C_{fin}\left(\mathrm{\Omega }\right):=\underset{n=1}{\overset{\mathrm{}}{}}𝔄_n.$$
An immediate application of Stone–Weierstraß shows that $`C_{fin}\left(\mathrm{\Omega }\right)`$ is uniformly dense in $`C\left(\mathrm{\Omega }\right)`$. Let $`xX`$, and $`fC_{fin}\left(\mathrm{\Omega }\right)`$. Suppose
$$f\left(\omega \right)=f(\omega _1,\mathrm{},\omega _n),$$
and set
$$P_x\left[f\right]=\underset{(\omega _1,\mathrm{},\omega _n)_N^n}{}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)f(\omega _1,\mathrm{},\omega _n).$$
(6.11)
Note that, if there is some $`hL^{\mathrm{}}\left(X\right)`$ such that
$$f(\omega _1,\mathrm{},\omega _n)=h\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right),$$
then
$`P_x\left[f\right]`$ $`={\displaystyle \underset{(\omega _1,\mathrm{},\omega _n)}{}}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)h\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)`$ (6.12)
$`={\displaystyle \underset{yX,\sigma ^ny=x}{}}W\left(\sigma ^{n1}y\right)\mathrm{}W\left(\sigma y\right)W\left(y\right)h\left(y\right)`$
$`=\left(R_W^nh\right)\left(x\right).`$
We now show that $`P_x\left[f\right]`$ is well-defined. This is the Kolmogorov consistency: we must check that the number $`P_x\left[f\right]`$ is the same when some $`f𝔄_n`$ ($`𝔄_{n+1}`$) is viewed also as an element in $`𝔄_{n+1}`$. Then
$`f\left(\omega \right)`$ $`=f(\omega _1,\mathrm{},\omega _n)`$
$`=f(\omega _1,\mathrm{},\omega _n,\omega _{n+1}),`$
and
$`P_x\left[f_{n+1}\right]`$ $`={\displaystyle \underset{\omega _1,\mathrm{},\omega _{n+1}}{}}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _{n+1}}\mathrm{}\tau _{\omega _1}x\right)f(\omega _1,\mathrm{},\omega _{n+1})`$
$`={\displaystyle \underset{\omega _1,\mathrm{},\omega _n}{}}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)`$
$`\left(\underset{=1\text{ by (}\text{6.6}\text{)}}{\underset{}{{\displaystyle \underset{\omega _{n+1}}{}}W\left(\tau _{\omega _{n+1}}\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)}}\right)f(\omega _1,\mathrm{},\omega _n)`$
$`={\displaystyle \underset{\omega _1,\mathrm{},\omega _n}{}}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)f(\omega _1,\mathrm{},\omega _n)`$
$`=P_x\left[f_n\right],`$
as claimed.
The consistency conditions may be stated differently in terms of conditional probabilities: for $`fC\left(\mathrm{\Omega }\right)`$, set
$`P_x^{\left(n\right)}\left[f\right]`$ $`=P_x\left[f𝔄_n\right]`$ (6.13)
$`={\displaystyle \underset{(\omega _1,\mathrm{},\omega _n)}{}}W\left(\tau _{\omega _1}x\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}x\right)f(\omega _1,\mathrm{},\omega _n).`$
We proved that
$$P_x^{\left(n\right)}\left[f\right]=P_x^{\left(n+1\right)}\left[f\right]\text{ for all }f𝔄_n.$$
Using now the theorems of Stone–Weierstraß and Riesz, we get the existence of the measure $`P_x`$ on $`\mathrm{\Omega }`$. It is clear that it has the desired properties. In particular, the property (6.7) results from applying (6.11) to the function
$$f\left(\omega \right):=\delta _{i_1,\omega _1}\mathrm{}\delta _{i_n,\omega _n},\omega \mathrm{\Omega },$$
(6.14)
when the point $`(i_1,\mathrm{},i_n)`$ is fixed.
These functions, in turn, span a dense subalgebra in $`C\left(\mathrm{\Omega }\right)`$ (by Stone–Weierstraß), so $`P_x`$ is determined uniquely by (6.7). ∎
## 7 Transition measures
We now compute the Markov transition measures $`P_x`$, as $`x`$ varies over the set $`X`$. The construction of $`P_x`$, and the probability space $`\mathrm{\Omega }`$, begins with a sequence of measures $`P_x^{\left(n\right)}`$, $`n=1,2,\mathrm{}`$, corresponding to finite paths of length $`n`$. Then we show that the infinite path-space measure, Lemma 6.2, results from an application of the Kolmogorov extension principle. Our technical analysis involves a certain transition operator $`R`$ which generalizes Lawton’s wavelet transition operator; see .
We already mentiond how the measures $`P_x`$ serve to prescribe sampling of functions on the real line, Lemma 4.1. However, a deeper understanding of this sampling viewpoint is facilitated by the introduction of the Perron–Frobenius–Ruelle operator $`R`$, Definitions 6.1, and an associated family of harmonic functions. For these harmonic functions, we note in Section 6 that there is a crucial analogue from classical harmonic analysis of Fatou-boundary function in the present discrete context. As noted after Definition 6.1 below, the boundary value functions take the form of certain cocycles, and the function $`h`$ from (3.11) corresponds to a special $`_0`$-cocycle. The existence of these cocycles is based on a martingale convergence theorem; see \[33, Theorem 2.7.1\] and for additional details.
A fundamental tightness condition for the random walk model is introduced. The transition probabilities $`P_x`$ live in a universal probability space $`\mathrm{\Omega }`$, but the essential convergence questions for the infinite products depend on the the $`P_x`$’s being supported on a certain copy of $`_0`$ (the natural numbers), or of $``$ (the integers). This refers to the natural embedded of $``$ in $`\mathrm{\Omega }`$ which we described above.
## 8 Kolmogorov’s consistency condition
For general reference, we now make explicit the extension principle of Kolmogorov in its function-theoretic form.
###### Lemma 8.1.
(Kolmogorov) Let $`N2`$ be fixed, and let
$$\mathrm{\Omega }=\{0,1,\mathrm{},N1\}^{}.$$
For $`n=1,2,`$, let
$$P^{\left(n\right)}:𝔄_n$$
be a sequence of linear functionals such that (i)–(iii) hold:
1. $`P^{\left(n\right)}\left[\text{1}\text{1}\right]=1`$, where 11 denotes the constant function $`1`$ on $`\mathrm{\Omega }`$,
2. $`f𝔄_n`$, $`f0`$ pointwise$`P^{\left(n\right)}\left[f\right]0`$,
and
1. $`P^{\left(n\right)}\left[f\right]=P^{\left(n+1\right)}\left[f\right]`$ for all $`f𝔄_n`$.
Then there is a unique Borel probability measure $`P`$ on $`\mathrm{\Omega }`$ such that
$$P\left[f\right]=P^{\left(n\right)}\left[f\right],f𝔄_n.$$
(8.1)
Specifically, for $`P`$, we have the implication
$$fC\left(\mathrm{\Omega }\right),f0\text{ pointwise}P\left[f\right]0.$$
(8.2)
###### Remark 8.2.
Here we have identified positive linear functionals $`P`$ on $`C\left(\mathrm{\Omega }\right)`$ with the corresponding Radon measures $`\stackrel{~}{P}`$ on $`\mathrm{\Omega }`$, i.e.,
$$P\left[f\right]=_\mathrm{\Omega }f𝑑\stackrel{~}{P}.$$
(8.3)
This identification $`P\stackrel{~}{P}`$ is based on an implicit application of Riesz’s theorem; see \[41, Chapter 1\].
###### Proof of Lemma 8.1.
The proof of Kolmogorov’s extension result may be given several forms, but we note that the argument we used above (in a special case), based on an application of the Stone–Weierstraß theorem, also works in general. ∎
### The probability space $`\mathrm{\Omega }`$
We note that the probability space $`\mathrm{\Omega }`$ itself carries mappings $`\sigma `$ and $`\tau _i`$, $`i=0,1,\mathrm{},N1`$. Stressing the $`\mathrm{\Omega }`$-dependence, we write
$$\sigma ^\mathrm{\Omega }\left(\omega \right)=(\omega _2,\omega _3,\mathrm{})\text{ and }\tau _i^\mathrm{\Omega }\left(\omega \right)=(i,\omega _1,\omega _2,\mathrm{})$$
(8.4)
for $`\omega =(\omega _1,\omega _2,\mathrm{})\mathrm{\Omega }`$.
###### Remark 8.3.
The connection between the cylinder sets in (6.5) and the iterated function systems (IFS) $`(X,\sigma ,\tau _0,\mathrm{},\tau _{N1})`$ may be spelled out as follows: the cylinder sets in $`\mathrm{\Omega }`$ generate the sigma-algebra of measurable subsets of $`\mathrm{\Omega }`$, and similarly the subsets $`\tau _{i_1}\mathrm{}\tau _{i_n}\left(X\right)X`$ generate a sigma-algebra of measurable subsets of $`X`$. When nothing further is specified, these will be the sigma-algebras which we refer to when discussing the measurable functions on $`\mathrm{\Omega }`$ and $`X.`$ In particular, we will denote by $`M\left(\mathrm{\Omega }\right)`$ and $`M\left(X\right)`$ the respective algebras of all bounded measurable functions on $`\mathrm{\Omega }`$, respectively $`X`$.
Note that if $`X=[\mathrm{\hspace{0.17em}0},1]`$, and if
$$\tau _ix=\frac{x+i}{N},0iN1,$$
then we recover the familiar $`N`$-adic subintervals:
$$\tau _{i_1}\mathrm{}\tau _{i_n}\left(X\right)=[\frac{i_1}{N}+\mathrm{}+\frac{i_n}{N^n},\frac{i_1}{N}+\mathrm{}+\frac{i_n}{N^n}+\frac{1}{N^n}].$$
(8.5)
###### Lemma 8.4.
There is a unique mapping $`\rho :M\left(\mathrm{\Omega }\right)M\left(X\right)`$ which satisfies
$$\rho \left(fg\right)=\rho \left(f\right)\rho \left(g\right)$$
(8.6)
and
$$\rho \left(\chi _{A(i_1,\mathrm{},i_n)}\right)=\chi _{\tau _{i_1}\mathrm{}\tau _{i_n}\left(X\right)}.$$
(8.7)
The mapping $`\rho `$ is an isomorphism of $`M\left(\mathrm{\Omega }\right)`$ onto $`M\left(X\right)`$.
###### Proof..
Recalling (6.14), we note that
$$\chi _{A(i_1,\mathrm{},i_n)}\left(\omega \right)=\delta _{i_1,\omega _1}\mathrm{}\delta _{i_n,\omega _n},\omega \mathrm{\Omega }.$$
As a result,
$$\chi _{A(i_1,\mathrm{},i_n)}\chi _{A(j_1,\mathrm{},j_n)}=\delta _{i_1,j_1}\mathrm{}\delta _{i_n,j_n}\chi _{A(i_1,\mathrm{},i_n)}.$$
(8.8)
We then define $`\rho `$ first on $`𝔄_n`$ by
$$\rho \left(\underset{i_1,\mathrm{},i_n}{}a_{i_1,\mathrm{},i_n}\chi _{A(i_1,\mathrm{},i_n)}\right)=\underset{i_1,\mathrm{},i_n}{}a_{i_1,\mathrm{},i_n}\chi _{\tau _{i_1}\mathrm{}\tau _{i_n}\left(X\right)},$$
where $`a_{i_1,\mathrm{},i_n}`$, and note that (8.6) is satisfied.
It is easy to check that the extension of $`\rho `$ from $`𝔄_n`$ to $`𝔄_{n+1}`$ is consistent. The final extension from $`_n𝔄_n`$ to $`M\left(\mathrm{\Omega }\right)`$ is done by Kolmogorov’s lemma, and it can be checked that $`\rho `$ has the properties stated in the conclusion of the present lemma. ∎
###### Theorem 8.5.
Let $`X`$, $`W`$, and $`N`$ be as described in the beginning of this chapter, and let $`\{P_xxX\}`$ be the process obtained in the conclusion of Lemma 6.2. Then
$$\underset{i=0}{\overset{N1}{}}W\left(\tau _ix\right)P_{\tau _ix}\left[f(i,)\right]=P_x\left[f\right]\text{ for all }fC\left(\mathrm{\Omega }\right).$$
(8.9)
Moreover, equation (8.9) determines $`\left(P_x\right)`$ uniquely.
###### Remark 8.6.
Stated informally, formula (8.9) is an assertion about the random walk: it says that if the walk starts at $`x`$, then with probability one, it makes a transition to one of the $`N`$ points $`\tau _0\left(x\right)`$, $`\mathrm{}`$, $`\tau _{N1}\left(x\right)`$. The probability of the move $`x\tau _ix`$ is $`W\left(\tau _ix\right)`$. Recall (6.6) asserts that $`_iW\left(\tau _ix\right)=1`$.
###### Proof of Theorem 8.5.
It follows from (6.11) and the arguments in the proof of Lemma 6.2 that it is enough to verify (8.9) for $`fC_{fin}\left(\mathrm{\Omega }\right)`$, or for $`f𝔄_n`$. Let $`f𝔄_n`$. Then
$`{\displaystyle \underset{i}{}}W\left(\tau _ix\right)P_{\tau _ix}\left[f(i,)\right]`$
$`={\displaystyle \underset{i}{}}{\displaystyle \underset{\omega _1,\mathrm{},\omega _n}{}}W\left(\tau _ix\right)W\left(\tau _{\omega _1}\tau _ix\right)\mathrm{}W\left(\tau _{\omega _n}\mathrm{}\tau _{\omega _1}\tau _ix\right)f(i,\omega _1,\mathrm{},\omega _n)`$
$`\underset{\text{by (}\text{6.11}\text{)}}{=}P_x\left[f\right],`$
which is the desired conclusion. Recall $`C_{fin}\left(\mathrm{\Omega }\right)`$ is norm-dense in $`C\left(\mathrm{\Omega }\right)`$.
Note that the formula (8.9) generalizes the familiar notion of selfsimilarity for measures introduced by Hutchinson in . In fact, (8.9) may be restated as
$$\underset{i=0}{\overset{N1}{}}W\left(\tau _ix\right)P_{\tau _ix}\left(\tau _i^\mathrm{\Omega }\right)^1=P_x.$$
(8.10)
The uniqueness part of the theorem follows then from Lemma 8.1. ∎
###### Acknowledgements.
The author is pleased to acknowledge numerous constructive discussions about infinite products with Professors Dorin Dutkay and Richard Gundy; with Akram Aldroubi, John Benedetto and other members of our NSF-Focused Research Group (FRG). And we thank Brian Treadway for outstanding typesetting and helpful suggestions.
|
warning/0506/hep-ph0506211.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Regge theory is widely used to describe the low $`p_T`$ high-energy interactions of hadrons, nuclei and (real and virtual) photons. The theory takes into account both Regge poles and cuts. The latter are related to the exchange of several reggeons in the $`t`$-channel. The status of this theory within QCD is reviewed, for example, in Refs. . The Pomeranchuk singularities (that is the Pomeron pole and the corresponding cuts) play a special role in this theory as they determine the high energy behaviour of diffractive processes and multiparticle production . It is important to understand the connection between the general results of reggeon theory and the space-time picture of hadronic interactions. This becomes possible due to the relation between regge theory and the parton model . Multiple pomeron exchanges are especially important if the intercept of the pomeron, $`\alpha _P(0)`$, is larger than unity, that is $`\mathrm{\Delta }\alpha _P(0)1>0`$. This so-called “supercritical” theory is favoured both by experimental data and by calculations in QCD perturbation theory . In this case the partonic interpretation becomes very non-trivial. The relation between the probabilistic partonic picture of the interaction and diagrams of reggeon theory has been studied in Refs. .
In this paper we discuss two simple analytic models of regge theory, which provide particular examples of the partonic picture of high energy hadronic collisions. These are the eikonal model and the Schwimmer model <sup>1</sup><sup>1</sup>1 The currently popular Balitski-Kovchegov equation is, from the partonic and space-time viewpoints, a generalization of the Schwimmer model., which are often used in phenomenological applications of regge theory. We use the AGK cutting rules to obtain the inelastic, diffractive and inclusive cross sections predicted by these models; and discuss the partonic interpretation of these results. Although some of these are known, it is informative to summarize them here. For the Schwimmer model, and its eikonal generalization, we obtain explicit formulae for the total, inelastic and diffractive cross sections. We also obtain the dependence of the differential cross section on the size of the rapidity gap.
Our ultimate goal is to use these results to improve the ‘global’ analysis of data for ‘soft’ high energy processes, see, for example, Ref. .
## 2 Multiparticle content of reggeon diagrams
An interpretation of reggeon diagrams in terms of their inner multiparticle structure was given in Ref.. It corresponds to the qualitative picture that a fast hadron of momentum $`p`$ (of rapidity $`y\mathrm{ln}2p/m`$) interacts with a target due to quantum-mechanical fluctuations containing slow particles. The structure of the fluctuations is rather specific and is usually called ‘multiperipheral’. Such a fluctuation contains $`\mathrm{ln}p`$ soft virtual particles ordered in their rapidities. For brevity we shall call these particles soft partons, or simply partons <sup>2</sup><sup>2</sup>2 Here we do not associate (soft) partons with definite objects like quarks, gluons, or pions, because only rather general features of them are relevant for our analysis. Hard partons of different spatial scales are considered in connection with deep inelastic scattering and other hard processes.. Only the slowest partons <sup>3</sup><sup>3</sup>3 These partons have momenta of the order of the typical hadronic scale $`\mu `$ of about several hundred MeV. have a chance to interact directly with a target. The faster partons of the fluctuation simply play the role of spectators. The cross section of the interaction is proportional to the number $`n(y)`$ of slow partons. In this scheme, slow partons originate from faster partons close in rapidity and $`n(y)`$ has the exponential behaviour $`\mathrm{exp}(\mathrm{\Delta }y)`$. The interaction of a single slow parton corresponds to regge-pole behaviour of cross section with $`\mathrm{\Delta }=\alpha _P(0)1`$, while the interactions of two or more partons with the target give rise to regge-cut-type contributions. If the (multiperipheral) evolution of a fast parton into slow ones is independent of the evolution of the other fast partons, then we obtain independent slow partons whose interactions correspond to ’non-enhanced’ reggeon diagrams of the eikonal approximation (see section 3).
For the supercritical pomeron, i.e. described by a regge pole with $`\mathrm{\Delta }>0`$, the number of slow partons increases exponentially with the initial rapidity $`y`$, i.e. in the course of the evolution of the parton fluctuation in rapidity space the number of partons multiplies, for example, by a splitting mechanism. As a consequence another type of reggeon diagrams will appear – that is ‘enhanced’ diagrams of the Schwimmer type occur (see section 4). We emphasize that for the parton dynamics to be consistent we require not only splitting, but also fusion of partons – though in special cases, we may, to a good approximation, neglect the latter process.
The Abramovsky-Gribov-Kancheli (AGK) cutting rules are a powerful tool for the investigation of the multiparticle structure of complicated reggeon diagrams. They were derived as the high-energy version of Cutkosky cutting rules . They give the discontinuity of the whole reggeon diagram in terms of the discontinuities of its component subdiagrams. Each reggeon diagram has various discontinuities which correspond to different ways of cutting the diagram and to different intermediate states. For example, cutting the regge pole diagram corresponds to the simple multiperipheral intermediate state. On the other hand, cutting a double-pomeron-exchange diagram leads to intermediate states of both double and single density, depending on the number of cut pomerons, and also on the state with a large rapidity gap obtained when the diagram is cut between pomerons. The AGK rules give relations between the contributions of given reggeon diagrams to different multiparticle cross sections. Examples of such relations will be discussed below.
The space-time picture of the interaction is another valuable tool in the description of high energy collisions. Pomeron exchange is a highly non-local process. It is characterised by longitudinal and time scales which are proportional to the initial energy. As a consequence, only reggeon diagrams with so called non-planar structure contribute at high energies. Partonic fluctuations for these diagrams develop simultaneously at comparable longitudinal distances. In contrast, fluctuations corresponding to planar diagrams develop succesively and will only contribute for a very extended target. However, particular discontinuities of planar graphs (which vanish when summed) can be asymptotically essential, and must therefore be taken into account in the analyses of cross sections of particular processes. For instance, the elastic cross section is evidently determined by cutting a planar diagram.
Gribov managed to present the total cross section in a very simple way through dispersion integrals of the ‘particle-particle $`n`$ pomeron’ amplitudes. He did this by rearranging the contributions of both non-planar and planar multireggeon diagrams using the reggeon unitarity condition. Keeping only the one-particle pole contributions to these amplitudes (Fig. 1(a)) we reproduce the formula of the well-known ‘eikonal approximation’ for high-energy scattering. However we should keep in mind that the space-time picture behind this formula does not correspond to successive elastic rescatterings. The genuine space-time picture has been lost under rearrangement of diagrams with different planarities.
## 3 The eikonal model
### 3.1 The eikonal $`\chi _P`$
The single pomeron-exchange amplitude has the form <sup>4</sup><sup>4</sup>4 The normalization of amplitude is $`\sigma ^{\mathrm{tot}}(s)=2\mathrm{I}\mathrm{m}M(s,0)`$.
$`M_P(s,t)=\left({\displaystyle \frac{s}{s_0}}\right)^{\alpha _P(t)1}\eta _P(\alpha _P(t))g_1(t)g_2(t),`$ (1)
where $`g_{1,2}(t)`$ are the couplings of the pomeron to the colliding hadrons, and
$`\eta _P(\alpha _P(t))={\displaystyle \frac{1+\mathrm{exp}(i\pi \alpha _P(t))}{\mathrm{sin}\pi \alpha _P(t)}}`$ (2)
is the signature factor, which for $`\alpha _P(0)=1`$ is equal to $`i`$. Here we shall neglect the real part of the pomeron amplitude, assuming that the value of $`\mathrm{\Delta }=1\alpha _P(0)`$ is small.
It is convenient to analyse high-energy interactions in terms of the impact parameter, $`b`$, by introducing the Fourier transformed amplitude <sup>5</sup><sup>5</sup>5 The normalization used for $`M(s,t)`$ corresponds to $`\sigma ^{\mathrm{tot}}=2d^2b\mathrm{Im}f(s,b)`$ and $`\sigma ^{\mathrm{el}}=d^2b|f(s,b)|^2`$.
$`f(s,b)={\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle d^2ke^{i𝐛𝐤}M(s,𝐤^2)},`$ (3)
where $`t=𝐤^2`$. For an exponential parametrization of residues, $`g_i(t)=g_i\mathrm{exp}(R_i^2𝐤^2)`$, $`i=1,2`$, and a linear parametrization of the pomeron trajectory, $`\alpha _P(t)=\alpha _P(0)+\alpha _P^{}t`$, we obtain the familiar regge-pole approximation of the amplitude in impact parameter space
$`f_P(Y,b)ig_1g_2{\displaystyle \frac{\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda }}\right)}{4\pi \lambda }}\mathrm{exp}(\mathrm{\Delta }Y)=i{\displaystyle \frac{\chi _P(Y,b)}{2}},`$ (4)
$`Y\mathrm{ln}(s/s_0),s_01\text{ GeV}^2,`$
$`\lambda =R_1^2+R_2^2+\alpha _P^{}Y.`$
Note that the amplitude $`f_P(Y,b)`$ increases as $`\mathrm{exp}(\mathrm{\Delta }Y)`$, and so violates unitarity as $`s\mathrm{}`$. (Recall that $`|f(Y,b)|2`$ due to unitarity.) If we include multi-pomeron exchanges, then this inconsistency is avoided.
### 3.2 Cross section formulae in terms of the eikonal $`\chi _P`$
The eikonal approximation is the simplest way to restore $`s`$-channel unitarity for elastic amplitudes. Summation of the eikonal diagrams gives the following well known expressions in impact parameter space
$`f(Y,b)`$ $`=i\left(1\mathrm{exp}\left({\displaystyle \frac{\chi _P(Y,b)}{2}}\right)\right),`$ (5)
$`\sigma ^{\mathrm{tot}}(Y)`$ $`=2{\displaystyle d^2b\left(1\mathrm{exp}\left(\frac{\chi _P(Y,b)}{2}\right)\right)},`$ (6)
where, recall, $`Y\mathrm{ln}(s/s_0)`$. At very high energies, $`\mathrm{Im}f(Y,b)1`$ (the black disc limit) in the region of $`b`$ where $`\chi _P(Y,b)`$ is large. From (4) we see that $`\chi _P`$ becomes small only in the region $`b^2>4\mathrm{\Delta }\alpha ^{}Y^2`$. Thus, for very large $`s`$, the radius of interaction increases as $`R^2(s)=4\mathrm{\Delta }\alpha ^{}\mathrm{ln}^2(s/s_0)`$, and the total cross section increases as $`\sigma ^{\mathrm{tot}}2\pi R^2(s)`$.
To obtain the inelastic cross section we must consider, according to the AGK cutting rules, all eikonal type diagrams in which at least one pomeron is cut. Then for each cut pomeron we have a factor $`\chi _P`$, and for each uncut pomeron a factor $`(\chi _P)`$, which takes into account the position of the uncut pomerons both on the left and on the right of the cutting plane (that is $`if_Pif_P^{}=\chi _P`$). If no pomerons are cut, then it is necessary to subtract the extra terms where all pomerons formally are on the same side of the cutting plane. This rule is valid both in the momentum and in the coordinate representation .
For instance, for the two-pomeron-exchange diagram, the discontinuities for zero, one and two cut pomerons give, respectively,
$`\sigma _0^{(2)}={\displaystyle \frac{1}{2!}}\left[(\chi _P)^22(\chi _P/2)^2\right]={\displaystyle \frac{1}{4}}\chi _P^2,\sigma _1^{(2)}={\displaystyle \frac{1}{2!}}2\chi _P(\chi _P)=\chi _P^2,\sigma _2^{(2)}(b)={\displaystyle \frac{\chi _P^2}{2!}},`$ (7)
which reproduces the known AGK ratios $`1:4:2`$ . It is easy to check that these contributions sum to the second term, $`2\mathrm{Im}f^{(2)}(b)`$, in the power series expansion of the eikonal formula (5).
In this model, the distribution in terms of the number $`k`$ of cut pomerons at fixed $`b`$ has the Poissonian form
$`\sigma _k(Y,b)`$ $`={\displaystyle \frac{(\chi _P(Y,b))^k}{k!}}\mathrm{exp}(\chi _P(Y,b)),`$ (8)
$`\sigma _0(Y,b)`$ $`=1+\mathrm{exp}(\chi _P(Y,b))2\mathrm{exp}\left({\displaystyle \frac{\chi _P(Y,b)}{2}}\right),`$ (9)
which leads to the following expressions for the inelastic and diffractive cross sections
$`\sigma ^{\mathrm{inel}}(Y)`$ $`=`$ $`{\displaystyle d^2b\left(1\mathrm{exp}(\chi _P(Y,b))\right)}`$ (10)
$`\sigma _0(Y)`$ $`=`$ $`{\displaystyle d^2b\left(1\mathrm{exp}\left(\frac{\chi _P(Y,b)}{2}\right)\right)^2}.`$ (11)
In the eikonal model of Fig. 1 only elastic intermediate states appear in the rescattering diagrams <sup>6</sup><sup>6</sup>6 For the interaction with a nucleus, the AGK rules can be applied in a more general situation in which every pomeron exchange is substituted by the hadron-nucleon amplitude $`f_{hN}(Y,b)`$, i.e. by the whole set of multi-pomeron exchanges. The amplitudes $`f_{hN}(Y,b)`$ have both inelastic and elastic discontinuities. As a result, (6) will contain $`\sigma _{hN}^{\mathrm{tot}}(Y,b)/2`$ instead of $`\chi _P/2`$, (10) will contain $`\sigma _{hN}^{\mathrm{inel}}(Y,b)`$ and (11) will include both the elastic and diffractive dissociation cross sections. . Hence it is natural that $`\sigma _0=\sigma ^{\mathrm{el}}`$.
The AGK cancellation theorem enables the inclusive cross sections to be calculated. For example, consider the single particle inclusive process $`12aX`$. In this case, at least one pomeron is cut. The others may be either cut (giving a contribution $`\chi _P`$) or uncut (giving a contribution $`\chi _P`$). Thus the multiple-pomeron-exchange contributions cancel, and the single-inclusive cross section is described by the diagram shown in Fig. 2. As a function of the particle rapidity it is given by
$$\frac{d\sigma ^a}{dy}=\lambda _ag_1(0)g_2(0)e^{\mathrm{\Delta }y}e^{\mathrm{\Delta }(Yy)}=\lambda _ag_1(0)g_2(0)e^{\mathrm{\Delta }Y},$$
(12)
where $`\lambda _a`$ is related to the rapidity density of hadron $`a`$ in events originating from single-pomeron-exchange.
Note that the impact parameter $`b`$ is conserved during the eikonal rescattering. As a consequence, all the formulae are valid, not only for integrated cross sections, but at any fixed $`b`$. In particular, the same increase with energy, $`(s/s_0)^\mathrm{\Delta }`$, occurs at fixed $`b`$ for the density of particles $`(d\sigma ^a/dy)/\sigma ^{\mathrm{inel}}`$ in the limit $`s\mathrm{}`$.
### 3.3 Partonic interpretation of the eikonal model
The formulae of the eikonal approximation can be interpreted in terms of the interactions of fast partons in colliding hadrons. If the distribution in the number of fast partons in a hadron has a Poisson form, and if the cross section of a parton-parton interaction is denoted by $`\widehat{\sigma }`$, then the summation of the diagrams <sup>7</sup><sup>7</sup>7 Note that in these diagrams a parton of hadron 1 interacts with only one parton of hadron 2. It resembles the Czyź - Maximon approximation for nucleus-nucleus collisions, when only single nucleon-nucleon interactions are taken into account. shown in Fig. 3 leads to the eikonal results of (5) and (10) with
$`\chi _P(Y,b)`$ $`={\displaystyle d^2b_1d^2b_2\rho _1^{(0)}(𝐛_1)\widehat{\sigma }(Y,𝐛𝐛_1+𝐛_2)\rho _2^{(0)}(𝐛_2)},`$ (13)
where $`\rho _i^{(0)}(b_i)`$ is the fast parton distribution <sup>8</sup><sup>8</sup>8 The distributions are normalized so that $`d^2b_i\rho _i^{(0)}(b_i)=n_i`$, where $`n_i`$ is the mean number of fast partons in the hadron $`i`$. in impact parameter space of colliding hadron $`i`$, with $`i=1,2`$. That is the contribution of single-pomeron-exchange is represented by the single particle densities $`\rho _i^{(0)}`$ and the cross section $`\widehat{\sigma }(Y)`$ of the parton-parton interaction. Similarly, the $`n`$-pomeron-exchange contribution equals to the probability of finding $`n`$ partons in each of the colliding hadrons (which, for independently distributed partons, are simply $`(\rho _i^{(0)})^n`$ with $`i=1,2`$) multiplied by a sign-alternating factor $`(1)^{n+1}(\widehat{\sigma }/2)^n/n!`$. The alternating sign is due to parton screening.
Note that (13) is written assuming that the partons are all the same. In the more realistic situation, with quarks and gluons as partons, the formalism may be more complicated to allow for different parton-parton amplitudes. It results in a straightforward generalization of the simple eikonal model, see, for example Ref. .
One of the drawbacks of this simple model is the lack of energy-momentum conservation . Indeed a very large number of interactions $`\chi _P(Y)`$ become important as $`Y\mathrm{}`$, and it is necessary to allow for energy-momentum conservation in the distributions of momenta in the partonic systems. This will lead to deviations from Poisson distributions. These deviations are usually taken into account in realistic models of high-energy interactions .
However, the interpretation is not self-consistent in the case of the supercritical pomeron ($`\mathrm{\Delta }>0`$). We see that in this case the origin of the increase of the amplitude $`\widehat{f}(Y)`$ with energy is $`\mathrm{𝑛𝑜𝑡}`$ explained in terms of partons. It is desirable to reformulate this approach without reference to pomeron exchange in parton-parton interaction. As we discussed in section 2, the increase of the pomeron-exchange amplitude can be explained as the increase in the number of slow partons. It is possible to rewrite (13) in such form that it will correspond to the interaction of two partonic showers viewed from the Lorentz frame at arbitrary rapidity $`y_1`$:
$$\chi _P(Y,b)=d^2b_1d^2b_2\rho _1(y_1,𝐛_1)\widehat{\sigma }_0(𝐛𝐛_1+𝐛_2)\rho _2(Yy_1,𝐛_2),$$
(14)
where the slow parton densities $`\rho _i`$ have a Regge form
$`\rho _i(y,b)={\displaystyle \frac{g_i}{4\pi \alpha ^{}y}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\alpha ^{}y}}\right)\mathrm{exp}(\mathrm{\Delta }y),`$ (15)
and the parton-parton interaction cross section $`\widehat{\sigma }_0`$ is local in rapidity space. It is easy to see that expression (14) does not depend on a choice of Lorentz frame, i.e. on the point $`y_1`$, due to particular form of reggeon densities (15).
Since only the products of quantities occur, we have been able to move the energy dependence (that is the $`s`$ or $`Y`$ dependence) from the parton-parton reaction cross section $`\widehat{\sigma }(Y)`$ to the incoming parton distributions $`\rho _i`$ in (14). Thus all the $`Y`$ dependence now occurs in the parton densities, while the function $`\widehat{\sigma }_0`$ describes the interactions of two partons with the same rapidity. In other words, in the alternative form (14) there is no reference to pomeron exchange, and all the Regge behaviour occurs in the densities – the increase of the densities as a function of $`y`$ is natural because of the cascade development of the two partonic systems. The average number of slow partons, $`n_i(y)=d^2b_i\rho _i(y,b_i)`$, is the product of the average number of fast partons, i.e. of partonic cascades, and the average number of slow partons in the cascade, $`m(y)\mathrm{exp}(\mathrm{\Delta }y)`$. However taking a Poisson distribution for the partons at each rapidity is a strong assumption of the eikonal model.
Thus, in the framework of the parton model, we have either a Regge form of $`\widehat{\sigma }(Y)`$ in (13), or the Regge increase of the partonic densities $`\rho _i(y)`$ in (14). In (14), $`\widehat{\sigma }_0`$ describes a local parton-parton interaction as a function of both rapidity and impact parameter, which avoids the highly non-local pomeron interaction which occurs in (13). Energy-momentum conservation in the partonic interpretation can be imposed by requirement of the energy-momentum sum rule for the parton distributions $`\rho _i`$.
### 3.4 Probabilistic interpretation of the inelastic cross section
In the eikonal model there is a clear probabilistic interpretation of the inelastic cross section . The single inelastic cross section at fixed $`b`$, $`\chi _P(Y,b)`$, corresponds to the product of the average numbers of partonic cascades at fixed $`b_{1,2}`$ and of probability for soft partons to interact. Similarly, the interaction of $`k`$ soft partons from different parton chains is determined by the formula (8)
$`\sigma _k^{\mathrm{inel}}(Y,b)`$ $`={\displaystyle \frac{(\chi _P(Y,b))^k}{k!}}\mathrm{exp}(\chi _P(Y,b)),`$
where the exponential factor $`\mathrm{exp}(\chi _P)`$ corresponds to the requirement that all other partons do not interact. The total inelastic cross section at fixed $`b`$ is, therefore,
$`\sigma ^{\mathrm{inel}}(Y,b)={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\sigma _k^{\mathrm{inel}}(Y,b)`$ $`=1\mathrm{exp}(\chi _P(s,b)`$
$`={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{m1}\chi _P^m(s,b)}{m!}}.`$
We can readily see the origin of the last, sign-alternating, expression for the inelastic cross section. It is just a mathematical formula for the probability of joint (not mutually excluded) events:
$$\mathrm{Prob}(\mathrm{A}_1\mathrm{A}_2\mathrm{}\mathrm{A}_\mathrm{n})=\underset{\mathrm{i}}{}\mathrm{Prob}(\mathrm{A}_\mathrm{i})\underset{\mathrm{i}<\mathrm{j}}{}\mathrm{Prob}(\mathrm{A}_\mathrm{i}\mathrm{A}_\mathrm{j})+\underset{\mathrm{i}<\mathrm{j}<\mathrm{k}}{}\mathrm{Prob}(\mathrm{A}_\mathrm{i}\mathrm{A}_\mathrm{j}\mathrm{A}_\mathrm{k})\mathrm{}.$$
An inelastic event corresponds to the interaction of at least one slow parton with the target (an event $`\mathrm{A}_\mathrm{i}`$). Then it is necessary to subtract from the formula for $`\sigma _1^{\mathrm{inel}}(Y,b)`$ the probability of the interactions of two slow partons, and to add the one for triple interaction, and so on.
It follows from the Poisson distribution (8) that the average number, $`k`$, of cut pomerons is equal to $`\chi _P(Y,b)`$. It, thus, increases like $`\mathrm{exp}(\mathrm{\Delta }Y)`$. Thus, the increase of the produced particle densities, (12), in this model is related to the very large number of partons in the hadronic fluctuations. For the partonic interpretation (14) with partonic cascading, the exponential increase of the particle densities is clearly consistent with the inclusive particle formula (12).
Thus, to interpret the case of $`\mathrm{\Delta }>0`$, we have either to introduce the energy dependence of the fast parton interactions (which is not self-consistent), or to include a mechanism for splitting the partons in the course of the evolution in $`y`$. In the latter case the number of slow partons increases as $`\mathrm{exp}(\mathrm{\Delta }Y)`$, and if only one of slow partons interacts with the target, then the exchange of the supercritical pomeron is reproduced. The eikonal approximation will arise if several partons interact, but no more than one parton of each independent (splitting) fluctuation (see Fig.4). The more consistent approach, in which we allow any parton resulting from the fluctuation to participate in the interaction, leads to the Schwimmer model and its eikonalized version.
## 4 The Schwimmer model
Schwimmer proposed a simple model of reggeon field theory based on the triple-pomeron interaction only. It sums up the set of fan diagrams of the type shown in Fig. 5.
It is evident that the Schwimmer model is not symmetric with respect to the incoming hadrons, 1 and 2. It was originally formulated for the interaction of hadrons with nuclei. It is expected to be a reasonable approximation when hadron 1 (or a virtual photon) has a size much smaller than hadron 2 (nucleus), that is $`g_1g_2`$, and the triple-pomeron coupling $`r`$ is also small. In the case of the interaction with a nucleus we may also neglect the dependence on the impact parameter, $`b`$, at energies when the interaction radius is much smaller than the nuclear size. Here we adopt this situation as a toy model <sup>9</sup><sup>9</sup>9 The introduction of the $`b`$ dependence is straightforward, but in this case there is no analytic solution., so that the amplitudes depend only on the rapidity $`Y`$.
### 4.1 Total cross section in the Schwimmer model
Following the Schwimmer model, we choose the $`n`$-pomeron-particle 2 amplitude to have an eikonal form, with $`G_n=g_2^n`$. Rather than using the amplitude of (5), it is convenient to work in terms of the ‘truncated’ amplitude $`\phi _{\mathrm{tot}}(Y)=\sigma ^{\mathrm{tot}}(Y)/(g_1g_2)=2\mathrm{I}\mathrm{m}f(Y)/(g_1g_2)`$, and to introduce a new pomeron ‘propagator’ $`P(Y)=\chi _P/(2g_1g_2)=\mathrm{exp}(\mathrm{\Delta }Y)`$. By construction, the function $`\phi _{\mathrm{tot}}(Y)`$ satisfies the following non-linear integral equation<sup>10</sup><sup>10</sup>10Note that (16) is written for $`\phi _{\mathrm{tot}}(Y)/2`$, since the amplitude $`f(Y)=ig_1g_2\phi (Y)_{\mathrm{tot}}/2`$.
$`\phi _{\mathrm{tot}}(Y)/2=e^{\mathrm{\Delta }Y}rg_2{\displaystyle _0^Y}𝑑y_1e^{\mathrm{\Delta }(Yy_1)}(\phi _{\mathrm{tot}}(y_1)/2)^2,`$ (16)
see Fig. 6. The differential form of the equation is
$`{\displaystyle \frac{d\phi _{\mathrm{tot}}(Y)}{dY}}=\mathrm{\Delta }\phi _{\mathrm{tot}}{\displaystyle \frac{rg_2}{2}}\phi _{\mathrm{tot}}^2.`$ (17)
To solve the equation it is convenient to make the substitution
$`\phi _{\mathrm{tot}}(Y)=2\tau u_{\mathrm{tot}}(\tau ),\tau =e^{\mathrm{\Delta }Y},`$ (18)
so that (17) becomes
$`{\displaystyle \frac{du_{\mathrm{tot}}(\tau )}{d\tau }}=ϵu_{\mathrm{tot}}^2,u_{\mathrm{tot}}(1)=1,\text{with }ϵ={\displaystyle \frac{rg_2}{\mathrm{\Delta }}}.`$ (19)
The solution
$`u_{\mathrm{tot}}={\displaystyle \frac{1}{1+ϵ(\tau 1)}}`$ (20)
gives the well known expression for $`\phi _{\mathrm{tot}}(Y)`$
$`\phi _{\mathrm{tot}}(Y)={\displaystyle \frac{2P(Y)}{1+ϵ[P(Y)P(0)]}}.`$ (21)
Note that the integration in (16) goes from $`y=0`$ to $`Y`$. If, however, the integration starts from $`y_{\mathrm{min}}`$, then for the corresponding solution $`\phi _{\mathrm{tot}}(Y;y_{\mathrm{min}})`$ we should replace $`P(0)`$ in (21) by $`P(y_{\mathrm{min}})`$. From (21), we see that the cross section $`g_1g_2\phi _{\mathrm{tot}}(Y)`$ at first increases as $`\mathrm{exp}(\mathrm{\Delta }Y)`$, and then tends to the finite limit $`2g_1\mathrm{\Delta }/r`$ for very large $`Y`$.
### 4.2 Inelastic and diffractive cross sections
To obtain the inelastic and diffractive amplitudes we use the AGK cutting rules just as we did in section 3. We substitute for the cut amplitude the corresponding cross section $`\sigma ^{\mathrm{inel}}`$ or $`\sigma ^\mathrm{D}`$, and for the uncut amplitude the factor $`(\sigma ^{\mathrm{tot}})`$. This results in integral equations similar to (16), but with a non-diagonal structure for the inelastic cross section:
$$\begin{array}{c}\phi _{\mathrm{inel}}(Y)\frac{\sigma ^{\mathrm{inel}}}{g_1g_2}=2e^{\mathrm{\Delta }Y}2rg_2_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{inel}}(y_1)\phi _{\mathrm{tot}}(y_1)\hfill \\ \hfill +rg_2_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{inel}}^2(y_1)+2rg_2_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{inel}}(y_1)\phi _\mathrm{D}(y_1),\end{array}$$
(22)
while for the diffractive cross section, corresponding to the production of a state accompanied by a rapidity gap,
$$\begin{array}{c}\phi _\mathrm{D}(Y)\frac{\sigma ^\mathrm{D}}{g_1g_2}=\frac{rg_2}{2}_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{tot}}^2(y_1)\hfill \\ \hfill 2rg_2_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _\mathrm{D}(y_1)\phi _{\mathrm{tot}}(y_1)+rg_2_0^Y𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _\mathrm{D}^2(y_1).\end{array}$$
(23)
Note that coefficients in (23), which result from the different cuttings, are in the same ratios, ($`1:4:2`$), as in (7). Similar equations have been obtained in Ref. in the framework of the Balitsky-Kovchegov equation.
Taking into account that $`\phi _{\mathrm{inel}}+\phi _\mathrm{D}=\phi _{\mathrm{tot}}`$, we obtain from (22) the differential equation for inelastic cross section
$`{\displaystyle \frac{du_{\mathrm{inel}}(\tau )}{d\tau }}`$ $`=2ϵu_{\mathrm{inel}}^2,`$ $`u_{\mathrm{inel}}(1)=1,`$ (24)
where, similar to (18), we use the substitutions
$`\phi _{\mathrm{inel}}(Y)=2\tau u_{\mathrm{inel}}(\tau ),\phi _\mathrm{D}(Y)=2\tau u_\mathrm{D}(\tau ).`$
It is remarkable that the equation for $`u_{\mathrm{inel}}`$, i.e. for $`\sigma ^{\mathrm{inel}}`$, is diagonal. It is a generalization of the similar result for the non-enhanced diagrams (see the footnote in section 6).
Thus, in analogy to (21), we obtain the solutions
$`u_{\mathrm{inel}}(\tau )`$ $`={\displaystyle \frac{1}{1+2ϵ(\tau 1)}},`$ (25)
$`u_\mathrm{D}(\tau )`$ $`u_{\mathrm{tot}}(\tau )u_{\mathrm{inel}}(\tau )={\displaystyle \frac{1}{1+ϵ(\tau 1)}}{\displaystyle \frac{1}{1+2ϵ(\tau 1)}}.`$ (26)
Note that in the limit $`ϵ\tau 1`$, these solutions reproduce the first reggeon graphs, and that in the saturation regime (where $`ϵ\tau 1`$) we have $`\phi _{\mathrm{inel}}=\phi _\mathrm{D}=\phi _{\mathrm{tot}}/2=1/ϵ`$.
Next, we obtain the dependence of diffractive production on the rapidity gap $`y`$, or on the mass of the produced system, where ln$`(M^2/M_0^2)=Yy`$. We introduce a function $`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})`$ corresponding to the cross section for production of the final state with a rapidity gap larger than $`y_{\mathrm{min}}`$:
$`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})={\displaystyle \frac{1}{g_1g_2}}{\displaystyle _{y_{\mathrm{min}}}^Y}𝑑y_1{\displaystyle \frac{d\sigma ^\mathrm{D}}{dy_1}},\phi _{\mathrm{gap}}(Y;0)=\phi _\mathrm{D}.`$ (27)
This cross section satisfies the same integral equations as the diffractive dissociation cross section $`\phi _\mathrm{D}`$ except that the integration over rapidity starts from $`y_{\mathrm{min}}`$ instead of 0. That is
$`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})`$ $`={\displaystyle \frac{rg_2}{2}}{\displaystyle _{y_{\mathrm{min}}}^Y}e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{tot}}^2(y_1)`$
$`2rg_2{\displaystyle _{y_{\mathrm{min}}}^Y}𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{gap}}(y_1;y_{\mathrm{min}})\phi _{\mathrm{tot}}(y_1)`$
$`+rg_2{\displaystyle _{y_{\mathrm{min}}}^Y}𝑑y_1e^{\mathrm{\Delta }(Yy_1)}\phi _{\mathrm{gap}}^2(y_1;y_{\mathrm{min}}),`$ (28)
see Fig. 7.
As before, we may write this in the differential form
$`{\displaystyle \frac{du_{\mathrm{gap}}(\tau ;\tau _{\mathrm{min}})}{d\tau }}=2ϵ({\displaystyle \frac{1}{2}}u_{\mathrm{tot}}^22u_{gap}u_{\mathrm{tot}}+u_{\mathrm{gap}}^2),u_{\mathrm{gap}}(\tau _{\mathrm{min}};\tau _{\mathrm{min}})=0,`$ (29)
where $`u_{\mathrm{tot}}`$ is the solution of (19), and
$`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})=2\tau u_{\mathrm{gap}}(\tau ;\tau _{\mathrm{min}}),\tau _{\mathrm{min}}=e^{\mathrm{\Delta }y_{\mathrm{min}}}.`$ (30)
In analogy to (26), the solution is
$`u_{\mathrm{gap}}(\tau ;\tau _{\mathrm{min}})`$ $`={\displaystyle \frac{1}{1+ϵ(\tau 1)}}{\displaystyle \frac{1}{1+ϵ(2\tau \tau _{\mathrm{min}}1)}},`$ (31)
or,
$`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})`$ $`={\displaystyle \frac{2e^{\mathrm{\Delta }Y}}{1+ϵ(e^{\mathrm{\Delta }Y}1)}}{\displaystyle \frac{2e^{\mathrm{\Delta }Y}}{1+ϵ(2e^{\mathrm{\Delta }Y}e^{\mathrm{\Delta }y_{\mathrm{min}}}1)}}.`$ (32)
Thus, we can calculate the cross section for a fixed gap $`y`$, that is for the diffractive production of a state of given mass $`M`$ (with the value of $`y_M=Yy`$ fixed). It is determined by the derivative of the second term of (32):
$`{\displaystyle \frac{d\sigma ^\mathrm{D}}{dy_M}}M^2{\displaystyle \frac{d\sigma ^\mathrm{D}}{dM^2}}`$ $`=g_1g_2{\displaystyle \frac{d\phi _{\mathrm{gap}}(Y;y)}{dy}}={\displaystyle \frac{2g_1g_2\mathrm{\Delta }ϵe^{\mathrm{\Delta }(2Yy_M)}}{[1+ϵ(2e^{\mathrm{\Delta }Y}e^{\mathrm{\Delta }(Yy_M)}1)]^2}}`$ (33)
$`{\displaystyle \frac{g_1\mathrm{\Delta }^2}{r}}{\displaystyle \frac{2\mathrm{exp}(\mathrm{\Delta }y_M)}{[2\mathrm{exp}(\mathrm{\Delta }y_M)1]^2}}\text{( for }ϵ\mathrm{exp}(\mathrm{\Delta }y_{\mathrm{min}})1\text{ )}.`$ (34)
This cross section in the Schwimmer model was first obtained in Ref. .
We see that the cross section (34) decreases with $`M^2`$, which provides convergence of the integral over the mass of the diffractively produced system. Indeed, in the region of large $`M^2`$, that is in the saturation domain with $`y_M1`$, we have
$$M^2\frac{d\sigma ^\mathrm{D}}{dM^2}(M^2)^\mathrm{\Delta }$$
(35)
Thus the $`M^2`$ distribution gives information on the intercept of the bare pomeron, $`\alpha _P(0)1+\mathrm{\Delta }`$. Although (35) was derived in the Schwimmer model, we shall see that the same behaviour is valid for its eikonal generalization.
Another way to get information on the bare intercept is to study the inclusive spectrum. Using the AGK cutting rules, we find that the particle rapidity distribution is
$`{\displaystyle \frac{d\sigma ^a}{dy}}=\lambda _ag_1g_2e^{\mathrm{\Delta }y}\phi _{\mathrm{tot}}(y_2),\mathrm{with}y_2=Yy.`$ (36)
In a frame where hadron 1 is moving fast, (36) can be interpreted as a Regge-like increase of partons. However, the partonic interpretation of this result is different in a frame where particle 2 is fast; see section 4.4.
### 4.3 The eikonalized Schwimmer model
Suppose, now, that there are several partons in the initial state at $`y=0`$ which split in the course of the evolution. In the absence of splitting this would correspond to the usual eikonal model (see sect. 3). However as a result of splitting, the evolution of each initial parton corresponds to the Schwimmer amplitude – and the whole amplitude is described by Fig. 8. The AGK rules for this set of diagrams are similar to the ones for the eikonal graphs of Fig.1 except that each Schwimmer-type amplitude contains, not only the inelastic discontinuity $`\phi _{\mathrm{inel}}`$ due to pomeron exchange, but also the discontinuity corresponding to gap production $`\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}})`$, with relations $`\phi _{\mathrm{inel}}+\phi _\mathrm{D}=\phi `$ and $`\phi _\mathrm{D}=\phi _{\mathrm{gap}}(Y;0)`$. Then, the set of formulae for the various cross sections will be similar to the (6), (10), (11), together with the ones resulting from the extra discontinuities of the amplitude
$`\sigma ^{\mathrm{tot}}(Y;b)`$ $`=2\left[1\mathrm{exp}(g_1g_2\phi _{\mathrm{tot}}(Y)/2)\right],`$ (37)
$`\sigma ^{\mathrm{el}}(Y;b)`$ $`=\left(1\mathrm{exp}(g_1g_2\phi _{\mathrm{tot}}(Y)/2)\right)^2,`$ (38)
$`\sigma ^{\mathrm{inel}}(Y;b)`$ $`=1\mathrm{exp}(g_1g_2\phi _{\mathrm{inel}}(Y)),`$ (39)
$`\sigma ^\mathrm{D}(Y;b)`$ $`=\mathrm{exp}(g_1g_2\phi _{\mathrm{inel}}(Y))\mathrm{exp}(g_1g_2\phi _{\mathrm{tot}}(Y)),`$ (40)
$`\sigma ^{\mathrm{gap}}(Y;y_{\mathrm{min}};b)`$ $`=\mathrm{exp}(g_1g_2\phi _{\mathrm{tot}}(Y))[(\mathrm{exp}(g_1g_2\phi _{\mathrm{gap}}(Y;y_{\mathrm{min}}))1],`$ (41)
where $`\phi _{\mathrm{tot}}`$, $`\phi _{\mathrm{inel}}`$, $`\phi _\mathrm{D}`$ and $`\phi _{\mathrm{gap}}`$ have been defined above. We see that the following relations hold
$`\sigma ^{\mathrm{tot}}(b)\sigma ^{\mathrm{el}}(b)=\sigma ^{\mathrm{inel}}(b)+\sigma ^\mathrm{D}(b)=1e^{g_1g_2\phi _{\mathrm{tot}}}.`$ (42)
Note that again we have a closed expression for $`\sigma ^{\mathrm{inel}}`$.
The differential cross section for the diffractive production of a state of mass $`M`$ is obtained by differentiation of (41) with respect to $`y_{\mathrm{min}}`$, which enters via $`\phi _{\mathrm{gap}}`$. That is via
$`{\displaystyle \frac{d\sigma ^{\mathrm{gap}}}{dy_M}}=\mathrm{exp}\left[g_1g_2(\phi _{\mathrm{tot}}(Y;b)\phi _{\mathrm{gap}}(Y,Yy_M;b))\right]{\displaystyle \frac{d\sigma _{\mathrm{Sch}}^\mathrm{D}}{dy_M}},`$ (43)
where $`d\sigma _{\mathrm{Sch}}^\mathrm{D}/dy_M=g_1g_2d\phi _{\mathrm{gap}}(Y;y)/dy`$ is defined by (33). In the saturation limit (34) we obtain
$`{\displaystyle \frac{d\sigma ^\mathrm{D}}{dy_M}}\mathrm{exp}\left({\displaystyle \frac{g_1\mathrm{\Delta }}{r}}\right){\displaystyle \frac{g_1\mathrm{\Delta }^2}{r}}{\displaystyle \frac{2\mathrm{exp}(\mathrm{\Delta }y_M)}{[2\mathrm{exp}(\mathrm{\Delta }y_M)1]^2}},`$ (44)
for $`1y_MY`$. Thus, again, the dependence shown in (35) is valid at large values of $`M`$.
We note that equations (40), (41), (43) differ from the results of , where absorptive effects were included by multiplication by the factor $`\mathrm{exp}(\phi _{\mathrm{tot}})`$. This procedure, however, does not allow for the simultaneous diffractive cuttings of several Schwimmer amplitudes. This difference is especially important in calculations of the survival probability which take into account absorptive effects in inelastic diffractive processes. For example, for the inclusive production of particles in large-mass diffraction, the survival probability has the form
$`S^2(Y,y_M;b)=\mathrm{exp}\left[g_1g_2(\phi _{\mathrm{tot}}(Y;b)\phi _{\mathrm{gap}}(Y,Yy_M;b))\right].`$ (45)
This result can be easily obtained by using the method of Ref. . We emphasize that, in contrast to the eikonal model, the survival probability depends not only on $`Y`$, but also on the mass of the produced system $`y_M`$.
### 4.4 Partonic interpretation of Schwimmer diagrams
As we discussed in 3.3, the supercritical pomeron requires a mechanism for parton splitting. In the Schwimmer model this occurs through a single parton cascading in terms of reggeon diagrams. On the other hand, in the eikonalized Schwimmer model it is described by the independent cascading of a Poisson set of initial fast partons. In both models the inclusive spectrum is described by similar formula (36). The increase of the spectrum with the rapidity of the inclusive particle is due to the partonic cascade, which leads to an exponential growth of partons with $`y`$. Note that there is no fusion of partons in this cascade, which would have inhibited its growth.
We stress that the model is not symmetric with respect to the colliding hadrons. In the frame where hadron 2 moves fast, the parton interpretation requires both splitting and fusion of partons like the first and the second terms in the r.h.s. of Eq.(17)). As a result, we first have a growth of the number of partons and then saturation to a constant value, due to recombination. This is the usual interpretation used in discussions of the saturation of parton densities in QCD cascade . This behaviour of parton density, in the case where hadron 2 moves fast, can be traced to the $`y_2`$ behaviour of the inclusive spectrum (36). Note, however, that only parton fusion producing tree reggeon diagrams is allowed in this approximation. This is justified for $`r1,g_21`$, as was discussed in section 4.
Note that this dependence of the partonic interpretation on the choice of the Lorentz frame is due to the special (non-symmetric) selection of reggeon diagrams related to particular process. This is reasonable in a limited region of rapidity, with $`r\mathrm{exp}(\mathrm{\Delta }Y)1`$. For higher rapidities, loop diagrams become important. Of course, if the complete set of diagrams of reggeon theory were to be used, then the parton dynamics would be identical in all Lorentz frames .
The multiplicity distribution in the Schwimmer model is not Poisson-like . There are huge fluctuations, leading, at high energies, to a large dispersion. Hence, according to the Good-Walker formalism , there is a large probability of diffractive dissociation.
## 5 Conclusions
We have investigated the two simplest models of reggeon theory, using the AGK cutting rules for the supercritical pomeron. We discussed the partonic interpretation of the models. A closed set of equations is obtained for $`\sigma ^{\mathrm{tot}}`$, $`\sigma ^{\mathrm{inel}}`$ and $`\sigma ^\mathrm{D}`$ in the Schwimmer model. It is important that the equation for $`\sigma ^{\mathrm{inel}}`$ is diagonal, as is the equation for $`\sigma ^{\mathrm{tot}}`$. Explicit formulae for rapidity gap production are obtained.
We note that, from the partonic viewpoint, both the eikonal and Schwimmer models are incomplete at asymptotic energies. In particular, the partonic interpretation of the Schwimmer model depends on the choice of the Lorentz frame. We note that, at asymptotic energies, partonic dynamics must be Lorentz invariant. From the viewpoint of reggeon field theory, this corresponds to the crucial role of the pomeron loops.
The extension of the multi-pomeron formalism carried out in this paper can lead to a better understanding of high energy dynamics and to an improvement of the analysis of data for soft high energy interactions. This is important, for example, in the calculation of probabilities of rapidity gaps in diffractive processes; see, for example, Ref. .
## Acknowledgements
We thank E. Levin for drawing our attention to Refs. , and for useful discussions. ABK and MGR would like to thank the IPPP at the University of Durham for hospitality, and ADM thanks the Leverhulme Trust for an Emeritus Fellowship. This work was supported by the Royal Society, the UK Particle Physics and Astronomy Research Council, by grants CRDF RUP2-2621-MO-04, RFBR 04-02-16073, 04-02-17263 and 03-02-04004, SS-1124.2003.2 and SS-1774.2003.2.
|
warning/0506/cond-mat0506510.html
|
ar5iv
|
text
|
# Asymmetric exchange between electron spins in coupled semiconductor quantum dots
## Abstract
We obtain a microscopic description of the interaction between electron spins in bulk semiconductors and in pairs of semiconductor quantum dots. Treating the $`𝐤`$$``$$`\widehat{𝐩}`$ band mixing and the Coulomb interaction on the same footing, we obtain in the third order an asymmetric contribution to the exchange interaction arising from the coupling between the spin of one electron and the relative orbital motion of the other. This contribution does not depend on the inversion asymmetry of the crystal and does not conserve the total spin. We find that this contribution is $``$$`10^3`$ of the isotropic exchange, and is of interest in quantum information. Detailed evaluations of the asymmetric exchange are given for several quantum dot systems.
Spin in quantum dots (QDs) are attractive candidates for qubits in quantum information in part because of their long coherence times Gupta99 . Controllable coupling between these spins is an essential requirement for two qubit quantum gates, and it has been the focus of much recent research Loss98 -Wu02 . Typically the spin-spin coupling is dominated by the isotropic exchange interaction $`J𝐒_1`$$``$$`𝐒_2`$, which arises from the Coulomb interaction and the Pauli principle. The isotropy of this part of the exchange implies that the total spin is conserved. This is important in gate operations, which involve pulsing $`J(t)`$ Stepanenko .
In general, spin-orbit coupling in solids also gives rise to terms that are asymmetric between the spins, and which do not conserve total spin. These terms can cause loss of fidelity in gate operations involving two spins. To date, the asymmetric terms that have been discussed involve spin-orbit coupling of individual electrons in the conduction band Kavokin01 ; Gorkov2003 . For systems with bulk inversion asymmetry, this coupling has the Dresselhaus form Dresselhaus ; Altshuler . In addition, heterostructure asymmetry can introduce the so-called Rashba coupling Rashba84 . These two couplings can give an additional interaction between the spins of the Dzyaloshinskii-Moryia (DM) form $`𝜷_{so}`$$``$$`(𝐒_1`$$`\times `$$`𝐒_2)`$$`/\sqrt{\beta _{so}^2+J^2}`$ DM1 ; DM2 where $`𝜷_{so}`$ is linear in the spin-orbit coupling strength Kavokin01 ; Gorkov2003 . In general the dephasing caused by the DM contribution cannot be totally eliminated, but suggestions have been made for gate pulse shapings that eliminate it to first order Burkard02 .
Here we give a different contribution to the spin-spin coupling of two electrons in III-V semiconductors that does requires neither bulk inversion asymmetry or Rashba coupling. It arises from the Coulomb interaction between two electrons and from the conduction-valence band mixing. It is similar to the exchange interaction between excitons BirPikus . We describe this asymmetric electron spin coupling in bulk materials, and we make detailed evaluations for spins in several coupled QD systems.
We use a $`8`$$`\times `$$`8`$ Kane Hamiltonian to represent the band structure of III-V semiconductors BirPikus , with a gap between an $`s`$-like conduction band and $`p`$-like valence bands. The band coupling in this Hamiltonian is obtained in $`𝐤`$$``$$`\widehat{𝐩}`$ effective mass perturbation theory, where $`\widehat{𝐩}`$ is the relevant momentum operator. The standard parameters in this approach are the band gap $`E_g`$, the energy of the split-off band $`\mathrm{\Delta }`$, and the valence-conduction band coupling $`P`$=$`(\mathrm{}/m_0)S|\widehat{p}_x|X`$, where $`|S`$ and $`|X`$ are the Bloch states of the conduction and valence bands KaneParam .
We consider two electrons of relative coordinate $`𝐫`$$`=`$$`𝐫_1`$$``$$`𝐫_2`$ in a semiconductor of dielectric permittivity $`\kappa `$. The unperturbed two-particle Hamiltonian in the conduction band is $`H^{(0)}`$$`=`$$`(\widehat{𝐩}_1^2`$$`+`$$`\widehat{𝐩}_2^2)/2m_0`$. The two band mixing terms $`𝐤_1`$$``$$`\widehat{𝐩}_1`$, $`𝐤_2`$$``$$`\widehat{𝐩}_2`$, and the Coulomb interaction $`U_C`$$`=`$$`e^2/\kappa 𝐫`$ between the electrons are treated on equal footing. The leading part of the spin-dependent interaction arises in third order from first order contributions of $`U_C`$ and of each of the band mixing terms $`𝐤_1`$$``$$`\widehat{𝐩}_1`$, $`𝐤_2`$$``$$`\widehat{𝐩}_2`$. The spin-dependent part of the Coulomb interaction arises from the coupling between the electron spins and their relative motion:
$`H_s^{(3)}`$ $`=`$ $`{\displaystyle \frac{e^2}{\kappa }}{\displaystyle \frac{2P^2}{3E_g^2}}{\displaystyle \frac{\mathrm{\Delta }(2E_g+\mathrm{\Delta })}{(E_g+\mathrm{\Delta })^2}}`$
$`\times \left[(𝐫\times 𝐩_1)𝐒_1(𝐫\times 𝐩_2)𝐒_2\right]/\mathrm{}^2r^3.`$
The coefficient of the square bracket in Eq.(Asymmetric exchange between electron spins in coupled semiconductor quantum dots) is the coupling constant, which we will call $`C`$. For GaAs $`C`$=$`5.7`$ meVnm<sup>3</sup>, and for InAs $`C`$=$`10`$ meVnm<sup>3</sup>. The remaining spin-independent terms in this order of the theory contribute to the isotropic exchange $`J`$. These include a ’local contact’ term of the form
$$H_c^{(3)}=\frac{4\pi e^2P^2}{3\kappa }\frac{(E_g+\mathrm{\Delta })^2+E_g^2}{E_g^2(E_g+\mathrm{\Delta })^2}\delta (𝐫).$$
(2)
Coupling terms similar to $`H_s^{(3)}`$ and $`H_c^{(3)}`$ are known for two electrons in the free space QED where they arise from electron-positron band mixing and Coulomb interaction. In the present case the effect is stronger because the energy gap $`E_g`$ between electrons and holes in a crystal is much smaller than the energy gap $`m_0`$$`c^2`$ between electrons and positrons. Also, the symmetry of interaction $`H_s^{(3)}`$ \[Eq.(Asymmetric exchange between electron spins in coupled semiconductor quantum dots)\] is different from that between electrons in free space QED , because the valence bands $`\mathrm{\Gamma }_8`$ and $`\mathrm{\Gamma }_7`$ have a symmetry different from the $`s`$-symmetry of positrons.
We note here that for itinerant electrons the interaction $`H_s^{(3)}`$ \[Eq.(Asymmetric exchange between electron spins in coupled semiconductor quantum dots)\] flips a spin in an electron-electron collision, providing a new mechanism for the relaxation of spin polarization in addition to other known mechanisms OpticalPumping . We evaluated the spin dephasing time for a 2D electron gas in GaAs quantum wells and found that it is of the order $``$$`1`$ ns for He temperature and of the order $``$$`1`$ ps for room temperatures. This suggests that the interaction $`H_s^{(3)}`$ can be important in spin transport in low-dimensional structures, such as quantum wells Awsch .
Here we will primarily address systems with spins on two centers, such as two QDs or two charged donors, where $`H_s^{(3)}`$ gives rise to a DM coupling. The two electrons can be in a singlet state with total spin $`S`$$`=`$$`0`$ or in a triplet state with total spin $`S`$$`=`$$`1`$. $`H_s^{(3)}`$ has non-zero matrix elements between states of different total spin. We take the lowest singlet and triplet states to be separated by an exchange energy $`J`$. Then in the Hilbert space of these two states the Hamiltonian can be written in the form
$$\stackrel{~}{H}=J𝐒_1𝐒_2+i𝜷(𝐒_1𝐒_2),$$
(3)
where $`i𝜷`$=$`C/\mathrm{}^2S_0|(𝐫/r^3)`$$`\times `$$`𝐩_1|T_0`$ is the matrix element between the lowest orbital singlet $`|S_0`$ and the triplet $`|T_0`$. A block diagonal form equivalent to Eq.(3) can be obtained by an orthogonal transformation in spin space, giving the so called ’twisted spin representation’ Levitov2003 :
$`\stackrel{~}{H}`$ $`=`$ $`J\mathrm{cos}\varphi 𝐒_1𝐒_2+2J\left(\mathrm{sin}{\displaystyle \frac{\varphi }{2}}\right)^2(\widehat{𝐧}𝐒_1)(\widehat{𝐧}𝐒_2)+`$ (4)
$`+`$ $`J\mathrm{sin}\varphi \widehat{𝐧}(𝐒_1\times 𝐒_2),`$
where the spin-0 and spin-1 states are mixed by the operator $`\mathrm{exp}\left[i(\varphi /2)\widehat{𝐧}(𝐒_1𝐒_2)\right]`$, with $`𝐧`$=$`𝜷/\beta `$ and $`\varphi `$$`=`$$`\mathrm{arctan}(\beta /J)`$. In Eq.(4) the first two terms are the symmetric isotropic and the symmetric anisotropic Heisenberg terms. The last term is the asymmetric exchange in the DM form $`𝜷`$$``$$`(𝐒_1`$$`\times `$$`𝐒_2)`$$`/\sqrt{\beta ^2+J^2}`$. The contributions from the Dresselhaus and Rashba couplings to the anisotropic and asymmetric exchange were given in the form of Eq.(4) in Ref.Kavokin01 .
We consider the lowest single-particle states $`|\phi _\pm `$ from the two confining centers. In the absence of spin-orbit coupling the two-particle wavefunctions can be written as products of orbital states and spin states. The axial vector $`𝜷`$ is nonzero only for inversion-asymmetric confining potentials. For such systems we take into account the possibility that both electrons occupy the same site, therefore we use the Hund-Mulliken description of the two-particle orbital states Wu02 . We obtain the ground state by diagonalizing analytically the tunnelling and the Coulomb and local contact interactions. We then focus on the Hilbert space determined by the lowest singlet ($`S`$$`=`$$`0`$) and triplet states ($`S`$$`=`$$`1`$) in which the Hamiltonian is given by Eq.(3). The importance of the resulting DM asymmetric exchange depends on the ratio $`\mathrm{tan}\varphi `$=$`\beta /J`$ between the coefficients of the third and the first terms in Eq.(4).
Consider first the ratio $`\beta /J`$ for electrons confined by two vertically coupled QDs, such as to those in InAs/GaAs systems SK . These dots generally have different sizes and can have shape asymmetries as well AsymDots . The band offset between GaAs and InAs ($`U_0`$$``$$`0.7`$ eV) is larger than the quantization energy in the lateral direction ($`\mathrm{}\omega `$$``$$`20`$ meV) by more than an order of magnitude. This allows us to decouple the vertical and lateral degrees of freedom. We describe these dots by the potential offsets along the growth axis $`z`$ and parabolic potentials in lateral directions $`x`$ and $`y`$, which in general can be anisotropic. These lateral directions are independent of the crystal axis. The potentials and wavefunctions then can be written straightforwardly VertDots ; GaussianModel . Thus, the asymmetry studied here results from differences in the lateral sizes and shapes of the QDs.
The dependence of the ratio $`\beta /J`$ of the asymmetric exchange to the symmetric exchange on the separation $`d`$ between the two cylindrical dots with a fixed lateral offset $`b`$ is shown in Fig.1(a). Its dependence on the lateral offset $`b`$ for two dots at a fixed $`d`$ is shown in Fig.1(b). The insets of Fig.1(a,b) show sketches of the lateral and vertical views of two cylindrical dots with vectors $`𝒃`$, $`𝒅`$ pointing from the smaller dot to the larger one. The vector $`𝜷`$ is oriented in the direction $`𝒃`$$`\times `$$`𝒅`$, and for small differences $`\mathrm{\Delta }a`$=$`a_+`$$``$$`a_{}`$ between the dot radii $`a_\pm `$, it is proportional to $`\mathrm{\Delta }a`$. It follows from Eq.(Asymmetric exchange between electron spins in coupled semiconductor quantum dots) that for two cylindrical dots of equal sizes $`\beta /J`$ vanishes because the system then has a center of inversion and the ground state is symmetric is no asymmetric exchange because the system has cylindrical symmetry ($`𝒃`$$`\times `$$`𝒅`$=0). In each of these figures there are two regions of behavior: (i) For small size differences $`\mathrm{\Delta }a`$ (the left sides of the peaks) $`\beta /J`$ increases roughly proportional to $`d`$ in Fig.1(a) and roughly proportional to $`b`$ in Fig.1(b) until it reaches $``$$`3.5`$$`\times `$$`10^3`$; in this region the two electrons are distributed almost symmetrically between the dots. (ii) For larger values of $`\mathrm{\Delta }a`$ (right sides of the peaks), $`\beta /J`$ first increases with $`d`$ to a small maximum and then decreases \[Fig.1(a)\], and similarly with $`b`$ Fig.1(b)\]. In this second region both electrons tend to occupy preferentially the larger dot, where the Coulomb energy is overcome by the difference between single-particle energies of the two dots.
In order to study shape asymmetry, we consider two vertically coupled dots with deviations from cylindrical symmetry. In Fig.2 we give results for two identical elliptical dots that are rotated by $`\pi /2`$ with respect to one another, as shown in the inset, with an offset $`b`$=$`4`$ nm. The dots of equal sizes have equal energies, which leads to an equal distribution of two electrons on them. The ratio $`\beta /J`$ reaches substantial values ($``$$`10^3`$) as a function of the separation $`d`$ and has a relatively weak dependence on the angle $`\gamma `$ between the axis of the dots. In this case $`\beta `$ arises from the shape asymmetry. The coupling depends on the angle $`\gamma `$ between the relative position vector and the principal axis (for cylindrical dots, $`𝜷`$ changes direction but is constant in magnitude). The orientation of $`𝜷`$ is given by $`𝒃`$$`\times `$$`𝒅`$. In the case where the axies of the dots are parallel ($`a_{,x}`$=$`a_{+,x}`$, $`a_{,y}`$=$`a_{+,y}`$), the asymmetric exchange is zero (the system then has an inversion center at $`(𝒃/2,𝒅/2)`$). In cases when the dots are different in size and not cylindrical, $`\beta /J`$ is even larger. For example, for a cylindrical dot with $`a_{}`$=$`5`$ nm coupled with an elliptical dot with $`a_{+,x}`$=$`4`$ nm, $`a_{+,y}`$=$`6.25`$ nm, the maxima of $`\beta /J`$ are $``$$`3.5`$$`\times `$$`10^3`$ for all $`\gamma `$.
We have calculated the contribution to the asymmetric exchange for these structures from the bulk Dresselhaus coupling $`H_{so}^D`$$`=`$$`i\gamma _{so}^D_x(_y^2_z^2)`$$`S_\alpha /\mathrm{}`$ (plus cyclic permutations of cartesian indices). The coupling constant $`\gamma _{so}^D`$ is $`47`$ meVnm<sup>3</sup> for GaAs, and of the order of $`100`$ meVnm<sup>3</sup> for InAs. In Fig.1(c) we compare the contribution $`\beta _D/J`$ to the asymmetric exchange from the Dresselhaus coupling with asymmetric exchange $`\beta /J`$ for a size difference $`\mathrm{\Delta }a`$=$`0.22`$ nm Compare\_AE\_D ; linearDresselhaus ; Rashbaparam\_InAs . We see that $`\beta /J`$ is larger for intermediate separations $`d`$, a region of particular interest for implementations for quantum information. As the difference in the dot sizes increases, $`\beta /J`$ becomes larger relative to $`\beta _D/J`$, and as the size difference decreases, $`\beta _D/J`$ becomes larger.
We have also considered laterally coupled QDs. To obtain a representation of the barrier between the QDs we use inverted-gaussian potentials ReviewQD02 , and we use again the material parameters for InAs. The wavefunctions are obtained variationally LatDots . In Fig.(3) we give results for two elliptical dots of equal sizes rotated with respect to each other by $`\pi /2`$. In this case the anisotropic exchange $`\beta /J`$ arises from the shape asymmetry. From the operator $`𝐫`$$`\times `$$`_1`$ in Eq.Asymmetric exchange between electron spins in coupled semiconductor quantum dots, the asymmetric exchange $`\beta `$ has a nonzero component only along the growth axis, and the dependence of the modulus $`\beta /J`$ is symmetric with respect to $`\gamma `$$`=`$$`\pi /4`$. In Fig.(3) once again $`\beta /J`$ reaches a maximum of $``$$`10^3`$. In cases when the Dresselhaus linearDresselhaus and Rashba Rashbaparam\_InAs couplings have equal coupling constants, their contribution is small for $`\gamma `$=$`\pi /4`$ and then the total asymmetric exchange is dominated by $`\beta /J`$.
In summary, we derived an asymmetric contribution to the exchange interaction between two electrons in III-V semiconductors that arises from the Coulomb interaction and the band mixing and does not require inversion asymmetry. For asymmetric coupled semiconductor QDs, this contribution depends on the geometry and is typically $`10^3`$ of the isotropic exchange $`J`$. This interaction also can play a role in the relaxation and dephasing of spin in transport processes in low-dimensional structures.
This work was supported by the DARPA QUIST program, by ONR and by the ONR Nanoscale Electronics Program.
|
warning/0506/hep-th0506027.html
|
ar5iv
|
text
|
# Quantum supersymmetric Toda-mKdV hierarchies
## 1 Introduction
During more than quarter of a century both classical and quantum affine Toda field theories and related generalized (m)KdV ((modified) Korteweg-de Vries) hierarchies were extensively studied (see e.g. -). These theories are integrable and have the L-A pair (or zero curvature) formulation. The most famous are sine-Gordon and $`A_2^{(2)}`$ models (see e.g. , ) and the associated KdV and reduced Boussinesq hierarchies. These theories allow the supersymmetric and fermionic generalizations - both called “super” because their underlying algebraic structures are affine Lie superalgebras.
The super-Toda field theory appears to be supersymmetric if and only if the associated superalgebra possesses the purely fermionic system of simple roots. One should note, however, that usually the superalgebras allow a few simple root systems , they correspond to different Toda theories and only purely fermionic system corresponds to the supersymmetric one.
The supersymmetric version of the Drinfeld-Sokolov reduction applied to the matrix L-operator of the Toda-mKdV theories gives the generators of the related super W-algebra with the commutation relations provided by the associated Hamiltonian structure (see e.g. , ).
In this paper we consider these theories from a point of view of quantum inverse scattering method (QISM) (see , ). There are two approaches of applying QISM to Toda-mKdV type models. The first one is more traditional and based on the quantization of the corresponding lattice systems (see e.g. ,) and the second one is based on the quantization in terms of continuous free field theory and was introduced in , .
Here we use the second approach, generalizing our results obtained in - for the affine superalgebras of rank 2 to the case of general affine superalgebra.
We build the quantum generalization of the monodromy matrix and prove that the related auxiliary $`𝐋`$-operators (which are equal to the monodromy matrix multiplied by the exponential of the elements from the Cartan subalgebra) satisfy the RTT-relation , , while the quantum counterpart of the monodromy matrix itself is shown to satisfy the specialization of the reflection equation (see e.g. ). This provides the quantum integrability relation for the supertraces of the monodromy matrix taken in different representations (transfer matrices). Moreover, it is proven that the auxiliary $`𝐋`$-operators are related with the universal R-matrix associated with the underlying quantum affine superalgebra, with the lower Borel subalgebra represented by the vertex operators from the corresponding Toda field theory. Using this relation in the case when the simple root system is purely fermionic, it is demonstrated that the transfer matrix is invariant under the supersymmetry transformation as it was shown on particular classical examples of Toda field theories , .
In the last two sections the above constructions are illustrated by means of two important examples of integrable hierarchies: quantum super-KdV and SUSY N=1 KdV . These hierarchies generate two integrable structures of the superconformal field theory, the second one is invariant under the SUSY transformation while the first one is not.
## 2 Bosonic Toda-mKdV hierarchies
Each mKdV hierarchy and related Toda field theory associated with affine Lie algebra are generated by the following L-operator :
$`=_u_u\varphi ^i(u)H^i({\displaystyle \underset{i=0}{\overset{r}{}}}e_{\alpha _i}),`$ (1)
where $`u`$ lies on a cylinder of circumference $`2\pi `$, $`\varphi ^i`$ are the scalar fields with the Poisson brackets:
$`\{_u\varphi ^i(u),_v\varphi ^j(v)\}=\delta ^{ij}\delta ^{}(uv)`$ (2)
with quasiperiodic boundary condition:
$`\varphi ^i(u+2\pi )=\varphi ^i(u)+2\pi ip^i.`$ (3)
$`e_{\alpha _i}`$ are the Chevalley generators of the underlying affine Lie algebra and $`H^i`$ $`(i=1,\mathrm{},r)`$ form a basis in the Cartan subalgebra of the corresponding simple Lie algebra:
$`[H^i,e_{\alpha _k}]=\alpha _k^ie_{\alpha _k},[e_{\alpha _k},e_{\alpha _l}]=\delta _{kl}h_{\alpha _k},ad_{e_{\pm \alpha _k}}^{1a_{kj}}e_{\pm \alpha _j}=0,`$ (4)
where $`a_{kj}`$ is a Cartan matrix and $`h_{\alpha _k}(\alpha _k,H)=\alpha _k^iH^i`$. In our case this algebra is considered in evaluation representations, when $`e_{\alpha _0}=\lambda e_\theta `$ (this corresponds to the case of untwisted affine Lie algebra, the twisted case is more complicated), $`\theta `$ is the highest root of the related simple Lie algebra. The classical monodromy matrix for the linear problem associated with the L-operator (1) can be expressed in the following way :
$`\pi _s(\lambda )(𝐌)𝐌_𝐬(\lambda )=e^{2\pi ip^kH^k}Pexp{\displaystyle _0^{2\pi }}u({\displaystyle \underset{i=0}{\overset{r}{}}}e^{(\alpha _i,\varphi (u))}e_{\alpha _i}),`$ (5)
where $`\pi _s`$ is some evaluation representation of the corresponding affine Lie algebra. Defining the auxiliary $`𝐋`$-matrix:
$`𝐋(\lambda )=e^{\pi ip^kH^k}𝐌(\lambda )`$ (6)
one can find that the quadratic Poisson bracket relation is valid:
$`\{𝐋(\lambda ),𝐋(\mu )\}=[𝐫(\lambda \mu ^1),𝐋(\lambda )𝐋(\mu )],`$ (7)
where $`𝐫(\lambda )`$ is the trigonometric r-matrix , related with the corresponding simple Lie algebra. The traces of the monodromy matrices in different evaluation representations $`\pi _s`$:
$`𝐭_𝐬(\lambda )=\pi _s(𝐌(\lambda ))`$ (8)
are in involution under the Poisson brackets:
$`\{𝐭_𝐬(\lambda ),𝐭_𝐬^{}(\mu )\}=0.`$ (9)
The quantization means that we move from the quadratic Poisson bracket relation (7) to the RTT-relation with the underlying affine Lie algebra deformed to the quantum affine algebra (see e.g. , and below). In this section we give the generalization of the constructions appeared in , .
First, let’s quantize the scalar fields $`\varphi ^i`$:
$`\varphi ^k(u)=iQ^k+iP^ku+{\displaystyle \underset{n}{}}{\displaystyle \frac{a_n^k}{n}}e^{inu},`$ (10)
$`[Q^k,P^j]={\displaystyle \frac{i}{2}}\beta ^2\delta ^{kj},[a_n^k,a_m^j]={\displaystyle \frac{\beta ^2}{2}}n\delta ^{kj}\delta _{n+m,0},`$
and define the vertex operators
$`\stackrel{~}{V}_{\alpha _k}(u):=:e^{(\alpha _k,\varphi (u))}:=\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha _k,a_n)}{n}}e^{inu}\right)`$
$`\mathrm{exp}\left(i((\alpha _k,Q)+(\alpha _k,P)u)\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha _k,a_n)}{n}}e^{inu}\right).`$
Then one can define the quantum generalizations of auxiliary $`𝐋`$-operators :
$`𝐋^{(q)}(\lambda )=e^{\pi iP^kH^k}Pexp{\displaystyle _0^{2\pi }}u({\displaystyle \underset{i=0}{\overset{r}{}}}:e^{(\alpha _i,\varphi (u))}:e_{\alpha _i}),`$ (11)
where now $`e_{\alpha _i}`$, $`H^k`$ are the generators of the corresponding quantum affine algebra:
$`[H^i,e_{\alpha _k}]=\alpha _k^ie_{\alpha _k},[e_{\alpha _k},e_{\alpha _l}]=\delta _{kl}[h_{\alpha _k}]_q,ad_{e_{\pm \alpha _k}}^{(q)^{1a_{kj}}}e_{\pm \alpha _j}=0,`$ (12)
where $`q=e^{i\pi \frac{\beta ^2}{2}}`$, $`[a]_q=(q^aq^a)/(qq^1)`$ and $`ad_{e_\alpha }^{(q)}e_\beta =e_\alpha e_\beta q^{(\alpha ,\beta )}e_\beta e_\alpha `$. The object (11) is defined in the interval $`0<\beta ^2<2b^1`$, $`b=max(|b_{ij}|)`$ $`(ij)`$, where $`b_{ij}`$ is a symmetrized Cartan matrix, but can be analytically continuated to a wider region. The quantum monodromy matrix is defined using the relation (6):
$`𝐌^{(q)}(\lambda )=e^{\pi iP^kH^k}𝐋^{(q)}(\lambda ).`$ (13)
It can be shown that the $`𝐋^{(q)}`$ operator satisfies the mentioned RTT-relation in the two ways , . The first way is to consider the following product:
$`(𝐋^{(q)}(\lambda )I)(I𝐋^{(q)}(\mu )).`$ (14)
Then, moving all the Cartan multipliers to the left we find the following:
$`e^{i\pi P^i\mathrm{\Delta }(H^i)}Pexp{\displaystyle _0^{2\pi }}u\stackrel{~}{K}_1(u)Pexp{\displaystyle _0^{2\pi }}uK_2(u),`$ (15)
$`K_1(u)={\displaystyle \underset{j=0}{\overset{r}{}}}:e^{(\alpha _j,\varphi (u))}:e_{\alpha _j}q^{h_{\alpha _j}},K_2(u)={\displaystyle \underset{j=0}{\overset{r}{}}}1:e^{(\alpha _j,\varphi (u))}:e_{\alpha _j},`$
where $`\mathrm{\Delta }(H^i)=H^iI+IH^i`$. The commutation relations between vertex operators on a circle:
$`\stackrel{~}{V}_{\alpha _k}(u)\stackrel{~}{V}_{\alpha _j}(u^{})=q^{b_{kj}}\stackrel{~}{V}_{\alpha _j}(u^{})\stackrel{~}{V}_{\alpha _k}(u),u>u^{},`$ (16)
lead to
$`[\stackrel{~}{K}_1(u),K_2(u^{})]=0,u<u^{}.`$ (17)
Due to this property one can unify two P-exponents into the single one which is equal to
$`\pi _s(\lambda )\pi _s^{}(\mu )\mathrm{\Delta }(𝐋^{(q)}),`$ (18)
where $`\pi _s`$ and $`\pi _s^{}`$ are some evaluation representations, and the coproduct $`\mathrm{\Delta }`$ of the quantum affine (super)algebra is defined by:
$`\mathrm{\Delta }(H^i)=H^iI+IH^i,\mathrm{\Delta }(e_{\alpha _j})=e_{\alpha _j}q^{h_{\alpha _j}}+1e_{\alpha _j},`$ (19)
$`\mathrm{\Delta }(e_{\alpha _j})=e_{\alpha _j}1+q^{h_{\alpha _j}}e_{\alpha _j}.`$
Next, considering opposite product of the $`𝐋`$-operators, one finds that it coincide with opposite coproduct of the $`𝐋`$-operators:
$`(I𝐋^{(𝐪)}(\mu ))(𝐋^{(𝐪)}(\lambda )I)=\pi _s(\lambda )\pi _s(\mu )\mathrm{\Delta }^{op}(𝐋^{(𝐪)}),`$ (20)
where $`\mathrm{\Delta }^{op}=\tau \mathrm{\Delta }`$ and the map $`\tau `$ is defined as follows: $`\tau (ab)=ba`$. Using the property of the universal R-matrix, namely $`𝐑\mathrm{\Delta }=\mathrm{\Delta }^{op}𝐑`$ , we arrive to the RTT-relation:
$`𝐑(\lambda \mu ^1)(𝐋^{(q)}(\lambda )I)(I𝐋^{(q)}(\mu ))=`$ (21)
$`(I𝐋^{(q)}(\mu ))(𝐋^{(q)}(\lambda )I)𝐑(\lambda \mu ^1).`$
Remembering the expression for the monodromy matrix (13) one obtains that the RTT-relation is no longer valid for the monodromy matrices, however, it is easy to see that multiplying both RHS and LHS of (21) by $`e^{i\pi \mathrm{\Delta }H^kP^k}`$ one obtains:
$`𝐑_{12}(\lambda \mu ^1)\stackrel{~}{𝐌}_1^{(q)}(\lambda )𝐌_2^{(q)}(\mu )=\stackrel{~}{𝐌}_2^{(q)}(\mu )𝐌_1^{(q)}(\lambda )𝐑_{12}(\lambda \mu ^1),`$ (22)
where we have denoted $`\stackrel{~}{𝐌}_1^{(q)}(\lambda )`$ the monodromy matrix with $`e_{\alpha _i}1`$ replaced by $`e_{\alpha _i}q^{h_{\alpha _i}}`$ and $`\stackrel{~}{𝐌}_2^{(q)}(\lambda )`$ the monodromy matrix with $`1e_{\alpha _i}`$ replaced by $`q^{h_{\alpha _i}}e_{\alpha _i}`$. Taking the trace, the above additional Cartan factors in $`\stackrel{~}{𝐌}^{(q)}`$ cancel and we obtain the quantum integrability condition for their traces (transfer-matrices):
$`[𝐭_𝐬(\lambda ),𝐭_𝐬^{}(\mu )]=0.`$ (23)
Another more universal way to obtain the RTT-relation is the correspondence between the reduced universal R-matrix (see below) and the P-exponential form of the auxiliary $`𝐋^{(q)}`$ operator. That is, let’s consider integrals of vertex operators:
$`V_{\alpha _k}(u_2,u_1)={\displaystyle \frac{1}{qq^1}}{\displaystyle _{u_1}^{u_2}}u\stackrel{~}{V}_{\alpha _k}(u).`$ (24)
Via the contour technique one can show that these objects satisfy the quantum Serre relations of the lower Borel subalgebra of the associated quantum affine algebra with simple roots $`\alpha _k`$. Using the structure of the reduced universal R-matrix (see e.g. and Appendix) we can write $`\overline{R}=K^1𝐑=\overline{R}(\overline{e}_{\alpha _i},\overline{e}_{\alpha _i})`$, where
$`\overline{e}_{\alpha _i}=e_{\alpha _i}1,\overline{e}_{\alpha _i}=1e_{\alpha _i},`$ (25)
$`𝐑`$ is a universal R-matrix and $`K`$ depends on the elements from Cartan subalgebra, because $`\overline{R}`$ is represented as a power series of these elements. Then, following and using the fundamental feature of the universal R-matrix:
$`(I\mathrm{\Delta })𝐑=𝐑^{13}𝐑^{12},`$ (26)
one can show that the reduced R-matrix has the following property:
$$\overline{R}(\overline{e}_{\alpha _i},e_{\alpha _i}^{}+e_{\alpha _i}^{\prime \prime })=\overline{R}(\overline{e}_{\alpha _i},e_{\alpha _i}^{})\overline{R}(\overline{e}_{\alpha _i},e_{\alpha _i}^{\prime \prime }),$$
(27)
where
$`e_{\alpha _i}^{}`$ $`+`$ $`e_{\alpha _i}^{\prime \prime }=(I\mathrm{\Delta })(1e_{\alpha _i}),`$ (28)
$`e_{\alpha _i}^{}`$ $`=`$ $`1q^{h_{\alpha _i}}e_{\alpha _i},e_{\alpha _i}^{\prime \prime }=1e_{\alpha _i}1,\overline{e}_{\alpha _i}=e_{\alpha _i}11.`$
Their commutation relations are
$`e_{\alpha _i}^{}\overline{e}_{\alpha _j}`$ $`=`$ $`\overline{e}_{\alpha _j}e_{\alpha _i}^{},e_{\alpha _i}^{\prime \prime }\overline{e}_{\alpha _j}=\overline{e}_{\alpha _j}e_{\alpha _i}^{\prime \prime },`$ (29)
$`e_{\alpha _i}^{}e_{\alpha _j}^{\prime \prime }`$ $`=`$ $`q^{b_{ij}}e_{\alpha _j}^{\prime \prime }e_{\alpha _i}^{}.`$
Now, denoting by $`\overline{𝐋}^{(q)}(u_2,u_1)`$ the reduced R-matrix with $`e_{\alpha _i}`$ represented by $`V_{\alpha _i}(u_2,u_1)`$ and using the above property of $`\overline{R}`$ with $`e_{\alpha _i}^{}`$ replaced by appropriate vertex operators we find:
$`\overline{𝐋}^{(q)}(u_3,u_1)=\overline{𝐋}^{(q)}(u_3,u_2)\overline{𝐋}^{(q)}(u_2,u_1),u_3u_2u_1.`$ (30)
Hence, $`\overline{𝐋}^{(q)}`$ has the property of P-exponent. When the interval $`\delta =[u_2,u_1]`$ is small enough one can show that
$`\overline{𝐋}^{(q)}(u_2,u_1)=1+{\displaystyle _{u_1}^{u_2}}u({\displaystyle \underset{i=0}{\overset{r}{}}}:e^{(\alpha _i,\varphi (u))}:e_{\alpha _i})+O(\delta ^2).`$ (31)
That is, we obtain that
$`\overline{𝐋}^{(q)}(u_2,u_1)=Pexp{\displaystyle _{u_1}^{u_2}}u({\displaystyle \underset{i=0}{\overset{r}{}}}:e^{(\alpha _i,\varphi (u))}:e_{\alpha _i})`$ (32)
and $`𝐋^{(q)}=e^{i\pi H^iP^i}\overline{𝐋}^{(q)}(2\pi ,0)`$ satisfies RTT relation by construction.
## 3 Quantum P-exponential and Toda-mKdV hierarchies based on superalgebras
Now let’s generalize the above results to the case when the underlying algebraic structures and integrable hierarchies are related to the affine Lie superalgebra. In the previous part we have moved from classical theory to the quantum one, here we will go in opposite direction, moving from the quantum version of the monodromy matrix and related auxiliary $`𝐋`$-operators, satisfying RTT-relation to their classical counterparts.
First, let’s introduce two types of vertex operators, bosonic and fermionic ones:
$`W_{\alpha _i}^F(u)`$ $``$ $`{\displaystyle \theta }:e^{(\alpha _i,\mathrm{\Phi }(u,\theta ))}:={\displaystyle \frac{i}{\sqrt{2}}}(\alpha _i,\xi (u)):e^{(\alpha _i,\varphi (u))}:`$ (33)
$`W_{\alpha _i}^B(u)`$ $``$ $`{\displaystyle }\theta \theta :e^{(\alpha _i,\mathrm{\Phi }(u,\theta ))}:=:e^{(\alpha _i,\varphi (u))}:,`$ (34)
where $`\mathrm{\Phi }^k`$ are the superfields: $`\mathrm{\Phi }^k(u,\theta )=\varphi ^k(u)\frac{i}{\sqrt{2}}\theta \xi ^k(u)`$ and $`\theta `$ is a Grassmann variable. Their commutation relations on a circle are:
$`W_{\alpha _i}^s(u)W_{\alpha _k}^s^{}(u^{})=(1)^{p(s)p(s^{})}q^{b_{ik}}W_{\alpha _k}^s^{}(u^{})W_{\alpha _i}^s(u),u>u^{},`$ (35)
where $`b_{kj}`$ is the symmetrized Cartan matrix for the corresponding affine Lie superalgebra, $`s`$, $`s^{}`$ are $`B`$, $`F`$ and $`p(F)=1,p(B)=0`$ .
The mode expansion for the bosonic fields is the same as in (10) and for fermionic fields $`\xi ^k(u)`$ is the following:
$`\xi ^l(u)=i^{1/2}{\displaystyle \underset{n}{}}\xi _n^le^{inu},\{\xi _n^k,\xi _m^l\}=\beta ^2\delta ^{kl}\delta _{n+m,0}.`$ (36)
These fermion fields may satisfy two boundary conditions periodic and antiperiodic $`\xi ^i(u+2\pi )=\pm \xi ^i(u)`$ corresponding to the two sectors of (S)CFT – Ramond (R) and Neveu-Schwarz (NS) (the supersymmetry operator appears only when all fermions are in the R sector).
It can be shown that the integrals of the introduced vertex operators as in the bosonic case satisfy the Serre and “non Serre” relations (see e.g. \- ) for the lower Borel subalgebra:
$$ad_{e_{\alpha _k}}^{(q)^{1a_{kj}}}e_{\alpha _j}=0,[[e_{\pm \alpha _r},e_{\pm \alpha _s}]_q,[e_{\pm \alpha _r},e_{\pm \alpha _p}]_q]_q=0,$$
(37)
where the $`ad^{(q)}`$ operator is defined in the following way: $`ad_{e_\alpha }^{(q)}e_\beta =e_\alpha e_\beta (1)^{p(\alpha )p(\beta )}q^{(\alpha ,\beta )}e_\beta e_\alpha `$, and $`e_{\pm \alpha _r}`$ is a so-called “grey” root (for more details see Appendix) of the quantum affine superalgebra with the corresponding bosonic and fermionic roots $`\alpha _i`$.
Substituting them with the appropriate multiplier $`(qq^1)^1`$ in the reduced R-matrix one can find (easily generalizing the results of section 2 to the case of superalgebra) that it satisfies the P-exponential multiplication property:
$`\overline{𝐋}^{(q)}(u_2,u_1)=Pexp^{(q)}{\displaystyle _{u_1}^{u_2}}u({\displaystyle \underset{f}{}}W_{\alpha _f}^F(u)e_{\alpha _f}+{\displaystyle \underset{b}{}}W_{\alpha _b}^B(u)e_{\alpha _b})`$ (38)
$`\overline{𝐋}^{(q)}(u_3,u_1)=\overline{𝐋}^{(q)}(u_3,u_2)\overline{𝐋}^{(q)}(u_2,u_1),u_3u_2u_1,`$
where indices $`f`$ and $`b`$ imply that we are summing over fermionic and bosonic simple roots. The letter $`q`$ over the $`Pexp`$ means that the object introduced above in some cases (more precisely when a number of fermionic roots is more than one) cannot be written as P-exponential for any value of the deformation parameter due to the singularities in the operator products generated by the fermion fields $`\xi ^i`$. Thus we call this object quantum P-exponential.
Defining then $`𝐋^{(q)}e^{\pi ip^iH^i}\overline{𝐋}^{(q)}(2\pi ,0)`$ we find (similarly to the purely bosonic case) that it satisfies the RTT relation (21) and defining the monodromy matrix $`𝐌^{(q)}e^{\pi ip^iH^i}𝐋^{(q)}`$ we again arrive to the property (22) and obtain again the quantum integrability condition (23) for $`𝐭^{(q)}=str𝐌^{(q)}`$. We mention here that the relation (22) can be rewritten in a more universal way, as a specialization of the reflection equation :
$`\stackrel{~}{𝐑}_{12}(\lambda \mu ^1)𝐌_1^{(q)}(\lambda )F_{12}^1𝐌_2^{(q)}(\mu )=𝐌_2^{(q)}(\mu )F_{12}^1𝐌_1^{(q)}(\lambda )𝐑_{12}(\lambda \mu ^1),`$ (39)
where $`F=K^1`$ the Cartan’s factor from the universal R-matrix (see Appendix), and $`\stackrel{~}{𝐑}_{12}(\lambda \mu ^1)=F_{12}^1𝐑_{12}(\lambda \mu ^1)F_{12}`$.
Now let’s analyse the classical limit of the defined objects. We will use the P-exponential property of $`\overline{𝐋}^{(q)}(2\pi ,0)`$. Let’s decompose $`\overline{𝐋}^{(q)}(2\pi ,0)`$ in the following way:
$$\overline{𝐋}^{(q)}(2\pi ,0)=\underset{N\mathrm{}}{lim}\underset{m=1}{\overset{N}{}}\overline{𝐋}^{(q)}(x_m,x_{m1}),$$
(40)
where we divided the interval $`[0,2\pi ]`$ into infinitesimal intervals $`[x_m,x_{m+1}]`$ with $`x_{m+1}x_m=ϵ=2\pi /N`$. Let’s find the terms that can give contribution of the first order in $`ϵ`$ in $`\overline{𝐋}^{(q)}(x_m,x_{m1})`$. In this analysis one needs the operator product expansion of fermion fields and vertex operators:
$`\xi ^k(u)\xi ^l(u^{})={\displaystyle \frac{i\beta ^2\delta ^{kl}}{(iuiu^{})}}+{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}c_p^{kl}(u)(iuiu^{})^p,`$ (41)
$`:e^{(\alpha _k,\varphi (u))}::e^{(\alpha _l,\varphi (u^{}))}:=`$
$`(iuiu^{})^{\frac{(\alpha _k,\alpha _l)\beta ^2}{2}}(:e^{(\alpha _k+\alpha _l,\varphi (u))}:+{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}d_p^{kl}(u)(iuiu^{})^p),`$
where $`c_p^{kl}(u)`$ and $`d_p^{kl}(u)`$ are operator-valued functions of $`u`$. Now one can see that only two types of terms can give the contribution of the order $`ϵ`$ in $`\overline{𝐋}^{(q)}(x_{m1},x_m)`$ when $`q1`$. The first type consists of operators of the first order in $`W_{\alpha _i}`$ and the second type is formed by the operators, quadratic in $`W_{\alpha _i}`$, which give contribution of the order $`ϵ^{1\pm \beta ^2}`$ by virtue of operator product expansion. Let’s look on the terms of the second type in detail.
The terms of the second type appear from the quadratic products of vertex operators arising from:
i)the composite roots (more precisely q-commutators of two fermionic roots),
ii)the quadratic terms of the q-exponentials which are present in the universal R-matrix.
At first we consider terms emerging from composite roots, which have the following form (see Appendix):
$`{\displaystyle \frac{1}{a(\alpha _i+\alpha _j)(qq^1)}}[e_{\alpha _j},e_{\alpha _i}]_{q^1}({\displaystyle _{x_{m1}}^{x_m}}u_1W_{\alpha _i}(u_1i0)`$
$`{\displaystyle _{x_{m1}}^{x_m}}u_2W_{\alpha _j}(u_2+i0)+q^{b_{ij}}{\displaystyle _0^{2\pi }}u_2W_{\alpha _j}(u_2i0){\displaystyle _0^{2\pi }}u_1W_{\alpha _i}(u_1+i0).`$ (42)
Using the fact that $`q=e^{i\pi \frac{\beta ^2}{2}}`$ and that in the limit $`\beta ^20`$, $`a(\alpha _i+\alpha _j)b_{ij}`$, one can rewrite this as follows (leaving only terms that can give contribution to the first order in $`ϵ`$:
$`(2\pi i)^1[e_{\alpha _j},e_{\alpha _i}]`$
$`{\displaystyle _{x_{m1}}^{x_m}}u_1{\displaystyle _{x_{m1}}^{x_m}}u_2\left({\displaystyle \frac{1}{u_2u_1+i0}}{\displaystyle \frac{1}{u_2u_1i0}}\right):e^{(\alpha _i+\alpha _j,\varphi (u_2))}:.`$ (43)
Now using the well known formula
$`{\displaystyle \frac{1}{x+i0}}{\displaystyle \frac{1}{xi0}}=2i\pi \delta (x),`$ (44)
we obtain that (3) in the classical limit gives
$`[e_{\alpha _j},e_{\alpha _i}]{\displaystyle _{x_{m1}}^{x_m}}u:e^{(\alpha _i+\alpha _j,\varphi (u))}:.`$ (45)
Next let’s consider the quadratic products arising from quadratic parts of q-exponentials of fermionic roots. They look as follows:
$$\frac{1}{(2)_{q_{\alpha _i}^1}}_{x_{m1}}^{x_m}u_1W_{\alpha _i}(u_1i0)_{x_{m1}}^{x_m}u_2W_{\alpha _i}(u_2+i0)e_{\alpha _i}^2.$$
(46)
One can rewrite this product via the ordered integrals:
$$\frac{q^{b_{ii}}1}{(2)_{q_{\alpha _i}^1}}_{x_{m1}}^{x_m}u_1W_{\alpha _i}(u_1)_{x_{m1}}^{u_1}u_2W_{\alpha _i}(u_2)e_{\alpha _i}^2.$$
(47)
In the limit $`\beta ^20`$ we obtain (forgetting about the terms that could give contribution of the order $`ϵ^2`$):
$$\frac{ib_{ii}\beta ^2}{2}_{x_{m1}}^{x_m}u_1_{x_{m1}}^{u_1}u_2(iu_1iu_2)^{\frac{b_{ii}\beta ^2}{2}1}e^{2(\alpha _i,\varphi (u_2))}e_{\alpha _i}^2.$$
(48)
Therefore the final contribution is:
$$_{x_{m1}}^{x_m}ue^{2(\alpha _i,\varphi (u))}e_{\alpha _i}^2.$$
(49)
Collecting now all the terms of order $`ϵ`$ we find:
$`\overline{𝐋}^{(q)}(x_m,x_{m1})=1+{\displaystyle _{x_{m1}}^{x_m}}u({\displaystyle \underset{f}{}}W_{\alpha _f}^F(u)e_{\alpha _f}+{\displaystyle \underset{b}{}}W_{\alpha _b}^B(u)e_{\alpha _b}+`$
$`{\displaystyle \underset{f_1f_2}{}}[e_{\alpha _{f_1}},e_{\alpha _{f_2}}]W_{\alpha _{f_1}+\alpha _{f_2}}^B(u))+O(ϵ^2).`$ (50)
Gathering the $`\overline{𝐋}^{(q)}(x_m,x_{m1})`$ it is easy to see that in the $`q1`$ limit
$`\overline{𝐋}^{(q)}(x_m,x_{m1})`$ is equal to:
$`\overline{𝐋}^{(cl)}(2\pi ,0)=Pexp{\displaystyle _0^{2\pi }}u({\displaystyle \underset{f}{}}W_{\alpha _f}^F(u)e_{\alpha _f}+{\displaystyle \underset{b}{}}W_{\alpha _b}^B(u)e_{\alpha _b}`$
$`{\displaystyle \underset{f_1f_2}{}}[e_{\alpha _{f_1}},e_{\alpha _{f_2}}]W_{\alpha _{f_1}+\alpha _{f_2}}^B(u)).`$ (51)
Defining then $`𝐋^{(cl)}e^{\pi ip^iH^i}\overline{𝐋}^{(cl)}(2\pi ,0)`$ we find that it satisfies the quadratic Poisson bracket relation (7) and defining the monodromy matrix $`𝐌^{(cl)}e^{\pi ip^iH^i}𝐋^{(cl)}`$ we again obtain the classical integrability condition (9).
Now let’s find the L-operator which corresponds to the monodromy matrix defined above. Let’s consider the following one:
$$_F=D_{u,\theta }D_{u,\theta }\mathrm{\Phi }^i(u,\theta )H^i(\underset{f}{}e_{\alpha _f}+\underset{b}{}\theta e_{\alpha _b}),$$
(52)
where $`D_{u,\theta }=_\theta +\theta _u`$ is a superderivative and $`\mathrm{\Phi }^i`$ are the classical superfields with the following Poisson brackets:
$$\{D_{u,\theta }\mathrm{\Phi }^i(u,\theta ),D_{u^{},\theta ^{}}\mathrm{\Phi }^j(u^{},\theta ^{})\}=\delta ^{ij}D_{u,\theta }(\delta (uu^{})(\theta \theta ^{})).$$
(53)
Making a gauge transformation of the above L-operator one can arrive to the fields, satisfying classical version of super W-algebras with the commutation relations provided by the Poisson brackets , .
The associated “fermionic” linear problem can be reduced to the “bosonic” one. The linear problem
$$_F\mathrm{\Psi }(u,\theta )=(D_{u,\theta }+N_1+\theta N_0)(\chi +\theta \eta ),$$
(54)
where $`\mathrm{\Psi }(u,\theta )=\chi +\theta \eta `$, $`N_1=\frac{i}{\sqrt{2}}\xi ^i(u)H^i_fe_{\alpha _f}`$, $`N_0=_u\varphi ^i(u)H^i_be_{\alpha _b}`$, can be reduced to the linear problem on $`\chi `$:
$$_B\chi (u)=(_u+N_1^2+N_0)\chi (u).$$
(55)
That is:
$$_B=_u_u\varphi ^i(u)H^i+(\frac{i}{\sqrt{2}}\xi ^i(u)H^i\underset{f}{}e_{\alpha _f})^2\underset{b}{}e_{\alpha _b}.$$
(56)
One can easily see that the monodromy matrix for the corresponding linear problem is that described above.
## 4 Integrals of Motion and Supersymmetry Invariance
It is well known that (both classical and quantum) integrability conditions lead to the involutive family of (both local and nonlocal) integrals of motion (IM). For super- versions of these systems it is also known that sometimes it is possible to include supersymmetry generator
$`G_0=\beta ^2\sqrt{2}i^{1/2}{\displaystyle _0^{2\pi }}𝑑u\varphi ^l(u)\xi ^l(u)={\displaystyle \underset{l=0}{\overset{r}{}}}{\displaystyle \underset{nZ}{}}\beta ^2\xi _n^la_n^l`$ (57)
in these series . Here we will show that the transfer matrix $`𝐭^{(q)}(\lambda )=str𝐌^{(q)}(\lambda )`$ commute with $`G_0`$ if the simple root system is purely fermionic, that is:
$`𝐭^{(q)}(\lambda )=str(\pi (\lambda )(e^{2i\pi P^kH^k}Pexp{\displaystyle _0^{2\pi }}u({\displaystyle \underset{f=0}{\overset{r}{}}}W_{\alpha _f}^F(u)e_{\alpha _f}))),`$ (58)
where $`\pi `$ denote some representation of the corresponding superalgebra in which the supertrace is taken.
We note the crucial property:
$`[G_0,W_\alpha ^F(u)]=_uW_\alpha ^B(u)`$ (59)
Integrating over $`u`$ and multiplying by the appropriate coefficient one obtains:
$`[G_0,e_{\alpha _i}]={\displaystyle \frac{W_{\alpha _i}^F(0)W_{\alpha _i}^F(2\pi )}{qq^1}},`$ (60)
where $`e_{\alpha _i}`$ is represented by the vertex operator (see previous Section). Next, we will use the important Proposition (Prop. 1, Sec. 3.1 of ): For the objects $`A_i`$, $`B_i`$, $`𝐈`$, satisfying the commutation relations
$`[𝐈,e_{\alpha _i}]={\displaystyle \frac{A_iB_i}{qq^1}},A_ie_{\alpha _j}=q^{b_{ij}}e_{\alpha _j}A_i,B_ie_{\alpha _j}=q^{b_{ij}}e_{\alpha _j}B_i,`$ (61)
the following relation holds:
$`[1𝐈,\overline{R}]=\overline{R}({\displaystyle \underset{i}{}}e_{\alpha _i}A_i)({\displaystyle \underset{i}{}}e_{\alpha _i}B_i)\overline{R}.`$ (62)
Applying this relations to our case we find (identifying $`A_i`$ with $`W_{\alpha _i}^F(0)`$, $`B_i`$ with $`W_{\alpha _i}^F(2\pi )`$ and $`𝐈`$ with $`G_0`$):
$`[G_0,\overline{𝐋}^{(q)}(2\pi ,0)]=\overline{𝐋}^{(q)}(2\pi ,0)W^B(0)W^B(2\pi )\overline{𝐋}^{(q)}(2\pi ,0),`$ (63)
where
$`W^B(u)={\displaystyle \underset{i=0}{\overset{r}{}}}W_{\alpha _i}^B(u)e_{\alpha _i}.`$ (64)
Now using the property (35), periodicity properties of vertex operators:
$`W_\alpha ^s(u+2\pi )=q^{(\alpha ,\alpha )}e^{2i\pi (\alpha ,P)}W_{\alpha _i}^s(u+2\pi )`$ (65)
(here $`s=B,F`$) and cyclic property of the supertrace one obtains, multiplying both sides of (63) by $`e^{2i\pi P^kH^k}`$ and taking the supertrace:
$`[G_0,𝐭^{(q)}]=0.`$ (66)
We note here that if there were bosonic simple roots in the construction of the transfer-matrix the above reasonings are no longer valid, because the corresponding bosonic vertex operators commuting with $`G_0`$ give the fermionic vertex operators associated with the same root vector (the same happens when we construct the superstring vertex operators ), but not the total derivative as in (59). It was already shown explicitly on the concrete examples that the hierarchies, based on the partly bosonic simple root systems are not invariant under the supersymmetry transformation (see e.g. , , ).
The affine superalgebras which allow such root systems are of the following type : $`A(m,m)^{(1)}=sl(m+1,m+1)^{(1)}`$, $`A(2m,2m)^{(4)}=sl(2m+1,2m+1)^{(4)}`$, $`A(2m+1,2m+1)^{(2)}=sl(2m+2,2m+2)^{(2)}`$, $`A(2m+1,2m)^{(2)}=sl(2m+2,2m+1)^{(2)}`$, $`B(m,m)^{(1)}=osp(2m+1,2m)^{(1)}`$, $`D(m+1,m)^{(1)}=osp(2m+2,2m)^{(1)}`$, $`D(m,m)^{(2)}=osp(2m,2m)^{(2)}`$, $`D(2,1,\alpha )^{(1)}`$.
The involutive family of the (both classical and quantum) IM in the Toda field theories have the property, that the commutators of IM with the corresponding vertex operators reduce to the total derivatives :
$`[I_l,W_{\alpha _k}(u)]=_u(:O_{\alpha _k}^{(l)}(u)W_{\alpha _k}(u):)=_u\mathrm{\Theta }_k^{(l)}(u),`$ (67)
where $`O_\pm ^{(k)}(u)`$ is the polynomial of $`_u\varphi ^i(u)`$, $`\xi ^i(u)`$ and their derivatives. In it was shown in the bosonic case (the proof is similar to the arguments above) that $`I_l`$ commute with the transfer matrix, the generalization of these arguments to the super-case is straightforward.
In the next two sections we will give two examples of the KdV hierarchies related to affine superalgebra $`B(0,1)^{(1)}osp(1|2)^{(1)}`$ (super-KdV) and twisted affine superalgebra $`D(1,1)^{(2)}C(2)^{(2)}sl(2|1)^{(2)}osp(2|2)^{(2)}`$ (SUSY N=1 KdV).
## 5 Example 1: super-KdV hierarchy
The super-KdV model , is based on the following L-operator:
$$D_{u,\theta }+D_{u,\theta }\mathrm{\Phi }(u,\theta )h_{\alpha _0}(e_{\alpha _1}+\theta e_{\alpha _0}),$$
(68)
where $`h_{\alpha _0}`$, $`e_{\alpha _1}`$, $`e_{\alpha _0}`$ are the Chevalley generators of the upper Borel subalgebra of the $`osp(1|2)^{(1)}`$ which are taken in the evaluation representation that is $`h_{\alpha _0}=h`$, $`e_{\alpha _1}=iv_+`$, $`e_{\alpha _0}=\lambda X_{}`$, where $`X_\pm `$, $`v_\pm `$ and $`h`$ are the generators of $`osp(1|2)`$ superalgebra with the following commutation relations:
$`[h,X_\pm ]`$ $`=`$ $`\pm 2X_\pm ,[h,v_\pm ]=\pm v_\pm ,[X_+,X_{}]=h,`$ (69)
$`[v_\pm ,v_\pm ]`$ $`=`$ $`\pm 2X_\pm ,[v_+,v_{}]=h,[X_\pm ,v_{}]=v_\pm ,[X_\pm ,v_\pm ]=0.`$
Here $`p(v_\pm )=1`$, $`p(X_\pm )=0`$, $`p(h)=0`$. The classical monodromy matrix is:
$`𝐌(\lambda )`$ $`=`$ $`e^{2\pi iph_{\alpha _0}}P\mathrm{exp}{\displaystyle _0^{2\pi }}u({\displaystyle \frac{i}{\sqrt{2}}}\xi (u)e^{\varphi (u)}e_{\alpha _1}e^{2\varphi (u)}e_{\alpha _1}^2`$
$`+`$ $`e^{2\varphi (u)}e_{\alpha _0}).`$
The involutive family of the integrals of motion which can be extracted from this monodromy matrix (more precisely they arise as a coefficients in the expansion in $`\lambda ^1`$ series of the trace of the logarithm of $`𝐌`$-matrix) can be expressed via the following fields:
$`U(u)=\varphi ^{\prime \prime }(u)\varphi ^2(u){\displaystyle \frac{1}{2}}\xi (u)\xi ^{}(u),\alpha (u)=\xi ^{}(u)+\xi (u)\varphi ^{}(u)`$ (71)
generating the classical limit of the superconformal algebra under the Poisson brackets:
$`\{U(u),U(v)\}`$ $`=`$ $`\delta ^{\prime \prime \prime }(uv)+2U^{}(u)\delta (uv)+4U(u)\delta ^{}(uv),`$ (72)
$`\{U(u),\alpha (v)\}`$ $`=`$ $`3\alpha (u)\delta ^{}(uv)+\alpha ^{}(u)\delta (uv),`$
$`\{\alpha (u),\alpha (v)\}`$ $`=`$ $`2\delta ^{\prime \prime }(uv)+2U(u)\delta (uv).`$
The integrals of motion are:
$`I_1^{(cl)}={\displaystyle U(u)u},`$ (73)
$`I_3^{(cl)}={\displaystyle \left(U^2(u)/2+\alpha (u)\alpha ^{}(u)\right)u},`$
$`I_5^{(cl)}={\displaystyle \left((U^{})^2(u)2U^3(u)+8\alpha ^{}(u)\alpha ^{\prime \prime }(u)+12\alpha ^{}(u)\alpha (u)U(u)\right)u},`$
$`...`$
The second one $`I_3^{(cl)}`$ gives the super-KdV equation:
$`U_t=U_{uuu}6UU_u6\alpha \alpha _{uu},\alpha _t=4\alpha _{uuu}6U\alpha _u3U_u\alpha .`$ (74)
However, the supersymmetry operator $`G_0=_0^{2\pi }u\alpha (u)`$ can not be included in the pairwise commuting IM that can be easily seen from the second equation of (74) (i.e. $`I_3^{(cl)}`$ does not commute with $`G_0`$). Moving to the quantum case we find that the quantum analogue of the monodromy matrix is:
$`𝐌^{(𝐪)}(\lambda )=e^{2\pi iPh_{\alpha _0}}P\mathrm{exp}{\displaystyle _0^{2\pi }}u({\displaystyle \frac{i}{\sqrt{2}}}\xi (u):e^{\varphi (u)}:e_{\alpha _1}+:e^{2\varphi (u)}:e_{\alpha _0}),`$ (75)
where $`h_{\alpha _0}`$, $`e_{\alpha _0}`$, $`e_{\alpha _1}`$ are now the Chevalley generators of the $`osp_q(1|2)^{(1)}`$. We did not put the letter $`q`$ over the P-exponential because this is the case when for values of $`\beta ^2`$ from the interval (0,2) one can write the above object as a “real” P-exponential (represented via ordered integrals). Due to the presence of bosonic root $`\mathrm{}_{\alpha _0}`$ the trace of the quantum monodromy matrix is not invariant under the supersymmetry transformation as it was in the classical case.
## 6 Example 2: SUSY N=1 KdV hierarchy
The L-operator corresponding to the SUSY N=1 KdV model , is the following one:
$`_F=D_{u,\theta }D_{u,\theta }\mathrm{\Phi }(u,\theta )h_{\alpha _1}(e_{\alpha _0}+e_{\alpha _1}),`$ (76)
where $`h_{\alpha _1}`$, $`e_{\alpha _0}`$, $`e_{\alpha _1}`$ are the Chevalley generators of the twisted affine Lie superalgebra $`C(2)^{(2)}`$ with such set of commutation relations:
$`[h_{\alpha _1},h_{\alpha _0}]=0,[h_{\alpha _0},e_{\pm \alpha _1}]=e_{\pm \alpha _1},[h_{\alpha _1},e_{\pm \alpha _0}]=e_{\pm \alpha _0},`$ (77)
$`[h_{\alpha _i},e_{\pm \alpha _i}]=\pm e_{\pm \alpha _i},[e_{\pm \alpha _i},e_{\alpha _j}]=\delta _{i,j}h_{\alpha _i},(i,j=0,1),`$
$`ad_{e_{\pm \alpha _0}}^3e_{\pm \alpha _1}=0,ad_{e_{\pm \alpha _1}}^3e_{\pm \alpha _0}=0`$
Here $`p(h_{\alpha _{0,1}})=0`$, $`p(e_{\pm \alpha _{0,1}})=1`$, i.e. both simple roots are fermionic. The classical monodromy matrix is:
$`𝐌`$ $`=`$ $`e^{2\pi iph_{\alpha _1}}P\mathrm{exp}{\displaystyle _0^{2\pi }}u({\displaystyle \frac{i}{\sqrt{2}}}\xi (u)e^{\varphi (u)}e_{\alpha _1}`$
$``$ $`{\displaystyle \frac{i}{\sqrt{2}}}\xi (u)e^{\varphi (u)}e_{\alpha _0}e_{\alpha _1}^2e^{2\varphi (u)}e_{\alpha _0}^2e^{2\varphi (u)}[e_{\alpha _1},e_{\alpha _0}]).`$
The series of the integrals of motion starts with the following ones:
$`I_1^{(cl)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle U(u)u},`$ (79)
$`I_3^{(cl)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \left(U^2(u)+\alpha (u)\alpha ^{}(u)/2\right)u},`$
$`I_5^{(cl)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \left(U^3(u)(U^{})^2(u)/2\alpha ^{}(u)\alpha ^{\prime \prime }(u)/4\alpha ^{}(u)\alpha (u)U(u)\right)u},`$
$`...`$
The fields $`U`$ and $`\alpha `$ are defined in terms of the free fields as in previous section, but now one can unify them into one superfield:
$$𝒰(u,\theta )D_{u,\theta }\mathrm{\Phi }(u,\theta )_u\mathrm{\Phi }(u,\theta )D_{u,\theta }^3\mathrm{\Phi }(u,\theta )=\theta U(u)i\alpha (u)/\sqrt{2}.$$
(80)
The second IM from (79) generates the first nontrivial evolution equation, the SUSY $`N`$=1 KdV:
$$𝒰_t=𝒰_{uuu}+3(𝒰D_{u,\theta }𝒰)_u,$$
(81)
or in components:
$`U_t=U_{uuu}6UU_u{\displaystyle \frac{3}{2}}\alpha \alpha _{uu},\alpha _t=4\alpha _{uuu}3(U\alpha )_u.`$ (82)
Now one can see that unlike the previous example the supersymmetry generator $`_0^{2\pi }u\alpha (u)`$ commutes with $`I_3^{(cl)}`$ and can be included in the involutive family of IM . Moreover, the results of the Section 2 yield that the quantum monodromy matrix have the following form:
$$𝐌=e^{2\pi iPh_{\alpha _1}}Pexp^{(q)}_0^{2\pi }u(W_{}(u)e_{\alpha _1}+W_+(u)e_{\alpha _0}),$$
(83)
where $`W_\pm =\theta :e^{\pm \mathrm{\Phi }(u,\theta )}:`$. Due to the fact that both roots are fermionic the supersymmetry generator commutes with the transfer matrix.
This SUSY $`N`$=1 KdV model was studied from a point of view of the Quantum Inverse Scattering Method in . There were constructed the analogues of Baxter’s Q-operator, providing the following functional relations with the transfer-matrices:
$`𝐭_{\frac{1}{4}}(\lambda )𝐐_\pm (\lambda )`$ $`=`$ $`\pm 𝐐_\pm (q^{\frac{1}{2}}\lambda )𝐐_\pm (q^{\frac{1}{2}}\lambda ),`$
$`𝐭_{\frac{1}{2}}(q^{\frac{1}{4}}\lambda )𝐐_\pm (\lambda )`$ $`=`$ $`𝐭_{\frac{1}{4}}(q^{\frac{1}{2}}\lambda )𝐐_\pm (q^{\frac{1}{2}}\lambda )+𝐐_\pm (q\lambda ),`$ (84)
and the fusion relations between transfer-matrices in different representations:
$`𝐭_j(q^{\frac{1}{4}}\lambda )𝐭_j(q^{\frac{1}{4}}\lambda )=𝐭_{j+\frac{1}{4}}(\lambda )𝐭_{j\frac{1}{4}}(\lambda )+(1)^{4j},`$ (85)
which for the case when $`q`$ is a root of unity can be transformed into the Thermodynamic Bethe Ansatz equations of $`D_{2N}`$ type .
It should be noted also that the associated Toda field theory is a well known $`N`$=1 SUSY sinh-Gordon model , with the action:
$$\frac{1}{\beta ^2}^2u^2\theta (D_{u,\theta }\mathrm{\Phi }\overline{D}_{\overline{u},\overline{\theta }}\mathrm{\Phi }+m^2cosh(\mathrm{\Phi }))$$
(86)
## Acknowledgements
We are grateful to Professors F.A. Smirnov and M.A. Semenov-Tian-Shansky for useful discussions. The work was supported by the CRDF (Grant No. RUMI-2622-ST-04) and the Dynasty Foundation.
## 7 Appendix
The affine Lie superalgebra in has the following commutation relations between its Chevalley generators ($`H^i`$ forms a basis in the Cartan subalgebra of the underlying simple Lie superalgebra and $`e_{\pm \alpha _k}`$ are the generators associated with positive and negative simple roots of the whole affine algebra) :
$`[H^i,e_{\alpha _k}]=\alpha _k^ie_{\alpha _k},[e_{\alpha _k},e_{\alpha _l}]=\delta _{kl}[h_{\alpha _k}]_q,ad_{e_{\pm \alpha _k}}^{1a_{kj}}e_{\pm \alpha _j}=0,`$ (87)
$`[[e_{\pm \alpha _r},e_{\pm \alpha _s}]_q,[e_{\pm \alpha _r},e_{\pm \alpha _p}]_q]_q=0`$
if $`(\alpha _r,\alpha _r)=(\alpha _s,\alpha _p)=(\alpha _r,\alpha _s+\alpha _p)=0`$, in this case it is usually said that $`\alpha _r`$ is a “grey” root, which is between two roots $`\alpha _s`$, $`\alpha _p`$ on the Dynkin diagram , . The definition of the super q-commutator is:
$`ad_{e_\alpha }^{(q)}e_\beta =e_\alpha e_\beta (1)^{p(\alpha )p(\beta )}q^{(\alpha ,\beta )}e_\beta e_\alpha ,`$ (88)
where $`p(\alpha )`$ is equal to 1 when $`\alpha `$ is a fermionic root, and to 0 if $`\alpha `$ is a bosonic root. The universal R-matrix for the contragradient Lie superalgebra of finite growth (affine algebra as a particular case) has the following structure:
$$𝐑=K\overline{R}=K(\underset{\alpha \mathrm{\Delta }_+}{\overset{}{}}R_\alpha ),$$
(89)
where $`\overline{R}`$ is a reduced R-matrix and $`R_\alpha `$ are defined by the formulae:
$$R_\alpha =exp_{q_\alpha ^1}((1)^{p(\alpha )}(qq^1)(a(\alpha ))^1(e_\alpha e_\alpha )$$
(90)
for real roots and
$$R_{n\delta }=exp((1)^{p(n\delta )}(qq^1)(\underset{i,j}{\overset{mult}{}}c_{ij}(n)e_{n\delta }^{(i)}e_{n\delta }^{(j)}))$$
(91)
for pure imaginary roots. Here $`\mathrm{\Delta }_+`$ is the reduced positive root system (the bosonic roots which are two times fermionic roots are excluded). The generators corresponding to the composite roots are defined according to the construction of the Cartan-Weyl basis given in . For example the generators of the type $`e_{\pm \alpha _{f_1}\pm \alpha _{f_2}}`$ are constructed by means of the following q-commutators:
$$e_{\alpha _{f_1}+\alpha _{f_2}}=[e_{\alpha _{f_2}},e_{\alpha _{f_1}}]_{q^1},e_{\alpha _{f_1}\alpha _{f_2}}=[e_{\alpha _{f_1}},e_{\alpha _{f_2}}]_q.$$
(92)
The $`a(\alpha )`$ coefficients are defined as follows:
$$[e_\gamma ,e_\gamma ]=a(\gamma )\frac{k_\gamma k_\gamma ^1}{qq^1}.$$
(93)
We will need the values of $`a(\gamma )`$ when $`\gamma `$ is equal to $`\alpha _{f_1}+\alpha _{f_2}`$, where $`\alpha _{f_1}`$ and $`\alpha _{f_2}`$ are fermionic simple roots:
$$a(\alpha _{f_1}+\alpha _{f_2})=\frac{q^{b_{f_1f_2}}q^{b_{f_1f_2}}}{qq^1}.$$
(94)
The q-exponentials in (89) are defined in the usual way:
$`exp_q(x)`$ $`=`$ $`1+x+{\displaystyle \frac{x^2}{(2)_q!}}+\mathrm{}+{\displaystyle \frac{x^n}{(n)_q!}}+\mathrm{}={\displaystyle \underset{n0}{}}{\displaystyle \frac{x^n}{(n)_q!}}`$ (95)
$`(a)_q`$ $``$ $`{\displaystyle \frac{q^a1}{q1}},q_\alpha (1)^{p(\alpha )}q^{(\alpha ,\alpha )}.`$
|
warning/0506/math0506548.html
|
ar5iv
|
text
|
# Inverting weak dihomotopy equivalence using homotopy continuous flow
## 1. Introduction
There are numerous uses of the notion of “category without identities”. For recent papers, see for example . An enriched version of this notion, in the sense of , over the category of general topological spaces can be found in . By considering “small categories without identities” enriched over the category of compactly generated topological spaces, that is weak Hausdorff k-spaces in the sense of , one obtains an object called a flow which allows a model categorical treatment of dihomotopy (directed homotopy). Indeed, a flow $`X`$ can model (the time flow of) a higher dimensional automaton as follows. A flow $`X`$ consists of
1. a set of states $`X^0`$;
2. for each pair of states $`(\alpha ,\beta )X^0\times X^0`$, there is a compactly generated topological space $`_{\alpha ,\beta }X`$ called the path space between $`\alpha `$ and $`\beta `$ representing the concurrency between $`\alpha `$ and $`\beta `$; each element of $`_{\alpha ,\beta }X`$ corresponds to a non-constant execution path from $`\alpha `$ to $`\beta `$; the emptiness of the space $`_{\alpha ,\alpha }X`$ for some state $`\alpha `$ means that there are no loops from $`\alpha `$ to itself; let
$$X=\underset{(\alpha ,\beta )X^0\times X^0}{}_{\alpha ,\beta }X.$$
3. for each triple of states $`(\alpha ,\beta ,\gamma )X^0\times X^0\times X^0`$, there is a strictly associative composition law $`_{\alpha ,\beta }X\times _{\beta ,\gamma }X_{\alpha ,\gamma }X`$ corresponding to the concatenation of non-constant execution paths.
The main problem to model dihomotopy is that contractions in the direction of time are forbidden. Otherwise in the categorical localization of flows with respect to the dihomotopy equivalences, the relevant geometric information is lost . Here is a very simple example. Take two non-constant execution paths going from one initial state to one final state. If contractions in the direction of time were allowed, then one would find in the same equivalence class a loop (cf. Figure 1): this is not acceptable.
Two kinds of deformations are of interest in the framework of flows. The first one is called weak S-homotopy equivalence: it is a morphism of flows $`f:XY`$ such that the set map $`f^0:X^0Y^0`$ is a bijection and such that the continuous map $`f:XY`$ is a weak homotopy equivalence. It turns out that there exists a model structure on the category of flows whose weak equivalences are exactly the weak S-homotopy equivalences ( and Section 4 of this paper). However, the identifications allowed by the weak S-homotopy equivalences are too rigid. So another kind of weak equivalence is required. The T-homotopy equivalences are generated by a set $`𝒯`$ of cofibrations obtained by taking the cofibrant replacement of the inclusions of posets <sup>1</sup><sup>1</sup>1Any poset $`P`$ can be viewed as a flow in an obvious way: the set of states is the underlying set of $`P`$ and there is a non-constant execution path from $`\alpha `$ to $`\beta `$ if and only if $`\alpha <\beta `$. Note that the inequality is strict. Indeed, the ordering of $`P`$ represents the direction of time and the flow associated with a poset must be loopless. of Definition 8. This approach of T-homotopy is presented for the first time in . The latter models “refinement of observation”. For instance, the inclusion of posets $`\{\widehat{0}<\widehat{1}\}\{\widehat{0}<A<\widehat{1}\}`$ corresponds to the identification of a directed segment $`U`$ going from the initial state $`\widehat{0}`$ to the final state $`\widehat{1}`$ with the composite $`U^{}U^{\prime \prime }`$ of two directed segments (cf. Figure 2).
The problem we face can then be presented as follows. We have:
1. A model structure on the category of flows $`\mathrm{𝐅𝐥𝐨𝐰}`$, called the weak S-homotopy model structure, such that the class of weak equivalences is exactly the class $`𝒮`$ of weak S-homotopy equivalences. One wants to invert the weak S-homotopy equivalences because two weakly S-homotopy equivalent flows are equivalent from an observational viewpoint. This model structure provides an implementation of Whitehead’s theorem for S-homotopy only.
2. A set of cofibrations $`𝒯`$ of generating T-homotopy equivalences one would like to invert because these maps model refinement of observation.
3. Three known invariants with respect to weak S-homotopy and T-homotopy: the underlying homotopy type functor , the branching homology and the merging homology .
4. Every model structure on $`\mathrm{𝐅𝐥𝐨𝐰}`$ which contains as weak equivalences the class of morphisms $`𝒮𝒯`$, contains weak equivalences which do not preserve the three known invariants . In particular, the category $`\mathrm{𝐅𝐥𝐨𝐰}[𝒮_𝒯^1]`$ below is not the Quillen homotopy category of a model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$. The left Bousfield localization of the weak S-homotopy model structure with respect to the set of cofibrations $`𝒯`$ is therefore not relevant here.
The negative result (4) prevents us from using the machinery of model category on the category $`\mathrm{𝐅𝐥𝐨𝐰}`$ for understanding the full dihomotopy equivalence relation. There are then several possibilities: reconstructing some pieces of homotopy theory in the framework of flows, finding new categories for studying S-homotopy and T-homotopy, or also relating dihomotopy on $`\mathrm{𝐅𝐥𝐨𝐰}`$ to other axiomatic presentations of homotopy theory. The possibility which is explored in this paper is the first one.
Indeed, the goal of this work is to prove that it is possible to find a full subcategory of the category of flows which is big enough to contain all dihomotopy types and in which the weak dihomotopy equivalences are exactly the invertible morphisms up to dihomotopy. The main theorem of the paper states as follows (cf. Section 4 for a reminder about flows):
###### Theorem 1.
(Theorem 3 and Theorem 4) Let $`J^{gl}`$ be the set of generating trivial cofibrations of the weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$. Let $`𝒯`$ be the set of generating T-homotopy equivalences. Let $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}`$ be the full subcategory of cofibrant flows. There exists a full subcategory $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}`$ of the category of cofibrant flows $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}`$, the one of homotopy continuous flows, a class of morphisms of flows $`𝒮_𝒯`$ and a congruence $`_𝒯`$ on the morphisms of $`\mathrm{𝐅𝐥𝐨𝐰}`$ such that the inclusion functors $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}\mathrm{𝐅𝐥𝐨𝐰}_{cof}\mathrm{𝐅𝐥𝐨𝐰}`$ induce the equivalences of categories
$$\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}/_𝒯\mathrm{𝐅𝐥𝐨𝐰}_{cof}[(𝒮\mathrm{𝐜𝐨𝐟}(J^{gl}T))^1]\mathrm{𝐅𝐥𝐨𝐰}[𝒮_𝒯^1].$$
Moreover, one has:
1. The class of morphisms $`𝒮_𝒯`$ contains the weak S-homotopy equivalences and the morphisms of $`\mathrm{𝐜𝐨𝐟}(J^{gl}T)`$ with cofibrant domains.
2. Every morphism of $`𝒮_𝒯`$ preserves the underlying homotopy type, the branching homology and the merging homology.
We now outline the contents of the paper. The purpose of Section 3 is to give the proof of the theorem above in a more abstract setting. The starting point is a model category $``$ together with a weak factorization system $`(,)`$ satisfying some technical conditions which are fulfilled by the weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$ and by the set of generating T-homotopy equivalences. Several proofs of Section 3 are adaptations of standard proofs . But since the existence of a convenient model structure for $`(,)`$ is not supposed <sup>2</sup><sup>2</sup>2that is: a model structure such that $``$ is the class of trivial cofibrations., there are some subtle differences and also new phenomena. The idea of considering the path object construction comes from the reading of Kurz and Rosický’s paper . In this paper, Kurz and Rosický have the idea of considering a cylinder object construction with any weak factorization system $`(,)`$. This allows them to investigate the categorical localization of the underlying category with respect to the class of morphisms $``$ viewed, morally speaking, as a class of trivial fibrations. The dual situation is explored in this section, with an underlying category which is not only a category but also a model category. The situation described in Section 3 makes one think of the notion of fibration category in the sense of Baues . However, we do not know how to construct a fibration category from the results of Section 3. The path object functor constructed in Section 3 cannot satisfy the whole set of axioms of a P-category in the sense of Baues since the associated homotopy relation is not transitive. In particular, it does not even seem to satisfy the pullback axiom. Next, Section 4 proves the theorem above as an application of Section 3.
### Link with the series of papers “T-homotopy and refinement of observation”
This paper is independent from the series of papers “T-homotopy and refinement of observation” except for the proof of Theorem 4 at the very end of this work in which Theorem 5.2 is used. This paper was written while the author was trying to understand whether the (categorical) localization $`\mathrm{𝐅𝐥𝐨𝐰}[\mathrm{𝐜𝐨𝐟}(𝒯)^1]`$ of the category of flows with respect to the T-homotopy equivalences introduced in is locally small. Indeed, the local smallness is not established in the series of papers “T-homotopy and refinement of observation”. The result we obtain is more subtle. In the “correct” localization, all morphisms of $`\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)`$ with cofibrant domains are inverted. This is enough for future application in computer science since the real concrete examples are all of them modelled by cofibrant flows. But it is not known whether the other morphisms of $`\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)`$ are inverted. If this fact should be true, then it would probably be a consequence of the left properness of the weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$ (which is proved in Theorem 6.4).
## 2. Prerequisites and notations
The initial object (resp. the terminal object) of a category $`𝒞`$, if it exists, is denoted by $`\mathrm{}`$ (resp. $`\mathrm{𝟏}`$).
Let $`i:AB`$ and $`p:XY`$ be maps in a category $`𝒞`$. Then $`i`$ has the left lifting property (LLP) with respect to $`p`$ (or $`p`$ has the right lifting property (RLP) with respect to $`i`$) if for every commutative square
there exists a morphism $`g`$ called a lift making both triangles commutative.
Let $`𝒞`$ be a cocomplete category. If $`K`$ is a set of morphisms of $`𝒞`$, then the class of morphisms of $`𝒞`$ that satisfy the RLP (right lifting property) with respect to every morphism of $`K`$ is denoted by $`\mathrm{𝐢𝐧𝐣}(K)`$ and the class of morphisms of $`𝒞`$ that are transfinite compositions of pushouts of elements of $`K`$ is denoted by $`\mathrm{𝐜𝐞𝐥𝐥}(K)`$. Denote by $`\mathrm{𝐜𝐨𝐟}(K)`$ the class of morphisms of $`𝒞`$ that satisfy the LLP (left lifting property) with respect to every morphism of $`\mathrm{𝐢𝐧𝐣}(K)`$. The cocompleteness of $`𝒞`$ implies $`\mathrm{𝐜𝐞𝐥𝐥}(K)\mathrm{𝐜𝐨𝐟}(K)`$. Moreover, every morphism of $`\mathrm{𝐜𝐨𝐟}(K)`$ is a retract of a morphism of $`\mathrm{𝐜𝐞𝐥𝐥}(K)`$ as soon as the domains of $`K`$ are small relative to $`\mathrm{𝐜𝐞𝐥𝐥}(K)`$ ( Corollary 2.1.15). An element of $`\mathrm{𝐜𝐞𝐥𝐥}(K)`$ is called a relative $`K`$-cell complex. If $`X`$ is an object of $`𝒞`$, and if the canonical morphism $`\mathrm{}X`$ is a relative $`K`$-cell complex, one says that $`X`$ is a $`K`$-cell complex.
A congruence $``$ on a category $`𝒞`$ consists of an equivalence relation on the set $`𝒞(X,Y)`$ of morphisms from $`X`$ to $`Y`$ for every object $`X`$ and $`Y`$ of $`𝒞`$ such that if $`f,g𝒞(X,Y)`$, then $`fg`$ implies $`ufuf`$ and $`fvgv`$ for any morphism $`u`$ and $`v`$ as soon as $`uf`$ and $`fv`$ exist.
Let $`𝒞`$ be a cocomplete category with a distinguished set of morphisms $`I`$. Then let $`\mathrm{𝐜𝐞𝐥𝐥}(𝒞,I)`$ be the full subcategory of $`𝒞`$ consisting of the object $`X`$ of $`𝒞`$ such that the canonical morphism $`\mathrm{}X`$ is an object of $`\mathrm{𝐜𝐞𝐥𝐥}(I)`$. In other terms, $`\mathrm{𝐜𝐞𝐥𝐥}(𝒞,I)=(\mathrm{}𝒞)\mathrm{𝐜𝐞𝐥𝐥}(I)`$.
It is obviously impossible to read this paper without some familiarity with model categories. Possible references for model categories are , and . The original reference is but Quillen’s axiomatization is not used in this paper. The Hovey’s book axiomatization is preferred. If $``$ is a cofibrantly generated model category with set of generating cofibrations $`I`$, let $`\mathrm{𝐜𝐞𝐥𝐥}():=\mathrm{𝐜𝐞𝐥𝐥}(,I)`$. Any cofibrantly generated model structure $``$ comes with a cofibrant replacement functor $`Q:\mathrm{𝐜𝐞𝐥𝐥}()`$. For every morphism $`f`$ of $``$, the morphism $`Q(f)`$ is a cofibration, and even an inclusion of subcomplexes. A set $`K`$ of morphisms of a model category permits the small object argument if the domains of the morphisms of $`K`$ are small relative to $`\mathrm{𝐜𝐞𝐥𝐥}(K)`$. For such a set $`K`$, one can use the small object argument. The small object argument is recalled in the proof of Proposition 7.
In this paper, the notation $``$ means weak equivalence or equivalence of categories, and the notation $``$ means isomorphism.
A partially ordered set $`(P,)`$ (or poset) is a set equipped with a reflexive antisymmetric and transitive binary relation $``$. A poset $`(P,)`$ is bounded if there exist $`\widehat{0}P`$ and $`\widehat{1}P`$ such that $`P[\widehat{0},\widehat{1}]`$ and such that $`\widehat{0}\widehat{1}`$. Let $`\widehat{0}=\mathrm{min}P`$ (the bottom element) and $`\widehat{1}=\mathrm{max}P`$ (the top element).
Every poset $`P`$, and in particular every ordinal, can be viewed as a small category denoted in the same way: the objects are the elements of $`P`$ and there exists a morphism from $`x`$ to $`y`$ if and only if $`xy`$. If $`\lambda `$ is an ordinal, a $`\lambda `$-sequence (or a transfinite sequence) in a cocomplete category $`𝒞`$ is a colimit-preserving functor $`X`$ from $`\lambda `$ to $`𝒞`$. We denote by $`X_\lambda `$ the colimit $`\underset{}{\mathrm{lim}}X`$ and the morphism $`X_0X_\lambda `$ is called the transfinite composition of the $`X_\mu X_{\mu +1}`$.
If $`𝒞`$ is a locally small category, and if $`\mathrm{\Sigma }`$ is a class of morphisms of $`𝒞`$, then we denote by $`𝒞[\mathrm{\Sigma }^1]`$ the (categorical) localization of $`𝒞`$ with respect to $`\mathrm{\Sigma }`$ . The category $`𝒞[\mathrm{\Sigma }^1]`$ is not necessarily locally small. If $``$ is a model category with class of weak equivalences $`𝒲`$, then the localization $`[𝒲^1]`$ is locally small and it is called the Quillen homotopy category of $``$. It is denoted by $`\mathrm{𝐇𝐨}()`$.
## 3. Localizing a model category w.r.t. a weak factorization system
###### Definition 1.
Let $`𝒞`$ be a category. A weak factorization system is a pair $`(,)`$ of classes of morphisms of $`𝒞`$ such that the class $``$ is the class of morphisms having the LLP with respect to $``$, such that the class $``$ is the class of morphisms having the RLP with respect to $``$ and such that every morphism of $`𝒞`$ factors as a composite $`r\mathrm{}`$ with $`\mathrm{}`$ and $`r`$. The weak factorization system is functorial if the factorization $`r\mathrm{}`$ can be made functorial.
In a weak factorization system $`(,)`$, the class $``$ (resp. $``$) is completely determined by $``$ (resp. $``$).
###### Definition 2.
Let $`𝒞`$ be a cocomplete category. A weak factorization system $`(,)`$ is cofibrantly generated if there exists a set $`K`$ of morphisms of $`𝒞`$ permitting the small object argument such that $`=\mathrm{𝐜𝐨𝐟}(K)`$ and $`=\mathrm{𝐢𝐧𝐣}(K)`$.
A cofibrantly generated weak factorization system is necessarily functorial. Definition 2 appears in in the context of locally presentable category as the notion of small weak factorization system.
The data for this section are:
1. a complete and cocomplete category $``$ equipped with a model structure denoted by $`(\mathrm{Cof},\mathrm{Fib},𝒲)`$ for respectively the class of cofibrations, of fibrations and of weak equivalences such that the weak factorization system $`(\mathrm{Cof}𝒲,\mathrm{Fib})`$ is cofibrantly generated : the set of generating trivial cofibrations is denoted by $`J`$.
2. a cofibrantly generated weak factorization system $`(,)`$ on $``$ satisfying the following property: $`\mathrm{Cof}𝒲\mathrm{Cof}`$. So there exists a set of morphisms $`K`$ such that $`=\mathrm{𝐜𝐨𝐟}(JK)`$ and $`=\mathrm{𝐢𝐧𝐣}(JK)`$ and such that $`JK`$ permits the small object argument. Therefore every morphism $`f`$ factors as a composite $`f=\beta (f)\alpha (f)`$ where $`\alpha (f)\mathrm{𝐜𝐞𝐥𝐥}(JK)`$ and where $`\beta (f)`$. The functorial factorization is supposed to be obtained using the small object argument. It is fixed for the whole section.
###### Definition 3.
Let $`X`$ be an object of $``$. The path object of $`X`$ with respect to $``$ is the functorial factorization
of the diagonal morphism $`(\mathrm{Id}_X,\mathrm{Id}_X):XX\times X`$ by the morphism $`\alpha (\mathrm{Id}_X,\mathrm{Id}_X):X\mathrm{Path}_{}(X)`$ of $``$ composed with the morphism $`\beta (\mathrm{Id}_X,\mathrm{Id}_X):\mathrm{Path}_{}(X)X\times X`$ of $``$.
###### Notation 1.
Let $`_{cof}`$ be the full subcategory of cofibrant objects of $``$.
The path object of $`X`$ with respect to $``$ is cofibrant as soon as $`X`$ is cofibrant since the morphism $`\alpha (\mathrm{Id}_X,\mathrm{Id}_X):X\mathrm{Path}_{}(X)`$ is a cofibration. So the path object construction yields an endofunctor of $`_{cof}`$.
The reader must notice that we do not assume here that $`\alpha (\mathrm{Id}_X,\mathrm{Id}_X)`$ is a weak equivalence of any kind, contrary to the usual definition of a path object. As in for the construction of the cylinder functor, we do use the functorial factorization and we do suppose that $`\alpha (\mathrm{Id}_X,\mathrm{Id}_X)`$ belongs to $``$. A morphism of $``$ being an isomorphism of $`_{cof}[(𝒲)^1]`$, our condition is stronger than the usual one for the construction of a path object in a model category.
###### Definition 4.
An object $`X`$ of $``$ is fibrant with respect to $``$ if the unique morphism $`f_X:X\mathrm{𝟏}`$, where $`\mathrm{𝟏}`$ is the terminal object of $``$, is an element of $``$.
An object which is fibrant with respect to $`\mathrm{Cof}𝒲`$ is a fibrant object in the usual sense.
###### Notation 2.
Let $`^{f,}`$ be the full subcategory of $``$ of fibrant objects with respect to $``$. Let $`_{cof}^{f,}`$ be the full subcategory of $`_{cof}`$ of fibrant objects with respect to $``$.
If $`X`$ is fibrant with respect to $``$, the morphism $`X\times XX\times \mathrm{𝟏}X`$ belongs to $``$. Therefore the composite
$$\mathrm{Path}_{}(X)X\times XX\mathrm{𝟏}$$
belongs to $``$ as well. So the path object $`\mathrm{Path}_{}(X)`$ is also fibrant with respect to $``$. Thus, the path object construction yields endofunctors of $`^{f,}`$ and of $`_{cof}^{f,}`$.
If $`f:XY`$ is a morphism of $`_{cof}^{f,}`$, then the functorial factorization $`(\alpha ,\beta )`$ yields a composite a priori in $`_{cof}`$ (since $`\mathrm{Cof}`$)
equal to $`f`$. The unique morphism $`Z\mathrm{𝟏}`$ is equal to the composite $`ZY\mathrm{𝟏}`$ of two morphisms of $``$. Therefore $`Z`$ is fibrant with respect to $``$ and the functorial weak factorization system $`(,)`$ restricts to a functorial weak factorization system of $`_{cof}^{f,}`$ denoted in the same way.
###### Definition 5.
Let $`f,g:XY`$ be two morphisms of $``$. A right homotopy with respect to $``$ from $`f`$ to $`g`$ is a morphism $`H:X\mathrm{Path}_{}(Y)`$ such that
$$\beta (\mathrm{Id}_Y,\mathrm{Id}_Y)H=(f,g).$$
This situation is denoted by $`f_{}^rg`$.
Note the binary relation $`_{}^r`$ does not depend on the choice of the functorial factorization $`(\alpha ,\beta )`$. Indeed, with another functorial factorization $`(\alpha ^{},\beta ^{})`$, and the corresponding path object functor $`\mathrm{Path}_{}^{}`$, one can consider for every object $`Y`$ of $``$ the commutative diagram
The lift $`k`$ exists since the arrow $`Y\mathrm{Path}_{}(Y)`$ is in $``$ and since the arrow $`\mathrm{Path}_{}^{}(Y)Y\times Y`$ is in $``$.
The morphism $`\alpha (\mathrm{Id}_Y,\mathrm{Id}_Y)f:X\mathrm{Path}_{}Y`$ yields a right homotopy from $`f`$ to $`f`$ with respect to $``$. If $`H:X\mathrm{Path}_{}(Y)`$ is a right homotopy from $`f`$ to $`g`$ with respect to $``$, then the usual way for obtaining a right homotopy from $`g`$ to $`f`$ with respect to $``$ consists of considering the commutative diagram:
with $`\tau (y,y^{})=(y^{},y)`$. The existence of the lift $`k`$ comes from the definition of the path object and of the fact that $`(,)`$ is a weak factorization system. So the binary relation $`_{}^r`$ is reflexive and symmetric.
This relation is not transitive in general. The pair $`(^{op},^{op})`$ is a weak factorization system of the opposite category $`^{op}`$ (the model structure of $``$ is forgotten for this paragraph only). The path object becomes a cylinder object and the binary relation $`_{}^r`$ becomes the homotopy relation of . Example 3.6 gives an example where the homotopy is not transitive. Thus, the opposite category with the opposite weak factorization system gives an example where $`_{}^r`$ is not transitive.
###### Notation 3.
Let us denote by $`_{}`$ the transitive closure of the binary relation $`_{}^r`$.
###### Proposition 1.
Let $`X`$ be an object of $`_{cof}`$. Let $`Y`$ be an object of $``$. Let $`f,g:XY`$ be two morphisms between them. Then $`f_{\mathrm{Cof}𝒲}g`$ if and only if $`f`$ and $`g`$ are right homotopic in the usual sense of model categories.
Notice that it is crucial in the proof for $`X`$ to be cofibrant.
###### Proof.
Indeed, two morphisms $`f,g:XY`$ with $`X`$ cofibrant are right homotopic in the usual sense if the pair $`(f,g)`$ is in the transitive closure of the following situation denoted by $`f^rg`$ (cf. p7):
1. Decompose the diagonal morphism $`(\mathrm{Id}_Y,\mathrm{Id}_Y):YY\times Y`$ into a weak equivalence $`YPY`$ of $``$ followed by a fibration $`(p_1,p_2):PYY\times Y`$ of $``$.
2. There exists $`H:XPY`$ such that $`(p_1,p_2)H=(f,g)`$.
Let us factor the weak equivalence $`YPY`$ as a composite $`YP^{}YPY`$ where $`YP^{}Y`$ is a trivial cofibration and where $`P^{}YPY`$ is a trivial fibration. Then one can lift the right homotopy $`H:XPY`$ to a morphism $`\overline{H}:XP^{}Y`$ since $`X`$ is cofibrant. But $`\overline{H}`$ is not yet a right homotopy from $`f`$ to $`g`$ with respect to $`\mathrm{Cof}𝒲`$ since $`P^{}Y`$ is not necessarily the functorial path object $`\mathrm{Path}_{\mathrm{Cof}𝒲}(Y)`$ ! Let us consider the commutative diagram
Since the arrow $`YP^{}Y`$ is a trivial cofibration and since the arrow $`\mathrm{Path}_{\mathrm{Cof}𝒲}(Y)Y\times Y`$ is a fibration, there exists a lift $`k`$. Then $`k\overline{H}`$ is a right homotopy with respect to $`\mathrm{Cof}𝒲`$ from $`f`$ to $`g`$.
Conversely, the path object with respect to $`\mathrm{Cof}𝒲`$ is a path object in the above sense of model categories. So a right homotopy from $`f`$ to $`g`$ with respect to $`\mathrm{Cof}𝒲`$ is a right homotopy in the usual sense of model categories. ∎
The following proposition gives a sufficient condition for the binary relation $`_{}^r`$ to be transitive.
###### Proposition 2.
Let us suppose that there exists a model structure $`(\mathrm{Cof}_{},\mathrm{Fib}_{},𝒲_{})`$ on $``$ such that $`=\mathrm{Cof}_{}𝒲_{}`$ and such that every cofibrant object of $``$ is a cofibrant object of $`(\mathrm{Cof}_{},\mathrm{Fib}_{},𝒲_{})`$. Let $`X`$ and $`Y`$ be two objects of $`_{cof}^{f,}`$. Then the binary relation $`_{}^r`$ is an equivalence relation on $`_{cof}^{f,}(X,Y)`$.
Notice that we do not need suppose in the proof of Proposition 1 that the weak factorization system $`(\mathrm{Cof}𝒲,\mathrm{Fib})`$ is cofibrantly generated. So we do not need this hypothesis in the proof of Proposition 2.
###### Proof.
By Proposition 1 applied to the model structure $`(\mathrm{Cof}_{},\mathrm{Fib}_{},𝒲_{})`$, the binary relation $`_{}^r`$ coincides with right homotopy for the model structure $`(\mathrm{Cof}_{},\mathrm{Fib}_{},𝒲_{})`$. Since $`Y`$ is fibrant for the latter model structure, one deduces that $`_{}^r`$ is transitive by Proposition 1.2.5. ∎
###### Corollary 1.
If $``$ is the class of trivial cofibrations of a left Bousfield localization of the model structure of $``$, then the binary relation $`_{}^r`$ on the set of morphisms $`(X,Y)`$ with $`X_{cof}`$ and with $`Y`$ fibrant with respect to $``$ is an equivalence relation.
###### Proposition 3.
(dual to Lemma 3.2) Let $`f,g:XY`$ be two morphisms of $``$. Let $`u:YU`$ and $`v:VX`$ be two other morphisms of $``$. If $`f_{}^rg`$, then $`uf_{}^rug`$ and $`fv_{}^rgv`$.
In other terms, the equivalence relation $`_{}`$ defines a congruence in the sense of .
###### Proof.
By considering the opposite of the category $``$, the proof is complete using Lemma 3.2. ∎
The proof of Proposition 3 does use the factorization $`(\alpha ,\beta )`$ and its functoriality. We could avoid using the functoriality since the morphism $`\alpha (\mathrm{Id}_Y,\mathrm{Id}_Y):Y\mathrm{Path}_{}(Y)`$ belongs to $``$ and since the morphism $`\beta (\mathrm{Id}_U\times \mathrm{Id}_U):\mathrm{Path}_{}(U)U\times U`$ belongs to $``$. But anyway, the proof of Proposition 3 cannot be adapted to the usual notion of right homotopy. This is once again a difference between our notion of right homotopy and the usual one on model category.
Proposition 3 allows to consider the quotients $`/_{}`$ (resp. $`_{cof}/_{}`$, $`^{f,}/_{}`$, $`_{cof}^{f,}/_{}`$) of the category $``$ (resp. $`_{cof}`$, $`^{f,}`$ $`_{cof}^{f,}`$) by the congruence $`_{}`$. By definition, the objects of $`/_{}`$ (resp. $`_{cof}/_{}`$, $`^{f,}/_{}`$, $`_{cof}^{f,}/_{}`$) are the objects of $``$ (resp. $`_{cof}`$, $`^{f,}`$, $`_{cof}^{f,}`$), and for any object $`X`$ and $`Y`$ of $``$ (resp. $`_{cof}`$, $`^{f,}`$, $`_{cof}^{f,}`$), one has $`/_{}(X,Y)=(X,Y)/_{}`$ (resp. $`_{cof}/_{}(X,Y)=_{cof}(X,Y)/_{}`$, $`^{f,}/_{}(X,Y)=^{f,}(X,Y)/_{}`$, $`_{cof}^{f,}/_{}(X,Y)=_{cof}^{f,}(X,Y)/_{}`$). Let
$`[]_{}:/_{}`$
$`[]_{}:_{cof}_{cof}/_{}`$
$`[]_{}:^{f,}^{f,}/_{}`$
$`[]_{}:_{cof}^{f,}_{cof}^{f,}/_{}`$
be the canonical functors.
###### Proposition 4.
(dual to Lemma 3.7) Let $`f:XY`$ be a morphism of $``$ belonging to $``$. Let us suppose that $`X`$ is fibrant with respect to $``$. Then there exists $`g:YX`$ such that $`fg_{}\mathrm{Id}_Y`$ and $`gf=\mathrm{Id}_X`$.
###### Proof.
Let us consider the commutative diagram of $``$
Since the left vertical arrow is in $``$ and since the right vertical arrow is in $``$ by hypothesis, there exists a lift $`g:YX`$. In other terms, $`gf=\mathrm{Id}_X`$. The diagram of $``$
is commutative since
$$\beta (\mathrm{Id}_Y,\mathrm{Id}_Y)\alpha (\mathrm{Id}_Y,\mathrm{Id}_Y)f=(\mathrm{Id}_Y,\mathrm{Id}_Y)f=(f,f)=(fg,\mathrm{Id}_Y)f.$$
Since $`f`$ by hypothesis and since $`\beta (\mathrm{Id}_Y,\mathrm{Id}_Y)`$, there exists $`H:Y\mathrm{Path}_{}(Y)`$ preserving the diagram above commutative. The morphism $`H`$ is by construction a right homotopy from $`fg`$ to $`\mathrm{Id}_Y`$ with respect to $``$. ∎
###### Proposition 5.
(almost dual to Theorem 3.9) One has the isomorphism of categories $`_{cof}^{f,}/_{}_{cof}^{f,}[^1]`$. In particular, this means that the category $`_{cof}^{f,}[^1]`$ is locally small.
The proof of Proposition 5 also shows the isomorphism of categories
$$^{f,}/_{}^{f,}[^1].$$
###### Proof.
We know that the pair $`(,)`$ restricts to a weak factorization system of $`_{cof}^{f,}`$. By considering the opposite category, the proposition is then a consequence of Theorem 3.9. ∎
###### Proposition 6.
(Detecting weak equivalences) A morphism $`f:AB`$ of $`_{cof}^{f,}`$ is an isomorphism of $`_{cof}^{f,}[^1]`$ if and only if for every object $`X`$ of $`_{cof}^{f,}`$, the map
$$(B,X)/_{}(A,X)/_{}$$
is bijective.
Note the “opposite” characterization $`(X,A)/_{}(X,B)/_{}`$ also holds. The statement of the theorem is chosen for having a characterization as close as possible to the characterization of weak equivalences in a left Bousfield localization.
###### Proof.
The condition means that the map
$$(_{cof}^{f,}/_{})(B,X)(_{cof}^{f,}/_{})(A,X)$$
is a bijection. By Yoneda’s lemma applied within the locally small category $`_{cof}^{f,}/_{}`$, the condition is equivalent to saying that $`f:AB`$ is an isomorphism of $`_{cof}^{f,}/_{}`$. By Proposition 5, the condition is equivalent to saying that $`f:AB`$ is an isomorphism of $`_{cof}^{f,}[^1]`$. ∎
###### Definition 6.
Let $`X`$ be an object of $`_{cof}`$. The fibrant replacement of $`X`$ with respect to $``$ is the functorial factorization
of the unique morphism $`f_X:X\mathrm{𝟏}`$.
The mapping $`XR_{}(X)`$ is functorial and yields a functor from $`_{cof}`$ to $`_{cof}^{f,}`$ since the morphism $`\alpha (f_X):XR_{}(X)`$ is a cofibration.
###### Lemma 1.
Let $`\lambda `$ be a limit ordinal. Let $`X:\lambda `$ and $`Y:\lambda `$ be two transfinite sequences. Let $`f:XY`$ be a morphism of transfinite sequences such that for any $`\mu <\lambda `$, $`f_\mu :X_\mu Y_\mu `$ belongs to $``$. Then $`f_\lambda :X_\lambda Y_\lambda `$ belongs to $``$. Moreover, if for any $`\mu <\lambda `$, $`f_\mu :X_\mu Y_\mu `$ belongs to $`\mathrm{𝐜𝐞𝐥𝐥}(JK)`$, then $`f_\lambda :X_\lambda Y_\lambda `$ belongs to $`\mathrm{𝐜𝐞𝐥𝐥}(JK)`$ as well.
Lemma 1 and Proposition 7 are very close to Proposition 12.4.7. The difference is that we do not suppose here that the underlying model category is cellular.
###### Proof.
Let $`T_0=X_\lambda `$. Let us consider the unique transfinite sequence $`T:\lambda `$ such that one has the pushout diagram
where the left vertical arrow is the composite $`X_\mu X_\lambda T_\mu `$ for any $`\mu <\lambda `$. Let $`Z`$ be an object of $``$ and let $`\varphi :Y_\lambda Z`$ be a morphism of $``$. The composite $`X_\lambda Y_\lambda Z`$ together with the composite $`Y_0Y_\lambda Z`$ yields with the pushout diagram above for $`\mu =0`$ a morphism $`T_1Z`$ since $`f`$ is a morphism of transfinite sequences. And by an immediate transfinite induction, one obtains a morphism $`\underset{}{\mathrm{lim}}_\mu T_\mu Z`$. So one has the isomorphism $`\underset{}{\mathrm{lim}}_\mu T_\mu Y_\lambda `$ since the two objects of $``$ satisfy the same universal property. Hence the result since the class of morphisms $``$ and $`\mathrm{𝐜𝐞𝐥𝐥}(JK)`$ are both closed under transfinite composition. ∎
###### Proposition 7.
One has $`R_{}()`$, and even $`R_{}(\mathrm{𝐜𝐞𝐥𝐥}(JK))\mathrm{𝐜𝐞𝐥𝐥}(JK)`$.
###### Proof.
A morphism $`f\mathrm{𝐜𝐨𝐟}(JK)`$ is a retract of a morphism $`g\mathrm{𝐜𝐞𝐥𝐥}(JK)`$ since $`JK`$ permits the small object argument. And the morphism $`R_{\mathrm{𝐜𝐨𝐟}(JK)}(f)`$ is then a retract of the morphism $`R_{\mathrm{𝐜𝐨𝐟}(JK)}(g)`$. Therefore it suffices to prove that $`f\mathrm{𝐜𝐞𝐥𝐥}(JK)`$ implies $`R_{\mathrm{𝐜𝐨𝐟}(JK)}(f)\mathrm{𝐜𝐞𝐥𝐥}(JK)`$. The functor $`R_{\mathrm{𝐜𝐨𝐟}(JK)}`$ is obtained by a transfinite construction involving the small object argument. Let $`X_0=X`$ and $`Y_0=Y`$ and $`f=f_0`$. For any ordinal $`\lambda `$, let $`\overline{Y_\lambda }`$ be the object of $``$ defined by the following commutative diagram:
Let us suppose $`f_\lambda :X_\lambda Y_\lambda `$ constructed for some $`\lambda 0`$ and let us suppose that the morphism $`\overline{Y_\lambda }Y_\lambda `$ is an element of $`\mathrm{𝐜𝐞𝐥𝐥}(JK)`$. The small object argument consists of considering the sets of commutative squares $`\{kf_{X_\lambda },kJK\}`$ and $`\{kf_{Y_\lambda },kJK\}`$ where $`f_{X_\lambda }:X_\lambda \mathrm{𝟏}`$ and $`f_{Y_\lambda }:Y_\lambda \mathrm{𝟏}`$ are the canonical morphisms from respectively $`X_\lambda `$ and $`Y_\lambda `$ to the terminal object of $``$. The morphism $`f_\lambda `$ allows the identification of $`\{kf_{X_\lambda },kJK\}`$ with a subset of $`\{kf_{Y_\lambda },kJK\}`$. And the morphism $`f_{\lambda +1}:X_{\lambda +1}Y_{\lambda +1}`$ is obtained by the diagram (where the notations $`\mathrm{dom}(k)`$ and $`\mathrm{codom}(k)`$ mean respectively domain and codomain of $`k`$):
Therefore $`f_{\lambda +1}:X_{\lambda +1}Y_{\lambda +1}`$ is an element of $`\mathrm{𝐜𝐞𝐥𝐥}(JK)`$. The proof is complete with Lemma 1. ∎
Note the same kind of argument as the one of Proposition 7 leads to the following proposition (worth being noticed, but useless for the sequel):
###### Proposition 8.
One has $`\mathrm{Path}_{}()`$, and even $`\mathrm{Path}_{}(\mathrm{𝐜𝐞𝐥𝐥}(JK))\mathrm{𝐜𝐞𝐥𝐥}(JK)`$.
Proposition 7 will be used in particular in the proof of Proposition 11 with the functorial weak factorization system $`(\mathrm{Cof}𝒲,\mathrm{Fib})`$ and in the proof of Proposition 9.
###### Proposition 9.
The inclusion functor $`_{cof}^{f,}_{cof}`$ induces an equivalence of categories $`_{cof}^{f,}[^1]_{cof}[^1]`$. In particular, this implies that the category $`_{cof}[^1]`$ is locally small.
###### Proof.
Since $`R_{}()`$ by Proposition 7, there exists a unique functor $`L(R_{})`$ making the following diagram commutative:
If $`i:_{cof}^{f,}_{cof}`$ is the inclusion functor, then there exists a unique functor $`L(i)`$ making the following diagram commutative:
There are two natural transformations $`\mu :\mathrm{Id}_{_{cof}}iR_{}`$ and $`\nu :\mathrm{Id}_{_{cof}^{f,}}R_{}i`$ such that for any $`X_{cof}`$ and any $`Y_{cof}^{f,}`$, the morphisms $`\mu (X)`$ and $`\nu (Y)`$ belong to $``$. So at the level of localizations, one obtains the isomorphisms of functors $`\mathrm{Id}_{_{cof}[^1]}L(i)L(R_{})`$ and $`\mathrm{Id}_{_{cof}^{f,}[^1]}L(R_{})L(i)`$. Hence the result. ∎
###### Proposition 10.
Let $`(^{},^{})`$ be another cofibrantly generated weak factorization of $``$ such that $`\mathrm{Cof}𝒲^{}\mathrm{Cof}`$. Let us suppose that $`^{}`$. Then the localization functor $`_{cof}_{cof}[^1]`$ factors uniquely as a composite
$$_{cof}_{cof}[^1]_{cof}[^1].$$
###### Proof.
One has $`^{}`$. ∎
###### Proposition 11.
The localization functor $`L:_{cof}_{cof}[^1]`$ sends the weak equivalences of $``$ between cofibrant objects to isomorphisms of $`_{cof}[^1]`$.
###### Proof.
The localization functor $`L:_{cof}_{cof}[^1]`$ factors uniquely as a composite
$$_{cof}_{cof}[(\mathrm{Cof}𝒲)^1]_{cof}[^1]$$
by Proposition 10 and since $`\mathrm{Cof}𝒲`$. By Proposition 9 and Proposition 7 applied to $`=\mathrm{Cof}𝒲`$, one has the equivalence of categories
$$_{cof}[(\mathrm{Cof}𝒲)^1]_{cof}^{f,\mathrm{Cof}𝒲}[(\mathrm{Cof}𝒲)^1].$$
By Proposition 5 applied to $`=\mathrm{Cof}𝒲`$, one has the isomorphism of categories
$$_{cof}^{f,\mathrm{Cof}𝒲}[(\mathrm{Cof}𝒲)^1]_{cof}^{f,\mathrm{Cof}𝒲}/_{\mathrm{Cof}𝒲}.$$
Therefore one obtains the equivalence of categories
$$_{cof}[(\mathrm{Cof}𝒲)^1]_{cof}^{f,\mathrm{Cof}𝒲}/_{\mathrm{Cof}𝒲}.$$
The category $`_{cof}^{f,\mathrm{Cof}𝒲}`$ is the full subcategory of cofibrant-fibrant objects of $``$. Since right homotopy with respect to $`\mathrm{Cof}𝒲`$ corresponds to the usual notion of right homotopy by Proposition 1, one has the equivalence of categories
$$_{cof}^{f,\mathrm{Cof}𝒲}/_{\mathrm{Cof}𝒲}\mathrm{𝐇𝐨}()$$
where $`\mathrm{𝐇𝐨}()=[𝒲^1]`$. Hence the result. ∎
###### Proposition 12.
The categories $`_{cof}[^1]`$ and $`_{cof}[(𝒲)^1]`$ are isomorphic. In particular, this implies that the category $`_{cof}[(𝒲)^1]`$ is locally small.
###### Proof.
By Proposition 11, there exists a unique functor
$$_{cof}[(𝒲)^1]_{cof}[^1]$$
such that the following diagram is commutative:
Since $`𝒲`$, there exists a unique functor $`_{cof}[^1]_{cof}[(𝒲)^1]`$ such that the following diagram is commutative:
Hence the result. ∎
In the same way, one can prove the
###### Proposition 13.
The categories $`_{cof}^{f,}[^1]`$ and $`_{cof}^{f,}[(𝒲)^1]`$ are isomorphic. In particular, this implies that the category $`_{cof}^{f,}[(𝒲)^1]`$ is locally small.
###### Notation 4.
Let
$$𝒲_{}=\{f:XY,Z_{cof}^{f,},(R_{}(Q(Y)),Z)/_{}(R_{}(Q(X)),Z)/_{}\}.$$
###### Proposition 14.
The inclusion functor $`_{cof}`$ induces an equivalence of categories $`_{cof}[(𝒲)^1][𝒲_{}^1]`$. In particular, this implies that the category $`[𝒲_{}^1]`$ is locally small.
###### Proof.
Let us consider the composite
where $`Q`$ is the cofibrant replacement functor of $``$. By definition of $`𝒲_{}`$ and by Proposition 6, by Proposition 9 and by Proposition 12, the functor $`LQ`$ sends the morphisms of $`𝒲_{}`$ to isomorphisms. Thus, there exists a unique functor $`L(Q)`$ making the following diagram commutative:
Let $`i:_{cof}`$ be the inclusion functor. Let $`f:XY𝒲`$ be a morphism of $`_{cof}`$. Then $`Q(f):Q(X)Q(Y)`$ is still invertible in $`_{cof}[(𝒲)^1]`$ since the morphism $`Q(X)X`$ and $`Q(Y)Y`$ are both weak equivalences of $``$ between cofibrant objects. So by Proposition 9, the morphism $`R_{}(Q(f)):R_{}(Q(X))R_{}(Q(Y))`$ is invertible in $`_{cof}^{f,}[(𝒲)^1]_{cof}^{f,}[^1]`$. So by Proposition 6, one deduces that $`f𝒲_{}`$. Thus, there exists a unique functor $`L(i)`$ making the following diagram commutative:
There exist two natural transformations $`\mu :Qi\mathrm{Id}_{_{cof}}`$ and $`\nu :iQ\mathrm{Id}_{}`$. If $`X`$ is an object of $`_{cof}`$, then $`\mu (X)`$ is a trivial fibration between cofibrant objects, i.e. $`\mu (X)𝒲`$. So one deduces that $`L(\mu (X))`$ is an isomorphism of $`_{cof}[(𝒲)^1]`$. Therefore one obtains the isomorphism of functors $`L(Q)L(i)\mathrm{Id}_{_{cof}[(𝒲)^1]}`$. If $`Y`$ is an object of $``$, then $`\nu (Y):Q(Y)Y`$ is a trivial fibration of $``$. Thus, $`Q(\nu (Y)):Q(Q(Y))Q(Y)`$ is a trivial cofibration of $``$ between cofibrant objects. So $`Q(\nu (Y))\mathrm{Cof}𝒲`$. Since $`R_{}()`$, one deduces that $`R_{}(Q(\nu (Y)))`$ is an isomorphism of $`_{cof}^{f,}[^1]`$. Again by Proposition 6, one deduces that $`\nu (Y)𝒲_{}`$ and one obtains the isomorphism of functors $`L(i)L(Q)\mathrm{Id}_{[𝒲_{}^1]}`$. The proof is complete. ∎
###### Theorem 2.
(Whitehead’s theorem for the localization of a model category with respect to a weak factorization system) The inclusion functors $`_{cof}^{f,}_{cof}`$ induce the equivalences of categories
$$_{cof}^{f,}/_{}_{cof}[(𝒲)^1][𝒲_{}^1].$$
The functor $`[𝒲_{}^1]_{cof}^{f,}/_{}`$ is given by the cofibrant-fibrant w.r.t. $``$ functor $`R_{}Q:_{cof}^{f,}`$. Moreover, the localization functor $`[𝒲_{}^1]`$ factors uniquely as a composite
$$[𝒲^1][𝒲_{}^1].$$
The equivalence of categories $`_{cof}[(𝒲)^1][𝒲_{}^1]`$ shows that up to weak equivalence and up to the 2-out-of-3 axiom, a morphism of $`𝒲_{}`$ is a morphism of $`𝒲`$ between cofibrant objects of $``$. This means that the class of morphisms $`𝒲_{}`$ is not too big.
###### Proof.
The equivalence of categories $`[𝒲_{}^1]_{cof}[(𝒲)^1]`$ is given by Proposition 14. By Proposition 14, the functor $`[𝒲_{}^1]_{cof}[(𝒲)^1]`$ is induced by the cofibrant replacement functor $`Q:_{cof}`$. Therefore every morphism of $`𝒲`$ is in $`𝒲_{}`$. At last, one has
$`_{cof}[(𝒲)^1]`$
$`_{cof}[^1]`$ by Proposition 12
$`_{cof}^{f,}[^1]`$ by Proposition 9
$`_{cof}^{f,}/_{}`$ by Proposition 5.
Theorem 2 says that the category $`[𝒲_{}^1]`$ inverts all weak equivalences of $``$ and all morphisms of $``$ with cofibrant domains. We do not know if all morphisms of $``$ (and not only the ones with cofibrant domain) are inverted in the category $`[𝒲_{}^1]`$. But there is a kind of reciprocal statement:
###### Proposition 15.
Let us suppose that $``$ is left proper. Let $`f`$ be a cofibration of $``$ such that $`Q(f)`$. Then $`f`$.
###### Proof.
Let $`p`$. Since $`\mathrm{Cof}𝒲`$, the morphism $`p`$ is a fibration of $``$. By hypothesis, $`p`$ satisfies the RLP with respect to $`Q(f)`$. Since $`f`$ is a cofibration and by Proposition 13.2.1, one deduces that $`p`$ satisfies the RLP with respect to $`f`$. So $`f`$. ∎
Before treating the case of T-homotopy equivalences in Section 4, let us give some examples of the situation explored in this section.
### Example 1
Let $``$ be a cofibrantly generated model category with set of generating cofibrations $`I`$ and with set of generating trivial cofibrations $`J`$ and with class of weak equivalences $`𝒲`$. Let $`(,)=(\mathrm{𝐜𝐨𝐟}(J),\mathrm{𝐢𝐧𝐣}(J))`$. The main theorem gives the equivalences of categories
$$_{cof}^{fib}/_{\mathrm{𝐜𝐨𝐟}(J)}_{cof}[\mathrm{𝐜𝐨𝐟}(J)^1][𝒲_{\mathrm{𝐜𝐨𝐟}(J)}^1]$$
where $`_{cof}^{fib}`$ is the full subcategory of cofibrant-fibrant objects. One can directly check that $`𝒲_{\mathrm{𝐜𝐨𝐟}(J)}^1=𝒲`$. This is not surprising since the category $`_{cof}^{fib}/_{\mathrm{𝐜𝐨𝐟}(J)}`$ is the category of cofibrant-fibrant objects of $``$ up to homotopy.
### Example 2
Let $``$ be a cofibrantly generated model category with set of generating cofibrations $`I`$ and with set of generating trivial cofibrations $`J`$ and with class of weak equivalences $`𝒲`$. Then one can consider the model structure $`(\underset{¯}{\mathrm{All}},\underset{¯}{\mathrm{All}},\underset{¯}{\mathrm{Iso}})`$ where all morphisms are a cofibration and a fibration and where the weak equivalences are the isomorphisms. Indeed, one has $`(\underset{¯}{\mathrm{Iso}},\underset{¯}{\mathrm{All}})=(\mathrm{𝐜𝐨𝐟}(\mathrm{Id}_{\mathrm{}}),\mathrm{𝐢𝐧𝐣}(\mathrm{Id}_{\mathrm{}}))`$. The main theorem applied with the latter model structure and with the weak factorization system $`(,)=(\mathrm{𝐜𝐨𝐟}(J),\mathrm{𝐢𝐧𝐣}(J))`$ gives the equivalences of categories
$$^{fib}/_{\mathrm{𝐜𝐨𝐟}(J)}[\mathrm{𝐜𝐨𝐟}(J)^1][\underset{¯}{\mathrm{Iso}}_{\mathrm{𝐜𝐨𝐟}(J)}^1]$$
where $`^{fib}`$ is the full subcategory of fibrant objects, where $`_{\mathrm{𝐜𝐨𝐟}(J)}`$ is a congruence on the morphisms of the full subcategory of fibrant objects of $``$. The functor $`R_{\mathrm{𝐜𝐨𝐟}(J)}`$ is a fibrant replacement functor of $``$ and the functor $`Q`$ is a cofibrant replacement functor of the model structure $`(\underset{¯}{\mathrm{All}},\underset{¯}{\mathrm{All}},\underset{¯}{\mathrm{Iso}})`$. That is one can suppose that $`Q=\mathrm{Id}_{}`$.
### Example 3
Let $``$ be a combinatorial model category (in the sense of Jeff Smith), that is a cofibrantly generated model category such that the underlying category is locally presentable . Let $`J`$ be the set of generating trivial cofibrations. Then for any set $`K`$ of morphisms of $``$, the pair $`(\mathrm{𝐜𝐨𝐟}(JK),\mathrm{𝐢𝐧𝐣}(JK))`$ is a cofibrantly generated weak factorization system by Proposition 1.3. Then the main theorem of this section applies. One obtains the equivalences of categories
$$_{cof}^{f,\mathrm{𝐜𝐨𝐟}(JK)}/_{\mathrm{𝐜𝐨𝐟}(JK)}_{cof}[(𝒲\mathrm{𝐜𝐨𝐟}(JK))^1][𝒲_{\mathrm{𝐜𝐨𝐟}(JK)}^1].$$
Assume Vopěnka’s principle ( chapter 6). If $``$ is left proper, then the left Bousfield localization $`L_K`$ of the model category $``$ with respect to the set of morphisms $`K`$ exists by a theorem of Jeff Smith proved in Theorem 1.7 and in Theorem 2.2. The category $`[𝒲_{\mathrm{𝐜𝐨𝐟}(JK)}^1]`$ is not necessarily equivalent to the Quillen homotopy category $`\mathrm{𝐇𝐨}(L_K)`$ of this Bousfield localization. But all morphisms inverted by $`[𝒲_{\mathrm{𝐜𝐨𝐟}(JK)}^1]`$ are inverted by the Bousfield localization. In other terms, the functor $`\mathrm{𝐇𝐨}(L_K)`$ factors uniquely as a composite $`[𝒲_{\mathrm{𝐜𝐨𝐟}(JK)}^1]\mathrm{𝐇𝐨}(L_K)`$.
### Example 4
Let $``$ be a model category with model structure denoted by $`(\mathrm{Cof},\mathrm{Fib},𝒲)`$ for respectively the class of cofibrations, of fibrations and of weak equivalences such that the weak factorization system $`(\mathrm{Cof}𝒲,\mathrm{Fib})`$ is cofibrantly generated : the set of generating trivial cofibrations is denoted by $`J`$. Let $`(,)`$ be a cofibrantly generated weak factorization system such that $`\mathrm{Cof}𝒲\mathrm{Cof}`$. Let us suppose that the left Bousfield localization $`L_{}`$ of $``$ with respect to $``$ exists and let us suppose that $``$ is the class of trivial cofibrations of this Bousfield localization. One obtains the equivalences of categories
$$_{cof}^{f,}/_{}_{cof}[(𝒲)^1][𝒲_{}^1].$$
The category $`_{cof}^{f,}`$ is the full subcategory of $``$ containing the cofibrant-fibrant object of $`L_{}`$. The congruence $`_{}`$ is the usual notion of homotopy in $``$ ( Proposition 3.5.3). Then the category $`_{cof}^{f,}/_{}`$ is equivalent to the Quillen homotopy category $`\mathrm{𝐇𝐨}(L_{})`$.
## 4. Application : homotopy continuous flow and Whitehead’s theorem
The category $`\mathrm{𝐓𝐨𝐩}`$ of compactly generated topological spaces (i.e. of weak Hausdorff $`k`$-spaces) is complete, cocomplete and cartesian closed (more details for this kind of topological spaces in , the appendix of and also the preliminaries of ). For the sequel, all topological spaces will be supposed to be compactly generated. A compact space is always Hausdorff. The category $`\mathrm{𝐓𝐨𝐩}`$ is equipped with the unique model structure having the weak homotopy equivalences as weak equivalences and having the Serre fibrations <sup>3</sup><sup>3</sup>3that is a continuous map having the RLP with respect to the inclusion $`𝐃^n\times 0𝐃^n\times [0,1]`$ for all $`n0`$ where $`𝐃^n`$ is the $`n`$-dimensional disk. as fibrations .
As already described in the introduction, the time flow of a higher dimensional automaton is encoded in an object called a flow. The category $`\mathrm{𝐅𝐥𝐨𝐰}`$ is equipped with the unique model structure such that :
* The weak equivalences are the weak S-homotopy equivalences, i.e. the morphisms of flows $`f:XY`$ such that $`f^0:X^0Y^0`$ is a bijection and such that $`f:XY`$ is a weak homotopy equivalence.
* The fibrations are the morphisms of flows $`f:XY`$ such that $`f:XY`$ is a Serre fibration.
This model structure is cofibrantly generated. The set of generating cofibrations is the set $`I_+^{gl}=I^{gl}\{R,C\}`$ with
$$I^{gl}=\{\mathrm{Glob}(𝐒^{n1})\mathrm{Glob}(𝐃^n),n0\}$$
where $`𝐃^n`$ is the $`n`$-dimensional disk, where $`𝐒^{n1}`$ is the $`(n1)`$-dimensional sphere, where $`R`$ and $`C`$ are the set maps $`R:\{0,1\}\{0\}`$ and $`C:\mathrm{}\{0\}`$ and where for any topological space $`Z`$, the flow $`\mathrm{Glob}(Z)`$ is the flow defined by $`\mathrm{Glob}(Z)^0=\{\widehat{0},\widehat{1}\}`$, $`\mathrm{Glob}(Z)=Z`$, $`s=\widehat{0}`$ and $`t=\widehat{1}`$, and a trivial composition law. The set of generating trivial cofibrations is
$$J^{gl}=\{\mathrm{Glob}(𝐃^n\times \{0\})\mathrm{Glob}(𝐃^n\times [0,1]),n0\}.$$
The weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$ has some similarity with the model structure on the category of small simplicial categories (with identities !) constructed in . The weak equivalences (resp. the fibrations) of the latter look like the weak equivalences (resp. the fibrations) of the model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$ with an additional condition. The weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$ has also some similarity with the model structure on the category of small simplicial categories (with identities again !) on a fixed set of objets $`O`$ constructed in . For the latter, the set maps $`R:\{0,1\}\{0\}`$ and $`C:\mathrm{}\{0\}`$ are not used since the set of objects is fixed.
###### Definition 7.
A flow $`X`$ is loopless if for any $`\alpha X^0`$, the space $`_{\alpha ,\alpha }X`$ is empty.
Recall that a flow is a small category without identities morphisms enriched over a category of topological spaces. So the preceding definition is meaningful.
A poset $`(P,)`$ can be identified with a loopless flow having $`P`$ as set of states and such that there exists a non-constant execution path from $`x`$ to $`y`$ if and only if $`x<y`$. The corresponding flow is still denoted by $`P`$. This defines a functor from the full subcategory of posets whose morphisms are the strictly increasing maps to the full subcategory of loopless flows. The category of finite bounded posets is essentially small. Let us choose a small subcategory of representatives.
###### Definition 8.
Let $`𝒯`$ be the set of cofibrations $`Q(f):Q(P_1)Q(P_2)`$ such that $`f:P_1P_2`$ is a morphism of posets satisfying the following conditions:
1. The posets $`P_1`$ and $`P_2`$ are finite and bounded.
2. The morphism of posets $`f:P_1P_2`$ is one-to-one; in particular, if $`x`$ and $`y`$ are two elements of $`P_1`$ with $`x<y`$, then $`f(x)<f(y)`$.
3. One has $`f(\mathrm{min}P_1)=\mathrm{min}P_2=\widehat{0}`$ and $`f(\mathrm{max}P_1)=\mathrm{max}P_2=\widehat{1}`$.
4. The posets $`P_1`$ and $`P_2`$ are objects of the chosen small subcategory of representatives of the category of finite bounded posets.
The set $`𝒯`$ is called the set of generating T-homotopy equivalences.
The set $`𝒯`$ is introduced in for modelling T-homotopy as a refinement of observation. By now, this is the best known definition of T-homotopy.
###### Definition 9.
A flow $`X`$ is homotopy continuous if the unique morphism of flows $`f_X:X\mathrm{𝟏}`$ belongs to $`\mathrm{𝐢𝐧𝐣}(J^{gl}𝒯)`$.
Notice that $`\mathrm{𝐢𝐧𝐣}(J^{gl}𝒯)=\mathrm{𝐢𝐧𝐣}(𝒯)`$ because all flows are fibrant for the weak S-homotopy model structure of $`\mathrm{𝐅𝐥𝐨𝐰}`$.
Let $`X`$ be a homotopy continuous flow. Then, for instance, consider the unique morphism $`Q(f):Q(\{\widehat{0}<\widehat{1}\})Q(\{\widehat{0}<A<\widehat{1})\}`$ of $`𝒯`$. For any commutative square
there exists $`k:Q(\{\widehat{0}<A<\widehat{1}\})X`$ making the triangle commutative. Therefore, the existence of $`k`$ ensures that any directed segment of $`X`$ can always be divided up to S-homotopy. Roughly speaking, a flow $`X`$ is homotopy continuous if it is indefinitely divisible up to S-homotopy.
###### Proposition 16.
The pair $`(\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯),\mathrm{𝐢𝐧𝐣}(J^{gl}𝒯))`$ is a cofibrantly generated weak factorization system. Moreover, it satisfies the conditions of Section 1, that is:
$$\mathrm{𝐜𝐨𝐟}(J^{gl})\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)\mathrm{𝐜𝐨𝐟}(I_+^{gl}).$$
###### Proof.
For every $`gJ^{gl}𝒯`$, the continuous map $`g`$ is a closed inclusion of topological spaces. Therefore by Proposition 11.5 and by Theorem 2.1.14, the small object argument applies. ∎
###### Notation 5.
$`𝒮_𝒯=𝒮_{\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)},`$
$`R_𝒯=R_{\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)},`$
$`_𝒯=_{\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)},`$
$`\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}=\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)}.`$
We can now apply Theorem 2 to obtain the theorem:
###### Theorem 3.
The inclusion functors $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}\mathrm{𝐅𝐥𝐨𝐰}_{cof}\mathrm{𝐅𝐥𝐨𝐰}`$ induce the equivalences of categories
$$\mathrm{𝐅𝐥𝐨𝐰}_{cof}^{f,𝒯}/_𝒯\mathrm{𝐅𝐥𝐨𝐰}_{cof}[(𝒮\mathrm{𝐜𝐨𝐟}(J^{gl}T))^1]\mathrm{𝐅𝐥𝐨𝐰}[𝒮_𝒯^1].$$
It remains to check the invariance of the underlying homotopy type and of the branching and merging homology theories:
###### Theorem 4.
A morphism of $`𝒮_𝒯`$ preserves the underlying homotopy type and the branching and merging homology theories.
###### Proof.
It has been already noticed above that up to weak S-homotopy and up to the 2-out-of-3 axiom, a morphism of $`𝒮_𝒯`$ is a morphism of $`𝒮\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)`$ between cofibrant flows. Formally, let $`f:AB`$ be an element of $`𝒮_𝒯`$. Then consider the commutative diagram:
The morphism $`Q(f)`$ is an isomorphism of $`\mathrm{𝐅𝐥𝐨𝐰}_{cof}[(𝒮\mathrm{𝐜𝐨𝐟}(J^{gl}T))^1]`$. Therefore it preserves the underlying homotopy type and the branching and merging homology theories because any morphism of $`𝒮`$ preserves these invariants by Proposition VII.2.5 and by Corollary 6.5 and Corollary A.11, and because any morphism of $`\mathrm{𝐜𝐨𝐟}(J^{gl}𝒯)`$ preserves these invariants by Theorem 5.2. The morphisms $`Q(A)A`$ and $`Q(B)B`$ are weak S-homotopy equivalences. So both preserve the underlying homotopy type Proposition VII.2.5 and the branching and merging homology theories Corollary 6.5 and Corollary A.11. Hence the result. ∎
|
warning/0506/astro-ph0506623.html
|
ar5iv
|
text
|
# Fully Degenerate Self-Gravitating Fermionic Dark Matter: Implications to the Density Profile of the Cluster of Galaxies A1689, and the Mass Hierarchy of Black Holes
## 1 Introduction
Recent precision observations have revealed that the unknown dark matter dominates the matter contents of the Universe. We wish to study the possible dynamical structures of their existence, especially in a universal form. If the dark matter is in the form of ordinary thermal gas, the structure and the dimension would be strongly dependent on the initial conditions and the environment of the expanding universe. On the other hand if the dark matter is almost degenerate, we will naturally find a universal structure. Moreover, we would like to know the possibility that the dark matter forms black holes.
To make the problem setting better defined, we start our consideration from the typical structure made from ordinary baryonic matter, in which an electron and a nucleon form the basic ingredient through the electromagnetic force. In the total energy $`E=\frac{1}{2}m_ev^2(e^2/r)`$, the second term represents the attractive force making the system collapse and the first term represents the pressure against it through the Heisenberg’s uncertainty principle. Actually putting the expression of the de Broglie wave length $`r=\mathrm{}/(m_ev)`$ in $`E`$, and extremizing it with respect to $`r`$, we obtain the ground state bounding energy of hydrogen $`E=m_ee^4/(2\mathrm{}^2)=2.19\times 10^{18}`$ J and the Bohr radius $`r=\mathrm{}^2/(e^2m_e)=5.28\times 10^9`$ cm. The corresponding energy density is $`\rho =m_p/(\frac{4\pi }{3}r^3)=2.71`$ g$``$cm<sup>-3</sup>. This is the basic ingredient of the structure formed by electromagnetism. The electron can collapse more. Setting $`vc`$ in de Broglie wave length yields the Compton wave length. The corresponding energy density becomes $`\rho =m_p/(\frac{4\pi }{3}r_c^3)=6.93\times 10^6`$ g$``$cm<sup>-3</sup>. This is the basic ingredient of white dwarfs. Further collapse makes the electron and the proton into a neutron, through the inverse beta decay process. Putting $`m_em_p`$, we have $`r=\mathrm{}/(m_pc)`$, and therefore the corresponding energy density is $`\rho =m_p^4c^3/\mathrm{}^3=4.29\times 10^{16}`$ g$``$cm<sup>-3</sup>. This is the basic ingredient of neutron stars.
In general, the material formed by fully degenerate fermion of mass $`m`$ would have the energy density $`\rho =m^4c^3/\mathrm{}^3`$. A structure of this density with radius $`R`$ has the mass $`M=R^3\rho `$. On the other hand the size $`R`$ should be smaller than the limiting scale, Schwarzschild radius, $`2GMR`$. This condition yields the maximum mass as
$$M_{fermi}=G^{3/4}m^2=m_{pl}^3/m^2$$
(1)
where the Planck mass is defined as $`m_{pl}=\left(\mathrm{}c/G\right)^{1/2}=2.18\times 10^5`$ g. Thus the quantum mechanics ($`\mathrm{}`$), gravity ($`G`$), and the particle physics ($`m`$) as well as relativity ($`c`$) characterize the universal structure of a fully degenerate fermion star (FDFS).
Incidentally, degenerate bosons form much smaller structures since bosons have no exclusion principle like fermions. Therefore only the Heisenberg uncertainty principle (quantum pressure) can support the structure against the collapse due to gravity. The sole characteristic length scale is the Compton wave length $`l_{Compton}=\mathrm{}/(mc)R`$, which must be larger than the limiting scale, Schwarzschild radius, $`2GMR`$. This condition yields the density $`\rho =m^2m_{pl}^2c^3/\mathrm{}^3`$ and
$$M_{bose}=m_{pl}^2/m,$$
(2)
which is known as the Kaup mass (Kaup, 1968). This structure is a boson star. This structure is smaller than that formed by fermions by a factor $`m_{pl}/m`$ and no further discussion will be given in this paper, apart from a brief comment regarding the equation of state in §2.
Now we proceed to consider the equilibrium structures made of fermionic dark matter. One may think that this problem is analogous to the white dwarf case and the Chandrasekhar mass is the limiting mass. However this is not the right answer. In the case of a white dwarf, the pressure is due to the degeneracy pressure of relativistic electrons, but the gravity is Newtonian because it is due to the rest mass of hadrons which are nonrelativistic. When we treat a degenerate star made purely of dark matter itself, the total rest mass is also due to dark matter particles themselves which can be relativistic. Therefore we need to consider the gravitational mass of the star in a general relativistic manner. This problem is therefore analogous to the case of a neutron star. In the case of the neutron star, there is a complication due to the fact that neutrons interact by nuclear forces. However in the case of weakly interacting fermionic dark matter, it can be considered as ideal gas and the equation of state is well defined. It turned out that this situation was treated by the classical paper by Oppenheimer & Volkoff (1939), since they assumed free neutrons and used the equation of state of ideal gas in their calculation. Basically fully degenerate fermions can form very compact dense structures including black holes (Bilić, Munyaneza & Viollier, 1999; Bilić, Tupper & Viollier, 2003).
On the other hand, from the observational side, we now have many candidates of dense structures at various scales. The most massive example is a huge dark matter distribution of the cluster of galaxies, A1689, reported recently by Broadhurst et al. (2005a, b). They obtained the mass column density distribution of the cluster of galaxies, A1689 by gravitational lensing. The column density profile has a flat-top and we suspected that this flat-top nature might be due to the degeneracy pressure of fermions. We derive a volume density profile from the observed column density profile assuming spherical symmetry and compare the observed 3D encircled mass profile with our model profile of an FDFS.
The condense structure made from the fully degenerate fermion is not restricted to a center of a cluster of galaxies. If we consider more massive neutrinos, such as sterile neutrinos, the similar structures are realized in scaled-down form as in Eq.(1). If this structure is universal, we will find groups of black holes which have typical masses directly characterized by fermion masses.
The paper is organized as follows. The formalism by Oppenheimer & Volkoff (1939) is introduced and equilibrium solutions are discussed in §2. Readers who are not interested in the derivation of the solutions, are advised to skip to §2.3 where the properties of the solutions are discussed. The application of a nonrelativistic FDFS to the cluster of galaxies A1689, is discussed in §3. The possible relation between the mass hierarchies of black holes and sterile neutrinos are considered in §4.
## 2 Formalism
Here we review the derivation of the general relativistic equilibrium equations following Tolman (1934) and Oppenheimer & Volkoff (1939).
### 2.1 Relativistic treatment of equilibrium
The most general static line element exhibiting spherical symmetry may be expressed in the form
$$ds^2=e^\lambda dr^2r^2d\theta ^2r^2\mathrm{sin}^2\theta d\varphi ^2+e^\nu dt^2,\lambda =\lambda (r),\nu =\nu (r).$$
(3)
If the matter supports no traverse stresses and has no mass motion, then its energy momentum tenor is given by
$$T_1^1=T_2^2=T_3^3=p,T_4^4=ϵ,$$
(4)
where $`p`$ and $`ϵ`$ are respectively the pressure and the macroscopic energy density measured in proper coordinates. Einstein’s field equations without the cosmological constant reduce to
$$8\pi p=e^\lambda \left(\frac{\nu ^{}}{r}+\frac{1}{r^2}\right)\frac{1}{r^2},$$
(5)
$$8\pi ϵ=e^\lambda \left(\frac{\lambda ^{}}{r}\frac{1}{r^2}\right)+\frac{1}{r^2},$$
(6)
$$\frac{dp}{dr}=\frac{p+ϵ}{2}\nu ^{},$$
(7)
where primes denote differentiation with respect to $`r`$. These three equations together with the equation of state of the material $`ϵ`$ = $`p(ϵ)`$ determine the mechanical equilibrium of the matter distribution as well as the dependence of the metric $`g_{\mu \nu }`$’s on $`r`$.
The boundary of the matter distribution is the value of $`r=r_b`$ for which $`p=0`$, and such that for $`r<r_b,p>0`$. For $`r<r_b`$ the solution depends on the equation of state of the material connecting $`p`$ and $`ϵ`$. For many equations of state a sharp boundary exists with a finite value of $`r_b`$.
In the empty space, $`p=ϵ=0`$, surrounding the spherically symmetric distribution of matter, the Schwarzschild’s exterior solution is obtained:
$$e^{\lambda (r)}=1\frac{2m}{r},e^{\nu (r)}=1\frac{2m}{r},$$
(8)
where $`m`$ is the Newtonian mass of the matter as calculated by a distant observer.
Inside the boundary, Eqs. (5), (6), and (7) may be rewritten as follows. Using the equation of state $`ϵ=ϵ(p)`$, Eq. (7) may be immediately integrated
$$\nu (r)=\nu (r_b)_0^{p(r)}\frac{2dp}{p+ϵ(p)},e^{\nu (r)}=e^{\nu (r_b)}\mathrm{exp}\left[_0^{p(r)}\frac{2dp}{p+ϵ(p)}\right].$$
(9)
The constant $`e^{\nu (r_b)}`$ is determined by making $`e^\nu `$ continuous across the boundary.
$$e^{\nu (r)}=(1\frac{2m}{r})\mathrm{exp}\left[_0^{p(r)}\frac{2dp}{p+ϵ(p)}\right].$$
(10)
Thus $`e^\nu `$ is known as a function of $`r`$ if $`p`$ is known as a function of $`r`$. Further in Eq.(6). introduce a new variable
$$u(r)=\frac{1}{2}r(1e^\lambda )\mathrm{or}e^\lambda =1\frac{2u}{r}.$$
(11)
Then Eq. (5) becomes:
$$\frac{du}{dr}=4\pi ϵ(p)r^2.$$
(12)
In Eq. (5) we replace $`e^\lambda `$ by its expression (11) and $`\nu ^{}`$ by its expression (7). Then it becomes
$$\frac{dp}{dr}=\frac{p+ϵ(p)}{r(r2u)}[4\pi pr^3+u].$$
(13)
Eqs. (12) and (13) form a system of two first-order equations in $`u`$ and $`p`$. Starting with some initial values $`u=u_0`$ and $`p=p_0`$ at $`r=0`$, the two equations are integrated simultaneously to the value $`r=r_b`$ where $`p=0`$, i.e., until the boundary of the matter distribution is reached. The value of $`u=u_b`$ at $`r=r_b`$ determines the value of $`e^{\lambda (r_b)}`$ at the boundary to the exterior solution, making
$$u_b=\frac{r_b}{2}[1e^{\lambda (r_b)}]=\frac{r_b}{2}\left[1(1\frac{2m}{r_b})\right]=m.$$
(14)
Thus the mass of this spherical distribution of matter as measured by a distant observer is given by the value $`u_b`$ of $`u`$ at $`r=r_b`$.
The following restrictions must be made on the choice of $`p_0`$ and $`u_0`$, the initial values of $`p`$ and $`u`$ at $`r=0`$:
(a) In accordance with its physical meaning as pressure, $`p_00`$.
(b) From Eq.(11) it is seen that for all finite values of $`e^\lambda `$, $`u_0=0`$. Since $`g_{11}=e^\lambda `$ must never be positive, $`u_00`$ for infinite values of $`e^\lambda `$ at the origin. However, it may be shown that of all the finite values of $`p_0`$ at the origin $`p_0=0`$ is the only one compatible with a negative value of $`u_0`$, and that for equations of state of the type occurring in this problem even this possibility is excluded, so that $`u_0`$ must vanish.
This can be seen from the following argument. Having chosen some particular value of $`p_0`$ one may usually represent the equation of state in that pressure range by $`ϵ=Cp^s`$ with some appropriate value of $`s`$. Using this equation of state and taking the approximate form of Eq (13) near the origin for the case $`u_0<0`$, and finite $`p_0`$, one obtains:
$$\frac{dp}{dr}=\frac{p+ϵ(p)}{2r}=\frac{p+Cp^s}{2r}.$$
(15)
Integration of this equation shows that for $`s<1`$, $`p_00`$ can not be satisfied, and for $`s1`$ only the value $`p_0=0`$ is possible. This immediately excludes the possibility that degenerate bosons form an equilibrium structure, because $`p`$ is independent of $`ϵ`$, if $`p`$ is solely due to thermal bosons (Landau & Lifshitz, 1980). As we mentioned in §1, a boson star is supported by the quantum pressure of ground-state bosons, which is outside of the scope of the consideration of the equation of state (Kaup, 1968). For the equations of state used for degenerate fermions, always $`s<1`$ holds. It is also be noted that the above equation together with Eq. (10) show that $`e^{\nu (r)}\mathrm{}`$ as $`r0`$.
(c) A special investigation for any particular equation of state must be made to see whether solutions exist in which $`0u_0\mathrm{}`$ and $`p\mathrm{}`$ as $`r0`$.
### 2.2 Equation of state for degenerate Fermi gas
If the matter consists of fermions of rest mass $`\mu _0`$ and statistical weight $`g`$, and their thermal energy and all forces between them are neglected, then a parametric form for the equation of state (Landau & Lifshitz, 1980) is,
$$ϵ=K(\mathrm{sinh}tt),$$
(16)
$$p=\frac{1}{3}K\left(\mathrm{sinh}t8\mathrm{sinh}\frac{t}{2}+3t\right),$$
(17)
where
$$K=\frac{\pi g\mu _0^4c^5}{8h^3},$$
(18)
and
$$t=4\mathrm{a}\mathrm{r}\mathrm{c}\mathrm{s}\mathrm{i}\mathrm{n}\mathrm{h}\frac{p_F}{\mu _0c},$$
(19)
where $`p_F`$ is the maximum momentum in the Fermi distribution and is related to the proper particle density $`N/V`$ by
$$\frac{N}{V}=\frac{4\pi g}{3h^3}p_F^3=\frac{4\pi g}{3}(\frac{\mu _0c}{h})^3\mathrm{sinh}^3\frac{t}{4}.$$
(20)
If we define the Fermi velocity $`v_F`$ by
$$\frac{p_F}{\mu _0c}=\frac{v_F}{\sqrt{1(v_F/c)^2}},$$
(21)
$$t=2\mathrm{log}\left(\frac{1+v_F/c}{1v_F/c}\right)\mathrm{or}v_F/c=\mathrm{tanh}\frac{t}{4}.$$
(22)
$`t`$ and $`v_F`$ are independent of particle properties.
Substituting the above expressions for $`p`$ and $`ϵ`$ into Eqs. (12) and (13) one obtains:
$$\frac{du}{dr}=4\pi r^2K(\mathrm{sinh}tt),$$
(23)
and
$`{\displaystyle \frac{dt}{dr}}`$ $`=`$ $`4/[r(r2u)]`$ (24)
$`\times (\mathrm{sinh}t2\mathrm{sinh}1/2t)/(\mathrm{cosh}t4\mathrm{cosh}1/2t+3)`$
$`\times \left[{\displaystyle \frac{4\pi }{3}}Kr^3(\mathrm{sinh}t8\mathrm{sinh}1/2t+3t)+u\right].`$
These equations are to be integrated from the values $`u=0,t=t_0`$ at $`r=0`$ to $`r=r_b`$ where $`t_b=0`$ (which makes $`p=0`$), and $`u=u_b`$.
So far, the equations are written in relativistic units, i.e., such that $`c=1,G=1`$. This determines the unit of time and the unit of mass in terms of still arbitrary unit of length. The unit of length is now fixed by the requirement that $`K=1/(4\pi )`$. From the dimensional analysis of Einstein’s field equations, this requirement fixes the unit of length to be
$$a=\frac{1}{\pi }\left(\frac{2}{g}\right)^{1/2}\left(\frac{h}{\mu _0c}\right)^{3/2}\frac{c}{(\mu _0G)^{1/2}},$$
(25)
and the unit of mass to be
$$b=\frac{c^2}{G}a=\frac{1}{\pi }\left(\frac{2}{g}\right)^{1/2}\left(\frac{h}{\mu _0c}\right)^{3/2}\frac{c^3}{(\mu _0G^3)^{1/2}}.$$
(26)
Eqs. (12) and (24) written in a dimensionless form become:
$$\frac{du}{dr}=r^2(\mathrm{sinh}tt),$$
(27)
$`{\displaystyle \frac{dt}{dr}}`$ $`=`$ $`4/[r(r2u)]`$ (28)
$`\times (\mathrm{sinh}t2\mathrm{sinh}1/2t)(\mathrm{cosh}t4\mathrm{cosh}1/2t+3)`$
$`\times \left[{\displaystyle \frac{1}{3}}r^3(\mathrm{sinh}t8\mathrm{sinh}1/2t+3t)+u\right].`$
For a given $`t_0`$, Eqs. (27) and (28) can be integrated. Fixing $`t_0`$ is equivalent to fixing $`v_{F0}`$. As long as dark matter particles are fully degenerate, the solution describes the equilibrium between the gravity and degeneracy pressure from Newtonian to general relativistic gravity and from nonrelativistic $`v_F`$ to relativistic $`v_F`$. Therefore the solution for a given $`t_0`$ (or $`v_{F0}`$) is independent of particle properties, while the units of length and mass ($`a`$ and $`b`$) are fixed by the particle properties $`\mu _0`$ and $`g`$.
### 2.3 Discussion on solutions
The Eqs.(27) and (28) are numerically integrated using the fourth-order Runge-Kutta method. Apart from the gravitational mass $`u`$, there is another mass indicator, the dimensionless total rest mass $`y_b`$, defined as
$$y_b=_0^{r_b}\frac{32}{3}\mathrm{sinh}^3\frac{t}{4}/\sqrt{(12u/r)}r^2dr,$$
(29)
which is the integral of the number density with the proper volume inside the radius $`r_b`$. $`t_0`$, $`v_{F0}/c`$, $`u_b`$, $`y_b`$, and $`r_b`$ are given for $`t_0=0.114.0`$ in Table 1. At $`t_0=0.1`$, degenerate particles are nonrelativistic ($`v_{F0}/c=0.025`$), while at $`t_0=14`$, particles are extremely relativistic ($`v_{F0}=0.998`$). In Fig. 1, the relation between the gravitational mass $`u_b`$ and outer radius $`r_b`$ is plotted. For $`t_03`$, $`u_b`$ is an increasing function of $`t_0`$, while $`r_b`$ is a decreasing function of $`t_0`$. The maximum of $`u_b=0.0766`$ is reached for $`t_0=3`$ and $`r=0.663`$. This is the maximum stable solution for which $`v_{F0}/c=0.635`$, which is modestly relativistic. In Fig.2, the radial profiles of the energy density, $`ϵ`$, and pressure, $`p`$, are plotted. The contribution of $`p`$ is relatively small compared to $`ϵ`$. This solution is termed as quasi-Newtonian by Oppenheimer & Volkoff (1939). In the case of self-gravity of degenerate fermions themselves, a stable configuration never reaches an extremely relativistic situation unlike the case of white dwarfs. Therefore the formula of the Chandrasekhar mass applied to this case over-estimates the maximum stable mass by one order of magnitude. For $`t_0>3.0`$, there is no stable solution (Oppenheimer & Volkoff, 1939). $`u_b`$ decreases takes a minimum value at ($`u_b`$,$`r_b`$)=(0.0395,0.364) for $`t_0=8.0`$. For a large $`t_0`$, ($`u_b`$,$`r_b`$) spirals to $`(0.041,0.29)`$.
The gravitational mass defect, $`\mathrm{\Delta }=u_by_b`$ is one measure of stability of a general relativistic equilibrium configuration. In Fig.3, both $`y_b`$ and $`u_b`$ are plotted against $`t_0`$ for $`0.1<t_0<14.0`$. For a small $`t_0`$, $`\mathrm{\Delta }`$ is negative and gradually decreases and takes a minimum value, $`\mathrm{\Delta }=0.0029`$ for the maximum mass stable solution ($`t_0=3.0`$) and then increases. The fractional mass defect, or the packing fraction, is $`f=\mathrm{\Delta }/y_b=0.036`$. $`|f|`$ is only 3.6% and may appear small. However, if we compare this value to a typical nuclear packing fraction, which is less than 0.1% (Fermi, 1950), we will find it to be significant. $`\mathrm{\Delta }<0`$ is a necessary condition for stability, but not sufficient. $`\mathrm{\Delta }<0`$ still holds for $`3<t_0<5`$, but these solutions correspond to unstable equilibria. When $`\mathrm{\Delta }<0`$, solutions are for stable equilibria for $`d\mathrm{\Delta }/dt_00`$, while they are for unstable equilibria for $`d\mathrm{\Delta }/dt_0>0`$.
The question of whether dark matter fermions can form a black hole would be answered in terms of the maximum gravitational mass, $`M_{max}`$ and the corresponding radius $`R_{max}`$, as functions of the particle mass $`\mu _0`$ and statistical weight $`g`$. The results are:
$`M_{max}`$ $`=`$ $`1.73\times 10^{51}{\displaystyle \frac{1}{\sqrt{g}}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^2(\mathrm{g})`$ (30)
$`=`$ $`8.70\times 10^{17}{\displaystyle \frac{1}{\sqrt{g}}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^2(\mathrm{M}_{}),`$ (31)
and
$`R_{max}`$ $`=`$ $`1.10\times 10^{24}{\displaystyle \frac{1}{\sqrt{g}}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^2(\mathrm{cm})`$ (32)
$`=`$ $`3.54\times 10^2{\displaystyle \frac{1}{\sqrt{g}}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^2(\mathrm{kpc}).`$ (33)
It should be noted that Landau & Lifshitz (1980) define the maximum rest mass, $`M_{rest}=1.037M_{max}`$ as the maximum mass, which is the total rest mass brought from infinity. In the case of a star made of free neutrons, $`\mu _0=939`$ MeV and $`g=2`$. Then $`M_{max}=0.70\mathrm{M}_{}`$, and $`R_{max}=8.8`$ km. It should also be noted that $`R_{max}`$ is sensitive to the actual outer boundary condition in the numerical calculation, which in principle should be $`p=0`$. In reality, integration must be stopped at a finite value of $`p`$ and $`r_b`$ is somewhat sensitive to this finite $`p`$. Physically speaking, it is unrealistic to assume the fully degenerate equation of state to a small $`p`$ and the outer boundary condition is not well defined. The total gravitational mass $`u_b`$ is not very sensitive to the choice of the actual outer boundary condition $`p`$ and is well defined. Although it is not explicitly stated in the original paper, Oppenheimer & Volkoff (1939) were probably aware of this limitation in the applicability of the equation of state, if we judge from the carefully chosen title “On Massive Neutron Cores”. We also admit that without this boundary condition, we cannot use the Schwarzschild’s exterior solution outside the boundary and the formulation becomes more complicated.
Before concluding this section, we comment on the nonrelativistic limit. For the nonrelativistic limit, it is expected that between the mass, $`M`$, and the size, $`R`$, of an FDFS, the relation,
$$MR^3=const.$$
(34)
should hold (Landau & Lifshitz, 1980).
For $`t_0<0.5`$, $`u_br_b^3`$ is indeed nearly constant and
$$u_br_b^3=0.135,$$
(35)
within 4% precision. In physical units,
$`MR^3`$ $`=`$ $`(u_bb)(r_ba)^3`$ (36)
$`=`$ $`1.4\times 10^{124}{\displaystyle \frac{1}{g^2}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^8(\mathrm{g}\mathrm{cm}^3)`$
$`=`$ $`2.3\times 10^{26}{\displaystyle \frac{1}{g^2}}\left({\displaystyle \frac{\mu _0}{\mathrm{eV}}}\right)^8(\mathrm{M}_{}\mathrm{kpc}^3).`$ (37)
Nonrelativistic solutions are similar and in Fig.4, the profiles of the normalized mass density and 3D encircled mass are plotted. In the next section, we will find the behavior of the 3D encircled mass profile interesting and the logarithmic profile of the normalized 3D encircled mass is given in Fig.5.
## 3 Can degenerate fermions be the dominant dark matter in the cluster of galaxies, A1689?
Recently Broadhurst et al. (2005a, b) reported a mass column density profile of the cluster of galaxies, A1689, obtained from gravitational lensing. One of the important properties of the profile is that it has a flat top. We propose that this flat-top column density profile might be explained by the effects of degeneracy pressure of fermionic dark matter. Here we analyze this proposal.
First we briefly introduce the main results of Broadhurst et al. (2005a, b). In their analysis, 1 corresponds to 129 kpc $`h^1`$. In Broadhurst et al. (2005a), the central 250 kpc $`h^1`$ in radius of multi-color HST/ACS images were analyzed. The mass column density profile, $`\mathrm{\Sigma }(r)`$, is not expressed as a single power law of radius. The mass column density profile flattens toward the center with a mean slope of dlog$`\mathrm{\Sigma }`$/dlogr $`0.55`$ within $`r<`$250 kpc $`h^1`$. Inside the Einstein radius ($`\theta _E50^{\prime \prime }`$), they obtained the slope of $`0.3`$ from the ratio between $`\theta _E`$ and the radius of the radial critical curve, $`\theta _r17^{\prime \prime }`$. They fit their results with an inner region of an NFW profile (Navarro, Frenk & White, 1996) with a relatively high concentration, $`C_{vir}=8.2`$.
The mass column density, $`\mathrm{\Sigma }(r)`$, is the integral of the volume density, $`\rho (r)`$, along the line of sight over the entire cluster scale of Mpc. In order to study the possibility of fermion degeneracy near the center of the cluster, we need information on the volume density, $`\rho (r)`$, instead of the column density, $`\mathrm{\Sigma }(r)`$. Broadhurst et al. (2005b) present the weak-lensing analysis of the wide field data obtained by Subaru and obtained the column density profile at $`r<2`$ Mpc $`h^1`$. They fit the combined profile of HST/ACS and Subaru with an NFW profile with a very high concentration, $`C_{vir}=13.7`$, significantly larger than theoretically expected value of $`C_{vir}4`$. They also fit the same observed column density profile with a power law profile with a core. They give this result in terms of the angular radius dependence of the convergence, $`\kappa `$, as
$$\kappa (\theta +\theta _C)^n.$$
(38)
$`\theta _C=1.65^{}`$ and $`n=3.16`$ give the best fit although $`\theta _C`$ and $`n=3.16`$ are mutually dependent and a finite range of the combination ($`\theta _C`$,$`n`$) gives equally good fits. In terms of $`\chi ^2`$ and the degrees of freedom, this core power law profile fits the observation better than the best-fit NFW profile and we use this profile for further discussion. Although Broadhurst et al. (2005b) do not claim so explicitly, the two facts that the best-fit NFW profile shows a much higher concentration than the value predicted by the CDM cosmology and the phenomenological profile, Eq.(38), fits better than the best-fit NFW profile, indicate some serious contradiction to the CDM cosmology.
We start our analysis from this core power-law profile, (38), for further discussion. We convert (38) to a column density profile, $`\mathrm{\Sigma }(r)`$, in physical units of length and mass using the relations, $`\kappa =\mathrm{\Sigma }/\mathrm{\Sigma }_{crit}`$, $`\mathrm{\Sigma }_{crit}`$ 0.95 $`\mathrm{g}\mathrm{cm}^2`$, and the normalization of 2D encircled mass inside the Einstein radius, $`r_E`$, $`_0^{r_E}\mathrm{\Sigma }(r)2\pi r𝑑r=\mathrm{\Sigma }_{crit}\pi r_E^2`$. The result is expressed as
$$\mathrm{\Sigma }(r)=25.2\left(r/r_E+2.2\right)^{3.16},(\mathrm{g}\mathrm{cm}^2)$$
(39)
where $`r_E=97`$ kpc $`h^1`$ corresponds to $`\theta _E=45^{\prime \prime }`$, the value used in Broadhurst et al. (2005b). The 2D encircled mass, $`M_2(r)=\mathrm{\Sigma }(r)2\pi r𝑑r`$, is analytically obtained and $`M_2(r)=1.3\times 10^{14}h^2M_{}`$ and $`1.1\times 10^{15}h^2M_{}`$ respectively for $`r=r_E`$ and $`r=\mathrm{}`$. Therefore a high concentration of the mass is expected on the scale of $`r_E`$. By assuming spherical symmetry, we wish to obtain the volume density $`\rho (r)`$ by solving
$$\mathrm{\Sigma }(x)=\rho (\sqrt{x^2+z^2})𝑑z,$$
(40)
but we were not able to obtain an analytic solution. Instead, we assumed another power law with a core for $`\rho (r)`$ and obtained the best fit parameters. The range of integration in $`z`$ is from $``$2Mpc $`h^1`$ to $`+`$2Mpc $`h^1`$. The result is
$$\rho (r)=1.60\times 10^{23}\left(r/r_E+1.28\right)^{3.71}h.(\mathrm{g}\mathrm{cm}^3)$$
(41)
Near the center, $`\rho (r)=6.4\times 10^{24}h`$ and $`7.5\times 10^{25}h`$ (g$``$ cm<sup>-3</sup>) respectively at $`r`$ = 0 and $`r_E`$. As for the case with the column density profile, the core radius and power-law index are mutually dependent and a finite range of their combination gives equally good fits. The above volume mass densities may appear small, but the degeneracy depends on the number density and de Broglie wavelength, both of which are heavily dependent on the rest mass of the particle. Before proceeding to the analysis by an FDFS, we first confirm that eV-mass fermions can be degenerate at these low mass densities. Since 1 eV corresponds to $`1.8\times 10^{33}`$ g, the number density, $`N/V10^{11}`$ cm<sup>-3</sup>, the mean inter-particle spacing, $`(N/V)^{1/3}2\times 10^4`$ cm. On the other hand, the de Broglie wavelength for a 1 eV particle with a velocity $`v`$ is, $`\mathrm{}/\mu _0v=\lambda _{Compton}(c/v)=2\times 10^5(c/v)`$ cm. Therefore for nonrelativistic particles with $`v<0.1c`$, the condition for degeneracy, $`(N/V)^{1/3}<\lambda _{(\mathrm{deBroglie})}`$, is satisfied.
More quantitative discussion is possible for a 3D encircled mass profile given in Fig.6. For the purpose of an explicit comparison, here we fix $`h=0.7`$. The 3D encircled mass profile also predicts the rotation curve profile as Fig.7, which can be compared with observations of kinematics of the galaxies in the cluster in the future. The column density profile, Eq.(39), was derived as a phenomenological formula without assuming any background physics, and so was the volume density profile, Eq.(41). Here we wish to explain the presence of the core in the power-law profile, which causes the flat-top nature of Eq.(39), in terms of an FDFS. The comparison must be made either in the volume density profiles or in the 3D encircled mass profiles. Here we choose the latter, because the observed 3D encircled mass profile is much less sensitive to the actual choice of the combination of the core radius and the power-law index in the original Eq.(39). As we mentioned above, the finite range of the parameter combinations gives equally good fits to the observations. The basic reason for this is that our observables are integrated quantities rather than the local densities. First, we compare the slopes of the 3D encircled mass profiles of the observation (Fig.6) and our model (Fig.5). The observed slope is 2.93 at $`r<`$10 kpc and 2.54 at 10 $`<r<`$ 100 kpc. On the other hand, the slope of a nonrelativistic FDFS is 2.90 between $`r=0.03r_b`$ and $`0.3r_b`$ and becomes shallower at $`0.3r_b<r<r_b`$. Note that the slope of the 3D encircled mass profile for a constant volume density is 3.0 and both the observed and model inner slopes of 2.9 are consistent with the flat-top nature of the volume density profile. Now we proceed to the physical scaling of the model in terms of the mass and length. For this purpose, it is most convenient to use the relation Eq.(37) for the mass and size of a nonrelativistic FDFS. The Eq.(37) can be rewritten as
$$\mathrm{log}M_{14}+3\mathrm{log}R_2=f(g,\mu _0),$$
(42)
where $`M_{14}=M/(10^{14}M_{})`$, $`R_2=R/(100\mathrm{k}\mathrm{p}\mathrm{c})`$ and $`f(g,\mu _0)=6.362\mathrm{log}g8\mathrm{log}(\mu _0/\mathrm{eV})`$. In Fig.8, the possible combinations of $`M`$ and $`R`$ are plotted for different values of $`f(g,\mu _0)`$. Intersections of the dotted lines and the solid line (3D encircled mass profile of A1689) correspond to solutions for representative values of $`f`$. Actually $`f`$ is continuous and the solutions are continuous. The probable range is $`f=62`$. For this range of $`f`$, the rest mass of a fermion ranges from 30 to 2 eV. The larger the rest mass, the smaller the degenerate structure. The possible range of $`g`$ is limited. For a particle with spin 1/2, $`g`$ is either 1 (Majorana particle) or 2 (Dirac particle).
The special case is the massive neutrinos with similar masses for which effective $`g=3`$ (Majorana) or $`g=6`$ (Dirac). Recent underground experiments have shown that the mass differences among three species of neutrinos are smaller than $`0.05`$ eV (Shirai, 2005). In order for massive neutrinos to have masses greater than 1 eV, they must have similar masses (degenerate in mass). The mean number density of relic neutrinos is cosmologically fixed and there is the well known relation for the contribution of the relic neutrinos to $`\mathrm{\Omega }`$,
$$\mathrm{\Omega }_\nu h^2=\mathrm{\Sigma }_i\mu _i/(93.5\mathrm{eV}),$$
(43)
where $`i=13`$ for Majorana neutrinos and $`i=16`$ for Dirac neutrinos. The maximum allowed neutrino mass can be estimated by setting $`\mathrm{\Omega }_\nu =0.3`$ and $`h=0.7`$, to be 4.7 or 2.3 eV respectively for Majorana or Dirac neutrinos. The former gives $`f=0`$ or the degenerate mass of $`10^{14}M_{}`$ in 100 kpc, and the latter gives $`f=2`$ or that of $`5\times 10^{14}M_{}`$ in 300 kpc. These are rather comfortable numbers for this cluster profile. If we question the CDM cosmology based on the possible inconsistency of the high concentration of the observed mass profile with the CDM predictions mentioned above, we might as well revive massive neutrinos (hot dark matter) as the candidate for the dominant dark matter.
## 4 Mass hierarchy of black holes from that of neutrino through FDFS?
In the previous section, we have studied a possibility that fully degenerate fermions forms a huge mass concentration at the center of a cluster of galaxies and we mainly examined the case for massive neutrinos. Is this possibility really true? In order to answer this question, we focus on the universality and the scalability of FDFS. We try to extend this idea of FDFS to other neutrinos and fermions with different masses. As we have already seen in Eq.(1,31), the rest mass of the fermion mostly determines the characteristic mass scale of the structure; more massive fermions form lighter structures.
The most extreme condensed structure is a black hole. We have already known that there exist many black holes of several species in the Universe (Gebhardt et al., 2002). They are the most familiar stellar mass black holes ($`M_{}`$), giant black holes at the center of a galaxy ($`10^7M_{}`$), and the intermediate mass black holes ($`10^3M_{}`$). Although these massive black holes are actively studied based on the bottom-up scenarios that they are formed by the coalescence of the stellar sized black holes, we try to propose yet another scenario based on the context of FDFS. The most prominent property of those black holes are that they appear to have a hierarchy in mass range. Most of the black holes are classified in the above three types and those in other mass ranges is rare. Therefore it would be natural to suspect any definite mechanism to construct such hierarchy from the fundamental level. Our hypothesis is that such fundamental mechanism is the mass hierarchy in fermions or possibly neutrinos. If those black holes are formed by the overweight FDFSs, we can estimate the masses of those neutrinos as in Table 2.
The corresponding masses of fermions, except $`\nu _{e,\mu ,\tau }`$, through Eq.(1) are far heavier than eV and we may identify those fermions as more massive sterile neutrinos. Lighter neutrino of mass $`10^110^3`$ eV would yield FDFSs much extended dilute structures whose sizes well exceed the Horizon size. In the above, we have approximate values. However, the lower limit of the black hole mass within a class yields the precise value of the corresponding neutrino mass. The existence of such a lower limit would be the key ingredient of the FDFS model.
## 5 Conclusions and Discussions
We have examined the possible structures formed by the fully degenerate self-gravitating fermions (FDFS) at various scales, such as the mass concentration of a cluster of galaxies, the giant black holes at the center of a galaxy and the intermediate black holes in galaxies. As an order estimation, their characteristic masses directly reflect the constituent fermion masses through the simple relation, $`M_{fermi}=G^{3/4}m^2=m_{pl}^3/m^2`$. For the purpose of a quantitative analysis, exact masses and detailed mass density profiles of FDFSs were examined from nonrelativistic to relativistic situation, using the formalism of Oppenheimer and Volkoff.
These results were applied to the cluster of galaxies, A1689, whose mass distribution has been observationally obtained. We converted the observed column density profile to a volume density profile assuming spherical symmetry, and compared the observed and model 3D encircled mass profiles. We found that the flat-top nature of the observed profile is reproduced by the model and the particle mass range is between 2 eV and 30 eV depending on the actual scale of the degenerate structure.
For about black holes, our scenario will provide alternative mechanism of the black hole formation. Most of the present theories assume the coalescence of stellar mass black holes in the gravitational potential. Such processes seem to be quite complex compared to the FDFS scenario, and therefore it would be much difficult to explain, for example, the observed universal relation between the black hole mass at the center of a galaxy and the bulge mass: $`M_{BH}/M_{bul}0.002.`$ This point will be discussed in detail in our future work.
We thank Nobuo Arimoto for comments on the manuscript. This study was motivated by a very interesting talk given by Tom Broadhurst in 2004 at NAOJ, Mitaka, Japan.
|
warning/0506/math0506144.html
|
ar5iv
|
text
|
# 1. Deformation theory
## 1. Deformation theory
### 1.1. Formal deformations
The most general approach to the question “how to deform $`A_0`$?” is the theory of formal deformations.
Let $`k`$ be a field and $`K:=k[[\mathrm{}_1,\mathrm{},\mathrm{}_{\mathrm{}}]]`$ the ring of formal power series in variables $`\mathrm{}_i`$. Let $`𝔪`$ be the maximal ideal in $`K`$.
A $`K`$-module $`M`$ is said to be topologically free if it is isomorphic to $`M_0[[\mathrm{}_1,\mathrm{},\mathrm{}_{\mathrm{}}]]`$ for some vector space $`M_0`$.
Let $`A_0`$ be an algebra over $`k`$.<sup>2</sup><sup>2</sup>2By “an algebra” we always mean an associative algebra with unit.
###### Definition 1.1.
An $`\mathrm{}`$-parameter flat formal deformation of $`A_0`$ is an algebra $`A`$ over $`K`$ which is topologically free as a $`K`$-module, together with an isomorphism of algebras $`\varphi :A/𝔪A_0`$.<sup>3</sup><sup>3</sup>3The word “flat” refers to the fact that $`A`$ is a (topologically) flat module over $`K`$, i.e. the functor of completed tensor product with this module is exact.
For simplicity we will mostly consider 1-parameter deformations. If $`A`$ is a 1-parameter flat formal deformation of $`A_0`$ then we can choose an identification $`AA_0[[\mathrm{}]]`$ as $`K`$-modules, which reduces to $`\varphi `$ modulo $`\mathrm{}`$. Then the algebra structure on $`A`$ transforms into a new $`K`$-linear multiplication law $`\mu `$ on $`A_0[[\mathrm{}]]`$. Such a multiplication law is determined by the product $`\mu (a,b)`$, $`a,bA_0A_0[[\mathrm{}]]`$, which is given by the formula
$$\mu (a,b)=\mu _0(a,b)+\mathrm{}\mu _1(a,b)+\mathrm{}^2\mu _2(a,b)+\mathrm{},a,bA_0,$$
where $`\mu _i:A_0A_0A_0`$ are linear maps, and $`\mu _0(a,b)`$ is the undeformed product $`ab`$ in $`A_0`$. Thus, to find formal deformations of $`A_0`$ means to find all such series $`\mu `$ which satisfy the associativity equation, modulo the automorphisms of the $`K`$-module $`A_0[[\mathrm{}]]`$ which are the identity modulo $`\mathrm{}`$. <sup>4</sup><sup>4</sup>4Note that we don’t have to worry about the existence of a unit in $`A`$ since a flat formal deformation of an algebra with unit always has a unit
The associativity equation $`\mu (\mu Id)=\mu (Id\mu )`$ reduces to a hierarchy of linear equations:
(1)
$$\underset{s=0}{\overset{N}{}}\mu _s(\mu _{Ns}(a,b),c)=\underset{s=0}{\overset{N}{}}\mu _s(a,\mu _{Ns}(b,c)).$$
(These equations are linear in $`\mu _N`$ if $`\mu _i`$, $`i<N`$, are known).
### 1.2. Hochschild cohomology
These equations can be analyzed using Hochschild cohomology. Let us recall its definition. Let $`M`$ be a bimodule over $`A_0`$. A Hochschild $`n`$-cochain of $`A_0`$ with coefficients in $`M`$ is a linear map $`A_0^nM`$. The space of such cochains is denoted by $`C^n(A_0,M)`$. The differential $`d:C^n(A_0,M)C^{n+1}(A_0,M)`$ is defined by the formula
$$df(a_1,\mathrm{},a_{n+1})=f(a_1,\mathrm{},a_n)a_{n+1}f(a_1,\mathrm{},a_na_{n+1})$$
$$+f(a_1,a_2a_3,\mathrm{},a_{n+1})\mathrm{}+(1)^nf(a_1a_2,\mathrm{},a_{n+1})+(1)^{n+1}a_1f(a_2,\mathrm{},a_{n+1}).$$
It is easy to show that $`d^2=0`$, and one defines the Hochschild cohomology $`H^{}(A_0,M)`$ to be the cohomology of the complex $`(C^{}(A_0,M),d)`$. If $`M=A_0`$, the algebra itself, then we will denote $`H^{}(A_0,M)`$ by $`H^{}(A_0)`$ (it is an algebra). For example, $`H^0(A_0)`$ is the center of $`A_0`$, and $`H^1(A_0)`$ is the quotient of the Lie algebra of derivations of $`A_0`$ by inner derivations.
The following are standard facts from deformation theory (due to Gerstenhaber \[Ge\]), which can be checked directly.
1. The linear equation for $`\mu _1`$ says that $`\mu _1`$ is a Hochschild 2-cocycle. Thus algebra structures on $`A_0[\mathrm{}]/\mathrm{}^2`$ deforming $`\mu _0`$ are parametrized by the space $`Z^2(A_0)`$ of Hochschild 2-cocycles of $`A_0`$ with values in $`M=A_0`$.
2. If $`\mu _1,\mu _1^{}`$ are two 2-cocycles such that $`\mu _1\mu _1^{}`$ is a coboundary, then the algebra structures on $`A_0[\mathrm{}]/\mathrm{}^2`$ corresponding to $`\mu _1`$ and $`\mu _1^{}`$ are equivalent by a transformation of $`A_0[\mathrm{}]/\mathrm{}^2`$ that equals the identity modulo $`\mathrm{}`$, and vice versa. Thus equivalence classes of multiplications on $`A_0[\mathrm{}]/\mathrm{}^2`$ deforming $`\mu _0`$ are parametrized by the cohomology $`H^2(A_0)`$.
3. The linear equation for $`\mu _N`$ says that $`d\mu _N`$ is a certain quadratic expression $`b_N`$ in $`\mu _0,\mu _1,\mathrm{},\mu _{N1}`$. This expression is always a Hochschild 3-cocycle, and the equation is solvable iff it is a coboundary. Thus the cohomology class of $`b_N`$ in $`H^3(A_0)`$ is the only obstruction to solving this equation.
### 1.3. Universal deformation
In particular, if $`H^3(A_0)=0`$ then the equation for $`\mu _n`$ can be solved for all $`n`$, and for each $`n`$ the freedom in choosing the solution, modulo equivalences, is the space $`H:=H^2(A_0)`$. Thus there exists an algebra structure over $`k[[H]]`$ on the space $`A_u:=A_0[[H]]`$ of formal functions from $`H`$ to $`A_0`$, $`a,b\mu _u(a,b)A_0[[H]]`$, ($`a,bA_0`$), such that $`\mu _u(a,b)(0)=abA_0`$, and every 1-parameter flat formal deformation $`A`$ of $`A_0`$ is given by the formula $`\mu (a,b)(\mathrm{})=\mu _u(a,b)(\gamma (\mathrm{}))`$ for a unique formal series $`\gamma \mathrm{}H[[\mathrm{}]]`$, with the property that $`\gamma ^{}(0)`$ is the cohomology class of the cocycle $`\mu _1`$.
Such an algebra $`A_u`$ is called a universal deformation of $`A_0`$. It is unique up to an isomorphism.
Thus in the case $`H^3(A_0)=0`$, deformation theory allows us to completely classify 1-parameter flat formal deformations of $`A_0`$. In particular, we see that the “moduli space” parametrizing formal deformations of $`A_0`$ is a smooth space – it is the formal neighborhood of zero in $`H`$.
### 1.4. Quantization of Poisson structures
If $`H^3(A_0)`$ is nonzero then in general the universal deformation parametrized by $`H`$ does not exist, as there are obstructions to deformations. In this case, the moduli space of deformations will be a closed subscheme of $`H`$, which is often singular. On the other hand, even when $`H^3(A_0)0`$, the universal deformation parametrized by $`H`$ may exist (although it may be more difficult to prove than in the vanishing case). In this case one says that the deformations of $`A_0`$ are unobstructed (since all obstructions vanish even though the space of obstructions doesn’t).
To illustrate these statements, consider the quantization theory of Poisson manifolds. Let $`M`$ be a smooth $`C^{\mathrm{}}`$-manifold or a smooth affine algebraic variety over $``$, and $`A_0`$ the structure algebra of $`M`$.
Remark. In the $`C^{\mathrm{}}`$-case, we will consider only local maps $`A_0^nA_0`$, i.e. those given by polydifferential operators, and all deformations and the Hochschild cohomology is defined using local, rather than general, cochains.
###### Theorem 1.2.
(Hochschild-Kostant-Rosenberg) \[HKR\] $`H^i(A_0)=\mathrm{\Gamma }(M,^iTM)`$ as a module over $`A_0=H^0(A_0)`$.
In particular, $`H^2`$ is the space of bivector fields, and $`H^3`$ the space of trivector fields. So the cohomology class of $`\mu _1`$ is a bivector field; in fact, it is $`\pi (a,b):=\mu _1(a,b)\mu _1(b,a)`$, since any 2-coboundary in this case is symmetric. The equation for $`\mu _2`$ says that $`d\mu _2`$ is a certain trivector field that depends quadratically on $`\pi `$. It is easy to show that this is the Schouten bracket $`[\pi ,\pi ]`$. Thus, for the existence of $`\mu _2`$ it is necessary that $`[\pi ,\pi ]=0`$, i.e. that $`\pi `$ be a Poisson bracket.
Suppose now that $`\pi `$ is a Poisson bracket, i.e. $`[\pi ,\pi ]=0`$. In this case the algebra $`A=A_0[[\mathrm{}]]`$ with the product $`\mu `$ is said to be a quantization of $`\pi `$, and $`(M,\pi )`$ the quasiclassical limit of $`(A,\mu )`$. So, is it possible to construct a quantization of $`\pi `$?
By the above arguments, $`\mu _2`$ exists (and a choice of $`\mu _2`$ is unique up to adding an arbitrary bivector). So there arises the question of existence of $`\mu _3`$ etc., i.e. the question whether there are other obstructions.
The answer to this question is yes and no. Namely, if you don’t pick $`\mu _2`$ carefully, you may be unable to find $`\mu _3`$, but you can always pick $`\mu _2`$ so that $`\mu _3`$ exists, and there is a similar situation in higher orders. This subtle fact is a consequence of the following deep theorem of Kontsevich:
###### Theorem 1.3.
\[K\] Any Poisson structure $`\pi `$ on $`A_0`$ can be quantized. Moreover, there is a natural bijection between products $`\mu `$ up to an isomorphism and Poisson brackets $`\pi _0+\mathrm{}\pi _1+\mathrm{}^2\pi _2+\mathrm{}`$, such that the quasiclassical limit of $`\mu `$ is $`\pi _0`$.
Remark. Note that, as was shown by O. Mathieu, a Poisson bracket on a general commutative $``$-algebra may fail to admit a quantization.
Let us consider the special case of symplectic manifolds, i.e. the case when $`\pi `$ is a nondegenerate bivector. In this case we can consider $`\pi ^1=\omega `$, which is a closed, nondegenerate 2-form (=symplectic structure) on $`M`$. In this case, Kontsevich’s theorem is easier, and was proved by De Wilde - Lecomte, and later Deligne and Fedosov (see e.g. \[F\]). Moreover, in this case there is the following additional result, also due to Kontsevich, \[K\].
###### Theorem 1.4.
If $`M`$ is symplectic and $`A`$ is a quantization of $`M`$, then the Hochschild cohomology $`H^i(A[\mathrm{}^1])`$ is isomorphic to $`H^i(M,((\mathrm{})))`$.
Remark. Here the algebra $`A[\mathrm{}^1]`$ is regarded as a (topological) algebra over the field of Laurent series $`((\mathrm{}))`$, so Hochschild cochains are, by definition, linear maps $`A_0^nA_0((\mathrm{}))`$.
###### Example 1.5.
The algebra $`B=A[\mathrm{}^1]`$ provides an example of an algebra with possibly nontrivial $`H^3(B)`$, for which the universal deformation parametrized by $`H=H^2(B)`$ exists. Namely, this deformation is attached through the correspondence of Theorem 1.3 (and inversion of $`\mathrm{}`$) to the Poisson bracket $`\pi =(\omega +t_1\omega _1+\mathrm{}+t_r\omega _r)^1`$, where $`\omega _1,\mathrm{},\omega _r`$ are closed 2-forms on $`M`$ which represent a basis of $`H^2(M,)`$, and $`t_1,\mathrm{},t_r`$ are the coordinates on $`H`$ corresponding to this basis.
### 1.5. Examples
###### Example 1.6.
Let $`V`$ be a symplectic vector space over $``$ with symplectic form $`\omega `$. Let $`\mathrm{Weyl}(V)`$ denote the Weyl algebra of $`V`$, which is the quotient of the free (=tensor) algebra on $`V`$ by the ideal generated by elements $`xyyx\omega (x,y)`$.
Let $`G`$ be a finite group acting symplectically on $`V`$. Then $`G`$ acts on $`\mathrm{Weyl}(V)`$, and one can form a semidirect product algebra $`A_0=G\mathrm{Weyl}(V)`$. Let us study deformations of $`A_0`$.
We say that an element $`gG`$ is a symplectic reflection in $`V`$ if $`\mathrm{rank}(g1)|_V=2`$. Let $`S`$ be the set of symplectic reflections in $`G`$.
###### Proposition 1.7.
\[AFLS\] $`H^i(A_0)`$ is the space of functions on the set of conjugacy classes of elements $`gG`$ such that $`\mathrm{rank}(g1)|_V=i`$. In particular, $`H^i(A_0)=0`$ if $`i`$ is odd, and $`H^2(A_0)=[S]^G`$.
###### Corollary 1.8.
There exists a universal deformation $`A_u=_c(V,G)`$ of $`A_0`$, which is parametrized by $`c[S]^G`$.
The algebra $`_c(V,G)`$ is called the symplectic reflection algebra (see \[EG\]). Such algebras were first considered by Drinfeld in 1986. If $`V=𝔥𝔥^{}`$, where $`𝔥`$ is a representation of $`G`$, and the symplectic form on $`G`$ is the pairing between $`𝔥`$ and $`𝔥^{}`$, then $`_c(V,G)`$ is called the rational Cherednik algebra. We will later construct $`_c(V,G)`$ explicitly.
###### Example 1.9.
Let $`X`$ be a smooth affine algebraic variety over $``$, with an action of a finite group $`G`$. Let $`D(X)`$ be the algebra of algebraic differential operators on $`X`$. Let $`A_0=GD(X)`$. Let us study deformations of $`A_0`$.
For every $`gG`$, the fixed set $`X^g`$ of $`g`$ in $`Y`$ is a smooth affine variety, which consists of connected components $`X_j^g`$, possibly of different dimensions. Such a component is said to be a reflection hypersurface if it has codimension $`1`$ in $`X`$. Let $`S`$ be the set of pairs $`(g,Y)`$, where $`gG`$, and $`YX^g`$ is a connected component which is a reflection hypersurface (i.e., has codimension 1).
###### Proposition 1.10.
\[E\] One has $`H^2(A_0)=(H^2(X,)[S])^G`$. Moreover, there exists a universal deformation of $`A_0`$ parametrized by $`H=H^2(A_0)`$.
This deformation $`_c[X,G]`$ is called the rational Cherednik algebra attached to $`(X,G)`$, and is described in \[E\]. If $`X`$ is a vector space $`𝔥`$ and $`G`$ acts linearly, then $`_c[𝔥,G]=_c(𝔥𝔥^{},G)`$ is the rational Cherednik algebra discussed above.
###### Example 1.11.
The following example from the paper \[DE\] (conjecturally) generalizes examples 1.5,1.6, and 1.9.
Let $`M`$ be a symplectic $`C^{\mathrm{}}`$-manifold (or affine complex algebraic variety). Let $`G`$ be a finite group acting on $`M`$ by symplectic transformations, and $`B`$ be a quantization of $`M`$ which is equivariant under $`G`$ (such a quantization always exists). Let $`A_0=GB[\mathrm{}^1]`$. Let us study deformations of $`A_0`$.
The Hochschild cohomology of $`A_0`$ is given by the following theorem. Let the fixed set $`M^g`$ be the union of connected components $`M_i^g`$, $`i=1,\mathrm{},N_g`$.
###### Theorem 1.12.
$`H^{}(A_0)`$ equals, as a vector space, the orbifold cohomology of $`M/G`$ with coefficients in $`((\mathrm{}))`$. Namely,
$$H^p(A_0)=(_{gG}_{i=1}^{N_g}H^{p\mathrm{codim}M_i^g}(M_i^g))^G.$$
(where the coefficients on the RHS are $`((\mathrm{}))`$).
Remark. Let $`S`$ be the set of pairs $`(g,Y)`$, where $`gG`$, and $`YM^g`$ is a connected component of codimension $`2`$. Theorem 1.12 implies that $`H^2(A_0)=(H^2(M)[S])^G`$.
Thus, we see that $`H^3(A_0)`$ does not always vanish. Nevertheless, we make the following conjecture.
###### Conjecture 1.13.
The deformations of the algebra $`A_0`$ are unobstructed. Thus there exists a universal deformation $`H_c`$ of this algebra parametrized by $`cH^2(A_0)`$.
Thus the conjecture implies that if $`S\mathrm{}`$, then there exist “interesting” deformations of $`A_0`$, i.e., ones not coming from $`G`$-invariant deformations of $`B`$.
Let us give a few examples in which this conjecture is true.
1. $`H^3(A_0)=0`$. This includes the following interesting case considered in \[EO\]: $`\mathrm{\Sigma }`$ is a smooth affine algebraic surface such that $`H^1(\mathrm{\Sigma },)=0`$, and $`M=\mathrm{\Sigma }^n`$, $`G=S_n`$. In this case there is one interesting deformation parameter corresponding to reflections in $`S_n`$.
2. $`G`$ is trivial (Example 1.5).
3. $`M=T^{}Y`$, where $`Y`$ is a smooth affine variety, and $`G`$ acts on $`Y`$ (Example 1.9).
4. If $`M=V`$ is a symplectic vector space and $`G`$ acts linearly (Example 1.6).
5. Let $`M=V/L`$, where $`V`$ is a symplectic vector space and $`L`$ a lattice in $`V`$ (i.e., $`L`$ is the abelian group generated by a basis of $`V`$). Thus $`M`$ is an algebraic torus with a symplectic form. We assume that the symplectic form is integral and unimodular on $`L`$. Let $`GSp(L)`$ be a finite subgroup; then $`G`$ acts naturally on $`M`$. In this case $`H_c`$ is an “orbifold Hecke algebra” defined in \[E\] (it will be discussed below).
## 2. Algebras given by generators and relations
### 2.1. Giving formal deformations by generators and relations
Another approach to exploring deformations of $`A_0`$ is defining deformations by generators and relations.
Let us first consider the setting of formal deformations, which we have discussed in the previous section. Namely, let $`A_0`$ be an algebra over a field $`k`$, generated by $`a_1,a_2,\mathrm{}`$ with defining relations $`R_j^0(a_1,a_2,\mathrm{})=0`$ (here $`R_j^0`$ are elements in the free $`k`$-algebra $`F`$ generated by $`a_i`$). Let us now define a formal deformation of $`A_0`$ as the algebra over $`K=k[[\mathrm{}]]`$ with the same generators and deformed relations $`R_j=R_j^0+\mathrm{}R_j^1+\mathrm{}^2R_j^2+\mathrm{}`$. That is, $`A`$ is the quotient of the free algebra $`F[[\mathrm{}]]`$ by the $`\mathrm{}`$-adically closed ideal generated by the relations $`R_j`$.
###### Example 2.1.
(The Weyl algebra.) Let $`A_0=[x,y]`$ be the algebra generated by $`x,y`$ with the defining relation $`yxxy=0`$. We can then define $`A`$ by the same generators and the deformed relation $`yxxy=\mathrm{}`$ (the Heisenberg indeterminacy relation). Then $`A`$ is indeed a 1-parameter flat formal deformation of $`A_0`$, which provides a quantization of the standard Poisson bracket $`\{y,x\}=1`$.
So, is $`A`$ always a 1-parameter flat formal deformation of $`A_0`$? In general the answer is no: the flatness property can fail. The following typical example of this is obtained by adding just one relation to the relations above.
###### Example 2.2.
Assume the algebra $`A_0`$ is defined by generators $`x,y`$ and defining relations
$$yxxy=0,x=0,$$
and $`A`$ is defined by generators $`x,y`$ and relations
$$yxxy=\mathrm{},x=0.$$
Then $`A`$ is not topologically free, as it contains $`\mathrm{}`$-torsion. Indeed, $`\mathrm{}1=yxxy=0`$ since $`x=0`$. On the other hand, $`10`$, since the algebra $`A_0=[y]`$ is nonzero.
In fact, it is easy to show that if we add any relation to $`xyyx=\mathrm{}`$, it will produce a non-flat deformation (unless the algebra to be deformed is zero to begin with). This shows that if one wants to secure flatness, one has to deform the relations in a very special way. In fact, it is usually rather difficult to do so, as well as to check that the resulting deformations are actually flat. Below I would like to show several situations when this task can be successfully completed.
### 2.2. Deformations of quadratic algebras
The first situation is deformation theory of quadratic algebras.
Let $`R`$ be a finite dimensional semisimple algebra (say over $``$). Let $`A`$ be a $`_+`$-graded algebra, $`A=_{i0}A[i]`$, such that $`A[0]=R`$. For simplicity assume that the spaces $`A[i]`$ are finite dimensional for all $`i`$.
###### Definition 2.3.
(i) The algebra $`A`$ is said to be quadratic if it is generated over $`R`$ by $`A[1]`$, and has defining relations in degree 2.
(ii) $`A`$ is Koszul if all elements of $`Ext^i(R,R)`$ (where $`R`$ is the augmentation module over $`A`$) have grade degree precisely $`i`$.
Remarks. 1. Thus, in a quadratic algebra, $`A[2]=A[1]_RA[1]/E`$, where $`E`$ is the subspace ($`R`$-subbimodule) of relations.
2. It is easy to show that a Koszul algebra is quadratic, since the condition to be quadratic is just the Koszulity condition for $`i=1,2`$.
3. Many important algebras, e.g. the free algebra, the polynomial algebra and the exterior algebra are Koszul.
Now let $`A_0`$ be a quadratic algebra, $`A_0[0]=R`$. Let $`E_0`$ be the space of relations for $`A_0`$. Let $`EA_0[1]_RA_0[1][[\mathrm{}]]`$ be a topologically free (over $`[[\mathrm{}]]`$) $`R`$-subbimodule which reduces to $`E_0`$ modulo $`\mathrm{}`$ (“deformation of the relations”). Let $`A`$ be the ($`\mathrm{}`$-adically complete) algebra generated over $`R[[\mathrm{}]]`$ by $`A[1]=A_0[1][[\mathrm{}]]`$ with the space of defining relations $`E`$. Thus $`A`$ is a $`_+`$-graded algebra.
Then we have the following fundamental result
###### Theorem 2.4.
(Koszul deformation principle,\[D\],\[BG\],\[PP\],\[BGS\]) If $`A_0`$ is Koszul then $`A`$ is a topologically free $`[[\mathrm{}]]`$ module if and only if so is $`A[3]`$.
Remark. Note that $`A[i]`$ for $`i<3`$ are obviously topologically free.
### 2.3. Symplectic reflection algebras.
We will now demonstate by an example how the Koszul deformation principle works.
Let $`V`$ be a finite dimensional symplectic vector space over $``$ with a symplectic form $`\omega `$, and $`G`$ be a finite group acting symplectically on $`V`$. For simplicity let us assume that $`(^2V)^G=\omega `$.
If $`sG`$ is a symplectic reflection, then let $`\omega _s(x,y)`$ be the form $`\omega `$ applied to the projections of $`x,y`$ to the image of $`1s`$ along the kernel of $`1s`$; thus $`\omega _s`$ is a skewsymmetric form of rank $`2`$ on $`V`$.
Let $`SG`$ be the set of symplectic reflections, and $`c:S`$ be a function which is invariant under the action of $`G`$. Let $`t`$.
###### Definition 2.5.
The symplectic reflection algebra $`H_{t,c}=H_{t,c}(V,G)`$ is the quotient of the algebra $`G𝕋(V)`$ by the ideal generated by the relation
(2)
$$[x,y]=t\omega (x,y)2\underset{sS}{}c_s\omega _s(x,y)s.$$
The following theorem shows that the algebras $`H_{t,c}(V,G)`$ satisfy a flatness property, and moreover, they are the only ones satisfying this property within a certain natural class.
###### Theorem 2.6.
Let $`\kappa :^2V[G]`$ be a $`G`$-equivariant function ($`G`$ acts on the target by conjugation). Define the algebra $`H_\kappa `$ to be the quotient of the algebra $`G𝕋(V)`$ by the relation $`[x,y]=\kappa (x,y)`$, $`x,yV`$. Put an increasing filtration on $`H_\kappa `$ by setting $`\mathrm{deg}(V)=1`$, $`\mathrm{deg}(G)=0`$, and define $`\xi :GSV\mathrm{gr}H_\kappa `$ to be the natural surjective homomorphism. Then $`\xi `$ is an isomorphism if and only if $`\kappa `$ has the form
$$\kappa (x,y)=t\omega (x,y)2\underset{sS}{}c_s\omega _s(x,y)s,$$
for some $`t`$ and $`G`$-invariant function $`c:S`$.
Before proving this theorem, let us point out a corollary. Denote by $`_c(V,G)`$ the algebra defined as $`H_{t,c}(V,G)`$, but with $`t=1`$ and $`c`$ being a formal parameter.
###### Corollary 2.7.
The algebra $`_c(V,G)`$ is a flat formal deformation of $`G\mathrm{Weyl}(V)`$, parametrized by $`[S]^G`$.
In fact, it turns out (see \[EG\]) that $`_c(V,G)`$ is the universal deformation of $`G\mathrm{Weyl}(V)`$, whose existence was proved in Example 1.6.
###### Proof.
(of Theorem 2.6) Let $`\kappa :^2V[G]`$ be an equivariant map. We write $`\kappa (x,y)=_{gG}\kappa _g(x,y)g`$, where $`\kappa _g(x,y)^2V^{}`$. To apply Theorem 2.4, let us homogenize our algebras. Namely, let $`A_0=(GSV)[u]`$. Also let $`\mathrm{}`$ be a formal parameter, and consider the deformation $`A=H_{\mathrm{}u^2\kappa }`$ of $`A_0`$. That is, $`A`$ is the quotient of $`G𝕋(V)[u][[\mathrm{}]]`$ by the relations $`[x,y]=\mathrm{}u^2\kappa (x,y)`$. This is a deformation of the type considered in Theorem 2.4, and it is easy to see that its flatness in $`\mathrm{}`$ is equivalent to Theorem 2.6. Also, the algebra $`A_0`$ is Koszul, because the polynomial algebra $`SV`$ is a Koszul algebra. Thus by Theorem 2.4, it suffices to show that $`A`$ is flat in degree 3.
The flatness condition in degree 3 is “the Jacobi identity”
$$[\kappa (x,y),z]+[\kappa (y,z),x]+[\kappa (z,x),y]=0,$$
which must be satisfied in $`GV`$. In components, this equation transforms into the system of equations
$$\kappa _g(x,y)(zz^g)+\kappa _g(y,z)(xx^g)+\kappa _g(z,x)(yy^g)=0$$
for every $`gG`$ (here $`z^g`$ denotes the result of the action of $`g`$ on $`z`$).
This equation, in particular, implies that if $`x,y,g`$ are such that $`\kappa _g(x,y)0`$ then for any $`zV`$ $`zz^g`$ is a linear combination of $`xx^g`$ and $`yy^g`$. Thus $`\kappa _g(x,y)`$ is identically zero unless the rank of $`(1g)|_V`$ is at most 2, i.e. $`g=1`$ or $`g`$ is a symplectic reflection.
If $`g=1`$ then $`\kappa _g(x,y)`$ has to be $`G`$-invariant, so it must be of the form $`t\omega (x,y)`$, where $`t`$.
If $`g`$ is a symplectic reflection, then $`\kappa _g(x,y)`$ must be zero for any $`x`$ such that $`xx^g=0`$. Indeed, if for such an $`x`$ there had existed $`y`$ with $`\kappa _g(x,y)0`$ then $`zz^g`$ for any $`z`$ would be a multiple of $`yy^g`$, which is impossible since $`Im(1g)|_V`$ is 2-dimensional. This implies that $`\kappa _g(x,y)=2c_g\omega _g(x,y)`$, and $`c_g`$ must be invariant.
Thus we have shown that if $`A`$ is flat (in degree 3) then $`\kappa `$ must have the form given in Theorem 2.6. Conversely, it is easy to see that if $`\kappa `$ does have such form, then the Jacobi identity holds. So Theorem 2.6 is proved. ∎
### 2.4. Deformation of representations
Another method of estabishing flatness of a deformation $`A`$ of $`A_0`$ defined by generators and relations is showing that a given faithful representation $`M_0`$ of the algebra $`A_0`$ (for example, the regular representation) can be deformed (flatly) to a representation $`M`$ of $`A`$. In this case it follows automatically that $`A`$ is flat. Let us give two examples of situations where this method can be applied.
###### Example 2.8.
(see \[E\]). Let $`X`$ be a connected, simply connected complex manifold, and $`G`$ a discrete group of automorphisms of $`X`$. In this case the quotient $`X/G`$ is a complex orbifold. Let $`X^{}X`$ be the set of points having trivial stabilizer (it is a nonempty open subset of $`X`$). Define the braid group $`\stackrel{~}{G}`$ of the orbifold $`X/G`$ to be the fundamental group of the manifold $`X^{}/G`$ with some base point $`x_0`$. We have a surjective homomorphism $`\varphi :\stackrel{~}{G}G`$, which corresponds to gluing back the points which have a nontrivial stabilizer. Let $`K`$ be the kernel of this homomorphism.
The kernel $`K`$ can be described by simple relations, corresponding to reflection hypersurfaces in $`X`$. Namely, given a reflection hypersurface $`YX`$, we have a conjugacy class $`C_Y`$ in $`\stackrel{~}{G}`$ which corresponds to the loop in $`X^{}/G`$ which goes counterclockwise around $`Y`$. Let $`T_Y`$ be a representative of $`C_Y`$. Also, let $`G_YG`$ be the stabilizer of a generic point on $`Y`$; this is a cyclic group of some order $`n_Y`$. Then it follows from basic topology that the elements $`T_Y^{n_Y}`$ belong to $`K`$, and $`K`$ is the smallest normal subgroup of $`\stackrel{~}{G}`$ containing all of them. In other words, the group $`G`$ is the quotient of the braid group $`\stackrel{~}{G}`$ by the relations
(3)
$$T_Y^{n_Y}=1.$$
Now let $`A_0=[G]`$, and let us define a deformation $`A`$ of $`A_0`$ to be the quotient of the group algebra of the braid group $`\stackrel{~}{G}`$ by a deformation of relations (3). Namely, for every reflection hypersurface $`YX`$ we introduce formal parameters $`\tau _{Y,j}`$, $`j=1,\mathrm{},n_Y`$ (which are conjugation invariant), and replace relations (3) by the relations
(4)
$$\underset{j=1}{\overset{n_Y}{}}(T_Ye^{2\pi ij/n_Y}e^{\tau _{Y,j}})=0.$$
The quotient $`A`$ of $`[\stackrel{~}{G}][[\tau ]]`$ by these relations is called the orbifold Hecke algebra of $`X/G`$, and denoted by $`_\tau (X,G)`$.
###### Theorem 2.9.
(\[E\]) If $`H^2(X,)=0`$ then $`_\tau (X,G)`$ is a flat deformation of $`[G]`$.
Remark. If $`X`$ is $`^n`$ and $`G=G_0L`$, where $`L`$ is a lattice of rank $`2n`$ and $`G_0`$ is a finite group acting on $`L`$ then $`_\tau (X,G)`$ is, essentially, the algebra which was mentioned in Example 1.11.
To illustrate the relevance of the condition $`H^2(X,)=0`$, let us consider the special case when $`G`$ is the triangle group $`F_{p,q,r}`$, generated by $`a,b,c`$ with defining relations
$$a^p=1,b^q=1,c^r=1,abc=1,$$
where $`p,q,r>1`$ are positive integers. The group $`G`$ is the group generated by rotations around the vertices of a triangle with angles $`\pi /p,\pi /q,\pi /r`$, by twice the angle at the vertex. Let $`S=\frac{1}{p}+\frac{1}{q}+\frac{1}{r}`$. The triangle lies on the sphere, Euclidean plane, or hyperbolic plane $`X`$ when $`S>1`$, $`S=1`$, and $`S<1`$, respectively. The deformation $`_\tau (X,G)`$ is generated by $`a,b,c`$ with defining relations
$$\underset{j=1}{\overset{p}{}}(a\alpha _j)=0,\underset{j=1}{\overset{q}{}}(b\beta _j)=0,\underset{j=1}{\overset{r}{}}(c\gamma _j)=0,abc=1,$$
where
$$\alpha _j=e^{2\pi ij/p}e^{\tau _{1j}},\beta _j=e^{2\pi ij/q}e^{\tau _{2j}},\gamma _j=e^{2\pi ij/r}e^{\tau _{3j}}.$$
Theorem 2.9 says that the deformation is flat for the Euclidean and hyperbolic plane, but says nothing about the sphere, i.e. the triples $`(p,q,r)`$ equal to $`(2,2,n)`$, $`(2,3,3),(2,3,4),(2,3,5)`$, in which case the group $`G`$ is finite. And indeed, in this case $`_\tau (X,G)`$ is actually not flat! To see this, note that in the sphere case $`_\tau `$, if it were flat, would have dimension $`|G|`$ (over $`[[\tau ]]`$). So we may take the determinant of the relation $`abc=1`$ (using the fact that the eigenvalues of $`a,b,c`$ are $`\alpha _j,\beta _j,\gamma _j`$, with equal multiplicities). This yields a nontrivial relation on $`\tau `$:
$$(\underset{j=1}{\overset{p}{}}\alpha _j)^{|G|/p}(\underset{j=1}{\overset{q}{}}\beta _j)^{|G|/q}(\underset{j=1}{\overset{r}{}}\gamma _j)^{|G|/r}=1,$$
which rules out flatness of $`_\tau `$.
###### Example 2.10.
(\[ER\]) Let $`W`$ be a Coxeter group of rank $`r`$ with generators $`s_i`$ and defining relations
$$s_i^2=1,(s_is_j)^{m_{ij}}=1\text{ for }m_{ij}<\mathrm{},i,j=1,\mathrm{},r,ij,$$
where $`m_{ij}=m_{ji}`$ are integers $`2`$ or $`\mathrm{}`$, defined for $`ij`$. Let $`W_+`$ be the subgroup of even elements of $`W`$. It is easy to see that $`W_+`$ is generated by the elements $`a_{ij}:=s_is_j`$, with defining relations
$$a_{ij}a_{ji}=1,a_{ij}a_{jk}a_{ki}=1,a_{ij}^{m_{ij}}=1.$$
Define a deformation of $`A_0=[W_+]`$ as follows. Introduce invertible parameters $`t_{ij,k}=t_{ji,k}^1`$, $`k/m_{ij}`$ for $`m_{ij}<\mathrm{}`$. Let $`R=[t_{ij,k}]`$, and $`A`$ be the $`R`$-algebra generated by $`a_{ij}`$ with defining relations
$$a_{ij}a_{ji}=1,a_{ij}a_{jk}a_{ki}=1,\underset{k=1}{\overset{m_{ij}}{}}(a_{ij}t_{ij,k})=0.$$
For any $`xW_+`$, fix a reduced word $`w(x)`$ representing $`x`$. Let $`T_{w(x)}`$ be the element of $`A`$ corresponding to this word.
###### Theorem 2.11.
\[ER\] (i) The elements $`T_{w(x)}`$ for $`xW_+`$ span $`A`$ over $`R`$.
(ii) These elements form a basis of $`A`$ over $`R`$ if and only if $`W`$ has no finite parabolic subgroups of rank $`3`$, i.e. iff for each $`i,j,l`$,
$$\frac{1}{m_{ij}}+\frac{1}{m_{jl}}+\frac{1}{m_{li}}1.$$
###### Corollary 2.12.
Let $`\widehat{A}`$ be the completion of $`A`$ with respect to the ideal generated by $`t_{ij,k}e^{2\pi k\sqrt{1}/m_{ij}}`$. Then $`\widehat{A}`$ is a flat deformation of $`A_0`$ iff $`W`$ has no finite parabolic subgroups of rank 3.
Remark. Note that triangle groups $`F_{p,q,r}`$ are groups $`W_+`$ for Coxeter groups of rank 3 (with $`m_{12}=p`$, $`m_{23}=q`$, $`m_{31}=r`$), so the “only if” part of Theorem 2.11 (and the “if” part in rank 3) follow from Example 2.8.
In both of these examples, flatness is established by showing, using geometric methods (D-modules or constructible sheaves), that the regular representation of $`A_0`$ can be flatly deformed to a representation of the deformation. Let us conclude by illustrating this in Example 2.8, in the case when $`X=E`$ is a complex vector space, and $`G`$ is a finite group acting linearly on $`E`$. In this case, Theorem 2.9 was proved by Broué, Malle, and Rouquier \[BMR\], following an idea of Cherednik. Let us sketch their proof.
The main idea of the proof is to introduce Dunkl operators $`D_a`$, $`aE`$, which act on functions on $`E`$ (with poles on the reflection hyperplanes $`Y`$):
$$D_a=_a+\underset{Y}{}\frac{\alpha _Y(a)}{\alpha _Y}(\underset{gG_Y}{}c_{Y,g}g),$$
where the summation is over all reflection hyperplanes $`Y`$, $`\alpha _Y`$ is the nonzero element of $`E^{}`$ vanishing on $`Y`$, and $`c_{Y,g}`$ is a conjugation invariant function of $`Y,g`$.
It can be shown that the Dunkl operators commute: $`[D_a,D_b]=0`$. This implies that the system of equations $`D_a\psi =0`$, $`aE`$, can be regarded as a local system with fiber $`G`$ on $`(EY)/G`$. The fundamental group of $`(EY)/G`$, by definition, is $`\stackrel{~}{G}`$, so we may consider the corresponding monodromy representation of this group. If $`c=0`$, the monodromy representation is the standard homomorphism $`\stackrel{~}{G}G`$. One may show that if $`c0`$, then the monodromy representation is a deformation of this standard homomorphism, which factors through the Hecke algebra $`_\tau (E,G)`$, for an appropriate linear change of variables $`c\tau `$. This implies the flatness of $`_\tau (E,G)`$.
|
warning/0506/hep-ph0506043.html
|
ar5iv
|
text
|
# Dihedral Flavor Symmetry from Dimensional Deconstruction
## I Introduction
The Yukawa sector of the standard model (SM) contains a large number of redundant parameters. The presence of the redundant parameters is not related to a symmetry in the SM. That is, they will appear in higher orders in perturbation theory even if they are set equal to zero at the level. These redundant parameters may become physical parameters when going beyond the SM, and, moreover, they can induce flavor changing neutral currents (FCNCs) and CP violating phenomena that are absent or strongly suppressed in the SM. One of the most well known examples is the case of the minimal supersymmetric model (MSSM). Since the SM can not control the redundant parameters, the size of the new FCNCs and CP violating phases may be unacceptably large unless there is some symmetry, or one fine tunes their values <sup>1</sup><sup>1</sup>1For recent reviews, see, for instance, susy and references therein. .
A natural guidance to constrain the Yukawa sector and to reduce the redundancy of this sector is a flavor symmetry. It has been recently realized that nonabelian discrete flavor symmetries, especially dihedral symmetries, can not only reduce the redundancy, but also partly explain the large mixing of neutrinos <sup>2</sup><sup>2</sup>2 Models based on dihedral flavor symmetries, ranging from $`D3(S_3)`$ to $`Q_6`$ and $`D_7`$, have been recently discussed in frampton1 -grimus4 .. When supersymmetrized, it has been found that the same flavor symmetries can suppress FCNCs that are caused by soft supersymmetry breaking terms kobayashi2 ; choi (see also hall2 -king2 ).
In this letter we address the question of the origin of dihedral flavor symmetries. We will find that dimensional deconstruction hamed1 ; hill is a possible origin of dihedral flavor symmetries.
## II Dihedral invariance in an Extra dimensional space
Consider an extra dimension which is compactified on a closed one-dimensional lattice with $`N`$ sites. We assume that the lattice has the form of a regular polygon with $`N`$ edges as it is illustrated in Fig. 1.
The regular polygon is invariant under the symmetry operations of the dihedral group $`D_N`$. The $`D_N`$ operations are $`2N`$ discrete rotations, where $`N`$ of $`2N`$ rotations are combined with a parity transformation. Clearly, a discrete polygon rotation of $`n\times \theta _N,n\{1,\mathrm{},N\}`$ corresponds to a discrete translation of the lattice sites of $`n\times a`$, where $`a`$ is the lattice spacing and
$`\theta _N`$ $`2\pi /N.`$ (1)
The coordinate of the extra dimension is denoted by $`y`$, and the $`N`$ sites are located at $`y=y_0,y_1,\mathrm{},y_{N1}`$. ($`y_{N+i}`$ is identified with $`y_i`$.) Under a $`D_N`$ transformation, the set of coordinates $`(y_0,y_1,\mathrm{},y_{N1})`$ changes to $`(y_0^{},y_1^{},\mathrm{},y_{N1}^{})`$, which we express in terms of a $`N\times N`$ real matrix. The matrix for the fundamental rotation (i.e., a rotation of $`\theta _N`$) is given by
$`R_N`$ $`=`$ $`\left(\begin{array}{ccccc}0& 0& 0& \mathrm{}& 1\\ 1& 0& 0& \mathrm{}& 0\\ 0& 1& 0& \mathrm{}& 0\\ & & \mathrm{}& & \\ 0& 0& \mathrm{}& 1& 0\end{array}\right),`$ (7)
and that for the parity transformation is
$`P_D`$ $`=`$ $`\left(\begin{array}{ccccc}1& 0& \mathrm{}& 0& 0\\ 0& \mathrm{}& 0& 0& 1\\ 0& \mathrm{}& 0& 1& 0\\ & & \mathrm{}& & \\ 0& 1& 0& \mathrm{}& 0\end{array}\right).`$ (13)
Then the $`2N`$ group elements of $`D_N`$ are:
$`𝒢_{D_N}=\{R_N,(R_N)^2,\mathrm{},(R_N)^N=\mathrm{𝟏},R_NP_D,(R_N)^2P_D,\mathrm{},(R_N)^NP_D=P_D\}.`$ (14)
Using the properties, $`P_D^2=\mathrm{𝟏}`$ and $`P_DR_NP_D=(R_N)^1`$, one can convince oneself that $`𝒢_{D_N}`$ is indeed a group.
There exist two-dimensional representations for $`\stackrel{~}{R}_N`$ and $`\stackrel{~}{P}_D`$ frampton1 ; babu4 :
$`\stackrel{~}{R}_N`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta _N& \mathrm{sin}\theta _N\\ \mathrm{sin}\theta _N& \mathrm{cos}\theta _N\end{array}\right),\stackrel{~}{P}_D=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (19)
which are useful representations in finding irreducible representations (irreps) of $`D_N`$ ($`\theta _N`$ is given in (1)). It follows that $`D_N`$ is a subgroup of $`SO(3)`$, which one sees if one embeds $`\stackrel{~}{R}_N`$ and $`\stackrel{~}{P}_D`$ into $`3\times 3`$ matrices
$`\stackrel{~}{R}_N`$ $``$ $`\left(\begin{array}{ccc}\mathrm{cos}\theta _N& \mathrm{sin}\theta _N& 0\\ \mathrm{sin}\theta _N& \mathrm{cos}\theta _N& 0\\ 0& 0& 1\end{array}\right),\stackrel{~}{P}_D\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).`$ (26)
Therefore, $`D_N`$ has only real representations.
$`SU(2)`$ is the universal covering group of $`SO(3)`$, and has pseudo real and real irreps. $`Q_{2N}`$ is a finite subgroup of $`SU(2)`$. It can be interpreted as the covering group of $`D_N`$ in the sense that the defining matrices $`\stackrel{~}{R}_{2N}`$ and $`\stackrel{~}{P}_Q`$ for $`Q_{2N}`$ satisfy
$`(\stackrel{~}{R}_{2N})^2=\stackrel{~}{R}_N,(\stackrel{~}{P}_Q)^4=(\stackrel{~}{P}_D)^2=\mathrm{𝟏},`$ (27)
where
$`\stackrel{~}{R}_{2N}`$ $`=`$ $`\left(\begin{array}{ccc}\mathrm{cos}\frac{\theta _N}{2}& \mathrm{sin}\frac{\theta _N}{2}& \\ \mathrm{sin}\frac{\theta _N}{2}& \mathrm{cos}\frac{\theta _N}{2}& \end{array}\right),\stackrel{~}{P}_Q=\left(\begin{array}{ccc}i& 0& \\ 0& i& \end{array}\right).`$ (32)
The set of $`4N`$ elements of $`Q_{2N}`$ is given by
$`𝒢_{Q_{2N}}=\{\stackrel{~}{R}_{2N},(\stackrel{~}{R}_{2N})^2,\mathrm{},(\stackrel{~}{R}_{2N})^{2N}=\mathrm{𝟏},\stackrel{~}{R}_{2N}\stackrel{~}{P}_Q,(\stackrel{~}{R}_{2N})^2\stackrel{~}{P}_Q,\mathrm{},(\stackrel{~}{R}_{2N})^{2N}\stackrel{~}{P}_Q=\stackrel{~}{P}_Q\}.`$ (33)
There exist only one- and two-dimensional irreps for $`D_N`$ and $`Q_{2N}`$. For $`Q_{2N}`$, there are $`N1`$ different two-dimensional irreps, which we denote by
$`\mathrm{𝟐}_{\mathrm{}},\mathrm{}=1,\mathrm{},N1.`$ (34)
$`\mathrm{𝟐}_{\mathrm{}}`$ with odd $`\mathrm{}`$ is a pseudo real representation, while $`\mathrm{𝟐}_{\mathrm{}}`$ with even $`\mathrm{}`$ is a real representation, where $`\mathrm{𝟐}_{\mathrm{}}`$ with even $`\mathrm{}`$ is exactly $`\mathrm{𝟐}_{\mathrm{}/2}`$ of $`D_N`$. Under the fundamental rotation (i.e., a rotation of $`\theta _N`$ which is defined in (1)), $`\mathrm{𝟐}_{\mathrm{}}`$ transforms with the matrix
$`\stackrel{~}{R}_{2N}(\mathrm{𝟐}_{\mathrm{}})`$ $`=`$ $`(\stackrel{~}{R}_{2N})^{\mathrm{}}=\left(\begin{array}{ccc}\mathrm{cos}(\mathrm{}\frac{\theta _N}{2})& \mathrm{sin}(\mathrm{}\frac{\theta _N}{2})& \\ \mathrm{sin}(\mathrm{}\frac{\theta _N}{2})& \mathrm{cos}(\mathrm{}\frac{\theta _N}{2})& \end{array}\right).`$ (37)
It is straightforward to calculate the Clebsch-Gordan coefficients for tensor products of irreps babu4 . There exist four different one-dimensional irreps of $`Q_{2N}`$. Because of the relation (27), each of them has a definite $`Z_4`$ charge. Further, under the fundamental rotation, they either remain unchanged or change their sign. Therefore, one-dimensional irreps can be characterized according to $`Z_2\times Z_4`$ charge:
$`\mathrm{𝟏}_{+,0},`$ $`\mathrm{𝟏}_{,0},\mathrm{𝟏}_{+,2},\mathrm{𝟏}_{,2}\text{for}N=2,4,6,\mathrm{},`$ (38)
$`\mathrm{𝟏}_{+,0},`$ $`\mathrm{𝟏}_{,1},\mathrm{𝟏}_{+,2},\mathrm{𝟏}_{,3}\text{for}N=3,5,7,\mathrm{},`$ (39)
where the $`\mathrm{𝟏}_{+,0}`$ is the true singlet of $`Q_{2N}`$, and only $`\mathrm{𝟏}_{,1}`$ and $`\mathrm{𝟏}_{,3}`$ are complex irreps. Note that all the real representations of $`Q_{2N}`$ are exactly those of $`D_N`$, which is one of the reasons why we would like to call $`Q_{2N}`$ as the covering group of $`D_N`$.
## III Field Theory with the dihedral invariance
Let us now discuss how to construct field theory models with a dihedral invariance. We denote the five-dimensional coordinate by
$`z^M`$ $`=`$ $`(x^\mu ,y)\text{with}\mu =0,\mathrm{},3.`$ (40)
The coordinates $`y_i`$ of the lattice sites transform to $`y_i^{}`$ with $`N\times N`$ matrices of $`D_N`$, which are given in (7) and (13). Then it is natural to assume <sup>3</sup><sup>3</sup>3$`D_N`$ may be understood as a twisted product of $`Z_N`$ and $`Z_2`$. Witten witten has considered this $`Z_N`$ (the symmetry of the boundary of a deconstructed disc) to solve the triplet-doublet splitting problem in GUTs. that the fields defined on the lattice are irreps of $`Q_{2N}`$ which is the covering group of $`D_N`$. That is <sup>4</sup><sup>4</sup>4Nonabelian discrete family symmetries appearing in extra dimension models of seidl1 ; altarelli , for instance, are not directly related to a symmetry of the extra dimension.,
$`\varphi (x,y)`$ $``$ $`\varphi ^{}(x,y)=\stackrel{~}{Q}_{2N}\varphi (x,\stackrel{~}{D}_N^1y),\stackrel{~}{Q}_{2N}Q_{2N}\text{and}\stackrel{~}{D}_ND_N.`$ (41)
In Table 1 explicit expressions of the matrices corresponding to the fundamental rotation and the parity transformation are given, where we assume that the gauge fields belong to the true singlet $`\mathrm{𝟏}_{+,0}`$.
| Irreps | $`\mathrm{𝟏}_{+,0}`$ | $`\mathrm{𝟏}_{+,2}`$ | $`\mathrm{𝟏}_{,0}`$ | $`\mathrm{𝟏}_{,1}`$ | $`\mathrm{𝟏}_{,2}`$ | $`\mathrm{𝟏}_{,3}`$ | $`\mathrm{𝟐}_{2\mathrm{}1}`$ | $`\mathrm{𝟐}_2\mathrm{}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| rotation | $`1`$ | $`1`$ | $`1`$ | $`1`$ | $`1`$ | $`1`$ | $`(\stackrel{~}{R}_{2N})^{2\mathrm{}1}`$ | $`(\stackrel{~}{R}_{2N})^2\mathrm{}`$ |
| parity | $`1`$ | $`1`$ | $`1`$ | $`i`$ | $`1`$ | $`i`$ | $`\stackrel{~}{P}_Q`$ | $`\stackrel{~}{P}_D`$ |
| reality | r | r | r | c | r | c | pr | r |
Table 1. Explicit expressions of the matrices corresponding to the fundamental rotation (i.e., a rotation of $`\theta _N`$ given in (1)) and the parity transformation. $`\stackrel{~}{R}_{2N}`$, $`\stackrel{~}{P}_Q`$ and $`\stackrel{~}{P}_D`$ are given in (32) and (19), respectively, where $`\mathrm{}𝐍`$ and $`(N1)/2`$, r=real, c=complex, pr=pseudo real. All the real irreps of $`Q_{2N}`$ are those of $`D_N`$. Complex one-dimensional irreps exist only for $`N=3,5,7,\mathrm{}`$, while the real one-dimensional irreps $`\mathrm{𝟏}_{,0}`$ and $`\mathrm{𝟏}_{,2}`$ exist only for $`N=2,4,6,\mathrm{}`$.
Given the details of the $`Q_{2N}`$ irreps, it is then straightforward to construct an invariant action hamed1 ; hill ; hamed2 . Supersymmetrization can also be straightforwardly done hamed2 .
## IV Orbifold boundary conditions and $`Q_{2N}`$ flavor symmetry
In the case of a continuos extra dimension, orbifold boundary conditions are used to suppress unnecessary light fields and also to obtain four-dimensional chiral fields. We shall discuss next how an internal $`Q_{2N}`$ flavor symmetry can appear even if orbifold boundary conditions break the dihedral invariance (41). Let $`\varphi (x,y)`$ be a generic field which satisfies the periodic boundary condition, $`\varphi (x,y)=\varphi (x,y+Na)`$. Then the field $`\varphi (x,y)`$ can be decomposed into the cosine and sine modes:
$`\varphi (x,y)`$ $`=`$ $`{\displaystyle \frac{\varphi (x)}{\sqrt{N}}}+{\displaystyle \underset{i=1}{\overset{i_{\mathrm{max}}}{}}}\varphi _{+,i}(x)cos(k_iy)+{\displaystyle \underset{i=1}{\overset{i_{\mathrm{max}}^{}}{}}}\varphi _{,i}(x)sin(k_iy),`$ (42)
where
$`\varphi (x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{N1}{}}}\varphi (x,y_n),`$ (43)
$`k_i`$ $`=`$ $`{\displaystyle \frac{2\pi i}{aN}},i𝐍,i_{\mathrm{max}}=\{\begin{array}{c}i_{\mathrm{max}}^{}+1=N/21\\ i_{\mathrm{max}}^{}=(N1)/2\end{array}\text{for}\{\begin{array}{c}\text{even}N\\ \text{odd}N\end{array}.`$ (48)
$`\varphi (x)`$ is the zero mode. As in the continuos case, we can drop the cosine or sine modes by imposing an appropriate boundary condition: Under the parity transformation (13), i.e.,
$`y_0`$ $`y_0^{}=y_0,y_1y_1^{}=y_{N1},\mathrm{},y_iy_i^{}=y_{Ni},\mathrm{},`$ (49)
the zero mode $`\varphi (x)`$ and the cosines modes are even, while the sine modes are odd.
Since the $`D_N`$ transformation mixes the cosine and sine modes, the orbifold boundary conditions break the dihedral invariance explicitly. However, the $`Q_{2N}`$ invariant construction of an action discussed in the previous section ensures that the $`Q_{2N}`$ invariance remains intact as a global, internal symmetry. This is because there is no derivative with respect to $`y`$ is used in the construction. So, the theory with orbifold boundary conditions is invariant under the internal transformation
$`\varphi (x,y)`$ $``$ $`\varphi ^{}(x,y)=\stackrel{~}{Q}_{2N}\varphi (x,y),\stackrel{~}{Q}_{2N}Q_{2N},`$ (50)
which should be compared with (41). The internal symmetry is nothing but a global flavor symmetry based on $`Q_{2N}`$.
## V An example
In what follows, we would like to discuss a concrete model. One of the successful Ansätze for the quark mass matrices is of a nearest neighbor interaction (NNI) type weinberg ; wilczek ; fritzsch1
$`M`$ $`=`$ $`\left(\begin{array}{ccc}0& C& 0\\ \pm C& 0& B\\ 0& B^{}& A\end{array}\right).`$ (54)
In babu4 it has been proposed to derive the mass matrix (54) solely from a dihedral symmetry, and concluded that two conditions should be met: (i) There should be real as well as pseudo real nonsinglet representations, and (ii) there should be the up and down type Higgs $`SU(2)_L`$ doublets (type II Higgs). The smallest finite group that allows both real and pseudo real nonsinglet representations is $`Q_6`$ as we have seen. Further, the Higgs sector of the MSSM fits the desired Higgs structure. Therefore, we assume supersymmetry in four dimensions. The $`D_3(S_3)`$ model of kubo1 with a $`Z_2`$ symmetry in the leptonic sector is one of the most predictive models for the leptonic sector. However, the $`Z_2`$ symmetry in the quark sector is broken, so that the $`Z_2`$ symmetry should be seen as an approximate symmetry in that model. It was found, however, that this leptonic sector can be reproduced in a supersymmetric $`Q_6`$ model without introducing an additional discrete symmetry into the leptonic sector babu4 . In Table 2 we write the $`Q_6`$ assignment of the quark and lepton chiral supermultiplets <sup>5</sup><sup>5</sup>5The same model exists for $`Q_{2N}`$ if $`N`$ is odd and a multiple of $`3`$.:
| | $`Q`$ | $`U^c,D^c,L,E^c,N^cH^u,H^d`$ | $`Q_3`$ | $`U_3^c,D_3^c,H_3^u,H_3^d`$ | $`L_3,E_3^c`$ | $`N_3^c`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`Q_6`$ | $`\mathrm{𝟐}_1`$ | $`\mathrm{𝟐}_2`$ | $`\mathrm{𝟏}_{+,2}`$ | $`\mathrm{𝟏}_{,1}`$ | $`\mathrm{𝟏}_{+,0}`$ | $`\mathrm{𝟏}_{,3}`$ |
Table 2. $`Q_6`$ assignment of the matter supermultiplets. $`Q,Q_3,L,L_3`$ and $`H^u,H_3^u,H^d,H_3^d`$ stand for $`SU(2)_L`$ doublets supermultiplets for quarks, leptons and Higgs bosons, respectively. Similarly, $`SU(2)_L`$ singlet supermultiplets for quarks, charged leptons and neutrinos are denoted by $`U^c,U_3^c,D^c,D_3^c,E^c,E_3^c`$ and $`N^c,N_3^c`$. This is an alternative assignment to the one given in the footnote of babu4 . The present assignment can more suppress the proton decay itou1 . The assignment for the mirror supermultiplets can be simply read off from Table 2.
We impose the following orbifold boundary conditions: All the mirror chiral supermultiplets are odd under the parity transformation (49). Similarly, the $`N=1`$ chiral supermultiplets, which are the $`N=2`$ superpartners of the $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge supermultiplets, are also odd. It is then clear that the zero modes of the gauge, matter and Higgs supermultiplets coincide with those of the supersymmetric $`Q_6`$ model of babu4 , and hence it is the low energy effective theory. The low energy Yukawa superpotential $`W_Y`$ is given by
$`W_Y`$ $`=`$ $`W_Q+W_L,`$ (55)
where <sup>6</sup><sup>6</sup>6 The Higgs sector of the model of babu4 possesses a permutation symmetry $`H_1^{u(d)}H_2^{u(d)}`$, which ensures the stability of the VEV $`<H_1^{u(d)}>=<H_2^{u(d)}>`$. The resulting mass quark matrices are equivalent to (54). The leptonic sector given in kubo1 can be obtained by the interchange $`12`$.
$`W_Q`$ $`=`$ $`Y_a^uQ_3U_3^cH_3^u+Y_b^uQ^T\sigma _1U_3^cH^uY_b^{}^uQ_3U^{cT}i\sigma _2H^u+Y_c^uQ^T\sigma _1U^cH_3^u`$ (56)
$`+Y_a^dQ_3D_3^cH_3^d+Y_b^dQ^T\sigma _1D_3^cH^dY_b^{}^dQ_3D^{cT}i\sigma _2H^d+Y_c^dQ^T\sigma _1D^cH_3^d,`$
$`W_L`$ $`=`$ $`Y_c^ef^{IJK}L_IE_J^cH_K^d+Y_b^{}^eL_3(H_1^dE_1^c+H_2^dE_2^c)+Y_b^e(L_1H_1^d+L_2H_2^d)E_3^c`$ (57)
$`+Y_a^\nu L_3N_3^cH_3^u+Y_c^\nu f^{IJK}L_IN_J^cH_K^u+Y_b^{}^\nu L_3(H_1^uN_1^c+H_2^uN_2^c),`$
and $`f^{122}=f^{212}=f^{222}=f^{111}=1`$. In babu4 it has been found that by introducing a certain set of gauge singlet Higgs supermultiplets it is possible to construct a Higgs sector in such a way that CP phases can be spontaneously induced. Therefore, all the parameters appearing in the Lagrangian including the soft supersymmetry breaking (SSB) sector are real. Consequently, no $`CP`$ violating processes induced by SSB terms are possible in this model, satisfying the most stringent experimental constraint coming from the EDM of the neutron and the electron edm . Since the Higgs sector is also $`Q_6`$ invariant, it is straightforward to derive it from dimensional deconstruction. Consequently, the quark sector contains only 8 real parameters with one independent phase to describe the quark masses and their mixing, and the leptonic sector contains only 6 real parameters with one independent phase to describe 12 independent physical parameters. Predictions in the $`|V_{ub}|\mathrm{sin}2\varphi _1`$ planes are shown in Fig. 2, while Fig. 3 shows the predictions in the $`\mathrm{sin}2\varphi _1\varphi _3`$ planes.
As we can see from Fig. 2 and 3, with accurate measurements of the Cabibbo-Kobayashi-Maskawa matrix elements, the predictions could be tested.
The predictions in the leptonic sector are summarized as follows <sup>7</sup><sup>7</sup>7Large mixing of neutrinos may be obtained in dimensional deconstruction models in a different mechanism. See, for instance, seidl1 ; balaji ; enkhbat .:
1. Inverted neutrino mass spectrum, i.e., $`m_{\nu _3}<m_{\nu _1},m_{\nu _2}`$.
2. $`m_{\nu _2}^2/\mathrm{\Delta }m_{23}^2=\frac{(1+2t_{12}^2+t_{12}^4rt_{12}^4)^2}{4t_{12}^2(1+t_{12}^2)(1+t_{12}^2rt_{12}^2)\mathrm{cos}^2\varphi _\nu }\mathrm{tan}^2\varphi _\nu (r=\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{23}^2,t_{12}=\mathrm{tan}\theta _{12})`$, where $`\varphi _\nu `$ is an independent phase.
3. $`\mathrm{sin}\theta _{13}m_e/\sqrt{2}m_\mu 3.4\times 10^3`$ and $`\mathrm{tan}\theta _{23}1(m_e/\sqrt{2}m_\mu )^2=1O(10^5)`$.
4. The prediction of $`<m_{ee}>`$ is shown in Fig. 4.
We emphasize that the smallness of $`\mathrm{sin}\theta _{13}`$ and the almost maximal mixing of the atmospheric neutrinos are consequences of the $`Q_6`$ flavor symmetry. $`\mathrm{sin}\theta _{13}`$ in the present model may be too small to be measured in a laboratory experiment minakata , but the tiny deviation from zero ($`\mathrm{sin}^2\theta _{13}m_e^2/2m_\mu ^210^5`$) are important in supernova neutrino oscillationsando .
## VI Conclusion
In this letter we have looked for a possible origin of dihedral symmetries. It has been recently realized that a flavor symmetry based on a dihedral group can be used to soften the flavor problem of the SM and the MSSM. We have considered an extra dimension compactified on a closed chain, which is assumed to have the form of a regular polygon. Since the symmetry group of the regular polygon is the dihedral group $`D_N`$, we assumed that the fields are irreps of the covering group of $`D_N`$, which is the binary dihedral group $`Q_{2N}`$. The construction of an action with the dihedral invariance is straightforward, and moreover we found that the $`Q_{2N}`$ symmetry remains as an intact, internal flavor symmetry even if the original dihedral invariance is broken by orbifold boundary conditions. We hope that with our finding we can come closer to a deep understanding of the origin of a flavor symmetry based on a nonabelian finite group.
|
warning/0506/quant-ph0506039.html
|
ar5iv
|
text
|
# Two-way quantum communication channels
## 1 Introduction
Suppose two parties, Alice and Bob, wish to exchange information. To do so, they must be connected by some physical interaction, or in information-theoretic language, a channel. One of the main problems of information theory is to determine the maximum rates (i.e., the capacities) for communication through such a channel.
A particularly well-studied type of interaction is the one-way channel, which models transporting some carrier of information from a fixed sender (say, Alice) to a fixed receiver (say, Bob), with the state encoding the information possibly modified in some way during transit. In other words, a one-way channel is a formal change of ownership of the state together with a state change.
Classical one-way channels, introduced by Shannon , are stochastic maps on probability distributions. More generally, in quantum mechanics, evolution is described by a quantum operation, i.e., a trace-preserving, completely positive map on quantum states represented by density matrices (positive semidefinite matrices of unit trace). Figure 1(a) shows a one-way quantum channel described by the quantum operation $``$. Alice prepares an input state $`\rho _\mathrm{A}`$ which is transported to Bob as the output state $`\rho _\mathrm{B}=(\rho _\mathrm{A})`$. The capacities to transmit classical and quantum messages through such channels have been extensively studied.
Although a one-way channel is an intuitive model for communication, it is only a special case of the possible interactions between Alice and Bob. The most general interaction is a joint quantum operation $`𝒩`$ acting on the joint Hilbert space of the two parties. If they are in possession of a joint density matrix $`\rho _{\mathrm{AB}}`$, then the action of the channel produces a joint output state $`\rho _{\mathrm{AB}}^{}=𝒩(\rho _{\mathrm{AB}})`$. In other words, Alice and Bob each provide an input to the two-way channel, which evolves their inputs jointly to produce a joint output shared by the two parties. Such a two-way quantum channel is shown in Fig. 1(b).
The two-way quantum channel model is the most general setting for two-party communication. For example, one-way channels are simply two-way channels with a zero-dimensional input for Bob and a zero-dimensional output for Alice (or equivalently from an operational standpoint, channels that discard Bob’s input and give Alice a fixed output).
Another subclass of two-way quantum channels consists of the classical two-way channels. Such channels were also first studied by Shannon , who gave outer and inner bounds on their capacity regions. Shannon’s bounds were subsequently improved by many others (see for example Refs. 13, 14, 15, 16, 17, 18, 19, 20).
In the quantum setting, another natural special class of two-way channels is the set of bipartite unitary interactions acting on systems of fixed dimension. The capacity question in this setting was formalized and studied in Ref. 21, and further capacity expressions were subsequently found in Refs. 22 and 23. For such channels, the problem of bidirectional communication is closely related to the problem of generating entanglement . Generalizing to allow systems whose input and output dimensions are different, one finds an especially simple class of interactions, the quantum feedback channels . Such channels take no input from Bob and evolve Alice’s input into a state shared between her and Bob.
A unifying viewpoint is that any two-way quantum channel can be viewed as an isometry of the two input states to three output states, discarding the third part of the system to an inaccessible environment. A variety of subclasses can be obtained simply by changing the dimensions of the terminals.
The various subclasses of two-way channels can exhibit remarkably different properties. For example, any (nonlocal) two-way unitary channel can communicate in either direction and generate entanglement at a nonzero rate , while in each of the other examples, some of these tasks are known to be impossible. Because of the wide variety of possible subclasses of two-way channels, and because calculating capacities is known to be difficult for several of the possibilities, we do not expect especially tight capacity results for the general two-way channel.
A communication protocol using two-way channels may yield classical or quantum communication in either or both directions. It can also create or consume classical or quantum correlations as auxiliary resources. In particular, providing Alice and Bob with enough free entanglement considerably simplifies and unifies the study of communication capacities of one-way quantum channels as well as two-way channels .
In this paper, we will mostly be considering the set of achievable rates of classical communication $`R_{}`$ from Alice to Bob and $`R_{}`$ from Bob to Alice, and the rate of producing entanglement $`R_\mathrm{e}`$ (which can either be positive, indicating that entanglement is produced, or negative, indicating that entanglement is consumed). The set of achievable rates $`(R_{},R_{},R_\mathrm{e})`$ forms a three-dimensional region. The boundary of this region represents the set of best achievable rates. Figure 2 shows a schematic diagram of the achievable rates of classical communication $`(R_{},R_{})`$ at fixed $`R_\mathrm{e}`$. Typically, there is a tradeoff between how much information Alice can send to Bob and how much information Bob can send to Alice. The end points of the optimal tradeoff curve, where $`R_{}=0`$ or $`R_{}=0`$, are called the one-way capacities of the two-way channel. The cases in which $`R_\mathrm{e}=0`$ and $`R_\mathrm{e}\mathrm{}`$, called entanglement-unassisted (or simply unassisted) and entanglement-assisted, respectively, are of particular interest.
For the case of entanglement-assisted one-way classical communication by unitary two-way channels, Ref. 21 gave a simplified single-letter expression for the capacity, showed that it is additive, and gave a protocol for achieving it. These results were subsequently extended in Ref. 22 to the case in which a fixed amount of entanglement is consumed or generated (a scenario dubbed finite entanglement assistance). In addition, tradeoff curves $`\{(R_{},0,R_\mathrm{e})\}`$ for one-way classical communication were related to analogous curves $`\{(R_{},0,R_\mathrm{e})\}`$ for quantum communication. In Ref. 23, the entire three-dimensional achievable regions for finite entanglement-assisted bidirectional classical or quantum communication were related.
In this paper, we consider bidirectional classical communication using a two-way quantum channel in the general (not necessarily unitary) case, possibly allowing entanglement assistance. We begin in Sec. 2 by introducing some basic notation and formalizing the notions of a protocol and the capacities it achieves. In Sec. 3, we extend techniques from Ref. 21 to provide inner and outer bounds on the achievable region. The outer bound generalizes Shannon’s bounds to the quantum setting. The inner bound is an extension of the protocol of Ref. 21 for simultaneous two-way communication, and also reduces to Shannon’s bound in the classical case. Furthermore, the bounds meet on the axes (as depicted in Fig. 2), giving a formula for the one-way capacity of a two-way quantum channel (extending the unitary result of Ref. 21). In Sec. 4, we describe two immediate applications of the inner and outer bounds. First, as a simple consequence of the fact that the inner and outer bounds meet on the axes, we recover the additivity of the entanglement-assisted classical capacity of one-way channels . We also provide a simple operational derivation of the additivity of the unassisted classical capacity of an entanglement-breaking channel, a special case of a result that was first proved in Ref. 27. Finally, in Sec. 5, we discuss some open questions and directions for future investigation.
## 2 Preliminaries
In this section, we describe our framework and define the notation used throughout the paper.
An ebit refers to a unit of shared quantum correlation, as quantified by an EPR pair of qubits $`|\mathrm{\Phi }_{\mathrm{AB}}:=\frac{1}{\sqrt{2}}_{x=0}^1|x_\mathrm{A}|x_\mathrm{B}`$, with the density matrix $`\mathrm{\Phi }:=|\mathrm{\Phi }\mathrm{\Phi }|`$. Throughout the paper, we omit the tensor product symbol, $``$, if no confusion will arise. The functions $`\mathrm{exp}`$ and $`\mathrm{log}`$ are always base $`2`$. The positive and negative parts of a real number $`x`$ are written as $`x^\pm `$: for $`x0`$, $`x^+:=x`$ and $`x^{}:=0`$; for $`x<0`$, $`x^+:=0`$ and $`x^{}:=x`$. We adopt the convention that a quantum operation $`𝒩`$ acts on an ensemble of quantum states $`=\{p_i,\rho _i\}`$ by acting on each state individually (preserving its probability), i.e., $`𝒩=\{p_i,𝒩(\rho _i)\}`$.
Our framework for communication using two-way channels extends the model of Refs. 21 and 23. A general two-way quantum channel $`𝒩`$ has two inputs of dimensions $`d_\mathrm{A}^{\mathrm{in}},d_\mathrm{B}^{\mathrm{in}}`$ and two outputs of dimensions $`d_\mathrm{A}^{\mathrm{out}},d_\mathrm{B}^{\mathrm{out}}`$. A general protocol $`𝒫_n`$ for two-way classical communication with $`n`$ uses of $`𝒩`$ consists of $`n`$ alternating steps of local processing and application of $`𝒩`$ followed by a final step of local processing, as depicted in Fig. 3. Let $`a,b`$ be the respective messages to be communicated from Alice to Bob and vice versa, and let $`𝒜_k,_k`$ be the local operations of Alice and Bob between the $`k`$th and $`(k+1)`$th use of $`𝒩`$. The operations $`𝒜_k,_k`$ may be arbitrary, but without loss of generality, we can assume that all ancillas are present at the beginning of the protocol and that all measurements are performed coherently, with no systems discarded by $`𝒜_k,_k`$. In other words, all of $`𝒜_k,_k`$ can be assumed to be unitary. Alice’s initial operation $`𝒜_0`$ has three input systems: $`\mathrm{A}_\mathrm{m}`$, which contains the message $`a`$; $`\mathrm{A}_\mathrm{a}`$, which contains a local ancilla; and $`\mathrm{A}_\mathrm{e}`$, which is maximally entangled with $`\mathrm{B}_\mathrm{e}`$. The operation $`𝒜_0`$ converts $`\mathrm{A}_{\mathrm{m},\mathrm{a},\mathrm{e}}`$ unitarily into two systems: $`\mathrm{A}`$, which is the input to $`𝒩`$, and $`\mathrm{A}^{}`$, which is not acted on by $`𝒩`$. Each of $`𝒜_k`$ for $`k=1,\mathrm{},n1`$ has two inputs $`\mathrm{A},\mathrm{A}^{}`$ and two outputs, which for simplicity will also be labeled as $`\mathrm{A},\mathrm{A}^{}`$, although in general they may have different dimensions than the input systems. The operation $`𝒜_k`$ changes the dimensions of $`\mathrm{A},\mathrm{A}^{}`$ if $`d_\mathrm{A}^{\mathrm{in}}d_\mathrm{A}^{\mathrm{out}}`$. Finally, $`𝒜_t`$ converts $`\mathrm{A},\mathrm{A}^{}`$ to three output systems: $`\mathrm{A}_\mathrm{m}`$, whose reduced state represents Bob’s message for Alice; $`\mathrm{A}_\mathrm{a}`$, which is an arbitrary ancillary system; and $`\mathrm{A}_\mathrm{e}`$, which is nearly maximally entangled with $`\mathrm{B}_\mathrm{e}`$. The situation is analogous on Bob’s side.
The goal is to find a family of protocols $`\{𝒫_n:n\}`$, where protocol $`𝒫_n`$ employs $`n`$ uses of $`𝒩`$, that will perform the communication task with accuracy that can be made arbitrarily good by taking $`n`$ sufficiently large—the output state of $`\mathrm{A}_\mathrm{m},\mathrm{B}_\mathrm{m}`$ should be close to $`|bb||aa|`$, the state of $`\mathrm{A}_\mathrm{e},\mathrm{B}_\mathrm{e}`$ should be nearly maximally entangled, and $`\mathrm{A}_\mathrm{a},\mathrm{B}_\mathrm{a}`$ should be nearly disentangled from $`\mathrm{A}_{\mathrm{m},\mathrm{e}},\mathrm{B}_{\mathrm{m},\mathrm{e}}`$ (but can depend on the messages $`a,b`$). The input dimensions of $`\mathrm{A}_\mathrm{m},\mathrm{B}_\mathrm{m}`$ determine the communication rates, while the change of the dimensions of $`\mathrm{A}_\mathrm{e},\mathrm{B}_\mathrm{e}`$ determines the amount of entanglement consumed or generated by $`𝒫_n`$.
To quantify the proximity of two quantum states $`\rho ,\sigma `$, we use the trace distance $`\frac{1}{2}\rho \sigma _1`$, where $`X_1:=tr\sqrt{X^{}X}`$. For example, for two pure states $`|\alpha ,|\beta `$, $`\frac{1}{2}|\alpha \alpha ||\beta \beta |_1=ϵ|\alpha |\beta |^2=1ϵ^2`$.
We are now ready to define the achievable region:
###### Definition 2.1 (Achievable rates).
We say that $`(R_{},R_{},R_\mathrm{e})`$ is achievable if there is a sequence of protocols $`\{𝒫_n\}`$, together with asymptotically vanishing sequences $`\{\delta _n\},\{ϵ_n\}`$ of nonnegative numbers and sequences of integers $`\{r_{,n}\},\{r_{,n}\},\{r_{\mathrm{e},n}\}`$ such that
$`r_{,n}`$ $`n(R_{}\delta _n),`$ (1)
$`r_{,n}`$ $`n(R_{}\delta _n),`$ (2)
$`r_{\mathrm{e},n}^+`$ $`n(R_\mathrm{e}^+\delta _n),`$ (3)
$`r_{\mathrm{e},n}^{}`$ $`n(R_\mathrm{e}^{}+\delta _n),`$ (4)
and for all messages $`a\{0,1\}^{r_{,n}}`$, $`b\{0,1\}^{r_{,n}}`$ the following success criteria hold. Let
$$\rho _n^{ab}:=𝒫_n(|00|_{\mathrm{A}_\mathrm{a}}|00|_{\mathrm{B}_\mathrm{a}}|aa|_{\mathrm{A}_\mathrm{m}}|bb|_{\mathrm{B}_\mathrm{m}}\mathrm{\Phi }_{\mathrm{A}_\mathrm{e}\mathrm{B}_\mathrm{e}}^{r_{\mathrm{e},n}^{}}),$$
(5)
and let $`\rho _{n,\mathrm{A}_\mathrm{m}\mathrm{B}_\mathrm{m}}^{ab}`$, $`\rho _{n,\mathrm{A}_\mathrm{e}\mathrm{B}_\mathrm{e}}^{ab}`$ be its reductions to the message and entanglement subspaces, respectively. Then the success criteria are
$`\frac{1}{2}\rho _{n,\mathrm{A}_\mathrm{m}\mathrm{B}_\mathrm{m}}^{ab}|bb|_{\mathrm{A}_\mathrm{m}}|aa|_{\mathrm{B}_\mathrm{m}}_1`$ $`ϵ_n,`$ (6)
$`\frac{1}{2}\rho _{n,\mathrm{A}_\mathrm{e}\mathrm{B}_\mathrm{e}}^{ab}\mathrm{\Phi }_{\mathrm{A}_\mathrm{e}\mathrm{B}_\mathrm{e}}^{r_{\mathrm{e},n}^+}_1`$ $`ϵ_n.`$ (7)
The achievable region is convex by convex combination (i.e., time sharing) of protocols. It is monotone since resources can always be discarded to give a lower rate of communication, lower rate of entanglement production, or higher rate of entanglement consumption.
Any quantum protocol for communicating classical information concludes by decoding the quantum state to produce a classical output. The final decoding can be viewed as a measurement of the final state, which in our model is implemented unitarily. The general problem of extracting classical information encoded in a quantum state (i.e., learning the index $`i`$ in a random draw from an ensemble $`=\{p_i,\rho _i\}`$ of quantum states) has been thoroughly studied in this context . The accessible information of $``$, denoted $`I_{\mathrm{acc}}()`$, is defined to be the maximum mutual information between the label $`i`$ and the outcome of any possible measurement. It is upper bounded as
$$I_{\mathrm{acc}}()\chi (),$$
(8)
where $`\chi ()`$ denotes the Holevo information of $``$,
$$\chi ():=S\left(\underset{i}{}p_i\rho _i\right)\underset{i}{}p_iS(\rho _i).$$
(9)
Here $`S(\rho ):=tr\rho \mathrm{log}\rho `$ is the von Neumann entropy of $`\rho `$.
In fact, by appropriate encoding, the Holevo information turns out to be asymptotically achievable in the following sense. Consider a CQ channel, i.e., a one-way channel with a classical input $`i`$ giving rise to a corresponding quantum output state $`\rho _i`$. Then we have
###### Theorem 2.2 (HSW Theorem ).
The capacity of the CQ channel $`\{i\rho _i\}`$ is given by $`sup_{\{p_i\}}\chi (\{p_i,\rho _i\})`$.
More specifically, for any fixed input probability distribution $`\{p_i\}`$, let $`=\{p_i,\rho _i\}`$ denote the corresponding ensemble. In the limit of large $`n`$, there is a set of $`\mathrm{exp}(n\chi ()\delta _n)`$ states $`\rho _{i_1}\rho _{i_2}\mathrm{}\rho _{i_n}`$ that can be distinguished with error bounded by some $`ϵ_n`$, where $`\delta _n,ϵ_n0`$ as $`n\mathrm{}`$. This is proved by showing that, with nonzero probability, a random code (in which each tensor component of each codeword is drawn independently from $``$) can be decoded with vanishing error.
The HSW theorem can be applied to an arbitrary one-way channel by viewing the states $`\rho _i`$ as the possible outputs of the channel. However, for a general two-way channel, the coding problem is complicated somewhat by the fact that the input and output ensembles are bipartite. In this case, communication in one direction is typically affected by the input to the channel from the opposite direction. Thus, we will need to prove a modified version of the HSW theorem when we derive protocols for bidirectional communication.
## 3 Entanglement-assisted capacity of two-way quantum channels
In this section, we present inner and outer bounds on the classical capacity region of an entanglement-assisted two-way quantum channel.
Let $`𝒫_n`$ be an arbitrary protocol that employs $`n`$ uses of the two-way channel $`𝒩`$. We assume without loss of generality that Alice and Bob retain a copy of their input messages throughout the protocol. Then, after $`t`$ uses of $`𝒩`$ followed by $`𝒜_t_t`$, Alice and Bob possess a joint state drawn from an ensemble
$$^{(t)}=\{p_{ab},|aa|\rho _{ab}^{(t)}|bb|\}$$
(10)
indexed by the classical messages $`a,b`$ to be communicated, where $`^{(t)}`$ is completely determined by $`𝒫_n`$.
For a bipartite ensemble $``$ over systems $`\mathrm{AA}^{}`$ and $`\mathrm{BB}^{}`$, we define the local Holevo information of $`\mathrm{BB}^{}`$ and $`\mathrm{AA}^{}`$, respectively, as
$`\chi _{}()`$ $`:=\chi (tr_{\mathrm{AA}^{}}),`$ (11)
$`\chi _{}()`$ $`:=\chi (tr_{\mathrm{BB}^{}}).`$ (12)
Following earlier work , our analysis will be based on examining the local Holevo information $`\chi _{}(^{(t)})`$, $`\chi _{}(^{(t)})`$ for $`t=0,1,\mathrm{},n`$ during the course of an $`n`$-use protocol. By the joint entropy theorem (see for example Eq. (1.58) in Ref. 28), for an ensemble $`^{(t)}`$ of the form of Eq. (10) (i.e., with local copies of the classical messages), the local Holevo information can be rewritten as
$`\chi _{}(^{(t)})`$ $`=H(\{p_b\})+{\displaystyle \underset{b}{}}p_b\chi (_b^{(t)}),`$ (13)
$`\chi _{}(^{(t)})`$ $`=H(\{p_a\})+{\displaystyle \underset{a}{}}p_a\chi (_a^{(t)}),`$ (14)
where $`p_a:=_bp_{ab}`$ and $`p_b:=_ap_{ab}`$ are the marginal distributions and $`_b^{(t)}:=\{p_{a|b},tr_{\mathrm{AA}^{}}\rho _{ab}^{(t)}\},_a^{(t)}:=\{p_{b|a},tr_{\mathrm{BB}^{}}\rho _{ab}^{(t)}\}`$ are the ensembles for Bob and Alice conditioned on their known inputs $`b`$ and $`a`$, respectively, where $`p_|`$ denotes conditional probability.
Equations (13) and (14) give natural interpretations of the local Holevo information: for example, for Bob, it is the sum of the information about $`b`$ that he already knows and the information about $`a`$ obtainable from $`_b^{(t)}`$, averaged over $`b`$. Removing the information that the sender already knows gives a quantity that is useful for obtaining bounds. Thus, we define the readjusted local Holevo information as
$`\overline{\chi }_{}(^{(t)})`$ $`:={\displaystyle \underset{b}{}}p_b\chi (_b^{(t)}),`$ (15)
$`\overline{\chi }_{}(^{(t)})`$ $`:={\displaystyle \underset{a}{}}p_a\chi (_a^{(t)}).`$ (16)
### 3.1 Additive outer bound
In this section, we obtain an additive outer bound on the capacity region for entanglement-assisted bidirectional classical communication using a two-way quantum channel.
Consider the differences in local Holevo information induced by an application of $`𝒩`$ to an ensemble $``$:
$`\mathrm{\Delta }\chi _{}()`$ $`:=\chi _{}(𝒩)\chi _{}()=\overline{\chi }_{}(𝒩)\overline{\chi }_{}(),`$ (17)
$`\mathrm{\Delta }\chi _{}()`$ $`:=\chi _{}(𝒩)\chi _{}()=\overline{\chi }_{}(𝒩)\overline{\chi }_{}().`$ (18)
Let $`[[x,y]]`$ denote the region $`[0,x]\times [0,y]^2`$, and let $`conv()`$ denote the convex hull. In terms of these quantities, we have the following outer bound on the achievable region:
###### Theorem 3.1.
If $`(R_{},R_{},\mathrm{})`$ is achievable, then
$$(R_{},R_{})conv\{[[\mathrm{\Delta }\chi _{}(),\mathrm{\Delta }\chi _{}()]]:\mathrm{arbitrary}\}.$$
(19)
###### Proof 3.2.
For any $`𝒫_n`$, the Holevo bound on the accessible information, Eq. (8), implies that the number of bits that can be faithfully transmitted forward and backward are no more than $`\overline{\chi }_{}(^{(n)})`$ and $`\overline{\chi }_{}(^{(n)})`$ respectively. In other words,
$`R_{}`$ $`\frac{1}{n}\overline{\chi }_{}(^{(n)}),`$ (20)
$`R_{}`$ $`\frac{1}{n}\overline{\chi }_{}(^{(n)}).`$ (21)
Expressing $`\overline{\chi }_{}(^{(n)})`$ as a telescopic sum and using the fact $`\overline{\chi }_{}(^{(0)})=0`$,
$`\overline{\chi }_{}(^{(n)})`$ $`={\displaystyle \underset{t=1}{\overset{n}{}}}[\overline{\chi }_{}(^{(t)})\overline{\chi }_{}(^{(t1)})].`$ (22)
$`={\displaystyle \underset{t=1}{\overset{n}{}}}\mathrm{\Delta }\chi _{}(^{(t1)}),`$ (23)
where Eq. (23) comes from the fact that $`^{(t)}=(𝒜_t_t)𝒩^{(t1)}`$, and $`𝒜_t`$, $`_t`$ leave $`\overline{\chi }_{}`$ invariant so that $`\overline{\chi }_{}(^{(t)})=\overline{\chi }_{}(𝒩^{(t1)})`$. Likewise,
$$\overline{\chi }_{}(^{(n)})=\underset{t=1}{\overset{n}{}}\mathrm{\Delta }\chi _{}(^{(t1)}).$$
(24)
Putting Eqs. (20), (21), (23), and (24) together, we find
$`(R_{},R_{})`$ $`[[\frac{1}{\mathrm{n}}\chi _{}(^{(n)}),\frac{1}{n}\chi _{}(^{(n)})]]`$ (25)
$`conv\{[[\mathrm{\Delta }\chi _{}(),\mathrm{\Delta }\chi _{}()]]:\mathrm{arbitrary}\},`$ (26)
which completes the proof.
As a side remark, we can also keep track of the entanglement $`E`$ of the ensemble $`^{(t)}`$ along with the local Holevo information. We define $`\mathrm{\Delta }E():=E(𝒩)E()`$, where $`E(\{p_i,\rho _i\}):=_ip_iE(\rho _i)`$ and $`E`$ is defined to be the distillable entanglement if $`R_\mathrm{e}0`$ and the entanglement cost if $`R_\mathrm{e}0`$. Thus, we have the following:
###### Theorem 3.3.
If $`(R_{},R_{},R_\mathrm{e})`$ is achievable, then
$$(R_{},R_{},R_\mathrm{e})conv\{[[\mathrm{\Delta }\chi _{}(),\mathrm{\Delta }\chi _{}()]]\times [\mathrm{\Delta }E()^{},\mathrm{\Delta }E()^+]:\text{arbitrary }\}.$$
(27)
However, in the remainder of Sec. 3, we will consider only the case of unlimited entanglement assistance.
### 3.2 Two-way communication protocols based on remote state preparation
Having derived outer bounds on the achievable region, we now turn to inner bounds based on protocols that achieve particular rates of communication. The main idea is to generalize the protocol from Sec. 4.3 of Ref. 21 (and see also generalizations in Refs. 22 and 29), originally designed for one-way communication with two-way unitary channels, to the general problem of two-way communication with possibly nonunitary two-way channels. The main ingredients of the protocol of Ref. 21 are the HSW Theorem and remote state preparation (RSP), which allows Alice to share states of her choice with Bob. Our two-way protocol is based on two-way analogs of the HSW Theorem and RSP. In this section, we describe these tools, present the generalized two-way protocol, and analyze its error rate and inefficiency. In the following section, we show that combining the general technique with a particular method of remote state preparation gives an explicit inner bound.
First, we describe a bidirectional version of the HSW Theorem. Such a tool is necessary since the effective channel through which Alice can send signals to Bob depends on what input Bob is using to send signals to Alice, and vice versa. Consider any two-way CQ channel, i.e., a channel with two classical inputs, $`i`$ for Alice and $`j`$ for Bob, giving rise to a joint quantum output state $`\rho _{ij}`$. In this setting, we have the following:
###### Lemma 3.4 (Bidirectional HSW inner bound).
For the two-way CQ channel $`\{i,j\rho _{i,j}\}`$, rates $`(R_{},R_{})`$ satisfying
$$(R_{},R_{})conv\{[[\overline{\chi }_{}(),\overline{\chi }_{}()]]:=\{p_iq_j,\rho _{ij}\}\}$$
(28)
are achievable.
###### Proof 3.5 ( (sketch)).
We omit the straightforward (but lengthy) generalization of the proof of Theorem 2.2 in Refs. 4 and 5. The basic idea is as follows. For each use of the two-way channel, Alice is unaware of Bob’s input (which defines the effective channel from Alice to Bob), but such information is available for Bob in his decoding operation. The same holds for communication from Bob to Alice. Then, it is possible to show that good random codes (chosen independently by Alice and Bob) exist, and allow communication at the above rates according to a packing lemma analogous to that in the original proof.
Note that in this lemma, we have assumed that Alice and Bob choose their signals independently, with encoding distributions $`\{p_i\}`$ for Alice and $`\{q_j\}`$ for Bob, so that we can view Bob’s encoding distribution as inducing a distribution over CQ channels from Alice to Bob, and vice versa. However, this is not the most general encoding distribution possible with many uses of the channel, so it is not clear whether this inner bound for CQ channels can be exceeded.
To prepare ensembles for bidirectional communication, we consider bipartite remote state preparation. Here the goal is to prepare a large number $`n`$ of states drawn from the bipartite ensemble $`=\{p_{ij},\rho _{ij}\}`$ with each party knowing one of the labels $`i,j`$. Note that Alice’s label may control Bob’s portion of the state as well as her own (indeed, the state may be entangled between their respective systems), and similarly for Bob’s label.
We will assume the existence of a (not necessarily optimal) protocol for bipartite remote state preparation with known asymptotic classical communication and entanglement costs $`C_{},C_{}`$ and $`C_\mathrm{e}`$. More specifically, suppose $`n(C_{}+\delta _n^{\text{rsp}})`$ forward classical bits, $`n(C_{}+\delta _n^{\text{rsp}})`$ backward classical bits, and $`n(C_\mathrm{e}+\delta _n^{\text{rsp}})`$ ebits are sufficient for Alice and Bob to prepare a state drawn from $`^n`$ with fidelity $`1ϵ_n^{\text{rsp}}`$, such that $`\delta _n^{\text{rsp}},ϵ_n^{\text{rsp}}0`$ as $`n\mathrm{}`$. The problem of optimizing such costs in general is quite difficult, and has only been solved for very special cases of $``$. But assuming the existence of such a protocol to prepare any particular ensemble of the form $`=\{p_iq_j,\rho _{ij}\}`$ (where the labels are chosen independently but the corresponding state may be arbitrary) at given costs, a corresponding point can be attained in the achievable region.
###### Lemma 3.6.
$`(R_{},R_{},\mathrm{})`$ is achievable for all
$`(R_{},R_{})conv\{[[\mathrm{\Gamma }_{}(),\mathrm{\Gamma }_{}()]]:`$ $`=\{p_aq_b,\rho _{ab}\}\text{ such that}`$
$`C_{}()\overline{\chi }_{}(𝒩)\mathrm{\Gamma }_{}(),`$
$`C_{}()\overline{\chi }_{}(𝒩)\mathrm{\Gamma }_{}()\}.`$ (29)
Here $`\overline{\chi }_{}(𝒩),\overline{\chi }_{}(𝒩)`$ represent achievable forward and backward communication rates for the ensemble $`𝒩`$ (according to Lemma 3.4). The quantities $`\mathrm{\Gamma }_{}(),\mathrm{\Gamma }_{}()`$ thus represent the amount of communication gained by one use of $`𝒩`$ on the ensemble $``$; that is, the communication rates of $`𝒩`$ minus the communication costs of preparing $``$.
###### Proof 3.7.
By convexity and monotonicity, we only need to show that $`R_{}=\overline{\chi }_{}(𝒩)C_{}()`$ and $`R_{}=\overline{\chi }_{}(𝒩)C_{}()`$ are achievable rates (given a sufficiently large amount of entanglement assistance). We do this by giving a communication protocol achieving those rates assuming the existence of an RSP protocol with the stated communication costs. Since the protocol is a generalization of that in Sec. 4.3 of Ref. 21 for one-way communication, some readers may wish to refer to the detailed description and schematic diagram therein.
The protocol is as follows. Alice and Bob preagree on sufficiently large values of $`n`$ and $`k`$ (to be determined later) and proceed with the following protocol, using $`𝒩`$ approximately $`nk`$ times to communicate $`k`$ messages $`a_1,a_2,\mathrm{},a_k`$, each consisting of $`n(R_{}\delta _n^{\text{hsw}}\delta _n^{\text{rsp}})`$ bits, in the forward direction, and $`k`$ messages $`b_1,b_2,\mathrm{},b_k`$, each consisting of $`n(R_{}\delta _n^{\text{hsw}}\delta _n^{\text{rsp}})`$ bits, in the backward direction.
1. Using bipartite RSP, prepare a state $`\rho _1`$ from $`^n`$ with fidelity at least $`1ϵ_n^{\text{rsp}}`$. This requires $`O(n)`$ uses of $`𝒩`$ to communicate $`n(C_{}()+\delta _n^{\text{rsp}})`$ bits from Alice to Bob and $`n(C_{}()+\delta _n^{\text{rsp}})`$ bits from Bob to Alice. (Note that this initial communication is always possible if the appropriate $`\mathrm{\Delta }\chi `$ is positive for some ensemble (and otherwise, it is not necessary). Suppose that although $`\mathrm{\Delta }\chi _{}>0`$ for some ensemble, $`\mathrm{\Delta }\chi _{}0`$ for all ensembles that can be created at zero cost. In this case the operation is semicausal from Alice to Bob, and hence is also semilocalizable , meaning that it can be simulated by a local operation by Bob, sending a quantum state to Alice, and a final local operation by Alice. But such an operation clearly has $`\mathrm{\Delta }\chi _{}0`$ for all ensembles, which is a contradiction. A similar argument applies to $`\mathrm{\Delta }\chi _{}`$.)
2. Apply $`𝒩^n`$ to $`\rho _1`$, which has been chosen such that local measurements by Alice and Bob on $`𝒩^n(\rho _1)`$ provide an $`n(\overline{\chi }_{}(𝒩)\delta _n^{\text{hsw}})`$-bit message for Bob and an $`n(\overline{\chi }_{}(𝒩)\delta _n^{\text{hsw}})`$-bit message for Alice with probability at least $`1ϵ_n^{\text{rsp}}ϵ_n^{\text{hsw}}`$, according to Lemma 3.4. Bob receives the message $`a_1`$ as well as the information needed to perform RSP of $`\rho _2^n`$ in the next step. Similarly, Alice receives $`b_1`$ and together with the information she needs for RSP in the next round.
3. Perform the $`2\text{nd},\mathrm{},k\text{th}`$ rounds of RSP.
Just as in Ref. 21, Alice and Bob must know all of $`a_1,a_2,\mathrm{},a_k`$ and $`b_1,b_2,\mathrm{},b_k`$ at the beginning of the protocol, perform their measurements for the $`k`$th round of RSP to obtain the RSP instructions for the $`k`$th round, encode them as part of the message of the $`(k1)`$th round of RSP, and proceed with their part of the $`(k1)`$th round of RSP, and so on, until the first round RSP messages are generated and sent by the initial $`O(n)`$ uses of $`𝒩`$. For simplicity, we only consider non-interactive RSP protocols (with only one round of communication from Alice to Bob and vice versa), which will be sufficient for our applications. (With one more level of block coding, one should be able to use interactive RSP protocols, but this would require a more detailed error analysis.)
Finally, we analyze the errors and inefficiencies to show achievability of the rates. Fix any desired $`\delta ,ϵ>0`$. The above protocol employs $`n(c+k)`$ uses of $`𝒩`$ (for some constant $`c`$) and communicates $`nk(\overline{\chi }_{}()C_{}()\delta _n^{\text{rsp}}\delta _n^{\text{hsw}})`$ and $`nk(\overline{\chi }_{}()C_{}()\delta _n^{\text{rsp}}\delta _n^{\text{hsw}})`$ bits forward and backward, respectively, with error $`k(ϵ_n^{\text{hsw}}+ϵ_n^{\text{rsp}})`$. We can choose $`k,n`$ independently large enough so that $`\frac{c}{k}+\delta _n^{\text{rsp}}+\delta _n^{\text{hsw}}<\delta `$ and then increase $`n`$ if needed, to ensure $`k(ϵ_n^{\text{hsw}}+ϵ_n^{\text{rsp}})ϵ`$.
### 3.3 Inner bound
Explicit inner bounds for the achievable region can be obtained from Lemma 3.6 together with known RSP protocols. In particular, we will make extensive use of the protocol for one-way RSP of entangled states:
###### Theorem 3.8 (RSP of entangled states ).
Asymptotically, the ensemble $`=\{p_a,|\psi _a_{\mathrm{AB}}\}`$ can be prepared with a communication cost from Alice to Bob of $`C_{}()=\chi _{}()`$ and a rate of entanglement consumption of $`C_\mathrm{e}()=_ap_aS(tr_\mathrm{A}|\psi _a\psi _a|)`$.
Using this RSP protocol in the general inner bound of Lemma 3.6, we find
###### Theorem 3.9.
$`(R_{},R_{},\mathrm{})`$ is achievable if
$$(R_{},R_{})conv\{[[\mathrm{\Delta }\chi _{}(),\mathrm{\Delta }\chi _{}()]]:=\{p_a,\rho _a\}\{q_b,\eta _b\}\}$$
(30)
In the ensemble of Eq. (30), the tensor product decomposition of $`\mathrm{AA}^{}\mathrm{BB}^{}`$ may be arbitrary: Alice can prepare joint states of any subspace of her system and Bob’s, as can Bob, so long as the two subspaces are disjoint. In other words, we can have $`\rho _a\mathrm{X},\eta _b\mathrm{Y}`$ for any fixed decomposition $`\mathrm{AA}^{}\mathrm{BB}^{}=\mathrm{X}\mathrm{Y}`$ into arbitrary complementary subspaces.
###### Proof 3.10.
Alice and Bob are each given the knowledge of $`a`$ and $`b`$. By Theorem 3.8, $`C_{}()=\chi (\{p_a,\rho _a\})`$ and $`C_{}()=\chi (\{q_b,\eta _b\})`$. Thus, the result follows from Lemma 3.6.
When one of the ensembles is trivial, the protocol performs one-way communication, with $`\overline{\chi }_{}()=0`$ or $`\overline{\chi }_{}()=0`$ as appropriate. For example, in the former case, the rate of forward communication $`R_{}=\mathrm{\Delta }\chi _{}()`$ is achievable for arbitrary $``$, so that the inner and outer bounds meet. Therefore, we find an expression for the entanglement-assisted one-way forward capacity of a two-way quantum channel. Similarly, we find $`R_{}=\mathrm{\Delta }\chi _{}()`$ for the one-way backward capacity. Indeed, this result is immediate from the fact that the protocol of Ref. 21 for one-way communication applies unchanged even when $`𝒩`$ is not unitary. Thus we have the following:
###### Corollary 3.11 (One-way capacity of a two-way channel).
$`R_{}^{\mathrm{max}}`$ $`:=sup\{R_{}:(R_{},0,\mathrm{})\text{ achievable}\}=\underset{}{sup}\mathrm{\Delta }\chi _{}()`$ (31)
$`R_{}^{\mathrm{max}}`$ $`:=sup\{R_{}:(0,R_{},\mathrm{})\text{ achievable}\}=\underset{}{sup}\mathrm{\Delta }\chi _{}()`$ (32)
where the supremum is over all ensembles $`\{p_i,\rho _{i,\mathrm{AA}^{}\mathrm{BB}^{}}\}`$ with ancillary systems $`\mathrm{A}^{},\mathrm{B}^{}`$.
In particular, we have
###### Corollary 3.12.
$`R_,^{\mathrm{max}}`$ are strongly additive. In other words, for any pair of two-way quantum channels $`𝒩,𝒩^{}`$, we have $`R_{}^{\mathrm{max}}(𝒩𝒩^{})=R_{}^{\mathrm{max}}(𝒩)+R_{}^{\mathrm{max}}(𝒩^{})`$ and $`R_{}^{\mathrm{max}}(𝒩𝒩^{})=R_{}^{\mathrm{max}}(𝒩)+R_{}^{\mathrm{max}}(𝒩^{})`$.
### 3.4 Relation to Shannon’s classical bounds
Both the inner and outer bounds given above reduce to Shannon’s bounds in the case of a two-way classical channel (in which case entanglement assistance clearly does not help).
Consider sending information through a two-way classical channel. Suppose the input symbols $`a,b`$ appear with the joint probability distribution $`p_{ab}`$. Then the output symbols $`a^{},b^{}`$ appear with the joint probability distribution $`p_{ab}p_{a^{}b^{}|ab}`$, where the conditional probabilities $`p_{a^{}b^{}|ab}`$ define the channel. Let $`I(\mathrm{X};\mathrm{Y}|\mathrm{Z})`$ denote the conditional mutual information,
$$I(\mathrm{X};\mathrm{Y}|\mathrm{Z}):=H(\mathrm{X}|\mathrm{Z})H(\mathrm{X}|\mathrm{YZ}),$$
(33)
where $`H(\mathrm{X}|\mathrm{Y})`$ denotes the conditional Shannon entropy of $`\mathrm{X}`$ given $`\mathrm{Y}`$. In terms of the conditional mutual information, Shannon proved the following inner and outer bounds on the capacity of a two-way classical channel:
###### Theorem 3.13 (Shannon ).
If $`(R_{},R_{})`$ is achievable, then
$$(R_{},R_{})conv\{[[I(\mathrm{A};\mathrm{B}^{}|\mathrm{B}),I(\mathrm{B};\mathrm{A}^{}|\mathrm{A})]]:p_{ab}\text{ arbitrary}\}.$$
(34)
Conversely, if
$$(R_{},R_{})conv\{[[I(\mathrm{A};\mathrm{B}^{}|\mathrm{B}),I(\mathrm{B};\mathrm{A}^{}|\mathrm{A})]]:p_{ab}=p_aq_b\},$$
(35)
then $`(R_{},R_{})`$ is achievable.
Now consider the corresponding outer and inner bounds from Theorems 3.1 and 3.9. Let $`=\{p_{ab},|aa_{\mathrm{AA}^{}}|bb_{\mathrm{BB}^{}}\}`$ where $`|a,|b`$ are mutually orthogonal states on systems $`\mathrm{A},\mathrm{B}`$ and as before, the senders retain copies of their inputs in systems $`\mathrm{A}^{},\mathrm{B}^{}`$. The action of the two-way classical channel $`𝒩`$ on this ensemble is
$$𝒩(|aa||bb|)=\underset{a^{},b^{}}{}p_{a^{}b^{}|ab}|a^{}a^{}||b^{}b^{}|.$$
(36)
It is straightforward to compute
$`\chi (tr_{\mathrm{AA}^{}})`$ $`=H(\mathrm{B})`$ (37)
$`\chi (tr_{\mathrm{AA}^{}}𝒩)`$ $`=H(\mathrm{BB}^{})H(\mathrm{BB}^{}|\mathrm{AB})=H(\mathrm{B})+I(\mathrm{B}^{};\mathrm{A}|B).`$ (38)
Therefore,
$`\mathrm{\Delta }\chi _{}()`$ $`=I(\mathrm{A};\mathrm{B}^{}|\mathrm{B}).`$ (39)
Thus, we see that in the classical case, the outer bound of Theorem 3.1 and the inner bound of Theorem 3.9 are identical to Shannon’s outer and inner bounds, respectively.
The equivalence to Shannon’s bounds shows that in general, the bounds of Theorems 3.1 and 3.9 are not tight; even in the classical case, the inner bound may be exceeded and the outer bound may not be achievable . Such results for classical channels might provide insight into how the general (quantum) bounds could be tightened.
## 4 Additivity results for one-way channels
Since one-way quantum channels are simply special cases of two-way channels, it is possible to obtain results about one-way channels by thinking of them as two-way channels. In this section, we use such an approach to rederive two previously known additivity results for one-way channels: the entanglement-assisted capacity of an arbitrary one-way quantum channel, and the entanglement-unassisted capacity of an entanglement-breaking one-way quantum channel.
In this section, $``$ denotes a one-way channel, and the classical capacity is simply the maximum value of $`R_{}`$ in the achievable region (at fixed $`R_\mathrm{e}`$). We use $`R()`$ and $`R^E()`$ to denote the classical capacity with no entanglement assistance and unlimited entanglement assistance, respectively.
### 4.1 Entanglement-assisted capacity of one-way channels
A general expression for $`R^E()`$ in terms of the quantum mutual information was found in Ref. 7. Furthermore, Ref. 7 proved that $`R^E`$ is strongly additive, i.e., $`R^E(_1_2)=R^E(_1)+R^E(_2)`$ for any pair of (one-way) quantum channels $`_1,_2`$.
The original proof of additivity used entropy inequalities to show that the explicit expression for $`R^E`$ is indeed additive . But specializing Corollary 3.12 to one-way quantum channels provides an immediate alternative proof of strong additivity.
These two proofs appear to be inequivalent. The simplicity of proving additivity via Corollary 3.12 seems to follow from the structure of the protocol of Ref. 21 (or equivalently, that in Lemma 3.6). The main idea of this protocol, to borrow a resource and later regenerate some or more of it, has recently found a number of applications in quantum information theory . Such a protocol gives rise to a coding structure very different from more standard, direct techniques, such as those used in Ref. 7.
From Corollary 3.11, we see that the capacity expression of Ref. 7 in terms of the quantum mutual information can be written as a supremum of $`\mathrm{\Delta }\chi `$. It is not obvious simply by looking at these two expressions that they are in fact equal.
### 4.2 Unassisted capacity of one-way entanglement-breaking channels
We now turn our attention to unassisted classical communication using a one-way channel. In particular, we consider entanglement-breaking channels, which are guaranteed to output a state that is unentangled between the sender and the receiver. Using the framework of two-way channels, we will prove a special case of the following result:
###### Theorem 4.1 (Shor ).
If $``$ is an arbitrary one-way quantum channel and $`^{}`$ is an entanglement-breaking one-way quantum channel, then
$$R(^{})=R()+R(^{}).$$
(40)
We will prove this result in the special case in which both $``$ and $`^{}`$ are entanglement-breaking. In particular, this includes the case $`=^{}`$, demonstrating the additivity of the Holevo capacity of an entanglement-breaking channel. The proof in terms of two-way channels for this special case is significantly simpler.
As in Shor’s proof , we use strong subadditivity of the von Neumann entropy in various guises. In particular, we will use the following lemma:
###### Lemma 4.2.
Let $`\{\sigma _i\},\{\eta _i\}`$ be sets of quantum states, and let $`\{p_i\}`$ be a probability distribution. Then
$$S(_ip_i\sigma _i\eta _i)S(_ip_i\sigma _i)+_ip_iS(\eta _i).$$
(41)
We give two proofs of this lemma: an operational proof and a proof that uses strong subadditivity directly.
###### Proof 4.3 ( 1 (Operational)).
Let $`_1=\{p_i,\sigma _i\}`$ and $`_2=\{p_i,\sigma _i\eta _i\}`$. Since $`_1`$ can be obtained from $`_2`$ by discarding the second system,
$`0`$ $`\chi (_2)\chi (_1)`$ (42)
$`=S(_ip_i\sigma _i\eta _i)_ip_iS(\eta _i)S(_ip_i\sigma _i),`$ (43)
where the last line is obtained by using the definition Eq. (9) and the fact that $`S(\sigma \eta )=S(\sigma )+S(\eta )`$.
###### Proof 4.4 ( 2 (Direct use of strong subadditivity)).
For the state $`\rho _{ABC}:=_ip_i\sigma _{i,A}\eta _{i,B}|ii|_C`$,
$`S_{ABC}`$ $`=H(\{p_i\})+_ip_iS(\sigma _i\eta _i)`$ (44)
$`S_{AB}`$ $`=S(_ip_i\sigma _i\eta _i)`$ (45)
$`S_{AC}`$ $`=H(\{p_i\})+_ip_iS(\sigma _i)`$ (46)
$`S_A`$ $`=S(_ip_i\sigma _i).`$ (47)
Equation (41) then follows from the strong subadditivity inequality $`S_{ABC}+S_AS_{AB}+S_{AC}`$.
###### Proof 4.5 ( of Theorem 4.1 for $`𝓜`$ entanglement-breaking).
The idea of the proof is to show that the states that can be output by either channel are of no use in enhancing the capacity of the other channel, and hence that the capacity of the joint channel is simply the sum of the individual capacities.
For any entanglement-unassisted protocol that uses only entanglement-breaking channels, we can rerun our proof of the outer bound in Sec. 3.1 restricting to ensembles of separable states. Thus we have an upper bound analogous to Eq. (24),
$$R()\mathrm{max}\{\mathrm{\Delta }\chi _{}():\text{separable }\}.$$
(48)
Applying the two-way channel formalism to a one-way channel $``$ with input system $`\mathrm{A}`$ and output system $`\mathrm{B}`$, the most general input and output ensembles are $`_{\mathrm{in}}=\{p_i,\rho _{i,\mathrm{A}^{}\mathrm{AB}^{}}\}`$ and $`_{\mathrm{out}}=\{p_i,(\rho _{i,\mathrm{A}^{}\mathrm{AB}^{}})\}`$. Without loss of generality, we can omit the system $`\mathrm{A}^{}`$. This system does not appear in the bound on the communication rate in terms of $`\mathrm{\Delta }\chi _{}`$, and there is no entanglement to be stored in $`\mathrm{A}^{}`$. Thus, the optimal input ensemble can be restricted to have the form $`_{\mathrm{in}}=\{p_i,\rho _{i,\mathrm{AB}^{}}\}`$ with $`_{\mathrm{out}}=\{p_i,(\rho _{i,\mathrm{AB}^{}})\}`$.
We will be interested in three cases where the form of the input ensemble is restricted to different extents. Let $`\mathrm{\Delta }\chi ^S`$, $`\mathrm{\Delta }\chi ^P`$, $`\mathrm{\Delta }\chi ^0`$ denote $`sup_{_{\mathrm{in}}}[\chi (_{\mathrm{out}})\chi (tr_\mathrm{A}_{\mathrm{in}})]`$ for $`_{\mathrm{in}}`$ ranging over
$`_{\mathrm{in}}^S`$ $`:=\{p_i,_jq_{ij}\rho _{ij,\mathrm{A}}\sigma _{ij,\mathrm{B}^{}}\},`$ (49)
$`_{\mathrm{in}}^P`$ $`:=\{p_i,\rho _{i,\mathrm{A}}\sigma _{i,\mathrm{B}^{}}\},`$ (50)
$`_{\mathrm{in}}^0`$ $`:=\{p_i,\rho _{i,\mathrm{A}}\}`$ (51)
for separable ensembles, product ensembles, and ensembles with $`\chi =0`$, respectively. It is clear that $`\mathrm{\Delta }\chi ^0\mathrm{\Delta }\chi ^P\mathrm{\Delta }\chi ^S`$; we will show that $`\mathrm{\Delta }\chi ^S\mathrm{\Delta }\chi ^P\mathrm{\Delta }\chi ^0`$, so that in fact, all three quantities are equal.
First we show that $`\mathrm{\Delta }\chi ^S\mathrm{\Delta }\chi ^P`$. For any separable ensemble in the form of Eq. (49), let $`\eta _{ij,\mathrm{A}}:=(\rho _{ij,\mathrm{A}})`$. Then we have
$`\chi (_{\mathrm{out}})\chi (tr_\mathrm{A}_{\mathrm{in}})`$ $`=S(_{ij}p_iq_{ij}\sigma _{ij,\mathrm{B}^{}}\eta _{ij,\mathrm{A}})_ip_iS(_jq_{ij}\sigma _{ij,\mathrm{B}^{}}\eta _{ij,\mathrm{A}})`$
$`S(_{ij}p_iq_{ij}\sigma _{ij,\mathrm{B}^{}})+_ip_iS(_jq_{ij}\sigma _{ij,\mathrm{B}^{}})`$ (52)
$`S(_{ij}p_iq_{ij}\sigma _{ij,\mathrm{B}^{}}\eta _{ij,\mathrm{A}})_{ij}p_iq_{ij}S(\sigma _{ij,\mathrm{B}^{}}\eta _{ij,\mathrm{A}})`$
$`S(_{ij}p_iq_{ij}\sigma _{ij,\mathrm{B}^{}})+_{ij}p_iq_{ij}S(\sigma _{ij,\mathrm{B}^{}})\mathrm{\Delta }\chi ^P`$ (53)
where Lemma 4.2 has been applied to each term in $`_ip_iS()`$. Thus, a separable ensemble is no better than a product ensemble.
Now we show that $`\mathrm{\Delta }\chi ^P\mathrm{\Delta }\chi ^0`$. For any product ensemble in the form of Eq. (50), let $`\eta _{i,\mathrm{A}}:=(\rho _{i,\mathrm{A}})`$. Then
$`\chi (_{\mathrm{out}})\chi (tr_\mathrm{A}_{\mathrm{in}})`$ $`=S(_ip_i\sigma _{i,\mathrm{B}^{}}\eta _{i,\mathrm{A}})_ip_iS(\sigma _{i,\mathrm{B}^{}}\eta _{i,\mathrm{A}})`$
$`S(_ip_i\sigma _{i,\mathrm{B}^{}})+_ip_iS(\sigma _{i,\mathrm{B}^{}})`$ (54)
$`S(_ip_i\eta _{i,\mathrm{A}})_ip_iS(\eta _{i,\mathrm{A}})\mathrm{\Delta }\chi ^0`$ (55)
where the inequality is due to subadditivity of $`S(_ip_i\sigma _i\eta _i)`$ and additivity of $`S(\sigma \eta )`$. Thus a product ensemble with $`\chi 0`$ is no better than one with $`\chi =0`$.
This argument shows that we can assume without loss of generality that the input ensemble is of the form of Eq. (51). But such an ensemble costs nothing to create, so by using the protocol of Lemma 3.6, we see that the capacity $`\mathrm{\Delta }\chi ^0`$ can be achieved for any ensemble $`_{\mathrm{in}}^0`$, and
$$R()=\underset{_{\mathrm{in}}^0}{sup}\mathrm{\Delta }\chi ^0.$$
(56)
Finally, consider the capacity of the combined channel $`^{}`$ where both $``$ and $`^{}`$ are entanglement-breaking. Without loss of generality, we can assume that the channels act sequentially. Each channel can only produce separable output states, which by the above argument are no better than states with $`\chi =0`$, which can be produced at zero cost. Therefore the capacity of the combined channel is simply the sum of the individual capacities.
## 5 Open questions
In this paper, we have established simple inner and outer bounds on the entanglement-assisted classical capacity region of a two-way quantum channel, and we have applied the framework of two-way channels to rederive two previous additivity results for one-way channels. However, since the two-way channel framework includes a wide variety of disparate communication scenarios as special cases, this work raises many more questions than it answers.
Calculating the capacity region for any particular channel can be a challenging problem, and has not been done except in a few particular special cases. In fact, even computing the inner and outer bounds given in this paper can be difficult, since the ancillary state spaces $`\mathrm{A}^{},\mathrm{B}^{}`$ may be arbitrarily large, and we do not know that low-dimensional ancillas are sufficient to achieve the capacity (even in the unitary case ).
Although calculating the precise capacity region may be difficult, a more modest goal is to improve upon the inner and outer bounds given in this paper. In particular, known classical bounds that improve upon Shannon’s bounds might be useful for finding improved quantum bounds. Also, it would be interesting to find conditions under which the inner and outer bounds coincide. Shannon showed that his inner and outer bounds coincide for certain kinds of symmetric two-way classical channels , so it is plausible to suppose that a similar result might hold in the quantum case.
We have primarily considered the case of unlimited entanglement assistance, but it would be interesting to consider unassisted communication as well as the general case of finite entanglement assistance. Recently, Shor has given a protocol for classical communication through a one-way quantum channel with limited entanglement assistance that interpolates between the HSW capacity and the entanglement-assisted capacity . In addition, Harrow has obtained an expression for the one-way classical communication capacity with finite entanglement assistance for unitary two-way channels . It would be interesting to generalize these results to arbitrary two-way channels.
Another approach is to consider particular families of channels to see whether the capacity region is simpler for those channels. One such family is the set of two-way entanglement-breaking channels. There are several possible definitions of a two-way entanglement-breaking channel, but perhaps the simplest is that the output state should be triseparable between Alice’s output, Bob’s output, and any ancillas. Unfortunately, it is not even clear how to characterize such channels (as can be done for one-way entanglement-breaking channels ). Some results have been obtained for other families of two-way quantum channels, such as unitary two-way channels and feedback channels . Another special class of two-way channels, those that simply distribute bipartite states, have been much better understood .
One way to obtain a better understanding of channel capacities is to consider the problem of simulating a channel using a certain amount of communication (in each direction) and entanglement. Such reverse theorems have been studied for one-way classical and quantum channels, and more recently for feedback channels with restricted input sources . In particular, reverse theorems can be useful for establishing bounds on capacities . However, simple reverse theorems for general two-way channels seem unlikely to exist. For example, the communication costs (in each direction) of any simulation must exceed the corresponding one-way capacities, since the simulated channel can be used to achieve the one-way capacity in either direction. As another example, the set of causal operations that are not localizable cannot produce any communication, but any such operation requires communication in at least one direction to simulate, even with entanglement assistance.
Finally, note that we have completely avoided the problems of communicating quantum information through a two-way quantum channel and of multi-way communication through $`k`$-partite quantum operations with $`k>2`$. These problems present further challenges for understanding the the capabilities of quantum communication channels.
## Acknowledgments
We thank Aram Harrow for many helpful discussions, especially regarding the protocol for achieving the entanglement-assisted one-way capacity of a two-way channel. Thanks also to Charles Bennett, whose insights on entanglement-assisted communication via entanglement breaking channels and the simulation of nonlocal boxes saved us from wandering in unfruitful directions. This work was initiated while AMC and DWL were at the IBM T. J. Watson Research Center and HKL was at Magiq Technologies. Part of this work was done while AMC was at the MIT Center for Theoretical Physics and while DWL was a visiting researcher at MSRI. AMC was supported in part by the Fannie and John Hertz Foundation, the Cambridge–MIT Foundation, the DOE under cooperative research agreement DE-FC02-94ER40818, and the NSA and ARDA under ARO contract DAAD19-01-1-0656. AMC and DWL received support from the NSF under Grant No. EIA-0086038. DWL received support from the Tolman Foundation and the Croucher Foundation. HKL received support from NSERC, the CRC Program, CFI, OIT, PREA, and CIPI.
|
warning/0506/cond-mat0506678.html
|
ar5iv
|
text
|
# Coulombian Disorder in Periodic Systems
## I Introduction
The effect of quenched disorder on various condensed elastic systems is one of the fascinating problems in statistical mechanics. Examples of physical systems range from domain walls in magnetic and ferroelectric materialslemerle\_domainwall\_creep ; tybell\_ferroelectric , contact lines of of a liquid meniscus on a rough substratemoulinet\_contact\_line , crack propagationbouchaud\_bouchaud , to vortex lattices in type II superconductorsblatter\_vortex\_review ; giamarchi\_vortex\_cargese , charge density waves (CDW)gruner\_book\_cdw ; nattermann\_brazovskii and Wigner crystalsdeville\_wigner ; williams\_wigner ; coupier\_wigner ; giamarchi\_varenna . In these systems, the competition between elastic interactions which tend to impose some long range order in the system and quenched disorder, leads to the formation of glassy phases. Two broad classes of elastic systems can be distinguished: random manifold systems such as domain walls, contact lines and cracks, and periodic systems such as vortex lattices, charge density waves, and Wigner crystals. The latter are characterized by a long range crystalline order in the absence of disorder and thermal fluctuations. For these systems, a crucial question is whether a weak disorder destroys entirely the crystalline order, or whether some remnants of the underlying periodic structure remain observable.
One of the earliest attempts to answer this question, was the pioneering work by Larkinlarkin\_70 on vortex lattices. Using a random-force model, he showed that due to the relevance of disorder in the renormalization group sense, long range order was entirely destroyed below four dimensions. Above four dimensions, long range order persists as disorder becomes irrelevant. A similar conclusion was reached by Sham and Patton for the case of a CDW with short range elasticity sham\_peierls\_disorder , where, using an Imry-Ma approach,imry\_ma they concluded that long range order was impossible in the presence of disorder below four dimensions. The problem of short range disorder in periodic systems with short ranged elasticity was revisited in Refs. villain\_cosine\_realrg, ; nattermann\_pinning, ; giamarchi\_vortex\_short, ; giamarchi\_vortex\_long, . It was argued that the periodicity present in systems like CDW and vortex lattices plays a pivotal role in determining the physics of the system in the presence of disorder. More precisely, it was shown that, though the disorder is relevant below four dimensions, due to the underlying periodicity of the system a quasi long-range order persisted for dimensions between two and four. This is in stark contrast to the earlier results which predicted a total destruction of order. The resulting phase, nicknamed Bragg glass phase, possesses both quasi long range order and metastability and glassy properties.giamarchi\_vortex\_short ; giamarchi\_vortex\_long It was further shown that the Bragg glass phase is stable to the formation of defects. giamarchi\_vortex\_long ; gingras\_dislocations\_numerics ; fisher\_bragg\_proof ; zeng\_fisher\_nobglass2d Recent neutron scattering experiments on vortex lattices have furnished clear evidence for the existence of such a phase.klein\_neutron
A complication arises in charged periodic systems due to the Coulomb repulsion, which renders the elasticity non-local.efetov\_larkin\_replicas ; lee\_coulomb\_cdw ; bergman\_coulomb\_disorder ; chitra\_wigner\_hall This non local elasticity tends to rigidify the system, so that short range correlated disorder could be irrelevant in dimension smaller than four.chitra\_vortex For instance, within the random force model, the correlation function of the displacement in three dimensions displays a logarithmic growth indicating quasi-long range order.efetov\_larkin\_replicas ; bergman\_coulomb\_disorder ; lee\_coulomb\_cdw In fact, when the periodic structure of the CDW is properly taken into account, the growth of the displacement correlation function is even weaker, increasing only as $`\mathrm{log}(\mathrm{log}(r))`$ with the distance $`r`$.rosso\_cdw\_long A second complication arising from Coulomb interaction is that the disorder induced by charged impurities has long-range correlations. This type of disorder can exist in certain doped CDW materialsrouziere\_friedel\_cdw such as $`\mathrm{KMo}_{1\mathrm{x}}\mathrm{V}_\mathrm{x}\mathrm{O}_3`$.
In this paper, we study the effect of the competition of the non-local elasticity produced by the Coulomb interaction with the long-range random potential resulting from the presence of charged impurities on the statics. The paper is organized as follows: in Sec. II, we introduce a decomposition of the Coulomb potential on the Fourier modes of the periodic structure. With this decomposition, we show that only the long-wavelength component of the random potential i.e., forward scattering disorder, possesses long range correlations. Using statistical tilt symmetrynarayan\_fisher\_cdw , we deduce that due to the short ranged nature of the backward scattering terms engendered by the disorder, the dynamical properties in the presence of charged impurities are not qualitatively different from those in the presence of neutral short ranged impurities. In Sec. III, we consider the problem of the CDW system, and we derive the non-local elastic Hamiltonian. In Sec. IV, we derive the static displacement correlation functions and x-ray intensity of the CDW with charged impurities and we highlight the similitude of the latter to the x-ray intensity of smectic-A liquid crystals subjected to thermal fluctuations.caille\_smectic\_xray In Sec. V, we discuss the experimental significance of our result and suggest that the smectic-like correlations should be observable in experiments on $`\mathrm{KMo}_{1\mathrm{x}}\mathrm{V}_\mathrm{x}\mathrm{O}_3`$. Finally, we summarize the possible behavior of the static correlators in a pinned Charge Density wave according to the local or non-local character of elasticity and the presence or absence of charged impurities.
## II Elasticity and Disorder in Periodic Systems
In this section, we discuss how Coulomb interactions affect elasticity and disorder in periodic systems. For a periodic elastic structure, the density can be written as:
$`\rho (𝐫)=\rho _0(𝐫)+{\displaystyle \underset{𝐆}{}}e^{i𝐆(𝐫𝐮(𝐫))},`$ (1)
where $`\rho _0(𝐫)=\rho _0(1𝐮)`$ describes the density fluctuation arising from the long wavelength deformation of the periodic structure and $`\rho _0`$ is the average density. In the second term, the vectors $`𝐆`$ belong to the reciprocal lattice of the perfect periodic structure, and $`𝐮(𝐫)`$ represents a slowly varying<sup>1</sup><sup>1</sup>1by slowly varying, we mean that the Fourier components of $`u`$ are different from zero only for wavevectors much smaller than $`|𝐆_{\text{min.}}|`$. elastic deformation of the structure.giamarchi\_vortex\_long The quantities $`e^{i𝐆(𝐫𝐮(𝐫))}`$ describe fluctuations of the density on the scale of a lattice spacing. The low energy physics of the periodic structure can be described in terms of a purely elastic Hamiltonian which has the generic form for isotropic systems
$$H_0=_𝐫\frac{c}{2}(𝐮)^2$$
(2)
where $`c`$ is the elastic coefficient and $`_𝐫`$ is a shorthand for $`d^3𝐫`$. This form can easily be generalized to anisotropic systems. Well known examples of charged periodic structures are the Wigner crystal,wigner ; wigner\_crystal ; deville\_wigner ; williams\_wigner ; coupier\_wigner charged colloidal crystals,colloids and the charge density waves.froehlich\_cdw ; peierls\_inst ; denoyer\_cdw\_ttf-tcnq In many charged systems, unscreened Coulomb interactions are present:
$`H_C={\displaystyle \frac{e^2}{2}}{\displaystyle _{𝐫,𝐫^{}}}{\displaystyle \frac{\rho (𝐫)\rho (𝐫^{})}{4\pi ϵ|𝐫𝐫^{}|}},`$ (3)
and strongly affect the elasticity and dispersion of the compression modes of the system. Moreover, in the presence of charged impurities, the original charge density on the lattice interacts with the charge impurity yielding:
$`H_{\text{dis.}}=e^2{\displaystyle _{𝐫,𝐫^{}}}{\displaystyle \frac{\rho (𝐫)\rho _{\text{imp.}}(𝐫^{})}{4\pi ϵ|𝐫𝐫^{}|}},`$ (4)
where $`\rho _{\text{imp.}}`$ denote the impurity density. Using the decomposition of the density (1), we now show that the Coulomb interactions fundamentally modify only the long wavelength components of the elasticity and of the disorder energy.
To better handle the periodicity of the elastic structure, it is convenient to use the decomposition of the Coulomb interaction in terms of the reciprocal lattice vectors $`G`$. In three dimensions, this decomposition reads
$`{\displaystyle \frac{1}{4\pi |𝐫|}}`$ $`=`$ $`{\displaystyle \frac{d^3𝐪}{(2\pi )^3}\frac{e^{i𝐪𝐫}}{𝐪^2}}`$ (5)
$`=`$ $`{\displaystyle \underset{𝐆}{}}e^{i𝐆𝐫}V_𝐆(𝐫)`$
where
$`V_𝐆(𝐫)={\displaystyle _{BZ}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}}{\displaystyle \frac{e^{i𝐪𝐫}}{(𝐪+𝐆)^2}},`$ (6)
and $`_{BZ}`$ indicates that the integral is restricted to the first Brillouin zone. It is straightforward to check that $`V_𝐆(𝐫)=V_𝐆^{}(𝐫)`$. Using Eq. (5), the interaction term $`H_C`$ can be rewritten as:
$`H_C`$ $`=`$ $`{\displaystyle \frac{e^2}{2ϵ}}{\displaystyle \underset{G0}{}}{\displaystyle _{𝐫,𝐫^{}}}V_𝐆(𝐫𝐫^{})e^{i𝐆(𝐮(𝐫)𝐮(𝐫^{}))}`$ (7)
$`+{\displaystyle \frac{e^2}{2ϵ}}{\displaystyle _{𝐫,𝐫^{}}}V_0(𝐫𝐫^{})\rho _0(𝐫)\rho _0(𝐫^{}).`$
Note that due to the slow variation of $`𝐮(𝐫)`$, terms involving the oscillatory factors $`e^{i(𝐆𝐆^{})𝐫}`$ can be dropped from the interaction. Let us first consider the term involving long wavelength fluctuations of the density. Since we are interested only in the long wavelength properties, we can replace the integration over the Brillouin zone in $`V_0(𝐫)`$. by a Gaussian integration:
$`{\displaystyle _{BZ}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}e^{a^2𝐪^2}},`$ (8)
with the parameter $`a`$ chosen so that $`\pi /a|𝐆_{\text{min.}}|`$, $`𝐆_{\text{min.}}`$ being the reciprocal lattice vector having the shortest length. In this case, $`V_0(𝐫)`$ can be obtained indirectly by solving the Poisson equation with a Gaussian charge density and is found to be
$`V_0(𝐫)={\displaystyle \frac{1}{4\pi r}}\mathrm{erf}\left({\displaystyle \frac{r}{2a}}\right),`$ (9)
In the limit $`ra`$, we recover the known result $`V_01/(4\pi r)`$. Clearly, the non-oscillating component of the Coulomb potential remains long-ranged and tends to rigidify the system.
It now remains to be seen whether the oscillating parts of the Coulomb interaction specified by $`V_𝐆`$ for $`𝐆`$ are long-ranged or not. We first note that the above trick of replacing the integration over the Brillouin zone by a Gaussian integral over the entire space is not applicable anymore, as it would introduce a spurious integration over a region where $`𝐆+𝐪=0`$. This would result in an (incorrect) $`1/r`$ behavior of $`V_{𝐆\mathrm{𝟎}}(r)`$. To obtain a correct estimate for $`V_𝐆`$ we replace the integral over the Brillouin zone by
$`{\displaystyle _{BZ}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}F_{BZ}(𝐪,ϵ)},`$ (10)
where $`F_{BZ}(q,ϵ)`$ is an indefinitely derivable function with a compact support contained in the first Brillouin zone (see App. A for an explicit form of $`F_{BZ}(q,ϵ)`$).<sup>2</sup><sup>2</sup>2A more straightforward approach would be to keep the hard cutoff at the edge of the Brillouin zone. Then, the function $`V_𝐆`$ would decay as $`1/r^2`$ with an oscillating prefactor. The same oscillation would be also obtained for a short ranged potential, and is only a consequence of the hard cutoff. Obviously, for $`𝐆0`$, $`F_{BZ}(q,ϵ)/|𝐪+𝐆|^2`$ is also an indefinitely differentiable function of compact support. A well-known theoremgelfand64\_theorem ; reed75\_mmmp then shows that the Fourier transform of $`F_{BZ}(q,ϵ)/|𝐪+𝐆|^2`$ is indefinitely differentiable and for $`r\mathrm{}`$ is $`o(1/r^n)`$ for any $`n>0`$. This implies that the function $`V_𝐆(r)`$ is short ranged. Incorporating the above results in Eq.(7), we see that while the non oscillating part of the Coulomb interaction modifies the long wavelength behavior of the elasticity, rendering it non-local, the short ranged nature of the oscillatory terms merely renormalizes the elastic coefficients. This is explicitly shown in App. B for the particular case of a CDW. The resulting non-local character of the elastic interactions modifies strongly the static and dynamic properties of the system.chitra\_wigner\_hall ; chitra\_wigner\_long ; rosso\_cdw\_long ; chitra\_wigner\_zero
To understand the nature of the interaction with the charged impurities, we use the above procedure to rewrite the random potential generated by the impurities as:
$`U_{\text{imp.}}(r)`$ $`=`$ $`{\displaystyle _{BZ}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}}{\displaystyle \frac{e^{i𝐪𝐫}\rho _{\text{imp.}}(𝐪)}{𝐪^2}},`$ (11)
$`=`$ $`{\displaystyle \underset{𝐆}{}}e^{i𝐆𝐫}U_𝐆(r).`$
Using this in Eq.(4), the interaction of the system with the random potential is given by
$`H_{\text{dis.}}`$ $`=`$ $`{\displaystyle \frac{e^2}{ϵ}}{\displaystyle \underset{G0}{}}{\displaystyle _{𝐫,𝐫^{}}}U_𝐆(𝐫)e^{i𝐆𝐮(𝐫)}`$ (12)
$`+{\displaystyle \frac{e^2}{ϵ}}{\displaystyle _{𝐫,𝐫^{}}}U_0(𝐫)\rho _0(𝐫),`$
In Eq. (12), the interaction of $`\rho _0`$ with the random potential $`U_0`$ is called forward scattering, and the terms containing $`e^{i𝐆𝐮(𝐫)}`$ are called backward scattering. This nomenclature originates in the theory of electrons in 1D random potential.giamarchi\_book\_1d To calculate the disorder correlation functions, we consider the case of Gaussian distributed impurities where itzykson\_stat2 $`\overline{\rho _{\text{imp.}}(𝐆+𝐪)\rho _{\text{imp.}}(𝐆^{}+𝐪^{})}=(2\pi )^3D\delta _{𝐆,𝐆^{}}\delta (𝐪+𝐪^{})`$, the parameter $`D`$ measuring the disorder strength. Consequently, we find that for $`𝐆0`$:
$$\overline{U_𝐆(𝐫)U_𝐆(𝐫^{})}=D_{BZ}\frac{d^3𝐪}{(2\pi )^3}\frac{e^{i𝐪(𝐫𝐫^{})}}{(𝐪+𝐆)^4},$$
(13)
Using the same arguments as before, we infer that the correlations of $`U_𝐆(𝐫)`$ are short ranged, as in the case of neutral impurities, for all $`𝐆`$ except $`𝐆=0`$. This implies that the backward scattering terms induced by disorder are short-ranged and the treatment of these terms within the replica or the Martin-Siggia-Rosemartin\_siggia\_rose ; dominicis\_dynamics methods is identical to the case of neutral or screened impurities. However, the $`𝐆=0`$ component
$`U_0(𝐫)={\displaystyle _{BZ}}{\displaystyle \frac{d^3𝐪}{(2\pi )^3}}{\displaystyle \frac{e^{i𝐪𝐫}}{𝐪^2}}\rho _{\text{imp.}}(𝐪),`$ (14)
manifests power law decay of the forward scattering correlations. This term however can be gauged out by the statistical tilt symmetrynarayan\_fisher\_cdw , and affects mainly the static properties of the periodic system. Typically, in periodic systems with both short range disorder and local elasticity, the contribution of the forward scattering disorder can be neglected and it is the backward scattering that induces collective pinning and Bragg glass features like a quasi order in the static correlation functions. Here, we have shown that even in the case of long range disorder, the backward scattering terms behave essentially like their short ranged (neutral impurities) counterparts. However, the effect of the forward scattering terms on the correlation functions has to be studied carefully. In the next section, we show that in the case of charged impurities in a charge density wave system, the forward scattering term strongly modifies the static correlation. Finally, we remark that our decomposition of the elastic energy and the impurity potential is not exclusive to the Coulomb potential and is applicable to other long-range potentials. As a result, the conclusions of the present sections are expected to be valid for more general long range potentials.
## III Charge Density Waves
In this section, we re-derive the elastic Hamiltonian for a $`d`$ dimensional CDW with screened Coulomb interactions between the density fluctuations at zero temperature. We consider an incommensurate CDW, in which the electron density is modulated by a modulation vector $`Q`$ incommensurate with the underlying crystal lattice. In this phase, the electron density has the following form nattermann\_brazovskii :
$$\rho (𝐫)=\rho _0+\frac{\rho _0}{Q^2}𝐐\varphi (𝐫)+\rho _1\mathrm{cos}(𝐐𝐫+\varphi (𝐫))),$$
(15)
where $`\rho _0`$, is the average electronic density (see App. C for details). The second term in Eq. (15) is the long wavelength density and corresponds to variations of the density over scales larger than $`Q^1`$ . The last oscillating term describes the sinusoidal deformation of the density at a scale of the order of $`Q^1`$ induced by the formation of the CDW with amplitude $`\rho _1`$ and phase $`\varphi `$.
In the absence of Coulomb interactions, the low energy properties of the CDW can be described by an effective Hamiltonian for phase fluctuations. For CDW aligned along the $`x`$ axis, i.e., $`𝐐=Q\widehat{x}`$, this phase only Hamiltonian readsfukuyama\_cdw\_pinning ; gruner\_revue\_cdw ; feinberg\_cdw ; maki\_phase\_hamiltonian :
$`H_0={\displaystyle \frac{\mathrm{}v_Fn_c}{4\pi }}{\displaystyle d^3𝐫\left[(_x\varphi )^2+\frac{v_y^2}{v_x^2}(_y\varphi )^2+\frac{v_z^2}{v_x^2}(_z\varphi )^2\right]},`$
where $`v_F`$ is the Fermi velocity and $`n_c`$ is the number of chains per unit surface that crosses a plane orthogonal to $`Q`$. The velocity of the phason excitations parallel to $`Q`$ is $`v_x=(m_e/m_{})^{1/2}v_F`$ with $`m^{}`$ the effective mass of the CDW and $`m_e`$ the mass of an electron. $`v_y`$ and $`v_z`$ denote the phason velocities in the transverse directions. A crucial observation is that a deformation of $`\varphi (r)`$ along $`Q`$ produces an imbalance of the electronic charge density which then augments the electrostatic energy due to Coulomb repulsion between density fluctuations. We evaluate the contribution of Coulomb interactions screened beyond the characteristic length $`\lambda `$ which accounts for the presence of free carriers. This length diverges in the limit $`T0`$.wong\_coulomb\_cdw ; virosztek\_collective\_cdw The electrostatic energy takes the form:
$`H_C={\displaystyle \frac{e^2}{8\pi ϵ}}{\displaystyle d^d𝐫d^d𝐫^{}e^{|𝐫𝐫^{}|/\lambda }\frac{\rho (𝐫)\rho (𝐫^{})}{|𝐫𝐫^{}|}},`$ (17)
where we have assumed for simplicity an isotropic dielectric permittivity $`ϵ`$ of the host medium.efetov\_larkin\_replicas ; lee\_coulomb\_cdw ; kurihara\_coulomb\_phasons ; wong\_coulomb\_cdw ; brazovskii\_longrange Due to the periodicity of the CDW system, we can use the decomposition of the Coulomb potential derived in Sec. II, obtaining:
$`H_C`$ $`=`$ $`{\displaystyle \frac{e^2\rho _0^2}{8\pi ϵQ^2}}{\displaystyle _{𝐫,𝐫^{}}}_x\varphi (𝐫){\displaystyle \frac{e^{|𝐫𝐫^{}|/\lambda }}{|𝐫𝐫^{}|}}_x^{}\varphi (𝐫^{})`$
$`+{\displaystyle \frac{e^2\rho _1^2}{2ϵ}}{\displaystyle _{𝐫,𝐫^{}}}\left[V_Q(𝐫𝐫^{})e^{i(\varphi (𝐫)\varphi (𝐫^{}))}+\text{c.c.}\right],`$
In Eq. (III), we have neglected the contribution of the higher harmonics of the CDW. Note that the oscillating terms, as discussed in App. B, only contribute to a renormalization of the coefficients in the short range elastic Hamiltonian (III) and thus can be neglected. However, the contribution of the long-wavelength term has more dramatic effects and reads
$`H_C`$ $`=`$ $`{\displaystyle \frac{e^2\rho _0^2}{2ϵQ^2}}{\displaystyle _{BZ}}{\displaystyle \frac{q_x^2}{\lambda ^2+q^2}}|\varphi (q)|^2,`$ (19)
in the three dimensional case. It is interesting to note that Coulomb interactions generate a non local elasticity i.e., a $`q`$-dispersion in the elastic constant.
The total Hamiltonian in $`d=3`$ now reads,
$`H_{\text{el.}}`$ $`=`$ $`H_0+H_C={\displaystyle \frac{1}{2}}{\displaystyle G^1(q)|\varphi (q)|^2},`$ (20)
$`G^1(q)`$ $`=`$ $`{\displaystyle \frac{n_c\mathrm{}v_F}{2\pi }}\left[{\displaystyle \frac{q_x^2}{(𝐪^2+\lambda ^2)\xi ^2}}+q_x^2+{\displaystyle \frac{v_y^2}{v_x^2}}q_y^2+{\displaystyle \frac{v_z^2}{v_x^2}}q_z^2\right]`$
where the lengthscale $`\xi `$ is defined by:
$`\xi ^2={\displaystyle \frac{n_c\mathrm{}v_F}{2\pi e^2\rho _0^2}}Q^2ϵ.`$ (21)
Depending on the ratio $`\lambda /\xi `$, two regimes of behavior can be identified: (i) Short ranged elasticity: when $`\lambda /\xi 1`$ the Coulomb correction to the short range elasticity is small even in the limit $`q0`$ and hence can be neglected.
(ii) Long range elasticity: for $`\lambda /\xi 1`$, the Coulombian correction to the short-range elasticity cannot be neglected. This regime is relevant at low temperatures, when the number of free carriers available to screen the Coulomb interaction is suppressed by the CDW gap.wong\_coulomb\_cdw ; virosztek\_collective\_cdw Mean field calculations show that this regime is obtained for temperatures $`T<0.2T_c`$ where $`T_c`$ is the Peierls transition temperature.virosztek\_collective\_cdw In the following, we focus on regime (ii), and accordingly, we take $`\lambda ^1=0`$ in Eq. (20).
## IV Forward Scattering
As discussed in Sec. I, the case of the short ranged elasticity has been studied by various authors. For charged periodic systems with short range disorder and a non-local elasticity generated by Coulomb interactions, it is known that $`d=3`$ becomes the upper critical dimension for disorder and the displacement correlations grow as $`B(r)=\mathrm{log}\mathrm{log}\mathrm{\Lambda }r`$.chitra\_vortex ; rosso\_cdw\_long Here, we study the effect of the long range disorder on the static correlations of a charged periodic system. Since, the backward scattering terms generated by such a disorder are short ranged, they lead to the same physics as that of short ranged disorder with the corresponding nonlocal elasticity. These terms contribute a $`\mathrm{log}\mathrm{log}r`$ term to the displacement correlations. However, in this case a simple dimensional analysis shows that the forward scattering terms generate the leading contribution to the correlation functions. In the following, we calculate the contribution of the forward scattering disorder to the displacement correlation function in the $`d=3`$ CDW.
### IV.1 Displacement correlation functions
The displacement correlation function is defined by:
$`B(r)`$ $`=`$ $`\overline{(\varphi (r)\varphi (0))^2}`$ (22)
$`=`$ $`{\displaystyle \frac{2}{L^6}}{\displaystyle \underset{𝐪}{}}\overline{\varphi (q)\varphi (q)}[1\mathrm{cos}𝐪𝐫].`$
The calculation of the correlation induced by the forward scattering disorder is analogous to the calculation of Larkin for the random force model.larkin\_70 Assuming an infinite screening length $`\lambda `$, the Hamiltonian reads:
$`H`$ $`=`$ $`H_{\text{el.}}+{\displaystyle \frac{e^2}{4\pi ϵ}}{\displaystyle d^3𝐫d^3𝐫^{}\frac{\rho _{\text{imp.}}(𝐫)\rho (𝐫^{})}{|rr^{}|}}`$
Using Eq. (15) in Eq. (IV.1), we obtain an expression of the form Eq. (12). Keeping only the forward scattering term we get:
$`H={\displaystyle \frac{d^3q}{(2\pi )^3}\left[\frac{G^1(q)}{2}|\varphi (q)|^2+\frac{i\rho _0e^2q_x}{Qϵq^2}\rho _{\text{imp.}}(q)\varphi (q)\right]}.`$
Shifting the field $`\varphi `$
$`\varphi (q)=\stackrel{~}{\varphi }(q)+{\displaystyle \frac{e^2\rho _0}{Qϵ}}{\displaystyle \frac{iq_xG(q)}{q^2}}\rho _{\text{imp.}}(q),`$ (25)
brings the Hamiltonian (IV.1) back to the form of Eq. (20). The average over disorder now yields:
$`\overline{\varphi (q)\varphi (q)}`$ $`=`$ $`\stackrel{~}{\varphi }(q)\stackrel{~}{\varphi }(q)`$ (26)
$`+{\displaystyle \frac{e^4\rho _0^2}{Q^2ϵ^2}}{\displaystyle \frac{q_x^2G(q)^2}{q^4}}\overline{\rho _{\text{imp.}}(q)\rho _{\text{imp.}}(q)}`$
$`=`$ $`L^3\left[TG(q)+{\displaystyle \frac{e^4\rho _0^2}{Q^2ϵ^2}}{\displaystyle \frac{q_x^2G(q)^2}{q^4}}D\right]`$
where $`\mathrm{}`$ and $`\overline{\text{}}`$ denote thermal average and disorder average respectively. Eq (25) shows that even in the presence of Coulombian disorder, the statistical tilt symmetrynarayan\_fisher\_depinning is preserved. This implies that in the presence of backward scattering disorder, the forward scattering term can be gauged out by (25), and the contribution of the forward scattering disorder $`B^{FS}`$ is simply added to the one obtained from the backward scattering disorder, $`B^{BS}`$.rosso\_cdw\_short ; rosso\_cdw\_long We conclude
$$B^{FS}(r)=2D\frac{e^4\rho _0^2}{Q^2ϵ^2}\frac{d^3q}{(2\pi )^3}\frac{q_x^2G(q)^2}{q^4}[1\mathrm{cos}(𝐪𝐫)].$$
(27)
We want to evaluate this integral for the case of $`v_y=v_z=v_{}`$. In the following we will use $`q_{}^2=q_y^2+q_z^2`$. To obtain the asymptotic behavior of $`B(r)`$ for $`r\mathrm{}`$ we need to consider the $`q0`$ limit of the integrand. The form of $`G(q)`$ suggests a scaling $`q_xq_{}^2`$ which then allows us to consider the integral:
$`F(𝐫)`$ $`=`$ $`{\displaystyle \frac{d^3𝐪}{(2\pi )^3}\frac{q_x^2}{[q_x^2+(\xi ^{}q_{}^2)^2]^2}(1\mathrm{cos}(𝐪𝐫))}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi \xi ^{}}}\{\mathrm{ln}[1+(\mathrm{\Lambda }_{}r_{})^2]+E_1\left({\displaystyle \frac{r_{}^2}{4|x|\xi ^{}}}\right)`$
$`+e^{r_{}^2/(4|x|\xi ^{})}\},`$
where $`\xi ^{}=\xi v_{}/v_x`$, $`r_{}^2=y^2+z^2`$ and $`\mathrm{\Lambda }_{}`$ is a momentum cut-off. A study of the limits of this function for $`r_{}\mathrm{}`$ and $`|x|\mathrm{}`$ shows that its asymptotic behavior is well described by:
$`F(𝐫){\displaystyle \frac{v_x}{16\pi v_{}\xi }}\mathrm{ln}\left({\displaystyle \frac{r_{}^2+4(v_{}|x|\xi /v_x)}{\mathrm{\Lambda }_{}^2}}\right).`$ (29)
Therefore, we have for $`r\mathrm{}`$:
$`B^{FS}(r)=\kappa \mathrm{ln}\left({\displaystyle \frac{r_{}^2+4(v_{}|x|\xi /v_x)}{\mathrm{\Lambda }_{}^2}}\right).`$ (30)
where:
$`\kappa ={\displaystyle \frac{DQ^2v_x}{16\pi \xi \rho _0^2v_{}}}.`$ (31)
The full asymptotic correlation function is given by the sum of the forward scattering contribution, Eq. (30), and the backward scattering contribution given in Eq. (51) of Ref. rosso\_cdw\_long, for the case of a short-range disorder and non-local elasticity:
$$B^{BS}(𝐫)=log(log(\mathrm{max}(\mathrm{\Lambda }|x|,(\mathrm{\Lambda }r_{})^2))).$$
(32)
Obviously, the contribution of the backward scattering terms is subdominant and can be neglected.
### IV.2 Analogy with smectics-A
We note that the result Eq. (29) can be obtained in the entirely different context of liquid crystals. If we consider a smectic-A liquid crystal, its elastic free energy readsdegennes\_liquid\_crystals ; landau\_elasticity ; chandrasekhar\_smectic ; pieranski\_book2 :
$`_{\text{el.}}={\displaystyle d^3𝐫\left[\frac{1}{2}B(_zu)^2+\frac{1}{2}k_{11}(\mathrm{\Delta }_{}u)^2\right]},`$ (33)
where $`u`$ represents the displacement of the smectic layers, $`B`$ is the compressibility, and $`k_{11}`$ measures the bending energy of the smectic layers. If we now assume a random compression force given by:
$`_{\text{dis.}}`$ $`=`$ $`{\displaystyle d^3𝐫\eta (𝐫)_zu(𝐫)},`$ (34)
$`\overline{\eta (𝐫)\eta (𝐫^{})}`$ $`=`$ $`D\delta (𝐫𝐫^{}),`$ (35)
a straightforward calculation shows that the displacement correlation function $`\overline{(u(𝐫)u(0))^2}`$ is given by Eq. (29). Smectics-A with disorder have been considered previously in Ref. radzihovsky99\_smectique, albeit with a different type of disorder coupling to $`_{}u`$. This yields a displacement correlation function superficially similar to $`F(r)`$ with $`q_{}^2`$ replacing $`q_x^2`$ in the numerator. The random compression force, which is not natural in the smectic-A context, is thus easily realized with charge density wave systems.
## V Experimental implications
In the preceding sections, we have shown that the forward scattering terms generated by charged impurities lead to smectic-like order in a charge density wave material. A frequently used technique to characterize positional correlations in CDW systems is x-ray diffraction.cowley\_x-ray\_cdw In the present section, we provide a calculation of the x-ray intensity resulting from such a smectic-like order, and we provide a quantitative estimate of the exponent $`\kappa `$.
### V.1 x-ray intensity
The intensity of the x-ray spectrum is given by guiner\_xray
$`I(q)={\displaystyle \frac{1}{L^3}}{\displaystyle \underset{i,j}{}}e^{iq(R_iR_j)}\overline{f_if_je^{iq(u_iu_j)}}.`$ (36)
$`u_i`$ is the atom displacement from the equilibrium position $`R_i`$, $`f_i`$ represents the total amplitude scattered by the atom at the position $`i`$ and depends exclusively on the atom type. We consider the simple case of a disordered crystal, made of one kind of atoms, characterized by the scattering factor $`\overline{f}\mathrm{\Delta }f/2`$, and containing impurities of scattering factor $`\overline{f}+\mathrm{\Delta }f/2`$. Since we are interested in the behavior of the scattering intensity near a Bragg peak ($`qK`$), we can use the continuum approximation.rosso\_cdw\_long In the case of the CDW, the lattice modulation is given by:
$$u(𝐫)=\frac{u_0}{Q}_x\left[\mathrm{cos}(Qx+\varphi (𝐫))\right],$$
(37)
It is well known that the presence of a CDW in the compound is associated with the appearance of two asymmetric satellites at positions $`qK\pm Q`$ around each Bragg peak.cowley\_x-ray\_cdw The intensity profiles of these satellites give access to the structural properties of the CDW. For this reason a lot of work has been done to compute and measure these intensities. ravy\_x-ray\_whiteline ; brazovskii\_x-ray\_cdwT ; rouziere\_friedel\_cdw ; rosso\_cdw\_short ; rosso\_cdw\_long By expanding Eq. (36) for low $`q(u_iu_j)`$, one finds an expression of the x-ray satellite intensity comprising a part $`I_\text{d}`$, which is symmetric under inversion around the Bragg vector $`K`$ and a part $`I_\text{a}`$ which is antisymmetric under the same transformation.rosso\_cdw\_long The symmetric part is given by the following correlation function:
$$I_\text{d}(𝐪)=\overline{f}^2q^2d^3𝐫e^{i\delta 𝐪𝐫}\overline{u(𝐫/2)u(𝐫/2)},$$
(38)
and the antisymmetric part by:
$$\frac{I_\text{a}(𝐪)}{𝐚(𝐛\times 𝐜)}=2q\mathrm{\Delta }f\mathrm{Im}d^3𝐫e^{i\delta 𝐪𝐫}\overline{\rho _{\text{imp}}(𝐫/2)u(𝐫/2)},$$
(39)
where: $`\delta q=(qK)Q`$, and $`𝐚(𝐛\times 𝐜)`$ is the volume of the unit cell of the crystal. After some manipulations, Eq. (38) can be rewritten as:
$`I_\text{d}(K+Q+k)`$ $`=`$ $`u_0^2\overline{f}^2K^2{\displaystyle d^3𝐫e^{i𝐤𝐫}C_\text{d}(𝐫)},`$ (40)
$`I_\text{a}(K+Q+k)`$ $`=`$ $`\overline{f}Ku_0\mathrm{\Delta }f\sqrt{𝒩D}{\displaystyle d^3𝐫e^{i𝐤𝐫}C_\text{a}(𝐫)}`$
where $`𝒩`$ is the number of impurities in the unit cell, and:
$`C_\text{d}(𝐫)`$ $`=`$ $`\overline{e^{i(\varphi (𝐫/2)\varphi (𝐫/2))}},`$ (41)
$`=`$ $`C_\text{d}^{\text{F.S}}(𝐫)C_\text{d}^{\text{B.S}}(𝐫),`$ (42)
$`C_\text{a}(𝐫)`$ $`=`$ $`\chi (𝐫)C_\text{d}(𝐫),`$ (43)
where $`\chi (𝐫)`$ is defined by Eq. (33) of Ref. rosso\_cdw\_long, . It is easy to show, using this definition and the statistical tilt symmetry that $`\chi (r)`$ is independent of the forward scattering disorder. In Eq. (41), $`C_\text{d}^{\text{B.S}}`$ is the backward scattering contribution which has been obtained in Ref. rosso\_cdw\_long, , and $`C_\text{d}^{\text{F.S}}`$ is the forward scattering contribution, given by:
$`C_\text{d}^{\text{F.S}}(𝐫)=\left({\displaystyle \frac{\mathrm{\Lambda }_{}^2}{r_{}^2+4(v_{}|x|\xi /v_x)}}\right)^\kappa ,`$ (44)
where we have used Eq. (30), assuming a Gaussian disorder. Using Eq. (32), one sees that the term $`C_\text{d}^{\text{B.S}}`$ gives only a logarithmic correction to (41). As a result, the symmetric structure factor $`I_\text{d}`$ is dominated by the contribution of the forward scattering disorder. To obtain the structure factor, we Fourier transform Eq. (44) to obtain
$`I_\text{d}(𝐪)`$ $`=`$ $`{\displaystyle d^2𝐫_{}e^{i𝐪_{}𝐫_{}}\left(\frac{\mathrm{\Lambda }_{}^1}{r_{}}\right)^{2\kappa }\frac{r_{}^2v_x}{2\xi v_{}}}`$
$`\times {\displaystyle _0^+\mathrm{}}{\displaystyle \frac{du}{(1+u)^\gamma }}\mathrm{cos}\left({\displaystyle \frac{|q_x|r_{}^2v_x}{4\xi v_{}}}u\right)`$
Using the following relation,
$`{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{du}{(1+u)^\gamma }}\mathrm{cos}(\lambda u)={\displaystyle \frac{\lambda ^{\gamma 1}}{\mathrm{\Gamma }(\gamma )}}{\displaystyle _0^+\mathrm{}}𝑑ve^{v\lambda }{\displaystyle \frac{v^\gamma }{v^2+1}}`$ (45)
we finally obtain
$`I_\text{d}(𝐪)`$ $`=`$ $`{\displaystyle \frac{\pi (|q_x|\mathrm{\Lambda }_{}^1)^{\kappa 2}}{2^{2(\kappa 1)}\mathrm{\Gamma }(\kappa )}}{\displaystyle \frac{\mathrm{\Lambda }_{}^{\kappa 2}}{(\xi v_{}/v_x)^{\kappa 1}}}`$ (46)
$`\times {\displaystyle _0^+\mathrm{}}dw{\displaystyle \frac{w^{1\kappa }}{w^2+1}}e^{w\frac{(\xi v_{}/v_x)q_{}^2}{2|q_x|}},`$
so that $`I_\text{d}(𝐪)(|q_x|)^{\kappa 2}`$ for $`q_{}^2(\xi v_{}/v_x)|q_x|`$ and $`I_\text{d}(𝐪)(|q_{}|)^{2(\kappa 2)}`$ otherwise. The intensity $`I_d(𝐪=0)`$ is divergent for for $`\kappa <2`$ but is finite for $`\kappa >2`$, i.e. for strong disorder. Next, we turn to the evaluation of $`I_\text{a}`$. From Ref. rosso\_cdw\_long, , we know that $`\chi (r)1/x`$ when $`x\xi r_{}^2`$ and $`\chi (r)1/r_{}^2`$ when $`|x|\xi r_{}^2`$. This implies that $`I_a`$ is subdominant in comparison with $`I_d`$. In particular, $`I_\text{a}(𝐪)(|q_x|)^{\kappa 1}`$ for $`q_{}^2(\xi v_{}/v_x)|q_x|`$ and $`I_\text{a}(𝐪)(|q_{}|)^{2(\kappa 1)}`$ otherwise. We illustrate the behavior of the x-ray intensities on Fig. 1.
We note that these intensities are remarkably similar to those of a disorder-free smectic-A liquid crystalcaille\_smectic\_xray at positive temperature. In fact, the expression of the exponent $`\kappa `$ Eq. (31) is analogous to the expression (5.3.12) in Ref. chandrasekhar\_smectic, , with the disorder strength $`D`$ playing the role of the temperature $`k_BT`$ in the smectic-A liquid crystal.
### V.2 Estimate of the exponent $`\kappa `$
Let us turn to an estimate of the exponent $`\kappa `$ appearing in the intensities to determine whether such smectic-like intensities are indeed observable in experiments. To do this, we first need to determine whether Coulomb interactions are unscreened by comparing the screening length with $`\xi `$ given by Eq. (21). This question is relevant only to a material with a full gap, in which free uncondensed electrons cannot screen charged impurities. A good candidate is the blue bronze material $`\mathrm{K}_{0.3}\mathrm{MoO}_3`$ which has a full gap, and is well characterized experimentally. We now evaluate the quantity $`\xi `$ for this material. Using the parameters of Ref. pouget\_bronzes, :
$`n_c=10^{20}\text{chains}/m^2,`$ (47)
$`v_F=1.3\times 10^5m.s^1,`$ (48)
$`\rho _0=3\times 10^{27}e^{}/m^3,`$ (49)
$`Q=6\times 10^9m^1,`$ (50)
and a relative permittivity of $`ϵ_{\mathrm{K}_{0.3}\mathrm{MoO}_3}=1`$, so that $`ϵ`$ in Eq. (21) is equal to the permittivity of the vacuum, we obtain: $`\xi 5`$Å. Therefore, the screening length can be large compared to $`\xi `$ at low temperature, and we expect that Coulombian effects will play an important role in this material. We can use this value of $`\xi `$ to evaluate the exponent $`\kappa `$ in Eq. (31). For the doped material $`\mathrm{K}_{0.3}\mathrm{Mo}_{1\mathrm{x}}\mathrm{V}_\mathrm{x}\mathrm{O}_3`$, we find that the disorder strength can be expressed as a function of the doping and obtain:
$$D=x(1x)\frac{\mathrm{\#}(\text{Mo atoms/unit cell})}{𝐚(𝐛\times 𝐜)}.$$
(51)
This formula is derived in App. E. For the crystal parameters, $`a=`$18.25Å, $`b=`$7.56Å, $`c=`$9.86Å, $`\beta =`$117.53<sup>o</sup> sato\_xray\_kmoo3 , with 20 Molybdenum atoms per unit cell, and a doping $`x=3\%`$, the disorder strength $`D=4.8\times 10^{26}\mathrm{m}^3`$. Moreover, using the experimental bounds of the velocities: $`3.6\times 10^2m.s^1<v_{}<1.6\times 10^3m.s^1`$ and $`v_x=3.7\times 10^3m.s^1`$, we find that $`\kappa `$ is in the range $`[0.160.8]`$. Therefore, the smectic-like order should be observable in x-ray diffraction measurements on this material.
## VI Conclusion
In this paper, we have introduced a decomposition of the disorder induced by charged impurities in terms of the reciprocal lattice vectors of a periodic charged elastic system. Using this decomposition, we have shown that only the long wavelength (forward scattering) component of the disorder was long-range correlated. Components with wavevectors commensurate with the reciprocal lattice of the periodic elastic system remain short ranged. The latter can thus be treated with the standard techniques developed for impurities producing short range forces.giamarchi\_vortex\_long We find that only the forward scattering is affected by the long-range character of the forces created by charged impurities. Due to the statistical tilt symmetry, this implies that only the statics of the periodic elastic system is modified by Coulombian disorder. This has allowed us to obtain a full picture of the statics of charge density wave systems in $`d=3`$ in the presence of charged and neutral impurities. The results are summarized in table 1.
A remarkable result is that in the case of charged impurities in a system with unscreened Coulomb elasticity, the x-ray intensity turns out to be identical to that produced by thermal fluctuations in a smectic-A liquid crystalcaille\_smectic\_xray , with the disorder strength playing the role of an effective temperature. This behavior of the scattering intensity should be observable in the blue bronze material K<sub>0.3</sub>MoO<sub>3</sub> doped with charged impurities such as Vanadium.
###### Acknowledgements.
This work was supported in part by the Swiss National Fund under MANEP and division II. A. R. thanks the University of Geneva for hospitality and support. We thank S. Brazovskii for enlightening discussions.
## Appendix A Decomposition of the Coulomb potential
In this appendix, we give more details on the decomposition of the Coulomb potential using infinitely differentiable functions of compact support.gelfand64\_partition First, let us discuss a simple decomposition in 1D. We consider the function $`F_ϵ(x)`$ such that:
$`F_ϵ(x)`$ $`=`$ $`F_ϵ(x)`$
$`F_ϵ(x)`$ $`=`$ $`1,\mathrm{\hspace{0.17em}0}x1`$
$`F_ϵ(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{1\mathrm{tanh}\left[{\displaystyle \frac{2(1x)}{(x1)^2+ϵ^2}}\right]\right\},|x1|<ϵ`$
$`F_ϵ(x)`$ $`=`$ $`0,x>1+ϵ`$ (52)
It is easy to check that $`F_ϵ`$ is continuous, infinitely differentiable, and that:
$`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}F_ϵ(x2n)=1.`$ (53)
Applying this formula to a one dimensional reciprocal space, we obtain:
$`{\displaystyle \underset{n_x=\mathrm{}}{\overset{\mathrm{}}{}}}F_ϵ\left({\displaystyle \frac{a}{\pi }}\left(q{\displaystyle \frac{2\pi n}{a}}\right)\right)=1,`$ (54)
i.e. we have constructed explicitly a partition of the unity.gelfand64\_partition The generalization to a cubic lattice in a three dimensional space is obvious:
$`{\displaystyle \underset{n_x,n_y,n_z=\mathrm{}}{\overset{\mathrm{}}{}}}F_ϵ\left({\displaystyle \frac{a}{\pi }}\left(q_x{\displaystyle \frac{2\pi n_x}{a}}\right)\right)F_ϵ\left({\displaystyle \frac{a}{\pi }}\left(q_y{\displaystyle \frac{2\pi n_y}{a}}\right)\right)`$
$`\times F_ϵ\left({\displaystyle \frac{a}{\pi }}(q_z{\displaystyle \frac{2\pi n_z}{a}})\right)={\displaystyle \underset{𝐆}{}}F_ϵ^{3D}(𝐪𝐆)=1.`$
The function $`F_ϵ^{3D}`$ has a compact support, and vanishes rapidly outside the Brillouin zone.
## Appendix B Contribution of the oscillating components of the density to the elastic Hamiltonian in the presence of Coulomb interaction
In this appendix, we calculate the contribution of the oscillating terms to the Hamiltonian of the charge density wave with unscreened Coulomb interactions, Eq. (17), and show that they only induce corrections to the short range elasticity. Inserting the expression of the density, Eq. (15), in Eq. (7), the contribution of the oscillating component of wavevector $`Q`$ is given by,
$`H_C^{\text{osc.}}`$ $`=`$ $`{\displaystyle \frac{e^2\rho _1^2}{2ϵ}}{\displaystyle \underset{𝐧,𝐧^{}}{}}{\displaystyle 𝑑x𝑑x^{}V_Q(xx^{},𝐧𝐧^{})}`$ (56)
$`\times \mathrm{cos}[\varphi (x,𝐧)\varphi (x^{},𝐧^{})],`$
where we have reestablished the discrete character of the transverse dimension $`𝐲`$. Both the intrachain ($`𝐧=𝐧^{}`$) and the interchain contributions ($`𝐧𝐧^{}`$) are short ranged. Let us first consider the case of $`𝐧𝐧^{}`$. We have to compute integrals:
$`V_Q(x,𝐧)={\displaystyle \frac{d^dq}{(2\pi )^d}F_ϵ^{3D}(𝐪)\frac{e^{i(q_xx+𝐪_{}𝐧)}}{(Q+q_x)^2+q_{}^2}},`$ (57)
To evaluate the above integral in closed form, we need to make some approximations. Since $`F_ϵ`$ vanishes for $`Q+q_x=0`$, we can neglect $`q_x`$ compared to $`Q`$ in this integral. Then, we can extend the integration over the whole reciprocal space without encountering any singularity. The $`q_x`$ integration produces a $`\delta (x)`$ function, and the $`q_{}`$ integration gives an exponential in 2D and a modified Bessel function $`K_0`$ in 3D. In the two-dimensional case, we find that:
$`V_Q(x,𝐧)={\displaystyle \frac{e^{Q\mathrm{}_yn}}{2Q}}\delta (x),`$ (58)
and in the three-dimensional case:
$`V_Q(x,𝐧)={\displaystyle \frac{K_0(Q\sqrt{(n_y\mathrm{}_y)^2+(n_z\mathrm{}_z)^2})}{2\pi }}\delta (x).`$ (59)
where, $`\mathrm{}_y`$ and $`\mathrm{}_z`$ are interchain spacings. Due to the exponential decay of the interchain interaction with the distance, it is justified to neglect interchain interactions beyond nearest neighbors. The logarithmic divergence in Eq. (59) for $`𝐧=0`$ is an artefact of the approximation we make when we integrate over the entire reciprocal space instead of the first Brillouin zone. A more refined estimate yields a finite, short ranged intrachain contribution. The short range contribution in the electrostatic energy thus reads:
$`H_C^{\text{osc.}}`$ $`=`$ $`{\displaystyle \underset{\alpha =y,z}{}}J_\alpha {\displaystyle 𝑑x\underset{𝐧}{}\mathrm{cos}[\varphi (x,𝐧+𝐞_\alpha )\varphi (x,𝐧)]}+`$
$`+{\displaystyle \underset{𝐧}{}}{\displaystyle V_Q(xx^{},0)𝑑x𝑑x^{}\mathrm{cos}[\varphi (x,𝐧)\varphi (x,𝐧)]},`$
where:
$$J_\alpha =\frac{e^2\rho _1^2}{2\pi ϵ}K_0(Q\mathrm{}_\alpha ).$$
(61)
In Eq. (B), we expand $`\mathrm{cos}(\varphi (x,𝐧)\varphi (x^{},𝐧))1(xx^{})^2/2(_x\varphi (x,𝐧))^2`$ to show that the backscattering term reduces to short ranged elastic forces. Making $`\varphi (x,𝐧)=\overline{\varphi }(x,𝐧)+n_y\pi +n_z\pi `$, we can make the sign in front of $`J_\alpha `$ negative, and obtain the ground state for $`\overline{\varphi }=0`$. In its ground state, the CDW is out of phase on two nearest neighbor chains. This ground state is represented on Fig. 3. The excitations above this ground state are described by the Lagrangian (D) derived in App. D.
## Appendix C Zero temperature limit of the charge density in a CDW
In the present appendix, we discuss the zero temperature limit of the expression of the charge density in a CDW. Let us consider the form of the charge density in the presence of a non-uniform $`\varphi `$. It is givenfukuyama\_cdw\_pinning ; lee\_coulomb\_cdw ; lee\_depinning ; rice\_lee\_cross by the expression:
$`\rho (𝐫)=\rho _0+{\displaystyle \frac{\rho _0\overline{\rho }_c}{Q}}_x\varphi (𝐫)+\rho _1|\psi |\mathrm{cos}(Qx+\varphi (𝐫))`$ (62)
where $`\rho _0`$ is the average electron density, $`\rho _1`$ is the condensate amplitude at $`T=0`$, $`|\psi |`$ takes into account the reduction of CDW order by thermal fluctuations ($`|\psi |=1`$ at $`T=0`$), the factor $`\overline{\rho }_c`$ takes into account the presence of non-condensed electronsrice\_lee\_cross at finite temperature (at $`T=0`$, $`\overline{\rho }_c=1`$) and $`Q=2k_F`$. Using the relation $`k_F=\frac{\pi }{2}\rho _0`$, valid in a one-dimensional system, one can see that this relation simplifies (at $`T=0`$) to:
$`\delta \rho ={\displaystyle \frac{_x\varphi }{\pi }}.`$ (63)
In the paper, we consider temperatures very low compared to the Peierls transition temperature, and we take $`|\psi |=1`$, $`\overline{\rho }_c=1`$. This yields Eq. (15).
## Appendix D Derivation of the Hamiltonian of a three-dimensional charge density wave
In this appendix, we provide a derivation of the Hamiltonian of a three dimensional CDW starting from the original Fukuyama one dimensional description. The Lagrangian of a CDW in a single chain is given by:fukuyama\_cdw\_pinning
$`={\displaystyle \frac{\mathrm{}v_F}{4\pi }}{\displaystyle 𝑑x\left[\frac{(_t\varphi )^2}{v_\varphi ^2}(_x\varphi )^2\right]}`$ (64)
Obviously, this Lagrangian describes phason waves propagating with the velocity $`v_\varphi `$. $`v_F`$ is the Fermi velocity of the electrons forming the CDW and $`\varphi `$ is the phase of the CDW. We define an effective mass $`m^{}`$ by:
$`{\displaystyle \frac{v_F^2}{v_\varphi ^2}}={\displaystyle \frac{m^{}}{m_e}}`$ (65)
where $`m_e`$ is the electron mass.
In a three-dimensional CDW with screened Coulomb interactions, the chains are coupled by a backscattering interaction. The resulting Lagrangian reads:
$``$ $`=`$ $`{\displaystyle \frac{\mathrm{}v_F}{4\pi }}{\displaystyle \underset{𝐧}{}}{\displaystyle 𝑑x\left[\frac{(_t\varphi )^2}{v_\varphi ^2}(_x\varphi )^2\right](x,𝐧)}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐧,𝐧^{}}{}}J(𝐧,𝐧^{}){\displaystyle 𝑑x\mathrm{cos}[\varphi (x,𝐧)\varphi (x,𝐧^{})]},`$
where $`J(𝐧,𝐧^{})`$ is short-ranged and is given by Eq. (61). Expanding the cosines,
$`\mathrm{cos}[\varphi (x,𝐧)\varphi (x,𝐧+𝐞_𝐲)]=`$
$`=1{\displaystyle \frac{(\varphi (x,𝐧)\varphi (x,𝐧+𝐞_𝐲))}{2}}+o(_y^2)`$
$`=1{\displaystyle \frac{\mathrm{}_y^2}{2}}(_y\varphi )^2+o(_y\varphi )^2,`$ (67)
and defining,
$`{\displaystyle \frac{\mathrm{}v_F}{4\pi }}{\displaystyle \frac{v_{y,z}^2}{v_\varphi ^2}}`$ $`=`$ $`J_{y,z}{\displaystyle \frac{\mathrm{}_{y,z}^2}{2}},`$ (68)
the Lagrangian in Eq. (D) can be rewritten aslee\_coulomb\_cdw :
$``$ $`=`$ $`{\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\mathrm{}v_F}{4\pi }}{\displaystyle }dx[{\displaystyle \frac{(_t\varphi )^2}{v_\varphi ^2}}(_x\varphi )^2`$ (69)
$`{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2]`$
The sum over lattice sites in the transverse direction can be replaced by an integral, by writing:
$``$ $`=`$ $`{\displaystyle \frac{\mathrm{}v_F}{4\pi \mathrm{}_y\mathrm{}_z}}{\displaystyle }d^3𝐫[{\displaystyle \frac{(_t\varphi )^2}{v_\varphi ^2}}(_x\varphi )^2`$ (70)
$`{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2]`$
The phason dispersion is now $`\omega (q)^2=v_\varphi ^2q_x^2+v_y^2q_y^2+v_zq_z^2`$. The momentum conjugate to $`\varphi `$ is obtained by the usual relation:
$$\mathrm{\Pi }=\frac{\delta }{\delta \left(_t\varphi \right)}=\frac{\mathrm{}v_F}{2\pi v_\varphi ^2\mathrm{}_y\mathrm{}_z}_t\varphi ,$$
(71)
yielding the Hamiltonian:
$`H`$ $`=`$ $`{\displaystyle \frac{\mathrm{}v_F}{4\pi \mathrm{}_y\mathrm{}_z}}{\displaystyle }d^3𝐫[{\displaystyle \frac{4\pi ^2v_\varphi ^2(\mathrm{}_y\mathrm{}_z)^2}{\mathrm{}^2v_F^2}}\mathrm{\Pi }^2+(_x\varphi )^2`$ (72)
$`+{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2+{\displaystyle \frac{v_y^2}{v_\varphi ^2}}(_y\varphi )^2].`$
From the Hamiltonian Eq. (72), it is straightforward to obtain the Debye-Waller factor associated with the zero point fluctuations of the phase $`\varphi `$. In the isotropic case, $`v_y=v_z=v_\varphi `$, one finds that $`\mathrm{cos}\varphi (x)_{T=0}e^{C(m_e/m^{})^{1/2}}`$, where $`C\pi /4`$ is a dimensionless constant of order $`1`$. Due to the smallness of the ratio $`m_e/m^{}10^2`$, the zero point motion can be neglected, and the kinetic term $`\mathrm{\Pi }^2`$ in Eq. (72) can be dropped. This leads to the Hamiltonian (III).
## Appendix E Estimation of the disorder strength
Here, we give an estimation of the disorder strength $`D`$ in doped $`\mathrm{KMo}_{1\mathrm{x}}\mathrm{V}_\mathrm{x}\mathrm{O}_3`$. We assume a binomial distribution of Vanadium impurities on the Molybdenum sites. The Vanadium impurities carry an extra electron compared to the Molybdenum ions. The resulting charge density fluctuation reads:
$$\delta \rho (𝐫)=\underset{i,\alpha }{}(x\sigma _{i,\alpha })\delta (𝐫𝐑_{i,\alpha }),$$
(73)
where $`i`$ is the index of the cell and $`\alpha `$ is the index of the Molybdenum site in a given cell. $`\sigma _{i,\alpha }=0`$ if the site is occupied by a Molybdenum ion, and $`\sigma _{i,\alpha }=1`$ if it is occupied by a Vanadium impurity. By construction, the expectation value of $`\delta \rho (𝐫)`$ is zero. We estimate the second moment of $`\delta \rho (𝐫)`$ as:
$`\overline{\delta \rho (𝐫)\delta \rho (𝐫^{})}=`$
$`={\displaystyle \underset{i,j,\alpha ,\beta }{}}\overline{(x\sigma _{i,\alpha })(x\sigma _{j,\beta })}\delta (𝐫𝐑_{i,\alpha })\delta (𝐫^{}𝐑_{j,\beta })`$
$`=x(1x){\displaystyle \underset{i,\alpha }{}}\delta (𝐫𝐑_{i,\alpha })\delta (𝐫^{}𝐫),`$ (74)
where we have used the property that $`\overline{(x\sigma _{i,\alpha })(x\sigma _{j,\beta })}=\delta _{i,j}\delta _{\alpha ,\beta }\overline{(x\sigma _{j,\beta })^2}`$. The expectation value of $`_{i,\alpha }\delta (𝐫𝐑_{i,\alpha })`$ is simply the number of Molybdenum ions per unit volume, leading to the formula Eq. (51).
|
warning/0506/hep-ph0506060.html
|
ar5iv
|
text
|
# The 𝑔 factor of an electron or muon bound by an arbitrary central potential
## I Introduction
Study of energy levels of atomic electron (muon) has a long history. Various relativistic and radiative corrections have been obtained in analytic or semi-analytic form taking into account a deviation of the potential from the Coulomb one. Recently interest to an accurate theory of the $`g`$ factor increased. Accurate theoretical approach needs to take into account relativistic and radiative corrections. Here we consider a case when such corrections correspond to some modification of a central potential. Such a problem is actual because of the Uehling potential, which is responsible for a dominant QED correction in muonic atoms, finite-nuclear-size effects and some others. The results are presented in a fully relativistic approach, i.e. exact in $`Z\alpha `$ for the Coulomb interaction and its modifications.
The methods of calculation of energy levels have been successfully developed over the decades. The energy levels can be calculated with high accuracy by solving certain equations analytically or numerically, or using the perturbation theory with respect to the corresponding correction to the hamiltonian. Calculation of the $`g`$ factor is a more complicated problem. In the traditional approach to the calculation of the correction to the $`g`$ factor due to modification of the potential $`\delta V`$, one considers both the magnetic field and $`\delta V`$ as a perturbation, thus starting with the second order of the perturbation theory. The numerical approaches are also not easy to apply because they usually are much more accurate in determination of the energy than the wave function. However, it is the latter which is needed to calculate the $`g`$ factor.
In this paper, we develop a framework which allows one to express the $`g`$ factor of a Dirac particle bound in arbitrary central potential via the binding energy. For a small deviation of this potential from the Coulomb one, we derive the general expression for the first correction to the $`g`$ factor due to this deviation. The results are valid both for muonic and electronic atoms.
## II General relations
In the weak homogenous magnetic field $`𝑩`$, the correction to energy levels of a Dirac particle bound by a central potential reads
$`\mathrm{\Delta }E`$ $`=`$ $`e{\displaystyle d^3r\overline{\mathrm{\Psi }}(𝒓)\left(𝜸𝐀\right)\mathrm{\Psi }(𝒓)}`$ (1)
$`=`$ $`{\displaystyle \frac{e𝑩𝑱}{2m}}g,`$
where
$$𝐀=\frac{1}{2}\left[𝑩\times 𝐫\right].$$
(2)
The relativistic units in which $`\mathrm{}=c=1`$ are applied throughout the paper; the charge of the bound particle is $`e`$, $`e`$ is a charge of a proton, $`e^2=\alpha =1/137`$ is the fine structure constant.
Writing the Dirac wave function $`\psi (𝒓)`$ in the form
$$\psi (𝒓)=\left(\begin{array}{c}f_1(r)\mathrm{\Omega }\\ if_2(r)\stackrel{~}{\mathrm{\Omega }}\end{array}\right),$$
(3)
where $`\mathrm{\Omega }`$ is the spherical spinor BLP with the angular momentum $`J`$ and orbital momentum $`L`$, $`\stackrel{~}{\mathrm{\Omega }}=(𝝈𝒏)\mathrm{\Omega }`$, we obtain for the $`g`$ factor of a bound Dirac particle
$$g=\frac{2m\kappa }{j(j+1)}𝑑rr^3f_1(r)f_2(r),$$
(4)
$`\kappa =j+1/2`$ for $`l>j`$ and $`\kappa =(j+1/2)`$ for $`l<j`$. The integral in the right-hand side of Eq. (4) can be presented as follows (see, e.g., rose )
$$𝑑rr^3f_1f_2=\frac{1}{4m}\left[12\kappa 𝑑rr^2\left(f_1^2f_2^2\right)\right].$$
(5)
This formula is valid for the Dirac equation with any central potential. Eq. (4) has been known for a while (see, e.g., rose ), but, to the best of our knowledge, it was never applied to a non-Coulomb problems.
Using the identity
$$𝑑rr^2\left(f_1^2f_2^2\right)=\gamma _0,$$
(6)
we have
$$g=\frac{\kappa [12\kappa \gamma _0]}{2j(j+1)}.$$
(7)
Assuming that the potential $`V(r)`$ is independent of the mass of the bound particle, we find
$$\gamma _0=\frac{H}{m}=\frac{E}{m},$$
(8)
where $`H`$ is the Dirac hamiltonian $`H=\gamma _0\left(𝜸𝒑\right)+\gamma _0m+V(r)`$.
Thus we arrive at an equation, which expresses the $`g`$ factor of the state via its binding energy
$$g=\frac{\kappa }{2j(j+1)}\left[12\kappa \frac{E}{m}\right].$$
(9)
This equation can be used even in the case when a potential substantially differs from the Coulomb one, as for heavy muonic atoms. In many cases, Eq. (9) essentially simplifies calculation of the corrections to the $`g`$ factor of a Dirac particle. In particular, it can be applied to such problems as finite-nuclear-size effect and vacuum polarization.
The equation (9) becomes essentially simpler if the potential is close to the Coulomb one, and the deviation $`\delta V(r)=V(r)V_C(r)`$ can be treated as a perturbation. For the pure Coulomb case, when $`E_\mathrm{C}/m=E_\mathrm{C}/m`$, we immediately obtain the well-known result (see, e.g., OS ; shabaev )
$$\mathrm{\Delta }g_\mathrm{C}=\frac{\kappa }{2j(j+1)}\left[12\kappa \frac{\gamma +n_r}{N}\right].$$
(10)
where
$$N=\sqrt{(\gamma +n_r)^2+(Z\alpha )^2},\gamma =\sqrt{\kappa ^2(Z\alpha )^2},$$
and $`n_r`$ is the radial quantum number. The correction to the energy is
$$\delta E=d^3r\mathrm{\Psi }^+(r)\delta V(r)\mathrm{\Psi }(r).$$
(11)
As follows from the dimensional reasons, the wave function in the Coulomb field can be presented in the form $`\mathrm{\Psi }(r)=m^{3/2}\stackrel{~}{\mathrm{\Psi }}(mr)`$, where $`\stackrel{~}{\mathrm{\Psi }}`$ is dimensionless. Passing to the variable $`\rho =mr`$, we find
$$\delta E=d^3\rho \stackrel{~}{\mathrm{\Psi }}^+(\rho )\delta V(\rho /m)\stackrel{~}{\mathrm{\Psi }}(\rho ).$$
(12)
Taking the derivative over $`m`$ and returning to the variable $`r`$, we obtain
$$\frac{\delta E}{m}=\frac{1}{m}d^3r\mathrm{\Psi }^+(r)r\frac{\delta V(r)}{r}\mathrm{\Psi }(r),$$
(13)
and
$$\delta g=\frac{\kappa }{2j(j+1)}\left[1+\frac{2\kappa }{m}r\frac{\delta V(r)}{r}\right].$$
(14)
Eqs. (9), (14) are the basis of our approach.
## III Anomalous magnetic moment
It may be interesting to generalize our results to the case of a particle with a non-vanishing anomalous magnetic moment. In particular, it is necessary for antiprotonic atoms. The anomalous magnetic moment of the antiproton is big because of the complicated internal structure of the particle. Therefore, it makes sense to consider the modification of Eqs. (9), (14) due to the anomalous magnetic moment separately from the radiative corrections.
As known, the modified Dirac hamiltonian for a particle with the anomalous magnetic moment in an arbitrary electromagnetic field is of the form (see, e.g., BLP )
$$H=\gamma _0\left(𝜸(𝒑+e𝑨)\right)+\gamma _0meA_0+\frac{i\mu _a}{2}\gamma _0\sigma _{\mu \nu }F^{\mu \nu },$$
(15)
where $`\sigma _{\mu \nu }=(1/2)[\gamma _\mu ,\gamma _\nu ]`$, $`\mu _a`$ is the dimensional anomalous magnetic moment, $`F^{\mu \nu }`$ is the tensor of the electromagnetic field; we remind that the charge of the bound particle is $`(e)`$.
As follows from Eq. (15), the correction to the energy of a bound particle in the homogenous magnetic field (2) is of the form (cf. Eq. (1))
$$\mathrm{\Delta }E=d^3r\overline{\mathrm{\Psi }}(𝒓)\left\{e\left(𝜸𝐀\right)\mu _a\left(𝚺𝐁\right)\right\}\mathrm{\Psi }(𝒓),$$
(16)
where $`\mathrm{\Sigma }_j=(i/2)ϵ_{jkl}\sigma _{kl}`$ is the spin operator. The wave function $`\mathrm{\Psi }(𝒓)`$ is an eigenfunction of the hamiltonian (15) with $`eA_0(r)=V(r)`$ and $`A_i=0`$. Since the potential $`A_0(r)`$ is spherically symmetric, the wave function still has the form (3), and the result for the $`g`$ factor reads (cf. Eq. (4))
$`g`$ $`=`$ $`{\displaystyle \frac{2m}{j(j+1)}}\times \{\kappa {\displaystyle }drr^3f_1f_2`$ (17)
$``$ $`{\displaystyle \frac{\mu _a}{2e}}[{\displaystyle }drr^2(f_1^2f_2^2)2\kappa ]\}.`$
The radial wave functions $`f_1`$ and $`f_2`$ satisfy the system of equations (cf. BLP )
$`f_2^{}+(EVm)f_1+\left({\displaystyle \frac{1\kappa }{r}}{\displaystyle \frac{\mu _a}{e}}V^{}\right)f_2=0,`$
$`f_1^{}(EV+m)f_2+\left({\displaystyle \frac{1+\kappa }{r}}+{\displaystyle \frac{\mu _a}{e}}V^{}\right)f_1=0.`$ (18)
Then we multiply the first equation by $`r^3f_2`$ and the second one by $`r^3f_1`$, sum up the results and take the integral over $`r`$. Integrating by parts the terms with the derivatives, we arrive at the following relation (cf. (5))
$`{\displaystyle 𝑑rr^3f_1f_2}`$ $`=`$ $`{\displaystyle \frac{1}{4m}}[12\kappa {\displaystyle }drr^2(f_1^2f_2^2)`$ (19)
$``$ $`{\displaystyle \frac{2\mu _a}{e}}{\displaystyle }drr^3V^{}(f_1^2f_2^2)].`$
Using this formula and Eqs. (6) and (8), we present the result for the $`g`$ factor in the form
$`g`$ $`=`$ $`{\displaystyle \frac{1}{2j(j+1)}}\{\kappa (1+2a)`$ (20)
$`+`$ $`(2\kappa ^2+a){\displaystyle \frac{E}{m}}`$
$``$ $`a{\displaystyle \frac{\kappa }{m}}{\displaystyle }drr^3V^{}(f_1^2f_2^2)\},`$
where $`a=2m\mu _a/e`$ is the dimensionless anomalous magnetic moment (for a free particle $`g=2(1+a)`$). This result is obtained for an arbitrary potential $`A_0`$ and is exact in the parameter $`a`$. We emphasize that the radial wave functions $`f_1`$ and $`f_2`$ and the binding energy $`E`$ depend on the anomalous magnetic moment because of the equations (III).
In the nonrelativistic approximation, Eq. (20) turns into
$`g`$ $`=`$ $`{\displaystyle \frac{1}{2j(j+1)}}\{(12\kappa )(a\kappa )`$ (21)
$``$ $`(2\kappa ^2+2\kappa a+a){\displaystyle \frac{1}{2m}}{\displaystyle }drr^3V^{}f_{nr}^2\},`$
where $`f_{nr}`$ is the radial part of the nonrelativistic wave function. This formula is valid even if $`a1`$.
It is interesting also to consider an expansion of the $`g`$ factor (20) in the parameter $`a`$ for arbitrary field strength (when the nonrelativistic approximation is not valid). The linear in this parameter term reads
$`\delta g_a`$ $`=`$ $`{\displaystyle \frac{a}{2j(j+1)}}\{2\kappa +{\displaystyle \frac{E}{m}}`$ (22)
$``$ $`{\displaystyle \frac{2\kappa ^2}{m}}{\displaystyle \frac{}{m}}{\displaystyle 𝑑rr^2V^{}f_1f_2}`$
$``$ $`{\displaystyle \frac{\kappa }{m}}{\displaystyle }drr^3V^{}(f_1^2f_2^2)\}.`$
Here we used the formula for the linear in $`\mu _a`$ correction to the energy,
$`\delta E_a={\displaystyle \frac{2\mu _a}{e}}{\displaystyle 𝑑rr^2V^{}f_1f_2},`$ (23)
and took into account that $`\mu _a/m=0`$. Let us consider Eq.(22) for the pure Coulomb potential. Strictly speaking, the equations (III) have no sense for a pure Coulomb field because of the terms $`1/r^2`$, which lead to the phenomenon of falling to the center. It is a consequence of the point-like source of the field and is absent if finite nuclear size is taken into account. However, for large quantum numbers which are mostly interesting for the experiments, the correction due to the finite nuclear size does not change essentially the result obtained from Eq.(22) for a pure Coulomb field. Using the radial matrix element from Ref.shabaev , we obtain
$`\delta g_a`$ $`=`$ $`{\displaystyle \frac{a}{2j(j+1)}}\{2\kappa +{\displaystyle \frac{\gamma +n_r}{N}}{\displaystyle \frac{\kappa (Z\alpha )^2}{N^2}}`$ (24)
$``$ $`{\displaystyle \frac{4\kappa ^2(Z\alpha )^4[2\kappa (\gamma +n_r)N]}{\gamma (4\gamma ^21)N^4}}\},`$
where $`N`$ and $`\gamma `$ are defined after Eq. (10). We see that the correction to the $`g`$ factor due to the anomalous magnetic moment has essentially more complicated dependence on quantum numbers than the leading term (10).
## IV Finite-nuclear-size effect
As an illustration of efficiency of our method, let us consider a finite-nuclear-size correction to the bound-electron g factor. For $`1s`$ state, this problem was solved analitically in Ref.Savelii2000 in the nonrelativistic approximation ($`Z\alpha 1`$), and numerically in Ref.Beier2000 for arbitrary value of the parameter $`Z\alpha `$. In Ref.Shabaev2002 , the result was obtained analytically in the next-to-leading approximation in the parameter $`Z\alpha `$, and for arbitrary states. The result of Ref.Shabaev2002 allows one to describe well the correction to the $`g`$ factor up to $`Z=20`$. On the other hand, in Ref.Shabaev1993 the approximate formulas for the finite-nuclear-size correction to the energy levels, $`\delta E_{fns}`$, were obtained, which gives the result with a relative error of about $`0.2\%`$ up to $`Z=100`$. Using Eq.(9) and known dependence of $`E_{fns}`$ on $`m`$, see Shabaev1993 , we obtain
$$\delta g_{fns}=\frac{\kappa ^2}{j(j+1)}\frac{E_{fns}}{m}=\frac{\kappa ^2(2\gamma +1)}{j(j+1)}\frac{E_{fns}}{m}.$$
(25)
Moreover, using our method we can calculate the radiative correction to $`\delta g_{fns}`$ for $`s_{1/2}`$ state and $`p_{1/2}`$ state (and arbitrary $`n_r`$) coming from the effect of vacuum polarization. For these states, the finite-nuclear-size effect is the most significant. The radiative correction $`\delta E_{fns}`$ to the energy $`E_{fns}`$ was considered in detail in Ref.MST2004 . The correction due to the vacuum polarization $`\delta E_{fns}^{VP}`$ can be represented in the form
$$\delta E_{fns}^{VP}=E_{fns}\mathrm{\Delta }^{VP},$$
(26)
where $`E_{fns}`$ was obtained in Ref. Shabaev1993 , and the explicit analytical form of $`\mathrm{\Delta }^{VP}`$ for $`s_{1/2}`$ state and $`p_{1/2}`$ state is derived in Ref.MST2004 . Eq. (26) with $`E_{fns}`$ from Ref.Shabaev1993 and $`\mathrm{\Delta }^{VP}`$ from Ref.MST2004 provides the high accuracy of $`\delta E_{fns}^{VP}`$ up to $`Z=100`$. However, the correction $`\mathrm{\Delta }^{VP}`$ was calculated in MST2004 for the mass of the bound particle $`m`$ being equal to the mass $`M`$ of the particle in the fermion loop. Therefore, in order to use Eq. (9) we should know the explicit dependence of $`\mathrm{\Delta }^{VP}`$ on $`m/M`$ (we can set $`M=m`$ only after the differentiation over $`m`$). The corresponding expression for $`\mathrm{\Delta }^{VP}`$ can be easily obtained following the derivation of $`\mathrm{\Delta }^{VP}`$ in MST2004 . For $`Z\alpha 1`$, the main contribution to $`\mathrm{\Delta }^{VP}`$ is given by the logarithmically enhanced term
$$\mathrm{\Delta }^{VP}\frac{2\alpha (Z\alpha )^2}{3\pi \gamma _1}\mathrm{ln}^2(mR),$$
(27)
where $`\gamma _1=\sqrt{1(Z\alpha )^2}`$, and $`R`$ is the nuclear radius. This term is the same for $`s_{1/2}`$ and $`p_{1/2}`$ states and comes from the distances $`1/m,\mathrm{\hspace{0.17em}1}/MrR`$. Therefore, it contains only the logarithmic dependence on the mass of the bound particle, in contrast to the power-like dependence of $`E_{fns}`$. Within the logarithmic accuracy, we have
$$\frac{}{m}\delta E_{fns}^{VP}\mathrm{\Delta }^{VP}\frac{}{m}E_{fns}=(2\gamma _1+1)\frac{\delta E_{fns}^{VP}}{m}.$$
(28)
Thus, for $`Z\alpha 1`$, we obtain
$$\delta g_{fns}^{VP}=\frac{4}{3}(2\gamma _1+1)\frac{\delta E_{fns}^{VP}}{m}.$$
(29)
For $`Z\alpha 1`$ the dependence of $`\mathrm{\Delta }^{VP}`$ on $`m/M`$ is more complicated and will be considered elsewhere.
## V Conclusion
We present a new effective approach to obtain various corrections to the bound particle $`g`$ factor using the corresponding corrections to the energy levels. It should be noted that this method is appropriate for the perturbations of potential-like type. There are also corrections which cannot be derived via this approach. The leading part of the non-potential contribution is due to the virtual Delbrück scattering and it was calculated in Ref. pra .
As other applications of our method, we can suggest the Uehling potential contribution to the $`g`$ factor of bound electrons and muons. In muonic atoms it is one of the dominant effects, while in electronic hydrogen-like atoms the accuracy is essentially higher (see, e.g., eions ).
For muonic atoms, the Uehling-potential correction to the energy levels was calculated in the nonrelativistic approximation for certain levels in pusto . The result exact in $`Z\alpha `$ was obtained in cjp98 for ground state of hydrogen-like atom and agrees with Eq. (9). The methods developed there can be easily applied to arbitrary levels and thus, using Eq. (9), one can find the correction to the $`g`$ factor due to the Uehling potential.
Using the known corrections to the energy levels coming from the high-order terms of vacuum polarization, we can obtain the corresponding contribution to the g factor. One correction comes from the Wichmann-Kroll potential (which accounts for the higher-order in $`Z\alpha `$ terms in the induced charge density), see Review BR1982 . Others are the second-order correction with respect to the Uehling potential, and first-order correction with respect to the radiative correction to the Uehling potential (of order of $`\alpha `$).
## Acknowledgements
We are grateful to V.G. Ivanov for useful discussions. A.I.M. thanks the Max-Planck-Institute for Nuclear physics, Heidelberg, and the Max-Planck-Institute for Quantum Optics, Garching, for hospitality during the visit. The work was supported in part by RFBR Grants No. 03-02-16510, 03-02-04029, and 03-02-16843, and DFG GZ: 436 Rus 113/769/0-1R.
|
warning/0506/gr-qc0506114.html
|
ar5iv
|
text
|
# A geometric description of the intermediate behaviour for spatially homogeneous models
## 1 Introduction
Due to their appealing geometric, kinematic and dynamical structure, Spatially Homogeneous (SH or Bianchi) models have received considerable attention in the last 3 decades. Apart from the obvious gain of a direct generalization of the standard Friedmann-Lemaître cosmological model, one of the main reasons for this interest is the fact that the Einstein’s Field Equations (EFE) are reduced to a coupled system of ordinary differential equations. Then by introducing the orthonormal frame formalism and expansion-normalized variables, in order to scale away the influence of the expansion from the overall evolution of the corresponding models, one exploits methods from the theory of dynamical systems to examine their behaviour at early, late and intermediate periods of their history. This approach has shown the significant role of the transitively self-similar SH models since they represent the past and future (equilibrium) states for the majority of evolving vacuum and $`\gamma `$law perfect fluid models .
However, one can supplement this discussion with a broader consideration of the issues concerning the evolutionary era of SH models. The equilibrium state, whenever it exists, of evolving SH models is geometrically described by a model that admits a Homothetic Vector Field (HVF or self-similarity of the first kind) or a Kinematic Self-Similarity (KSS or self-similarity of the second kind) $`𝐗`$ which is defined according to:
$$_𝐗u_a=\delta u_a_𝐗h_{ab}=2\alpha h_{ab}$$
(1.1)
where $`h_{ab}=g_{ab}+u_au_b`$ is the projection operator normal to the timelike unit vector field $`u^au_a=1`$ and $`\alpha ,\delta =`$const. essentially represent the (constant) time amplification and space dilation .
Therefore it will be enlightening to complete the geometric picture and find a relevant way to invariantly (although not uniquely) characterize the intermediate behaviour of the associated models, or to put it equivalently, one is tempted to ask: what was the nature of the generator of the self-similarity (either of the first or second kind) at the time of the intermediate evolution i.e. what kind of symmetry invariantly describes the intermediate behaviour of evolving vacuum or perfect fluid SH models? The natural and intuitive answer (although using heuristic arguments) is that one should expect that this type of symmetry must involve a generalization of the self-similarity. Consequently, the symmetry will represent the smooth transition mechanism between the evolving and the equilibrium states of SH models i.e. its “asymptotic behaviour” (into the past and the future) will be one of the well known symmetries.
In order to achieve the above goal one could consider and study general symmetries taking into account the associated local diffeomorphisms and the intrinsic geometric structure of the spacetime manifold . However, from a dynamical point of view, it appears natural to pursue a different direction by making use of the full set of EFE in order to augment and enrich the study of geometric symmetries. Clearly, this approach will provide us the necessary set of compatibility equations which can be studied in each subclass of SH (or even less symmetric) models. Motivated by the above discussion and the great success of studying SH models using elements from the theory of dynamical systems, in the present article we propose an alternative technique to study geometric symmetries that meets the aforementioned scope. This method fully incorporates the expansion-normalized orthonormal frame approach and give a transparent picture of how a specific symmetry “assumption” can be consistently endowed within a class of models or will produce further constraints, thus losing certain features of the general case. An illustrative and interesting example to which the preceding discussion applies is a recently presented new type of symmetry namely the concept of *Kinematic Conformal Symmetry* (KCS). This type of symmetry is regarded as a consistent, with geometry, generalization of the KSS or the case of conformal transformations preserving at the same time the causal structure of the spacetime manifold . Therefore, as a first step towards an invariant and geometric description of SH models, it seems adequate to study the implications from the existence of a KCS.
The structure of the paper is as follows: in section 2, after a brief review of the basic elements of the 1+3 orthonormal frame scheme and the introduction of expansion-normalized variables, we specialize our study to SH geometries and give the tetrad/expansion-normalized form of the kinematical quantities that define a generic symmetry. The determination of the corresponding expressions for the case of a KCS is treated in section 3. Since any type of symmetry assumption lead to geometric constraints which pass, through the EFE, into dynamics, we also determine the dynamical restrictions as follow from the associated integrability conditions. This set of equations, together with the expansion-normalized form of the EFE, constitute a close system of compatibility equations which can be used to check the consistent existence of a KCS in all the SH cosmologies. In section 4, special attention is given to some models of class A in which we study the consistency of the system of equations. We find that, *at least* models which are Locally Rotationally Symmetric (LRS) or lying within the invariant subset satisfying $`N_\alpha ^\alpha =0`$, always admit a KCS. A side result is that, at the asymptotic regimes, the KCS reduces to HVF or a KSS which shows that, in principle, *the existence of a four dimensional group of KCS invariantly characterizes the intermediate behaviour of the associated evolving vacuum and perfect fluid models*. Finally in section 5 we give some concluding remarks and comments and we discuss possible extensions and further advances which can be made in the direction of a satisfactory and conclusive answer to the “geometrisation” of the evolutionary behaviour of SH (or less symmetric) models.
Throughout this paper, the following conventions have been used: spatial frame indices are denoted by lower Greek letters $`\alpha ,\beta ,\mathrm{}=1,2,3`$, lower Latin letters denote spacetime indices $`a,b,\mathrm{}=0,1,2,3`$ and we use geometrized units such that $`8\pi G=c=1`$.
## 2 The generic symmetry in SH models
### 2.1 The 1+3 orthonormal frame formalism
Spatially Homogeneous models are specified in geometric terms by requiring the existence of a $`𝒢_3`$ Lie algebra of Killing Vector Fields (KVFs) $`𝐗_\alpha `$ with three-dimensional spacelike orbits $`𝒮`$. This implies the existence of a uniquely defined unit timelike vector field $`u^a`$ ($`u^au_a=1`$) normal to the spatial foliations $`𝒮`$:
$$u_{[a;b]}=0=u_{a;b}u^b\frac{1}{2}_𝐮g_{ab}=u_{a;b}=\sigma _{ab}+\frac{\theta }{3}h_{ab}$$
(2.1)
where $`\sigma _{ab},\theta =u_{a;b}g^{ab},h_{ab}=g_{ab}+u_au_b`$ are the kinematical quantities associated with the $`u^a`$ according to the standard 1+3 decomposition of an arbitrary timelike congruence . Because $`u^a`$ is irrotational and geodesic, there exists a time function $`t(x^a)`$ such that $`u^a=\delta _t^a`$ i.e. each value of $`t`$ essentially represents the hypersurfaces $`𝒮`$.
As far as the dynamical structure is concerned, from a cosmological point of view, it is sufficient to focus our study on SH models with an orthogonal perfect fluid assuming a $`\gamma `$law equation of state. Therefore the EFE become:
$$R_{ab}=\frac{3\gamma 2}{2}\rho u_au_b+\frac{2\gamma }{2}\rho h_{ab}$$
(2.2)
It follows that the timelike vector field $`u^a`$ is identified with the average fluid flow velocity.
On using the orthonormal frame formalism in SH cosmologies the starting point is to introduce a set of four linearly independent vector fields $`\{𝐞_a\}`$ and their dual 1-forms $`\{\omega ^a\}`$ which are invariant under the three-parameter group of isometries:
$$[𝐗_\alpha ,𝐞_a]=0,[𝐗_\alpha ,\omega ^a]=0.$$
(2.3)
Then, by performing a time-dependent rescaling of $`\{𝐞_a\}`$, we can write (locally) the metric tensor in a manifestly Minkowskian form:
$$ds^2=\eta _{ab}\omega ^a(t,x^\alpha )\omega ^b(t,x^\alpha ).$$
(2.4)
In this case the invariant vector fields $`\{𝐞_a\}`$ and the connection forms $`\mathrm{\Gamma }_{bc}^a\omega ^c`$ (where $`_{𝐞_c}𝐞_b=\mathrm{\Gamma }_{bc}^a𝐞_a`$) satisfy the commutation and the first Cartan structure equations:
$$[𝐞_a,𝐞_b]=\gamma _{ab}^c(t)𝐞_c,d\omega ^a=\mathrm{\Gamma }_{bc}^a\omega ^b\omega ^c=\frac{1}{2}\gamma _{bc}^a\omega ^b\omega ^c$$
(2.5)
where $`d`$ and $``$ are the usual exterior derivative and exterior product of 1-forms respectively.
It follows that the commutation functions $`\gamma _{ab}^c`$ and the connection coefficients $`\mathrm{\Gamma }_{ab}^c`$ are related via:
$$\mathrm{\Gamma }_{abc}=\frac{1}{2}\left[\eta _{ad}\gamma _{cb}^d+\eta _{bd}\gamma _{ac}^d\eta _{cd}\gamma _{ba}^d\right]\gamma _{bc}^a=\left(\mathrm{\Gamma }_{bc}^a\mathrm{\Gamma }_{cb}^a\right).$$
(2.6)
Furthermore, the requirement of the constancy of the metric under covariant differentiation implies that, $`\mathrm{\Gamma }_{(ab)c}=0`$ where $`\mathrm{\Gamma }_{abc}\eta _{da}\mathrm{\Gamma }_{bc}^d`$.
The above definitions suggest that the tetrad form of the covariant derivative of every tensor field is written in the well-known way :
$$_cK_i\mathrm{}^a\mathrm{}=𝐞_c\left(K_i\mathrm{}^a\mathrm{}\right)+\mathrm{\Gamma }_{dc}^aK_i\mathrm{}^d\mathrm{}+\mathrm{}\mathrm{\Gamma }_{ic}^dK_d\mathrm{}^a\mathrm{}\mathrm{}$$
(2.7)
where $`𝐞_c\left(K_i\mathrm{}^a\mathrm{}\right)`$ is regarded as the directional derivative of the functions $`K_i\mathrm{}^a\mathrm{}`$ along the vector field $`𝐞_c`$.
Because in SH (vacuum or perfect fluid) models there always exists the preferred and well defined timelike vector field $`u^a`$ associated with a congruence of curves normal to the three-dimensional spacelike orbits $`𝒮`$ of homogeneity, it is natural to select it as the timelike frame vector i.e. $`𝐞_0=𝐮`$. It follows from equation (2.5) that the kinematical quantities of the timelike congruence $`𝐮`$ are directly related with the commutation functions $`\gamma _{ab}^c`$ according to :
$$\gamma _{0\alpha }^0=\dot{u}_\alpha =0,\gamma _{\alpha \beta }^0=2\epsilon _{\alpha \beta \gamma }\omega ^\gamma =0$$
(2.8)
$$\gamma _{0\alpha }^\beta =\frac{1}{3}\theta \delta _\alpha ^\beta \sigma _\alpha ^\beta +\epsilon _{\alpha \gamma }^\beta \mathrm{\Omega }^\gamma $$
(2.9)
where $`\mathrm{\Omega }^\gamma `$ is the local angular velocity of the spatial frame with respect to a Fermi-propagated frame along $`𝐞_0`$. On the other hand, the spatial components of $`\gamma _{ab}^c`$ are decomposed as:
$$\gamma _{\beta \gamma }^\alpha =a_\beta \delta _\gamma ^\alpha a_\gamma \delta _\beta ^\alpha +\epsilon _{\beta \gamma \rho }n^{\alpha \rho }$$
(2.10)
leading to Bianchi class A and B models according to whether the quantity $`a_\beta `$ vanishes or not.
Combining equations (2.6) and (2.8)-(2.10) we easily find:
$$\mathrm{\Gamma }_{\beta 0\alpha }=\frac{\theta }{3}\delta _{\alpha \beta }+\sigma _{\alpha \beta },\mathrm{\Gamma }_{0\alpha 0}=0,\mathrm{\Gamma }_{\alpha \beta 0}=\epsilon _{\alpha \beta \gamma }\mathrm{\Omega }^\gamma $$
(2.11)
$$\mathrm{\Gamma }_{\alpha \beta \gamma }=2a_{[\alpha }\delta _{\beta ]\gamma }+\epsilon _{\gamma \rho [\alpha }n^\rho {}_{\beta ]}{}^{}+\frac{1}{2}\epsilon _{\alpha \beta \rho }n^\rho _\gamma $$
(2.12)
The complete set of the gravitational field equations is expressed in terms of the shear components $`\sigma _{\alpha \beta }`$, the expansion $`\theta `$ and the spatial curvature quantities $`a_\beta `$, $`n_{\alpha \beta }`$, by utilizing the Ricci identity for $`u^a`$, the Jacobi identities for the frame vector fields and the Bianchi identities.
### 2.2 Expansion normalized variables
Of particular importance in the exploration of the asymptotic dynamics of SH models, is the reformulation of the complete set of orthonormal frame equations, as an autonomous system of first order ordinary differential equations. This can be done by defining a set of expansion-normalized (dimensionless) variables which results the decoupling of the evolution equation of $`H=\theta /3`$ from the rest of the evolution equations:
$$\frac{dt}{d\tau }=\frac{1}{H},\frac{dH}{d\tau }=\left(1+q\right)H$$
(2.13)
where $`q,H`$ are the deceleration and Hubble parameter respectively and $`\tau `$ is the dimensionless time variable.
The complete set of equations can be written in the form :
$$\mathrm{\Sigma }_{\alpha \beta }^{}=\left(2q\right)\mathrm{\Sigma }_{\alpha \beta }+2ϵ_{(\alpha }^{\mu \nu }\mathrm{\Sigma }_{\beta )\mu }R_\nu S_{\alpha \beta }$$
(2.14)
$$N_{\alpha \beta }^{}=qN_{\alpha \beta }+2\mathrm{\Sigma }_{(\alpha }^\mu N_{\beta )\mu }+2ϵ_{(\alpha }^{\mu \nu }N_{\beta )\mu }R_\nu $$
(2.15)
$$A_\alpha ^{}=qA_\alpha \mathrm{\Sigma }_\alpha ^\mu A_\mu +ϵ_\alpha ^{\mu \nu }A_\mu R_\nu $$
(2.16)
$$\mathrm{\Omega }^{}=\mathrm{\Omega }\left[2q\left(3\gamma 2\right)\right]$$
(2.17)
where:
$`S_{\alpha \beta }`$ $`=`$ $`2N_\alpha ^\gamma N_{\beta \gamma }N_\gamma ^\gamma N_{\alpha \beta }{\displaystyle \frac{1}{3}}\left(2N_\alpha ^\gamma N_{\beta \gamma }N_\gamma ^\gamma N_{\alpha \beta }\right)\delta ^{\alpha \beta }`$ (2.18)
$`2ϵ_{(\alpha }^{\mu \nu }N_{\beta )\mu }A_\nu `$
and a prime denotes derivative w.r.t. $`\tau `$.
The above system is subjected to the algebraic constraints:
$$\mathrm{\Omega }=1\frac{1}{6}\mathrm{\Sigma }^{\alpha \beta }\mathrm{\Sigma }_{\alpha \beta }K$$
(2.19)
$$3\mathrm{\Sigma }_\alpha ^\beta A_\beta ϵ_\alpha ^{\mu \nu }\mathrm{\Sigma }_\mu ^\beta N_{\beta \nu }=0$$
(2.20)
where:
$$K=\frac{1}{12}\left(2N_\alpha ^\gamma N_{\beta \gamma }N_\gamma ^\gamma N_{\alpha \beta }\right)\delta ^{\alpha \beta }+A_\gamma A^\gamma $$
(2.21)
and the deceleration parameter is given by the relation:
$`q`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{\Sigma }^{\alpha \beta }\mathrm{\Sigma }_{\alpha \beta }+{\displaystyle \frac{1}{2}}\mathrm{\Omega }\left[\left(3\gamma 2\right)\right]`$ (2.22)
$`=`$ $`2\left(1K\right)+{\displaystyle \frac{1}{2}}\mathrm{\Omega }\left[3\left(\gamma 2\right)\right].`$
### 2.3 The generic symmetry in expansion-normalized variables
The folklore for investigating the implications of the existence of geometric symmetries in General Relativity can be divided into two main categories. The first category is a geometric approach in which we study symmetries taking into account the holonomy group structure of the spacetime manifold together with the associated local diffeomorphisms . In the second category one formulates the necessary and sufficient conditions, coming from the existence of the symmetry, in a covariant way and study their consequences in the kinematics and dynamics of the corresponding model . Of course one could also deal directly with the resulting system of partial differential equations (pdes), presupposing a specific geometrical and dynamical configuration which render the symmetry equations to be more tractable. Obviously this approach undergo many disadvantages and pathologies. One of the serious stumbling blocks is the fact that as the generality of a model is increased (i.e. the underlying geometric structure of the model is less symmetric than the SH geometry) the symmetry pdes are progressively non-linear and very often lead to solutions of the EFE which are immediately ruled out physically.
Here we suggest an alternative approach for the study of geometric symmetries which fully exploits the well-established orthonormal frame formalism in terms of expansion-normalized variables. Although this technique will be applied to a specific symmetry “assumption” (as we shall see in the next section this is not really an assumption, at least for a class of models, but a consequence of the complete set of EFE) it can be used in a more general context for the study of other types of important symmetries .
Let us consider an arbitrary vector field $`𝐗`$ and express its components in terms of the frame vector fields $`𝐞_a`$:
$$𝐗=X^a𝐞_a=\lambda 𝐞_0+X^\alpha 𝐞_\alpha $$
(2.23)
where $`\lambda ,X^\alpha `$ are continuously differential functions of the spacetime manifold.
The first derivatives of $`𝐗`$ are decomposed into irreducible symmetry kinematical parts $`\{\psi ,H_{ab},F_{ab}\}`$ in the standard way:
$$_bX_a=\psi g_{ab}+H_{ab}+F_{ba}$$
(2.24)
where
$$4\psi _kX^k,$$
(2.25)
$$H_{ab}=\left[_{(b}X_{a)}\frac{1}{4}\left(_kX^k\right)g_{ab}\right]$$
(2.26)
$$F_{ab}=_{[b}X_{a]}$$
(2.27)
are the conformal factor, the traceless symmetric part and the antisymmetric part respectively.
Using the definition (2.7) and equations (2.11), (2.12), (2.23)-(2.27), we find the tetrad analogue of the kinematical quantities. If we further invoke the expansion-normalized differential operators $`_a𝐞_a/H`$ in the general expressions we finally obtain:
$$4\psi =H\left[_0\left(\lambda \right)+_\alpha (X^\alpha )+3\lambda 2A_\alpha X^\alpha \right]$$
(2.28)
$$H_{00}=\frac{H}{4}\left[3_0(\lambda )+_\alpha (X^\alpha )+3\lambda 2A_\alpha X^\alpha \right]$$
(2.29)
$$2H_{0\alpha }=H\left[_\alpha (\lambda )+_0(X_\alpha )\left(\delta _{\alpha \gamma }+\mathrm{\Sigma }_{\alpha \gamma }\right)X^\gamma +\epsilon _{\alpha \beta \gamma }R^\gamma X^\beta \right]$$
(2.30)
$`H_{\alpha \beta }`$ $`=`$ $`H\{_{(\beta }X_{\alpha )}+(\delta _{\alpha \beta }+\mathrm{\Sigma }_{\alpha \beta })\lambda [A_\gamma \delta _{\alpha \beta }A_{(\alpha }\delta _{\beta )\gamma }]X^\gamma `$ (2.31)
$`+\epsilon _{\gamma \delta (\alpha }N_{\beta )}^\delta X^\gamma \}\psi \delta _{\alpha \beta }`$
$$2F_{0\alpha }=H\left[_\alpha (\lambda )+_0(X_\alpha )+\left(\delta _{\alpha \gamma }+\mathrm{\Sigma }_{\alpha \gamma }\right)X^\gamma \epsilon _{\alpha \beta \gamma }R^\gamma X^\beta \right]$$
(2.32)
$$F_{\alpha \beta }=H\left\{_{[\beta }X_{\alpha ]}\left[A_{[\alpha }\delta _{\beta ]\gamma }+\frac{1}{2}\epsilon _{\alpha \beta \delta }N_\gamma ^\delta \right]X^\gamma \right\}.$$
(2.33)
Equations (2.28)-(2.33) represent the *symmetry kinematical quantities in terms of expansion-normalized variables* in SH geometries and can be used in order to have a first hint of how the dynamics affects on the geometry of SH models (or vice-versa). We note that, one could choose to define expansion-normalized symmetry kinematical quantities in a similar way we do for the shear and spatial curvature variables. This would be convenient for high symmetries where first or second derivatives of $`\{\psi ,H_{ab},F_{ab}\}`$ are involved. However for the cases we are interested in the present article the use of (2.28)-(2.33) will be satisfactory.
We conclude this section by pointing out that, because any type of symmetry assumption is described in terms of geometric constraints, these are inherited by the dynamics through the EFE (2.2). Therefore in order to visualize how the symmetry further interacts with the dynamics, it is necessary to determine the effect of the former on the Ricci tensor. By using the well-known commutation relation between the connection and the Lie derivative :
$$_𝐗\mathrm{\Gamma }_{bc}^a=\frac{1}{2}g^{ar}\left[_c\left(_𝐗g_{br}\right)+_b\left(_𝐗g_{cr}\right)_r\left(_𝐗g_{bc}\right)\right]$$
$$_𝐗R_{ab}=_c\left[\mathrm{\Gamma }_{\left(ab\right)}^c\right]_{(b}[\mathrm{\Gamma }^c{}_{a)c}{}^{}]$$
and the defining equation (2.24), we can show after a straightforward calculation that:
$$_𝐗R_{ab}=2_b_a\psi g_{ab}_c^c\psi +2_k_{(b}H^k{}_{a)}{}^{}_c^cH_{ab}.$$
(2.34)
Essentially, the last equation represents a set of integrability conditions which can be used to check the consistent existence of *every type* of symmetry assumption.
## 3 Generalized Conformal Symmetries in SH models
Although there exists (and can be defined) a sufficiently large number of symmetries, the most important type of them (up to date) appears to concern the constant scale invariance of the metric represented by the existence of a proper HVF . For SH vacuum and perfect fluid models this is indeed the case due to the profound relevance of homothetic models with the equilibrium points of the SH state space .
Recently, a new type of symmetry has been suggested, the so-called *bi-conformal transformations* which can be interpreted as generalizing the concepts of the self-similarity and the conformal motions. In the present work we will concern with an interesting subcase, that is, the so-called *Kinematic Conformal Symmetry* (KCS). In particular a smooth vector field $`𝐗`$ is the generator of a KCS *iff* the following relations hold :
$$_𝐗u_a=\delta u_a,_𝐗h_{ab}=2\alpha h_{ab}$$
(3.1)
or, in terms of the metric:
$$_𝐗g_{ab}=2\alpha g_{ab}+2\left(\alpha \delta \right)u_au_b$$
(3.2)
where $`\alpha ,\delta `$ are smooth functions that we shall call *symmetry scales* and $`h_{ab}=g_{ab}+u_au_b`$ is the projection operator perpendicular to the timelike congruence $`u_a`$.
Combining equations (2.24) and (3.2) we express the symmetry kinematical parts in the form:
$$\psi =\frac{3\alpha +\delta }{4},H_{ab}=\frac{\alpha \delta }{4}\left(g_{ab}+4u_au_b\right).$$
(3.3)
It can be easily observed that when $`\alpha =\delta `$ the KCS reduces to a Conformal Vector Field (CVF) which, due to the equation (3.1), is necessarily inheriting i.e. the integral curves of $`u^a`$ are mapped conformally by the CVF $`𝐗`$ . As a result the Lie algebra $``$ of inheriting CVFs is always a subalgebra of the Lie algebra of KCS which shall be denoted as $``$. Moreover when the symmetry scales $`\alpha ,\delta `$ are both (different) constants we recover the case of a Kinematic Self-Similarity.
### 3.1 Symmetry constraints
Clearly there exists a direct dependence between the existence of a KCS, as well as any other type of symmetry assumption, and a specific cosmological model. This mutual influence is reflected in the induced geometric, kinematic and dynamic constraints which are imposed due to the intrinsic nature of the symmetry and/or as a consequence of the specific geometric and dynamical structure of the spacetime manifold. In the case of a KCS these constraints are derived from the general relations (2.28)-(2.31) and the symmetry assumptions (3.3):
$$3\alpha =H\left[_\alpha (X^\alpha )+3\lambda 2A_\alpha X^\alpha \right]$$
(3.4)
$$\left[4H_0(\lambda )+3\alpha +\delta \right]=3\left(\alpha \delta \right)\delta =H_0(\lambda )$$
(3.5)
$$_\alpha (\lambda )+_0(X_\alpha )\left(\delta _{\alpha \gamma }+\mathrm{\Sigma }_{\alpha \gamma }\right)X^\gamma +\epsilon _{\alpha \beta \gamma }R^\gamma X^\beta =0$$
(3.6)
$`H\{_{(\beta }X_{\alpha )}+(\delta _{\alpha \beta }+\mathrm{\Sigma }_{\alpha \beta })\lambda [A_\gamma \delta _{\alpha \beta }A_{(\alpha }\delta _{\beta )\gamma }]X^\gamma `$
$`+\epsilon _{\gamma \delta (\alpha }N_{\beta )}^\delta X^\gamma \}\alpha \delta _{\alpha \beta }=0.`$ (3.7)
We should emphasize that the above set of constraints must be augmented with the associated Jacobi identities satisfied by the generators $`\{𝐗_\alpha ,𝐗\}`$ that constitute the Lie Algebra $``$ of KCS. This will require the determination of the dimension of $``$ and the assumption that the former is finite since there are cases for which $``$ is infinite dimensional . Nevertheless we shall not pursue the problem in full generality and we will restrict our considerations to the case where the dimension of $``$ is finite and equal to four (together with the $`𝒢_3`$ of KVFs which can be seen as “trivial” KCS). In fact this assumption seems reasonable and not even restrictive, from a geometrical point of view, because we intend to describe the “intermediate behaviour” of a (proper) self-similarity of the first or second kind in a way that is identified with the presence of a KCS. In this case it can be shown from Jacobi identities that the Lie bracket of a KCS with the KVFs is always a linear combination of the later. It turns out that the scalar $`\lambda `$ and the symmetry scales $`\alpha ,\delta `$ (due to equation (3.1)) are spatially homogeneous:
$$_\alpha \lambda =_\alpha \alpha =_\alpha \delta =0.$$
(3.8)
In addition we point out that equation (3.2) is a consequence of (3.1) since the former suppresses the scale amplification of both time and space which is essentially represented by the latter. However the spatial homogeneity condition (3.8) ensures that the definitions (3.2) and (3.1) are equivalent.
### 3.2 Integrability conditions
Due to the geometric character of the KCS to preserve the causal structure of the spacetime manifold, it is natural to expect that the associated transformation group will respect, to some level, the intrinsic properties of SH models. Therefore one should expect that the presence of a KCS will induce weaker constraints than the case of CVFs or the KSS. In order to confirm this, it will be required to determine the integrability conditions adapted to the case of SH models filled, in general, with a $`\gamma `$law non-tilted perfect fluid.
As we have shown, the scale functions and the trace $`\psi `$ are both spatially homogeneous which implies that equation (2.34), due to (3.3), can be written:
$`_𝐗R_{ab}`$ $`=`$ $`3\left[H\left(\dot{\delta }2\dot{\alpha }\right)\ddot{\alpha }\right]u_au_b`$ (3.9)
$`+\left[2\left(\alpha \delta \right)\left(3H^2+\dot{H}\right)+\ddot{\alpha }+H\left(6\dot{\alpha }\dot{\delta }\right)\right]h_{ab}`$
$`+\left[\left(3\dot{\alpha }\dot{\delta }\right)\sigma _{ab}+\left(\alpha \delta \right)\left(6H\sigma _{ab}+2\dot{\sigma }_{ab}\right)\right]`$
where a dot “$``$” denotes differentiation w.r.t. $`u^a`$.
On the other hand the EFE (2.2) and the relations (3.1) give:
$$_𝐗R_{ab}=\frac{3\gamma 2}{2}\left[𝐗\left(\rho \right)+2\delta \rho \right]u_au_b+\frac{2\gamma }{2}\left[𝐗\left(\rho \right)+2\alpha \rho \right]h_{ab}.$$
(3.10)
The complete set of integrability conditions, coming from the existence of a KCS, follow by equating (3.9) and (3.10):
$$3\left[H\left(\dot{\delta }2\dot{\alpha }\right)\ddot{\alpha }\right]=\frac{3\gamma 2}{2}\left[𝐗\left(\rho \right)+2\delta \rho \right]$$
(3.11)
$$2\left(\alpha \delta \right)\left(3H^2+\dot{H}\right)+\ddot{\alpha }+H\left(6\dot{\alpha }\dot{\delta }\right)=\frac{2\gamma }{2}\left[𝐗\left(\rho \right)+2\alpha \rho \right]$$
(3.12)
$$\left(3\dot{\alpha }\dot{\delta }\right)\sigma _{ab}+\left(\alpha \delta \right)\left(6H\sigma _{ab}+2\dot{\sigma }_{ab}\right)=0.$$
(3.13)
We observe that in the case of a CVF where $`\alpha =\delta `$, the last equation implies the well-known result $`\dot{\alpha }\sigma _{ab}=0`$ i.e. either the (necessary inheriting) CVF reduces to a HVF or the spacetime is the Robertson-Walker spacetime .
The system of equations (3.11)-(3.13) can be conveniently reformulated in expansion-normalized variables in order to append them in the autonomous set (2.14)-(2.17) and the symmetry equations (3.4)-(3.7). For simplification purposes we define the dimensionless symmetry scale functions:
$$\stackrel{~}{\alpha }=\frac{\alpha }{H},\stackrel{~}{\delta }=\frac{\delta }{H}.$$
(3.14)
Then, taking into account equations (2.14) and (2.17), an appropriate combination of (3.11)-(3.13) eliminates the second order time-derivatives and give evolution equations for $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\delta }`$:
$`4\left[\stackrel{~}{\alpha }^{}\left(q+1\right)\stackrel{~}{\alpha }\right]`$ $`=`$ $`2\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)\left(2q\right)`$ (3.15)
$`+\left[\left(3\gamma 2\right)\stackrel{~}{\delta }+3\left(2\gamma \right)\stackrel{~}{\alpha }6\lambda \gamma \right]\mathrm{\Omega }`$
$$\left[\left(3\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)^{}\left(q+1\right)\left(3\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)\right]\mathrm{\Sigma }_{\alpha \beta }2\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)S_{\alpha \beta }=0.$$
(3.16)
It is interesting to note that equation (3.16) *excludes the existence of a proper KSS in Bianchi vacuum or perfect fluid cosmologies with $`S_{\alpha \beta }0`$*. The constraint $`S_{\alpha \beta }=0`$ is identically satisfied in Kasner type I models which is well known to admit a (proper) self-similarity of the second kind . This result implies that a KSS fails to be considered as a generic candidate to describe the intermediate behaviour of general SH models, but rather one should explore the possibility for a symmetry, representing a direct generalization of the conformal motions, in such a way that its asymptotic behaviour is the self-similarity of the first or second kind. We shall demonstrate in the next section that the case of a KCS provides an evidence towards a satisfactory (but not conclusive) answer to this question, for a significant subclass of evolving SH vacuum and perfect fluid models.
## 4 Application to SH models of class A
Spatially homogeneous cosmologies of class A are defined by the condition $`A^\alpha =0`$, which due to equation (2.20), implies that the shear and spatial curvature matrices $`\mathrm{\Sigma }_{\alpha \beta },N_{\alpha \beta }`$ commute. Therefore $`\mathrm{\Sigma }_{\alpha \beta },N_{\alpha \beta }`$ have a common eigenframe and we can write:
$$\mathrm{\Sigma }_{\alpha \beta }=\text{diag}(\mathrm{\Sigma }_{11},\mathrm{\Sigma }_{22},\mathrm{\Sigma }_{22}),N_{\alpha \beta }=\text{diag}(N_{11},N_{22},N_{22}).$$
(4.1)
Furthermore the evolution equations (2.14), (2.15) and equation (4.1) show that $`R^\alpha =0`$ which means that the common eigenframe of $`\mathrm{\Sigma }_{\alpha \beta },N_{\alpha \beta }`$ is also Fermi-propagated. Under these conditions the study of the class A models is considerably simplified and permits one to investigate the set of consistency equations constituting of the EFE (2.14)-(2.17), the geometric constraints (3.4)-(3.7) and the integrability conditions (3.15)-(3.16) in a straightforward way. To illustrate the method that could be used for the consistency checking we outline some applications by restricting our study to type I, II and VI<sub>0</sub> models.
However, before we proceed, it will be useful to give the corresponding analysis for the flat Friedmann-Lemaître model $``$ since it represents the past or future attractor for several SH models of class A.
### 4.1 Flat Friedmann-Lemaître $`\gamma `$law perfect fluid models
It is well known that the (flat) Robertson-Walker spacetime:
$$ds^2=dt^2+S^2(t)\left(dx^2+dy^2+dz^2\right)$$
(4.2)
has a variety of ways for an invariant characterization of its structure. For example, kinematically, is defined by the vanishing of the shear, vorticity and acceleration of the preferred timelike congruence $`u^a`$ which, due to equations (2.19) and (2.22), implies:
$$\mathrm{\Omega }=1,q=\frac{3\gamma 2}{2}$$
(4.3)
On the other hand one can use geometric terms and describe the Friedmann-Lemaître model by the existence of nine proper CVFs, one of which is parallel to $`u^a`$ . Although in the case of a $`\gamma `$law perfect fluid model a proper HVF and KSS exists (lowering the dimension of the Lie algebra of conformal motions to eight) this does not mean that we have exhausted all the possible (geometric) ways for the description of Friedmann-Lemaître models. As a result one should expect that a KCS will also exists without imposing extra geometrical or dynamical restrictions (which is often the case for other types of symmetries). Indeed from the constraints (3.4)-(3.7) we find:
$$3(\stackrel{~}{\alpha }\lambda )=_\alpha (X^\alpha ),\stackrel{~}{\delta }=\lambda ^{}$$
(4.4)
$$_0(X_\alpha )X_\alpha =0,_{(\beta }X_{\alpha )}+\left(\lambda \stackrel{~}{\alpha }\right)\delta _{\alpha \beta }=0.$$
(4.5)
Using the set of equations (4.3)-(4.5) and the integrability condition (3.15) we can show that a KCS *always exists* in a Friedmann-Lemaître model and the symmetry scale $`\stackrel{~}{\alpha }`$ is given in terms of the temporal component:
$$\left(\lambda \stackrel{~}{\alpha }\right)^{}=\left(q+1\right)\left(\lambda \stackrel{~}{\alpha }\right)=\frac{3\gamma }{2}\left(\lambda \stackrel{~}{\alpha }\right).$$
(4.6)
From the above equations we determine the exact form of the KCS in the Robertson-Walker spacetime:
$$𝐗=\lambda (t)_t+c\left(x_x+y_y+z_z\right)$$
(4.7)
where:
$$\alpha (t)=\frac{2\lambda (t)+3c\gamma t}{3t\gamma }$$
(4.8)
and the scale factor is given by $`S(t)=t^{2/3\gamma }`$. We note that, after a suitable change of the basis of the KCS Lie algebra, we can set the constant $`c=0`$ which implies that $`𝐗`$ is also parallel to the fluid velocity $`u^a`$.
### 4.2 Type I models
In this case $`N_{\alpha \beta }=0=S_{\alpha \beta }`$ and equation (3.16) gives:
$$3\stackrel{~}{\alpha }\stackrel{~}{\delta }=\stackrel{~}{c}\frac{c}{H}3\alpha \delta =c$$
(4.9)
where $`c`$ is an arbitrary constant.
In addition, equation (3.5) implies that:
$$\stackrel{~}{\delta }=\lambda ^{}.$$
(4.10)
We should pointed out that the symmetry constraints are necessary to ensure the existence of a KCS. However this does not imply that they will be preserved along the integral curves of the timelike vector field $`𝐞_0`$. Therefore we must propagate equations (3.4)-(3.7) in order to retain the existence of a KCS in *every* spacelike hypersurface $`𝒮`$. After a short calculation and the use of the commutator relations (2.9), we obtain:
$$\lambda ^{}\left(q+1\right)\lambda =\stackrel{~}{\alpha }^{}\left(q+1\right)\stackrel{~}{\alpha }$$
(4.11)
$$\left[\left(q2\right)\lambda +\lambda ^{}\left(q+1\right)\lambda \right]\mathrm{\Sigma }_{\alpha \beta }+\mathrm{\Lambda }_{\alpha \beta }=0$$
(4.12)
where
$$\mathrm{\Lambda }_{\alpha \beta }=\mathrm{\Sigma }_{\gamma (\alpha }X_{,\beta )}^\gamma \mathrm{\Sigma }_{\gamma (\alpha }X_{\beta )}^{,\gamma }$$
(4.13)
and $`X_{,\beta }^\gamma _\beta X^\gamma `$.
From the $`\alpha \alpha `$component of (4.12) it follows that:
$$\left(q2\right)\lambda +\lambda ^{}\left(q+1\right)\lambda =0\lambda ^{}=3\lambda $$
(4.14)
while the equation $`\mathrm{\Lambda }_{\alpha \beta }=0`$ expresses the spatial components of the KCS in terms of the shear variables.
In summary we have shown that *every type I* $`\gamma `$*law perfect fluid model always admits a four dimensional group of Kinematic Conformal Symmetries*. An interesting feature of this result concerns the past “asymptotic behaviour” of the KCS. In particular, from equation (4.14) we observe that $`\lambda ^{}\left(q+1\right)\lambda =\left(2q\right)\lambda `$ hence at the equilibrium point ($`q=2`$) we have $`\lambda ^{}\left(q+1\right)\lambda =0`$ i.e. $`\alpha ,\delta =`$const. and the KCS reduces to a proper KSS with $`\lambda =\text{const.}\times H^1`$. The future state of the KCS is treated similarly. The type I models, approach at late times the Friedmann-Lemaître model $``$, so equation (4.6) is trivially satisfied and the temporal component is given in equation (4.14). This implies that $`\alpha ^{}\left(q+1\right)\alpha 0\alpha ,\delta `$const. as expected. Therefore, one could argue that *the existence of a proper KCS describes, in a geometric fashion, the intermediate behaviour of the evolving type I models*.
We conclude the type I case by giving the exact form of the KCS in local coordinates. The general $`\gamma `$law perfect fluid solution can be conveniently written in the form :
$$ds^2=A^{2(\gamma 1)}dt^2+\underset{\alpha }{}t^{2p_\alpha }A^{2\left(2/3p_\alpha \right)}\left(dx^\alpha \right)^2$$
(4.15)
where the function $`A^{2\gamma }=k+m^2t^{2\gamma }`$ and $`k,m`$ are constants.
It follows that the (local) coordinate form of the KCS is:
$$𝐗=\lambda (t)_t+\left(ck+c_3\right)x_x+c_2y_y+c_3z_z$$
(4.16)
where:
$$\lambda (t)=\frac{tkcA}{\left(p_2+2p_31\right)\left(AtA_{,t}\right)}$$
(4.17)
and the constant $`c_2`$ is given by:
$$c_2=\frac{c_3\left(p_2+2p_31\right)+kc\left(p_3p_2\right)}{p_2+2p_31}$$
(4.18)
where we have used the well-known relations:
$$p_1+p_2+p_3=1,p_1^2+p_2^2+p_3^2=1$$
(4.19)
satisfied by the Kasner exponents $`p_\alpha `$.
We also note that, at early times $`\lambda (t)t`$ and $`\alpha ,\delta =`$const. whereas at late times, $`k=0`$ and $`\alpha ,\delta `$const. which confirm the reduction of the KCS to a proper KSS and to the KCS (4.7) respectively.
### 4.3 Type II models
The Bianchi type II invariant set is characterized by the conditions $`N_1>0`$ and $`N_2=N_3=0`$. We find convenient to collect the consistency equations as follow from (2.14)-(2.17), (3.4)-(3.7) and (3.15)-(3.16):
$$\lambda ^{}\left(q+1\right)\lambda =\stackrel{~}{\alpha }^{}\left(q+1\right)\stackrel{~}{\alpha }$$
(4.20)
$`0`$ $`=`$ $`\left(\lambda ^{}3\lambda \right)\mathrm{\Sigma }_{\alpha \beta }+\mathrm{\Lambda }_{\alpha \beta }+\left[2\mathrm{\Sigma }_{11}X^\gamma +\mathrm{\Sigma }_\epsilon ^\gamma X^\epsilon \right]N_1\epsilon _{\gamma 1(\alpha }\delta _{\beta )}^1`$ (4.21)
$`\left({\displaystyle \frac{2}{3}}\delta _\alpha ^1\delta _\beta ^1{\displaystyle \frac{1}{3}}\delta _\alpha ^2\delta _\beta ^2{\displaystyle \frac{1}{3}}\delta _\alpha ^3\delta _\beta ^3\right)N_1^2\lambda `$
$`4\left[\stackrel{~}{\alpha }^{}\left(q+1\right)\stackrel{~}{\alpha }\right]`$ $`=`$ $`2\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)\left(2q\right)`$ (4.22)
$`+\left[\left(3\gamma 2\right)\stackrel{~}{\delta }+3\left(2\gamma \right)\stackrel{~}{\alpha }6\lambda \gamma \right]\mathrm{\Omega }`$
$`0`$ $`=`$ $`\left[\left(3\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)^{}\left(1+q\right)\left(3\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)\right]\mathrm{\Sigma }_{\alpha \beta }`$ (4.23)
$`\left({\displaystyle \frac{2}{3}}\delta _\alpha ^1\delta _\beta ^1{\displaystyle \frac{1}{3}}\delta _\alpha ^2\delta _\beta ^2{\displaystyle \frac{1}{3}}\delta _\alpha ^3\delta _\beta ^3\right)N_1^2\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)`$
$$\stackrel{~}{\delta }=\lambda ^{}$$
(4.24)
$$1\frac{1}{12}N_1^2\mathrm{\Sigma }^2\mathrm{\Omega }=0,q=2\mathrm{\Sigma }^2+\frac{3\gamma 2}{2}\mathrm{\Omega }$$
(4.25)
$$\mathrm{\Sigma }_{\alpha \beta }^{}=(q2)\mathrm{\Sigma }_{\alpha \beta }\left(\frac{2}{3}\delta _\alpha ^1\delta _\beta ^1\frac{1}{3}\delta _\alpha ^2\delta _\beta ^2\frac{1}{3}\delta _\alpha ^3\delta _\beta ^3\right)N_1^2$$
(4.26)
$$N_1^{}=\left(q+2\mathrm{\Sigma }_{11}\right)N_1.$$
(4.27)
Note that, in complete analogy with the type I models, equations (4.20) and (4.21) are the result of the propagation of the symmetry constraints (3.4) and (3.7) along $`_0`$.
A quick observation can be made, due to the form of (4.21) or (4.23). For example, the first equation implies<sup>1</sup><sup>1</sup>1We recall the traceless property of $`\mathrm{\Sigma }_{\alpha \beta }`$ i.e. $`\mathrm{\Sigma }_{11}+\mathrm{\Sigma }_{22}+\mathrm{\Sigma }_{33}=0`$.:
$$\left(\lambda ^{}3\lambda \right)\mathrm{\Sigma }_{22}=\left(\lambda ^{}3\lambda \right)\mathrm{\Sigma }_{33}=\frac{1}{3}N_1^2\lambda $$
(4.28)
where for $`\mathrm{\Sigma }_{22}=0`$ the KCS turns to the KVF of Bianchi type II models.
Equation (4.28) means that the existence of a KCS is compatible *only with type II models which are LRS* i.e. only within the invariant subset $`S_1^+(II)`$. In order to determine the symmetry scale $`\stackrel{~}{\alpha }`$ we make use of the $`\alpha \beta `$components of equation (4.21). Then, from the spatial commutators (2.10) we get $`X_{2,1}=X_{3,1}=0`$ and the consistency of the remaining set of equations is assured provided that:
$$\stackrel{~}{\alpha }=\lambda \left(4\mathrm{\Sigma }_{22}+1\right).$$
(4.29)
We have proved that *every type II evolving vacuum or* $`\gamma `$*law perfect fluid model that belongs to the invariant subset* $`S_1^+(II)`$ *can be invariantly characterized by the existence of a four dimensional group of Kinematic Conformal Symmetries*.
Regarding the “asymptotic behaviour” of the KCS in type II models, it is straightforward to show that in $`S_1^+(II)|_{\mathrm{\Omega }=0}`$ the following relations hold:
$$\stackrel{~}{\alpha }^{}\left(q+1\right)\stackrel{~}{\alpha }\mathrm{\Sigma }_{22}^21$$
$$\left(3\stackrel{~}{\alpha }+\stackrel{~}{\delta }\right)^{}\left(q+1\right)\left(3\stackrel{~}{\alpha }+\stackrel{~}{\delta }\right)\mathrm{\Sigma }_{22}^21.$$
Therefore, at the equilibrium points $`\mathrm{\Sigma }_{22}=\pm 1`$ and $`N_1=0`$, the KCS becomes a proper KSS as expected, since at the asymptotic regimes, all models within $`S_1^+(II)|_{\mathrm{\Omega }=0}`$ approach some vacuum Kasner model .
On the other hand, it is well known that non-vacuum models $`S_1^+(II)|_{\mathrm{\Omega }>0}`$ are all future asymptotic to the homothetic Collins model $`P_1^+(II)`$ for which we have proved that does not admit a proper KSS, and past asymptotic to a Kasner model $`𝒦`$ or to the Friedmann-Lemaître model $``$. Consequently, every orbit in $`S_1^+(II)|_{\mathrm{\Omega }>0}`$ joins two (first and second kind) self-similar equilibrium points, hence we expect that the KCS will be reduced to a proper HVF and a KSS, except to the case where the orbit lies in the one-dimensional unstable manifold of $``$. This model has past attractor the point $``$ and the KCS will reduce to the associated KCS of the Robertson-Walker spacetime (equation (4.7)).
Indeed, using equations (4.24)-(4.29) we find:
$$\stackrel{~}{\alpha }\stackrel{~}{\delta }=\frac{\lambda \left[N_1^2+6\mathrm{\Sigma }_{22}\left(2\mathrm{\Sigma }_{22}1\right)\right]}{3\mathrm{\Sigma }_{22}}.$$
From the last equation we can show easily that $`\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)|_{P_1^+(II)}=0`$ which implies that $`\underset{\tau +\mathrm{}}{lim}\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)=0`$. Similarly we can show that, at the equilibrium point $`𝒦`$ the symmetry scales $`\alpha ,\delta `$ become constants with $`\alpha |_𝒦\delta |_𝒦`$ and at $``$ we have $`\alpha =\lambda H`$, $`\lambda ^{}=3\lambda `$.
The exact form of the KCS in local coordinates is found to be:
$$𝐗=\lambda (t)_t+\left(2cx+by\right)_x+cy_y+\left(cz1\right)_z$$
(4.30)
where $`c`$ is constant and:
$$\lambda (t)=\frac{2cBC}{BC_{,t}CB_{,t}}.$$
(4.31)
The LRS type II metric is:
$$ds^2=A^2dt^2+B\left(dx+bzdy\right)^2+C\left(dy^2+dz^2\right)$$
(4.32)
for smooth functions $`B(t),C(t)`$ of their argument.
For example, in the LRS vacuum model (a special case of the general solution found by Taub ) the metric is :
$$ds^2=A^2dt^2+t^{2p_1}A^2\left(dx+4p_1bzdy\right)^2+t^{2p_3}A^2\left(dy^2+dz^2\right)$$
(4.33)
where $`A=\left(1+b^2t^{4p_1}\right)^{1/2}`$ and $`p_\alpha `$ satisfy (4.19).
The temporal component of the KCS is given by:
$$\lambda (t)=\frac{tA}{2tA_{,t}\left(p_1p_3\right)A}$$
(4.34)
and reduces to a KSS for small and large values of the time coordinate (the KSS in plane symmetric Bianchi models has also been found in ).
As a final remark we note that by defining the new time coordinate $`\stackrel{~}{t}t_0=C(t)/B(t)`$, the KCS becomes $`𝐗=2c\stackrel{~}{t}_{\stackrel{~}{t}}+\left(2cx+by\right)_x+cy_y+\left(cz1\right)_z`$. Then at $`\stackrel{~}{t}=0`$ we have $`\sigma _{ab}=0`$ and equation (3.2) implies $`c=0`$ i.e. the KCS is reduced to the KVF of type II models.
### 4.4 Type VI<sub>0</sub> models
Let us consider now the Bianchi type VI<sub>0</sub> invariant set which is defined by $`N_1=0`$ and $`N_2N_3<0`$. Propagating equation (3.7) and using (2.9) we obtain:
$`0`$ $`=`$ $`\left(\lambda ^{}3\lambda \right)\mathrm{\Sigma }_{\alpha \beta }+\mathrm{\Lambda }_{\alpha \beta }+\left[2\mathrm{\Sigma }_{22}X^\gamma +\mathrm{\Sigma }_\epsilon ^\gamma X^\epsilon \right]N_2\epsilon _{\gamma 2(\alpha }\delta _{\beta )}^2`$ (4.35)
$`+`$ $`\left[2\mathrm{\Sigma }_{33}X^\gamma +\mathrm{\Sigma }_\epsilon ^\gamma X^\epsilon \right]N_3\epsilon _{\gamma 3(\alpha }\delta _{\beta )}^3+\lambda S_{\alpha \beta }`$
where $`S_{\alpha \beta }`$ and $`\mathrm{\Lambda }_{\alpha \beta }`$ are given in (2.18) and (4.13) respectively.
The $`\alpha \beta `$components of (4.35), after a short calculation, give:
$$\left(\mathrm{\Sigma }_{22}\mathrm{\Sigma }_{33}\right)\left[2_{[3}X_{2]}+X_1\left(N_2+N_3\right)\right]=0$$
(4.36)
$$\left(2\mathrm{\Sigma }_{22}+\mathrm{\Sigma }_{33}\right)\left(2_{[2}X_{1]}+X_3N_2\right)=0=\left(2\mathrm{\Sigma }_{33}+\mathrm{\Sigma }_{22}\right)\left(2_{[3}X_{1]}+X_2N_3\right)$$
(4.37)
which implies that $`\mathrm{\Sigma }_{22}=\mathrm{\Sigma }_{33}`$ since it can be shown that for $`\mathrm{\Sigma }_{22}\mathrm{\Sigma }_{33}`$ the KCS is becoming a proper HVF.
In this case $`N_3=N_2`$ and the system of equations (2.14)-(2.17), (3.4)-(3.7) and (3.15)-(3.16) is identically satisfied, provided that:
$$\stackrel{~}{\alpha }=\lambda \left(12\mathrm{\Sigma }_{22}\right)$$
(4.38)
and the function $`\lambda `$ is given in terms of the shear and spatial curvature variables ($`\mathrm{\Sigma }_{22}0`$):
$$\lambda ^{}=\frac{\lambda \left(2N_2^2+9\mathrm{\Sigma }_{22}\right)}{3\mathrm{\Sigma }_{22}}.$$
(4.39)
Note that for $`\mathrm{\Sigma }_{22}=0`$ the above system of equations implies $`\lambda =0=\alpha =\delta `$ and the KCS reduces to an isometry. Moreover, it is not difficult to show that:
$$\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)|_{\mathrm{\Omega }=0}=\frac{2\lambda \left(\mathrm{\Sigma }_{22}+1\right)}{\mathrm{\Sigma }_{22}}$$
(4.40)
$$\left(\stackrel{~}{\alpha }\stackrel{~}{\delta }\right)|_{\mathrm{\Omega }>0}=\frac{2\lambda \left[N_2^2+3\mathrm{\Sigma }_{22}\left(\mathrm{\Sigma }_{22}+1\right)\right]}{3\mathrm{\Sigma }_{22}}.$$
(4.41)
Due to equations (4.40) and (4.41), the analysis of the asymptotic behaviour of the KCS is straightforward. The past and future attractors of the vacuum invariant subset are the equilibrium points $`\mathrm{\Sigma }_{22}=1,N_2=0`$ (the arc $`𝒦_1`$) and $`\mathrm{\Sigma }_{22}=1,N_2=0`$ (the Taub point $`𝒯_1`$) respectively and the KCS reduces to a proper KSS (with $`3\alpha +\delta =4\psi =0`$) and a proper HVF. Regarding the non-vacuum models it has been shown that the future attractor is the Collins homothetic model $`P_1^+(VI_0)`$ for which $`(\stackrel{~}{\alpha }\stackrel{~}{\delta })|_{P_1^+(VI_0)}=0`$ and the KCS reduces to a proper HVF. The past equilibrium state of the invariant set $`S_1^+(VI_0)`$ is either the LRS Kasner point $`𝒬_1`$ where the KCS becomes a proper KSS or the Friedmann-Lemaître model $``$ and the KCS is given in equation (4.7) with $`\lambda ^{}=3\lambda `$.
These results suggest that, *every type VI*<sub>0</sub> *evolving vacuum or* $`\gamma `$*law perfect fluid model lying in the invariant subset* $`S_1^+(VI_0)`$ *always admits a four dimensional group of KCS that reduces, at the asymptotic regimes, to a self-similarity group of the first or second kind except from a set of measure zero for which the KCS preserves its nature.*
For completeness we also give the form of the KCS for the general vacuum solution satisfying $`N_\alpha ^\alpha =0`$ (an arbitrary multiplication constant has been omitted) :
$$ds^2=t^{1/2}e^{t^2}\left(dt^2+dx^2\right)+t\left(e^{2x}dy^2+e^{2x}dz^2\right)$$
(4.42)
$$𝐗=\frac{4t\left(c_1+c_2\right)}{4t^23}_t+c_1_x+c_2y_y+\left(2c_1+c_2\right)z_z$$
(4.43)
where the symmetry scales are:
$$\alpha =\frac{\left(c_1+c_2\right)\left(4t^21\right)}{\left(4t^23\right)},\delta =\frac{\left(c_1+c_2\right)\left(4t^29\right)\left(4t^2+1\right)}{\left(4t^23\right)^2}.$$
(4.44)
We can verify that, at early and late times, we have $`3\alpha +\delta 0`$ and $`\alpha \delta 0`$, so the expected limiting cases of the KCS is a self-similarity of the second or first kind<sup>2</sup><sup>2</sup>2We recall here that the time coordinate $`t`$ does not represent the proper clock time as $`t0^+`$ or $`t+\mathrm{}`$. However, the freedom of choosing the time gauge, ensures a similar asymptotic behaviour for the symmetry scales.. Also to be noted is the apparent “singular” behaviour (due to the specific time gauge) of the KCS as $`t\sqrt{3}/2`$ ($`\sigma _{ab}0`$). However, similarly to the previous case, we can show that at this value, the temporal component of $`𝐗`$ vanishes and the symmetry equations imply that $`c_1+c_2=0`$ i.e. the KCS $`𝐗_xy_y+z_z`$ reduces to the KVF of Bianchi type VI<sub>0</sub> models.
## 5 Discussion
As we have mentioned, it is customary to study a specific symmetry assumption from a geometrical point of view without taking into consideration the kinematical and dynamical structure of the corresponding model. As a consequence, the effects on the dynamics of the symmetry constraints are hidden and usually produce unphysical results. Up to date, self-similarity of the first kind appears to be the only symmetry condition with a transparent physical nature since it represents *a geometrization of the asymptotic (equilibrium) state of general models*. Therefore, in order to complete the dynamical picture, it was of interest to seek and find a symmetry that could be used effectively as a consistent tool for the invariant description of the intermediate behaviour of general models.
In the present article we have proposed a new technique of studying geometric symmetries by fully exploiting the 1+3 orthonormal frame scheme and the introduction of expansion-normalized variables. The analysis of the complete set of consistency equations (2.14)-(2.17), (3.4)-(3.7) and (3.15)-(3.16) has revealed a *novel feature* of a large class of SH models that is summarized in the following:
Proposition 1 *Evolving SH models of Bianchi type I, LRS models of Bianchi type II and type VI<sub>0</sub> models within the associated invariant subset* $`N_\alpha ^\alpha =0`$*, are geometrically described by the existence of a four dimensional group of KCS that is reduced to a self-similarity group of the first or second kind at the asymptotic regimes, except from a set of measure zero for which the properness of the KCS is preserved.*
However, as we have seen, even in the exceptional cases the corresponding equilibrium state is the Friedmann-Lemaître model $``$ which we have proved that always admits a proper KCS, supplementing the geometric properties of the standard cosmological model.
At first sight, the existence of a proper KCS in SH models appears somewhat surprising, at least from a dynamical point of view. This is mainly because the interaction mechanism between the geometric “assumption” of a KCS and the dynamical behaviour of SH models is not, conceptually, apparent. However a closer look on the structural properties of SH models indicates a possible qualitative interpretation of this interaction. In particular, the full set of non-linear EFE can be seen as the perturbed version of the associated linearization of equations (2.14)-(2.17), at the vicinity of a hyperbolic equilibrium point. Accordingly we may interpret the generator of a KCS as representing the perturbation (to some order) of the corresponding generator of the self-similarity transformation group of the first or second kind. Eventually, this observation will enable us to geometrize the majority of the concepts and techniques that are used in the theory of dynamical systems with a view to optimize and efficiently elaborate the results from the qualitative study of general cosmological models.
Although the existence of a KCS signifies the physical ground for a first promising attempt towards a “geometrisation” of the evolutionary behaviour of SH models, we do not allegate that the KCS (uniquely) characterizes the intermediate epoch of the totality of SH models. In fact, because a KCS possess two arbitrary spatially homogeneous functions (the symmetry scales), we expect that *only* SH models with *two* essential degrees of freedom will exhibit a proper KCS. This conjecture is confirmed by the Proposition 1 in which all the exact solutions found to admit a proper KCS belong to this class. The main reason, that provides an interpretation for the possible connection between exact solutions and the existence of KCS, appears to involve the so-called “hidden” symmetries of the SH cosmologies. Indeed, using the Hamilton-Jacobi reformulation of the EFE in which all the dynamical picture is encoded in one geometric object, it has been shown that all known exact solutions with two degrees of freedom are associated with a specific kind of “hidden” symmetry namely the existence of a Killing tensor symmetry of the Jacobi metric that generalizes the corresponding cyclic variables and the Hamilton-Jacobi separability . Therefore it will be of interest to extent the analysis to the rest of the Bianchi models in order to see if the existence of a KCS is a general feature of the two-dimensional SH invariant subset i.e. if it is related with the above type of “hidden” symmetry of the Bianchi cosmological models .
We should remark that, although we have mainly focused our study to Bianchi types with clear and simple past and future equilibrium states, one could apply the approach presented in this paper to models with more complicated dynamical structure e.g. models with oscillating or diverging asymptotic behaviour near to the past or future attractor. The non-vacuum Bianchi VII<sub>0</sub> invariant subset provides an interesting example since it is well known that the associated kinematical variables are unbounded and do not approach any equilibrium point into the future i.e. non-vacuum Bianchi type VII<sub>0</sub> models are not asymptotically self-similar . As a consequence one should expect that a KCS does not exist in those models due to the non-existence of a self-similar model as future attractor. However, a preliminary analysis has shown that a KCS does exist in LRS type VII<sub>0</sub> models which is never reduced to a HVF or a KSS , suggesting that the concept of the KCS represents not simply a perturbed version of the self-similarity group but a generic property (in the spirit of ) of the two-dimensional invariant subset of the SH cosmological models.
Clearly, the case of higher dimensional invariant subsets requires further investigation. Assuming the existence of several proper KCS, will not solve the problem for models with three or more essential degrees of freedom, since the condition (3.1) implies that the symmetry scales are always spatially homogeneous restricting the dimension of the Lie algebra of KCS in SH models to four. Therefore the question of determining the symmetry which invariantly describes the whole set of SH models is still open. Nevertheless, the implications of the above results, enforce the important role which may play a specific (still unknown) general symmetry as an effective geometric implement for the invariant description of general cosmological models and not only as a simplification rule towards the determination of exact solutions with ambiguous (or even without any) physical meaning. A closely related issue is how the constraints, coming from the presence of the general symmetry, could reveal a path of constructing the general (whenever is possible) solution of the corresponding cosmological model. We expect that, the approach of studying generic geometric symmetries in spacetime developed in this article, can be applied to more general geometric setups leading to a more efficient qualitative and analytical study of general vacuum and perfect fluid models.
Acknowledgments
The author would like to thank the anonymous referees for suggesting changes that have improved the presentation of the manuscript. This work is supported through the research programme “Pythagoras” of the Greek Ministry of National Education, contract No 70-03-7310. Also the author gratefully acknowledges the award of a postdoctoral fellowship from the Spanish “Ministerio de Educación y Ciencia”, grant No SB2004-0110.
|
warning/0506/math-ph0506025.html
|
ar5iv
|
text
|
# Poisson involutions, Spin Calogero-Moser systems associated with symmetric Lie subalgebras and the symmetric space spin Ruijsenaars-Schneider models
## 1. Introduction
In the last few years, a groupoid-theoretic scheme based on the coboundary dynamical Poisson groupoids and their corresponding Lie bialgebroids was introduced in the study of certain integrable Hamiltonian systems and their solutions \[LX\],\[L1\],\[L2\],\[L3\]. As is well-known, these geometric objects are naturally associated with so-called classical dynamical r-matrices \[EV\],\[BKS\] which first appeared in the context of Wess-Zumino-Witten (WZW) conformal field theory \[BDF\],\[F\].
In this paper, we shall continue to use these geometric objects to unify the study of a variety of Hamiltonian systems known under the general name of spin Calogero-Moser systems and spin Ruijsenaars-Schneider models. More specifically, we shall consider in this work examples of such systems which turn out to be realizable in the stable loci of the geometric objects mentioned above under Poisson involutions. As we know through the work in \[X\], the stable locus of a Poisson involution is an example of a class of submanifolds with induced Poisson structures which the author in \[X\] called Dirac submanifolds.(The fact that the stable locus of a Poisson involution carries a natural induced Poisson structure was also noted by the authors in \[FV\].) Indeed, as we shall explain below, it is advantageous to formulate several of our results in this broader framework.
We now give an outline of our approach. As starting point, we consider certain spin Calogero-Moser systems (resp. spin Ruijsenaars-Schneider models) which can be realized in the dual bundles of symmetric coboundary dynamical Lie algebroids (resp. symmetric coboundary dynamical Poisson groupoids). For these Hamiltonian systems, the underlying Poisson manifolds as well as their realization spaces are both Hamiltonian $`H`$-spaces which carry natural Poisson involutions. The construction of the integrable systems of interests then proceeds in two stages. In the first stage, we apply Dirac reduction (which will be developed here) to reduce the initial realization maps to ones between the stable loci of Poisson involutions. In this way, we obtain the realization of the Dirac reduction of the afore-mentioned systems. In general, these reduced systems are not integrable systems (see Theorem 3.15, Theorem 3.18 and Section 4.2 for exceptions). However, as it turns out, the stable loci are Hamiltonian $`D`$-spaces for some subgroup $`D`$ of $`H`$. Moreover, the natural invariant functions Poisson commute on certain fibers of the equivariant momentum maps. Consequently, we can apply Poisson reduction (and this is the second stage) to obtain the associated integrable systems.
There are several motivations for this work. One of these has come from the desire to understand the Hamiltonian formulation as well as the integrability of the equations of motion which arise from the so-called level dynamics approach in random matrix theory \[Y\],\[HKS\],\[NM\],\[GRMN\]. This connection is reflected in our choice of examples in Section 4. On the other hand, there is a well-known correspondence between the $`N`$-soliton solutions of the $`A_n^{(1)}`$ affine Toda field theory and some spin-generalized Ruijensaars-Schneider equations \[BH\]. (Some of the variables in these equations actually depend on the choice of eigenvectors of a certain skew-Hermitian matrix $`V`$.) However, the Hamiltonian formulation of these equations has remained open. We shall give a solution to this problem in Section 5 below. As the reader will see, these equations are related to a symmetric coboundary dynamical Poisson groupoid $`(\mathrm{\Gamma },\mathrm{\Sigma })`$, where $`\mathrm{\Gamma }`$ is associated to a hyperbolic dynamical r-matrix, and $`\mathrm{\Sigma }`$ is a Poisson involution on $`\mathrm{\Gamma }.`$ More precisely, they can be obtained from a Hamiltonian system on the stable locus $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ of $`\mathrm{\Sigma }`$ by restricting the equations of motion to an appropriate fiber of the momentum map. Consequently, the system which is invariant under the gauge freedom in picking the eigenvectors of $`V`$ is an integrable Hamiltonian system on a Poisson reduction of $`\mathrm{\Gamma }^\mathrm{\Sigma }.`$ Finally we remark that in the process of assembling the necessary machinery in order to tackle the above problems, we will give the explicit expression for the Poisson structure on the stable locus of a Poisson involution on a coboundary dynamical Lie algebroid (resp. coboundary dynamical Poisson groupoid). Thus this answers a question raised in \[X\].
The paper is organized as follows. In section 2, we begin by recalling some basic facts about coboundary dynamical Lie algebroids and coboundary dynamical Poisson groupoids which will be used throughout the paper. In particular, we will discuss a subclass of such Lie algebroids defined by so-called classical dynamical r-matrices with spectral parameter. We will also recall what we mean by spin Calogero-Moser systems associated with this subclass of coboundary dynamical Lie algebroids. In section 3, the main goal is to develop a general scheme of constructing integrable systems based on realization in symmetric coboundary dynamical Poisson groupoids and the dual bundles of symmetric coboundary dynamical Lie algebroids. As we already mentioned above in the context of specific examples, the construction proceeds in two stages. For Dirac reduction, our main tool comes from an elementary result which shows how to reduce a Poisson map between two Poisson manifolds to one between their respective Dirac submanifolds (Theorem 3.2 and Corollary 3.5). From this, we also obtain a condition under which a Dirac submanifold $`Q`$ of a Hamiltonian $`G`$-space $`P`$ is Hamiltonian $`H`$-space for some Lie subgroup $`H`$ of $`G`$ (Proposition 3.6 and Corollary 3.7). There are two reasons for formulating our results in terms of Dirac submanifolds. First, the notion offers a better conceptual framework. Secondly, when formulated in this broader framework, the results are also applicable to the cosymplectic submanifolds \[W1\] (when $`P`$ is symplectic, the cosymplectic submanifolds of $`P`$ are precisely its symplectic submanifolds). In the special case when the Dirac submanifold is given by the stable locus of a Poisson involution on the dual bundle of a coboundary dynamical Lie algebroid (resp. coboundary dynamical Poisson groupoid), we also derive the intrinsic expression for the induced Poisson structure which is essential for our purpose here. In Section 4, we introduce several examples of spin Calogero-Moser systems associated with real symmetric Lie algebras. Then we show how the reduction procedure developed in Section 3 can be carried out to obtain the associated integrable systems of interests. In the special case when the Lie algebra $`𝔤`$ is $`gl(N,)`$, we also provide a sketch of the Liouville integrability of the associated integrable models. Note that our goal of this section is suggestive rather than exhaustive in the sense that we have made no attempt to give a classification of systems which can be treated by our method. Finally, in Section 5, we consider the spin Ruijsenaars-Schneider models associated with a symmetric coboundary dynamical Poisson groupoid $`(\mathrm{\Gamma },\mathrm{\Sigma })`$. In this case, the realization map is just the identity map and it is easy to show how the scheme in Section 3 can be implemented. As we mentioned earlier, our goal here is explain how the spin-generalized Ruijsenaars-Schneider equations in \[BH\] are obtained from an invariant Hamiltonian system on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ which is a special case of what we call symmetric space Ruijsenaars-Schneider models here.
To close, we remark that a factorization theory also exists for the solution of the Hamiltonian systems treated here (provided the classical dynamical r-matrix satisfies the modified dynamical Yang-Baxter equation), as is clear from assumptions A5 and G5 in Section 3 and the development in \[L1\],\[L2\]. For this reason, we do not give any details here.
Acknowledgments. The author would like to thank Ping Xu for the reference \[FV\] when this work was in its final stage of preparation.
## 2. Preliminaries
The purpose of this section is to recall some basic results about coboundary dynamical Lie algebroids and coboundary dynamical Poisson groupoids. For our applications in this work, we will pay special attention to a subclass of such Lie algebroids which are associated with so-called classical dynamical r-matrices with spectral parameter. We will also recall what we mean by spin Calogero-Moser systems associated with this subclass of coboundary dynamical Lie algebroids.
Let $`G`$ be a connected Lie group, and $`HG`$ a connected Lie subgroup. We shall denote by $`𝔤`$ and $`𝔥`$ the corresponding Lie algebras and let $`\iota :𝔥𝔤`$ be the Lie inclusion. In what follows, the Lie groups and Lie algebras can be real or complex unless we specify otherwise.
We begin by recalling a fundamental construction in \[EV\] which gives a geometric interpretation of dynamical r-matrices in terms of Poisson groupoids. We shall, however, follow the formulation in \[L1\] and in particular we shall give the explicit expression for the Poisson structure which is essential for our purpose here. Let $`U𝔥^{}`$ be a connected $`Ad_H^{}`$-invariant open subset, we say that a smooth (resp. holomorphic) map $`R:UL(𝔤^{},𝔤)`$ (here and henceforth we denote by $`L(𝔤^{},𝔤)`$ the set of linear maps from $`𝔤^{}`$ to $`𝔤`$) is a classical dynamical r-matrix if and only if it is pointwise skew-symmetric:
$$<R(q)(A),B>=<A,R(q)B>$$
$`(2.1)`$
and satisfies the classical dynamical Yang-Baxter condition
$$\begin{array}{cc}& [R(q)A,R(q)B]+R(q)(ad_{R(q)A}^{}Bad_{R(q)B}^{}A)\hfill \\ \hfill +& dR(q)\iota ^{}A(B)dR(q)\iota ^{}B(A)+<dR(q)()(A),B>=\chi (A,B),\hfill \end{array}$$
$`(2.2)`$
where $`<dR(q)()(A),B>`$ is the element in $`𝔥`$ whose pairing with $`\lambda 𝔥^{}`$ is given by $`<dR(q)(\lambda )(A),B>`$ and $`\chi :𝔤^{}\times 𝔤^{}𝔤`$ is $`G`$-equivariant, that is,
$$\chi (Ad_{g^1}^{}A,Ad_{g^1}^{}B)=Ad_g\chi (A,B)$$
$`(2.3)`$
for all $`A,B𝔤^{}`$, $`gG`$, and all $`qU`$. The dynamical $`r`$-matrix is said to be $`H`$-equivariant if and only if
$$R(Ad_{h^1}^{}q)=Ad_hR(q)Ad_h^{}$$
$`(2.4)`$
for all $`hH,qU`$.
We shall equip $`\mathrm{\Gamma }=U\times G\times U`$ with the trivial Lie groupoid structure over $`U`$ \[M\] with target and source maps
$$\alpha (u,g,v)=u,\beta (u,g,v)=v$$
$`(2.5)`$
and multiplication map
$$m((u,g,v),(v,g^{},w))=(u,gg^{},w).$$
$`(2.6)`$
For a smooth (resp. holomorphic) function $`\phi `$ on $`\mathrm{\Gamma }`$, we define its partial derivatives and its left and right gradients (with respect to $`G`$) by
$$\begin{array}{cc}& <\delta _1\phi ,u^{}>=\frac{d}{dt}|_{t=0}\phi (u+tu^{},g,v),<\delta _2\phi ,v^{}>=\frac{d}{dt}|_{t=0}\phi (u,g,v+tv^{}),u^{},v^{}𝔥^{}\hfill \\ & <D\phi ,X>=\frac{d}{dt}|_{t=0}\phi (u,e^{tX}g,v),<D^{}\phi ,X>=\frac{d}{dt}|_{t=0}\phi (u,ge^{tX},v),X𝔤.\hfill \end{array}$$
###### Theorem 2.1
(a) The bracket
$$\begin{array}{cc}\hfill \{\phi ,\psi \}_R(u,g,v)=& <u,[\delta _1\phi ,\delta _1\psi ]><v,[\delta _2\phi ,\delta _2\psi ]>\hfill \\ & <\iota \delta _1\phi ,D\psi ><\iota \delta _2\phi ,D^{}\psi >\hfill \\ & +<\iota \delta _1\psi ,D\phi >+<\iota \delta _2\psi ,D^{}\phi >\hfill \\ & +<R(v)D^{}\phi ,D^{}\psi ><R(u)D\phi ,D\psi >\hfill \end{array}$$
$`(2.7)`$
defines a Poisson structure on $`\mathrm{\Gamma }`$ if and only if $`R:UL(𝔤^{},𝔤)`$ is an $`H`$\- equivariant classical dynamical $`r`$-matrix. (b) The trivial Lie groupoid $`\mathrm{\Gamma }`$ equipped with the Poisson bracket $`\{,\}_R`$ is a Poisson groupoid. Moreover, it is a Hamiltonian $`H`$-space under the natural left and right $`H`$-actions with equivariant momentum maps given by $`\alpha `$ and $`\beta `$ respectively.
We shall call the pair $`(\mathrm{\Gamma },\{,\}_R)`$ the coboundary dynamical Poisson groupoid associated to $`R`$. Note that the explicit expression for $`\{\phi ,\psi \}_R`$ in Theorem 2.1(a) above can be derived from the characterizing properties in \[EV\] and the corresponding expression for the general dynamical case can be found in \[LP\].
Let $`A\mathrm{\Gamma }=_{qU}\{0_q\}\times 𝔤\times 𝔥^{}TU\times 𝔤`$ be the Lie algebroid of the trivial Lie groupoid $`\mathrm{\Gamma }`$. (See \[CdSW\],\[M\] for details.) Then associated with $`(\mathrm{\Gamma },\{,\}_R)`$ is a Lie algebroid structure on the dual bundle $`A^{}\mathrm{\Gamma }=_{qU}\{0_q\}\times 𝔤^{}\times 𝔥T^{}U\times 𝔤^{}`$ as a consequence of Weinstein’s coisotropic calculus \[W2\].(See \[LP\] for a more general discussion and \[BKS\] for a different approach.) The anchor map $`a_{}:A^{}\mathrm{\Gamma }TU`$ of $`A^{}\mathrm{\Gamma }`$ of this Lie algebroid is given by
$$a_{}(0_q,A,Z)=(q,\iota ^{}Aad_Z^{}q)$$
$`(2.8)`$
while the bracket $`[,]_{A^{}\mathrm{\Gamma }}`$ on $`Sect(U,A^{}\mathrm{\Gamma })`$ has the following form \[BKS\],\[L2\]:
$$\begin{array}{cc}& [(0,\xi ,Z),(0,\xi ^{},Z^{})]_{A^{}\mathrm{\Gamma }}(q)\hfill \\ \hfill =& (0_q,d\xi ^{}(q)(\iota ^{}\xi (q)ad_{Z(q)}^{}q)d\xi (q)(\iota ^{}\xi ^{}(q)ad_{Z^{}(q)}^{}q)\hfill \\ & ad_{R(q)\xi (q)Z(q)}^{}\xi ^{}(q)+ad_{R(q)\xi ^{}(q)Z^{}(q)}^{}\xi (q),\hfill \\ & dZ^{}(q)(\iota ^{}\xi (q)ad_{Z(q)}^{}q)dZ(q)(\iota ^{}\xi ^{}(q)ad_{Z^{}(q)}^{}q)\hfill \\ & [Z,Z^{}](q)+<dR(q)()\xi (q),\xi ^{}(q)>)\hfill \end{array}$$
$`(2.9)`$
where $`\xi ,\xi ^{}:U𝔤^{}`$, $`Z,Z^{}:U𝔥`$ are smooth (resp. holomorphic) maps and $`<dR(q)()\xi (q),\xi ^{}(q)>`$ is the element in $`𝔥`$ whose pairing with $`\lambda 𝔥^{}`$ is $`<dR(q)(\lambda )\xi (q),\xi ^{}(q)>.`$ We shall call $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ the coboundary dynamical Lie algebroid associated to $`R.`$
Now, for any Lie algebroid $`(A,[,]_A,a_A)`$ over a smooth manifold $`M`$, recall that there exists a Lie-Poisson structure on the dual bundle $`A^{}`$ \[CDW\] which is uniquely determined by the property
$$\{l_X,l_Y\}=l_{[X,Y]_A}$$
$`(2.10)`$
where for $`X,YSect(M,A)`$, $`l_X`$ and $`l_Y`$ are the corresponding linear functions on $`A^{}`$. The following result was obtained in \[L2\].
###### Theorem 2.2
(a) The Lie-Poisson structure on the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ is given by
$$\begin{array}{cc}& \{\phi ,\psi \}_{A\mathrm{\Gamma }}(q,\lambda ,X)\hfill \\ \hfill =& <\lambda ,[\delta _2\phi ,\delta _2\psi ]>+<dR(q)(\lambda )\delta \phi ,\delta \psi >\hfill \\ & +<X,ad_{R(q)\delta \phi \delta _2\phi }^{}\delta \psi +ad_{R(q)\delta \psi \delta _2\psi }^{}\delta \phi >\hfill \\ & <q,[\delta _2\phi ,\delta _1\psi ]+[\delta _1\phi ,\delta _2\psi ]>+<\delta _1\psi ,\iota ^{}\delta \phi ><\delta _1\phi ,\iota ^{}\delta \psi >.\hfill \end{array}$$
$`(2.11)`$
(b) With the action of $`H`$ on $`A\mathrm{\Gamma }`$ defined by the formula
$$h(q,\lambda ,X)=(Ad_{h^1}^{}q,Ad_{h^1}^{}\lambda ,Ad_hX),$$
$`(2.12)`$
the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ equipped with the Lie-Poisson structure is a Hamiltonian $`H`$-space with equivariant momentum map
$$\gamma :A\mathrm{\Gamma }𝔥^{},(q,\lambda ,X)\lambda .$$
$`(2.13)`$
In the rest of the section, we shall assume that $`𝔤`$ is a Lie algebra with a nondegenerate invariant pairing $`(,)`$ and $`𝔥𝔤`$ is a non-degenerate (i.e. $`(,)_{𝔥\times 𝔥}`$ is nondegenerate) abelian Lie subalgebra. Then we can make the identifications $`𝔤^{}𝔤`$, $`𝔥^{}𝔥`$, $`ad^{}ad`$, $`\iota ^{}\mathrm{\Pi }_𝔥`$, where $`\mathrm{\Pi }_𝔥`$ is the projection map to $`𝔥`$ relative to the direct sum decomposition $`𝔤=𝔥𝔥^{}`$. Hence we can regard $`R(q)`$ as taking values in $`End(𝔤)`$. In this case, an important sufficient condition for an $`H`$-equivariant map $`R`$ to define a coboundary dynamical Poisson groupoid is given by the modified dynamical Yang-Baxter equation (mDYBE):
$$\begin{array}{cc}& [R(q)X,R(q)Y]R(q)([R(q)X,Y]+[X,R(q)Y])\hfill \\ \hfill +& dR(q)\mathrm{\Pi }_𝔥X(Y)dR(q)\mathrm{\Pi }_𝔥Y(X)+(dR(q)()X,Y)\hfill \\ \hfill =& c^2[X,Y],\hfill \end{array}$$
$`(2.14)`$
where $`c`$ is a nonzero constant.
We now turn our attention to a subclass of $`(\mathrm{\Gamma },\{,\}_R)`$ and $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ which are associated with so-called classical dynamical r-matrices with spectral parameter. In the following, we shall assume that $`𝔤`$ is a Lie algebra over $``$.
###### Definition Definition 2.3 \[EV\],\[LX\]
A classical dynamical r-matrix with spectral parameter is a meromorphic map $`r:𝔥\times 𝔤𝔤`$ with a simple pole at $`z=0`$ satisfying the following conditions for all $`(q,z)𝔥\times `$ away from the poles of $`r`$, 1. the zero weight condition:
$$[h1+1h,r(q,z)]=0,$$
$`(2.15)`$
for all $`h𝔥`$, 2. the generalized unitarity condition:
$$r^{12}(q,z)+r^{21}(q,z)=0,$$
$`(2.16)`$
3. the residue condition:
$$Res_{z=0}r(q,z)=\mathrm{\Omega },$$
$`(2.17)`$
where $`\mathrm{\Omega }(S^2𝔤)^𝔤`$ is the Casimir element corresponding to $`(,)`$, 4. the classical dynamical Yang-Baxter equation (CDYBE) with spectral parameter:
$$\begin{array}{cc}\hfill Alt(dr)& +[r^{12}(q,z_{12}),r^{13}(q,z_{13})+r^{23}(q,z_{23})]\hfill \\ & +[r^{13}(q,z_{13}),r^{23}(q,z_{23})]=0,\hfill \end{array}$$
$`(2.19)`$
where $`z_{ij}=z_iz_j`$.
Let $`L𝔤`$ be the loop algebra consisting of Laurent series with coefficients in $`𝔤.`$ If $`r`$ is a classical dynamical r-matrix with spectral parameter, we define
$$(R(q)X)(z)=p.v.\frac{1}{2\pi i}_C(r(q,wz),X(w)1)𝑑w,XL𝔤,$$
$`(2.20)`$
where $`C`$ is a small circle centered at $`0`$ with positive orientation, and $`p.v.`$ denotes the principal value of the improper integral. We have the following result \[LX\].
###### Theorem 2.4
(a) $`R`$ is an $`H`$-equivariant classical dynamical r-matrix which satisfies the mDYBE with $`c=\frac{1}{4}.`$ (b) For $`XL𝔤`$, we have the formula
$$(R(q)X)(z)=\frac{1}{2}X(z)+\underset{k0}{}\frac{1}{k!}(\frac{^kr}{z^k}(q,z),X_{(k+1)}1).$$
We now fix an open connected set $`U𝔥`$ on which $`R`$ is holomorphic. Let $`A\mathrm{\Omega }=U\times 𝔥\times 𝔤`$ be the trivial Lie algebroid over $`U`$ with vertex algebra $`𝔤`$. We shall identify its dual bundle $`A^{}\mathrm{\Omega }`$ with $`U\times 𝔥\times 𝔤`$ and equip it with the Lie-Poisson structure. On the other hand, we can use $`R`$ in Theorem 2.4 above to construct the associated coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }U\times 𝔥\times L𝔤`$. Therefore, we can equip its dual bundle $`A\mathrm{\Gamma }`$ with the corresponding Lie-Poisson structure. For each $`qU`$, we now define a map $`r_{}^\mathrm{\#}(q):𝔤L𝔤`$ by the formula
$$((r_{}^\mathrm{\#}(q)\xi )(z),\eta )=(r(q,z),\eta \xi )$$
$`(2.21)`$
where $`\xi `$, $`\eta 𝔤.`$
###### Theorem 2.5 \c{LX}
The map $`\rho :A^{}\mathrm{\Omega }A\mathrm{\Gamma }`$ given by
$$(q,p,\xi )(q,\mathrm{\Pi }_𝔥\xi ,p+r_{}^\mathrm{\#}(q)\xi )$$
is an $`H`$-equivariant Poisson map, where $`H`$ acts on $`A^{}\mathrm{\Omega }`$ by $`h(q,p,\xi )=(q,p,Ad_h\xi )`$.
###### Definition Definition 2.6
Let $`r`$ be a classical dynamical r-matrix with spectral parameter and let $`L=Pr_3\rho `$, where $`\rho `$ is the realization map in Theorem 2.5. Then the (complex holomorphic) Hamiltonian system on $`A^{}\mathrm{\Omega }`$ generated by the Hamiltonian function
$$^{}(q,p,\xi )=\frac{1}{2}_C(L(q,p,\xi ),L(q,p,\xi ))\frac{dz}{2\pi iz}$$
$`(2.22)`$
is called the spin Calogero-Moser system associated with $`r`$. Here, $`C`$ is a small circle centered at $`0`$ with the positive orientation.
Remark 2.7. Actually we will use the real version of Theorem 2.5 in our application in Section 4 below.
## 3. Dirac reduction of Poisson maps and geometric construction of inte- fak grable systems via successive reductions
The goal of this section is to develop a general scheme of constructing integrable systems based on realization in symmetric coboudary dynamical Lie algebroids and symmetric coboundary dynamical Poisson groupoids. In order to do this, we have to consider the method of Dirac reduction. We begin by recalling the notion of a Dirac submanifold as recently introduced in \[X\]. It is a generalization of the notion of cosymplectic submanifolds of Weinstein \[W1\]. For convenience, we shall formulate our results in this section in the differentiable category, but it will be clear that the results are also valid for the holomorphic category.
###### Definition Definition 3.1
Let $`(P,\pi )`$ be a Poisson manifold. A submanifold $`Q`$ of $`P`$ is a Dirac submanifold iff there exists a Whitney sum decomposition
$$T_QP=TQV_Q$$
$`(3.1)`$
where $`V_Q^{}`$ is a Lie subalgebroid of the cotangent Lie algebroid $`T^{}P`$.
If $`Q`$ is a Dirac submanifold of $`(P,\pi )`$, then necessarily $`Q`$ carries a natural Poisson structure $`\pi _Q`$ whose symplectic leaves are given by the intersection of $`Q`$ with the symplectic leaves of $`P`$. Indeed, $`\pi _Q^{\mathrm{}}:T^{}QTQ`$ is just the anchor map of the Lie subalgebroid $`T^{}QV_Q^{}`$ of $`T^{}P`$. Moreover, from the knowledge of the injective Lie algebroid morphism $`T^{}QT^{}P`$, it is easy to show that
$$\pi _Q^{\mathrm{}}=pr\pi ^{\mathrm{}}_Qpr^{}$$
$`(3.2a)`$
where $`pr:T_QPTQ`$ is the projection map induced by the decomposition in (3.1) and $`pr^{}`$ is its dual. Alternatively, we have
$$\pi _Q=\pi _Q+\pi ^{}$$
$`(3.2b)`$
where $`\pi ^{}Sect(^2V_Q).`$
We shall call $`Q`$ equipped with the induced Poisson structure a Dirac reduction of $`P.`$ Remark 3.2. An important class of Dirac submanifolds is given by the cosymplectic submanifolds of Weinstein \[W1\], in which case $`V_Q=\pi ^{\mathrm{}}_Q((TQ)^{}).`$ Note that when $`P`$ is symplectic, the cosymplectic submanifolds of $`P`$ are precisely its symplectic submanifolds. We shall give another important class of examples in Proposition 3.5 below. Since we will be dealing with realization maps into the dual bundles of symmetric coboundary dynamical Lie algebroids (resp. symmetric coboundary dynamical Poisson groupoids), the following result is fundamental in reducing such maps.
###### Theorem 3.3
Let $`\varphi :P_1P_2`$ be a Poisson map and let $`Q_1P_1`$, $`Q_2P_2`$ be Dirac submanifolds with respective Whitney sum decompositions
$$T_{Q_1}P_1=TQ_1V_{Q_1},T_{Q_2}P_2=TQ_2V_{Q_2}.$$
Then under the assumptions that
(i) $`\varphi (Q_1)Q_2`$, (ii) $`T_x\varphi (V_{Q_1})_x(V_{Q_2})_{\varphi (x)},xQ_1`$, the map $`\varphi Q_1:Q_1Q_2`$ is a Poisson map, when $`Q_1`$ and $`Q_2`$ are equipped with the induced Poisson structures.
###### Demonstration Proof
Let $`\pi _{P_i}^{\mathrm{}}`$, $`\pi _{Q_i}^{\mathrm{}}`$ be the bundle maps associated with the Poisson structures on $`P_i`$, $`Q_i`$, $`i=1,2`$. Since $`Q_1`$ is a Dirac submanifold of $`P_1`$, we have
$$\pi _{P_1}^{\mathrm{}}(x)(\alpha )=\pi _{Q_1}^{\mathrm{}}(x)(\alpha T_xQ_1)+\stackrel{~}{\pi }_{Q_1}^{\mathrm{}}(x)(\alpha (V_{Q_1})_x)()$$
for all $`xQ_1`$, $`\alpha T_x^{}P_1`$, where $`\stackrel{~}{\pi }_{Q_1}^{\mathrm{}}:V_{Q_1}^{}V_{Q_1}.`$ Similarly,
$$\pi _{P_2}^{\mathrm{}}(y)(\beta )=\pi _{Q_2}^{\mathrm{}}(y)(\beta T_yQ_2)+\stackrel{~}{\pi }_{Q_2}^{\mathrm{}}(y)(\beta (V_{Q_2})_y)()$$
for all $`yQ_2`$, $`\beta T_y^{}P_2`$, where $`\stackrel{~}{\pi }_{Q_2}^{\mathrm{}}:V_{Q_2}^{}V_{Q_2}.`$ Now, it follows from (\*) that
$`T_x\varphi \pi _{P_1}^{\mathrm{}}(x)T_x^{}\varphi (\beta )`$
$`=`$ $`T_x\varphi \pi _{Q_1}^{\mathrm{}}(x)(T_x^{}\varphi (\beta )T_xQ_1)+T_x\varphi \stackrel{~}{\pi }_{Q_1}^{\mathrm{}}(x)(T_x^{}\varphi (\beta )(V_{Q_1})_x)`$
for all $`xQ_1`$, $`\beta T_{\varphi (x)}^{}P_2`$. From assumption (i), we have $`T_x\varphi \pi _{Q_1}^{\mathrm{}}(x)(T_x^{}\varphi (\beta )T_xQ_1)T_{\varphi (x)}Q_2.`$ On the other hand, assumption (ii) implies that $`T_x\varphi \stackrel{~}{\pi }_{Q_1}^{\mathrm{}}(x)(T_x^{}\varphi (\beta )(V_{Q_1})_x)(V_{Q_2})_{\varphi (x)}.`$ Since $`\varphi `$ is Poisson, it follows from (\**) above that we also have
$`T_x\varphi \pi _{P_1}^{\mathrm{}}(x)T_x^{}\varphi (\beta )`$
$`=`$ $`\pi _{Q_2}^{\mathrm{}}(\varphi (x))(\beta T_{\varphi (x)}Q_2)+\stackrel{~}{\pi }_{Q_2}^{\mathrm{}}(\varphi (x))(\beta (V_{Q_2})_{\varphi (x)})`$
for all $`xQ_1`$, $`\beta T_{\varphi (x)}^{}P_2.`$ Therefore, upon equating the two expressions for $`T_x\varphi \pi _{P_1}^{\mathrm{}}(x)T_x^{}\varphi (\beta )`$, we obtain
$`T_x\varphi \pi _{Q_1}^{\mathrm{}}(x)(T_x^{}\varphi (\beta )T_xQ_1)`$
$`=`$ $`\pi _{Q_2}^{\mathrm{}}(\varphi (x))(\beta T_{\varphi (x)}Q_2)`$
which shows that $`\varphi Q_1:Q_1Q_2`$ is Poisson, as desired. $`\mathrm{}`$
###### Definition Definition 3.4
The map $`\varphi Q_1:Q_1Q_2`$ in the theorem above will be called a Dirac reduction of the Poisson map $`\varphi :P_1P_2`$.
The following result gives an important class of Dirac submanifolds which plays a key role in this work.
###### Proposition 3.5 \c{X}
Let $`\sigma :PP`$ be a Poisson involution, i.e., an involution which is also a Poisson map. Then its stable locus $`Q`$ is a Dirac submanifold of $`P`$ with $`V_Q=_{xQ}ker(T_x\sigma +1).`$
As a consequence of this result, the stable locus of a Poisson involution carries a natural Poisson structure. This fact was also noted in \[FP\] and was implicit in the earlier work of several authors \[Bon\], \[Boal\]. (See also p.194 of \[RSTS\].)
In the special case when the Dirac submanifolds in Theorem 3.3 are the stable loci of Poisson involutions, we have the following result.
###### Corollary 3.6
Let $`\sigma _1:P_1P_1`$, $`\sigma _2:P_2P_2`$ be Poisson involutions with stable loci given by $`Q_1`$ and $`Q_2`$ respectively. If $`\varphi :P_1P_2`$ is a Poisson map which commutes with $`\sigma _1`$, $`\sigma _2`$,i.e. $`\sigma _2\varphi =\varphi \sigma _1`$, then $`\varphi Q_1:Q_1Q_2`$ is a Poisson map, when $`Q_1`$ and $`Q_2`$ are equipped with the induced structures.
###### Demonstration Proof
Under the assumption that $`\varphi `$ commutes with the Poisson involutions, it is easy the check that the conditions in Theorem 3.3 are satisfied with
$$V_{Q_1}=\underset{xQ_1}{}ker(T_x\sigma _1+1),V_{Q_2}=\underset{yQ_2}{}ker(T_y\sigma _2+1).$$
Hence the assertion follows. $`\mathrm{}`$
We next consider the problem of reducing a Hamiltonian $`G`$-space $`P`$ to a Dirac submanifold $`Q`$. Note that $`Q`$ is in general not a $`G`$-space. So a natural question is: under what condition is $`Q`$ a Hamiltonian $`H`$-space for some Lie subgroup $`HG`$ ?
###### Proposition 3.7
Let $`\mathrm{\Phi }:G\times PP`$ be a Hamiltonian group action of $`G`$ on the Poisson manifold $`(P,\pi )`$, and let $`Q`$ be a Dirac submanifold of $`P`$ with Whitney sum decomposition $`T_QP=TQV_Q`$. If $`H`$ is a Lie subgroup of $`G`$ with $`Lie(H)=𝔥`$ and if the action $`\mathrm{\Phi }`$ induces an action of $`H`$ on $`Q`$ satisfying
$$T_x\mathrm{\Phi }_h(V_x)V_{\mathrm{\Phi }_h(x)},hH,xQ,$$
where $`V_Q=_{qQ}V_q,`$ then the $`H`$-action on $`Q`$ is also a Hamiltonian group action. Moreover, if $`J:P𝔤^{}`$ is a $`G`$-equivariant momentum map for $`\mathrm{\Phi }`$, then the map $`J_Q=i^{}(JQ):Q𝔥^{}`$ is a $`H`$-equivariant momentum map for the $`H`$-action on $`Q`$. Here, $`i^{}`$ is the dual map of the Lie inclusion $`i:𝔥𝔤.`$
###### Demonstration Proof
For $`hH`$, the assertion that $`\mathrm{\Phi }_hQ:QQ`$ is Poisson is a consequence of Theorem 3.3. To show that $`J_Q`$ is an equivariant momentum map for the $`H`$ action on $`Q`$, note that
$$\frac{d}{dt}|_{t=0}\mathrm{\Phi }_{e^{tZ}}(x)=\pi _Q^{\mathrm{}}(x)(d\widehat{J}(Z)(x)T_xQ)$$
for all $`Z𝔥`$ and $`xQ`$. But it is easy to check that the map $`\widehat{J}_Q:𝔥C^{\mathrm{}}(Q)`$ defined by $`\widehat{J}_Q(Z)(x)=<J_Q(x),Z>`$, $`Z𝔥`$, $`xQ`$ satisfies $`d\widehat{J}_Q(x)=d\widehat{J}(Z)(x)T_xQ`$. This completes the proof. $`\mathrm{}`$
###### Corollary 3.8
Let $`\mathrm{\Phi }:G\times PP`$ be a Hamiltonian group action of $`G`$ on the Poisson manifold $`(P,\pi )`$ and suppose $`\sigma :PP`$ is a Poisson involution. If $`H`$ is a Lie subgroup of $`G`$ such that
$$\mathrm{\Phi }_h\sigma =\sigma \mathrm{\Phi }_h,hH,$$
then $`\mathrm{\Phi }`$ induces a Hamiltonian group action of $`H`$ on the stable locus $`P^\sigma `$. Moreover, if $`J:P𝔤^{}`$ is a $`G`$-equivariant momentum map for $`\mathrm{\Phi }`$, then the map $`\stackrel{~}{J}=i^{}(JP^\sigma ):P^\sigma 𝔥^{}`$ is a $`H`$-equivariant momentum map for the $`H`$-action on $`P^\sigma .`$
###### Demonstration Proof
It follows from the assumption $`\mathrm{\Phi }_h\sigma =\sigma \mathrm{\Phi }_h,hH`$ that $`\mathrm{\Phi }`$ induces an action of $`H`$ on $`Q=P^\sigma `$. Let $`V_Q=_{xQ}ker(T_x\sigma +1)`$, we want to show that $`T_x\mathrm{\Phi }_h(V_x)V_{\mathrm{\Phi }_h(x)},hH,xQ.`$ For this purpose, take any $`vV_x.`$ Then we have
$`(T_{\mathrm{\Phi }_h(x)}\sigma +1)T_x\mathrm{\Phi }_h(v)`$
$`=`$ $`T_x(\mathrm{\Phi }_h\sigma )(v)+T_x\mathrm{\Phi }_h(v)(\text{since}\mathrm{\Phi }_h\sigma =\sigma \mathrm{\Phi }_h)`$
$`=`$ $`0`$
where we have used the property $`T_x\sigma (v)=v`$ in the last step. Hence the assertion follows from Proposition 3.7. $`\mathrm{}`$
We next discuss Poisson involutions on $`(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }})`$, where $`\{,\}_{A\mathrm{\Gamma }}`$ is the Lie-Poisson structure in (2.11). The following result was obtained in \[X\] using Lie bialgebroid theory \[MX\]. We will give an elementary proof based on the explicit formula in (2.11).
###### Proposition 3.9
Let $`s:𝔤𝔤`$ be an involutive Lie algebra anti-morphism which preserves $`𝔥`$ and assume that $`sR(q)s^{}=R(s_𝔥^{}q)`$ for all $`qU`$ where $`s_𝔥=s_𝔥.`$ Then the map
$$\sigma :(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }})(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }}),(q,\lambda ,X)(s_𝔥^{}(q),s_𝔥^{}(\lambda ),s(X))$$
$`(3.3)`$
is a Poisson involution.
###### Demonstration Proof
From the property that $`sR(q)s^{}=R(s_𝔥^{}q)`$, it follows that $`s(dR(q)(\lambda )s^{}\xi )=dR(s_𝔥^{}q)(s_𝔥^{}\lambda )(\xi ).`$ The rest of the proof is plain. $`\mathrm{}`$
Remark 3.10. The virtue of our direct verification in the above proof lies in the fact that it extends to more general constructions in which the Lie bialgebroid structure is lost. We shall call $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{},\sigma ^{})`$ a symmetric coboundary dynamical Lie algebroid. In order to compute the Poisson structure on the stable locus, we shall introduce some notation which we shall use in the rest of the section. Let $`(P,\{,\}_P)`$ be a Poisson manifold and suppose $`\tau :PP`$ is a Poisson involution with stable locus $`P^\tau `$. Then for $`\phi C^{\mathrm{}}(P)`$, we put $`\stackrel{~}{\phi }=\phi P^\tau `$ and $`\phi ^\tau =\frac{1}{2}(\phi +\tau ^{}\phi ).`$ Since for $`Q=P^\tau `$, we have $`V_Q=_{xQ}ker(T_x\tau +1)`$ in the Whitney sum decomposition for $`T_QP`$. Hence it follows from (3.2) that the induced Poisson structure on $`P^\tau `$ is given by the formula
$$\{\stackrel{~}{\phi },\stackrel{~}{\psi }\}_{P^\tau }(x)=\{\phi ^\tau ,\psi ^\tau \}_P(x)$$
$`(3.4)`$
for $`xP^\tau `$ and $`\phi `$, $`\psi C^{\mathrm{}}(P).`$
###### Proposition 3.11
The Poisson structure on the stable locus $`A\mathrm{\Gamma }^\sigma `$ of the Poisson involution in (3.3) is given by
$$\begin{array}{cc}& \{\stackrel{~}{F}_1,\stackrel{~}{F}_2\}_{A\mathrm{\Gamma }^\sigma }(q,\lambda ,X)\hfill \\ \hfill =& <\lambda ,[\delta _2\stackrel{~}{F}_1,\delta _2\stackrel{~}{F}_2]>+<dR(q)(\lambda )\delta \stackrel{~}{F}_1,\delta \stackrel{~}{F}_2>\hfill \\ & +<X,ad_{R(q)\delta \stackrel{~}{F}_1\delta _2\stackrel{~}{F}_1}^{}\delta \stackrel{~}{F}_2+ad_{R(q)\delta \stackrel{~}{F}_2\delta _2\stackrel{~}{F}_2}^{}\delta \stackrel{~}{F}_1>\hfill \\ & <q,[\delta _2\stackrel{~}{F}_1,\delta _1\stackrel{~}{F}_2]+[\delta _1\stackrel{~}{F}_1,\delta _2\stackrel{~}{F}_2]>+<\delta _1\stackrel{~}{F}_2,\iota ^{}\delta \stackrel{~}{F}_1><\delta _1\stackrel{~}{F}_1,\iota ^{}\delta \stackrel{~}{F}_2>\hfill \end{array}$$
$`(3.5)`$
for $`F_1,F_2C^{\mathrm{}}(A\mathrm{\Gamma })`$, $`(q,\lambda ,X)A\mathrm{\Gamma }^\sigma `$ where
$$\begin{array}{cc}& \delta _1\stackrel{~}{F}_i:=\frac{1}{2}(\delta _1F_i+s_𝔥(\delta _1F_i)),\delta _2\stackrel{~}{F}_i:=\frac{1}{2}(\delta _2\stackrel{~}{F}_is_𝔥(\delta _2\stackrel{~}{F}_i)),\hfill \\ & \delta \stackrel{~}{F}_i:=\frac{1}{2}(\delta \stackrel{~}{F}_i+s^{}(\delta \stackrel{~}{F}_i)),i=1,2.\hfill \end{array}$$
$`(3.6)`$
Moreover, the Hamiltonian vector field on $`A\mathrm{\Gamma }^\sigma `$ generated by $`\stackrel{~}{F}`$ is of the form
$$\begin{array}{cc}& X_{\stackrel{~}{F}}(q,\lambda ,X)\hfill \\ \hfill =& (\iota ^{}\delta \stackrel{~}{F}ad_{\delta _2\stackrel{~}{F}}^{}q,ad_{\delta _2\stackrel{~}{F}}^{}\lambda +\iota ^{}ad_X^{}\delta \stackrel{~}{F}ad_{\delta _1\stackrel{~}{F}}^{}q,\hfill \\ & [X,R(q)\delta \stackrel{~}{F}\delta _2\stackrel{~}{F}]+dR(q)(\lambda )\delta \stackrel{~}{F}\delta _1\stackrel{~}{F}+R(q)ad_X^{}\delta \stackrel{~}{F}).\hfill \end{array}$$
$`(3.7)`$
###### Demonstration Proof
The bundle map $`\pi ^{\mathrm{}}`$ corresponding to the Poisson bracket $`\{,\}_{A\mathrm{\Gamma }}`$ in (2.11) is given by
$`\pi ^{\mathrm{}}(q,\lambda ,X)(Z_1,Z_2,\xi )`$
$`=`$ $`(\iota ^{}\xi ad_{Z_2}^{}q,ad_{Z_2}^{}\lambda +\iota ^{}ad_X^{}\xi ad_{Z_1}^{}q,`$
$`[X,R(q)\xi Z_2]+dR(q)(\lambda )\xi Z_1+R(q)(ad_X^{}\xi ))`$
where $`(q,\lambda ,X)A\mathrm{\Gamma }`$ and $`(Z_1,Z_2,\xi )𝔥\times 𝔥\times 𝔤^{}T_{(q,\lambda ,X)}^{}(A\mathrm{\Gamma }).`$ Let $`Q=A\mathrm{\Gamma }^\sigma `$ and let $`\iota _Q:QA\mathrm{\Gamma }`$ be the canonical inclusion. In this case, the bundle $`V_Q`$ in the vector bundle decomposition $`T(A\mathrm{\Gamma })=TQV_Q`$ is just the bundle of $`1`$ eigenspaces of $`T\sigma `$. Therefore, if $`(q,\lambda ,X)Q`$, and $`(Z_1,Z_2,\xi )𝔥\times 𝔥\times 𝔤^{}T_{(q,\lambda ,X)}^{}(A\mathrm{\Gamma }),`$ it follows that the bundle map of the induced Poisson structure on $`Q`$ is given by
$`\pi _Q^{\mathrm{}}(q,\lambda ,X)(T_{(q,\lambda ,X)}^{}\iota _Q(Z_1,Z_2,\xi ))`$
$`=`$ $`(\iota ^{}\stackrel{~}{\xi }ad_{\stackrel{~}{Z}_2}^{}q,ad_{\stackrel{~}{Z}_2}^{}\lambda +\iota ^{}ad_X^{}\stackrel{~}{\xi }ad_{\stackrel{~}{Z}_1}^{}q,`$
$`[X,R(q)\stackrel{~}{\xi }\stackrel{~}{Z}_2]+dR(q)(\lambda )\stackrel{~}{\xi }\stackrel{~}{Z}_1+R(q)(ad_X^{}\stackrel{~}{\xi }))`$
where
$`\stackrel{~}{Z}_1={\displaystyle \frac{1}{2}}(Z_1+s_𝔥(Z_1)),\stackrel{~}{Z}_2={\displaystyle \frac{1}{2}}(Z_2s_𝔥(Z_2)),`$
$`\stackrel{~}{\xi }={\displaystyle \frac{1}{2}}(\xi +s^{}(\xi )).`$
Since $`d\stackrel{~}{F}(q,\lambda ,X)=T_{(q,\lambda ,X)}^{}\iota _QdF(q,\lambda ,X)`$, the formula for the vector field $`X_{\stackrel{~}{F}}`$ is immediate from the above expression. On the other hand, the formula for the Poisson bracket is a consequence of (2.3) and (3.4) as we have $`dF_i^\sigma (q,\lambda ,X)=(\delta _1\stackrel{~}{F}_i,\delta _2F_i,\delta \stackrel{~}{F}_i),i=1,2.`$ $`\mathrm{}`$
We now turn to corresponding results for the coboundary dynamical Poisson groupoids. The following result was also obtained in \[X\] by invoking Lie bialgebroid theory. We can, of course, verify the assertion in a direct way by using the formula in (2.7).
###### Proposition 3.12
Let $`R`$ be an $`H`$-equivariant classical dynamical r-matrix such that $`sR(q)s^{}=R(s_𝔥^{}q)`$ for all $`qU`$ where $`s`$ is as in Proposition 3.9. If $`(\mathrm{\Gamma },\{,\}_R)`$ is the coboundary dynamical Poisson groupoid associated to $`R`$, then the map
$$\mathrm{\Sigma }:(\mathrm{\Gamma },\{,\}_R)(\mathrm{\Gamma },\{,\}_R),(u,g,v)(s_𝔥^{}(v),S(g),s_𝔥^{}(u))$$
$`(3.8)`$
is a Poisson involution, where $`S:GG`$ is the group anti-morphism which integrates $`s`$.
We shall call $`(\mathrm{\Gamma },\{,\}_R,\mathrm{\Sigma })`$ a symmetric coboundary dynamical Poisson groupoid.
###### Proposition 3.13
With the involution $`\mathrm{\Sigma }`$ in (3.8), the induced Poisson structure on its stable locus $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ is given by
$$\begin{array}{cc}& \{\stackrel{~}{\phi },\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}(u,g,s_𝔥^{}(u))\hfill \\ \hfill =& 2<u,[\delta _1\stackrel{~}{\phi },\delta _1\stackrel{~}{\psi }]>2<\iota \delta _1\stackrel{~}{\phi },D\stackrel{~}{\psi }>+2<\iota \delta _1\stackrel{~}{\psi },D\stackrel{~}{\phi }>\hfill \\ & +<R(s_𝔥^{}(u))D^{}\stackrel{~}{\phi },D^{}\stackrel{~}{\psi }><R(u)D\stackrel{~}{\phi },D\stackrel{~}{\psi }>\hfill \\ \hfill =& 2<u,[\delta _1\stackrel{~}{\phi },\delta _1\stackrel{~}{\psi }]>2<\iota \delta _1\stackrel{~}{\phi },D\stackrel{~}{\psi }>+2<\iota \delta _1\stackrel{~}{\psi },D\stackrel{~}{\phi }>\hfill \\ & 2<R(u)D\stackrel{~}{\phi },D\stackrel{~}{\psi }>\hfill \end{array}$$
$`(3.9)`$
for $`\phi ,\psi C^{\mathrm{}}(\mathrm{\Gamma })`$, $`(u,g,s_𝔥^{}(u))\mathrm{\Gamma }^\mathrm{\Sigma }`$, where
$$\begin{array}{cc}& \delta _1\stackrel{~}{\phi }:=\frac{1}{2}(\delta _1\phi +s_𝔥(\delta _2\phi )),D\stackrel{~}{\phi }:=\frac{1}{2}(D\phi +s^{}(D^{}\phi )),\hfill \\ & D^{}\stackrel{~}{\phi }:=\frac{1}{2}(D^{}\phi +s^{}(D\phi )),\hfill \end{array}$$
$`(3.10)`$
and similarly for $`\stackrel{~}{\psi }.`$ Here, $`\delta _1\phi `$,$`\delta _2\phi `$ are the partial derivatives of $`\phi `$ with respect to the variables in $`U`$ and $`D^{}\phi `$, $`D\phi `$ are the left and right gradients of $`\phi `$ with respect to the variable in $`G`$. Hence the Hamiltonian vector field on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ generated by $`\stackrel{~}{\phi }`$ is of the form
$$\begin{array}{cc}& X_{\stackrel{~}{\phi }}(u,g,s_𝔥^{}(u))\hfill \\ \hfill =& (ad_{\delta _1\stackrel{~}{\phi }}^{}u+\iota ^{}D\stackrel{~}{\phi },T_er_g\delta _1\stackrel{~}{\phi }T_el_gs_𝔥^{}(\delta _1\stackrel{~}{\phi })+T_el_gR(s_𝔥^{}(u))s^{}(D\stackrel{~}{\phi })\hfill \\ & T_er_gR(u)D\stackrel{~}{\phi },s_𝔥^{}(ad_{\delta _1\stackrel{~}{\phi }}^{}u+\iota ^{}D\stackrel{~}{\phi })).\hfill \end{array}$$
$`(3.11)`$
###### Demonstration Proof
According to (2.7) and (3.4), we can express the bracket $`\{\stackrel{~}{\phi },\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}`$ in the following form:
$$\begin{array}{cc}\hfill \{\stackrel{~}{\phi },\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}(u,g,s_𝔥^{}(u))=& <u,[\delta _1\stackrel{~}{\phi },\delta _1\stackrel{~}{\psi }]><s_𝔥^{}(u),[\delta _2\stackrel{~}{\phi },\delta _2\stackrel{~}{\psi }]>\hfill \\ & <\iota \delta _1\stackrel{~}{\phi },D\stackrel{~}{\psi }><\iota \delta _2\stackrel{~}{\phi },D^{}\stackrel{~}{\psi }>\hfill \\ & +<\iota \delta _1\stackrel{~}{\psi },D\stackrel{~}{\phi }>+<\iota \delta _2\stackrel{~}{\psi },D^{}\stackrel{~}{\phi }>\hfill \\ & +<R(s_𝔥^{}(u))D^{}\stackrel{~}{\phi },D^{}\stackrel{~}{\psi }><R(u)D\stackrel{~}{\phi },D\stackrel{~}{\psi }>\hfill \end{array}$$
where $`\delta _2\stackrel{~}{\phi }:=\frac{1}{2}(\delta _2\phi +s_𝔥(\delta _1\phi ))`$ and the other derivatives are defined in (3.10). However, it is easy to show that
$$<s_𝔥^{}(u),[\delta _2\stackrel{~}{\phi },\delta _2\stackrel{~}{\psi }]>=<u,[\delta _1\stackrel{~}{\phi },\delta _1\stackrel{~}{\psi }]>,<\iota \delta _2\stackrel{~}{\phi },D^{}\stackrel{~}{\psi }>=<\iota \delta _1\stackrel{~}{\phi },D\stackrel{~}{\psi }>$$
and $`<R(s_𝔥^{}(u))D^{}\stackrel{~}{\phi },D^{}\stackrel{~}{\psi }>=<R(u)D\stackrel{~}{\phi },D\stackrel{~}{\psi }>.`$ Therefore, the above expression for the bracket simplifies to the ones in the statement of the theorem. The computation of the vector field proceeds as in the proof of Proposition 3.11 and so we skip the details. $`\mathrm{}`$
Now let us recall from \[L2\] that the Lie-Poisson structure on the dual bundle $`A^{}\mathrm{\Gamma }`$ of the trivial Lie algebroid $`A\mathrm{\Gamma }`$ is given by $`\{,𝒢\}_{A^{}\mathrm{\Gamma }}(q,p,\xi )=<\delta _2,\delta _1𝒢><\delta _1,\delta _2𝒢>+<\xi ,[\delta ,\delta 𝒢]>.`$
We shall leave the proof of the next result to the reader.
###### Proposition 3.14
Let $`b:𝔥𝔥`$ be an involutive linear map and suppose $`c:𝔤𝔤`$ is an involutive Lie algebra morphism. Then the map
$$\theta :(A^{}\mathrm{\Gamma },\{,\}_{A^{}\mathrm{\Gamma }})(A^{}\mathrm{\Gamma },\{,\}_{A^{}\mathrm{\Gamma }}),(q,p,\xi )(b^{}(q),b(p),c^{}(\xi ))$$
$`(3.12)`$
is a Poisson involution. Moreover, the induced Poisson structure on the stable locus $`A^{}\mathrm{\Gamma }^\theta `$ is given by
$$\begin{array}{cc}& \{\stackrel{~}{},\stackrel{~}{𝒢}\}_{A^{}\mathrm{\Gamma }^\theta }(q,p,\xi )\hfill \\ \hfill =& <\delta _2\stackrel{~}{},\delta _1\stackrel{~}{𝒢}><\delta _1\stackrel{~}{},\delta _2\stackrel{~}{𝒢}>+<\xi ,[\delta \stackrel{~}{},\delta \stackrel{~}{𝒢}]>\hfill \end{array}$$
$`(3.13)`$
for $`,𝒢C^{\mathrm{}}(A^{}\mathrm{\Gamma })`$ and $`(q,p,\xi )A^{}\mathrm{\Gamma }^\theta `$ where
$$\begin{array}{cc}& \delta _1\stackrel{~}{}:=\frac{1}{2}(\delta _1+b(\delta _1)),\delta _2\stackrel{~}{}:=\frac{1}{2}(\delta _2+b^{}(\delta _2)),\hfill \\ & \delta \stackrel{~}{}:=\frac{1}{2}(\delta +c(\delta )),\hfill \end{array}$$
$`(3.14)`$
and similarly for $`\stackrel{~}{𝒢}`$.
Thus the induced structure $`\{,\}_{A^{}\mathrm{\Gamma }^\theta }`$ is still a product structure. Indeed, under the natural isomorphism between $`(𝔤^{})^c^{}`$ and $`(𝔤^c)^{}`$, we can identify the bracket on the stable locus $`(𝔤^{})^c^{}`$ with the Lie-Poisson structure on $`(𝔤^c)^{}`$.
We are now ready to formulate the main results of this section. In what follows, let $`X`$ be a Hamiltonian $`H`$-space (the $`H`$-action will be denoted by $`𝒞`$) with equivariant momentum map $`J:X𝔤^{}`$, and let $`\kappa :XX`$ be a Poisson involution on $`X`$. Beginning with $`H`$-invariant Hamiltonian systems on $`X`$ which admit either a realization in $`(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }})`$ or $`(\mathrm{\Gamma },\{,\}_R)`$, we shall show how reduction to Dirac submanifolds followed by Poisson reduction can lead us to integrable systems. Case 1. The case of realization in $`(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }})`$ Let $`\rho :XA\mathrm{\Gamma }`$ be a realization of the Poisson manifold $`X`$ in the dual bundle $`A\mathrm{\Gamma }`$ of the Lie algebroid $`A^{}\mathrm{\Gamma }`$; i.e., $`\rho `$ is a Poisson map. Let us recall from Theorem 2.2 that $`A\mathrm{\Gamma }`$ with the action
$$𝒜:H\times A\mathrm{\Gamma }A\mathrm{\Gamma },𝒜_h(q,\lambda ,\xi )=(Ad_{h^1}^{}q,Ad_{h^1}^{}\lambda ,Ad_h\xi )$$
$`(3.15)`$
is a Hamiltonian $`H`$-space with equivariant momentum map
$$\gamma :A\mathrm{\Gamma }𝔥^{},(q,\lambda ,\xi )\lambda .$$
$`(3.16)`$
We begin by making the following assumption:
A1. there exists a Poisson involution
$$\sigma :(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }})(A\mathrm{\Gamma },\{,\}_{A\mathrm{\Gamma }}),(q,\lambda ,\xi )(s_𝔥^{}(q),s_𝔥^{}(\lambda ),s(\xi ))$$
$`(3.17)`$
on $`A\mathrm{\Gamma }`$ (where $`s`$ satisfies the assumptions in Proposition 3.9) such that
$$\sigma \rho =\rho \kappa .$$
$`(3.18)`$
Then according to Corollary 3.6, the map $`\rho `$ restricts to a Poisson map $`\stackrel{~}{\rho }:X^\kappa A\mathrm{\Gamma }^\sigma `$, when $`X^\kappa `$ and $`A\mathrm{\Gamma }^\sigma `$ are equipped with the induced structures. Thus the stable locus $`X^\kappa `$ admits a realization in $`A\mathrm{\Gamma }^\sigma U_s\times 𝔥_s^{}\times 𝔤^s`$, where
$$U_s=\{qUs_𝔥^{}(q)=q\},$$
$`(3.19)`$
$$𝔥_s^{}=\{\lambda 𝔥^{}s_𝔥^{}(\lambda )=\lambda \},$$
$`(3.20)`$
and $`𝔤^s`$ is the fixed point set of $`s.`$ Let $`I(𝔤)`$ be the ring of ad-invariant functions on $`𝔤`$, and let $`I(𝔤^s)`$ consists of the restrictions of functions in $`I(𝔤)`$ to $`𝔤^s.`$ If $`Pr_3`$ denote the projection map from $`A\mathrm{\Gamma }^\sigma U_s\times 𝔥_s^{}\times 𝔤^s`$ to the factor $`𝔤^s`$, then a natural family of invariant functions on $`A\mathrm{\Gamma }^\sigma `$ is $`Pr_3^{}I(𝔤^s)`$. Our first result on Poisson commuting functions will have application in Section 4.2 below.
###### Theorem 3.15
Let $`\sigma `$ be a Poisson involution of the form in (3.17) on $`A\mathrm{\Gamma }`$ (where $`s`$ satifies the assumptions in Proposition 3.9) and suppose $`𝔥_s^{}=\{0\}`$, then the functions in $`Pr_3^{}I(𝔤^s)`$ Poisson commute in $`A\mathrm{\Gamma }^\sigma `$. Consequently, if we assume in addition that A1 is valid, then $`\stackrel{~}{\rho }^{}Pr_3^{}I(𝔤^s)`$ is a Poisson commuting family of functions on $`X^\kappa `$, where $`\stackrel{~}{\rho }=\rho X^\kappa `$.
###### Demonstration Proof
Let $`f_1`$, $`f_2I(𝔤)`$, and let $`\stackrel{~}{f}_1`$, $`\stackrel{~}{f}_2`$ be their restrictions to $`I(𝔤^s)`$. Then from (3.5), we have
$`\{Pr_3^{}\stackrel{~}{f}_1,Pr_3^{}\stackrel{~}{f}_2\}_{A\mathrm{\Gamma }^\sigma }(q,0,\xi )`$
$`=`$ $`\{\stackrel{~}{Pr_3^{}f_1},\stackrel{~}{Pr_3^{}f_2}\}_{A\mathrm{\Gamma }^\sigma }(q,0,\xi )`$
$`=`$ $`<\xi ,ad_{R(q)\delta \stackrel{~}{Pr_3^{}f_1}}^{}\delta \stackrel{~}{Pr_3^{}f_2}+ad_{R(q)\delta \stackrel{~}{Pr_3^{}f_2}}^{}\delta \stackrel{~}{Pr_3^{}f_1}>`$
where in the last two lines, we have used the same symbol $`Pr_3`$ to denote the projection map from $`A\mathrm{\Gamma }`$ to $`𝔤`$ and $`\stackrel{~}{Pr_3^{}f_i}`$ denote the restriction of $`Pr_3^{}f_i`$ to $`A\mathrm{\Gamma }^\sigma ,i=1,2.`$ Now, by direct calculation, we find
$$\delta \stackrel{~}{Pr_3^{}f_i}=\frac{1}{2}(df_i+s^{}(df_i)),i=1,2.$$
Therefore, upon substituting into the above expression, we obtain
$`\{Pr_3^{}\stackrel{~}{f}_1,Pr_3^{}\stackrel{~}{f}_2\}_{A\mathrm{\Gamma }^\sigma }(q,0,\xi )`$
$`=`$ $`{\displaystyle \frac{1}{4}}<[\xi ,R(q)(df_1+s^{}(df_1))],df_2+s^{}(df_2)>(12)`$
$`=`$ $`{\displaystyle \frac{1}{4}}<R(q)(df_1+s^{}(df_1)),ad_\xi ^{}df_2+ad_\xi ^{}s^{}(df_2)>(12).`$
But as $`\xi 𝔤^s`$, we have $`ad_\xi ^{}s^{}=s^{}ad_\xi ^{}.`$ Hence the first assertion follows from the fact that $`ad_\xi ^{}df_i=0`$, $`i=1,2.`$ The second assertion is now clear as assumption A1 implies that $`\stackrel{~}{\rho }`$ is Poisson by Corollary 3.6. $`\mathrm{}`$
In the general case when $`𝔥_s^{}\{0\}`$, the functions in $`Pr_3^{}I(𝔤^s)`$ is no longer a Poisson commuting family on $`A\mathrm{\Gamma }^\sigma `$. Indeed, by a computation similar to the one in the proof of the above theorem, we have
$$\begin{array}{cc}& \{Pr_3^{}\stackrel{~}{f}_1,Pr_3^{}\stackrel{~}{f}_2\}_{A\mathrm{\Gamma }^\sigma }(q,\lambda ,\xi )\hfill \\ \hfill =& \frac{1}{4}<dR(q)(\lambda )(df_1+s^{}(df_1)),df_2+s^{}(df_2)>\hfill \end{array}$$
$`(3.21)`$
for $`f_1,f_2I(𝔤).`$ Nevertheless, it is clear from this expression that when we restrict to the submanifold $`U_s\times \{0\}\times 𝔤^s`$ of $`A\mathrm{\Gamma }^\sigma `$, the bracket vanishes. Note that in general neither $`X^\kappa `$ nor $`A\mathrm{\Gamma }^\sigma `$ are Hamiltonian $`H`$-spaces. We now discuss a situation where we can obtain Poisson commuting functions on a reduced phase space. Motivated by our application in Section 4.1 below, we shall make the following assumptions to prepare the way for Poisson reduction: A2. the realization map $`\rho `$ is $`H`$-equivariant,
A3. for some Lie subgroup $`D`$ of $`H`$,
$$𝒜_d\sigma =\sigma 𝒜_d,𝒞_d\kappa =\kappa 𝒞_d,dD,$$
$`(3.22)`$
A4. if $`𝔡=Lie(D)`$ and $`𝔥_s^{}`$ is as in (3.20), we assume
$$𝔥_s^{}𝔡^{}.$$
$`(3.23)`$
###### Proposition 3.16
Under assumptions A1-A4, the stable loci $`X^\kappa `$, $`A\mathrm{\Gamma }^\sigma `$ are Hamiltonian $`D`$-spaces with equivariant momentum maps $`\stackrel{~}{J}=i_𝔡^{}(JX^\kappa )`$ and $`\stackrel{~}{\gamma }=\gamma A\mathrm{\Gamma }^\sigma `$ respectively, where $`i_𝔡:𝔡𝔥`$ in the natural inclusion and $`i_𝔡^{}`$ is its dual map. Moreover, the map
$$\stackrel{~}{\rho }=\rho X^\kappa :X^\kappa A\mathrm{\Gamma }^\sigma $$
$`(3.24)`$
is a $`D`$-equivariant Poisson map.
###### Demonstration Proof
From A3 and Corollary 3.8, it follows that the actions $`𝒜`$ and $`𝒞`$ induce Hamiltonian group actions of $`D`$ on the stable loci $`A\mathrm{\Gamma }^\sigma `$ and $`X^\kappa `$ respectively. Using the second part of the same corollary and A4, we can easily obtain the equivariant momentum maps of these induced actions. Finally, that the map $`\stackrel{~}{\rho }`$ is well-defined and Poisson is a consequence of A1 and Corollary 3.6, and its $`D`$-equivariance is obvious from A2. $`\mathrm{}`$
The above proposition completes the first stage of our reduction process in the general case and prepares the way for Poisson reduction. In order to obtain Poisson commuting functions in this general case, we shall make an additional assumption: A5. for some regular value $`\mu 𝔡^{}`$ of $`\stackrel{~}{J},`$
$$\stackrel{~}{\rho }(\stackrel{~}{J}^1(\mu ))\stackrel{~}{\gamma }^1(0)=U_s\times \{0\}\times 𝔤^s.$$
$`(3.25)`$
With this additional assumption, we will construct the integrable systems and their realizations by Poisson reduction of the map $`\stackrel{~}{\rho }`$ in Proposition 3.16. Let $`D_\mu `$ be the isotropy subgroup of $`D`$ for the $`D`$-action on $`X^\kappa `$, then by Poisson reduction \[MR\],\[OR\], the variety $`X_\mu ^\kappa =\stackrel{~}{J}^1(\mu )/D_\mu `$ inherits a unique Poisson structure $`\{,\}_{X_\mu ^\kappa }`$ satisfying
$$\pi _\mu ^{}\{f_1,f_2\}_{X_\mu ^\kappa }=i_\mu ^{}\{f_1^{},f_2^{}\}_{X^\kappa }.$$
$`(3.26)`$
Here, $`i_\mu :\stackrel{~}{J}^1(\mu )X^\kappa `$ is the inclusion map, $`\pi _\mu :\stackrel{~}{J}^1(\mu )X_\mu ^\kappa `$ is the canonical projection, $`f_1`$, $`f_2C^{\mathrm{}}(X_\mu ^\kappa )`$, and $`f_1^{}`$, $`f_2^{}`$ are (locally defined) smooth extensions of $`\pi _\mu ^{}f_1`$, $`\pi _\mu ^{}f_2`$ with differentials vanishing on the tangent spaces of the $`D`$-orbits. Similarly, we have the Poisson variety
$$\left(A\mathrm{\Gamma }_0^\sigma =\stackrel{~}{\gamma }^1(0)/D,\{,\}_{A\mathrm{\Gamma }_0^\sigma }\right),$$
$`(3.27)`$
with the inclusion map $`i_D:\stackrel{~}{\gamma }^1(0)A\mathrm{\Gamma }^\sigma `$ and the canonical projection $`\pi _D:\stackrel{~}{\gamma }^1(0)A\mathrm{\Gamma }_0^\sigma .`$ If $`Pr_i`$ denotes the projection map onto the $`i`$-th factor of $`U_s\times 𝔥_s^{}\times 𝔤^sA\mathrm{\Gamma }`$, $`i=1,2,3,`$ we put
$$m=Pr_1\stackrel{~}{\rho }:X^\kappa U_s,$$
$`(3.28)`$
$$\tau =Pr_2\stackrel{~}{\rho }:X^\kappa 𝔥_s^{},$$
$`(3.29)`$
$$L=Pr_3\stackrel{~}{\rho }:X^\kappa 𝔤^s.$$
$`(3.30)`$
Clearly, functions in $`i_D^{}Pr_3^{}I(𝔤^s)C^{\mathrm{}}(\stackrel{~}{\gamma }^1(0))`$ are $`D`$-invariant, hence they descend to functions in $`C^{\mathrm{}}(A\mathrm{\Gamma }_0^\sigma )`$. On the other hand, it follows from Proposition 3.16 that the functions in $`i_\mu ^{}L^{}I(𝔤^s)C^{\mathrm{}}(\stackrel{~}{J}^1(\mu ))`$ drop down to functions in $`C^{\mathrm{}}(X_\mu ^\kappa ).`$ Now,by Proposition 3.16 and assumption A5, it follows from Theorem 2.14 of \[OR\] that $`\stackrel{~}{\rho }`$ induces a unique Poisson map (called the reduction of $`\stackrel{~}{\rho }`$)
$$\widehat{\rho }:X_\mu ^\kappa A\mathrm{\Gamma }_0^\sigma =(U_s\times \{0\}\times 𝔤^s)/D$$
$`(3.31)`$
characterized by $`\pi _D\stackrel{~}{\rho }i_\mu =\widehat{\rho }\pi _\mu .`$ Hence $`X_\mu ^\kappa `$ admits a realization in the Poisson variety $`A\mathrm{\Gamma }_0^\sigma `$.
We shall use the following notation. For $`fI(𝔤)`$, the unique function in $`C^{\mathrm{}}(A\mathrm{\Gamma }_0^\sigma )`$ determined by $`i_D^{}Pr_3^{}\stackrel{~}{f}`$ ($`\stackrel{~}{f}=f𝔤^s`$) will be denoted by $`\overline{f}`$; while the unique function in $`C^{\mathrm{}}(X_\mu ^\kappa )`$ determined by $`i_\mu ^{}L^{}\stackrel{~}{f}`$ will be denoted by $`_\mu `$. From the definitions, we have
$$_\mu \pi _\mu =(\widehat{\rho }^{}\overline{f})\pi _\mu =i_\mu ^{}L^{}\stackrel{~}{f}.$$
$`(3.32)`$
###### Theorem 3.17
If $`𝔥_s^{}\{0\}`$, then under assumptions A1-A5, the map $`\stackrel{~}{\rho }=\rho X^\kappa :X^\kappa A\mathrm{\Gamma }^\sigma `$ induces a unique Poisson map $`\widehat{\rho }:X_\mu ^\kappa A\mathrm{\Gamma }_0^\sigma `$ such that
(a) functions $`_\mu =\widehat{\rho }^{}\overline{f}`$, $`fI(𝔤)`$, Poisson commute in $`(X_\mu ^\kappa ,\{,\}_{X_\mu ^\kappa })`$, (b) if $`\psi _t`$ is the induced flow on $`\stackrel{~}{\gamma }^1(0)=U_s\times \{0\}\times 𝔤^s`$ generated by the Hamiltonian $`Pr_3^{}\stackrel{~}{f}`$, $`fI(𝔤)`$, and $`\varphi _t`$ is the Hamiltonian flow of $`=L^{}\stackrel{~}{f}`$ on $`X^\kappa `$, then under the flow $`\varphi _t`$, we have
$`{\displaystyle \frac{d}{dt}}m(\varphi _t)={\displaystyle \frac{1}{2}}\iota ^{}(df(L(\varphi _t))+s^{}(df(L(\varphi _t)))),`$
$`{\displaystyle \frac{d}{dt}}\tau (\varphi _t)=0,`$
$`{\displaystyle \frac{d}{dt}}L(\varphi _t)={\displaystyle \frac{1}{2}}[L(\varphi _t),R(m(\varphi _t))(df(L(\varphi _t))+s^{}(df(L(\varphi _t)))]`$
$`+dR(m(\varphi _t))(\tau (\varphi _t))(df(L(\varphi _t))+s^{}(df(L(\varphi _t))))`$
where the term involving $`dR`$ drops out on $`\stackrel{~}{J}^1(\mu )`$. Moreover, the reduction $`\varphi _t^{red}`$ of $`\varphi _ti_\mu `$ on $`X_\mu ^\kappa `$ defined by $`\varphi _t^{red}\pi _\mu =\pi _\mu \varphi _ti_\mu `$ is a Hamiltonian flow of $`_\mu =\widehat{\rho }^{}\overline{f}`$ and $`\widehat{\rho }\varphi _t^{red}(\pi _\mu (x))=\pi _D\psi _t(\stackrel{~}{\rho }(x))`$, $`x\stackrel{~}{J}^1(\mu )`$.
###### Demonstration Proof
(a) Let $`f_1`$, $`f_2I(𝔤)`$, then it is easy to check that $`Pr_3^{}\stackrel{~}{f}_1`$,$`Pr_3^{}\stackrel{~}{f}_2`$ are extensions of $`\pi _D^{}\overline{f}_1`$, $`\pi _D^{}\overline{f}_2`$ with differentials vanishing on the tangent spaces of the $`D`$\- orbits. Therefore, if $`x\stackrel{~}{J}^1(\mu ))`$, we have $`\tau (x)=0`$ by assumption A5 and hence
$`\{\widehat{\rho }^{}\overline{f}_1,\widehat{\rho }^{}\overline{f}_2\}_{X_\mu ^\kappa }\pi _\mu (x)`$
$`=`$ $`\{\overline{f}_1,\overline{f}_2\}_{A\mathrm{\Gamma }_0^\sigma }\pi _D(\stackrel{~}{\rho }(x))`$
$`=`$ $`\{Pr_3^{}\stackrel{~}{f}_1,Pr_3^{}\stackrel{~}{f}_2\}_{A\mathrm{\Gamma }^\sigma }(\stackrel{~}{\rho }(x))`$
$`=`$ $`{\displaystyle \frac{1}{4}}<[L(x),R(m(x))(df_1+s^{}(df_1))],df_2+s^{}(df_2)>(12)`$
$`=`$ $`{\displaystyle \frac{1}{4}}<R(m(x))(df_1+s^{}(df_1)),ad_{L(x)}^{}df_2s^{}ad_{L(x)}^{}df_2>(12)`$
$`=`$ $`0.`$
(b) The equations of motion is a consequence of Proposition 3.11 and the fact that $`\stackrel{~}{\rho }`$ is Poisson. On the other hand, the assertion on $`\varphi _t^{red}`$ is a corollary of Theorem 2.16 in \[OR\] and the relation $`\stackrel{~}{\rho }\varphi _ti_\mu =\psi _t\stackrel{~}{\rho }i_\mu `$. $`\mathrm{}`$
Case 2. The case of realization in $`(\mathrm{\Gamma },\{,\}_R)`$
Let $`𝒫:X\mathrm{\Gamma }`$ be a realization map of $`X`$ in the coboundary dynamical Poisson groupoid $`(\mathrm{\Gamma },\{,\}_R)`$. Recall from \[L1\] that $`\mathrm{\Gamma }`$ equipped with the action
$$:H\times \mathrm{\Gamma }\mathrm{\Gamma },_h(u,g,v)=(Ad_{h^1}^{}u,hgh^1,Ad_{h^1}^{}v)$$
$`(3.33)`$
is a Hamiltonian $`H`$-space with equivariant momentum map
$$\alpha \beta :\mathrm{\Gamma }𝔥^{},(u,g,v)uv.$$
$`(3.34)`$
We begin with the following assumption: G1. there exists a Poisson involution
$$\mathrm{\Sigma }:(\mathrm{\Gamma },\{,\}_R)(\mathrm{\Gamma },\{,\}_R),(u,g,v)(s_𝔥^{}(v),S(g),s_𝔥^{}(u))$$
$`(3.35)`$
on $`\mathrm{\Gamma }`$ (where $`s`$ satisfies the assumptions in Proposition 3.9) such that
$$\mathrm{\Sigma }𝒫=𝒫\kappa .$$
$`(3.36)`$
Under G1, the map $`\stackrel{~}{𝒫}=𝒫X^\kappa :X^\kappa \mathrm{\Gamma }^\mathrm{\Sigma }`$ is a well-defined Poisson map by Corollary 3.6, when the stable loci are equipped with the induced structures. Let $`I(G)`$ be the ring of central functions in $`G`$ and let $`I(G^S)`$ consists of restrictions of functions in $`I(G)`$ to the stable locus $`G^S`$ of $`S`$. If $`𝙿r_2`$ denote the projection map $`\mathrm{\Gamma }^\mathrm{\Sigma }G^S,(u,g,s_𝔥^{}(u))g`$, a natural family of invariant functions on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ is $`𝙿r_2^{}I(G^S)`$. As in the algebroid case, we begin with a special situation.
###### Theorem 3.18
If $`s_𝔥^{}(u)=u`$ for all $`uU`$ so that $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ coincides with the gauge group bundle $`\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$, then the functions in $`𝙿r_2^{}I(G^S)`$ Poisson commutes in $`\mathrm{\Gamma }^\mathrm{\Sigma }U\times G^S`$. Therefore, under the additional assumption that G1 is satisfied, $`\stackrel{~}{𝒫}^{}𝙿r_2^{}I(G^S)`$ is a Poisson commuting family of functions on $`X^\kappa .`$
###### Demonstration Proof
Let $`\phi `$, $`\psi I(G)`$ and let $`\stackrel{~}{\phi }=\phi G^S`$, $`\stackrel{~}{\psi }=\psi G^S`$. Then on using the first expression in (3.9), we have
$`\{𝙿r_2^{}\stackrel{~}{\phi },𝙿r_2^{}\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}(u,g,u)`$
$`=`$ $`<R(u)D^{}\stackrel{~}{𝙿r_2^{}\phi },D^{}\stackrel{~}{𝙿r_2^{}\psi }><R(u)D\stackrel{~}{𝙿r_2^{}\phi },D\stackrel{~}{𝙿r_2^{}\psi }>`$
where in the second line of the above formula, we have used the same symbol $`𝙿r_2`$ to denote the projection map from $`\mathrm{\Gamma }`$ to $`G`$. Now, by a direct computation, we can check that
$$D^{}\stackrel{~}{𝙿r_2^{}\phi }=D\stackrel{~}{𝙿r_2^{}\phi }=\frac{1}{2}(D\phi +s^{}(D\phi )).$$
Hence the two terms in the above expression cancel out. The second assertion is now clear as $`\stackrel{~}{𝒫}`$ is Poisson under G1. $`\mathrm{}`$
In the general case when the assumption in the above theorem is not satisfied, we have
$$\begin{array}{cc}\hfill \{𝙿r_2^{}\stackrel{~}{\phi },𝙿r_2^{}\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}(u,g,s_𝔥^{}(u))=& \frac{1}{4}<R(s_𝔥^{}(u))(D\phi +s^{}(D\phi )),D\psi +s^{}(D\psi )>\hfill \\ & \frac{1}{4}<R(u)(D\phi +s^{}(D\phi )),D\psi +s^{}(D\psi )>\hfill \end{array}$$
$`(3.37)`$
for $`\phi `$, $`\psi I(G)`$. Therefore, $`𝙿r_2^{}I(G^S)`$ is no longer a Poisson commuting family of functions on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$. However, the two terms in (3.37) above do cancel out on $`\mathrm{\Gamma }^\mathrm{\Sigma }\mathrm{\Gamma }=\{(u,g,u)uU_s,gG^S\}`$ where $`U_s`$ is defined in (3.19). Analogous to Case 1, we now describe a situation where we can construct Poisson commuting functions on a reduced phase space. To prepare the way for Poisson reduction, we shall make the following assumptions in addition to G1: G2. the realization map $`𝒫`$ is $`H`$-equivariant,
G3. for some Lie subgroup $`D`$ of $`H`$,
$$_d\mathrm{\Sigma }=\mathrm{\Sigma }_d,𝒞_d\kappa =\kappa 𝒞_d,dD,$$
$`(3.38)`$
G4. $`us_𝔥^{}(u)𝔡^{}`$ for all $`uU,`$ where $`𝔡=Lie(D).`$
###### Proposition 3.19
Under assumptions G1-G4, the stable loci $`X^\kappa `$, $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ are Hamiltonian $`D`$-spaces with equivariant momentum maps $`\stackrel{~}{J}=i_𝔡^{}(JX^\kappa )`$ and $`\stackrel{~}{\alpha }\stackrel{~}{\beta }=\alpha \beta \mathrm{\Gamma }^\mathrm{\Sigma }`$ respectively. Moreover, the map
$$\stackrel{~}{𝒫}=𝒫X^\kappa :X^\kappa \mathrm{\Gamma }^\mathrm{\Sigma }$$
$`(3.39)`$
is a $`D`$-equivariant Poisson map.
###### Demonstration Proof
The assertion follows from Corollaries 3.6 and 3.8, as in Proposition 3.16. $`\mathrm{}`$
In order to obtain Poisson commuting functions in the general case, it is natural (in view of the remark after (3.37)) to make the following additional assumption: G5. for some regular value $`\mu 𝔡^{}`$ of $`\stackrel{~}{J},`$
$$\stackrel{~}{𝒫}(\stackrel{~}{J}^1(\mu ))(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)U_s\times G^S.$$
$`(3.40)`$
Analogous to the algebroid case, we have the Poisson variety
$$\left(\mathrm{\Gamma }_0^\mathrm{\Sigma }=(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)/D,\{,\}_{\mathrm{\Gamma }_0^\mathrm{\Sigma }}\right)$$
$`(3.41)`$
with the inclusion map $`𝚒_D:(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)\mathrm{\Gamma }^\mathrm{\Sigma }`$ and the canonical projection $`𝚙r_D:(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)\mathrm{\Gamma }_0^\mathrm{\Sigma }.`$ Moreover, under G5, the map $`\stackrel{~}{𝒫}`$ in Proposition 3.19 induces a Poisson map
$$\widehat{𝒫}:X_\mu ^\kappa \mathrm{\Gamma }_0^\mathrm{\Sigma }(U_s\times G^S)/D.$$
$`(3.42)`$
We shall use the following notation. For $`\phi I(G)`$, the unique function in $`C^{\mathrm{}}(\mathrm{\Gamma }_0^\mathrm{\Sigma })`$ determined by $`𝚒_D^{}𝙿r_2^{}\stackrel{~}{\phi }`$ ($`\stackrel{~}{\phi }=\phi G^S`$) will be denoted by $`\overline{\phi }.`$ Also, we set
$$𝙻=𝙿r_2\stackrel{~}{𝒫}:X^\kappa G^S,$$
$`(3.43)`$
$$𝚖_1=\stackrel{~}{\alpha }\stackrel{~}{𝒫}:X^\kappa U,$$
$`(3.44)`$
and
$$𝚖_2=\stackrel{~}{\beta }\stackrel{~}{𝒫}:X^\kappa U,$$
$`(3.45)`$
i.e. $`\stackrel{~}{𝒫}=(𝚖_1,𝙻,𝚖_2)`$.
###### Theorem 3.20
If $`\mathrm{\Gamma }^\mathrm{\Sigma }\mathrm{\Gamma }`$, then under assumptions G1-G5, there exists a unique Poisson structure $`\{,\}_{X_\mu ^\kappa }`$ on the reduced space $`X_\mu ^\kappa =\stackrel{~}{J}^1(\mu )/D_\mu `$ and a unique Poisson map $`\widehat{𝒫}:X_\mu ^\kappa \mathrm{\Gamma }_0^\mathrm{\Sigma }`$ such that
(a) functions $`\widehat{𝒫}^{}\overline{\phi }`$, $`\phi I(G)`$, Poisson commute in $`(X_\mu ^\kappa ,\{,\}_{X_\mu ^\kappa })`$, (b) if $`\psi _t`$ is the induced flow on $`(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)\mathrm{\Gamma }`$ generated by the Hamiltonian $`𝙿r_2^{}\stackrel{~}{\phi }`$, $`\phi I(G)`$ and $`\varphi _t`$ is the Hamiltonian flow of $`𝙻^{}\stackrel{~}{\phi }`$ on $`X^\kappa `$, then under the flow $`\varphi _t`$, we have
$`{\displaystyle \frac{d}{dt}}𝚖_1(\varphi _t)={\displaystyle \frac{1}{2}}\iota ^{}(D\phi +s^{}(D\phi ))`$
$`{\displaystyle \frac{d}{dt}}𝙻(\varphi _t)={\displaystyle \frac{1}{2}}T_el_{𝙻(\varphi _t)}R(s_𝔥^{}(𝚖_1(\varphi _t))(D\phi +s^{}(D\phi ))`$
$`{\displaystyle \frac{1}{2}}T_er_{𝙻(\varphi _t)}R(𝚖_1(\varphi _t))(D\phi +s^{}(D\phi ))`$
$`{\displaystyle \frac{d}{dt}}𝚖_2(\varphi _t)={\displaystyle \frac{1}{2}}\iota ^{}(D\phi +s^{}(D\phi )).`$
###### Demonstration Proof
(a) Let $`\phi _1`$, $`\phi _2I(G)`$, then from the invariance properties of these functions, we can check that $`𝙿r_2^{}\stackrel{~}{\phi }_1`$, $`𝙿r_2^{}\stackrel{~}{\phi }_2`$ are extensions of $`𝚙r_D^{}\overline{\phi }_1`$, $`𝚙r_D^{}\overline{\phi }_2`$ with differentials vanishing on the tangent spaces of the $`D`$-orbits. Since $`\widehat{𝒫}`$ is Poisson, by making use of this fact, it follows that for $`x\stackrel{~}{J}^1(\mu )`$, we have
$`\{\widehat{𝒫}^{}\overline{\phi }_1,\widehat{𝒫}^{}\overline{\phi }_2\}_{X_\mu ^\kappa }\pi _\mu (x)`$
$`=`$ $`\{\overline{\phi }_1,\overline{\phi }_2\}_{\mathrm{\Gamma }_0^\mathrm{\Sigma }}𝚙r_D(\stackrel{~}{𝒫}(x))`$
$`=`$ $`\{𝙿r_2^{}\stackrel{~}{\phi }_1,𝙿r_2^{}\stackrel{~}{\phi }_2\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}\stackrel{~}{𝒫}(x)`$
$`=`$ $`0`$
where in the last step we have invoked the formula in (3.37) and assumption G5.
(b) This follows from Proposition 3.13 and the fact that $`\stackrel{~}{𝒫}`$ is Poisson.
## 4. Spin Calogero-Moser systems associated with symmetric Lie subal- fak gebras
A symmetric Lie algebra is a Lie algebra equipped with a Lie algebra involution. If $`(𝔤,\eta )`$ is a symmetric Lie algebra, then the fixed point set $`𝔤^\eta `$ will be called a symmetric Lie subalgebra of $`𝔤.`$ In this section, we shall show that the general scheme in Section 3 can be applied to several examples of spin Calogero-Moser systems associated with real symmetric Lie algebras. Because the spin variables of the Dirac reduction belong to symmetric Lie subalgebras, we shall called the reduced systems spin Calogero-Moser systems associated with symmetric Lie subalgebras. In the following, we shall restrict ourselves to the trigonometric case. It will be clear that the rational case and the elliptic case can also be handled in a similar fashion and for this reason, we shall omit the details. (See Remark 4.1.12(b) and Remark 4.2.7(b) in this connection.)
## 4.1 Compact real forms of some spin Calogero-Moser systems
We begin by introducing a number of Lie-theoretic objects which we will use throughout the present and the next subsections.
Let $`𝔤`$ be a complex simple Lie algebra of rank $`N`$ with Killing form $`(,).`$ We fix a Cartan subalgebra $`𝔥`$ and let $`𝔤=𝔥_{\alpha \mathrm{\Delta }}𝔤_\alpha `$ be the root space decomposition of $`𝔤`$ with respect to $`𝔥.`$ For each $`\alpha \mathrm{\Delta }`$, denote by $`H_\alpha `$ the element in $`𝔥`$ which corresponds to $`\alpha `$ under the isomorphism between $`𝔥`$ and $`𝔥^{}`$ induced by the Killing form $`(,)`$. On the other hand, for each $`\alpha \mathrm{\Delta }`$, we choose root vectors $`e_\alpha 𝔤_\alpha `$ such that for all $`\alpha ,\beta \mathrm{\Delta }`$, (i) $`[e_\alpha ,e_\alpha ]=H_\alpha `$,
(ii) the constants $`N_{\alpha ,\beta }`$ in the relations
$$[e_\alpha ,e_\beta ]=N_{\alpha ,\beta }e_{\alpha +\beta },\alpha ,\beta ,\alpha +\beta \mathrm{\Delta }$$
are real and satisfy $`N_{\alpha ,\beta }=N_{\alpha ,\beta }.`$
With the notation introduced above, we define
$$𝔥_0=\underset{\alpha \mathrm{\Delta }}{}H_\alpha .$$
$`(\mathrm{4.1.1})`$
Then $`(,)_{𝔥_0\times 𝔥_0}`$ is positive definite and each root is real-valued on $`𝔥_0`$. We shall fix an orthonormal basis $`(x_i)_{1iN}`$ of $`𝔥_0`$ in what follows.
The Lie algebra $`𝔤`$ has two standard real forms, namely, the normal real form
$$𝔤_0=𝔥_0+\underset{\alpha \mathrm{\Delta }}{}e_\alpha $$
$`(\mathrm{4.1.2})`$
and the compact real form
$$𝔲_0=i𝔥_0+\underset{\alpha \mathrm{\Delta }}{}(e_\alpha e_\alpha )+\underset{\alpha \mathrm{\Delta }}{}i(e_\alpha +e_\alpha ).$$
$`(\mathrm{4.1.3})`$
Therefore, if $`𝔤^{}`$ denote the algebra $`𝔤`$ regarded as a real Lie algebra, we have $`𝔤^{}=𝔤_0i𝔤_0=𝔲_0i𝔲_0.`$ We shall denote the conjugation of $`𝔤`$ with respect to $`𝔤_0`$ and $`𝔲_0`$ by $`\upsilon `$ and $`\tau `$ respectively. For simplicity of notation, we also write
$$\upsilon (q)=\overline{q},q𝔥.$$
$`(\mathrm{4.1.4})`$
The pairing on $`𝔤^{}`$ will be taken to be the Killing form on $`𝔤^{}`$ scaled by the factor $`\frac{1}{2}`$, and will be denoted by $`(,)_{}`$.
In addition to the finite dimensional Lie algebras above, we will also need their corresponding loop algebras $`L𝔤`$, $`L𝔤^{}`$ and so on. For $`XL𝔤`$, we shall write $`X(z)=_{\mathrm{}}^{\mathrm{}}X_nz^n`$ with coefficients $`X_n𝔤`$ and similarly for the other loop algebras. Using the Killing form $`(,)`$ on $`𝔤`$, we can define a non-degenerate invariant pairing on $`L𝔤`$:
$$(X,Y)_{L𝔤}=\underset{j}{}(X_j,Y_{(j+1)}),X,YL𝔤.$$
$`(\mathrm{4.1.5})`$
Similarly, we have a pairing on $`L𝔤^{}`$, given by
$$(X,Y)_{L𝔤^{}}=\underset{j}{}(X_j,Y_{(j+1)})_{},X,YL𝔤^{}.$$
$`(\mathrm{4.1.6})`$
In what follows, the connected and simply-connected Lie groups which integrate the Lie algebras $`𝔤^{}`$, $`𝔥^{}`$, $`𝔤_0`$, $`𝔥_0`$ and $`i𝔥_0`$ will be denoted by $`G^{}`$, $`H^{}`$, $`G_0`$, $`H_0`$ and $`T`$ respectively. On the other hand, $`U`$ will denote a fixed connected component of $`\{q𝔥\mathrm{sin}(\frac{1}{2}\alpha (q))0\text{for all}\alpha \mathrm{\Delta }\}`$. Using the non-degeneracy of $`(,)_{}_{𝔥^{}\times 𝔥^{}}`$, $`(,)_{i𝔥^0\times i𝔥^0}`$ and the pairings above, we shall make the identifications $`(𝔤^{})^{}𝔤^{}`$, $`(𝔥^{})^{}𝔥^{}`$, $`(i𝔥_0)^{}i𝔥_0`$, $`𝔥_0^{}𝔥_0`$, $`L𝔤_{}^{}{}_{}{}^{}L𝔤^{}`$, where $`L𝔤_{}^{}{}_{}{}^{}`$ is the restricted dual.
Consider the trigonometric dynamical r-matrix with spectral parameter (which is gauge equivalent to the one in \[EV\]):
$$r(q,z)=\left(c(z)+\frac{1}{12}z\right)\underset{i}{}x_ix_i+\underset{\alpha \mathrm{\Delta }}{}\varphi _\alpha (q,z)e^{\frac{z}{12}\alpha (q)}e_\alpha e_\alpha $$
$`(\mathrm{4.1.7})`$
where
$$c(z)=\frac{1}{2}\mathrm{cot}\left(\frac{1}{2}z\right)$$
$`(\mathrm{4.1.8})`$
and
$$\varphi _\alpha (q,z)=(c(z)+c(\alpha (q)))$$
$`(\mathrm{4.1.9})`$
Then for each $`qU`$, we can define a map $`r_{}^\mathrm{\#}(q):𝔤L𝔤`$ by the formula
$$((r_{}^\mathrm{\#}(q)\xi )(z),\eta )=(r(q,z),\eta \xi )$$
$`(\mathrm{4.1.10})`$
where $`\xi `$, $`\eta 𝔤.`$ Therefore, if we write $`\xi =_i\xi _ix_i+_{\alpha \mathrm{\Delta }}\xi _\alpha e_\alpha `$ for $`\xi 𝔤`$, we have explicitly that
$$(r_{}^\mathrm{\#}(q)\xi )(z)=d(z)\underset{i}{}\xi _ix_i+\underset{\alpha \mathrm{\Delta }}{}\varphi _\alpha (q,z)e^{\frac{z}{12}\alpha (q)}\xi _\alpha e_\alpha $$
$`(\mathrm{4.1.11})`$
where $`d(z)=c(z)+\frac{1}{12}z`$. We can also construct the associated classical dynamical r-matrix $`R:UL(L𝔤,L𝔤)`$ for the pair $`(L𝔤,𝔥).`$ To do so, we will need to use the following formula:
$$\begin{array}{cc}& \frac{^kr}{z^k}(q,z)\hfill \\ \hfill =& d^{(k)}(z)\underset{i}{}x_ix_i+\underset{\alpha \mathrm{\Delta }}{}\underset{j=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{j}\right)\varphi _\alpha ^{(kj)}(q,z)\left(\frac{1}{12}\alpha (q)\right)^je^{\frac{z}{12}\alpha (q)}e_\alpha e_\alpha .\hfill \end{array}$$
$`(\mathrm{4.1.12})`$
###### Proposition 4.1.1
The classical dynamical r-matrix $`R`$ associated with the trigonometric dynamical r-matrix with spectral parameter in (4.1.7) is given by
$$\begin{array}{cc}& (R(q)X)(z)\hfill \\ \hfill =& \frac{1}{2}X(z)+\underset{k=0}{\overset{\mathrm{}}{}}\frac{d^{(k)}(z)}{k!}\mathrm{\Pi }_𝔥X_{(k+1)}\hfill \\ & +\underset{k0}{}\frac{1}{k!}\underset{\alpha \mathrm{\Delta }}{}\underset{j=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{j}\right)\varphi _\alpha ^{(kj)}(q,z)\left(\frac{1}{12}\alpha (q)\right)^je^{\frac{z}{12}\alpha (q)}(X_{(k+1)})_\alpha e_\alpha .\hfill \end{array}$$
$`(\mathrm{4.1.13})`$
###### Demonstration Proof
This follows upon substituting the expression for $`\frac{^kr}{z^k}(q,z)`$ in (4.1.12) into the formula for $`R(q)X`$ in Theorem 2.4(b). $`\mathrm{}`$
Clearly, $`R`$ induces a map $`UL(L𝔤^{},L𝔤^{})`$ which we shall also denote by $`R`$.
###### Proposition 4.1.2
The map $`R:UL(L𝔤^{},L𝔤^{})`$ is a solution of the mDYBE for the pair $`(L𝔤^{},𝔥^{})`$ with $`c=\frac{1}{4}`$. Moreover, for $`qU`$, $`\xi 𝔤^{}`$, we have $`r_{}^\mathrm{\#}(q)\xi L𝔤^{}.`$
###### Demonstration Proof
For $`qU`$, $`X`$, $`YL𝔤`$, the term $`(dR(q)()X,Y)_{L𝔤}`$ which appears in the mDYBE for the pair $`(L𝔤,𝔥)`$ is the unique element in $`𝔥`$ whose pairing with $`Z𝔥`$ is given by $`(dR(q)(Z)X,Y)_{L𝔤}`$. On the other hand, the term $`(dR(q)()X,Y)_{L𝔤^{}}𝔥^{}`$ has a similar meaning. But now it is easy to show from the non-degeneracy of $`(,)_{}_{𝔥^{}\times 𝔥^{}}`$ that $`(dR(q)()X,Y)_{L𝔤}=(dR(q)()X,Y)_{L𝔤^{}}.`$ Hence it follows from this argument and Theorem 2.4(a) that the map $`R:UL(L𝔤^{},L𝔤^{})`$ is a solution of the mDYBE for the pair $`(L𝔤^{},𝔥^{})`$. The assertion involving $`r_{}^\mathrm{\#}(q)\xi `$ is trivial. $`\mathrm{}`$
We now introduce the trivial Lie groupoids
$$\mathrm{\Omega }=U\times G^{}\times U,\mathrm{\Gamma }=U\times LG^{}\times U.$$
$`(\mathrm{4.1.14})`$
By Proposition 4.1.2 and (2.9), we can use the map $`R:UL(L𝔤^{},L𝔤^{})`$ to construct the associated coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }U\times 𝔥^{}\times L𝔤^{}`$. Hence its dual bundle $`A\mathrm{\Gamma }U\times 𝔥^{}\times L𝔤^{}`$ has a Lie-Poisson structure. On the other hand, we shall equip the the dual bundle $`A^{}\mathrm{\Omega }U\times 𝔥^{}\times 𝔤^{}`$ of the trivial Lie algebroid $`A\mathrm{\Omega }`$ with the corresponding Lie-Poisson structure. The Poisson manifolds $`A^{}\mathrm{\Omega }`$ and $`A\mathrm{\Gamma }`$ are Hamiltonian $`H^{}`$-spaces with actions
$$𝒞_h(q,p,\xi )=(q,p,Ad_h\xi ),hH^{},(q,p,\xi )A^{}\mathrm{\Omega }$$
$`(\mathrm{4.1.15})`$
and
$$𝒜_h(q,p.X)=(q,p,Ad_hX),hH^{},(q,p,X)A\mathrm{\Gamma }$$
$`(\mathrm{4.1.16})`$
and the corresponding equivariant momentum maps are respectively given by
$$J:A^{}\mathrm{\Omega }𝔥^{},(q,p,\xi )\mathrm{\Pi }_𝔥^{}\xi $$
$`(\mathrm{4.1.17})`$
and
$$\gamma :A\mathrm{\Gamma }𝔥^{},(q,p,X)p$$
$`(\mathrm{4.1.18})`$
where $`\mathrm{\Pi }_𝔥^{}`$ is the projection map to $`𝔥^{}`$ relative to the splitting $`𝔤^{}=𝔥^{}(𝔥^{})^{}`$.
In view of our discussion above, the following result is just a real analog of what we have in Theorem 2.5.
###### Proposition 4.1.3
The map
$$\rho :A^{}\mathrm{\Omega }A\mathrm{\Gamma },(q,\mathrm{\Pi }_𝔥^{}\xi ,p+r_{}^\mathrm{\#}(q)\xi )$$
$`(\mathrm{4.1.19})`$
is an $`H^{}`$-equivariant Poisson map.
The trigonometric spin Calogero-Moser system which we consider for our purpose here is the Hamiltonian system on $`A^{}\mathrm{\Omega }`$ generated by the Hamiltonian
$$(q,p,\xi )=Re\left\{\frac{1}{2}\underset{i}{}p_i^2\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)\xi _\alpha \xi _\alpha \right\}.$$
$`(\mathrm{4.1.20})`$
Let $`Q`$ be the quadratic function
$$Q(X)=\frac{1}{2}Re_C(X(z),X(z))\frac{dz}{2\pi iz},XL𝔤^{}$$
$`(\mathrm{4.1.21})`$
where $`C`$ is a small circle around the origin with the positive orientation.
The relation between $``$ and $`Q`$ is given in our next result which is obtained by a simple residue calculation.
###### Proposition 4.1.4
$`=\rho ^{}Pr_3^{}Q`$, where $`Pr_3`$ is the projection map onto the third factor of $`A\mathrm{\Gamma }.`$ Thus the Hamiltonian system generated by $``$ can be realized in $`A\mathrm{\Gamma }.`$
We next examine the phase space underlying $``$.
###### Proposition 4.1.5
The map
$$\kappa :A^{}\mathrm{\Omega }A^{}\mathrm{\Omega },(q,p,\xi )(\overline{q},\overline{p},\tau (\xi ))$$
$`(\mathrm{4.1.22})`$
is a Poisson involution with stable locus
$$A^{}\mathrm{\Omega }^\kappa =(U𝔥_0)\times 𝔥_0\times 𝔲_0.$$
$`(\mathrm{4.1.23})`$
###### Demonstration Proof
According to Proposition 3.14, we have to check that $`\tau ^{}`$ is a Lie algebra morphism. But from the orthogonality of $`𝔲_0`$ and $`i𝔲_0`$ under $`(,)_{}`$, we can show that $`\tau ^{}=\tau `$. As $`\tau `$ is a Cartan involution, the assertion follows. $`\mathrm{}`$
As can be easily verified, the real Hamiltonian system generated by $``$ has the same equations of motion as the one on $`U\times 𝔥\times 𝔤`$ with (complex holomorphic) Hamiltonian
$$^{}(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)\xi _\alpha \xi _\alpha .$$
$`(\mathrm{4.1.24})`$
This latter Hamiltonian system, on the other hand, has a compact real form, corresponding to $`q_i`$, $`p_i`$, $`i=1,\mathrm{},N`$ and $`\xi _\alpha =\overline{\xi }_\alpha `$, $`\alpha \mathrm{\Delta }`$. More precisely, the compact real form of $`^{}`$ is the Hamiltonian system on $`A^{}\mathrm{\Omega }^\kappa `$ generated by
$$\begin{array}{cc}\hfill \stackrel{~}{}(q,p,\xi )=& (A^{}\mathrm{\Omega }^\kappa )(q,p,\xi )\hfill \\ \hfill =& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)|\xi _\alpha |^2.\hfill \end{array}$$
$`(\mathrm{4.1.25})`$
In the rest of the subsection, we shall consider the realization of this compact real form and its associated integrable model. To start with, we introduce the map
$$s:L𝔤^{}L𝔤^{},s(X)(z)=\tau (X(\overline{z}))=\underset{j}{}\tau (X_j)(z)^j.$$
$`(\mathrm{4.1.26})`$
###### Proposition 4.1.6
The map
$$\sigma :A\mathrm{\Gamma }A\mathrm{\Gamma },(q,p,X)(\overline{q},\overline{p},s(X))$$
$`(\mathrm{4.1.27})`$
is a Poisson involution with stable locus
$$A\mathrm{\Gamma }^\sigma =(U𝔥_0)\times i𝔥_0\times (L𝔤^{})^s$$
$`(\mathrm{4.1.28})`$
where
$$(L𝔤^{})^s=\{XL𝔤^{}X_{2j+1}𝔲_0,X_{2j}i𝔲_0\text{for all}j\}.$$
$`(\mathrm{4.1.29})`$
Hence $`(A^{}\mathrm{\Gamma },\sigma ^{})`$ is a symmetric coboundary dynamical Lie algebroid.
###### Demonstration Proof
It is clear from (4.1.19) that $`s`$ is an involutive Lie algebra anti-morphism. If $`Z𝔥^{}`$, we have $`\tau (Z)=\overline{Z}`$ from which it follows that $`s_𝔥^{}(Z)=\overline{Z}`$. Next, we show that $`s^{}=s.`$ To do so, take $`X`$, $`YL𝔤^{}`$, then
$`(s^{}(X),Y)_{L𝔤^{}}`$
$`=`$ $`{\displaystyle \underset{j}{}}(X_j,\tau (Y_{(j+1)})(1)^{j+1})_{}`$
$`=`$ $`{\displaystyle \underset{j}{}}((1)^j\tau (X_j),Y_{(j+1)})_{}`$
$`=`$ $`({\displaystyle \underset{j}{}}\tau (X_j)(z)^j,Y)_{L𝔤^{}}`$
$`=`$ $`(s(X),Y)_{L𝔤^{}},`$
as required. We are now ready to compute $`sR(q)s^{}`$ for $`qU`$. Instead of using the explicit expression in (4.1.13), we will do this using the relationship between $`R`$ and $`r`$. This is more illuminating as the property of $`R`$ should follow from that of $`r`$. To start with, we have
$`(R(q)s(X))(z)`$
$`=`$ $`{\displaystyle \frac{1}{2}}s(X)(z)+{\displaystyle \underset{k0}{}}{\displaystyle \frac{1}{k!}}({\displaystyle \frac{^kr}{z^k}}(q,z),(1)^k\tau (X_{(k+1)})1).`$
Therefore,
$`(sR(q)s^{})(X)(z)`$
$`=`$ $`\tau (R(q)s(X)(\overline{z}))`$
$`=`$ $`{\displaystyle \frac{1}{2}}X(z)+{\displaystyle \underset{0}{}}{\displaystyle \frac{1}{k!}}\tau ({\displaystyle \frac{^kr}{z^k}}(q,\overline{z}),(1)^k\tau (X_{(k+1)})1).`$
To simplify the above expression, note that $`\overline{(a,\tau (\xi ))}=(\tau (a),\xi )`$ for all $`a`$,$`\xi 𝔤`$. From this relation, we find
$$\tau (ab,\tau (\xi )1)=(\tau ^2(ab),\xi 1)$$
for all $`a`$,$`b`$,$`\xi 𝔤`$. Consequently,
$`\tau ({\displaystyle \frac{^kr}{z^k}}(q,\overline{z}),(1)^k\tau (X_{(k+1)})1)`$
$`=`$ $`(\tau ^2\left({\displaystyle \frac{^kr}{z^k}}(q,\overline{z})\right),(1)^kX_{(k+1)}1).`$
But from (4.1.12), we can verify that
$$\tau ^2\left(\frac{^kr}{z^k}(q,\overline{z})\right)=(1)^k\frac{^kr}{z^k}(\overline{q},z).$$
Substitute this into the above expression, we obtain
$`\tau ({\displaystyle \frac{^kr}{z^k}}(q,\overline{z}),(1)^k\tau (X_{(k+1)})1)`$
$`=`$ $`({\displaystyle \frac{^kr}{z^k}}(\overline{q},z),X_{(k+1)}1)`$
and hence that
$$(sR(q)s^{})(X)(z)=(R(\overline{q})X)(z).$$
Therefore, we can now conclude from Proposition 3.9 that the map $`\sigma `$ is a Poisson involution. $`\mathrm{}`$
###### Proposition 4.1.7
(a) $`\sigma \rho =\rho \kappa `$. (b) For all $`dT`$,
$$𝒜_d\sigma =\sigma 𝒜_d,𝒞_d\kappa =\kappa 𝒞_d.$$
(c) $`\{Z𝔥^{}s_𝔥^{}^{}(Z)=Z\}=𝔱.`$
###### Demonstration Proof
(a) For any $`(q,p,\xi )A^{}\mathrm{\Omega }`$, we have
$$\rho \kappa (q,p,\xi )=(\overline{q},\mathrm{\Pi }_𝔥^{}\tau (\xi ),\overline{p}+r_{}^{\mathrm{}}(\overline{q})\tau (\xi )).$$
On the other hand,
$$\sigma \rho (q,p,\xi )=(\overline{q},\overline{\mathrm{\Pi }_𝔥^{}\xi },\overline{p}+s(r_{}^{\mathrm{}}(q)\xi )).$$
Now, from the fact that $`\tau `$ preserves $`𝔥^{}`$, we have
$$\mathrm{\Pi }_𝔥^{}\tau (\xi )=\overline{\mathrm{\Pi }_𝔥^{}\xi }.$$
Therefore, it remains to show that $`s(r_{}^{\mathrm{}}(q)\xi )=r_{}^{\mathrm{}}(\overline{q})\tau (\xi ).`$ To do so, we invoke the explicit expression for $`r_{}^{\mathrm{}}(q)(\xi )`$, according to which
$`(r_{}^{\mathrm{}}(\overline{q})\tau (\xi ))(z)`$
$`=`$ $`d(z){\displaystyle \underset{i}{}}\tau (\xi )_ix_i+{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}\varphi _\alpha (\overline{q},z)e^{\frac{z}{12}\alpha (\overline{q})}(\tau (\xi ))_\alpha e_\alpha `$
$`=`$ $`d(z){\displaystyle \underset{i}{}}\overline{\xi }_ix_i{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}\varphi _\alpha (\overline{q},z)e^{\frac{z}{12}\alpha (\overline{q})}\overline{\xi }_\alpha e_\alpha .`$
But on the other hand,
$`s(r_{}^{\mathrm{}}(q)\xi )(z)`$
$`=`$ $`d(z){\displaystyle \underset{i}{}}\overline{\xi }_i\tau (x_i){\displaystyle \underset{\alpha \mathrm{\Delta }}{}}\varphi _\alpha (\overline{q},z)e^{\frac{z}{12}\alpha (\overline{q})}\overline{\xi }_\alpha \tau (e_\alpha )`$
$`=`$ $`d(z){\displaystyle \underset{i}{}}\overline{\xi }_ix_i{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}\varphi _\alpha (\overline{q},z)e^{\frac{z}{12}\alpha (\overline{q})}\overline{\xi }_\alpha e_\alpha `$
and so we have equality.
(b) From the definition of $`\kappa `$ and $`𝒞`$, it is easy to see that $`𝒞_d\kappa =\kappa 𝒞_d`$ for $`dT`$ if and only if $`\tau (Ad_d\xi )=Ad_d\tau (\xi )`$ for $`dT.`$ But the latter is clear as $`\tau =1`$ on $`𝔲_0`$ and $`i𝔥_0𝔲_0`$. The validity of the other assertion also boils down to the same condition above, as can be verified from the definition of $`s`$.
(c) From the proof of Proposition 4.1.6, we have $`s_𝔥^{}(Z)=\overline{Z}`$ for $`Z𝔥^{}`$ from which the assertion clearly follows. $`\mathrm{}`$
Combining Proposition 4.1.3 and Proposition 4.1.7, we conclude that assumptions A1-A4 are satisfied. Hence Proposition 3.16 and the explicit form of $`\rho `$ give the following result.
###### Corollary 4.1.8
The stable loci $`A^{}\mathrm{\Omega }^\kappa `$, $`A\mathrm{\Gamma }^\sigma `$ are Hamiltonian $`T`$-spaces with equivariant momentum maps $`\stackrel{~}{J}:A^{}\mathrm{\Omega }^\kappa 𝔱,(q,p,\xi )\mathrm{\Pi }_𝔥^{}\xi ,`$ and $`\stackrel{~}{\gamma }=\gamma A\mathrm{\Gamma }^\kappa `$ respectively. Moreover, the map
$$\stackrel{~}{\rho }=\rho A^{}\mathrm{\Omega }^\kappa :(U𝔥_0)\times 𝔥_0\times 𝔲_0(U𝔥_0)\times i𝔥_0\times (L𝔤^{})^s$$
$`(\mathrm{4.1.30})`$
is a $`T`$-equivariant Poisson map and $`\stackrel{~}{\rho }(\stackrel{~}{J}^1(0))\stackrel{~}{\gamma }^1(0).`$
Remark 4.1.9. At this juncture, it is tempting to reduce $`\stackrel{~}{\rho }`$ further to the map in (4.2.5) in the next subsection. However, it is not hard to show that this cannot be achieved as we cannot find appropriate Poisson involutions which commute with $`\stackrel{~}{\rho }`$. We now introduce $`L=Pr_3\stackrel{~}{\rho }`$, as in (3.30). Also, let $`\stackrel{~}{Q}=Q(L𝔤^{})^s.`$
###### Proposition 4.1.10
(a) $`\stackrel{~}{}=L^{}\stackrel{~}{Q}`$ and is invariant under the canonical $`T`$-action on $`A^{}\mathrm{\Omega }^\kappa `$. (b) The restriction of the Hamiltonian equations of motion generated by $`\stackrel{~}{}`$ on $`A^{}\mathrm{\Omega }^\kappa `$ to the invariant submanifold $`\stackrel{~}{J}^1(0)`$ are given by
$$\begin{array}{cc}& \dot{q}=p,\hfill \\ & \dot{p}=\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\frac{\mathrm{cot}\frac{1}{2}\alpha (q)}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}|\xi _\alpha |^2H_\alpha ,\hfill \\ & \dot{\xi }=[\xi ,\frac{1}{4}\underset{\alpha \mathrm{\Delta }}{}\frac{\xi _\alpha }{\mathrm{sin}^2\frac{1}{2}\alpha (q)}e_\alpha ].\hfill \end{array}$$
$`(\mathrm{4.1.30})`$
Moreover, under the Hamiltonian flow, we have
$$\dot{L}(q,p,\xi )=[L(q,p,\xi ),R(q)M(q,p,\xi )]$$
$`(\mathrm{4.1.31})`$
on the invariant submanifold $`\stackrel{~}{J}^1(0)`$ where
$$M(q,p,\xi )(z)=L(q,p,\xi )(z)/z.$$
$`(\mathrm{4.1.32})`$
###### Demonstration Proof
(a) The assertion is clear from Proposition 4.1.4.
(b) From the expression for the Poisson bracket in (3.13), the Hamiltonian equations of motion are given by $`\dot{q}=\frac{1}{2}(\delta _2+\overline{\delta _2})`$, $`\dot{p}=\frac{1}{2}(\delta _1+\overline{\delta _1})`$, and $`\dot{\xi }=[\xi ,\frac{1}{2}(\delta +\tau (\delta )].`$ Therefore, (4.1.30) follows by a direct computation. On the other hand, (4.1.31) is a consequence of the last equation in Theorem 3.17(b). $`\mathrm{}`$
We now restrict ourselves to a smooth component of $`\stackrel{~}{J}^1(0)/T=(U𝔥_0)\times 𝔥_0\times (𝔲_0𝔥^{})/T.`$ For this purpose, we introduce the following open submanifold of $`𝔲_0`$:
$$𝒰=\{\xi 𝔲_0\xi _{\alpha _i}=(\xi ,e_{\alpha _i})0,i=1,\mathrm{},N\}.$$
$`(\mathrm{4.1.33})`$
Clearly, $`𝒰`$ is dense in $`𝔲_0`$ and is stable under the $`T`$-action. Therefore, $`(U𝔥_0)\times 𝔥_0\times 𝒰`$ is a Poisson submanifold of $`A^{}\mathrm{\Omega }^\kappa `$, and we can check that the $`T`$-action on $`A^{}\mathrm{\Omega }^\kappa `$ induces a locally free Hamiltonian group action on $`(U𝔥_0)\times 𝔥_0\times 𝒰.`$ Consequently, the corresponding momentum map is given by the restriction of the one in Corollary 4.1.8. To simplify notation, we shall denote this momentum map also by $`\stackrel{~}{J}`$ so that $`\stackrel{~}{J}^1(0)=(U𝔥_0)\times 𝔥_0\times (𝒰𝔥^{}).`$ Now observe that under the $`T`$-action, all the isotropy subgroups of the elements of $`\stackrel{~}{J}^1(0)`$ are identical. Since $`T`$ is compact, it follows from the above observation that the orbit space $`\stackrel{~}{J}^1(0)/T=(U𝔥_0)\times 𝔥_0\times (𝒰𝔥^{}/T)`$ is a smooth manifold. We shall fix a branch of the argument function. Then the formula
$$h(\xi )=exp\left(i\underset{j,k=1}{\overset{N}{}}C_{kj}arg\xi _{\alpha _k}h_{\alpha _j}\right)$$
$`(\mathrm{4.1.34})`$
defines a map $`h:𝒰𝔥^{}T`$ where $`C=(C_{jk})`$ is the inverse of the Cartan matrix and $`h_{\alpha _j}=\frac{2}{(\alpha _j,\alpha _j)}H_{\alpha _j}`$, $`j=1,\mathrm{},N`$. Note that for $`\xi 𝒰𝔥^{}`$, the element $`Ad_{h(\xi )^1}\xi `$ is such that $`(Ad_{h(\xi )^1}\xi ,e_{\alpha _i})>0,i=1,\mathrm{},N.`$ We shall henceforth identify the reduced phase space $`\stackrel{~}{J}^1(0)/T=(U𝔥_0)\times 𝔥_0\times (𝒰𝔥^{}/T)`$ with $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$ under the identification map $`(q,p,[\xi ])(q,p,Ad_{h(\xi )^1}\xi )`$ where
$$(𝔲_0)_{red}=\{f𝒰𝔥^{}f_{\alpha _i}=(f,e_{\alpha _i})>0,i=1,\mathrm{},N\}.$$
$`(\mathrm{4.1.35})`$
By Poisson reduction, the reduced manifold $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$ has a unique Poisson structure where the factor $`(𝔲_0)_{red}`$ is equipped with the reduction (at 0) of the Lie-Poisson structure on $`𝒰`$ by the $`T`$-action. Moreover, the reduction of the Hamiltonian $`\stackrel{~}{}`$ on $`A^{}\mathrm{\Omega }^\kappa `$ is given by
$$\stackrel{~}{}_0(q,p,f)=\frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)|f_\alpha |^2.$$
$`(\mathrm{4.1.36})`$
Combining Corollary 4.1.8, Proposition 4.1.10, Theorem 3.17 together with our discussion above, we therefore obtain the following result.(See the discussion preceding (3.32) for notations.)
###### Theorem 4.1.11
The map $`\stackrel{~}{\rho }`$ in (4.1.30) restricted to $`(U𝔥_0)\times 𝔥_0\times 𝒰`$ induces a unique Poisson map $`\widehat{\rho }:(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}A\mathrm{\Gamma }_0^\sigma `$ such that
(a) functions $`_0=\widehat{\rho }^{}\overline{f}`$, $`fI(L𝔤)`$, Poisson commute in $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$ and provide a family of conserved quantities in involution for the Hamiltonian $`\stackrel{~}{}_0`$, (b) Under the Hamiltonian flow generated by $`\stackrel{~}{}_0`$ on $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$,
$$\dot{L}(q,p,f)=[L(q,p,f),R(q)M(q,p,f)+]$$
where
$$\frac{=\frac{i}{4}\underset{j,k}{}\frac{C_{kj}}{f_{\alpha _k}}}{\alpha \mathrm{\Delta }\alpha _k\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }\frac{Im(f_\alpha f_{\alpha _k\alpha })}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}h_{\alpha _j}.}$$
###### Demonstration Proof
Only (b) requires a proof. First of all, from (4.1.31) and the $`T`$-equivariance of $`\stackrel{~}{\rho }`$ and $`R`$, we have (for $`f=Ad_{h(\xi )^1}\xi `$):
$`\dot{L}(q,p,f)`$
$`=`$ $`[L(q,p,f),R(q)M(q,p,f)+T_{h(\xi )}l_{h(\xi )^1}{\displaystyle \frac{d}{dt}}h(\xi )].`$
Now, differentiating $`h(\xi )`$ yields
$$T_{h(\xi )}l_{h(\xi )^1}\frac{d}{dt}h(\xi )=i\underset{j,k}{}C_{kj}(arg\xi _{\alpha _k})^.h_{\alpha _j}.()$$
But
$`\dot{\xi }_{\alpha _k}`$ $`=([\xi ,{\displaystyle \frac{1}{4}}{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}{\displaystyle \frac{\xi _\alpha }{\mathrm{sin}^2\frac{1}{2}\alpha (q)}}e_\alpha ],e_{\alpha _k})`$
$`{\displaystyle \frac{={\displaystyle \frac{1}{4}}{\displaystyle }}{\alpha \mathrm{\Delta }}}`$
$`\alpha _k\alpha \mathrm{\Delta }N_{\alpha ,\alpha _k\alpha }{\displaystyle \frac{\xi _\alpha \xi _{\alpha _k\alpha }}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}}.`$
Therefore, upon dividing both sides of the above expression by $`e^{iarg\xi _{\alpha _k}}`$ and taking the imaginary part of both sides, we find
$$\frac{f_{\alpha _k}(arg\xi _{\alpha _k})^.=\frac{1}{4}}{\alpha \mathrm{\Delta }\alpha _k\alpha \mathrm{\Delta }N_{\alpha ,\alpha _k\alpha }\frac{Im(f_\alpha f_{\alpha _k\alpha })}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}}$$
where we have used the reality of $`N_{\alpha ,\beta }`$ and $`\alpha (q)`$ together with the fact that $`f=_{\alpha \mathrm{\Delta }}\xi _\alpha e^{i_km_\alpha ^karg\xi _{\alpha _k}}e_\alpha `$ (the $`m_\alpha ^k`$ are defined by $`\alpha =_km_\alpha ^ke_{\alpha _k}`$). Consequently, when we substitute this in (\*), the desired expression for $``$ follows. $`\mathrm{}`$
Remark 4.1.12. (a) The Hamiltonian $`\stackrel{~}{}_0`$ is in fact completely integrable in the sense of Liouville on generic symplectic leaves of the reduced phase space. The same remark also applies to the integrable spin Calogero-Moser systems in \[LX2\] for all simple Lie algebras. A unifying and representation independent method to establish the Liouville integrability of such systems for all simple Lie algebras will be given in a forthcoming paper. For $`𝔤=gl(N,)`$ with $`𝔥`$ taken to be the set of diagonal matrices in $`𝔤`$, a sketch of the proof will be given below. (b) For the rational dynamical r-matrix $`\frac{\mathrm{\Omega }}{z}+_{\alpha \mathrm{\Delta }}\frac{1}{\alpha (q)}e_\alpha e_\alpha `$ and the elliptic dynamical r-matrix $`\zeta (z)_ix_ix_i_{\alpha \mathrm{\Delta }}l(\alpha (q),z)e_\alpha e_\alpha `$ (here $`l(w,z)=\frac{\sigma (w+z)}{\sigma (w)\sigma (z)}`$), recall that we can associate the corresponding (complex holomorphic) spin Calogero-Moser systems \[LX2\]. We remark that the compact real forms of these Hamiltonian systems can also be treated in the same way. Indeed, with the corresponding $`r_{}^{\mathrm{}}(q)`$ and $`R(q)`$, our analysis above can be repeated and everything goes through just the same as before. Note that the explicit form of $`r(q,z)`$ is only used in checking
$$\tau ^2\left(\frac{^kr}{z^k}(q,\overline{z})\right)=(1)^k\frac{^kr}{z^k}(\overline{q},z)$$
and in verifying $`s(r_{}^{\mathrm{}}(q)\xi )=r_{}^{\mathrm{}}(\overline{q})\tau (\xi ).`$ Finally we remark that a version of the rational spin Calogero-Moser system similar to the compact real form of our rational case has been obtained in \[AKLM\] by reducing a free Hamiltonian system on a cotangent bundle. However, it is not clear how this method can be generalized to handle the elliptic case. In the remainder of the subsection, we shall give a brief sketch of the Liouville integrability for the reductive case where $`𝔤=gl(N,).`$ Indeed, it is easy to see that we can repeat the same analysis above for this case with $`𝔥`$ taken to be the set of all diagonal matrices in $`𝔤`$ and with the trigonometric dynamical r-matrix with spectral parameter
$$r(q,z)=\left(c(z)+\frac{1}{12}z\right)\underset{i}{}e_{ii}e_{ii}+\underset{ij}{}(c(z)+c(q_iq_j))e^{\frac{z}{12}(q_iq_j)}e_{ij}e_{ji}$$
where $`e_{ij}`$ is the $`N\times N`$ matrix with a $`1`$ in the $`(i,j)`$ entry and zeros elsewhere and $`c(z)`$ is as in (4.1.8). In this case, $`𝔥_0`$ and $`𝔲_0`$ are the subalgebras of $`𝔤^{}`$ consisting of real diagonal matrices and skew-Hermitian matrices respectively and we take $`U`$ to be a fixed connected component of $`\{q𝔥\mathrm{sin}(\frac{q_iq_j}{2})0\text{for all}ij\}`$. Moreover, $`𝔤_0=gl(N,)`$ and $`T`$ is the maximal torus of the unitary group $`U(N)`$ consisting of unitary diagonal matrices. For our analysis below, we also have to introduce the torus $`T^{}T`$ consisting of matrices of the form $`diag(1,e^{i\theta _2},\mathrm{},e^{i\theta _N})`$. Clearly, the above results for simple Lie algebras have obvious analogs in this case. In particular, we can obtain the Hamiltonian
$$\stackrel{~}{}(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{ij}{}\left(\frac{1}{\mathrm{sin}^2(\frac{q_iq_j}{2})}\frac{1}{3}\right)|\xi _{ij}|^2$$
on $`A^{}\mathrm{\Omega }^\kappa =(U𝔥_0)\times 𝔥_0\times 𝔲_0`$ and its associated realization map
$$\stackrel{~}{\rho }:(U𝔥_0)\times 𝔥_0\times 𝔲_0(U𝔥_0)\times i𝔥_0\times (L𝔤^{})^s$$
by Dirac reduction. For the Hamiltonian system generated by $`\stackrel{~}{}`$, we note that the restriction of its equations of motion to the invariant submanifold $`\stackrel{~}{J}^1(0)`$ are given by
$$\begin{array}{cc}& \dot{q}=p,\hfill \\ & \dot{p}=\frac{1}{4}\underset{ij}{}\frac{\mathrm{cot}(\frac{q_iq_j}{2})}{\mathrm{sin}^2(\frac{q_iq_j}{2})}|\xi _{ij}|^2e_{ii},\hfill \\ & \dot{\xi }=[\xi ,\frac{1}{4}\underset{ij}{}\frac{\xi _{ij}}{\mathrm{sin}^2(\frac{q_iq_j}{2})}e_{ij}].\hfill \end{array}$$
$`(\mathrm{4.1.37})`$
Thus the equations coincide exactly with the ones derived in \[H\] (cf. also \[HKS\] and \[NM\]) for the eigenphases and (essentially) the eigenvectors of the unitary Floquet operator $`F=e^{i\lambda V}e^{iH_0}`$ (as a function of $`\lambda `$) associated with a periodically kicked quantum system if we take the time variable in (4.1.37) to be the kick strength $`\lambda .`$ More importantly, we have the Lax operator $`L=Pr_3\stackrel{~}{\rho }`$ whose restriction to $`\stackrel{~}{J}^1(0)=(U𝔥_0)\times 𝔥_0\times (𝔲_0𝔥^{})`$ is given explicitly by
$$L(q,p,\xi )(z)=p+\underset{ij}{}(c(z)+c(q_iq_j))e^{\frac{z}{12}(q_iq_j)}\xi _{ij}e_{ij}$$
$`(\mathrm{4.1.38})`$
and we can establish Liouville integrability for the reduced Hamiltonian system using $`L`$. To do so, we introduce the open submanifold of $`𝔲_0`$:
$$𝒰=\{\xi 𝔲_0\xi _{i,i+1}0,i=1,\mathrm{},N\}.$$
$`(\mathrm{4.1.39})`$
Then $`𝒰`$ is dense in $`𝔲_0`$ and is stable under the $`T`$-action. Hence $`(U𝔥_0)\times 𝔥_0\times 𝒰`$ is a Poisson submanifold of $`A^{}\mathrm{\Omega }^\kappa `$ and the $`T`$-action on $`A^{}\mathrm{\Omega }^\kappa `$ induces a Hamiltonian group action on $`(U𝔥_0)\times 𝔥_0\times 𝒰.`$ Denote the momentum map of this action also by $`\stackrel{~}{J}`$. Then $`\stackrel{~}{J}^1(0)=(U𝔥_0)\times 𝔥_0\times (𝒰𝔥^{}).`$ Since the isotropy subgroups of the elements of $`\stackrel{~}{J}`$ under the $`T`$-action are all equal to the center of $`U(N)`$ ($`=\{e^{i\theta }I\}S^1`$) and $`T`$ is compact, we conclude that the orbit space $`\stackrel{~}{J}^1(0)/T=(U𝔥_0)\times 𝔥_0\times (𝒰𝔥^{}/T)`$ is a smooth manifold. Indeed, $`\stackrel{~}{J}^1(0)/T=\stackrel{~}{J}^1(0)/T^{}`$ and we can check that the action of $`T^{}`$ on $`\stackrel{~}{J}^1(0)`$ is free. Let
$$(𝔲_0)_{red}=\{f𝒰𝔥^{}f_{i,i+1}>0,i=1,\mathrm{},N1\}.$$
$`(\mathrm{4.1.40})`$
Note that for $`\xi 𝒰𝔥^{}`$, there exists unique $`h(\xi )T^{}`$ such that $`Ad_{h(\xi )^1}\xi (𝔲_0)_{red}.`$ We shall henceforth identify $`\stackrel{~}{J}^1(0)/T`$ with $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}.`$
Now the dimension of the generic symplectic leaves of $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$ is given by
$$2N+dim_{}𝔲_0N2(N1)=N(N1)+2.$$
$`(\mathrm{4.1.41})`$
Therefore, in order to show that the reduced Hamiltonian $`\stackrel{~}{}_0`$ is completely integrable in the sense of Liouville on these generic symplectic leaves, we have to exhibit $`1+\frac{N(N1)}{2}`$ functionally independent conserved quantities in involution. To do so, we define $`\stackrel{~}{L}`$ by $`\stackrel{~}{L}(q,p,f)(z)=Ad_{e^{\frac{z}{12}q}}L(q,,p,f)(z)`$ for $`(q,p,f)(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}.`$ Then the characteristic polynomial of $`\stackrel{~}{L}(q,p,f)(z)`$ has the form
$$det(\stackrel{~}{L}(q,p,f)(z)w)=\underset{r=0}{\overset{N}{}}\underset{k=0}{\overset{r}{}}I_{rk}(q,p,f)c(z)^kw^{Nr}.$$
$`(\mathrm{4.1.42})`$
Now observe that $`(\stackrel{~}{L}(q,p,f)(z))^{}=\stackrel{~}{L}(q,p,f)(\overline{z})`$. Hence the functions $`I_{r,2k}(q,p,f)`$, $`iI_{r,2k+1}(q,p,f)`$ are real valued and provide the conserved quantities in involution. Clearly, among these are the $`N1`$ Casimirs $`I_{2k,2k}`$ and $`iI_{2k+1,2k+1}`$ ,$`k1`$ (note that $`I_{11}=0`$ on $`(U𝔥_0)\times 𝔥_0\times (𝔲_0)_{red}`$). On the other hand, it follows from the explicit expression of $`\stackrel{~}{L}`$ that
$$Ad_{e^{iq}}\stackrel{~}{L}(i\mathrm{})=\stackrel{~}{L}(i\mathrm{}).$$
$`(\mathrm{4.1.43})`$
Consequently, we obtain the relations
$$\underset{k=0}{\overset{\left[\frac{r1}{2}\right]}{}}(1)^kI_{r,2k+1}(q,p,f)=0$$
$`(\mathrm{4.1.44})`$
for $`r=1,\mathrm{},N`$. Hence the total number of independent nontrivial integrals equals
$$1+\underset{r=2}{\overset{N}{}}(r1)=1+\frac{1}{2}N(N1),$$
$`(\mathrm{4.1.45})`$
as required.
## 4.2 Normal compact forms of some spin Calogero-Moser systems
It is clear from (4.1.13) that $`R`$ also induces a map $`U𝔥_0L(L𝔤_0,L𝔤_0)`$ which we will also denote by $`R`$.
###### Proposition 4.2.1
The map $`R:U𝔥_0L(L𝔤_0,L𝔤_0)`$ is a solution of the mDYBE for the pair $`(L𝔤_0,𝔥_0)`$ with $`c=\frac{1}{4}.`$ Moreover, for $`qU𝔥_0`$, $`\xi 𝔤_0`$, we have $`r_{}^\mathrm{\#}(q)\xi L𝔤_0.`$
###### Demonstration Proof
For $`qU𝔥_0`$, $`X`$, $`YL𝔤_0`$, the element $`(dR(q)()X,Y)`$ must lie in $`𝔥_0`$ because for $`Z𝔥_0`$, $`(dR(q)(Z)X,Y)`$. The other assertion is clear from (4.1.11) as $`\alpha (q)`$ for $`qU𝔥_0`$. $`\mathrm{}`$
We next introduce the trivial Lie groupoids
$$\mathrm{\Omega }=(U𝔥_0)\times G_0\times (U𝔥_0),\mathrm{\Gamma }=(U𝔥_0)\times LG_0\times (U𝔥_0).$$
$`(\mathrm{4.2.1})`$
By Proposition 4.2.1 and (2.9), we can equip the dual bundle $`A^{}\mathrm{\Gamma }(U𝔥_0)\times 𝔥_0\times L𝔤_0`$ of $`A\mathrm{\Gamma }`$ with the Lie algebroid structure associated to $`R:U𝔥_0L(L𝔤_0,L𝔤_0)`$. Hence its dual bundle $`A\mathrm{\Gamma }(U𝔥_0)\times 𝔥_0\times L𝔤_0`$ has a Lie-Poisson structure. We shall also equip the the dual bundle $`A^{}\mathrm{\Omega }(U𝔥_0)\times 𝔥_0\times 𝔤_0`$ of the trivial Lie algebroid $`A\mathrm{\Omega }`$ with the corresponding Lie-Poisson structure. The Poisson manifolds $`A^{}\mathrm{\Omega }`$ and $`A\mathrm{\Gamma }`$ are Hamiltonian $`H_0`$-spaces. Indeed, the actions are defined by expressions identical to (4.1.15) and (4.1.16) provided that we change $`H^{}`$ to $`H_0`$ and use the definitions of $`\mathrm{\Omega }`$ and $`\mathrm{\Gamma }`$ in (4.2.1). On the other hand, the corresponding equivariant momentum maps are given by
$$J:A^{}\mathrm{\Omega }𝔥_0,(q,p,\xi )\mathrm{\Pi }_{𝔥_0}\xi $$
$`(\mathrm{4.2.2})`$
and
$$\gamma :A\mathrm{\Gamma }𝔥_0,(q,p,X)p$$
$`(\mathrm{4.2.3})`$
where $`\mathrm{\Pi }_{𝔥_0}`$ is the projection map to $`𝔥_0`$ relative to the decomposition $`𝔤_0=𝔥_0(𝔥_0)^{}`$.
As in Propositions 4.1.3 and 4.1.4, we have the following result in this case.
###### Proposition 4.2.2
The map
$$\rho :A^{}\mathrm{\Omega }A\mathrm{\Gamma },(q,\mathrm{\Pi }_{𝔥_0}\xi ,p+r_{}^\mathrm{\#}(q)\xi )$$
$`(\mathrm{4.2.4})`$
is an $`H_0`$-equivariant Poisson map. Moreover, the map $`\rho `$ gives a realization of the spin Calogero-Moser system on $`A^{}\mathrm{\Omega }`$ with Hamiltonian
$$(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)\xi _\alpha \xi _\alpha $$
$`(\mathrm{4.2.5})`$
in $`A\mathrm{\Gamma }.`$
The Hamiltonian system in (4.2.5) will be called the normal real form of the complex holomorphic system $`^{}`$ in (4.1.24).
In the rest of the section, we shall reduce this normal real form to what we call the normal compact form. As the reader will see, the normal compact form has a natural family of conserved quantities in involution.
For this purpose, we introduce
$$𝔨_0=\underset{\alpha \mathrm{\Delta }}{}(e_\alpha e_\alpha ),𝔭_0=𝔥_0+\underset{\alpha \mathrm{\Delta }}{}(e_\alpha +e_\alpha ).$$
$`(\mathrm{4.2.6})`$
Then
$$𝔤_0=𝔨_0+𝔭_0$$
$`(\mathrm{4.2.7})`$
is a Cartan decomposition of $`𝔤_0`$. Let $`\theta `$ be the corresponding involution.
###### Proposition 4.2.3
The map
$$\kappa :A^{}\mathrm{\Omega }A^{}\mathrm{\Omega },(q,p,\xi )(q,p,\theta (\xi ))$$
$`(\mathrm{4.2.8})`$
is a Poisson involution with stable locus
$$A^{}\mathrm{\Omega }^\kappa =(U𝔥_0)\times 𝔥_0\times 𝔨_0.$$
$`(\mathrm{4.2.9})`$
###### Demonstration Proof
Since $`\theta `$ is a Cartan involution, it follows that $`𝔨_0`$ and $`𝔭_0`$ are orthogonal under $`(,)_{𝔤_0\times 𝔤_0}`$. Using this property, we can show that $`\theta ^{}=\theta .`$ Consequently, we conclude from Proposition 3.14 that $`\kappa `$ is a Poisson involution. $`\mathrm{}`$
As $`𝔨_0=𝔲_0𝔤_0`$, we shall call $`𝔨_0`$ the normal compact form of $`𝔤.`$ Note, however, that $`𝔨_0`$ is not at all a real form of $`𝔤`$ because its complexification is different from $`𝔤.`$ In view of this terminology, we define the normal compact form of $`^{}`$ to be the Hamiltonian system on $`A^{}\mathrm{\Omega }^\kappa `$ generated by
$$\begin{array}{cc}\hfill \stackrel{~}{}(q,p,\xi )=& (A^{}\mathrm{\Omega }^\kappa )(q,p,\xi )\hfill \\ \hfill =& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{\alpha \mathrm{\Delta }}{}\left(\frac{1}{\mathrm{sin}^2\frac{1}{2}\alpha (q)}\frac{1}{3}\right)|\xi _\alpha |^2.\hfill \end{array}$$
$`(\mathrm{4.2.10})`$
In order to discuss the realization of this Hamiltonian system, we introduce
$$s:L𝔤_0L𝔤_0,s(X)(z)=\underset{j}{}\theta (X_j)(z)^j.$$
$`(\mathrm{4.2.11})`$
###### Proposition 4.2.4
The map
$$\sigma :A\mathrm{\Gamma }A\mathrm{\Gamma },(q,p,X)(q,p,s(X))$$
$`(\mathrm{4.2.12})`$
is a Poisson involution with stable locus
$$A\mathrm{\Gamma }^\sigma =(U𝔥_0)\times \{0\}\times (L𝔤_0)^s$$
$`(\mathrm{4.2.13})`$
where
$$(L𝔤_0)^s=\{XL𝔤_0X_{2j+1}𝔨_0,X_{2j}𝔭_0\text{for all}j\}.$$
$`(\mathrm{4.2.14})`$
Consequently, $`Pr_3^{}I((L𝔤_0)^s)`$ is a Poisson commuting family of functions on $`A\mathrm{\Gamma }^\sigma `$.
###### Demonstration Proof
Since $`\theta =1`$ on $`𝔭_0`$ and $`𝔥_0𝔭_0`$, we have $`s_{𝔥_0}=\theta _{𝔥_0}=id_{𝔥_0}.`$ Therefore, $`s_{𝔥_0}^{}=id_{𝔥_0}`$. The rest of the proof of the first assertion is similar to the one for Proposition 4.1.6. On the other hand, the second assertion is just a consequence of (4.2.13) and Theorem 3.15. $`\mathrm{}`$
Our next result shows that assumption A1 is satisfied. Its proof is similar to the one for Proposition 4.1.7.
###### Proposition 4.2.5
$`\sigma \rho =\rho \kappa `$.
From this proposition and (4.2.13), we see that the assumptions in Theorem 3.15 are satisfied. Hence we obtain the following result.
###### Theorem 4.2.6
(a) The map
$$\stackrel{~}{\rho }=\rho A^{}\mathrm{\Omega }^\kappa :(U𝔥_0)\times 𝔥_0\times 𝔨_0(U𝔥_0)\times \{0\}\times (L𝔤_0)^s$$
$`(\mathrm{4.2.15})`$
is a Poisson map. (b) $`\stackrel{~}{}=L^{}\stackrel{~}{Q}`$ and admits $`L^{}I((L𝔤_0)^s)`$ as a family of conserved quantities in involution. Here, $`L=Pr_3\stackrel{~}{\rho }`$ and $`\stackrel{~}{Q}=Q(L𝔤_0)^s.`$
Remark 4.2.7. (a) The Hamiltonian system generated by $`\stackrel{~}{}`$ is completely integrable in the sense of Liouville on generic symplectic leaves of $`A^{}\mathrm{\Omega }^\kappa `$. This will also be treated in the forthcoming paper which we mentioned in the previous subsection. (b) The normal compact forms of the rational and the elliptic spin calogero-Moser systems (corresponding to the dynamical r-matrices in Remark 4.1.12(b)) can also treated in a similar fashion. Note that for the elliptic case, in order for the corresponding $`R:UL(L𝔤,L𝔤)`$ to induce a map $`U𝔥_0L(L𝔤_0,L𝔤_0),`$ we have to make an additional assumption, namely, we have to restrict to periods $`2\omega _1`$, $`2\omega _2`$ (of the elliptic functions) for which the invariants $`g_2=60_{\omega \mathrm{\Lambda }\backslash \{0\}}\omega ^4`$ and $`g_3=140_{\omega \mathrm{\Lambda }\backslash \{0\}}\omega ^6`$ are real, where $`\mathrm{\Lambda }=2\omega _1+2\omega _2`$. As in Section 4.1, we shall close this subsection with a sketch of the Liouville integrability for $`𝔤=gl(N,)`$. In this case, $`U`$, $`𝔥_0`$ and $`r(q,z)`$ are the same objects which appear at the end of section 4.1 and we have $`𝔤_0=gl(N,)`$. Thus the factors $`𝔨_0`$ and $`𝔭_0`$ in the Cartan decomposition are respectively the set of skew-symmetric matrices and symmetric matrices in $`𝔤_0`$. Clearly the results above for simple Lie algebras have obvious analogs in this case. In particular, for the Hamiltonian system on $`A^{}\mathrm{\Omega }^\kappa =(U𝔥_0)\times 𝔥_0\times 𝔨_0`$ generated by
$$\stackrel{~}{}(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{ij}{}\left(\frac{1}{\mathrm{sin}^2(\frac{q_iq_j}{2})}\frac{1}{3}\right)|\xi _{ij}|^2,$$
$`(\mathrm{4.2.16})`$
its Hamiltonian equations of motion are given by the same expressions in (4.1.37) but with $`\xi 𝔨_0`$. In this case, the equations are associated with an orthogonal Floquet operator $`F=e^{i\lambda V}e^{iH_0}`$ and the Lax operator $`L=Pr_3\stackrel{~}{\rho }`$ on $`A^{}\mathrm{\Omega }^\kappa `$ has the same form as the one in (4.1.38) but with $`\xi 𝔨_0`$. Now the dimension of the generic symplectic leaves of $`A^{}\mathrm{\Omega }^\kappa `$ is given by $`2N+\frac{N(N1)}{2}\left[\frac{N}{2}\right]=2N+2\left[\frac{(N1)^2}{4}\right].`$ In order to show $`\stackrel{~}{}`$ is completely integrable in the sense of Liouville on these leaves, we have to exhibit $`N+\left[\frac{(N1)^2}{4}\right]`$ functionally independent conserved quantities in involution. Put $`\stackrel{~}{L}(q,p,\xi )(z)=Ad_{e^{\frac{z}{12}q}}L(q,,p,\xi )(z),`$ then it is easy to check that $`(\stackrel{~}{L}(q,p,\xi )(z))^T=\stackrel{~}{L}(q,p,\xi )(z)`$. Therefore the characteristic polynomial of $`L(q,p,\xi )(z)`$ is an even function of $`z`$. Hence we have
$$det(\stackrel{~}{L}(q,p,\xi )(z)w)=\underset{r=0}{\overset{N}{}}\underset{k=0}{\overset{\left[\frac{r}{2}\right]}{}}I_{rk}(q,p,\xi )c(z)^{2k}w^{Nr}$$
$`(\mathrm{4.2.17})`$
and the $`I_{rk}`$’s are conserved quantities in involution. Clearly, the functions $`I_{2k,k},`$$`k=1,\mathrm{},\left[\frac{N}{2}\right]`$, are Casimirs. Therefore, the total number of nontrivial integrals is given by
$$\underset{r=1}{\overset{N}{}}\left(\left[\frac{r}{2}\right]+1\right)\left[\frac{N}{2}\right]=N+\left[\frac{(N1)^2}{4}\right].$$
$`(\mathrm{4.2.18})`$
## 5. Symmetric space spin Ruijsenaars-Schneider models and soliton dy- fak namics of affine Toda field theory
There is a well-known correspondence between the $`N`$-soliton solutions of the $`A_n^{(1)}`$ affine Toda field theory and some spin-generalized Ruijensaars-Schneider equations \[BH\]. The goal of this section is to resolve a long-standing problem regarding the Hamiltonian formulation and the integrability of such equations.
Let $`𝔤=gl(N,)`$, and let $`𝔥`$ be the Cartan subalgebra of $`𝔤`$ consisting of diagonal matrices. We shall denote by $`𝔤^{}`$ (resp. $`𝔥^{}`$) the algebra $`𝔤`$ (resp. $`𝔥`$) regarded as a real Lie algebra. It is well-known that
$$𝔲(N)=\{N\times N\text{skew-Hermitian matrices }\}$$
$`(5.1)`$
is a compact real form of $`𝔤`$. We shall denote by $`\tau `$ the conjugation of $`𝔤`$ with respect to $`𝔲(N)`$. (Explicitly, $`\tau (\xi )=\xi ^{}`$ for $`\xi 𝔤.`$) Clearly, the map
$$s=\tau :𝔤^{}𝔤^{}$$
$`(5.2)`$
is an involutive Lie algebra anti-morphism satisfying $`s(𝔥)=𝔥.`$ In the following, the connected and simply-connected Lie groups which integrate $`𝔤^{}`$, $`𝔥^{}`$ will be denoted by $`G^{}`$ and $`H^{}`$ respestively. Then the Lie group anti-morphism $`S:G^{}G^{}`$ corresponding to $`s`$ is given by $`S(g)=g^{}`$ for $`gG^{}`$.
In what follows, $`U`$ will denote a fixed connected component of
$$\{q=diag(q_1,\mathrm{}q_N)𝔥\mathrm{sinh}\left(\frac{1}{2}(q_iq_j)\right)0\text{for all i and j}\}.$$
We consider the solution $`R:UL(𝔤,𝔤)`$ of the mDYBE, given by
$$R(q)\xi =\frac{1}{2}\underset{ij}{}\mathrm{coth}\left(\frac{1}{2}(q_iq_j)\right)\xi _{ij}e_{ij},$$
$`(5.3)`$
From this formula, it is clear that $`R`$ induces a map $`UL(𝔤^{},𝔤^{})`$ which is a solution of the mDYBE for the pair $`(𝔤^{},𝔥^{})`$. We shall denote this map also by $`R`$ and from now onwards we shall only consider $`R`$ as a map $`UL(𝔤^{},𝔤^{})`$. We now equip the trivial Lie groupoid
$$\mathrm{\Gamma }=U\times G^{}\times U$$
$`(5.4)`$
with the coboundary dynamical Poisson structure associated to $`R`$. Since $`H^{}`$ is abelian, its action on $`\mathrm{\Gamma }`$ is given by
$$:H^{}\times \mathrm{\Gamma }\mathrm{\Gamma },_h(u,g,v)=(u,hgh^1,v).$$
$`(5.5)`$
###### Proposition 5.1
The map
$$\mathrm{\Sigma }:(\mathrm{\Gamma },\{,\}_R)(\mathrm{\Gamma },\{,\}_R),(u,g,v)(\overline{v},g^{},\overline{u})$$
$`(5.6)`$
is a Poisson involution with stable locus
$$\mathrm{\Gamma }^\mathrm{\Sigma }=\{(u,g,\overline{u})\mathrm{\Gamma }g=g^{}\}.$$
$`(5.7)`$
Hence $`(\mathrm{\Gamma },\{,\}_R,\mathrm{\Sigma })`$ is a symmetric coboundary dynamical Poisson groupoid.
###### Demonstration Proof
Using the pairing $`(\xi ,\eta )_{}=2Retr(\xi \eta )`$ on $`𝔤^{}`$, it is straight forward to show that $`s_𝔥^{}^{}(u)=s_𝔥^{}(u)=\overline{u}`$ for $`uU`$. By a similar calculation, we also have $`s^{}=s`$. From this, it is easy to show that $`sR(q)s^{}=R(\overline{q})`$. Hence it follows from Proposition 3.12 that $`\mathrm{\Sigma }`$ is a Poisson involution. $`\mathrm{}`$
Let $`T`$ be the subgroup of $`H^{}`$ consisting of unitary diagonal matrices and let $`𝔱=Lie(T)`$.
###### Proposition 5.2
(a) For all $`dT`$, $`_d\mathrm{\Sigma }=\mathrm{\Sigma }_d.`$ (b) $`q\overline{q}𝔱`$ for all $`qU.`$
###### Demonstration Proof
(a) From the definition of $``$ and $`\mathrm{\Sigma }`$, we have $`_d\mathrm{\Sigma }=\mathrm{\Sigma }_d`$ for $`dT`$ iff $`dg^{}d^{}=(dg^{}d^{})^{}`$. But the latter is obvious. (b) This assertion is clear. $`\mathrm{}`$
Since we are dealing with the case in which the realization map is the identity map on $`\mathrm{\Gamma },`$ it follows from Propositions 5.1 and 5.2 that assumptions G1-G4 in Section 3 are satisfied. Hence we have the following result by Proposition 3.19.
###### Corollary 5.3
The stable locus $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ is a Hamiltonian $`T`$-space with equivariant momentum map $`\stackrel{~}{\alpha }\stackrel{~}{\beta }=\alpha \beta \mathrm{\Gamma }^\mathrm{\Sigma }`$.
###### Definition Definition 5.4
The spin Ruijsenaars-Schneider models associated to $`R`$ are the Hamiltonian systems on $`\mathrm{\Gamma }`$ generated by functions in $`𝙿r_2^{}I(G^{})`$. The symmetric space spin Ruijsenaars-Schneider models are the corresponding Hamiltonian systems on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ generated by functions in $`𝙿r_2^{}I((G^{})^S).`$ Here we have used the same symbol $`𝙿r_2`$ to denote the projection map from $`\mathrm{\Gamma }`$ to $`G^{}`$ and its restriction from $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ to $`(G^{})^S.`$
Note that in the case under consideration, we have $`(\stackrel{~}{\alpha }\stackrel{~}{\beta })^1(0)U_s\times (G^{})^S,`$ where $`U_s`$ consists of real diagonal matrices in $`U`$ and $`(G^{})^S`$ consists of Hermitian matrices in $`G^{}.`$ As we are identifying $`(𝔤^{})^{}`$ with $`𝔤^{}`$ using the pairing $`(,)_{}`$, it follows from (3.7) that the Poisson structure on $`\mathrm{\Gamma }^\mathrm{\Sigma }`$ is given by
$$\begin{array}{cc}\hfill \{\stackrel{~}{\phi },\stackrel{~}{\psi }\}_{\mathrm{\Gamma }^\mathrm{\Sigma }}(u,g,\overline{u})=& 2(\iota \delta _1\stackrel{~}{\phi },D\stackrel{~}{\psi })_{}+2(\iota \delta _1\stackrel{~}{\psi },D\stackrel{~}{\phi })_{}2(R(u)D\stackrel{~}{\phi },D\stackrel{~}{\psi })_{}\hfill \end{array}$$
$`(5.8)`$
where
$$\delta _1\stackrel{~}{\phi }:=\frac{1}{2}(\delta _1\phi +\overline{\delta _2\phi }),D\stackrel{~}{\phi }:=\frac{1}{2}(D\phi +(D^{}\phi )^{}).$$
$`(5.9)`$
Hence we obtain the following result by applying Proposition 3.13.
###### Proposition 5.5
Let $`f(g)=2Retr(g)`$, $`gG^{}`$ and let $`F=Pr_2^{}f`$. Then the restriction of the Hamiltonian equations of motion generated by $`\stackrel{~}{F}`$ to the invariant manifold $`U_s\times (G^{})^S`$ are given by
$`\dot{q}=\mathrm{\Pi }_𝔥^{}g,`$
$`\dot{g}=g(R(q)g)(R(q)g)g.`$
In terms of the components of $`q`$ and $`g`$, these read:
$$\ddot{q}_i=\dot{g}_{ii}=\frac{1}{2}\underset{ki}{}\mathrm{coth}((q_iq_k)/2)g_{ik}g_{ki},$$
$`(5.10a)`$
$$\begin{array}{cc}\hfill \dot{g}_{ij}=& \frac{1}{2}\mathrm{coth}((q_iq_j)/2)g_{ij}(g_{jj}g_{ii})\hfill \\ & +\frac{1}{2}\underset{ki,j}{}\left(\mathrm{coth}((q_iq_k)/2)\mathrm{coth}((q_kq_j)/2)\right)g_{ik}g_{kj},ij\hfill \end{array}$$
$`(5.10b)`$
Note that the equations in (5.10) for some special choice of $`g`$ are exactly the ones derived by Braden and Hone in \[BH\] from the $`N`$-soliton solutions of the $`A_n^{(1)}`$ affine Toda field theory with purely imaginary coupling constant (cf.\[KZ\]). For the convenience of the reader, let us recall that the equations of motion of the $`A_n^{(1)}`$ affine Toda field theory (with imaginary coupling constant $`i\beta `$) are given by
$$_+_{}\varphi _j+\frac{m^2}{2i\beta }\left(e^{i\beta (\varphi _j\varphi _{j+1})}e^{i\beta (\varphi _{j1}\varphi _j)}\right)=0,j=0,\mathrm{},n.$$
$`(5.11)`$
Here, $`_\pm `$ denotes differentiation with respect to the light-cone coordinates $`x_\pm =\frac{1}{\sqrt{2}}(t\pm x)`$ and the indices $`j`$ are taken modulo $`n+1`$. In \[BH\], the authors were dealing with the solitonic sector of the theory, so they assumed in addition that $`_{j=0}^n\varphi _j=0`$. Starting from the $`N`$-soliton solutions of (5.10) as derived by Hollowood \[H\]:
$$e^{i\beta \varphi _j}=\frac{\tau _{j+1}}{\tau _j},$$
$`(5.12)`$
where $`\tau _j`$ are the tau functions, Braden and Hone began by rewriting $`\tau _j`$ in determinantal form. As it turned out, they found that
$$\tau _j=det\left(1+e^{ij\mathrm{\Theta }/2}Ve^{ij\mathrm{\Theta }/2}\right)$$
$`(5.13)`$
where $`V`$ is an invertible skew-Hermitian matrix depending on $`x_\pm `$ (and the $`2N`$ soliton parameters) satisfying
$$\dot{V}=\frac{1}{2}(\mathrm{\Lambda }V+V\mathrm{\Lambda }),=_\pm $$
$`(5.14)`$
In (5.13) and (5.14), $`\mathrm{\Theta }=diag(\theta _1,\mathrm{},\theta _N)`$ and
$$\mathrm{\Lambda }=diag(\pm \sqrt{2}m\mathrm{exp}(\eta _1)sin(\theta _1/2),\mathrm{},\pm \sqrt{2}m\mathrm{exp}(\eta _N)sin(\theta _N/2))$$
$`(5.15)`$
where $`\theta _j`$ are discrete parameters associated with the solitons taking values in $`\{\frac{2\pi k}{n+1}k=1,\mathrm{},n\}`$, and $`\eta _j`$ are the rapidities. As $`V`$ is skew-Hermitian and invertible, there exists a unitary matrix $`U`$ (unique up to transformations $`U\delta U`$, where $`\delta T`$) which diagonalizes $`V`$, i.e.
$$ie^q=UVU^{}$$
$`(5.16)`$
where $`q`$ is real diagonal. Using $`U`$, define an invertible Hermitian matrix $`g`$ by:
$$g=U\mathrm{\Lambda }U^{}.$$
$`(5.17)`$
Then under the evolution as defined by (5.14), the authors in \[BH\] showed that $`q`$ and $`g`$ satisfy (5.10)! Of course, in this context, the variable $`g`$ depends on the choice of $`U`$. Clearly, the system which is independent of such a choice is the corresponding reduced system on the Poisson quotient $`U_s\times ((G^{})^S/T).`$ As a consequence of Theorem 3.20, we therefore conclude that the reduced system has $`N`$ Poisson commuting integrals. Remark 5.7. There are several questions which we have not addressed in our discussion above. One of these has to do with the nature of the transformation between the $`2N`$ soliton parameters and the dynamical variables in $`U_s\times ((G^{})^S/T).`$ On the other hand, there should be corresponding results for the soliton solutions associated with the other affine Lie algebras in \[OTU\]. Of course, there is also the question of Liouville integrability. We hope to return to these questions in future work.
|
warning/0506/hep-th0506082.html
|
ar5iv
|
text
|
# Multi-(super)graviton theory on topologically non-trivial backgrounds
## 1 Introduction
Renewed interest in the study of multi-graviton theories owes, in particular, to the fact that these formulations resemble higher-dimensional gravities in the presence of discrete dimensions. These classes of discretized Kaluza-Klein theories are now in fact under the focus of attention due to their primary importance for the realization of the dimensional deconstruction program . Moreover, multigravitons can be also related with discretized brane-world models .
In spite of the absence of a consistent interaction among the gravitons, one can think of possible couplings in the theory space. In particular, in a recent paper , a multi-graviton theory with nearest-neighbor couplings in the theory space has been proposed. As a result, a discrete mass spectrum appears. The theory seems to be equivalent to Kaluza-Klein gravity with a discretized dimension.
In a previous paper concerning multi-graviton theory , we have shown by means of an explicit example, namely a discretized Randall-Sundrum (RS) brane-world , that the induced cosmological constant becomes positive provided the number of massive gravitons is sufficiently large.
In the present paper, we would like to show that an alternative mechanisms can also give rise to positive contributions to the cosmological constant. In particular we shall consider a multi-supergraviton example with few gravitons, in a manifold (bulk) with non trivial topology. We shall show that in such a model a positive cosmological constant $`\mathrm{\Lambda }`$ can be generated, due to the presence of positive topological contributions. Moreover, by a suitable choice of the topological parameters, the number obtained for $`\mathrm{\Lambda }`$ can reach a value perfectly in accordance with result obtained from recent cosmological observations .
## 2 The multi-graviton and multi-supergraviton models
### 2.1 The graviton model
We start by considering the Lagrangian for the spin-two field $`h_{\mu \nu }`$ with mass $`m`$
$`_m`$ $`=`$ $`_0{\displaystyle \frac{m^2}{2}}\left(h_{\mu \nu }h^{\mu \nu }h^2\right)2\left(mA^\mu +^\mu \phi \right)\left(^\nu h_{\mu \nu }_\mu h\right)`$ (2.1)
$`{\displaystyle \frac{1}{2}}\left(_\mu A_\nu _\nu A_\mu \right)\left(^\mu A^\nu ^\nu A^\mu \right),`$
where $`_0`$ is the Lagrangian of the massless spin-two field (graviton) $`h_{\mu \nu }`$ ($`hh_\mu ^\mu `$)
$`_0={\displaystyle \frac{1}{2}}_\lambda h_{\mu \nu }^\lambda h^{\mu \nu }+_\lambda h_\mu ^\lambda _\nu h^{\mu \nu }_\mu h^{\mu \nu }_\nu h+{\displaystyle \frac{1}{2}}_\lambda h^\lambda h,`$ (2.2)
while $`A_\mu `$ and $`\phi `$ are Stückelberg fields .
The multi-graviton model is defined by taking $`N`$copies of (2.1) with graviton $`h_{n\mu \nu }`$ and Stückelberg fields $`A_{n\mu }`$ and $`\phi _n`$. Here, we propose a theory defined by a Lagrangian which is taken to be a generalization of the one in . It reads
$``$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{N1}{}}}[{\displaystyle \frac{1}{2}}_\lambda h_{n\mu \nu }^\lambda h_n^{\mu \nu }+_\lambda h_{n\mu }^\lambda _\nu h_n^{\mu \nu }_\mu h_n^{\mu \nu }_\nu h_n+{\displaystyle \frac{1}{2}}_\lambda h_n^\lambda h_n`$ (2.3)
$`{\displaystyle \frac{1}{2}}\left(m^2\mathrm{\Delta }h_{n\mu \nu }\mathrm{\Delta }h_n^{\mu \nu }\left(\mathrm{\Delta }h_n\right)^2\right)2\left(m\mathrm{\Delta }^{}A_n^\mu +^\mu \phi _n\right)\left(^\nu h_{n\mu \nu }_\mu h_n\right)`$
$`{\displaystyle \frac{1}{2}}(_\mu A_{n\nu }_\nu A_{n\mu })(^\mu A_n^\nu ^\nu A_n^\mu )].`$
The $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ are difference operators, which operate on the indices $`n`$ as
$`\mathrm{\Delta }\varphi _n{\displaystyle \underset{k=0}{\overset{N1}{}}}a_k\varphi _{n+k},\mathrm{\Delta }^{}\varphi _n{\displaystyle \underset{k=0}{\overset{N1}{}}}a_k\varphi _{nk},{\displaystyle \underset{k=0}{\overset{N1}{}}}a_k=0,`$ (2.4)
where the $`a_k`$ are $`N`$ constants and the $`N`$ variables $`\varphi _n`$ can be identified with periodic fields on a lattice with $`N`$ sites if the periodic boundary conditions $`\varphi _{n+N}=\varphi _n`$ are imposed. The latter condition in (2.4) assures that $`\mathrm{\Delta }`$ becomes the usual differentiation operator in a properly defined continuum limit.
The eigenvalues and eigenvectors for $`\mathrm{\Delta }`$ are given by<sup>1</sup><sup>1</sup>1Please note that here we use a different notation with respect to one used in Refs. and . In fact, in order to avoid confusion with masses, we have replaced the eigenvalue $`m`$ with $`\mu `$ and the index $`M`$ with $`p`$.
$`\mathrm{\Delta }\varphi _n^p=i\mu _p\varphi _n^p,i\mu _p={\displaystyle \underset{n=0}{\overset{N1}{}}}a_ne^{2\pi inp/N},`$ (2.5)
$`\varphi _n^p={\displaystyle \frac{e^{2\pi inp/N}}{\sqrt{N}}}.p=0,1,2,3,\mathrm{}`$ (2.6)
By using (2.4) in the latter equation and assuming $`a_n`$ to be real one gets the relations
$`\mu _0=0,\mu _p=\mu _{Np},\mu _{Np}=\mu _p^{},`$ (2.7)
which, for any fixed $`N`$, permit to obtain the whole spectrum of the theory.
Then we see that the Lagrangian (2.3) describes a massless graviton and $`N1`$ massive gravitons, with masses $`M_p=|\mu _p|`$ ($`p=1,2,\mathrm{},N1`$). It must be pointed out that the massive gravitons always appear in pairs which share a common mass and, moreover, the complex mass parameter $`\mu _p`$ can be arbitrarily chosen, just by properly selecting the coefficients $`a_k`$ in (2.5) . As discussed in , the multigraviton model can be regarded as corresponding to a Kaluza-Klein theory where the extra dimension lives in a lattice.
As a specific example, we now consider the two-brane Randall-Sundrum (RS) model (for a recent review see ). In this model, the masses of the Kaluza-Klein modes are given by
$`M_p={\displaystyle \frac{\pi p}{z_c}},z_c=l\left(e^{\pi r_c/l}1\right),p=0,1,2,\mathrm{}`$ (2.8)
where $`l`$ is the length parameter of the five-dimensional AdS space and $`\pi r_c`$ the geodesic distance between the two branes.
Motivated by this last equation (2.8), we consider an $`N=2N^{}+1`$ graviton model, with the graviton masses being given by
$`\mu _p=\{\begin{array}{cc}\frac{\pi p}{z_c},\hfill & p=0,1,\mathrm{},N^{},\hfill \\ \frac{\pi (Np)}{z_c},\hfill & p=N^{}+1,N^{}+2,\mathrm{},N1=2N^{}.\hfill \end{array}`$ (2.10)
Those are solutions of Eq. (2.5), with the choice $`a_0=0`$ and, for any $`n1`$,
$`a_n`$ $`=`$ $`{\displaystyle \frac{2\pi }{(2N^{}+1)z_c}}Im\left\{{\displaystyle \frac{\left(1e^{i\frac{2\pi N^{}n}{2N^{}+1}}\right)e^{i\frac{2\pi n}{2N^{}+1}}}{1e^{i\frac{2\pi n}{2N^{}+1}}}}\right\}`$ (2.11)
$`=`$ $`{\displaystyle \frac{(1)^n\mathrm{\hspace{0.25em}2}\pi }{Nz_c}}{\displaystyle \frac{\mathrm{sin}^2\left(\frac{\pi nN^{}}{N}\right)}{\mathrm{sin}\left(\frac{\pi n}{N}\right)}}.`$
We see that $`N`$ plays here the role of a cutoff of the Kaluza-Klein modes.
In previous models of deconstruction , mainly nearest neighbor couplings between the sites of the lattice have been considered. As a consequence, on imposing a periodic boundary condition, the lattice then looks as a circle. Departing from this standard situation, in the model considered here we have introduced non-nearest-neighbor couplings among the sites. That is, a site links to a number of other ones in a rather complicated way. In this sense, the lattice in the present model is no more a simple circle but it looks more like, say, a mesh or a net. Let us assume that the sites on the lattice would correspond to points in a brane. If the codimension of the spacetime is one, the brane should be ordered, resembling the sheets of a book. One brane can only interact (directly) with the two neighboring branes. However, if the spacetime is more complicated and/or the codimension is two or more, the brane can directly interact with three or more branes, an interaction that will be perfectly described by our model corresponding to this case. For example, a site on a tetrahedron connects directly with three neighboring sites. In this way, the non-nearest-neighbor couplings we here consider may quite adequately reflect the structure of the extra dimension. In this respect our model is very general and opens a number of interesting possibilities.
### 2.2 The supergravity case
By using the same sort of techniques described above, the multi-graviton model can be generalized to the supergravity case, just by starting with a supergravity theory in 5-dimensions and implementing deconstruction by way of replacing the fifth spacelike dimension with a one-dimensional lattice containing $`N`$-points. A multi-supergravity model of this kind has been proposed in Ref. , to which the interested reader is addressed for details. Here we shall only write down the essential aspects which will be used in what follows.
In the 5-dimensional linearized supergravity theory, the number of bosonic degrees of freedom is 8, 5 due to the massless graviton and 3 due to the massless vector (gauge) field and the number of fermionic degrees of freedom is 8 too, due to the the complex Rarita-Schwinger field ($`4\times 2`$).
The deconstruction process now consists in replacing the fifth dimension in the action of spin two+vector+Rarita-Schwinger fields with $`N`$points and the derivatives with respect to the corresponding variable with the operator $`\mathrm{\Delta }`$ as in Eq. (2.5). In this way one gets a complicated action in 4 dimensions, similar to the one in Eq. (2.3), but with vector and fermion parts too. It contains a spin-2 field $`h_{\mu \nu }`$ (the graviton), but also scalar, vector and fermionic fields. More precisely, in the massless sector one has 8 degrees of freedom due to bosons (graviton (2 d.o.f.), gauge and Stückelberg vectors (2+2 d.o.f.), a Stückelberg scalar and the fifth component of the gauge field (1+1 d.o.f.) and 8 degrees of freedom due to fermions (complex Dirac and Rarita-Schwinger fields), while in the massive sector one has again 8+8 degrees of freedom, but only due to a massive graviton, vector and Rarita-Schwinger fields. As in the pure-gravity case, one has $`N`$ copies of such fields and their masses —obtained by imposing periodic boundary conditions— are always given by means of Eq. (2.5), that is
$`\mathrm{\Delta }\varphi _n^p=i\mu _p\varphi _n^p,i\mu _p={\displaystyle \underset{n=0}{\overset{N1}{}}}a_ne^{2\pi inp/N},`$ (2.12)
$`\varphi _n^p=\varphi _{n+N}^p={\displaystyle \frac{e^{2\pi inp/N}}{\sqrt{N}}}.p=0,1,2,3,\mathrm{}`$ (2.13)
On the other hand, for fermion fields anti-periodic boundary conditions could also be considered. In such case one gets a different spectrum, given by means of the following equations
$`\mathrm{\Delta }\stackrel{~}{\varphi }_n^p=i\stackrel{~}{\mu }_p\varphi _n^p,i\stackrel{~}{\mu }_p={\displaystyle \underset{n=0}{\overset{N1}{}}}a_ne^{2\pi in(p+1/2)/N},`$ (2.14)
$`\stackrel{~}{\varphi }_n^p=\stackrel{~}{\varphi }_{n+N}^p={\displaystyle \frac{e^{2\pi in(p+1/2)/N}}{\sqrt{N}}}.p=0,1,2,3,\mathrm{}`$ (2.15)
It has to be noted that with boundary conditions of this sort, there are no massless fermions and this is a consequence of the explicitly breakdown of global supersymmetry.
## 3 The induced cosmological constant
We now turn to the evaluation of the induced cosmological constant for the $`N`$ graviton and super-graviton models discussed in the previous section. To this aim —the main one in the present paper— we shall compute the one-loop effective potential using zeta-function regularization ; needless to say, other regularization schemes could work as well. First of all, we compute the effective potential for a free scalar field with mass $`M`$, since this corresponds to the contribution of each degree of freedom to the one-loop effective potential of our theories.
In the zeta-function regularization method, the one-loop contribution to the effective potential is given by
$`V_{eff}^{(1)}={\displaystyle \frac{1}{2V}}\zeta ^{}(0|L/\mu ^2)={\displaystyle \frac{1}{2V}}\zeta ^{}(0|L){\displaystyle \frac{1}{2V}}\zeta (0|L)\mathrm{log}\mu ^2,`$ (3.1)
$`V`$ being the volume of the manifold and $`\zeta (s|L)`$ the zeta function corresponding to the Laplacian-like operator $`L=\mathrm{\Delta }^2+M^2`$, with $`M`$ a positive constant. The arbitrary parameter $`\mu `$ has to be introduced for dimensional reasons. It will be fixed by renormalization at the end of the process.
The manifold we are considering in the present paper is a flat one with non trivial topology of the kind $`=\text{I}\text{R}\times T^3`$. The simplest case $`=\text{I}\text{R}^4`$ has been already considered in .
The operator $`L`$ has the form
$`L={\displaystyle \frac{d}{d\tau ^2}}+L_3,L_3=\mathrm{\Delta }_3+M^2,`$ (3.2)
$`\mathrm{\Delta }_3`$ being the Laplace operator on $`T^3`$. The zeta-function is expressed in terms of the heat trace via the Mellin representation. The heat traces read
$`Tre^{tL}=V𝒦(t|L),Tre^{tL_3}=V_3𝒦(t|L_3),𝒦(t|L)={\displaystyle \frac{𝒦(t|L_3)}{\sqrt{4\pi t}}}.`$ (3.3)
As a result
$`\zeta (s|L)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{s1}Tre^{tL}={\displaystyle \frac{V}{\sqrt{4\pi }\mathrm{\Gamma }(s)}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{s3/2}𝒦(t|L_3)`$ (3.4)
$`=`$ $`{\displaystyle \frac{V\mathrm{\Gamma }(s1/2)}{\sqrt{4\pi }\mathrm{\Gamma }(s)}}\stackrel{~}{\zeta }(s1/2|L_3),`$
$`\stackrel{~}{\zeta }(s1/2|L_3)`$ being the zeta-function density on $`T^3`$ and $`V_3=(2\pi r)^3`$ the “volume” of the torus with “radius” $`r`$. The heat kernel and zeta function on $`T^3`$ are well known. In the Appendix A, for the reader’s convenience, we summarize some useful representations that will be used in what follows (for a review, see ).
Using expressions (3.4) and (A.1) one realizes that the zeta function can be written as the sum of two terms, that is
$`\zeta (s|L)=\zeta _0(s|L)+\zeta _T(s|L),`$ (3.5)
where $`\zeta _0`$ is the same one has on $`\text{I}\text{R}^4`$, namely
$`\zeta _0(s|L)={\displaystyle \frac{V\mathrm{\Gamma }(s2)M^{42s}}{16\pi ^2\mathrm{\Gamma }(s)}}={\displaystyle \frac{VM^{42s}}{16\pi ^2(s1)(s2)}},`$ (3.6)
while $`\zeta _T`$ represents the contribution due to the non-trivial topology, which explicitly depends on the topological parameter $`r`$. Expression (3.6) is also the leading contribution to the whole zeta function in a power series expansion for large values of $`M`$.
Recalling now (A.4), we obtain
$`\zeta _T(s|L)={\displaystyle \frac{V\mathrm{\Gamma }(s3/2)\mathrm{cos}\pi sM^{42s}}{8\pi ^{5/2}\mathrm{\Gamma }(s)}}{\displaystyle _1^{\mathrm{}}}𝑑uG(Mru)(u^21)^{3/2s}.`$ (3.7)
Observe that the topological contribution vanishes at $`s=0`$, and this means that
$`\zeta (0|L)=\zeta _0(0|L)={\displaystyle \frac{VM^4}{32\pi ^2}}.`$ (3.8)
Using (3.1), for the one-loop effective potential we finally have
$`V_{eff}^{(1)}={\displaystyle \frac{M^4}{64\pi ^2}}\left(\mathrm{log}{\displaystyle \frac{M^2}{\mu ^2}}{\displaystyle \frac{3}{2}}\right){\displaystyle \frac{M^4}{12\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(Mru)(u^21)^{3/2}.`$ (3.9)
It is interesting to note that for scalar fields, in the large mass case the topological contribution is always negative, and it is negligible with respect to the standard Coleman-Weinberg term.
As we have anticipated above, the parameter $`\mu `$ has to be fixed by a renormalization condition. To this aim, here we follow Ref. . The total one-loop effective potential is of the form
$`V_{eff}=V_R(\mu )+V_{eff}^{(1)}(\mu ),`$ (3.10)
$`V_R(\mu )`$ being the renormalized vacuum energy. For physical reasons, the last expression has to be independent of $`\mu `$, and this means that
$`\mu {\displaystyle \frac{dV_{eff}}{d\mu }}=0,`$ (3.11)
from which it follows that
$`V_R(\mu )=V_R(\mu _R)+{\displaystyle \frac{M^4}{64\pi ^2}}\mathrm{log}{\displaystyle \frac{\mu _R^2}{\mu ^2}},`$ (3.12)
$`\mu _R`$ being the renormalization point which has to be fixed by the condition $`V_R(\mu _R)=0`$. In this way, we finally get
$`V_{eff}={\displaystyle \frac{M^4}{64\pi ^2}}\left(\mathrm{log}{\displaystyle \frac{M^2}{\mu _R^2}}{\displaystyle \frac{3}{2}}\right)+V_T(r),`$ (3.13)
$`V_T(r)={\displaystyle \frac{M^4}{12\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(Mru)(u^21)^{3/2}={\displaystyle \frac{M^2}{16\pi ^4r^2}}{\displaystyle \underset{n\text{Z}\text{Z}^3;n0}{}}{\displaystyle \frac{K_2(2\pi Mr|\stackrel{}{n}|)}{n^2}}.`$ (3.14)
$`V_T(r)`$ represents the contribution coming from the non-trivial topology, which for scalar fields is always negative. We also note that, as a function of the topological parameter $`r`$, $`V_T(r)`$ can reach, in principle, any negative value. In Eq. (3.14), $`K_\nu `$ are the MacDonald’s (or modified Bessel) functions.
Before proceeding with the computation of the induced cosmological constant corresponding to the models we have discussed in Sect. 2, we first analyse here the behavior of $`V_T(r)`$ as a function of $`r`$. To this aim, we consider the two different regimes $`Mr1`$ and $`Mr1`$.
For the case $`Mr1`$, using (A.5) in (3.14) we get
$`V_T(r)`$ $`=`$ $`{\displaystyle \frac{M^4}{12\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(Mru)(u^21)^{3/2}`$
$`=`$ $`{\displaystyle \frac{1}{12\pi ^2r^4}}{\displaystyle _{Mr}^{\mathrm{}}}𝑑xG(x)(x^2M^2r^2)^{3/2}`$
$`=`$ $`{\displaystyle \frac{1}{12\pi ^2r^4}}[{\displaystyle _{Mr}^1}dx(1{\displaystyle \frac{1}{\pi ^2x^3}})(x^2M^2r^2)^{3/2}`$ (3.15)
$`+{\displaystyle _0^1}dxG_0(x)x^3+{\displaystyle _1^{\mathrm{}}}dxG(x)x^3+O(M^2r^2)]`$
$`=`$ $`{\displaystyle \frac{1}{32\pi ^2r^4}}\left[{\displaystyle \underset{n\text{Z}\text{Z}^3;|\stackrel{}{n}|0}{}}{\displaystyle \frac{1}{\pi ^4|\stackrel{}{n}|^4}}+O(M^2r^2)\right]{\displaystyle \frac{1}{64\pi ^2r^4}}+O(M^2/r^2).`$
We thus see that in this limit the leading term does not depend on $`M`$, and that it can be arbitrarily large, with a suitable choice of the parameter $`r`$. The series in the latter equation has been computed numerically.
On the contrary, in the opposite regime, $`Mr1`$, using (3.14) and the asymptotic expansion for the Bessel function, we obtain
$`V_T(r)`$ $`=`$ $`{\displaystyle \frac{M^2}{16\pi ^4r^2}}{\displaystyle \underset{n\text{Z}\text{Z}^3;n0}{}}{\displaystyle \frac{K_2(2\pi Mr|\stackrel{}{n}|)}{n^2}}`$
$``$ $`{\displaystyle \frac{3M^4}{32\pi ^4(Mr)^{5/2}}}e^{2\pi Mr}+\mathrm{}`$ (3.16)
In this limit the topological contribution could indeed be arbitrarily small. In Fig. 1 the whole behavior of the topological part of the effective potential is drawn. In order to work with dimensionless variables we have introduced the function $`\stackrel{~}{V}_T(y)r^4V_T(r)`$ of the dimensionless variable $`yMr`$. The graphic corresponds to the exact expression for $`\stackrel{~}{V}_T(y)`$, as given e.g. by the first lines of Eq. (3.15), multiplied by $`364\pi ^2`$. A very smooth transition from the behavior corresponding to $`Mr1`$, Eq. (3.15), to the one for $`Mr1`$, Eq. (3.16), is revealed.
In Fig. 2, the corresponding graphic of the full effective potential, Eq. (3.13), is drawn, again as a function of $`y`$ and setting $`\mu _Rr=1`$.
At this point, the effective potential $``$and, as a consequence, the induced cosmological constant for the models we are interested in$``$ can be obtained by adding up several contributions of the kind (3.13).
### 3.1 The multi-graviton model
We start with the explicit example of multi-graviton model given by (2.10), in which there is a single massless graviton and $`(N1)/2`$ couples of massive gravitons, with masses given by
$`M_0=|\mu _0|=0,M_p=|\mu _p|={\displaystyle \frac{\pi p}{z_c}},p=1,2,\mathrm{},{\displaystyle \frac{N1}{2}}.`$ (3.17)
On the manifold $`=\text{I}\text{R}\times T^3`$, the massless graviton gives no contribution to the effective potential, while it does appear explicitly on manifolds with a non-vanishing curvature. Since the massive gravitons always show up in pairs, in order to perform the computation of the effective potential, it is sufficient to consider only one half of the whole massive spectrum. Moreover, we have to take into account that each massive graviton contributes with five scalar degrees of freedom. After these considerations have been properly taken into account, for the effective potential we get the following expression
$`V_{eff}`$ $`=`$ $`10{\displaystyle \underset{p=1}{\overset{\frac{N1}{2}}{}}}{\displaystyle \frac{M_p^4}{64\pi ^2}}\left(\mathrm{ln}{\displaystyle \frac{M_p^2}{\mu _R^2}}{\displaystyle \frac{3}{2}}\right)`$ (3.18)
$`10{\displaystyle \underset{p=1}{\overset{\frac{N1}{2}}{}}}{\displaystyle \frac{M_p^4}{12\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(M_pru)(u^21)^{3/2}.`$
One can see that, as for the non-compact flat case (see Ref. for details), in order to have a (small) positive induced cosmological constant one has to consider a large value of $`N`$, that is, a huge number of massive gravitons. In this respect, the torus topology does not improve the situation. As we are now going to show, this is no longer the case for the multi-supergraviton model.
### 3.2 The multi-supergraviton model
Here we have to distinguish two cases: the first one corresponds to the choice of periodic boundary conditions in both the bosonic and fermionic sectors. In such situation, the degrees of freedom due to bosons exactly compensate the degrees of freedom due to fermions. Moreover, for any massive boson there is a fermion with the same mass and, since the contribution to the effective potential of any fermionic degree of freedom is opposite to the contribution of a bosonic degree of freedom, it turns out that the induced cosmological constant vanishes, independently of the mass spectrum.
In the second case, that is when anti-periodic boundary conditions are imposed in the fermionic sector, the situation changes completely, since the fermionic mass spectrum becomes really different with respect to the bosonic one. For example, by choosing $`N=3`$ , the solutions of Eqs. (2.12) and (2.14) give
$`M_0=0,M_1=M_2=m,\text{for bosons},`$ (3.19)
$`\stackrel{~}{M}_0=\stackrel{~}{M}_2={\displaystyle \frac{m}{\sqrt{3}}},\stackrel{~}{M}_1={\displaystyle \frac{2m}{\sqrt{3}}},\text{for fermions},`$ (3.20)
$`m`$ being an arbitrary mass parameter.
By taking into account the number of degrees of freedom, the one-loop effective potential becomes, in this case
$`V_{eff}`$ $`=`$ $`{\displaystyle \frac{M_1^4}{4\pi ^2}}\left(\mathrm{ln}{\displaystyle \frac{M_1^2}{\mu _R^2}}{\displaystyle \frac{3}{2}}\right){\displaystyle \frac{4M_1^4}{3\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(M_1ru)(u^21)^{3/2}`$ (3.21)
$`{\displaystyle \frac{\stackrel{~}{M}_0^4}{4\pi ^2}}\left(\mathrm{ln}{\displaystyle \frac{\stackrel{~}{M}_0^2}{\mu _R^2}}{\displaystyle \frac{3}{2}}\right)+{\displaystyle \frac{4\stackrel{~}{M}_0^4}{3\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(\stackrel{~}{M}_0ru)(u^21)^{3/2}`$
$`{\displaystyle \frac{\stackrel{~}{M}_1^4}{8\pi ^2}}\left(\mathrm{ln}{\displaystyle \frac{\stackrel{~}{M}_1^2}{\mu _R^2}}{\displaystyle \frac{3}{2}}\right)+{\displaystyle \frac{2\stackrel{~}{M}_1^4}{3\pi ^2}}{\displaystyle _1^{\mathrm{}}}𝑑uG(\stackrel{~}{M}_1ru)(u^21)^{3/2}`$
$`=`$ $`{\displaystyle \frac{m^4}{36\pi ^2}}\mathrm{log}{\displaystyle \frac{2^{16}}{3^9}}+V_T,`$
$`V_T`$ being the sum of all the topological contributions. As one sees, the first term on the right-hand side of the latter equation is always negative, but the whole effective potential could be positive due to the presence of the topological term. For example, in the regime $`mr1`$, from (3.15) one has
$`V_T{\displaystyle \frac{1}{8\pi ^2r^4}}V_{eff}>0\text{for}mr<\left({\displaystyle \frac{2}{9}}\mathrm{log}{\displaystyle \frac{2^{16}}{3^9}}\right)^{1/4}1.4,`$ (3.22)
while in the opposite regime, $`mr1`$, by using (3.16) one can see that the topological contribution although still positive it is negligible, and thus the effective potential remains negative.
In Fig. 3, the corresponding graphic of the full effective potential, Eq. (3.21), is drawn, again as a function of $`ymr`$.
## 4 Conclusions
In this paper, we have computed the effective potential corresponding to a multi-graviton model with supersymmetry in the case where the bulk is a flat manifold with the topology of a torus (more precisely $`\text{I}\text{R}\times T^3`$), and we have shown that the induced cosmological constant could be positive due to topological contributions.
In previous papers multi-graviton and multi-supergraviton models have been considered in $`\text{I}\text{R}^4`$. It has been shown that in the multi-graviton model the induced cosmological constant can be positive, but only if the number of massive gravitons is sufficiently large, while in the supersymmetric case the cosmological constant can be positive if one imposes anti-periodic boundary conditions in the fermionic sector. Note that the topological effects discussed above may also be relevant in the study of electroweak symmetry breaking in models with a similar type of non-nearest-neighbour couplings, for the deconstruction issue .
In the case of the torus topology, the topological contributions to the effective potential have always a fixed sign, depending on the boundary conditions one imposes. In fact, they are negative for periodic fields and positive for anti-periodic fields. This means that the torus topology provides a mechanism which, in a most natural way, permits to have a positive cosmological constant in the multi-supergravity model with anti-periodic fermions, being the value of such constant regulated by the corresponding size of the torus.<sup>1</sup><sup>1</sup>1A more crude analysis for the pure scalar case already hinted towards this conclusion. However, the sign issue was there not easy to fix , the reason being now clear. In this situation one can most naturally use the minimum number, $`N=3`$, of copies of bosons and fermions.
We finish with the remark that it would be interesting to apply the deconstruction scheme of Ref. also for the case of two latticized extra dimensions, which in the continuous limit would contain the orbifold singularity. This analysis might have a quite interesting impact on brane running coupling calculations .
## Acknowledgments
We thank Sergei D. Odintsov for useful discussions and suggestions. Support from the program INFN(Italy)-CICYT(Spain), from DGICYT (Spain), project BFM2003-00620, and from SEEU grant PR2004-0126 (EE), is gratefully acknowledged.
## A Zeta function on the torus
Eigenvalues of the Laplacian on the 3-dimensional torus are of the form $`\lambda _n=n^2`$, $`n\text{Z}\text{Z}^3`$, and thus the corresponding heat kernel is given by
$`𝒦(t|L_3)={\displaystyle \frac{e^{tM^2}}{V_3}}{\displaystyle \underset{n\text{Z}\text{Z}^3}{}}e^{tn^2/r^2}={\displaystyle \frac{e^{tM^2}}{(4\pi t)^{3/2}}}{\displaystyle \underset{n\text{Z}\text{Z}^3}{}}e^{\pi ^2n^2r^2/t},`$ (A.1)
being $`V_3=(2\pi r)^3`$ the “volume” of $`T_3`$. Using the expression above, the zeta function of this Laplacian can be written as
$`\zeta (s|L_3)=\zeta _0(s|L_3)+\zeta _T(s|L_3),`$ (A.2)
where the contribution
$`\zeta _0(s|L_3)={\displaystyle \frac{V_3\mathrm{\Gamma }(s3/2)M^{3/22s}}{(4\pi )^{3/2}\mathrm{\Gamma }(s)}},`$ (A.3)
comes from the $`n=0`$ term and it is the same one has for $`\text{I}\text{R}^3`$, while $`\zeta _T`$ corresponds to the contribution due to the non-trivial topology. Such term can be written in different ways, for instance, as an infinite sum of Bessel functions.
In Refs. one can find many interesting results concerning zeta functions and heat kernels corresponding to operators on manifolds with constant curvature. In particular, on the torus one has the very nice representation
$`\stackrel{~}{\zeta }_T(z|L_3)`$ $`=`$ $`{\displaystyle \frac{M^{32z}\mathrm{sin}\pi z}{4\pi ^2(1z)}}{\displaystyle _1^{\mathrm{}}}𝑑uG(Mru)(u^21)^{1z}`$ (A.4)
$`=`$ $`{\displaystyle \frac{\pi ^{z2}}{4\mathrm{\Gamma }(z)}}{\displaystyle \underset{n\text{Z}\text{Z}^3;n0}{}}\left({\displaystyle \frac{M}{r|\stackrel{}{n}|}}\right)^{3/2z}K_{3/2z}(2\pi Mr|\stackrel{}{n}|),`$
where $`G(x)`$ is given by
$`G(x)={\displaystyle \underset{n\text{Z}\text{Z}^3;n0}{}}e^{2\pi |\stackrel{}{n}|x}`$ $`=`$ $`1+{\displaystyle \frac{x}{\pi ^2}}{\displaystyle \underset{n\text{Z}\text{Z}^3}{}}{\displaystyle \frac{1}{(n^2+x^2)^2}}`$ (A.5)
$`=`$ $`1+{\displaystyle \frac{1}{\pi ^2x^3}}+{\displaystyle \frac{x}{\pi ^2}}G_0(x),`$
$`G_0(x)`$ being the regular function
$`G_0(x)={\displaystyle \underset{n\text{Z}\text{Z}^3;n0}{}}{\displaystyle \frac{1}{(n^2+x^2)^2}}.`$ (A.6)
|
warning/0506/astro-ph0506520.html
|
ar5iv
|
text
|
# Dark matter annihilation: the origin of cosmic gamma-ray background at 1–20 MeV
## Abstract
The origin of the cosmic $`\gamma `$-ray background at 1–20 MeV remains a mystery. We show that $`\gamma `$-ray emission accompanying annihilation of 20 MeV dark matter particles explains most of the observed signal. Our model satisfies all of the current observational constraints, and naturally provides the origin of “missing” $`\gamma `$-ray background at 1–20 MeV and 511 keV line emission from the Galactic center. We conclude that $`\gamma `$-ray observations support the existence of 20 MeV dark matter particles. Improved measurements of the $`\gamma `$-ray background in this energy band undoubtedly test our proposal.
What is the origin of the cosmic $`\gamma `$-ray background? It is usually understood that the cosmic $`\gamma `$-ray background is a superposition of unresolved astronomical $`\gamma `$-ray sources distributed in the universe. Active Galactic Nuclei (AGNs) alone explain most of the background light in two energy regions: ordinary (but obscured by intervening hydrogen gas) AGNs account for the low-energy ($`0.5`$ MeV) spectrum comastri ; zdziarski/etal:1995 ; ueda , whereas beamed AGNs (known as Blazars) account for the high-energy ($`20\mathrm{MeV}`$) spectrum salamon/stecker:1994 ; stecker/salamon:1996 ; pavlidou/fields:2002 . There is, however, a gap between these two regions. While historically supernovae have been a leading candidate for the background up to 4 MeVclayton/ward:1975 ; the ; zdziarski:1996 ; watanabe/etal:1999 , recent studiesstrigari/etal:2005 ; ahn/komatsu/hoeflich:2005 show that the supernova contribution is an order of magnitude lower than observed. The spectrum at 4–20 MeV also remains unexplained (for a review on this subject, see stecker/salamon:2001 ). It is not very easy to explain such high-energy background light by astronomical sources without AGNs or supernovae.
So, what is the origin of the cosmic $`\gamma `$-ray background at 0.5–20 MeV? On energetics, a decay or annihilation of particles having mass in the range of $`0.5\mathrm{MeV}m_X20\mathrm{MeV}`$ would produce the background light in the desired energy band. Since both lower- and higher-energy spectra are already accounted for by AGNs almost entirely, too lighter or too heavier (e.g., neutralinos) particles should be excluded. Is there any evidence or reason that such particles should exist? The most compelling evidence comes from 511 keV line emission from the central part of our Galaxy, which has been detected and mapped by the SPI spectrometer on the INTErnational Gamma-Ray Astrophysics Laboratory (INTEGRAL) satellitespi ; spi2 . This line should be produced by annihilation of electron-positron pairs, and one of the possible origins is the dark matter particles annihilating into electron-positron pairsmev\_dm . This proposal explains the measured injection rate of positrons as well as morphology of the signal extended over the bulge region. Intriguingly, popular astronomical sources such as supernovae again seem to fail to satisfy the observational constraintsdwarf . Motivated by this idea, in the previous paperahn/komatsu:2005 we have calculated the $`\gamma `$-ray background of redshifted 511 keV lines from extragalactic halos distributed over a large redshift range. We have shown that the annihilation signal makes a substantial contribution to the low-energy spectrum at $`<0.511`$ MeV, which constrains $`m_X`$ to be heavier than 20 MeV in order for the sum of the AGN and annihilation contributions not to exceed the observed signal.
In this paper, we extend our previous analysis to include continuum emission accompanying annihilation. The emerging continuum spectrum should of course depend on the precise nature of dark matter particles, which is yet to be determined. Recently, an interesting proposal was made by beacom/bell/bertone:2004 : radiative corrections to annihilation, $`XXe^+e^{}`$, should lead to emission of $`\gamma `$-rays via the internal bremsstrahlung, the emission of extra final-state photons during a reaction, $`XXe^+e^{}\gamma `$. They have calculated the spectrum of the internal bremsstrahlung expected for annihilation in the Galactic center, compared to the Galactic $`\gamma `$-ray data, and obtained a constraint on mass as $`m_X20`$ MeV. A crucial assumption in their analysis is that the cross section of internal bremsstrahlung is linearly proportional to the annihilation cross section, and the constant of proportionality is independent of the nature of annihilation, as is found for related processes crittenden/walker/ballam:1961 ; martyn ; berends/bohm ; bergstrom/etal:2005 . More specifically, they assumed that the cross section of $`XXe^+e^{}\gamma `$ would be calculated by that of $`e^+e^{}\mu ^+\mu ^{}\gamma `$ with the muon mass replaced by the electron mass. Although the equivalence between these two processes/cross-sections has not been demonstrated as yet, we adopt their procedure into our calculations.
We calculate the background intensity, $`I_\nu `$, as peacock
$$I_\nu =\frac{c}{4\pi }\frac{dzP_\nu ([1+z]\nu ,z)}{H(z)(1+z)^4},$$
(1)
where $`\nu `$ is an observed frequency, $`H(z)`$ is the expansion rate at redshift $`z`$, and $`P_\nu (\nu ,z)`$ is the volume emissivity (in units of energy per unit time, unit frequency and unit proper volume):
$$P_\nu =\frac{1}{2}h\nu \sigma vn_X^2\left[\frac{4\alpha }{\pi }\frac{g(\nu )}{\nu }\right],$$
(2)
where $`\alpha 1/137`$ is the fine structure constant, $`n_X`$ is the number density of dark matter particles, and $`\sigma v`$ is the thermally averaged annihilation cross section. To fully account for WMAP’s determination of mass density of dark matterwmap , $`\mathrm{\Omega }_Xh^2=0.113`$, by cold relics from the early universe, one finds $`\sigma v=[3.9,2.7,3.2]\mathrm{\hspace{0.17em}10}^{26}\mathrm{cm}^3\mathrm{s}^1`$ for $`m_X=[1,\mathrm{\hspace{0.17em}10},\mathrm{\hspace{0.17em}100}]\mathrm{MeV}`$, respectively (e.g., see Eq. in bes ). We have assumed that $`\sigma v`$ is velocity-independent (S-wave annihilation). One might add a velocity-dependent term (such as P-wave annihilation) to the cross-section; however, such terms add more degrees of freedom to the model, making the model less predictable. While Bœhm et al.mev\_dm argue that the S-wave cross section overpredicts the $`\gamma `$-ray flux from the Galactic center, we have shown in the previous paperahn/komatsu:2005 that it is still consistent with the data for $`m_X20`$ MeV and the Galactic density profile of $`\rho r^{0.4}`$ or shallower. (We shall discuss an issue regarding the density profile later.) Finally, a dimensionless spectral function, $`g(\nu )`$, is defined by
$$g(\nu )\frac{1}{4}\left(\mathrm{ln}\frac{s^{}}{m_e^2}1\right)\left[1+\left(\frac{s^{}}{4m_X^2}\right)^2\right],$$
(3)
where $`s^{}4m_X(m_Xh\nu )`$. This function is approximately constant for $`h\nu <m_X`$, and then sharply cuts off at $`h\nu m_X`$. Thus, one may approximate it as
$$g(\nu )\mathrm{ln}\left(\frac{2m_X}{m_e}\right)\vartheta (m_Xh\nu )$$
(4)
for the sake of an order-of-magnitude estimation. (Note that we have also assumed $`m_Xm_e`$.)
Since the number density is usually unknown, we use the mass density, $`\rho _Xn_X/m_X`$, instead. After multiplying by $`\nu `$, one obtains
$`\nu I_\nu `$ $`=`$ $`{\displaystyle \frac{\alpha h\nu \sigma v}{2\pi ^2m_X^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dzcg[(1+z)\nu ]}{H(z)}}\rho _X^2_z`$ (5)
$``$ $`3.800\mathrm{keV}\mathrm{cm}^2\mathrm{s}^1\mathrm{str}^1`$
$`\times \left({\displaystyle \frac{\sigma v}{10^{26}\mathrm{cm}^3\mathrm{s}^1}}\right)\left({\displaystyle \frac{h\nu 1\mathrm{MeV}}{m_X^2}}\right)`$
$`\times {\displaystyle }dz{\displaystyle \frac{g[(1+z)\nu ](1+z)^2(\mathrm{\Omega }_Xh^2)^2}{\sqrt{\mathrm{\Omega }_mh^2(1+z)^3+\mathrm{\Omega }_\mathrm{\Lambda }h^2}}}{\displaystyle \frac{C_X(z)}{10^3}},`$
where $`\rho _X^2_z`$ is the average of $`\rho _X^2`$ over proper volume at $`z`$, and $`C_X(z)\rho _X_z^2/\rho _X^2_z`$ is the dark matter clumping factor. (We have used $`\rho _X_z=10.54\mathrm{\Omega }_Xh^2(1+z)^3\mathrm{keV}\mathrm{cm}^3`$.) While equation (5) is exact, one may obtain a better analytical insight of this equation by using the approximation to $`g(\nu )`$ (Eq. ),
$`\nu I_\nu `$ $``$ $`3.800\mathrm{keV}\mathrm{cm}^2\mathrm{s}^1\mathrm{str}^1`$ (6)
$`\times \mathrm{ln}\left({\displaystyle \frac{2m_X}{0.511\mathrm{MeV}}}\right)\left[{\displaystyle \frac{(\mathrm{\Omega }_Xh^2)^2}{\sqrt{\mathrm{\Omega }_mh^2}}}\right]\left({\displaystyle \frac{\sigma v}{10^{26}\mathrm{cm}^3\mathrm{s}^1}}\right)`$
$`\times \sqrt{{\displaystyle \frac{1\mathrm{MeV}^2}{h\nu m_X}}}{\displaystyle _{h\nu /m_X}^1}𝑑yy^{1/2}{\displaystyle \frac{C_X[(m_X/h\nu )y]}{10^3}},`$
where $`yh\nu (1+z)/m_X`$. Here, we have also assumed that the integral is dominated by $`1+z(\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_m)=2.3`$.
We follow the method developed in our previous paper ahn/komatsu:2005 for calculating the clumping factor of dark matter, $`C_X(z)`$. We have shown that $`C_X(z)`$ at $`z20`$ is approximately a power law,
$$C_X(z)=C_X(0)(1+z)^\beta ,$$
(7)
and $`\beta `$ depends on adopted dark matter halo profiles. For example, a cuspy profile such as the Navarro-Frenk-White (NFW) profile nfw , $`\rho _X(r)r^1`$, gives $`C_X(0)10^5`$ and $`\beta 1.8`$, while a flat profile such as the Truncated Isothermal Sphere (TIS) tis , $`\rho _X(r)r^0`$, gives $`C_X(0)10^3`$ and $`\beta 0`$ (see Figure 2 of ahn/komatsu:2005 ). Using a power-law evolution of $`C_X(z)`$, one obtains an approximate shape of the spectrum as
$$\nu I_\nu \frac{h\nu \mathrm{ln}(2m_X/m_e)}{(\beta 3/2)m_X^2}\left[1\left(\frac{h\nu }{m_X}\right)^{\beta 3/2}\right]\vartheta (m_Xh\nu ),$$
(8)
for $`m_Xm_e`$. If $`\beta <3/2`$ (e.g., TIS), $`\nu I_\nu (h\nu )^{\beta 1/2}(\mathrm{ln}m_X)/m_X^{\beta +1/2}\vartheta (m_Xh\nu )`$, whereas if $`\beta >3/2`$ (e.g., NFW), $`\nu I_\nu h\nu [\mathrm{ln}(2m_X/m_e)]/m_X^2\vartheta (m_Xh\nu )`$. Note that the shape of the spectrum becomes insensitive to halo profiles for the latter case (while the amplitude still depends on profiles).
Henceforth we shall adopt the NFW profile as the fiducial model, as it fits the mean central halo profiles in numerical simulations well. Following the previous paper, we take into account a scatter in halo profiles by integrating over a probability distribution of halo concentration; thus, our model effectively incorporates significantly less concentrated (such as our Galaxy) or more concentrated profiles than the average NFW. One might argue that our model based on the NFW profile is unable to explain $`\gamma `$-ray emission from the Galactic center, which requires $`\rho r^{0.4}`$ (or shallower). If desired, one might use this profile and recalculate the $`\gamma `$-ray background spectrum; however, we continue to use the NFW profile, assuming that our Galaxy is not a “typical” halo in the universe. If there are so many more galaxies which obey the NFW profile, then the signal should be dominated by those typical halos. Of course, real universe does not have to be the same as numerical simulations, and one way to incorporate the uncertainty of halo profiles into our analysis would be to treat $`C_X(0)`$ and $`\beta `$ as free parameters. We shall come back to this point at the end of this paper. Figure 1 shows the predicted cosmic $`\gamma `$-ray background from dark matter annihilation, including lineahn/komatsu:2005 and continuum emission, for $`m_X=10`$, 20, and 50 MeV. The shape of the internal bremsstrahlung is described well by the approximate formula (Eq. ) with $`\beta =1.8`$. As expected, the continuum spectrum extends up to $`h\nu m_X`$, whereas line emission contributes only at $`<0.511`$ MeV.
Now let us add extra contributions from known astronomical sources and compare the total predicted spectrum with the observational data. Figure 2 compares the sum of dark matter annihilation, AGNsueda and Type Ia supernovaeahn/komatsu/hoeflich:2005 with the data points of HEAO-1heao , SMMsmm , and COMPTELcomptel experiments. We find that $`m_X20`$ MeV fits the low-energy spectrumahn/komatsu:2005 and explains about a half of the spectrum at 1–20 MeV. Therefore, the internal bremsstrahlung from dark matter annihilation is a very attractive source of the cosmic $`\gamma `$-ray background in this energy region. It is remarkable that such a simple model provides adequate explanations to two completely different problems: 511 keV line emission from the Galactic centermev\_dm , and missing $`\gamma `$-ray light at 1–20 MeV. (The regular Blazars would dominate the spectrum beyond 20 MeVsalamon/stecker:1994 ; stecker/salamon:1996 ; pavlidou/fields:2002 .)
If desired, one might try to improve agreement with the data in the following way. The continuum (combined with the other contributions) can fully account for the SMM and COMPTEL data, if the clumping factor is twice as large as predicted by the NFW profile. This could be easily done within uncertainty in our understanding of the structure of dark matter halos: for example, a slightly steeper profile, or the presence of substructuresubstruct . However, a larger clumping factor also increases 511 keV line emission by the same amount, which would exceed the HEAO-1 and SMM data. How do we reduce line emission independent of continuum? The line emission is suppressed by up to a factor of 4, if $`e^+e^{}`$ annihilation occurs predominantly via positronium formation. Once formed, a positronium decays into either two 511 keV photons or three continuum photons. As the branching ratio of the former process is only 1/4, line emission is suppressed by a factor of 4 if all of annihilation occurs via positronium formation. If a fraction, $`f`$, of annihilation occurs via positronium, then line is suppressed by $`13f/4`$beacom/bell/bertone:2004 ; thus, we can cancel the effect of doubling the clumping by requiring that 2/3 of line emission be produced via positronium. Figure 3 shows our “best-fit” model, which assumes (a) $`m_X=20`$ MeV, (b) the mean clumping factor is twice as large, and (c) line emission is solely produced via positronium ($`f=1`$). Note that this is a reasonable extension of the minimal model and makes the model more realistic: we know from simulations that there must exist substructures in halos. Some fraction of line emission must be produced via positronium, as it has been known that more than 90% of 511 keV emission from the Galactic center is actually produced via positronium formation kinzer/etal:2001 ; churazov/etal:2005 . While the model seems to slightly exceed the HEAO-1 and SMM data at low energy, we do not take it seriously as the discrepancy would be smaller than the uncertainty of the AGN model. The AGN model presented here assumes a high-energy cut-off energy of $`E_{\mathrm{cut}}=0.5`$ MeVueda . Since current data of AGNs in such a high energy band are fairly limited, uncertainty in $`E_{\mathrm{cut}}`$ is more than a factor of 2. Even a slight reduction in $`E_{\mathrm{cut}}`$ would make our model fit the low-energy spectrum.
The best-fit model is consistent with and supported by all of the current observational constraints: it fits the Galactic $`\gamma `$-ray emission as well as the cosmic $`\gamma `$-ray emission. It might also account for a small difference between theory and the experimental data of the muon and electron anomalous magnetic momentboehm/ascasibar:2004 . We stress here that, to the best of our knowledge, all of these data would remain unexplained otherwise. There is, however, one potential conflict with a new analysis of the SPI data by ascasibar/etal:2005 , which shows that a NFW density profile does provide a good fit to 511 keV line emission from the Galactic center, as opposed to the previous analysis by mev\_dm , which indicated a shallower profile than NFW. This new model would have much higher dark-matter clumping and require a substantially (more than an order of magnitude) smaller annihilation cross-section than $`\sigma v3\times 10^{26}\mathrm{cm}^3\mathrm{s}^1`$ to fit the Galactic data. Is our Galaxy consistent with NFW? This is a rather complicated issue which is still far from settled (e.g., binney/evans:2001 ; klypin/zhao/somerville:2002 ), and more studies are required to understand the precise shape of density profile of our Galaxy. If our Galaxy is described by a steep profile such as NFW, then the dark matter annihilation probably makes a negligible contribution to the $`\gamma `$-ray background, unless dark matter clumping is significantly increased by substructuresubstruct , compensating a small cross section. On the other hand, if it were confirmed that our Galaxy has a shallow density profile and the contribution of the dark matter annihilation to the $`\gamma `$-ray background is negligible, it would be difficult to explain the Galactic $`\gamma `$-ray signal solely by annihilation of light dark matter particles.
As shown in Figure 3, dark matter annihilation produces a distinctive $`\gamma `$-ray spectrum at 0.1–20 MeV. More precise determinations of the cosmic $`\gamma `$-ray background in this energy band will undoubtedly test our proposal. If confirmed, such measurements would shed light on the nature of dark matter, and potentially open a window to new physics: one implication is that neutralinos would be excluded from a candidate list of dark matter. Phenomenologically, our model may be parameterized by four free parameters: (1) dark matter mass, $`m_X`$, (2) a dark matter clumping factor at present, $`C_X(0)`$, (3) redshift evolution of clumping, $`\beta `$, and (4) a positronium fraction, $`f`$. When more precise data are available in the future, it might be possible to perform a full likelihood analysis and constrain properties of dark matter particles as well as dark matter halos.
Finally, the angular power spectrum of anisotropy of the $`\gamma `$-ray background at 1–20 MeV would also offer a powerful diagnosis of the detected signal (see zhang/beacom:2004 for the contribution from Type Ia supernovae). Our model predicts that the angular power spectrum should be given by the trispectrum (the Fourier transform of the four-point correlation function) of dark matter halos projected on the sky, as the signal is proportional to $`\rho ^2`$. More specifically, the power spectrum should follow precisely that of the dark matter clumping factor. More high-quality data of the cosmic $`\gamma `$-ray background in this energy band are seriously awaited.
We would like to thank D.E. Gruber for providing us with the HEAO-1 and COMPTEL data, K. Watanabe for providing us with the SMM data, Y. Ueda for providing us with the AGN predictions, C. Bœhm for sharing her results on the modeling of SPI data with us, and C. Bœhm and J. Beacom for valuable comments on early versions of this paper. We would also like to thank G. Bertone and P.R. Shapiro for discussion. K. A. was partially supported by NASA Astrophysical Theory Program grants NAG5-10825, NAG5-10826, NNG04G177G, and Texas Advanced Research Program grant 3658-0624-1999.
*Note added in proof* – A recent article by Rasera et al. (rasera/etal:2005 ) argues that our predictions for the $`\gamma `$-ray background from the redshifted 511 keV line (ahn/komatsu:2005 ) were too large because annihilation of electrons and positrons cannot take place in halos less massive than $`10^7M_{}`$, in which baryons cannot collapse. While this effect reduces the intensity of the line contribution, it does not affect the continuum emission (i.e., the internal bremsstrahlung), as the continuum emission is produced before annihilation. Since the major contribution to the $`\gamma `$-ray background at 1-20 MeV comes from the continuum emission, our conclusion in this paper is not affected by the results in rasera/etal:2005 .
|
warning/0506/astro-ph0506372.html
|
ar5iv
|
text
|
# Local Supercluster Dynamics: External Tidal Impact of the PSC𝑧 sample traced by Optimized Numerical Least Action Method
## 1 Introduction
Migration flows of cosmic matter are one of the major physical manifestations accompanying the emergence and growth of structure out of the virtually homogeneous primordial Universe. The cosmic flows displace matter towards the regions where ever more matter accumulates, ultimately condensing into the objects and structures we observe in the Universe.
Within the gravitational instability scenario of structure formation, the displacements are the result of the cumulative gravitational force exerted by the inhomogeneous spatial matter distribution of continuously growing density surpluses and deficits throughout the Universe. This establishes a direct causal link between gravitational force and the corresponding peculiar velocities. Given a suitably accurate measurement of peculiar velocities within a well-defined “internal” region of space, $`V_{int}`$, we may invert these velocities and relate them to the inducing gravitational force. Hence, the source of the measured motions may be traced and possibly even reconstructed. In principle, it may even allow us to infer the total amount of mass involved and thus estimate the value of the cosmological density parameter $`\mathrm{\Omega }_m`$ and other fundamental cosmological parameters.
The practical execution of such studies of cosmic velocity flows is ridden by various complicating factors. One major complication is that the cosmic regions in which peculiar velocities have been determined to sufficient accuracy may have a substantially smaller size than what may be deemed appropriate for a dynamically representative volume. Ideally, in order to account for almost the complete flow in our local cosmic neighbourhood we should have probed the density field in a sufficiently large cosmic volume. This should involve a region of space substantially superseding that of the characteristic scale of the largest coherent structures in the Universe. Only then the magnitude of the gravitational influence of inhomogeneities at larger distances will represent a negligible contribution and, as well, start to even out against each other.
The size of this dynamically effective volume depends sensitively on the structure formation scenario which is prevailing in our Universe. Hence, it will be closely affiliated to the spatial distribution, characteristic size and coherence scale of cosmic structures, and its size will therefore be in the order of the scale of the largest pronounced structures in the Universe. Within the conventional structure formation models, based on Gaussian initial density and velocity fields, this is fully specified through the scenario’s fluctuations power spectrum $`P(k)`$. When the power spectrum involves a substantial large-scale component and the survey volume is rather limited we have to be aware of significant external influences. Although not yet exactly determined, observational evidence suggests its size to be in the range of $`100200h^1\mathrm{Mpc}`$ (where $`h`$ denotes the Hubble constant in units of $`100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$).
An equally important consideration concerns the spatial resolution at which the velocity field is studied. Available samples of galaxy peculiar velocities extend out to reasonable depth of $`60h^1\mathrm{Mpc}`$. Yet, they involve a rather coarsely and inaccurately sampled cosmic velocity field. By absence of precise distance estimators, more accurately and densely sampled velocity information is therefore mostly confined to a rather limited region in and around the Local Supercluster \[LS\]. As a consequence, most analyses of large-scale cosmic flows are necessarily confined to spatial scales at which the evolving cosmic structures are still residing in a linear phase of development. The dynamics in more advanced stages of cosmic structure formation are as yet poorly constrained by measurements.
### 1.1 Cosmic Force Fields and Supercluster Dynamics
In this work we wish to extend the analysis of cosmic flows to the more advanced evolutionary stages pertaining within supercluster regions. Only within the local cosmic neighbourhood of our Local Supercluster, the quality, quantity and spatial coverage of the peculiar velocity data are sufficiently good to warrant an assessment of the cosmic velocity field and the corresponding dynamics at a sufficiently high spatial resolution. On these quasi-linear or mildly nonlinear scales we hope to find traces of the onset towards the more advanced stages of cosmic structure formation. In order for this to yield a meaningful and successful analysis, two major questions have to be addressed. Both form the main focus of this contribution.
The first issue, that of the rather restricted sample volume, constitutes the major incentive behind this work. The volume of the galaxy catalog that best samples our local cosmic neighbourhood, the Nearby Galaxy Catalog (Tully tullycat (1988), hereafter \[NBG\]), is certainly substantially smaller than what may be considered dynamically representative. Any analysis of the (internal) velocity field in our Local cosmic neighbourhood should therefore take into account the impact of external gravitational influences.
We focus on two related problems. In the first place, there is the need to quantify the effect of neglecting the external gravitational influence $`𝐠_{ext}`$ when modeling the cosmic velocity field on scales comparable to that of the Local Supercluster. Various studies have attempted to determine cosmological parameters on the basis of a comparison between modelled versus measured velocity field in the Local Universe (Tonry & Davis tonry (1981), Tully & Shaya tully (1984), Shaya, Peebles & Tully shaya (1995), Tonry et al. tonry00 (2000)). For this it is crucial to understand in how far local density perturbations may account for the local peculiar gravity field within the Local Supercluster. Directly related to this is the need to have a sufficiently accurate description of the external force field $`𝐠_{ext}`$, in terms of its nature and spatial extent, in order to properly model the total measured gravity field $`𝐠_{tot}`$. For studies intent on a comparison of modelled versus observed peculiar velocities on scales larger than the Local Supercluster this is an essential requirement (Faber & Burstein faber (1988), Han & Mould han (1990), Yahil et al.yahil (1991), Webster et al. webster (1997), Branchini et al. branchini99 (1999)). Similar considerations are equally relevant for the inverse problem, in which one attempts to infer the external gravitational influence $`𝐠_{ext}`$ from peculiar velocities measured within the Local Supercluster (e.g. Lilje, Yahil & Jones lilje (1986), Lynden-Bell & Lahav lynden (1988), Kaiser kaiser (1991), Hoffman et al. hoffman01 (2001)). Indeed, both problems concerning the restricted galaxy sample volume have figured prominently in previous cosmic velocity field studies and were addressed in a variety of publications. However, usually these tend to discard the fact that the local cosmic region in which we have access to high quality velocity data has already reached an advanced quasi-linear dynamical state.
Referring to the latter, the second major issue concerns the innovative way in which we evaluate the dynamical state of superclusters. These structures reside in a mildly nonlinear evolutionary stadium, having evolved significantly beyond their initial linear phase. Unlike the vast majority of previous studies, we seek to probe into the more detailed and informative kinematic aspects of these structures. A conventional linear analysis will not be able to provide an adequate description, and for the most evolved circumstances not even the Zel’dovich approximation (Zel’dovich zeld70 (1970)) may be expected to do so. In order to be able to optimally exploit the available velocity information – without suffering the loss of valuable high-resolution information through a filtering process – we apply the Least Action Principle \[LAP\] formalism (Peebles peebles89 (1989)) for dealing with the individual galaxy velocities. To that end, an optimal implementation developed by Nusser & Branchini (nusbran (2000), hereafter \[NB\]), the Fast Action Minimization \[FAM\] proved an essential tool.
Elaborating on the first issue of the external gravitational influence, one of the as yet undecided issues is the extent to which a LAP analysis of a cosmological self-gravitating system is dependent on a proper representation of the external gravitational influence. Various strategies have been followed, ranging from a complete neglect of external forces (Peebles peebles89 (1989)), or taking account of the influence of merely a few nearby objects (Peebles peebles89 (1989), peebles90 (1990), Dunn & Laflamme dunnlaf93 (1993), peebles2001 (2001)), towards methods involving approximate descriptions of external influences. The latter mostly incorporated the wider external influence through a frozen, linearly evolving, external tidal field estimated on the basis of the present-day locations of an extended sample of objects deemed representative for the external matter distribution (e.g. Shaya, Peebles & Tully shaya (1995), Schmoldt & Saha schmoldtsaha (1998), Sharpe et al. sharpe (2001)). This study does include the influence of force fields, but does so in a fully systematic and self-consistent fashion, enabled by the FAM method to take into account the evolution of the full sample of external matter concentrations.
The principal conclusion of our study is that the gravitational forces exacted by the matter inhomogeneities encapsulated by the IRAS-PSC$`z`$ redshift survey sample (Saunders et al. saunders (2000)) are indeed able to account for all motions within our local Universe. In addition, we demonstrate that its external influence may almost exclusively be ascribed to the bulk and shear flow components.
### 1.2 Strategy
This study is based on a number of artificial galaxy samples mimicking the properties of genuine catalogs. They consist of several well-defined and well-selected model catalogs of galaxies and galaxy peculiar velocities. These mock samples are extracted from a set of extensive $`N`$-body simulations: for the nearby Universe models they adhere to the characteristics of the Nearby Galaxy Catalog, for the deep galaxy redshift samples they are modelled after the IRAS PSC$`z`$ catalog. These model samples allow us to thoroughly investigate the various strategies forwarded for a successful and conclusive analysis.
Two sets of realistic mock catalogs of galaxies are extracted. The first set, the “local” one, is meant to mimic the mass distribution within the LS as traced by galaxies in the Nearby Galaxy Catalog of Tully (tullycat (1988)). It consists of a volume-limited galaxy sample within a (spherical) interior region with radius $`30h^1\mathrm{Mpc}`$. Each of these interior samples is embedded within a larger mock sample, the “extended” sample which covers a larger cosmic region. In addition to the interior volume-limited sample in the inner $`30h^1\mathrm{Mpc}`$ they contain an enclosing outer flux-limited sample covering the surrounding spherical region located between $`30h^1\mathrm{Mpc}<x<100h^1\mathrm{Mpc}`$. This exterior sample mimics a flux-limited galaxy catalog whose characteristics are modelled after the IRAS PSC$`z`$ sample.
For both the “local” and “extended” mock samples we model the peculiar velocity field at the positions of the particles in the “local” cosmic region, i.e. for the objects out to a radius $`x<30h^1\mathrm{Mpc}`$. These model predictions result from the application of Fast Action Minimization method (Nusser & Branchini nusbran (2000)). As our FAM reconstruction procedure only takes into account the gravitational forces between the particles in the mock samples – i.e. does not include contributions from outside objects – the differences in predicted velocities between the “local” and “extended” samples will reflect the influence of the mass concentrations in the surrounding region $`30h^1\mathrm{Mpc}<r<100h^1\mathrm{Mpc}`$. The comparison with the corresponding $`N`$-body velocities, representing the “real” velocities, will inform us in how far “sky-covering” samples of galaxies with a depth of $`100h^1\mathrm{Mpc}`$ may be expected to represent a proper cosmic region as far as its dynamics are concerned. The strategy of analyzing and comparing the velocity models obtained from the small “local” mock catalogs, the large “extended” mock catalogs, and the “real” velocities in the original $`N`$-body simulations, will yield a solid understanding of the effect of neglecting the externally induced peculiar gravitational acceleration $`𝐠_{ext}`$ (see Eqn. (1)).
When analyzing a dataset of galaxy peculiar velocities in the local Universe. The analysis of the larger PSC$`z`$ mimicking catalogs should elucidate if and to what extent such samples will be able to account for $`𝐠_{ext}`$. If the galaxies in these samples indeed appear to be responsible for the major share of the external forces, we may feel reassured to use the PSC$`z`$ sample of galaxies for a proper representation of $`𝐠_{ext}`$.
One aspect of this question concerns the investigation of the question whether the external tidal influence may be explicitly framed in an analytical approximation consisting of a dipolar and quadrupolar term. Our fully self-consistent FAM reconstructions, in which the “extended” mock catalogs are processed with the inclusion of all external matter concentrations, enable us to estimate the bulk and shear components in the induced “local” galaxy motions. By comparing the resulting velocity fields in the “local” and “extended” samples we will be able to judge the quality of the approximate methods, and quantify and investigate the possible presence of systematic trends throughout the “local” cosmos.
To account for possible systematic effects due to global cosmology, the mock galaxy catalogs are extracted from $`N`$-body simulations in two different cosmic structure formation scenarios. One involves a $`\mathrm{\Lambda }CDM`$ Universe with a characteristic large-scale dominated power spectrum, while the other concerns a $`\tau CDM`$ cosmology. The more small-scale dominated character of the latter leads to a different character of its gravitational field fluctuations, the smaller coherence scale of the density field fluctuations yielding a comparatively smaller influence of the external (quadrupolar) tidal field (the induced bulk flows are similar, as the smaller $`\tau CDM`$ fluctuations are exactly compensated by the larger mass involved). The resulting comparisons of FAM velocity field reconstructions are expected to reflect these velocity field differences.
In the end, this study of artificial galaxy samples should allow us to appreciate the manifestations of the real physical effects we wish to grasp. In this, we also should learn how to deal with the complications due to the host of measurement uncertainties which beset the observational data. The scope is to quantify the systematic errors which might have affected similar, local, comparisons based on real data and to judge whether the information on the external mass distribution available to these analyses is indeed sufficient to account for $`𝐠_{ext}`$.
### 1.3 Outline
In the next section, we will elaborate on the astrophysical background of this study, the study of velocity flows on cosmological scales, and in particular the issue of internal and external gravitational influences. Ensuingly, we address the specific problem of treating the dynamics and related cosmic motions within mildly nonlinear structures such as the Local Supercluster. This brings us to a brief exposition on the LAP analysis for dealing with the complications of mildly nonlinear orbits and the technical issue of the FAM technique which allows us to apply this to a system composed of many objects. Special emphasis is put on the inclusion of external gravitational influences within the LAP/FAM formalism. In section 4 we describe the cosmological setting of the simulations on which this study has been based. As a guidance towards interpreting our calculations, we address a variety of theoretical aspects and predictions concerning cosmic velocity fields in these cosmological scenarios. The basis of this work is the set of two “parent” N-body simulations and the mock catalogs extracted from these simulations, forty in total. They are presented in section 5. In the subsequent sections we present the results obtained from the various FAM computations. In section 6 we analyze the velocity vector maps for the FAM reconstructions. These maps allow a direct and visually illuminating appreciation of the effects we wish to address. This is followed by a first quantitative assessment in section 7. This consists of a comparison between the FAM velocity field reconstructions of the Local Supercluster volume($`r<30h^1\mathrm{Mpc}`$), the FAM reconstructions for the corresponding PSC$`z`$ sample and the complete “real world” N-body velocity field. The comparison is mainly based on a point-by-point evaluation through scatter diagrams of velocity-related quantities. To encapsulate these results into a spatially coherent description of the large scale external velocity and gravity field, in section 8 we turn to a decomposition of the peculiar velocity field into multipolar components. In particular, we demonstrate that a restriction to its dipolar and quadrupolar components, i.e. the bulk flow and velocity shear, does represent a good description. Thus having looked at the issue of cosmic velocity fields from different angles, the summary of section 9 will focus on the repercussions of our analysis and its relation to the study of the (relatively nearby) surrounding matter distribution. On the basis of these conclusions we provide a description of the various projects which follow up on this work, together with some suggestions for additional future work.
## 2 Cosmic Flows: <br>probes of cosmic matter distribution
### 2.1 the Large-Scale Universe: linear flows
Over the past two decades a major effort has been directed towards compiling large samples of galaxy peculiar velocities. These samples made it possible to obtain a rather impressive view of cosmic dynamics on scales
>
10h1Mpc
>
10superscript1Mpc\lower 2.69548pt\hbox{\vbox{\hbox{$>$}\hbox{$\sim$}} }10h^{-1}{\rm Mpc}. In particular the Mark III catalog, with an effective depth $`60h^1\mathrm{Mpc}`$, stands as a landmark achievement (Willick et al. 1997a , also see Dekel dekel (1994) and Strauss & Willick strauss (1995)). Further progress has been enabled by the assembly of additional and partially complementary samples of galaxy peculiar velocities, of which the SFI late-type galaxy and ENEAR early-type galaxy samples are noteworthy examples. The SFI Catalog of Peculiar Velocities of Galaxies (Giovanelli et al. 1997a , Giovanelli et al. 1997b , Haynes et al. 1999a and Haynes et al. 1999b ) consists of around 1300 spiral galaxies with I-band Tully-Fisher (TF) distances, out to $`cz<7500\mathrm{km}\mathrm{s}^1`$. The ENEAR sample (da Costa et al. costaenear (2000)) is an equivalent sample of around 1600 early-type galaxies, out to a distance $`cz<7000\mathrm{km}\mathrm{s}^1`$, with $`D_n\sigma `$ distance estimates available for nearly all of them. Tracing cosmic motions over larger volumes of space is a rather more cumbersome affair and attempts to do so are mainly based upon the peculiar motions of galaxy clusters. The claim of a puzzlingly large flow over scales of $`150h^1\mathrm{Mpc}`$ by Lauer & Postman laupost (1994) could not be corroborated. Nonetheless, flows on such large scales may indeed be a reality, as has been inferred from the far better defined “Streaming Motions of Abell Clusters” (SMAC) sample of Hudson et al. hudsiii (2001). They did recover a bulk flow in the order of $`687\pm 203\mathrm{km}\mathrm{s}^1`$, of which $`225\mathrm{km}\mathrm{s}^1`$ may arise from sources at a distance larger than $`100h^1\mathrm{Mpc}`$ (Hudson et al. hudsv (2004)). One prime objective of most analyses of these large samples of peculiar velocities has been the determination of the cosmological mass density parameter $`\mathrm{\Omega }_m`$ (Davis, Nusser & Willick 1996, Willick et al. 1997b , Willick & Strauss willick (1998), Nusser et al. nusseretal (2000), Branchini et al. branchini01b (2001)). Such assessments are based on a comparison of observed velocities to a model velocity field. A basic requirement for obtaining self-consistent estimates of $`\mathrm{\Omega }_m`$ is that the velocity samples concern a “representative” volume of space.
However, even while such studies appear to succeed in attuning the large-scale matter distributions and velocity fields in a reasonably self-consistent fashion, doubts remain with respect to a variety of practical and systematic problems. Firstly, in these comparisons the random errors on the observed velocities are substantial, much larger than those in the structure formation models. Considerable effort has been directed towards quantifying and minimizing errors on the observed peculiar velocities (e.g. Dekel dekel (1994), Strauss & Willick strauss (1995), and references therein). These involve random measurement errors as well as more subtle systematic, yet reasonably well understood, errors. Secondly, there remain various systematic effects which have not been addressed and corrected for in an equally convincing fashion. Even though they also tend to play a role with respect to the model predictions they are often overlooked.
A major systematic factor concerns the incomplete information on the spatial mass distribution within the region of the sample itself. This prevents an adequate treatment of artifacts due to the incomplete sky coverage and limited depth of the available samples, and effects systematic errors stemming from luminosity and density effects. These systematic errors are usually accounted for by using large, all-sky redshift surveys, such as the Optical Redshift Survey of Santiago et al. (santiago (1995)) or the 1.2 Jy and PSC$`z`$ surveys of IRAS galaxies (Fisher et al. fisher (1995), Saunders et al. saunders (2000)). In particular when using IRAS based surveys the effects of incomplete sky coverage are greatly reduced.
Even more problematic for a successful handling of luminosity and density related effects is our incomplete knowledge with respect to the relationship between the observable galaxy distribution and the underlying mass distribution. By absence of a compelling theory of galaxy formation this “galaxy bias” is usually encapsulated in heuristic formulations. The rather ad-hoc and possibly unrealistic or inadequate nature of the latter may seriously affect the significance of the inferred conclusions. Most studies make the simplifying assumption of a galaxy population fairly tracing the underlying density field. This is usually embodied in a global and linear “galaxy bias” factor. A large variety of results suggest that this may be a reasonable approximation on scales in excess of a few Megaparsec. Moreover, while this bias may be problematic in the case of early-type galaxies, it has proved to be quite successful with respect to the later type galaxies which figure prominently in IRAS based samples (Verde et al. verde2002 (2002)).
### 2.2 Internal and External influences
Unlike most studies of cosmic flows which seek to assess and analyze the nature and source of dynamical influences within a confined region of space, we try to get an impression of the cosmic dynamics on mildly nonlinear scales of only a few Megaparsec. We focus on the Local Supercluster region and its immediate neighbourhood. The galaxy sample of the NBG catalog is taken to be representative for this region. Because the catalog entails a volume which is substantially smaller than what may be considered dynamically representative, the peculiar velocities of the galaxies are partially due to the gravitational action by outside matter concentrations. That is, the peculiar velocities are not only due to the gravitational force induced by the matter concentrations within the “internal” survey volume $`V_{int}`$, but also reflect the gravitational influence by the “external” matter density distribution, $`𝐠_{ext}`$. Because it does not constitute a truly representative volume of the Universe, a meaningful dynamic analysis of the Local Universe on the basis of the NBG sample is substantially complicated by its limited depth, which is one of the major systematic problems besetting the analysis of virtually all available surveys of galaxy peculiar velocities. Theoretical models of peculiar velocities nearly always involve the implicit assumption of the mass being homogeneously distributed outside $`V_{int}`$, so that its gravitational effect may be neglected. Even in the case of having a sufficiently large volume at one’s disposal, this approximation is only valid in the central part of $`V_{int}`$, certainly not near its edges.
The distinction between external versus internal gravitational force may be best appreciated by noting that the total (peculiar) gravity field $`𝐠_{tot}(𝐱)`$ is the netto sum of the individual contributions by all patches of matter throughout the visible Universe. At any position within the internal volume $`V_{int}`$, we may then decompose the full gravitational field into an “internally” induced component $`𝐠_{int}`$ and an “externally” generated contribution $`𝐠_{ext}`$,
$$𝐠_{tot}(𝐱)=𝐠_{int}(𝐱)+𝐠_{ext}(𝐱).$$
(1)
In this way we have defined the internal gravitational force $`𝐠_{int}`$ as the integrated contribution from the density fluctuations $`\delta (𝐱)`$ within the volume $`V_{int}`$, while the external gravitational force $`𝐠_{ext}`$ concerns the combined gravitational force generated by the density fluctuations outside the realm of $`V_{int}`$, so that
$`𝐠_{tot}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{3\mathrm{\Omega }H^2}{8\pi }}{\displaystyle _{V_{int}}}d𝐱^{}\delta (𝐱^{},t){\displaystyle \frac{(𝐱^{}𝐱)}{|𝐱^{}𝐱|^3}}+`$ (2)
$`{\displaystyle \frac{3\mathrm{\Omega }H^2}{8\pi }}{\displaystyle _{V_{ext}}}d𝐱^{}\delta (𝐱^{},t){\displaystyle \frac{(𝐱^{}𝐱)}{|𝐱^{}𝐱|^3}}.`$
The peculiar velocities of galaxies within $`V_{int}`$ bear the mark of both the acceleration due to the matter concentrations within the volume itself, $`𝐠_{int}`$, as well as that of the combined gravitational influence of the external mass distribution, $`𝐠_{ext}`$. A comparison of predicted internally induced velocities with the observed local velocity field should therefore enable us to infer the magnitude and nature of the external field $`𝐠_{ext}`$. This analysis is usually facilitated by the fact that the fine details of the external force contribution are largely negligible. The contributions by the various external matter concentrations to the combined gravitational force mostly average out such that what remains noticeable is mainly confined to the low order components of the multipole decomposition of $`𝐠_{ext}`$. This can be most readily appreciated from a description of the external gravitational force field in terms of its successive multipole components. When we expand $`𝐠_{ext}`$ around some central location in $`V_{int}`$ – here defined to be the origin of the coordinates $`𝐱`$ – we find that to second order
$`g_{ext,i}(𝐱)`$ $`=`$ $`g_{bulk,i}{\displaystyle \underset{j=1}{\overset{3}{}}}𝒯_{ij}x_j.`$ (3)
In this, we assume that the additional divergence term $`\frac{1}{3}(𝐠_{ext})x_i`$ has been embedded into the (zeroth) order monopole term. In essence, it corresponds to a “breathing mode” affecting the “local” Hubble expansion within the volume, and therefore can not possibly be inferred from the local measurement of the internal gravity field $`𝐠_{int}`$.
The leading term in the overall external gravitational acceleration is the bulk gravity term $`g_{bulk,i}`$. This dipole term constitutes the uniform acceleration of the matter within $`V_{int}`$,
$`𝐠_{bulk}={\displaystyle \frac{3\mathrm{\Omega }H^2}{8\pi }}{\displaystyle _{V_{ext}}}d𝐱^{}\delta (𝐱^{},t){\displaystyle \frac{𝐱^{}}{|𝐱^{}|^3}}`$ (4)
Evidently, when considering peculiar velocities relative to the centre of mass inside the volume $`V_{int}`$ instead of absolute velocities this constant vector disappears. The first term whose magnitude and configuration is independent of the reference frame is the quadrupolar term $`𝒯_{ij}`$.
If the contribution to the (peculiar) gravitational potential by the external mass inhomogeneities is $`\varphi _{ext}`$, the quadrupolar tidal tensor $`𝒯_{ij}`$ is the trace-free part of $`^2\varphi _{ext}/x_ix_j`$, evaluated at the centre of $`V_i`$. It is determined by the external matter distribution through
$$𝒯_{ij}(t)=\frac{3\mathrm{\Omega }H^2}{8\pi }_{V_{ext}}d𝐱^{}\delta (𝐱^{},t)\left\{\frac{3x_i^{}x_j^{}|𝐱^{}|^2\delta _{ij}}{|𝐱^{}|^5}\right\}$$
(5)
The integral expressions for the dipole and quadrupole components of the external gravity field (eqn. 4 and 5), illustrate that it is unfeasible to exploit the observed local cosmic velocity field to recover the detailed and complete spatial distribution of the external matter inhomogeneities. On the other hand, it does indicate how it is that we can infer some overall characteristics of the external matter distribution from an analysis of the local velocity field. From this we may extract interesting and significant information on the nature and even distribution of the large scale cosmic matter distribution and set constraints on the values of some of the fundamental cosmological parameters. The pioneering work by Lilje, Yahil & Jones (lilje (1986)) in which the velocity field of the Local Supercluster was exploited to infer the presence of a major external source of gravitational attraction has shown the potential of this approach. Ultimately, it inspired the analysis of Lynden-Bell et al. (samur88 (1988)) that lead to the discovery of the Great Attractor.
## 3 Cosmic Flows: <br>the mildly nonlinear dynamics of superclusters
Even though a structure’s evolution may have progressed to a dynamical stage at which a first-order description of cosmic velocity fields will no longer be adequate, it may still be possible to find a direct link to the structure’s initial configuration. This is in particularly true for the early and mildly nonlinear phases of evolution. The exemplary archetype of a structure in which such mildly nonlinear circumstances are prevalent is that of superclusters, the filamentary or wall-shaped elements of the cosmic foamlike matter distribution.
Over the past two decades intriguing foamlike patterns have gained prominence as a prime characteristic of the cosmic matter distribution. The first indications for the actual existence of a foamlike galaxy distribution were provided by CfA2 redshift slices (de Lapparent, Geller & Huchra lappcfa (1986)) and established as a universal cosmic pattern with the Las Campanas redshift survey (Shectman et al. lcrs (1996)). With the arrival of the large recent and ongoing systematic galaxy redshift surveys, the 2dF galaxy redshift survey ($`250,000`$ redshifts, Colless et al. coll2df (2003), also see e.g. Colless coll2dfcosm (2004) and Tegmark et al. teg2df (2002) for a discussion on clustering in the 2dFGRS) and the Sloan Digital Sky Survey (SDSS, will determine $`1,000,000`$ redshifts, see e.g. Zehavi et al. zehasdss (2002) and Tegmark et. al. tegsdss (2004) for an overview of present-day status wrt. galaxy clustering), we may hope to have entered the stage in which we will be enabled to explore the formation and the dynamics of these characteristic spatial structures in the cosmic matter distribution. The typical elements of the cosmic foam – filamentary and wall-shaped superclusters – are precisely at the youthful yet mildly nonlinear phase of development mentioned earlier. They were identified as such within the context of Zel’dovich’ “pancake” theory of cosmic structure formation (see e.g. Shandarin & Zel’dovich shandzeld (1989)). The significance of the cosmic foamlike network for the understanding of the process of cosmic structure formation has since been generally recognized. This may be appreciated from the widespread use of the concept of the ‘cosmic web’, coined by Bond, Kofman & Pogosyan (cosmweb (1996)) in their study of the dynamics underlying its formation (see Van de Weygaert weyfoam (2002) for a recent general review).
Mildly nonlinear cosmic features such as superclusters have recently turned their initial co-expansion into a genuine physical contraction (or are on the brink of doing so), marking the emerging structure as a truely identifiable entity. Once it has surpassed this “turn-around” stadium the complexity of its internal kinematics quickly increases. At first moderately, and ultimately dramatically as the virialization process advances, the matter orbits inside the structure become more and more complex. Even in the more moderate early phases of this process, an appropriately sophisticated treatment of the mildly nonlinear dynamics appears to be a necessary requisite for any study based upon kinematic information. In and around emerging nonlinear structures a simple linear analysis for reconstructing initial conditions will therefore no longer suffice. In other words, a sufficiently detailed and profound understanding of the emergence of these key elements in the cosmic matter distribution cannot be obtained without the development of a more elaborate technique for the analysis of cosmic velocity fields.
### 3.1 Structure formation: mildly nonlinear dynamics
A linear analysis simplifies the dynamical evolution of a system into an initial conditions problem. It implies the reconstruction of the primordial density and velocity field by means of a simple linear inversion of the observed matter distribution and galaxy peculiar velocity field. Such an approach may even be followed towards a slightly more advanced stage. The Zel’dovich formalism, a Lagrangian first-order approximation for gravitationally evolving systems, has been remarkably successful in describing the early nonlinear evolution of a supercluster (for a review, see Shandarin & Zel’dovich shandzeld (1989)). Substantially surpassing its formal linear limitations, it proved to be a highly versatile medium for analyzing and explaining the overall spatial morphology and characteristics of emerging structures. The Zel’dovich approximation elucidated and explained qualitatively the fundamental tendencies of gravitational contraction in an evolving cosmos. Perhaps most noteworthy this concerned the tendency of gravitational collapse to proceed anisotropically, together with its predictive power with respect to location and timescales of the first phase of collapse into planar mass concentrations, “pancakes”. This offered the basic explanation for the foamlike morphology of the cosmic matter distribution, stressing its existence many years in advance of its discovery through observational programs to map the galaxy distribution (for an extensive review of various nonlinear approximation schemes seeking to expand upon the Zel’dovich approximation see Sahni & Coles sahnicoles (1995)).
In line with the above, the Zel’dovich approximation proved a highly versatile tool for the analysis of the cosmic matter flows. It was successfully applied to the nonlinear situation of mixed boundary conditions – tested and calibrated using $`N`$-body simulations – by Nusser et al. (1991) and Nusser & Dekel (1992). However, its validity remains restricted to the early stages of nonlinearity at which there is still a linear and direct relation between velocity and gravity field. Once matter inside the emerging structures starts to reach densities so high that local interactions become dominant, the Zel’dovich scheme quickly ceases to lose its applicability. Once matter elements start to cross each each others path the interaction between the nonlinear matter concentrations within the realm of the contracting structure will more and more deflect the orbits away from their initial linear trajectory. The linear kinematics of the Zel’dovich approximation will therefore no longer be able to follow the orbits of the matter elements. Higher order approximations based on perturbation theory have been advocated in order to follow such more advanced nonlinear circumstances. However, the improvement over simple first order Zel’dovich approximation turns out to be quite limited and not warranting the effort invested at each successive perturbation step. This is particularly so as with the onset of nonlinearity the rate at which successive perturbative orders terms become significant rapidly accelerates.
### 3.2 Least Action Principle in Cosmology
In more advanced nonlinear circumstances we encounter a more generic dynamical situation than a simple initial value problem. Typically, one seeks to compute the velocity field consistent with an observed density structure at the present epoch or, reversely, one deduces the density from the measured peculiar galaxy velocities. In the case of generic systems, the dynamical evolution represents a mixed boundary condition problem. This implies the system to be sufficiently constrained by complementing the incomplete dynamical information regarding the initial conditions with that pertaining to the dynamical state of the system at a different epoch. While $`N`$-body codes are particularly concerned with the ideal circumstances usually embodied in terms of initial value problems, a different kind of technique needs to be invoked to exploit the typical mixed boundary information yielded by observations.
A more profound and direct exploitation of the available information to follow the physics of such a cosmological nonlinear system was forwarded by Peebles (peebles89 (1989), peebles90 (1990)). He pointed out that finding the orbits that satisfy initial homogeneity – and by implication vanishing initial peculiar velocities – and match the (present-day) observed distribution of mass tracers constitutes a mixed-boundary value problem. Such problems lend themselves naturally to an application of Hamilton’s principle. This naturally leads to the formulation of the Least Action Principle (also known as “Numerical Action Method”), based on a variational analysis of the action $`S`$ of an isolated system of $`M`$ particles, which at a cosmic expansion factor $`a(t)`$ is given by
$$S=_0^{t_0}L𝑑t=_0^{t_0}𝑑t\underset{i}{}[\frac{1}{2}m_ia^2\dot{𝐱}_i^2m_i\varphi (𝐫_i)],$$
(6)
in which $`L`$ is the Lagrangian for the orbits of particles with masses $`m_i`$ and comoving coordinates $`x_i`$ and corresponding peculiar gravitational potential $`\varphi (𝐱)`$. For a system of particles interacting by gravity alone, embedded within a uniform cosmological background of density $`\rho _b(t)`$, this yields the following explicit expression for the action $`S`$,
$`S={\displaystyle _0^{t_0}}𝑑t`$ $`[`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{m_ia^2}{2}}\left({\displaystyle \frac{d𝐱_i}{dt}}\right)^2+{\displaystyle \frac{G}{a}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{m_im_j}{|𝐱_i𝐱_j|}}+`$ (7)
$`+{\displaystyle \frac{2}{3}}\pi G\rho _ba^2{\displaystyle }m_i𝐱_i^2]`$
The exact equations of motion for the particles are then obtained from finding the stationary trajectories amongst the variations of the action $`S`$ subject to fixed boundary conditions at both the initial and final time.
Confining oneself to a feasible approximate evaluation in this Least Action Principle approach, one describes the orbits of particles, $`𝐱_i(t)`$, as a linear combination of suitably chosen universal functions of time with unknown coefficients specific to each particle presently located at a position $`𝐱_{i,0}`$. For instance, by using the linear growth mode $`D(t)`$ as time variable (Giavalisco et al. giav93 (1993), Nusser & Branchini nusbran (2000)), one can parametrize the orbit $`𝐱_i(D)`$ of a particle as
$$𝐱_i(D)=𝐱_{i,0}+\underset{n=1}{\overset{N_f}{}}q_n(D)𝐂_{i,n},$$
(8)
where the functions $`q_n(D)`$ form a set of $`N_f`$ time-dependent basis functions. The factors $`𝐂_{i,n}`$ are then a set of free parameters, whose value is determined from evaluating the stationary variations of the action.
The functions $`q_n(D)`$ satisfy both two orbital constraints, necessary and sufficient to define solutions in agreement with evolution in the context of the Gravitational Instability theory for the formation of structure in the Universe: $`q_n(1)=0`$ ensures that at the present time the galaxies are located at their observed positions $`𝐱_i(1)=𝐱_{i,0}`$ and $`lim_{D0}D^{3/2}q_n(D)\theta (D)=0`$ guarantees vanishing peculiar velocities at early epochs which, in turns, ensures initial homogeneity.
### 3.3 Fast Action Minimization
The successful application of the Least Action Principle towards probing the kinematics and dynamics of an evolving cosmological system depends to a large extent on the specific implementation. This will be dictated by the characteristics of the physical system. In order to enable a meaningful LAP analysis of large samples of galaxies, like the Local Universe samples studied in this work, an optimized procedure is necessary for dealing with the large number of objects. Nusser & Branchini (nusbran (2000)) developed an optimized version of Peebles’ original LAP formalism, the Fast Action Minimization method. The various optimization aspects of the FAM implementation proved to be crucial for our purposes. The relevant optimization hinges on four major aspects of the FAM scheme.
The first FAM improvement involves the choice of time basis functions $`q_n(D)`$. Its convenient choice of time basis functions yields a simple expression for the action of the system and for its derivatives with respect to $`𝐂_{i,n}`$. Both quantities relate to the internal gravity term $`𝐠_{int}`$ of the system. Once the action and its derivatives are evaluated numerically, the minimum of the action is determined by means of the conjugate gradient method (Press et al. press (1992)). The orbits of the system are then fully specified through the set of parameters $`𝐂_{i,n}`$ found in correspondence to the minimum.
Closely related to the first aspect is that of tuning the choice of the time basis functions $`q_n(D)`$ such that only a limited number $`N_f`$ of basis functions is needed to successfully parameterize the orbits of the system. This is in particular beneficial to the the physical configuration we are studying here, involving Megaparsec scale dynamics characterized by quasi-linear or mildly nonlinear motions.
Note that using the growth factor $`D`$ as time variable makes the equations of motions almost independent of the value of $`\mathrm{\Omega }_m`$ (Nusser and Colberg 1998). As a consequence FAM orbits and peculiar velocities in a generic $`\mathrm{\Omega }_m`$ universe can be obtained by appropriate scaling those assuming a flat cosmology.
A final major aspect of the FAM implementation involves the efficient computation of the internal (self-consistent) gravity $`𝐠_{int}`$ from the particle distribution in the sample. To this end, the gravitational forces acting on the particles at the different epochs are computed by means of the TREECODE technique (Bouchet & Hernquist bouchet (1988)). By proceeding in this fashion, the FAM method is able to reconstruct the orbits of $`10^410^5`$ mass tracing objects back in time. This makes FAM numerically fast enough to perform a large number of orbit reconstructions, essential for performing the intended statistical analysis presented in the following sections.
In this work we use $`N_f=6`$ basis functions to parameterize the orbits, choosing a tolerance parameter $`tol=10^4`$ to search for the minimum of the action $`S`$ and setting a softening parameter of $`0.27h^1\mathrm{Mpc}`$ to smooth the gravitational force in the TREECODE. Orbit searching in dynamically relaxed systems is a difficult exercise since one has to choose among the many solutions found at the extrema of the action. However, since the purpose of FAM is to investigate large scale dynamics dominated by coherent flows rather than virial motions, our evaluations translates into an orbit search restricted to solutions which do not deviate too much from the Hubble flow i.e. to the simplest orbits that represents the minima of the action. Therefore, we set the initial guess for $`𝐂_{i,n}`$ according to linear theory prescription and search for the minimum of the action to avoid multiple solutions found a stationary points which typically describe passing orbits (Peebles peebles94 (1994)). We have checked that this choice of parameters is optimal in the sense that decreasing $`tol`$, increasing $`N_f`$ or changing the input set of $`𝐂_{i,n}`$ does not modify the final results appreciably.
Distortions in the resulting FAM-predicted peculiar velocities mainly arise from two systematic artifacts (Branchini, Eldar & Nusser branchini02 (2002)). One is the discrete sampling of the mass distribution within $`V_{int}`$. The second, and overriding one, is the failure of FAM in reproducing the virial motions within high-density regions that is a direct consequence of having considered solutions that represents perturbations to the Hubble flow. This deficiency of the FAM reconstructions is clearly illustrated by the residual velocity vector maps (see eq. 20) in Fig. 7 and Fig. 8 (bottom row). These show the velocity vector differences between the “real” measured, i.e. $`N`$-body, velocities and the corresponding FAM reconstructions (here based on either the galaxy distribution in a $`30h^1\mathrm{Mpc}`$ central region or the extended $`100h^1\mathrm{Mpc}`$ region). The maps show how the largest residuals are the ones found in the high density regions: although the FAM<sub>30</sub> and FAM<sub>100</sub> velocity fields do show pronounced velocities near these regions they are not the proper “real” virialized velocities they should have been. The residual fields thus underline the fact that FAM’s inaptitude to deal with regions characterized by large virial motions. Instead, in those situations it will lead to a false prediction of coherent inward streaming velocities, an effect pointed out by Nusser & Branchini (nusbran (2000)) and which can be also noted in our images when carefully studying them.
Finally, for practical reasons, since we are merely interested in measuring the effect of external gravity fields we make a further simplifying hypotheses. We ignore redshift distortion effects by working in real space In this respect, we should point out that extensions of the action principle method allowing a direct processing of redshift space information have been proposed and shown to work (Phelps phelps2000 (2000), Phelps phelps2002 (2002), also see Sharpe et al. sharpe (2001)).
### 3.4 The role of biasing
In this work we perform orbit reconstructions by assuming that all the mass of the systems is associated to point mass objects. More explicitly, we are making two different hypotheses. The first one is that we are able to identify a set of objects that trace the underlying mass density field in an unbiased way. The second one is that the internal structure of these objects is irrelevant for our reconstruction purposes.
The first assumption hardly applies to real galaxies that are most likely to be biased tracers of the mass distribution, as indicated by the relative bias between galaxies with different luminosities, colors and morphological type (Loveday et al. 1995). However, if galaxies and mass particles share the same velocity field so that the biasing relation remains constant along the streamlines, then the problem can be easily circumvented by specifying the biasing scheme at the present epoch (Nusser and Branchini 2000).
Within the standard lore of galaxies embedded in a virialized halo of dark matter that grow through hierarchical merging of smaller systems, neglecting the internal structure of objects is an assumption that is best justified a posteriori by showing how well Numerical Action methods can reproduce N-body velocities. Although the goodness of this assumption has been quantified by a number of numerical tests (e.g. Nusser and Branchini (2000) and Branchini Eldar & Nusser (2002)) little effort has been devoted to understand why numerical action methods can accurately reconstruct the velocity field of a large N-body simulation.
One of the reason for this success is that only $`5\%`$ of the points used in our reconstructions, that were randomly selected from the N-body simulation, belong to virialized regions where FAM reconstruction fails. Fortunately, the locality of this “virial effect” allows us to partially circumvent this problem by applying a modest spherical tophat smoothing of $`2h^1\mathrm{Mpc}`$ to the FAM-predicted velocities. This tophat filter has been specifically important for the quantitative aspects of our study, where such systematic problems may sort distorting conclusions. This smoothing has been invoked in quantitative comparisons between FAM and $`N`$-body velocities presented in this work, in particular in the regression analyses.
Little is known about the ability of numerical action methods to reconstruct the orbits of virialized systems. Indeed, when applied to extended objects rather than point masses, numerical action methods follow a single center of mass point per virialized objects, completely neglecting its merging history. Some argument can be given to back our choice of neglecting the internal structure of virialized objects. First of all, after tracing back the merging history of virialized halos in N-body experiments a simple visual inspection reveals that particles ending up in the same halo at z=0 are contained within regions with simple boundaries at high redshifts. As a consequence, high order terms in the gravity potential about the halo center of mass are probably rather small. This probably minimize the role of major mergers whose rate for galaxy-size halos peaks in the redshift range 2-4 (Volonteri, Haardt and Madau 2003) while peculiar motions mostly develop at $`z<2`$ (Branchini and Carlberg 1994). These qualitative arguments clearly need to be confirmed by appropriate numerical analyses similar to that of Branchini and Carlberg (1994) but extending out to scales of cosmological interest.
### 3.5 Ordered reconstructions
To obtain an idea of the level of improvement obtained through the use of successively higher order FAM evaluations, Figure 1 depicts 2D projections of the corresponding FAM particle orbit reconstructions within a local spherical volume of $`30h^1\mathrm{Mpc}`$. The black dots indicate the positions for each object in the sample, while the lines emanating from each dot represent the computed trajectories followed by these objects as they moved towards their present location. The illustrated configuration is taken from one of constructed mock catalogs, and resembles that of the Local Universe (see section 5.2.1). Each successive FAM reconstruction is based on the same (present-day) particle distribution. The four frames correspond to successively higher order FAM approximations, involving an increasing number $`N_f`$ of basis functions $`q_n(D)`$. The top-left panel shows FAM reconstructed orbits with $`N_f=1`$, which in fact corresponds to the conventional first order Zel’dovich approximation and thus represent the orbits that would have been obtained by the PIZA method (Croft & Gaztañaga croftgazt (1997)). These are followed by panels with $`N_f=2`$ (top right), $`N_f=3`$ (bottom left) and $`N_f=6`$ (bottom right). They show a clear and steady improvement towards the $`N_f=6`$ FAM evaluation. Testing proved that even higher order computations did not yield improvements significant enough to warrant the extra computational effort.
In summary, the galaxy orbits in our Local Universe environment are found at a minimum of the action which is not too far, yet different, from linear theory predictions. The FAM technique thus yields a significant modification of the recovered galaxy orbits and peculiar velocities for configurations that evolved well beyond the linear regime (see e.g. Figure 1). Potentially its ability to deal with nonlinear circumstances might even prove of benefit to recover sets of cosmological initial conditions satisfying nonlinear observational constraints at the present day, which indeed has been suggested by Goldberg & Spergel goldsperg (2000).
### 3.6 LAP and External forces
The original cosmological Least Action Principle formulation by Peebles ( peebles89 (1989)) considered a fully self-consistent, i.e. isolated, system of point masses. For practical reasons, the original implementation had to be restricted to systems of at most a few dozen objects. Almost exclusively, the Local Group of galaxies formed the focus of these LAP studies (Peebles peebles89 (1989)peebles90 (1990)peebles94 (1994), Dunn & Laflamme dunnlaf93 (1993)).
While these studies did indeed yield a substantial amount of new insight into the dynamical evolution of the Local Group, the issue of incorporating the dynamical influence exerted by external mass concentrations remained a major unsettled question. External forces do represent a significant component of the dynamics of the Local Group, as had been shown by Raychaudhury & Lynden-Bell ( raychlynd (1989)). They established beyond doubt that the Local Group cannot be considered a tidally isolated entity, and demonstrated that the Local Group is acted upon by an appreciable quadrupolar tidal force. The resulting tidal torque appears to be responsible for the large angular momentum of the Local Group as a whole, as Dunn & Laflamme (dunnlaf93 (1993)) showed in an elegant and pioneering analysis using orbits computed by the LAP variational method. They confirmed that the tidal influence of the external matter distribution is indeed essential to explain its present angular momentum.
In the course of time various strategies emerged to include external dynamical influences. The nature of these methods are mainly set by the character of the physical system under consideration, and to some extent was dependent on the available computational resources. Three strategies are outlined below.
#### 3.6.1 Directly Including External Masses
To incorporate the external tidal influence within the LAP analysis the work by Peebles (peebles89 (1989), peebles90 (1990), peebles94 (1994)), Peebles et al. peebles2001 (2001) and Dunn & Laflamme (dunnlaf93 (1993)) attempted to identify a few principal external mass concentrations which would be responsible for the major share of the external gravitational influence. While in his first LAP study Peebles (peebles89 (1989)) considered the Local Group mainly as an isolated system, sequel studies (Peebles peebles90 (1990), Peebles et al. peebles2001 (2001)) attempted to assess the possible external influence by neighbouring matter concentrations. In Peebles peebles90 (1990) he attempted to condense the external tidal force into two nearby mass concentrations, the Sculptor and Maffei group, each modeled as a single mass. Both were incorporated as 2 extra particles, with properly scaled masses, within the action $`S`$ in order to take them along in a fully self-consistent variational evaluation. A similar approach was followed by Dunn & Laflamme (dunnlaf93 (1993)), be it that they included five galaxies/groups in the local cosmic neighbourhood which arguably contribute a significant torque on the Local Group. Also in a later application (Peebles et al. peebles2001 (2001)) this approach was followed, be it with an extensive outer region between $`4h_{75}^1`$Mpc and $`20h_{75}^1`$Mpc whose mass distribution was condensed into a coarse sample of some 14 major external objects.
This “self-consistent” strategy is feasible to pursue within the context of the original, computationally intensive, LAP implementation. This approach may therefore be followed in LG resembling situations in which a few objects suffice to represent the main aspects of a system’s dynamical evolution. On the other hand, cosmic systems of a considerably larger scale than the Local Group would in general be too demanding for. Supercluster sized regions, with scales of up to a few tens of Megaparsec, count many more individual objects than a galaxy group. These systems have also not yet reached a formation stage so advanced that they have already largely decoupled from the global Hubble expansion, so the resulting external gravitational influence is usually the shared responsibility of a large number of external matter concentrations. Accounting for the large-scale tidal field would quickly become prohibitively expensive in terms of the computational effort for conventional LAP analyses.
#### 3.6.2 Inserting External Tidal Potential
An alternative strategy is to incorporate the external gravity in the LAP scheme via an approximate expression for the external contribution. This may be most directly achieved by inserting an extra external tidal potential term $`\mathrm{\Phi }_{tidal}(t)`$ in the action $`S`$ (eqn. 6). As on sufficiently large, Megaparsec, scales we may expect this term to evolve according to linear gravitational instability perturbation growth,
$$\mathrm{\Phi }_{tidal}(t)=D_\varphi (t)\mathrm{\Phi }_{tidal}(t_0)=\frac{D(t)}{a(t)}\mathrm{\Phi }_{tidal}(t_0),$$
(9)
in which $`\mathrm{\Phi }_{tidal}(t_0)`$ is the tidal term at the present cosmic epoch and $`D_\varphi =(D/a)`$ the linear growth term for the gravitational potential (the growth factor $`D_\varphi `$ and cosmic expansion factor $`a(t)`$ are set to be equal to unity at the present epoch, $`D(t_0)=a(t_0)=1.`$). Thus, instead of evolving it self-consistently along with the considered system, the external field is determined at one epoch – usually the present one – and then incorporated as an extra linearly growing gravity field term $`\mathrm{\Phi }_{tidal}(t)`$ in the action $`S`$:
$`S={\displaystyle _0^{t_0}}𝑑t`$ $`[`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{m_ia^2}{2}}\left({\displaystyle \frac{d𝐱_i}{dt}}\right)^2+{\displaystyle \frac{G}{a}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{m_im_j}{|𝐱_i𝐱_j|}}+`$ (10)
$`+{\displaystyle \frac{2}{3}}\pi G\rho _ba^2{\displaystyle }m_i𝐱_i^2\mathrm{\Phi }_{tidal}(t)].`$
There are various possibilities to compute the tidal potential term $`\mathrm{\Phi }_{tidal}`$, usually from the current mass distribution. One option is to compute it directly from a sample of $`M_{ext}`$ external objects which is deemed responsible and representative for the major share of the external tidal force field,
$$\mathrm{\Phi }_{tidal}(t_0)=\underset{i}{}m_i\left\{\frac{G}{a}\underset{j=1}{\overset{M_{ext}}{}}\frac{m_j}{|𝐱_i𝐲_j|}\right\}.$$
(11)
Note that none of these external objects ($`j=1,\mathrm{},M_{ext}`$) is taken into account as far as the action of the system and the computation of their orbits is concerned, except for their “passive” role in determining $`\mathrm{\Phi }_{tidal}`$. An alternative approach is to insert an approximate analytical expression for $`\mathrm{\Phi }_{tidal}`$, in particular one including the dipolar and quadrupolar contributions, $`𝐝`$ and $`𝐓`$, to the tidal potential,
$$\mathrm{\Phi }_{tidal}(t_0)=\underset{i}{}𝐝𝐱_i+\frac{1}{2}𝐱_i𝐓𝐱_i.$$
(12)
Equivalently, one may chose to insert the corresponding expressions directly into the expression for the derivative of the action with respect to an expansion coefficient, $`S/C_{i,n}^\alpha =0`$, evidently equal to zero within this variational approach.
The first, “direct”, procedure (eqn. 11) was followed by Shaya, Peebles & Tully (shaya (1995)), who for the purpose of studying the velocity field within the surrounding $`30h^1\mathrm{Mpc}`$ modelled the relevant external mass distribution after the distribution of rich Abell clusters from Lauer & Postman (laupost (1994)). To some extent, Sharpe et al. (sharpe (2001)) operated along the same lines, be it that they added the resulting tidal term directly to the reconstructed velocities produced by the LAP procedure. However, while in principle exact, such a concentrated and static mass distribution may involve considerable uncertainties and can be highly sensitive to the uncertainties in the location of a few dominant point masses. As this spatial point distribution is supposed to form a suitable model for the underlying large scale matter distribution this may be even more worrisome.
Potentially more elegant may therefore be the modelling of a smooth tidal field along the line of the second procedure (eqn. 12), as suggested by Schmoldt & Saha (schmoldtsaha (1998)). The corresponding dipolar and quadrupolar term may then be based on the best available determinations of these parameters. On the other hand, when the LAP volume is comparatively large, the analytical approximation may represent an oversimplification of the force field, neglecting potentially important local variations within the external force field.
#### 3.6.3 Selfconsistent and Direct FAM approach
The indirect “potential” approach which we described above (eqn. 11 or eqn. 12) may not properly account for the temporal evolution of the external field in the case of nonlinearly evolving systems. The formalism assumes a static, merely linearly evolving, gravitational potential. However, the matter concentrations which generate the external tidal forces will themselves get displaced as the cosmos evolves. These displacements may be relatively minor for distant masses, but for the more nearby entities this may be entirely different. A detailed treatment of the external mass distribution will be necessary when the influence of the nearby external objects on the evolution of small “interior” regions is comparable to or even dominant over the selfgravity of the region. It will be equally crucial to follow the detailed whereabouts of nearby matter concentrations in the case of a large “interior” region in which a marked contrast between the central regions and the outer realms may result in a significantly different dynamic evolution.
This prompted us to follow an alternative and direct approach, a fully self-consistent strategy in which also the external matter concentrations are accounted for in the computation of the system of evolving particle orbits. Alongside that in the “local” region for which we seek to reconstruct the velocity field, also the system of objects in the exterior regions ($`30h^1\mathrm{Mpc}<r<100h^1\mathrm{Mpc}`$) are considered. Non-uniform manifestations of the external influence can only be included by pursuing such a direct and systematic approach. It is only through the availability of the FAM technology that we were enabled to do so for a Megaparsec system consisting of a large number of objects.
## 4 Cosmological scenarios
The mock catalogs on which we apply our Fast Action Minimization analysis are extracted from $`N`$-body simulations in two different cosmological settings. Their characteristics, in terms of their relevant parameters, are listed in Table LABEL:table:parameters. The table also lists the simulation specifications. The first scenario is a flat $`\mathrm{\Lambda }`$CDM model with a cosmological constant term $`\mathrm{\Omega }_{\mathrm{\Lambda },0}=0.7`$ ($`\mathrm{\Omega }_0=0.3,\mathrm{\Omega }_0+\mathrm{\Omega }_{\mathrm{\Lambda },0}=1.0,\mathrm{\Gamma }=0.25,n=1`$). The second model is a $`\tau `$CDM Einstein-de Sitter ($`\mathrm{\Omega }_0=1.0,\mathrm{\Omega }_{\mathrm{\Lambda },0}=0,\mathrm{\Gamma }=0.25,n=1`$) model, motivated by the decaying particle model proposed by Bond & Efstathiou (bondef91 (1991)). Both scenarios were chosen to be viable with respect to the current observational constraints, implying similarities in many overall properties and appearances, yet with some significant differences with respect to their dynamical repercussions. This may provide indications on whether the galaxy motions in our local cosmic neigbourhood do contain information on the structure formation scenario.
In both cases the amplitude of density fluctuations is normalized on the basis of the observed abundance of rich galaxy clusters in the local universe. This abundance depends on the magnitude of the matter field fluctuations on the mass scale characteristic for galaxy clusters. This translates into a dependence on the amplitude of density fluctuations on cluster scales modulated by the mean global matter density. A variety of studies (e.g. White, Efstathiou & Frenk wef93 (1993), also see Eke, Cole & Frenk eke (1996)) found that in order to yield the present-day cluster abundance the amplitude of density fluctuations in spheres of radius $`8h^1\mathrm{Mpc}`$, $`\sigma _8`$, and $`\mathrm{\Omega }_0`$ are related by
$$\sigma _8=0.55\mathrm{\Omega }_0^{0.6},$$
(13)
The resulting power spectra are depicted in Figure 2 (top left). On all scales, the density fluctuations in the $`\tau `$CDM scenario, represented by the dotted lines (for both $`P(k)`$ (green lines) and $`k^3P(k)`$ (blue lines)), are less pronounced than those of the $`\mathrm{\Lambda }`$CDM scenario: the two power spectra have a similar shape and differ by a simple scaling factor over the entire wavelength range. Visually, this is immediately reflected in the stark differences between the spatial galaxy distribution in the resulting mock catalogs. Figure 3 provides such a visual comparison. It shows the “external” PSC$`z`$ catalog mimicking galaxy distribution in three mutually perpendicular central slices in the case of the $`\mathrm{\Lambda }`$CDM scenario (top row), together with the same set of frames for a $`\tau `$CDM mock galaxy catalog (bottom row). On all scales, the $`\tau `$CDM galaxy distribution looks considerably more uniform than that in the $`\mathrm{\Lambda }`$CDM Universe. Not only is the clustering of galaxies in the $`\mathrm{\Lambda }`$CDM scenario more pronounced, it also delineates considerably larger structures, a manifestation of the power spectrum’s amplitude at the corresponding large wavelengths.
Because the higher average matter density in the $`\mathrm{\Omega }_0=1`$ $`\tau `$CDM Universe does almost fully compensate for the lower amplitude of the density fluctuations the resulting gravity and velocity perturbation fields in the $`\mathrm{\Lambda }`$CDM and $`\tau `$CDM scenarios are very similar. The velocity power spectra $`k^3P_v(k)`$ are shown in the bottom lefthand panel of Figure 2: their functional dependence is the same over the entire wavelength range. The larger mass corresponding to a given density excess in the $`\tau `$CDM Universe evidently effects a stronger gravitational force. The resulting large scale motions scale as $`f(\mathrm{\Omega }_0)\mathrm{\Omega }_0^{0.6}`$. This happens to be almost exactly the inverse of the average density perturbation amplitude scaling (eq. 13), which is proportional to $`\mathrm{\Omega }_0^{0.6}`$ (eqn. 13). While this is exactly the factor involved in the normalization of the power spectrum, in terms of $`\sigma _8`$, the lower level of density fluctuations gets precisely compensated by the higher amount of mass involved with them. This can be directly observed from the velocity power spectra $`P_v(k)`$ for the two scenarios (Fig. 2, lower lefthand frame). The velocity power spectra for both scenarios are exactly equal over the entire wavelength range, both in functional dependence as well as in amplitude. Note that also the gravity perturbations in the $`\mathrm{\Lambda }`$CDM scenario are substantially stronger than those in the $`\tau `$CDM cosmology: because they scale with $`\frac{3}{2}\mathrm{\Omega }H_0^2`$ and the amplitude of the density perturbations, which according to eq. 13 is $`\mathrm{\Omega }^{0.6}`$, the average peculiar gravitational acceleration is proportional to $`\mathrm{\Omega }^{0.4}H_0^2`$.
The comparison between $`k^3P(k)`$ (Fig. 2, top panel, blue lines) and $`k^3P_v(k)`$ (Fig. 2, bottom panel) in the same figure shows the shift of the velocity perturbations, with respect to the density perturbations, towards a more large-scale dominated behaviour. This follows directly from the continuity equation, connecting the velocity and density perturbations such that the velocity power spectrum relates to $`P(k)`$ through $`P_v(k)P(k)k^2`$.
The large-scale behaviour of the (linear) velocity perturbation field immediately illuminates the difficulty in tracing the full array of matter inhomogeneities responsible for the cosmic motions within a specific cosmic region. To account for all noticeable contributions it is necessary to probe out to large depth. This is manifestly evident for the first order component in the externally induced flow, the “bulk flow” $`v_{bulk}`$. A measure for the expected bulk flow within a (tophat) spherical region of size $`R_{TH}`$,
$$𝐯_{bulk}(𝐱)_Vd𝐱^{}𝐯(𝐱^{})W_{TH}(𝐱𝐱^{},R_{TH}),$$
(14)
is represented by the (root square) average value $`\sigma _v`$, whose value may be inferred from the Fourier integral
$`\sigma _v(R_{\mathrm{TH}})`$ $`=`$ $`H_0f(\mathrm{\Omega }_0)\sigma _1`$
$``$ $`H_0f(\mathrm{\Omega }_0)\sqrt{{\displaystyle \frac{\mathrm{d}𝐤}{(2\pi )^3}P(𝐤)\widehat{W}_{TH}^2(𝐤)k^2}}.`$
In these relations, $`W_{TH}(𝐱,R_{TH})`$ and $`\widehat{W}_{TH}(𝐤)`$ are the expressions for the tophat window filter, spatially and in Fourier space, and $`\sigma _\mathrm{𝟏}`$ is the spectral moment $`\sigma _j`$ for $`j=1`$ (see Bardeen et al. bbks (1986), henceforth BBKS).
How substantial the large scale origin of the bulk flow is may be readily appreciated from figure 2 (top centre). Because the linear character of fluctuations on large scales, the spectral $`\sigma _v`$ (eq. 4) does provide a reasonable order-of-magnitude estimate of the magnitude of the large-scale bulk motions. The figure shows the estimated bulk flow amplitudes, $`\sigma _v`$, as a function of the (tophat) window radius: the bulk flow is clearly a large scale phenomenon, converging only very slowly towards large spatial scales. In both the $`\mathrm{\Lambda }`$CDM scenario and the $`\tau `$CDM scenario the externally induced bulk flow on a scale of $`30h^1\mathrm{Mpc}`$ will be in the order of $`200\mathrm{km}\mathrm{s}^1`$. Of this overall bulk flow, more than $`100\mathrm{km}\mathrm{s}^1`$ has to be ascribed to inhomogeneities on scales exceeding $`200h^1\mathrm{Mpc}`$ ! When assessing the motions in a local volume of $`30h^1\mathrm{Mpc}`$ radius, in terms of relative external contributions, inhomogeneities on a scale larger than $`100h^1\mathrm{Mpc}`$ still contribute more than $`25\%`$ of the total while the ones larger than $`200h^1\mathrm{Mpc}`$ are still responsible for more than $`10\%`$ (see Fig. 2, centre bottom). We should therefore expect to find substantial external contributions in the $`\mathrm{\Lambda }`$CDM and $`\tau `$CDM simulations. Note that this relative contribution to the bulk flow, the “bulk convergence”, is defined as the relative contribution by matter perturbations within a radius $`R_{TH}`$ to the externally induced bulk flow on a scale of $`30h^1\mathrm{Mpc}`$ (the size of the NBG volume):
$$_{\mathrm{bulk}}\mathrm{\hspace{0.17em}1}\frac{\sigma _v(R_{TH})}{\sigma _v(30h^1\mathrm{Mpc})}.$$
(16)
The second order aspect of the velocity field which we seek to study is the induced velocity shear $`s_{ij}`$,
$$s_{ij}\frac{1}{2}\left\{\frac{v_i}{x_j}+\frac{v_j}{x_i}\right\}\frac{1}{3}(𝐯)\delta _{ij}.$$
(17)
Also the velocity shear reveals interesting and distinguishing differences between the $`\mathrm{\Lambda }`$CDM and the $`\tau `$CDM scenario. In the linear regime the expected magnitude of the shear tensor $`s_{ij}`$, on a tophat scale $`R_{TH}`$, may be evaluated through its direct proportionality to the tidal shear $`𝒯_{ij}`$. Quantifying $`s_{ij}`$ by means of its (root square ) average $`\sigma _s`$ (van de Weygaert & Bertschinger rienbert (1996)), we find
$$\sigma _s=H_0f(\mathrm{\Omega }_0)\sigma _0(R_{TH})\sqrt{\frac{1\gamma ^2}{15}}$$
(18)
in which the (dimensionless) spectral parameter $`\gamma \sigma _1^2/\sigma _0\sigma _2`$ is defined through the $`0^{th}`$, $`1^{st}`$ and $`2^{nd}`$ spectral moments $`\sigma _j`$ (see BBKS bbks (1986)). The predictions for the two cosmological scenarios are shown in topright frame of Fig. 2. With respect to the bulk flow there is a marked difference in coherence scale: the major contributors to the tidal shear are located at considerably closer distances than the sources of the bulk flow (fig. 2, cf. lower right with lower centre). Most of the shear inducing matter inhomogeneities are found within a radius of $`100h^1\mathrm{Mpc}`$, accounting for more than $`95\%`$ of its value (fig. 2, lower right). On a scale of $`30h^1\mathrm{Mpc}`$ we expect an external tidal shear of $`7\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ for both the $`\mathrm{\Lambda }`$CDM scenario and the $`\tau `$CDM model.
## 5 Mock catalog Construction and Analysis
### 5.1 $`N`$-body simulations
The two $`N`$-body simulations used in this work were carried out by Cole et al. (cole (1998)) within the context of an extensive study of PSCz catalogue resembling galaxy mock samples in a large variety of cosmological structure formation scenarios. They consist of $`192^3`$ particles in a computational box of $`345.6h^1\mathrm{Mpc}`$. They are dynamically evolved using an AP$`3`$M code in which the force is smoothed with a softening parameter of $`0.27h^1\mathrm{Mpc}`$.
The purpose of this study is a demanding task for truely representative $`N`$-body simulations. The $`N`$-body simulations should provide an optimal compromise between a high mass resolution on the small scale side and, on the large-scale side, a cosmic volume large enough to be dynamically representative. The large dynamic range requirement involves a mass resolution refined enough to resolve mass entities comparable to galaxies. This translates into an average inter-particle separation that needs to be smaller or comparable to that of galaxies in real observational catalogues. On the other hand, the simulations have to extend over a cosmic volume which is large enough to incorporate the major share of the gravitational influence exerted by the inhomogeneous cosmic matter distribution. Given the slow convergence of the bulk flow and its large coherence scale this is particularly challenging, and will be in the order of several hundreds of Megaparsec (see discussion in the previous section and fig. 2). Although hardly any current $`N`$-body simulations would fully fulfill the dynamic range requirements, the used $`N`$-body simulations do appear sufficiently adequate for a meaningful investigation of the relevant systematic trends and effects. This remains true in a qualitative sense, even though on the basis theoretical arguments (see e.g. Fig. 2) and observational indications (e.g. Hudson et al. hudsv (2004)) we know there may be substantial bulk flow contributions stemming from even larger spatial scales.
In this respect it is important to note is that the mock catalog realizations in this work are constrained by the finite size of the simulation box. The practical repercussions of being confined to a limited simulation volume may be inferred from the dashed curves in Fig. 2. They show the corrections to the expected bulk flow and velocity shear predictions (solid curves) when only the inhomogeneities in the restricted volume of the $`345.6h^1\mathrm{Mpc}`$ simulation box are incorporated. Because perturbations on scales exceeding the fundamental scale of the box are absent, the realized power spectrum has a rather sharp and artificial large-scale cutoff: the limited boxsize $`L_{box}`$ implies a cutoff in the power spectrum at low wavenumber $`k_{box}=2\pi /L_{box}`$. From Fig. 2 we can conclude that this correction is particularly apt for the bulk flow, predictions for the velocity shear seem hardly affected. As a consequence, on scales over $`100h^1\mathrm{Mpc}`$ the bulk flows in the realized $`N`$-body simulations will be severely repressed and far from representative. Although large-scale mode adding procedures have been proposed to partially remedy this situation (Tormen & Bertschinger torbert (1996) and Cole cole97 (1997)), our $`\tau `$CDM and $`\mathrm{\Lambda }`$CDM simulations did not include such MAP (mode adding procedure) extensions. Conclusions with respect to the convergence of the FAM reconstructed velocity flows should therefore be referred to with respect to the suppressed velocity power spectrum indigenous to our $`N`$-body simulations (notice that this dynamic range issue is truely cumbersome to nearly any study attempting to assess velocity flows in computer simulations).
### 5.2 Mock Catalog Construction
From the full $`N`$-body simulations we extract mock catalogs made to resemble the local Universe. The $`\mathrm{\Lambda }`$CDM and $`\tau `$CDM $`N`$-body simulations are processed through specified observational masks to imprint the required characteristics on the resulting mock catalogs. We distinguish two types of mock catalogs. From each $`N`$-body simulation we extract ten different “local” mock catalogs mimicking the NBG catalogue and, with these “local” samples representing their interior, ten different “extended” samples resembling the PSC$`z`$ catalogue.
The “local” class of mock samples is meant to sample the mass distribution within a $`30h^1\mathrm{Mpc}`$ region in and immediately around the Local Supercluster. These catalogs constitute volume-limited galaxy samples mimicking the Nearby Galaxy Catalog of Tully (tullycat (1988)). Mock catalogs of the second type are designed to account for the mass distribution out to distances of $`100h^1\mathrm{Mpc}`$. These “extended” samples represent flux-limited samples, for which we take the IRAS PSC$`z`$ galaxy catalog (Saunders et al. saunders (2000)) as template. The PSCz sample is not only ideal for our purposes in that it covers one of the largest volumes of the Universe amongst the available galaxy redshift surveys, but also in that it concerns a survey covering a large fraction of the sky and involves a well-defined uniformity of selection. Assuming that on large linear scales IRAS galaxies define an unbiased tracer of the underlying dark matter, Hamilton etal hamilton2000 (2000) found that its real-space power spectrum is consistent with that of a COBE-normalized, untilted, flat $`\mathrm{\Lambda }`$CDM model with $`\mathrm{\Omega }_m=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$.
In both flux-limited and volume limited samples the mass of mock galaxies have been rescaled to the value of $`\mathrm{\Omega }_m=0`$ of the parent N-body simulations, listed in Table LABEL:table:parameters.
In Table 2 we have listed the main characteristics of all mock catalogs used in this work. .
In constructing the mock samples galaxies were identified with N-body particles selected randomly, exclusively according to the catalog selection criteria. Therefore, we did not attempt to include bias descriptions to model possibly relevant differences in the spatial distribution of dark matter and galaxies. This is different from the original use of the simulations (Cole et al. cole (1998)), in which various bias descriptions were invoked to construct artificial galaxy samples whose two-point correlation function and large-scale power spectrum largely matched that of the APM survey (Maddox et al. apmmaddox (1996)). The analysis of the small-scale nonlinear power spectrum of the PSCz by Hamilton & Tegmark ( hamilton2002 (2002)) even implies the bias on small scales to be very complex, involving a scale-dependent galaxy-to-mass bias. We, however, prefer not to include an extra level of modelling prescriptions. Our interest concerns the kinematics and dynamics of the matter distribution in the Local Universe, and the velocities of galaxies are thought to reflect these almost perfectly: they are mere probes moving along with the underlying dark matter flows, irrespective of their particular bias relation with respect to the dark matter distribution. The sole strict assumption is therefore that of having no velocity bias (Carlberg et al. carl90 (1990)), which on the large-scale Megaparsec scales at hand should be a more than reasonable approximation.
#### 5.2.1 Mock NBG catalogs
The mock NBG catalogs are obtained by extracting spherical volumes of $`30h^1\mathrm{Mpc}`$ from the $`N`$-body simulation particle distribution. The positions of the spheres in the parent simulations are not random but chosen to mimic as close as possible the characteristics of a Local Group look-alike region. Therefore, each mock catalog is centered on a particle moving at a speed of $`625\pm 25`$ km s<sup>-1</sup>, residing in a region in which the shear within $`5h^1\mathrm{Mpc}`$ is smaller than 200 km s<sup>-1</sup>, where the fractional overdensity measured within the same region ranges between -0.2 and 1.0.
The velocity vector of the central particles defines a Galactic coordinate system and a Zone of Avoidance. Particles within the Zone of Avoidance are removed and substituted with a population of synthetic objects distributed using a random-cloning technique (Branchini et al. branchini99 (1999)). The Zone of Avoidance \[ZA\] in the mock samples is designed to mimic that of the PSC$`z`$ catalog (Saunders et al. saunders (2000)) and is smaller than the one of the real NBG catalog (Tully tullycat (1988)).
Each spherical region contains on average $`2\times 10^4`$ particles. This set of particles is randomly resampled in order to produce an unbiased catalog of around 2800 objects, a number that matches that of the galaxies in the real NBG catalog (Table 2). This procedure preserves, within shot noise errors, density fluctuations and thus does not alter the orbit reconstruction.
These NBG mimicking mock catalogs define volume-limited galaxy samples, so that the number of objects within a distance $`x`$ therefore increases as $`x^3`$. This is indeed what the resulting realizations yield, as may be discerned from the central part of the corresponding histogram in Figure 4 ($`x30h^1\mathrm{Mpc}`$).
#### 5.2.2 Mock PSC$`z`$ catalogs
The second set of mock catalogs was obtained by carving out spherical regions of radius $`100h^1\mathrm{Mpc}`$ from the $`N`$-body simulations. Each of these new mock samples is centered on the same central position as that of corresponding NBG mock catalogs, with which they share the objects within the central $`30h^1\mathrm{Mpc}`$.
While the central $`30h^1\mathrm{Mpc}`$ region coincides with the NBG mock sample, the particle distribution in the external region ($`30h^1\mathrm{Mpc}<x<100h^1\mathrm{Mpc}`$) is supposed to mimic that of galaxies in the flux-limited IRAS PSC$`z`$ catalog. To achieve this the objects beyond $`30h^1\mathrm{Mpc}`$ were selected from the $`N`$-body particle samples according to the PSC$`z`$ selection function used by Branchini et al. (branchini99 (1999)):
$$\psi (x)=Ax^{2\alpha }\left(1+\frac{x^2}{x_{}^2}\right)^\beta \text{if }x>x_s.$$
(19)
In this expression $`x`$ is the distance of the galaxy in $`h^1\mathrm{Mpc}`$, while the parameters have the values $`\alpha =0.53`$, $`\beta =1.8`$, $`x_s=10.9h^1\mathrm{Mpc}`$, and $`x_{}=84h^1\mathrm{Mpc}`$ (Branchini branchini99 (1999)). The selection function $`\psi (x)`$ defines the relative number density of the galaxy sample with respect to the number density of the $`N`$-body particle simulation. While the number density of PSC$`z`$ galaxies is equal to that of the particles in the simulation at $`10.9h^1\mathrm{Mpc}`$, the amplitude $`A`$ of $`\psi (x)`$ is normalized such that $`\psi (10.9)=1`$.
Two additional steps concern the treatment of “Zone of Avoidance” objects and the evening of the matter density throughout the full external sample volume. A first step is the processing of sampled objects in the Zone of Avoidance such that the resulting sample conforms to a reality resembling situation. The ZA “removal$`+`$substitution” is implemented in the same way as in the case of the NBG mock catalog construction, with the replacement achieved with the same random-cloning technique. Finally, in order to guarantee a uniform average mass density throughout the volume, the mass of the objects in the flux-limited external object sample ($`30h^1\mathrm{Mpc}<x<100h^1\mathrm{Mpc}`$) has been scaled by the inverse of the selection function $`\varphi (x)`$.
#### 5.2.3 Mock catalog realizations
From both figure 5 and figure 6 one can obtain an impression of the spatial context of the local NBG mock sample within the wider environment of the surrounding $`100h^1\mathrm{Mpc}`$ PSC$`z`$ sample. To visually appreciate the selection criteria of the catalogs, and their interrelationship, figure 5 shows a three-dimensional view of one set of the $`\mathrm{\Lambda }`$CDM mock catalogs, extracted from the particle distribution in the $`N`$-body simulation of structure formation in a $`\mathrm{\Lambda }`$CDM scenario. Emanating from the full PSC$`z`$ \+ NBG mimicking galaxy samples in the top righthand cube is a row of two cubes showing the content of the external PSC$`z`$ mock catalog (righthand) and the content of the central NBG mock catalog (lefthand). Figure 6 elaborates on this, and shows the projected particle distribution in the same PSC$`z`$ \+ NBG mock catalog (left panel) while focusing in on the central region (right panel). The circle (left panel) indicates the boundary of the volume-limited region comprised by the mock NBG galaxy sample, which in the right panel has been enlarged to show the corresponding velocity field within this NBG region. Velocities of objects within a 10$`h^1\mathrm{Mpc}`$ thick slice are shown by means of arrows whose size is proportional to the amplitude of the galaxy velocity components within this slice.
As a matter of test, we checked the distance distribution of the resulting mock galaxy samples. The histograms of the resulting mock catalog distributions are shown in Figure 4. The upper part of Figure 4 shows the number of galaxies – averaged over all PSC$`z`$ mock catalogs for both two cosmological models – as a function of distance $`x`$ over the full range $`x100h^1\mathrm{Mpc}`$. Clearly visible is the discontinuity at $`30h^1\mathrm{Mpc}`$, marking the transition from volume-limited NBG-like region to the PSC$`z`$-like flux-limited outer region. For comparison, the solid line shows the theoretically expected counts (Eq. 19). The generated mock samples appear to match the expected distance distribution rather well. This is further underlined by the fractional difference between observed and expected counts, $`\mathrm{N}_{\mathrm{obs}}/\mathrm{N}_{\mathrm{exp}}1`$, shown in the lower part of Figure 4. The fractional difference between mock samples displays a perfect featureless scatter pattern: Poisson noise free of systematic effects.
### 5.3 Mock Catalog Analysis
#### 5.3.1 FAM velocity field reconstructions
On the basis of the FAM reconstructions of the galaxy velocities and the comparison with the true velocities – i.e. those in the original $`N`$-body simulation – we assess to what extent the matter distribution within the confines of each different mock galaxy sample does contribute to the total velocity of the galaxies. In these idealized circumstances of the $`N`$-body world the galaxy positions and velocities are known to perfect accuracy, thus circumventing the need to investigate the effects of measurement errors and deceptive systematic biases in the galaxy peculiar velocities. This should provide us with a better understanding of the nature and magnitude of genuine physical influences.
Three different velocities are accorded to each galaxy located within the “local” spherically shaped $`30h^1\mathrm{Mpc}`$ NBG region. The first velocity is that of the “true”, $`N`$-body velocity. For each of the in total 20 NBG mimicking galaxy mock catalogs, the FAM reconstructions produce two additional velocity estimates. One FAM velocity results from the application of the FAM analysis to the restricted inner $`30h^1\mathrm{Mpc}`$ NBG-like region itself. The second FAM based velocity is obtained on the basis of the FAM analysis on the extended, “full”, $`100h^1\mathrm{Mpc}`$ PSC$`z`$ survey resembling sample (in which the “local” NBG sample occupies the interior central region). In the following, we will indicate these FAM velocities by the names of FAM<sub>30</sub> and FAM<sub>100</sub> velocities.
#### 5.3.2 FAM<sub>30</sub> versus FAM<sub>100</sub> reconstructions
The mutual comparison between each of the three different galaxy velocities – the FAM<sub>30</sub>, the FAM<sub>100</sub> and the full $`N`$-body velocities – is expected to yield abundant information on the dynamics and development of the structure in the interior $`30h^1\mathrm{Mpc}`$ region:
The FAM<sub>30</sub> velocities are the galaxy velocities which would have been the product of the combined gravitational interaction of – solely – the matter concentrations within the central $`30h^1\mathrm{Mpc}`$ volume. Any deficiency with respect to the “real” $`N`$-body velocity of each galaxy has to be ascribed to the gravitational impact of matter inhomogeneities outside the local NBG region.
By tracing the mass distribution further out to a distance of $`100h^1\mathrm{Mpc}`$, invoking the matter distribution in the complete PSC$`z`$ mimicking mock samples, we will then evaluate the extent to which matter inhomogeneities within a 100$`h^1\mathrm{Mpc}`$ scale are able to account for the motions within the local 30$`h^1\mathrm{Mpc}`$ region. From this we can infer in how far the external influence over the local region can be ascribed to matter fluctuations situated between a radius of 30$`h^1\mathrm{Mpc}`$ and 100$`h^1\mathrm{Mpc}`$.
In this study we also have to take into account the fact that a single 30$`h^1\mathrm{Mpc}`$ region cannot be considered representative for the whole Universe, and generic conclusions on the basis of the kinematics within a single 30$`h^1\mathrm{Mpc}`$ volume cannot be drawn. This is also true for the the NBG mock samples in this work, even though they were selected according to some strict criteria (see section 5.2). Analysis and conclusions will therefore be based on a straightforward average over the 10 different $`30h^1\mathrm{Mpc}`$ mock samples which were constructed for each cosmological scenario. The dispersion in the extracted parameter values will provide a reasonable estimate for their significance.
#### 5.3.3 Analysis of Reconstructions
The basic product of the FAM reconstructions are velocity maps, in essence a velocity vector at the location of each galaxy in the sample. Our analysis consists of three different but complementary tracks. The first and most straightforward one is the visual inspection of the resulting velocity vector maps. It provides a direct impression of the extent to which a FAM reconstructed field reproduces the true velocities. Also, it will provide a direct impression of a spatial coherence in the differences between true and reconstructed field, which is an incisive way to uncover systematic contributions like e.g. a bulk flow component.
The second examination is a strictly local analysis, a pure point-to-point comparison between the velocities predicted by the FAM reconstructions on the one hand and the “true” $`N`$-body velocity of the same object on the other hand. To some extent, the analysis by means of scatter plots is the most direct and objective quantitative comparison between two fields. Various velocity related quantities will be assessed in this fashion. Note that these localized comparisons cannot address the presence of spatial coherence in the cosmic flows (even though they may uncover systematic effects caused by external influences).
Finally, the third track is targeted towards a factual description of the spatial coherence within the velocity fields or, rather, in the residual fields between the “true” velocities and the reconstructed velocities. Systematic trends in these residual fields are interpreted as manifestations of external forces. Of these we shall determine the first-order – bulk flow – and second order – velocity shear – components.
## 6 FAM velocity vector maps
For reasons of consistency and to achieve optimal transparency the illustrated velocity vector maps in the following discussion all concern the same mock sample of NBG calculations. For the illustration of the FAM<sub>30</sub> (Fig. 7) and the FAM<sub>100</sub> (Fig. 8) reconstructions we use one of the $`\mathrm{\Lambda }`$CDM $`30h^1\mathrm{Mpc}`$ NBG mock catalogs. It is the same galaxy sample that was shown in 3-D in Fig. 5 and in projection along the “x-y” plane in Fig. 6.
The vector maps in Figure 7 and Fig. 8 depict the projections of the raw unsmoothed galaxy velocities, for galaxies within a central slice of $`10h^1\mathrm{Mpc}`$. The size of the arrows is proportional to the amplitude of the peculiar velocity component within this slice, each arrow starting at the location of the galaxy. Both figures consist of three successive rows. The velocity maps in the first row correspond to the “real” world of the $`N`$-body simulation. The second row depicts the velocity maps for the FAM reconstructions, the FAM<sub>30</sub> reconstruction in Fig. 7 and the FAM<sub>100</sub> reconstruction in Fig. 8. The last row shows the resulting residual velocity vector fields,
$$𝐯_{\mathrm{res}}𝐯_{\mathrm{FAM}}𝐯_{\mathrm{Nbody}},$$
(20)
the vector difference between the $`N`$-body velocities and the corresponding FAM velocity reconstructions, \[$`N`$-body - FAM<sub>30</sub>\] and \[$`N`$-body - FAM<sub>100</sub>\].
Each row has three panels, containing the vector maps in the three mutually perpendicular “central” slices. Each plane is identified by means of the index combination “x-y”, “x-z” or “y-z” (top figure), the index pair identifying the horizontal and vertical axis along which the panel is seen. Imagining these three planes passing through the centre of the $`30h^1\mathrm{Mpc}`$ NBG volume provides a spatial impression of the full 3-D velocity field. Note that here the choice of Cartesian coordinate system does not have any special significance, arbitrarily set by the axes of the total $`345.6h^1\mathrm{Mpc}`$ simulation box (the “fundamental” box) from which the mock catalogs were distilled. This is unlike vector maps (e.g. Fig. 13) in some later sections.
### 6.1 $`N`$-body sample: the “observed” velocities
The velocity vectors in the top row vector maps depict the “real” $`N`$-body velocities of the “galaxies” located within the three “central” slices (the same for Fig. 7 and Fig. 8). The galaxy distribution is characterized by a few dense, massive and virialized clumps, visible as high concentrations of large and randomly directed velocity vectors. The truely massive concentration visible in the lower left of the $`xy`$ panel is part of a superstructure extending beyond the boundaries of the NBG region. It represents a major and dominant source for the motions in this area. This may be appreciated from the observed velocity flow towards this clump and the overall distortion of the flow in its vicinity. The large configuration visible in the “y-z” slice contains several dense compact regions embedded in a ridge-like structure running curvedly from the lower righthand corner to a location slightly left from the centre. At least partially related to this mass concentration in and around the ridge is the bulk flow along the right-to-left direction.
Overall, the “x-y” and “y-z” vector maps indicate the presence of a dominant coherent “bulk flow” pattern which can be traced throughout the whole NBG volume. By coincidence, the orientation of the coordinate axes is such that the direction of the bulk flow is almost perfectly aligned along the “y”-axis: for this particular mock sample the “y”-axis does represent a physically significant direction defined by the streaming pattern itself. The bulk flow seems to be directed towards some (fictitious) point outside the local $`30h^1\mathrm{Mpc}`$ region.
A dominant and conspicuous coherent flow pattern also characterizes the $`xz`$ velocity vector map. While the flow in the two other planes seems to be almost exclusively dominated by a bulk flow, here the pattern has a more complex geometry, readily recognizable as a typical “velocity shear” pattern. The specific shearing motion in this plane consists of a compressional component along the top lefthand to lower righthand direction, in combination with a dilational stretch along the perpendicular direction from the lower lefthand towards upper righthand corner.
### 6.2 NBG samples: FAM<sub>30</sub> velocity vector maps
The role of the local cosmic matter distribution on the motions in the local Universe is assessed on the based of the “FAM<sub>30</sub> velocities”. They are the peculiar velocities computed by FAM on the basis of the local matter distribution, supposedly reflected by the galaxies within the NBG catalogs. The corresponding reconstructed velocities are shown in the second panel row of Figure 7. With their final position as boundary condition, each velocity vector is located at the same galaxy position as in the $`N`$-body maps (top row). Note that the vector maps in Figure 7 and Figure 8, and also the later ones in Fig 12, Fig 13 and Fig 14, show the pure unsmoothed velocity vectors (and do not “correct” for the virialized regions).
The FAM<sub>30</sub> velocity maps are distinctly different from the corresponding $`N`$-body velocity maps (top row): a coherent flow pattern is almost entirely absent. The FAM<sub>30</sub> reconstructions obviously did not recover the strong bulk flow observed in the $`N`$-body velocity maps, nor the striking shear pattern in the $`xz`$ plane. Because the FAM<sub>30</sub> velocity field reconstructions solely relate to the matter distribution within the inner 30$`h^1\mathrm{Mpc}`$ NBG region, this indicates that the major share of coherent bulk flow and the velocity shear are due to the matter distribution outside the central 30$`h^1\mathrm{Mpc}`$. This is most readily apparent in the velocity residual maps \[$`N`$-body - FAM<sub>30</sub>\], the difference between the $`N`$-body and the FAM<sub>30</sub> velocity vector fields (bottom row of Fig. 7). In the residual field \[$`N`$-body - FAM<sub>30</sub>\] we recognize the same characteristic flow patterns, strong spatial correlation, long-range coherence and overall morphology as in the full $`N`$-body velocity field. This represents convincing evidence for the external origin of the large-scale “bulk” and “shear” component in the local velocity flow.
Prominently visible in the residual velocity field is the strong bulk flow along the “y”-axis. Overall, the spatial pattern of the residual bulk flow appears to reproduce that of the $`N`$-body flow field. However, some minor yet significant differences between the residual and the full $`N`$-body bulk flow can be discerned. The amplitude of the corresponding velocities in the residual map is somewhat smaller than the equivalent $`N`$-body velocities: apparently part of the bulk flow is induced by the local NBG matter distribution. This does not seem to be true for the velocity shear: the shear patterns in the “x-z” plane of the residual and $`N`$-body velocity fields are almost identical (except for the virialized motions in high-density clumps). Apparently, the velocity shear component is almost exclusively due to external matter distribution. As a locally flattened matter configuration would induce an internal shear flow, this appears to imply a local matter distribution whose geometry is hardly flattened or elongated.
Closer inspection of the FAM<sub>30</sub> velocity field provides a more detailed view of the small-scale flow pattern mentioned above. In the “x-y” plane the large-scale ($`N`$-body) bulk flow has virtually completely disappeared. Instead, the dominant motion in the “x-y” plane is a streaming flow towards a prominent matter concentration within this region (lower left). On the other hand, in the “y-z” slice a trace of the $`N`$-body bulk flow along the “y” axis remains, be it that the corresponding velocities have considerably smaller amplitudes than their $`N`$-body counterparts. These local motions appear to be effected by the matter located along the lower ridge, supporting the impression that this feature is a local extension or outlier of the large-scale matter configurations responsible for the full bulk flow. Examination of the panels in Fig. 5 and Fig 6 indeed seems to suggest that the density ridge in the lower half of the “y-z” plane is indeed connected to structures just outside the NBG volume, while this perhaps may be true for the massive matter clump in the “x-y” plane too. This may not come as a surprise: the local matter distribution will to some extent be correlated with the external matter configuration so that the locally induced bulk flow is expected to reflect at least partially the full $`N`$-body bulk flow.
In summary, the inability of the FAM<sub>30</sub> reconstruction to recover the large-scale bulk flow and velocity shear is a consequence of the fact that they are a result of the action of the mass distribution on scales larger than the internal $`30h^1\mathrm{Mpc}`$ size region while the FAM<sub>30</sub> velocities are entirely and self-consistently determined by the mass distribution within this interior region. The residual \[$`N`$-body - FAM<sub>30</sub>\] maps, which are a model for the possible findings of a real-world observational campaign, provide the most elucidating illustration of their “external” origin. Even though they do provide convincing evidence for their external nature, they do not provide sufficient information to infer the identity and nature of the main source of the flow patterns. In principle, however, we may deduce a substantial amount of information on the basis of a careful quantitative analysis: the work by Lilje, Yahil & Jones (lilje (1986)) still sets a prime example. To this end, we will investigate the external matter distribution in the PSC$`z`$ $`100h^1\mathrm{Mpc}`$ sized regions.
As a final note, we point to the rather artificial nature of velocity vectors in the vicinity of the massive clump in the “x-y” slice as indicative for the self-consistent nature of the FAM reconstructions. Its location near the edge of the NBG volume even appears to have generated the rather contrived infall motions along the rim of the NBG sphere.
### 6.3 PSC$`z`$ samples: FAM<sub>100</sub> velocity vector maps
The contribution by the relatively nearby external matter agglomerations, within a distance of $`100h^1\mathrm{Mpc}`$, to the motions in the local Universe is investigated on the basis of the “FAM<sub>100</sub> velocities”. FAM produces these peculiar galaxy velocities on the basis of the galaxy sample in the full mock PSC$`z`$ galaxy sample, extending out to $`100h^1\mathrm{Mpc}`$ around the center of our local region. The corresponding reconstructed velocities are shown in the second panel row of Figure 8. It is the analogy for the “FAM<sub>100</sub> velocities” of Figure 7, and concerns the same $`30h^1\mathrm{Mpc}`$ central region (the NBG region is the central subregion of the PSC$`z`$ mimicking catalog). The FAM<sub>30</sub> maps showed the dominant influence of externally induced forces on the motions in the local 30$`h^1\mathrm{Mpc}`$ NBG region: on the basis of the FAM<sub>100</sub> maps we seek to assess whether the major share of the responsible external matter agglomerations may be identified within the realm of a PSC$`z`$ like volume.
Comparison of the first and the second row of panels in Fig. 8 shows the large degree of similarity between the FAM<sub>100</sub> velocities (panels 2<sup>nd</sup> row) and the $`N`$-body velocities (panels top row). Unlike the FAM<sub>30</sub> maps in Fig. 7 we find that the FAM<sub>100</sub> maps successfully reproduce most of the large-scale behaviour and most of the finer details of the $`N`$-body velocity field. The degree of similarity is particularly evident in the corresponding residual velocity field \[$`N`$-body - FAM<sub>30</sub>\] (bottom row panels). With the exception of the high-density virialized regions the residual velocities are very small and mostly randomly oriented: no significant spatial correlations and spatial coherence can be detected.
The detailed similarity between the $`N`$-body and the FAM<sub>100</sub> maps shows that it is sufficient to take account of the mass distribution out to $`100h^1\mathrm{Mpc}`$ for explaining, in considerable detail, the velocity flows in the local NBG volume. Moreover, the detailed rendering of the velocity field by FAM is a convincing demonstration of the capacity of the FAM technique to accurately describe the dynamics implied by the observed local galaxy distribution. The quantitative comparisons in the following sections will provide ample support to this claim.
Of course, the above conclusion is partially related to the realizations of the cosmological scenarios we have studied. The behaviour of the power spectrum $`P(k)`$ on large scales will considerably influence the generality of our findings. A power spectrum with more power on large scales would modify our findings: potentially it may be so that we need a representation of the matter distribution out to larger radii than $`100h^1\mathrm{Mpc}`$. In this respect it is important to note that the used $`N`$-body velocity fields do not have any contributions from wavelengths larger than $`175h^1\mathrm{Mpc}`$ (both for $`\tau `$CDM as well as $`\mathrm{\Lambda }`$CDM simulations, see Fig. 2). This merely for the technical reason of the simulation box imposing an upper limit to the scale on which we can represent $`P(k)`$. The extent to which this may influence our conclusions may be readily appreciated from Figure 2 (right column, top and bottom panel: compare solid lines with dashed ones). The velocity field perturbations of $`\tau `$CDM and $`\mathrm{\Lambda }`$CDM carry out considerably further than the fundamental scale of the simulation box, in particularly affecting the resulting bulk flows.
## 7 Point-to-point comparison
Scatter diagrams are used to assess the point-to-point comparisons between quantitative aspects of the “real” galaxy velocities in the original $`N`$-body samples and the computed velocities in the FAM<sub>30</sub> and FAM<sub>100</sub> reconstructions. This analysis is meant to be a direct, in principal local, assessment of systematic trends in the velocity flows in volume of the NBG sample. The comparisons involve a component of the “true” $`N`$-body velocity (abscissa) versus the equivalent quantity for either the FAM<sub>30</sub> or FAM<sub>100</sub> velocities, or of the corresponding residuals (ordinate). Since the objects had been artificially added to the Zone of Avoidance, any particles inside this region (see sec. 5.2.2) are excluded from these diagrams.
### 7.1 Systematics
If we neglect the small-scale sources in the deficiencies of FAM reconstructions, the differences between FAM<sub>30</sub> and FAM<sub>100</sub> scatter plots are mainly to be ascribed to the corresponding differences in the external gravitational influence acting over the two corresponding sample volumes. In the external gravitational influence the the corresponding leading velocity terms are the bulk flow $`\mathrm{v}_{\mathrm{bulk}}`$ and the velocity shear $`s_{ij}`$,
$$v_{\mathrm{Nbody},i}v_{\mathrm{FAM},i}+v_{\mathrm{bulk},i}+\underset{j=1}{\overset{3}{}}s_{ij}x\widehat{x}_j+\mathrm{}$$
(21)
In the above $`i,j`$ denotes the Cartesian component indices. The vectors $`\widehat{x}_j`$ represent the vector components along the Cartesian $`j`$ direction of the spatial unity vector oriented along the object position vector $`𝐱`$. Following this definition, $`x\widehat{x}_j`$ is the $`j`$-component of the position vector $`𝐱`$, with $`x`$ the distance of the object.
Systematic differences in FAM velocity-$`N`$-body velocity scatter diagrams are therefore to be attributed to differences in bulk flow, shear and possibly higher order contributions. Because each of these large-scale phenomena will manifest themselves in distinctly different ways, we seek to identify them from the scatter diagrams. An horizontal offset in the scatter diagram would be the trademark for a bulk flow component. Velocity shear would manifest itself as a distinctly characteristic correlation between residuals and velocities, although the prominence of this signal will be dictated by shear magnitude, configuration, and orientation with respect to the reference system (as is true for the bulk flow). In reality, the situation will be more intricate. Subtle correlations between small-scale and large-scale contributions will bring about a change in the slope of the scatter diagram of FAM reconstructed velocity components against their full $`N`$-body values (see 7.3.2).
### 7.2 Velocity Scatter Diagrams Analysis
Scatter diagrams are presented in three successive figures. The depicted scatter diagrams all relate to a $`\mathrm{\Lambda }`$CDM mock catalogue, and each of these point-to-point analyses relates to a different aspect of the velocity field reconstructions. Figure 9 contains four different panels, of which each contains two scatter diagrams: FAM<sub>30</sub> versus $`N`$-body quantity (left) and the equivalent FAM<sub>100</sub> versus $`N`$-body quantity (right). The diagrams in Fig. 10 focus on the correlations between these quantities and the scatter around regression relations. The figure addresses three velocity-related quantities, each taking one column of each 2 panels: the top one for the comparison of the FAM<sub>30</sub> components with their $`N`$-body counterparts, and the same for the FAM<sub>100</sub> components in the bottom frame.
A straightforward comparison is that between the Cartesian velocity components $`v_i`$ of the FAM reconstructed velocities and the $`N`$-body velocities (Fig 10, righthand column). Complementary regressions involve coordinate system independent aspects of galaxy velocities. These involve the velocity amplitude $`|v_{\mathrm{FAM}}|`$ (Fig. 9, top lefthand panel), the component of each FAM velocity parallel to the corresponding $`N`$-body velocity, $`v_{}`$ ($``$ v<sub>proj</sub>), and the additional perpendicular component $`v_{}`$ (Fig 10, first column). Misalignments between the real $`N`$-body velocity and the FAM velocity reconstructions should indicate in how far a reconstruction has been failing to take into account all relevant gravitational forces along the path of a particle. Systematic misalignments reveal themselves in the scatter diagram of the angle $`\theta `$ between the FAM velocity and the galaxies’ $`N`$-body velocity $`𝐯_{\mathrm{Nbody}}`$ (in Fig. 9 we plot $`\mu \mathrm{cos}(\theta )`$, bottom righthand panel). In terms of the character and systematics of the underlying physics and dynamics the residual velocities, \[$`N`$-body - FAM\], represent highly informative aspects in our analysis. They are assessed in (Fig. 9, top righthand panel) and Figure 11.
The significance and strength of correlations between the $`N`$-body and FAM velocity components in the scatter diagrams are analyzed by means of a linear regression and correlation analysis. To circumvent excessive pollution of the computed parameters by the virialized motions in high-density regions (see sect. 3.3), the galaxy velocity components in these regression analyses involve $`2h^1\mathrm{Mpc}`$ tophat filtered velocity fields. The resulting numerical values of the correlation parameters are listed in Table 3. Table 3 is organized in two separate sections, one for the regression analysis results of the $`\mathrm{\Lambda }`$CDM mock samples (top section) $`\tau `$CDM samples (bottom section). For both the $`\mathrm{\Lambda }`$CDM and the $`\tau `$CDM section we list the results for four velocity related quantities, each separately for the FAM<sub>30</sub> and the FAM<sub>100</sub> reconstructions.
The presence of significant correlations between FAM reconstructions and their $`N`$-body counterparts is evaluated on the basis of the nonparametric Spearman correlation coefficient $`R_{Spear}`$. The linear regression parameter $`R_{lin}`$ quantifies the linearity of the relation. Prevailing in most situations, the linear regression parameters are used to characterize the relation between reconstructed and real $`N`$-body velocities: the zero-point (offset) $`a_0`$, the slope $`a_{lrg}`$ and the dispersion $`\sigma _{lrg}`$ around the linear regression relation. We assume equal errors in FAM and $`N`$-body velocities, as both are affected by similar shot noise errors (while $`2h^1\mathrm{Mpc}`$ top hat smoothing significantly reduces the impact of virial motions on FAM velocity predictions). In addition, we also list the rms scatter of the parameters, estimated on the basis of the results for the 10 different mock catalogs (for each of the four different configurations).
### 7.3 Inventory
#### 7.3.1 Velocity Amplitude
In the top lefthand frame of Fig. 9 the FAM<sub>30</sub> and FAM<sub>100</sub> velocity amplitudes are compared with their $`N`$-body counterpart $`|v_{\mathrm{Nbody}}|`$. The FAM<sub>30</sub> diagram differs considerably from the FAM<sub>100</sub> diagram: the FAM<sub>30</sub> velocities are systematically smaller than their FAM<sub>100</sub> counterparts. Also, while the latter have a strong one-to-one correlation to the $`N`$-body velocities, the FAM<sub>30</sub> diagram shows a systematic offset with respect to this relation (the solid line) and a somewhat larger scatter.
While the FAM<sub>100</sub> diagram tapers out to higher velocities and even shows a few points with $`|v_{\mathrm{FAM}}|>2000\mathrm{km}\mathrm{s}^1`$, there is a firm ceiling of $`|v_{\mathrm{FAM}}|1300\mathrm{km}\mathrm{s}^1`$ for the FAM<sub>30</sub> velocities. It is a direct reflection of the FAM<sub>30</sub> reconstructions missing out on the gravitational force contributions by the external mass distribution. The asymmetric nature of the scatter in both diagrams is due to particles in high density regions.
To a good approximation, the correlation between the FAM<sub>100</sub> velocity amplitudes and $`|v_{\mathrm{Nbody}}|`$ is that of a linear identity relation: the solid line, $`|v_{\mathrm{Nbody}}|=|v_{\mathrm{FAM}}|`$, forms a good fit to the scatter diagram (see Table 3: $`a_{lrg}1`$). The significantly higher value of Spearman’s correlation coefficient (Table 3: $`R_{Spear}0.68`$ vs. $`R_{Spear}0.54`$ for FAM<sub>30</sub>) indicates and confirms the visual impression of Fig. 8) of the tight correspondence between the FAM<sub>100</sub> and $`N`$-body vector velocity fields. The FAM<sub>30</sub> results stand in marked contrast: the majority of the FAM<sub>30</sub> velocities have a systematically lower amplitude than their $`N`$-body counterparts. It results in a relation with a significantly shallower slope than that of the identity relation $`|v_{\mathrm{Nbody}}|=|v_{\mathrm{FAM}}|`$ (also see table 3): objects with a higher velocity have a larger discrepancy. The contribution by the missing large-scale velocity component $`𝐯_{\mathrm{lss}}`$ to the amplitude of the FAM velocity includes a cross-term ($`𝐯_{\mathrm{Nbody}}𝐯_{\mathrm{lss}}`$), a term dependent on the velocity $`𝐯_{\mathrm{Nbody}}`$ of the galaxy. Most of the missing large-scale velocity component $`𝐯_{\mathrm{lss}}`$ is due to the absence of a bulk flow term in the FAM<sub>30</sub> reconstructions. Subtle and/or higher order external gravitational effects play an additional role: the velocity vector diagrams did already reveal that the presence of shear should be one of the main contributors (cf. eq. 21).
When comparing the $`\mathrm{\Lambda }`$CDM FAM velocity amplitudes with those of the $`\tau `$CDM reconstructions it is evident that in the case of the FAM<sub>30</sub> reconstructions the latter adhere considerably better to the corresponding $`N`$-body values. The linear fitting slope $`a_{lrg}`$ (see table 3) is considerably closer to unity for the $`\tau `$CDM samples than for the $`\mathrm{\Lambda }`$CDM samples. Over a 30$`h^1\mathrm{Mpc}`$ volume the external density inhomogeneities in the $`\mathrm{\Lambda }`$CDM cosmology will induce considerably higher bulk flows than the more moderate $`\tau `$CDM perturbations, which is entirely in line with the theoretical expectation (fig 2). In the case of the FAM<sub>100</sub> reconstructions the qualitative differences are far less prominent. On the scale of $`100h^1\mathrm{Mpc}`$ the mass distribution in both the $`\tau `$CDM and $`\mathrm{\Lambda }`$CDM simulation volumes have converged to homogeneity and no major bulk flows are to be expected.
#### 7.3.2 Velocity Decompositions
In the two righthand frames of Fig. 10 we show the scatter diagrams for the x-component of the FAM<sub>30</sub> (top) and FAM<sub>100</sub> velocities (bottom). Although with a significant level of scatter, the FAM<sub>100</sub> diagram can be fitted quite well by a straight line with a slope close to unity (linear regression line: dashed,unity line: solid). That the equivalent FAM<sub>30</sub> diagram may also be fitted by a straight line, be it with a slope significantly smaller than unity is not entirely straightforward. It stems from an intricate interplay between the small scale velocity field and its larger scale contributions, which in most circumstances are not uncorrelated (Berlind, Narayanan & Weinberg bernarw2000 (2000)).
Table 3 lists the linear regression parameters. Although the average best fitting slopes for the FAM<sub>100</sub> velocities are either larger ($`\tau `$CDM) or smaller ($`\mathrm{\Lambda }`$CDM) than unity, the deviation from unity is considerably smaller than that for the FAM<sub>30</sub> velocities, well within the $`1\sigma `$ uncertainty interval. In all, these regression results do adhere to the expected and noted trend of FAM<sub>100</sub> velocities accounting for practically all contributions to the local velocity field and FAM<sub>30</sub> velocities systematically neglecting significant external contributions. Notice that the scatter around the regression lines for the FAM<sub>30</sub> and FAM<sub>100</sub> reconstructions is of comparable magnitude (as may be inferred from the superposed number density contours).
While the choice of any Cartesian coordinate system is an arbitrary one we have also addressed the decomposition of the particle velocities in one defined by the system itself. The FAM velocities are decomposed in a component projected along the corresponding $`N`$-body velocity, $`v_{proj}`$ (or $`v_{}`$) and the complementary perpendicular component, $`v_{}`$. The second column of figure 10 contains the scatter diagrams for the parallel component of the FAM<sub>30</sub> (top) and FAM<sub>100</sub> (low) velocities. Qualitatively, the behaviour of both diagrams resembles that of the velocity amplitude scatter diagrams in fig. 9. A $`11`$ relation between FAM<sub>100</sub> velocities and $`N`$-body velocity amplitude represents a reasonable fit (solid line: slope $`a_{lrg}0.99`$). The FAM<sub>30</sub> diagram not only appears to deviate strongly from such a $`11`$ relation, it may even fail to fit any linear relation. Also, none of the projected FAM<sub>30</sub> velocity components appears to supersede a value of $`1200\mathrm{km}\mathrm{s}^1`$. Given the fact that the equivalent FAM<sub>100</sub> component even surpasses values of $`2000\mathrm{km}\mathrm{s}^1`$, this confirms the systematic deficiency of the gravitational field in the FAM<sub>30</sub> evaluations.
From the scatter diagrams for the perpendicular FAM velocity components, $`v_{}`$, one can infer that almost all systematic effects are confined to the parallel components $`v_{}`$. For both FAM<sub>30</sub> and FAM<sub>100</sub> the complementary perpendicular component lacks a systematic correlation with the $`N`$-body velocity. It mainly represents unrelated scatter, with a magnitude concentrated around values of $`200250\mathrm{km}\mathrm{s}^1`$. The only difference between the FAM<sub>30</sub> and FAM<sub>100</sub> reconstructions is that for the latter $`v_{}`$ involves considerably higher values, reflecting the higher amplitude of the FAM<sub>100</sub> velocities. FAM<sub>30</sub> velocities, on the other hand, involve stronger misalignments (section 7.3.3).
#### 7.3.3 Velocity Alignments
Misalignments between the reconstructed FAM velocities and the $`N`$-body velocity vectors are the result of a few effects. A major source is that of localized small-scale effects. These are not expected to lead to systematic offsets: they will have a noisy character and reflect random motions in highly nonlinear environments, in particular those of dense virialized regions. Because these have no preferred direction, they behave like randomly oriented “residual” velocities wrt. to the real $`N`$-body velocities of galaxies. Of an entirely different nature are misalignments stemming from the systematic neglect of the external gravitational forces. Because the resulting residual velocity vectors comprise systematic components along one or a few preferred directions, a distinctly anisotropic distribution is the result. This reflect itself as a systematic trend for total $`N`$-body galaxy velocities to be aligned along the residual velocity components.
For both the FAM<sub>30</sub> and FAM<sub>100</sub> reconstructions we computed the angles between the $`N`$-body velocity vector and the FAM velocities. The lower righthand panel in Figure 9 confirms that the alignments of “residuals” and total velocity is indeed considerably stronger for the FAM<sub>30</sub> reconstructions than the FAM<sub>100</sub> ones. The figure plots, for each galaxy in the sample, the misalignment angle $`\theta `$ (or, rather, $`\mu =\mathrm{cos}(\theta )`$) versus the $`N`$-body velocity magnitude $`v_{\mathrm{Nbody}}`$. For the FAM<sub>100</sub> velocities we see a near isotropic distribution of angles. With the exception of a minor concentration near perfect alignment, $`\mathrm{cos}(\theta )=1`$, the distribution is sweeping out nearly uniformly over the full range of $`\mu \mathrm{cos}(\theta )=11`$. If at all there is a trend in velocity amplitude, it appears to be the weak tendency for large velocities to be better aligned.
The above results reflect the observation that FAM<sub>100</sub> residual velocities do mainly consist of small-scale random effects. The FAM<sub>30</sub> residuals form a telling contrast. They are heavily aligned along the full $`N`$-body velocities, with a very strong concentration near $`\theta =0`$. Although occasionally there are serious misalignments, their occurrence diminishes rapidly towards large $`\theta `$. When they occur it almost exclusively concerns small velocities, mostly corresponding to serious misalignments between the locally induced velocity and the added external velocity component.
#### 7.3.4 Velocity Residuals
The residuals accumulate all systematic physical effects as well as random artifacts. They are therefore an excellent source of information on the dynamical role of matter concentrations in the various galaxy sample volumes. If there are large external contributions to the galaxies’ velocity these will constitute a major part of the residuals. On the other hand, if most of those influences are contained within the sample volume treated by FAM, the residuals may mainly reflect localized nonlinearities and artifacts of the FAM method. Scatter diagrams involving the residual velocities will indicate systematic trends and are well suited for elucidating the character and underlying dynamics of external influences.
Figure 11 elaborates on this observation. In two successive rows, the top one for the FAM<sub>30</sub> reconstructions and the bottom one for the FAM<sub>100</sub> reconstructions, it displays the residuals $`𝐯_{res}`$ for each of the three Cartesian velocity components, $`v_x`$, $`v_y`$ and $`v_z`$. Each panel plots the velocity component residual as a function of the corresponding $`N`$-body velocity component.
The mark of a bulk velocity is a constant offset of the scatter diagram, a translation of all FAM velocities by a constant term. This is indeed what is observed in the FAM<sub>30</sub> $`v_y`$ scatter diagram: the vast majority of points is located beneath the $`v_y=0\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ line. It is a telling confirmation of the impression yielded by the corresponding velocity vector fields in Fig. 7. The velocity vector field revealed the presence of strong bulk flow oriented almost perfectly along the $`y`$-axis: clearly visible in the $`N`$-body velocity field, hardly present in the corresponding FAM<sub>30</sub> velocity field, and representing a major component of the residual field \[$`N`$-body - FAM<sub>30</sub>\]. When turning to the equivalent FAM<sub>100</sub> diagram, the indicative offset for a bulk flow has almost completely disappeared. This implies that the source(s) for the bulk flow should be found within the region between $`30h^1\mathrm{Mpc}`$ and $`100h^1\mathrm{Mpc}`$. The equivalent $`v_x`$ and $`v_z`$ FAM<sub>30</sub> residual scatter diagrams do confirm the visual impression of there hardly being a bulk flow contribution along the $`x`$\- and $`y`$-directions.
Additional systematic behaviour is readily apparent in Fig 11: the diagrams show an almost linear increase of residual velocity with $`N`$-body velocity. Also, we find that the $`v_x`$ scatter is skewed towards negative $`v_{res}`$ values while the $`v_z`$ diagram is skewed towards positive $`v_{res}`$ values. In the equivalent FAM<sub>100</sub> scatter diagrams the linear increase of $`v_{res}`$ and the asymmetry in the $`v_x`$ and $`v_z`$ diagram has almost disappeared: the dense core of points has turned into a compact and nearly horizontal bar symmetrically distributed around the $`v_{res}=0\mathrm{km}\mathrm{s}^1`$ line. To a large extent this is explained by the much smaller contribution of external tidal shear to the flows over $`100h^1\mathrm{Mpc}`$ volumes (cf. Fig. 2).
The mark for external shear is a near linear increase of residual velocities as a function of their $`N`$-body (or measured) velocity. Depending on the location $`𝐱`$ of a galaxy within the sample volume and with respect to the shear configuration its participation in a shear flow will involve a velocity component $`v_ss_{ij}x_j`$. This may involve a negative or a positive contribution. With such shear contributions representing a non negligible component to the total velocity, its systematic contribution to a largely random local residual signal reshuffles the velocities such that on average the largest velocity involves the largest residual contribution.
With prominent large-scale bulk and shear motions at large, the FAM<sub>100</sub> residual scatter diagram has largely transformed into a featureless and purely random point distribution. The residuals mainly involve uncorrelated small-scale effects and are nearly independent of the amplitude of the $`N`$-body velocity. Some additional artifacts are seen upon closer inspection: the presence of diffuse “S”-shaped point clouds in both the FAM<sub>30</sub> and FAM<sub>100</sub> residual diagrams, tapering off towards a steep tail at both the negative and positive side of the plots.
The corresponding scatter diagrams for the velocity residual amplitudes $`|𝐯_{\mathrm{res}}|=|𝐯_{\mathrm{FAM}}𝐯_{\mathrm{Nbody}}|`$ represents a summary of the systematic trends (Fig. 10). The FAM<sub>30</sub> velocity residuals show a near linear increase as a function of the $`N`$-body velocity, starting with an offset, indicative of the ingredients of bulk and shear flow in the residuals. The lack of any clear correlation between $`|𝐯_{\mathrm{res}}|`$ and $`|𝐯_{\mathrm{Nbody}}|`$ in the case of the FAM<sub>100</sub> residuals confirms the absence of such systematic components. More clearly than in the case of the individual Cartesian components, the presence of local nonlinear motions may be discerned from the extensive surrounding clouds of outliers.
### 7.4 Power Spectrum Dependence
The contrast between FAM<sub>30</sub> and FAM<sub>100</sub> scatter diagrams is more pronounced in the case of the $`\mathrm{\Lambda }`$CDM mock catalogs than in those assembled for the $`\tau `$CDM universes. This clearly reflects the fact that within the $`\mathrm{\Lambda }`$CDM scenario cosmic structure is characterized by a larger coherence scale. It implies the presence of larger and more coherent structures whose size exceeds $`30h^1\mathrm{Mpc}`$. Their combined gravitation impact will yield a stronger systematic impact in the velocity-velocity comparisons. On the other hand, the dispersions listed in table 3 also show that it would hardly be possible to infer information on the cosmological scenario on the basis of one individual realization. The large dispersion around the average slopes, in particular in the case of the $`\mathrm{\Lambda }`$CDM Universe, show that the magnitude of the external dynamical effects may vary appreciably as a function of the location of the (mock) NBG sample within the simulation box. Local measurements will therefore be unable to separate cosmological effects from those stemming from local variations.
### 7.5 Nonlinearities
The point-to-point diagrams discussed above all contain a substantial level of scatter around the inferred regression relations. With a few exceptions the scatter of velocity quantities is in the order of $`200250\mathrm{km}\mathrm{s}^1`$, for both the $`\mathrm{\Lambda }`$CDM as well as the $`\tau `$CDM FAM reconstructions. The main source for this scatter are the virial motions in the high density and mildly nonlinear environments. Also shot noise provides a substantial additional contribution. In the case of small filter radii, another source of scatter is formed by spurious very close pairs of points in the parent $`N`$-body catalog which for artificial reasons failed to collapse into a single object (Branchini, Eldar & Nusser branchini02 (2002)). Scatter may also be due to higher order multipole components in the external gravity field. An inspection of the particle configurations and the velocity vector maps does unmistakably show significant systematic variations on top of dipolar and quadrupolar components. However, tests restricting the analysis to points in the central regions of the sample produced no substantial decrease in level of scatter. This seems to argue for a minor role of such contributions.
## 8 Bulk Flow and Tidal Shear: <br>Velocity Flow Multipole Components
In the previous sections we have found that in order to obtain a good representation of the local cosmic velocity field it is necessary to take into account the external gravitational influence. This was accomplished through the incorporation of the fully detailed external mass distribution contained in the (flux-limited) galaxy catalogs. This involved the galaxy distribution out to distances of $`100h^1\mathrm{Mpc}`$. The reconstructions showed that modelling of velocity fields by FAM with the inclusion of matter concentrations on such large scales is indeed rewarding.
In nearly all situations where the local volume $`V_{int}`$ is suitably large, the small-scale details of the external mass configuration are rather irrelevant for constructing an appropriate model of the flows in the local Universe. An appropriate approximate expression for the the gravitational potential $`\mathrm{\Phi }_{ext}(𝐫)`$ inside the internal volume $`V_{int}`$ due to the surrounding external matter distribution follows from its expansion in multipole contributions. Assuming a spherical local volume with radius $`R_{int}`$, the potential $`\mathrm{\Phi }_{ext}`$ may be written in terms of a multipole expansion of spherical harmonics $`Y_{lm}(\theta ,\varphi )`$ (see e.g. Jackson jackson (1975))
$`\mathrm{\Phi }_{ext}(𝐫)`$ $`=`$ $`{\displaystyle _{R_{int}}^{\mathrm{}}}{\displaystyle \frac{G\rho (𝐱^{})}{|𝐱𝐱^{}|}}𝑑𝐱^{}`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{m=l}{}}}{\displaystyle \frac{4\pi G}{2l+1}}𝒬_{lm}Y_{lm}(\theta ,\varphi )r^l,`$
in which the multipole moments $`𝒬_{lm}`$ relate to the external density field $`\rho (𝐱^{})`$ as
$$𝒬_{lm}=_r^{\mathrm{}}\rho (𝐱^{})r^{l3}Y_{lm}^{}(\theta ^{},\varphi ^{})𝑑𝐱^{}.$$
(23)
Most contributions to the external gravity $`𝐠_{ext}`$ will be confined to these dipole and quadrupole components, induced by the corresponding large-scale constellations in which the surrounding matter concentrations have grouped themselves. Here we will assess the approximation in which the potential expansion (eqn. 8) is restricted to the monopole term $`l=0`$, the dipole term $`l=1`$ and the quadrupole term $`l=2`$,
$$\mathrm{\Phi }_{ext}(𝐫)\mathrm{\Phi }_0(𝐫)+\mathrm{\Phi }_1(𝐫)+\mathrm{\Phi }_2(𝐫)$$
with
$$\mathrm{\Phi }_l(𝐫)=\underset{m=l}{\overset{m=l}{}}\frac{4\pi G}{2l+1}𝒬_{lm}Y_{lm}(\theta ,\varphi )r^l.$$
(24)
To explore the nature of the external component in the total gravitational field in the Local Universe we proceed by probing it through the resulting peculiar velocity field. The amplitude of higher order terms may be assumed to be so small that one cannot expect to deduce any significant value, given the sizeable errors in the available galaxy peculiar velocity datasets, we may expect this to be a reasonable approximation. We investigate the velocity field by decomposing the residual velocity field – i.e. the component in the velocity field which could not be accounted for in the FAM<sub>30</sub> reconstruction and supposedly induced by external influences – into its multipole components. Once we have determined the bulk flow component and shear tensor components in the tidal velocity field, we will assess whether we can indeed relate this to the external gravitational (“tidal”) influence within the local Universe.
Restricting the description of the external gravitational influence to the first few orders of its multipole expansion has several advantages. The large-scale external dipole and quadrupole gravity perturbations retain a largely linear character, simplifying the velocity field analysis and thus retaining the direct linear relation between gravity and velocity field. Also, by discarding its small-scale fluctuating contributions a physically more transparent image of the velocity field is obtained. This allows a straightforward relation and translation towards the corresponding large-scale pattern of the surrounding mass distribution. A final practical issue of some importance is the fact that the dipole and quadrupole characterization is particularly suited for an implementation in FAM. Restricting the external force field to these moments alleviates the need to take into account a large sample of external galaxies. Not only is the latter computationally expensive, in practice it is even not always feasible.
### 8.1 Velocity Field Multipole Decomposition
In the multipole analysis we restrict ourselves to the externally induced velocity components, $`𝐯_{ext}`$, which in the following we frequently designate by the term “tidal”<sup>1</sup><sup>1</sup>1in the following we regularly use the word “tidal” to shortly indicate the externally induced component of a gravity or velocity field. As it includes a dipolar contribution, strictly speaking this is not an appropriate term. Also see sect. 3.6.2.. For each object, the “tidal” velocity vector is determined by subtracting the internally induced velocity field, $`𝐯_{int}`$, from the object’s full velocity. The latter is usually the $`N`$-body velocity of the mock galaxy, although we will assess the possibility of using the FAM<sub>100</sub> velocity as a reasonable alternative. The internal velocity $`𝐯_{int}`$ is deduced by evaluating, through our FAM computations, the impact of the internal matter distribution within the internal catalog volume $`V_{int}`$. The resulting (residual) peculiar velocity vector $`𝐯_{ext}`$ field may then be expressed in terms of a Taylor series description as function of spatial position $`𝐱`$.
For the practical implementation, we follow the general scheme described by Kaiser (kaiser (1991)). The velocity field Taylor expansion is truncated at the quadratic term and is restricted to the dipole and quadrupole moments (and a minor monopole term). The tidal velocity field $`𝐯_{ext}`$, is then modeled by the the first two components, a bulk flow vector, $`\stackrel{~}{u}_i`$, and a quadratic shear tensor contribution, $`\stackrel{~}{s}_{ij}`$,
$$v_{ext,i}=\stackrel{~}{u}_i+\stackrel{~}{s}_{ij}x\widehat{x}_j,\mathrm{where}i,j=\{1,2,3\},$$
(25)
in which $`i,j`$ denotes the Cartesian component indices. As in eqn. 21, the vectors $`\widehat{x}_j`$ represent the vector components along the Cartesian $`j`$ direction of the spatial unity vector oriented along the object position vector $`𝐱`$. Using these notations, we can easily reconfigure eqn. 25 and express the $`i`$-component of the velocity of object $`n`$ into a product of the vectors $`F_{n,I}`$ and $`V_{Ii}`$,
$`v_{n,i}={\displaystyle \underset{I=1}{\overset{4}{}}}F_{n,I}(𝐱)V_{Ii}`$
in which the data 4-vector $`F_I`$ and the velocity field component 4-vector $`V_{Ii}`$ are defined as
$`F_{n,I}`$ $`=`$ $`\{1,x\widehat{x}_1,x\widehat{x}_2,x\widehat{x}_3\}`$
$`V_{Ii}`$ $`=`$ $`\{\stackrel{~}{u}_i,\stackrel{~}{s}_{i1},\stackrel{~}{s}_{i2},\stackrel{~}{s}_{i3}\}.`$ (26)
The “dipolar” bulk flow components $`\stackrel{~}{u}_i`$ and “quadrupolar” velocity shear components $`\stackrel{~}{s}_{ij}`$ can then be obtained by solving for the vectors $`V_{Ii}`$ on the basis of a fitting analysis (to be precise, $`\stackrel{~}{s}_{ij}`$ also includes a minor residual “monopole” expansion/contraction term). We accomplish this by computing for each Cartesian component $`i`$ the values for the multipole elements $`\stackrel{~}{u}_i`$ and $`\stackrel{~}{s}_{ij}`$ which minimize $`\chi ^2`$
$$\chi ^2=\underset{n=1}{\overset{N_{objs}}{}}\left(v_{n,ext,i}\underset{I}{}F_{n,I}(𝐱)V_{Ii}\right)^2,$$
(27)
to be evaluated on the basis of the data sample of $`N_{obj}`$ objects at locations $`𝐱_n`$ and with inferred “external” velocities $`𝐯_{n,ext}`$. The bulk flow and velocity shear in the externally induced velocity component $`𝐯_{ext}`$, along with a residual expansion term,
$$v_{ext,i}v_{exp,i}+v_{bulk,i}+v_{shear,i}$$
(28)
will follow directly from the inferred values of the 4-vector components $`V_{Ii}`$:
$`v_{exp,i}`$ $`=`$ $`{\displaystyle \frac{1}{3}}Tr(\stackrel{~}{s})x\widehat{x}_i`$
$`v_{bulk,i}`$ $`=`$ $`\stackrel{~}{u}_i;v_{shear,i}={\displaystyle \underset{j=1}{\overset{3}{}}}\stackrel{~}{s}_{}^{}{}_{ij}{}^{}x\widehat{x}_j`$ (29)
in which $`Tr(s)`$ is the trace of the tensor $`s_{ij}`$ and $`\stackrel{~}{s}_{}^{}{}_{ij}{}^{}`$ the traceless shear tensor
$`Tr(\stackrel{~}{s})`$ $``$ $`\stackrel{~}{s}_{11}+\stackrel{~}{s}_{22}+\stackrel{~}{s}_{33}`$
$`\stackrel{~}{s}_{}^{}{}_{ij}{}^{}`$ $`=`$ $`\stackrel{~}{s}_{ij}{\displaystyle \frac{1}{3}}Tr(\stackrel{~}{s})\delta _{ij}.`$ (30)
### 8.2 Velocity Multipole Analysis: results
The results of our analysis are summarized in Table 4. It lists the average quantities for the tidal bulk flow and shear components for the two cosmological scenarios discussed in this work. The table has been organized in four (horizontal) sections. Each corresponds to another “differential” velocity field, the difference between two differently processed velocity fields.
For both the $`\mathrm{\Lambda }`$CDM and the $`\tau `$CDM model each of the quoted values in Table 4 involve the average and standard deviation determined on the basis of ten different realizations. This adds up to 8 configurations, two cosmologies per section. For each of the 8 configurations, in the third column the table lists the dipole component of the external velocity field, the bulk flow $`v_{bulk}`$. Subsequently, the velocity shear is specified in terms of the three eigenvalues $`s_1`$, $`s_2`$ and $`s_3`$ of the traceless shear tensor. This is preceded in the fourth column by the amplitude $`s`$ of the shear. Note that shear is quoted in two units. First, units of $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, followed by the equivalent velocity differential in $`\mathrm{km}\mathrm{s}^1`$ over a volume of $`30h^1\mathrm{Mpc}`$ radius. The intention of the latter is to offer a directly appreciable comparison between the relative importance of bulk flow and shear contributions.
Each of the four sections specifies the values of the computed dipole and quadrupole moments of the velocity field of the corresponding sample. The first section relates to a multipole analysis of the differential velocity field between the full $`N`$-body velocity field and the FAM<sub>30</sub> velocity reconstructions of the inner $`30h^1\mathrm{Mpc}`$ region, $`N`$-body - FAM<sub>30</sub>. The resulting residual velocity field has been generated by the mass distribution beyond a radius of $`30h^1\mathrm{Mpc}`$. On these linear scales the inferred dipole and quadrupole components of the velocity field may be directly related to the moments of the surrounding mass distribution.
The second section of Table 4 does the same for the larger $`100h^1\mathrm{Mpc}`$ region. The outcome of similar analyses are presented in the third and fourth section. The third section repeats the analysis of the first section, except that the external tidal influences are determined on the basis of the difference between the FAM velocity reconstructions within the large $`100h^1\mathrm{Mpc}`$ region and the inner $`30h^1\mathrm{Mpc}`$ region. Earlier, in Section 6.3, we have found that the major share of the origin of the external tidal field is confined to this region and that it therefore may well be determined from the residuals between FAM<sub>100</sub> and FAM<sub>30</sub>. The comparison between the inferred multipole moments of the velocity differences between FAM<sub>100</sub>-FAM<sub>30</sub> in the third section and those in the first section are therefore expected to be rather similar, any systematic differences originating in tidal effects generated beyond a radius of $`100h^1\mathrm{Mpc}`$. The fourth section in Table 4 refers to the values of the residual tidal velocity field between $`N`$-body – FAM<sub>mpl</sub>, the FAM sample after having accounted for the missing external tidal contributions, FAM<sub>30</sub> \+ tidal bulk + tidal shear (see sec. 8.5). If indeed all significant contributions can be characterized by their dipolar and quadrupolar contributions, the multipole values in this section are expected to be negligible.
### 8.3 Velocity Multipole Contributions: Maps
For a direct visual appreciation of the various multipole contributions to the tidal velocity field we assess the “tidal” velocity field $`N`$-body - FAM<sub>30</sub>, the velocity field generated by the mass distribution beyond a radius of $`30h^1\mathrm{Mpc}`$, for one of the $`\mathrm{\Lambda }`$CDM catalogs. The presented maps concern the same $`\mathrm{\Lambda }`$CDM catalog as those presented in the maps of Figures 7. The map of the projection of this “tidal” velocity flow onto three central planes is shown in the top row of Figure 13.
#### 8.3.1 Dipolar component: bulk flow
The externally generated velocity flow is dominated by its bulk flow component. This is in general true for both cosmologies. The large impact of the bulk flow over the local 30$`h^1\mathrm{Mpc}`$ volume can be immediately inferred from the values in the first section of Table 4, revealing contributions in excess of $`200\mathrm{km}\mathrm{s}^1`$.
To facilitate visual appreciation of this observation we have have reoriented the reference system in Figure 13 such that the $`x`$-axis is oriented along the bulk flow. While the original Cartesian system is an arbitrary one and thus lacks a physical context, the “bulk flow reference system” confines the inferred bulk flow $`\stackrel{~}{u}`$ exclusively to the $`x`$-direction. As a result there are no bulk flow components in the corresponding $`y`$\- and $`z`$-direction (note that within the $`yz`$ plane their direction is arbitrarily defined). The pre-eminence of the bulk flow component can be immediately seen in the $`xy`$ and $`xz`$ frames in the top row of Figure 13. Note that the same velocity maps, mostly so the $`yz`$ frame, reveal a clear shear pattern.
#### 8.3.2 Quadrupolar component: velocity shear
Seeking to assess the quadrupolar term in the external velocity field we first remove the remaining expansion term from $`\stackrel{~}{s}_{ij}`$. Diagonalization of the resulting traceless shear tensor $`\stackrel{~}{s}_{}^{}{}_{ij}{}^{}`$ yields the shear eigenvalues and eigenvectors. The eigenvalues $`s_1`$, $`s_2`$ and $`s_3`$ are indicative for the strength of the tidal force field induced by the surrounding matter distribution, while the principal directions of this quadrupolar velocity perturbation field are indicated by the corresponding eigenvectors $`\widehat{e}_{s,i}`$.
The “shear ellipsoid”, the quadratic surface defined by the shear tensor $`\stackrel{~}{s}_{}^{}{}_{ij}{}^{}`$ with principal axes aligned along the eigenvectors and with axis size set by the corresponding eigenvalue $`s_i`$, defines a natural reference system to assess the tidal shear flow field. The coordinate axes of this “shear reference frame” are identified with the orthonormal basis defined by the (normalized) eigenvectors. The $`x`$-axis is chosen to be aligned along the major axis of the “shear ellipsoid”, the direction defined by the largest (positive) eigenvalue $`s_1`$ and directed along the strongest dilational (stretching) motion incited by the external tidal field. Likewise the $`z`$-axis is chosen to coincide with the lowest (negative) eigenvalue $`s_3`$, aligned along the strongest “compressional” component of the tidal velocity flow. This leaves the $`y`$-axis as the one coinciding with the intermediate eigenvalue $`s_2`$.
The imprint of the shearing motions can be discerned within the $`yz`$ plane and, most prominently, along the “x-z” projection of the “bulk flow reference system”. After subtraction of the bulk flow component, i.e. $`N`$-body-FAM<sub>30</sub>-$`v_{bulk}`$, the quadrupolar component of the externally induced velocity flow represents its principal constituent (Fig. 13, lower row). This is confirmed by the values quoted in Table 4 for the shear contribution. In particular when stated in the velocity equivalent unit of $`\mathrm{km}\mathrm{s}^1`$ these shear values suggest that the quadrupolar shear contributions are of a comparable magnitude to those of the bulk flow. The maps in the lower row of Figure 13 suggest that there are strong dilational and compressional motions within the $`yz`$ plane. By contrast, the shear motions in the $`x`$-direction appear to be uncommonly weak. Given the “bulk flow reference system”, it implies that for this particular realization we see a bulk flow directed almost perpendicular to the shear flow motions.
Figure 14 depicts the same $`\mathrm{\Lambda }`$CDM mock sample as presented in Figs. 7, 8 & 13, here in the “shear reference frame”. The top row shows the full externally induced flow field, $`N`$-body - FAM<sub>30</sub>, in this reference system. The tidal shear flows are almost exclusively confined to the $`xz`$ plane. This is most evidently illustrated in the central row of frames showing the velocity field without its bulk flow component: hardly any systematic flow is noticeable in the $`y`$-direction of the intermediate shear eigenvalue.
### 8.4 Multipole Scale Dependence
When turning to the external influences over a large $`100h^1\mathrm{Mpc}`$ region, we may conclude from the second section of the table that most of the external contributions are accounted for, both bulk flow and shear are at least a factor of 3-4 smaller than for the inner $`30h^1\mathrm{Mpc}`$ region. The third and fourth section show that the explicit contributions from the regions between $`100h^1\mathrm{Mpc}`$ and $`30h^1\mathrm{Mpc}`$ and those beyond $`100h^1\mathrm{Mpc}`$ are indeed significantly different, those beyond $`100h^1\mathrm{Mpc}`$ tending towards zero contributions and as far as the shear is concerned almost an order of magnitude smaller than the equivalent contributions by the $`30100h^1\mathrm{Mpc}`$ region.
A similar graphical assessment involving the FAM<sub>100</sub> reconstructions emphasizes the minor significance of tidal contributions stemming from density fluctuations beyond a radius of $`100h^1\mathrm{Mpc}`$. No coherent velocity pattern can be recognized in the residual velocity field between full $`N`$-body and FAM<sub>100</sub> reconstruction. The comparison between this residual velocity field with the velocity maps including the contributions of the inferred bulk flow and shear flow do hardly show any difference. In all cases the velocity fields are dominated by the same thermal motions.
### 8.5 Multipole Velocity Flow Model
Following our argument that the externally induced velocity flow within the inner $`30h^1\mathrm{Mpc}`$ mainly consists of a bulk flow and shear contribution, we may expect that the effect of the external gravity field can be sufficiently accounted for by adding these components to a local velocity field model based on the mass distribution in and around the Local Superclusters.
By separating the “internal” FAM velocity field from the “external” multipole contributions of the (monopole,) dipole and quadrupole components of the “tidal” velocity field and adding the two, we obtain a total “FAM-multipole” model velocity $`v_{fammpl,i}`$,
$$v_{fammpl,i}v_{FAM,i}+v_{exp,i}+v_{bulk,i}+v_{shear,i}$$
(31)
A visual impression of the extent of the successive multipole contributions may be obtained from Fig. 12. The vector plots of the four velocity contributions to $`v_{fammpl}`$ (Eq. 31) are depicted in four successive rows, each within the mutually perpendicular three central slices (wrt. the shear reference system). The top row concerns the FAM<sub>30</sub> velocity field reconstruction, followed successively by the expansion/contraction term (monopole), the bulk flow (dipole) and velocity shear (quadrupolar).
From fig. 13 and fig. 14 we conclude that the differences between the “full” $`N`$-body velocities and $`v_{fammpl}`$, the total sum of the internal FAM<sub>30</sub> and external dipole and quadrupole contributions, do not appear to show systematic trends as it can be noticed from the residual bulk and shear components in section 4 of Table 4. Wherever there are large deviations, these are mainly confined to the high density virialized regions.
### 8.6 Point-to-Point Comparison
A quantitative quality assessment of the “FAM-multipole” model is offered by the point-to-point comparison between the full $`N`$-body velocity and its difference with respect to the successive modes of the “FAM-multipole” velocity in Fig. 15. The $`v_x`$, $`v_y`$ and $`v_z`$ of the various velocity components refer to the “bulk flow reference system”. The top row, plotting $`v_{Nbody}`$ vs. the residual $`N`$-body-FAM<sub>30</sub>, reveals the expected systematic differences due to missing externally induced contributions. Given the fact that the bulk flow in this reference system is confined to the $`x`$-component, we may note the uniform systematic shift of the $`x`$ residuals with respect to the zeropoint $`(v_{Nbody},v_{res})=(0,0)`$ (top lefthand frame). The subsequent addition of the dipolar bulk flow contribution to FAM<sub>30</sub> leads to a systematic downward uniform vertical shift of $`N`$-body-FAM<sub>30</sub>-Bulk (middle row Fig. 15): also the residuals in the $`x`$-direction now center on $`v_{res}=0`$ (note that by virtue of the bulk flow the $`N`$-body velocities in the $`x`$-direction are also skewed to values larger than $`v_{Nbody}=0`$).
The three point-to-point diagrams in the middle row of Fig. 15 show that even while the bulk flow is taken into account systematic motions remain in all three directions. The point-to-point comparisons still follow a strong correlation with respect to the the $`N`$-body velocities. It mainly involves the presence of the quadrupolar velocity shear component (in addition to a minor ingredient contributed by the monopole expansion/contraction term). This can be immediately inferred from the comparison between the diagrams in the central and lower row of Fig. 15: once the quadrupole component “Shear” has been added to the “FAM<sub>30</sub>+Bulk” velocities the systematic effects seem to have largely vanished. What remains in the residuals is mainly random scatter, centered on the $`v_{res}=0\mathrm{km}\mathrm{s}^1`$ line, with some exceptional outliers originating in the virialized regions.
We have quantified the point-to-point comparisons by performing linear regressions similar to those presented in Section 7. Table 5 summarizes the results of this comparison for all catalog samples for both cosmologies. In both cosmological models the slope of the best fitting line is consistent with unity at $`1\sigma `$ confidence level. As expected, the scatter around the fit is similar to that of all previous analyses (see Table 3). Offsets around the zero-point are consistent with zero, although with a large dispersion. The strength of the point-to-point correlations has increased considerably with respect to their FAM<sub>30</sub> counterpart (Table 3) and it is very similar to the FAM<sub>100</sub> case.
### 8.7 Surrounding Matter Distribution: Tidal Source
The surrounding external matter distribution is the source for the tidal velocity field which we inferred in the previous sections. For various purposes we wish to relate the computed dipolar bulk flow and quadrupolar shear flow components to the surrounding matter distribution which induced them.
The induced tidal velocities involve spatial scales ranging from $`30h^1\mathrm{Mpc}`$ to $`100h^1\mathrm{Mpc}`$. Over this range the linear theory of gravitational instability holds to good approximation. This translates into a direct linear relationship between induced velocity $`𝐯_{ext}`$ and the cumulative external gravitational force $`𝐠_{ext}`$,
$$𝐯_{ext}(𝐱,t)=\frac{2f(\mathrm{\Omega },\mathrm{\Lambda })}{3H\mathrm{\Omega }}𝐠_{ext}(𝐱,t),$$
(32)
with $`f(\mathrm{\Omega },\mathrm{\Lambda })`$ the linear velocity growth factor. This linear relationship also holds for every component of the velocity and gravity fields, and thus also for the individual dipolar and quadrupolar components of the externally induced velocity field. They are directly proportional to equivalent dipole and quadrupolar components of the gravity field:
$`𝐯_{bulk}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{2f(\mathrm{\Omega },\mathrm{\Lambda })}{3H\mathrm{\Omega }}}𝐠_{bulk}(𝐱,t)`$
$`s_{ij}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{2f(\mathrm{\Omega },\mathrm{\Lambda })}{3H\mathrm{\Omega }}}𝒯_{ij}(𝐱,t)`$ (33)
with the (external) gravitational tidal shear tensor $`𝒯_{ij}`$ is defined as (see van de Weygaert & Bertschinger rienbert (1996)),
$$𝒯_{ij}=\frac{^2\mathrm{\Phi }_{tidal}}{x_ix_j}.$$
(34)
Notice that because of its external nature, the term $`\frac{1}{3}^2\mathrm{\Phi }_{tidal}\delta _{ij}`$ is always equal to zero.
Ideally, we would like to infer the external tidal potential $`\mathrm{\Phi }_{tidal}`$ directly from the galaxy distribution in a sufficiently large surrounding region. This is specifically true for its dipolar and quadrupolar moments, with the intention to insert these terms directly into the expression for the FAM potential (eqn. 12 and eqn. 10). The required externally induced bulk flow velocity and velocity shear should be the result. The comparison of the FAM computed velocities for the local volume, in combination with the computed tidal velocities (Sec. 8.2), and the observed and measured velocities would then enable us to determine the amount of mass and average density in the local volume.
To determine the gravitational influence of the surrounding matter distribution, we set out to assess the sky distribution of galaxies in the local NBG volume, out to $`r_{NBG}=30h^1\mathrm{Mpc}`$, along with the external mass distribution in radial shells out to a distance $`r_{PSCz}100h^1\mathrm{Mpc}`$. A prominent dipolar matter configuration in the sky distribution will translate into a strong bulk gravity force. Similarly, quadrupolar anisotropies will translate into an effective tidal shear force. In figure 16 we have plotted the galaxies in one of our $`\mathrm{\Lambda }`$CDM mock catalogs in five successive distance shells. Aitoff projections of the angular positions of the galaxies, as seen from the centre of the local NBG volume, provide an impression of the level of anisotropy in the mass distribution at successive radii.
The first sky plot (top sphere) depicts the sky position of the galaxies in the local NBG-mimicking mock sample. It involves a highly flattened distribution, perhaps reminiscent of the Supergalactic Plane. The four subsequent shells correspond to successive cuts through the PSCz mimicking samples, at sampling depths $`d_{sur}=[030],[3055],[5570]`$ and $`[85100]h^1\mathrm{Mpc}`$. The first and direct observation is the diminishing sample density as a function of survey depth, in accordance with the selection function (eqn. 19). Structure is most prominent in the first shell, at $`d_{sur}=[3055]h^1\mathrm{Mpc}`$ (central left sphere). The structure contained in this shell also shows a clear affiliation with the matter distribution in the local NBG volume. The compact massive concentration at $`l220^{}`$ is clearly connected to a dense region in the local “plane”. A superficial inspection of the angular galaxy distribution reveals the presence of strong dipolar and quadrupolar components, effecting considerable tidal forces. Note that both external shells display a rather strong concentration of galaxies in their southern hemisphere, in the vicinity of $`l180200^{}`$. Similar but weaker contributions can also be recognized from the galaxy distribution in the shell between $`d_{sur}7085h^1\mathrm{Mpc}`$. Beyond $`d_{sur}>85h^1\mathrm{Mpc}`$, however, the angular pattern appear to be considerably less pronounced. This is in line with the earlier findings that there were hardly noticeable tidal contributions from large distances.
To see to what extent the depicted galaxy distribution can indeed be held responsible for most of the inferred tidal bulk flow and tidal shear, we have determined the corresponding bulk force $`𝐠_{bulk}`$ (eqn. 4) and tidal shear $`𝒯_{ij}`$ (eqn. 5) evoked by the external galaxy distribution ($`r>30h^1\mathrm{Mpc}`$). Since we do not have a continuous density field but the positions of a finite number of objects in our galaxy flux-limited and full mass distribution catalogs, the bulk acceleration on the LG is computed from the discrete equivalent. For a sample of galaxies at locations $`𝐱_i`$, with an average number density $`n`$ of selected objects, this leads to
$$𝐠_{bulk}=\frac{Hf(\mathrm{\Omega },\mathrm{\Lambda })}{4\pi n}\underset{i}{}\frac{1}{\psi (x_k)}\frac{𝐱_k}{|𝐱_k|^3}.$$
(35)
where $`\psi (x_k)`$ is the sample selection function at distance $`x_k`$, whose inverse functions as weighting factor. For practical reasons, comparison with the inferred bulk flow $`𝐯_{bulk}`$, we have translated the bulk acceleration into equivalent velocity units by means of the transformation $`Hf(\mathrm{\Omega },\mathrm{\Lambda })/\frac{3}{2}\mathrm{\Omega }H^2`$. The equivalent “discrete” expression for the external tidal shear is
$$𝒯_{ij}=\frac{Hf(\mathrm{\Omega },\mathrm{\Lambda })}{4\pi n}\underset{i}{}\frac{1}{\psi (x_k)}\frac{3x_{ki}x_{kj}}{|𝐱_k|^5}.$$
(36)
For the $`\mathrm{\Lambda }`$CDM mock galaxy sample depicted in Fig 16 we determine the gravity dipole by computing for a set of spherical external shells the resulting bulk flow acceleration (eq. 35) and the gravity quadrupole by computing external tidal shear (eq. 36). Recently, a similar approach was followed by Teodoro et al. teod2004 (2004). The spherical shell volumes are defined by an inner radius $`r_{inn}=30h^1\mathrm{Mpc}`$ and an outer radius $`r_{out}`$. The width of the shell is gradually enlarged by increasing $`r_{out}`$ from $`r_{out}=30h^1\mathrm{Mpc}`$ to $`r_{out}=100h^1\mathrm{Mpc}`$. The convergence of the resulting gravity dipole direction on the sky can be observed in Fig. 17. The small red diamonds are consistently located near $`l230^{}`$, and converge at a sky location close to the direction of the velocity dipole (large red diamond). To get an idea of the amplitudes involved, Fig. 19 (top panel) shows the development of the cumulative gravity dipole as a function of external distance $`d_{sur}`$ ($`30h^1\mathrm{Mpc}<d_{sur}<100h^1\mathrm{Mpc}`$). By means of symbols the corresponding velocity dipole, for each of the three directions $`x`$, $`y`$ and $`z`$, are inserted at the outer radius of $`d_{sur}100h^1\mathrm{Mpc}`$. Note that we have restricted ourselves to a case study example. A more extensive and proper assessment, including a proper error estimate of both gravity dipole and quadrupole as well as the bulk and shear flow, is beyond the scope of the present argument. This issue, involving the shot noise influence on gravity dipole and quadrupole and the role of FAM uncertainties on the velocity flow components, will be treated in more detail in a forthcoming study.
We see that in the $`x`$-, $`y`$ and $`z`$-directions of the gravity and velocity dipoles are in reasonable agreement, within a margin of $`30\mathrm{km}\mathrm{s}^1`$. This observation justifies our expectation that the dipole can be estimated to sufficient accuracy from the surrounding external galaxy distribution. The dipole may then be estimated from the surrounding external galaxy distribution, so that the latter can be invoked to correct for the influence of the external tidal field in the dynamics of the local volume.
The situation is comparable for the cumulative tidal shear, in terms of its three eigenvalues and eigenvectors. Also the gravity quadrupole appears to converge relatively smoothly towards the velocity shear. This may be inferred from the plotted directions of the eigenvectors $`\widehat{𝐞}_{𝒯1}`$, $`\widehat{𝐞}_{𝒯2}`$ and $`\widehat{𝐞}_{𝒯3}`$ of the tidal shear $`𝒯_{ij}`$. They are indicated by means of three symbols, the triangle corresponding to the stretching component $`𝒯_1`$, the star the middle component $`𝒯_2`$ and the square the compressional component $`𝒯_3`$. The tidal shear tensor wanders extensively across the “sky” as we push the outer radius of the external shell outward, as is shown by the paths of the corresponding eigenvectors. Interestingly, once the shell radius starts to approach $`100h^1\mathrm{Mpc}`$, each of the eigenvectors appear to converge near the location of the corresponding stretching, central and compressing velocity shear tensor eigenvectors. However, also here we notice significant deviations in individual cases.
For comparison, we can appreciate the role of the external tidal field on local dynamics for the case of the $`\tau `$CDM cosmology. Figure 18 combines the galaxy sky distribution for a $`\tau `$CDM mock galaxy sample, in the same radial shells as in Fig. 16. The final frame shows the Aitoff projection of the gravity dipole and gravity quadrupole eigenvectors for a set of gradually increasing radial shells. From the galaxy sky distribution in the four slices at sampling depths $`d_{sur}=[030],[3055],[5570]`$ and $`[85100]h^1\mathrm{Mpc}`$ we notice that these involve considerably more isotropic distributions. Hardly any prominent patterns can be discerned in the sky distribution. This is expressed in a more erratic wandering of gravity dipole and quadrupole directions (lower frame, fig. 16). This also implies a more substantial contribution of shot noise effects. The latter represent a major source for deviations between the velocity dipole and shear flow quadrupole and the gravity dipole and quadrupole. The smaller coherence length of the $`\tau `$CDM fluctuations and the more randomly oriented contributions by the individual external matter concentrations may therefore be directly related to the lower level of coincidence between velocity and gravity directions than in the case of the more prominent anisotropies in the $`\mathrm{\Lambda }`$CDM cosmology. In this, we have to realize that the amplitude of dipole and quadrupolar contributions between the two scenarios are not too different (cf. table 5). The less prominent anisotropies in the $`\tau `$CDM catalogues are therefore compensated by a higher average matter density.
### 8.8 Multipole Components: Summary
The above results reassure the fact that the external tidal field can be well characterized by its main multipole components, the bulk flow and velocity shear. In terms of multipole amplitude convergence, these results show a better agreement for the $`\mathrm{\Lambda }`$CDM model than for the $`\tau `$CDM one. This is due to the intrinsic characteristics of both cosmic models. As has been discussed in Section 4, and may be directly appreciated from Fig. 2, the relatively lower amplitude of the $`\tau `$CDM perturbations is compensated by a higher mass content. It leads to an equally strong external gravitational influence. On the other hand, the smaller spatial coherence of density features in the $`\tau `$CDM scenario causes the orientation of the gravity dipole and quadrupoles to be rather jittery. The direction of the cumulative gravitational force in the $`\tau `$CDM scenario wanders erratically over the sky as we move further out from the local volume. This differs from the situation in the $`\mathrm{\Lambda }`$CDM samples, where we observe a consistent, systematic and coherent convergence towards the final dipole direction.
The above results confirm the fact that the external tidal velocity field can be well characterized by its main multipole components, the bulk flow and velocity shear. This depends to some extent on the cosmology. In terms of multipole amplitude convergence, these results show a better agreement for the $`\mathrm{\Lambda }`$CDM than for the $`\tau `$CDM model.
## 9 Conclusions
In this work we have applied the FAM technique to construct model velocity fields using mock catalogs resembling the NBG and IRAS-PSC$`z`$ galaxy catalogs. The mock catalogs were extracted from $`N`$-body simulations in which the central observer mimics some of the properties of the Local Group environment. Comparing FAM velocities obtained from the NBG mock catalogs with those obtained from the larger PSC$`z`$ mock catalogs and, finally, to the $`N`$-body velocities, allowed us to quantify the importance of the gravity field generated by the mass distribution within and beyond the LS.
Neglecting the mass distribution outside the LS leads to a systematic underestimate of the gravity field. The amplitude of this bias depends on the amount of power on scales larger than the LS, and thus on the cosmological models. In a $`\tau `$CDM universe model peculiar velocities are $`20\%`$ smaller than the true ones. In the case of a $`\mathrm{\Lambda }`$CDM model, which has more power on large scales, model velocities are underestimated by $`35\%`$.
The results of the described FAM analyses are encouraging in the sense that the presently available all-sky, flux limited catalogs such as PSC$`z`$ appear to be capable of accounting for the major share of the velocity field on the scale of the Local Supercluster. While the $`30h^1\mathrm{Mpc}`$ restricted NBG sample showed a substantial deficiency in its capacity to generate the local cosmic motions, in particular in the case of the $`\mathrm{\Lambda }`$CDM Universe models, in both cases the $`100h^1\mathrm{Mpc}`$ mock samples appear to embody nearly all matter concentrations responsible for the generated velocities in our local (NBG catalog) neighbourhood.
Also we notice a telling difference between the performance of both FAM<sub>30</sub> and FAM<sub>100</sub> reconstructions for the case of the $`\mathrm{\Lambda }`$CDM cosmology catalogs on the one hand and the $`\tau `$CDM model catalogs on the other hand. The fact that the $`\mathrm{\Lambda }`$CDM model involves substantially more power on large scales, $`r>30h^1\mathrm{Mpc}`$, than the $`\tau `$CDM model is reflected in the better quality of the FAM<sub>30</sub> reconstructions for the $`\tau `$CDM catalogs. The presence of substantial mass inhomogeneities with a scale in excess of that of the local Universe regions implies a larger external contribution to the local velocity field. This is also borne out by the fact that for the $`\mathrm{\Lambda }`$CDM catalogs we see a considerable improvement in velocity field reconstruction quality going from the FAM<sub>30</sub> to the FAM<sub>100</sub> reconstructions (see Table 4), while this is far less so for the $`\tau `$CDM catalogs.
Of course, whether the resulting models do indeed form an unbiased representation of the actual velocity field will to some extent also depend on whether the galaxy distribution in the flux limited galaxy catalogs does represent an unbiased reflection of the actual (external) mass distribution surrounding the Local Supercluster resembling region. The results of recent studies (Verde et al. verde2002 (2002), Lahav et al. lahav2002 (2002), Tegmark, Zaldarriaga & Hamilton tegzalham2001 (2001) and Branchini, Dekel & Sigad brandeksig2002 (2002)) are quite encouraging in this respect. They seem to indicate, certainly on scales larger than $`5h^1\mathrm{Mpc}`$, that both IRAS and 2dF galaxies trace the underlying mass distribution in an unbiased fashion.
Nonetheless, observations along the lines of the presented mock catalogs seem to suggest that a proper analysis of Local Universe dynamics based on a combination of information of local small-scale (peculiar) galaxy velocities and a rough yet well-founded idea of the matter distribution on scales of a few hundred $`100h^1\mathrm{Mpc}`$ may help us towards acquiring far more insight into the dynamical history of the emergence and assembly of the striking nonlinear patterns we have discovered in the large scale matter distribution. Moreover, we have uncovered evidence that a meticulous point-to-point analysis of such velocity samples may help towards modelling the total local force field, including a proper model for the external forces.
When modeling the peculiar velocity of a LS look-alike region by only considering the matter distribution within $`30h^1\mathrm{Mpc}`$, the end product is a biased velocity field lacking of any large scale signature. This bias can be eliminated by accounting for the mass distribution beyond the LS. Our experiments demonstrate that sampling the mass distribution out to scales of $`100h^1\mathrm{Mpc}`$, in a flux limited fashion, is sufficient to account for the large scale contribution to the peculiar velocities in our cosmological neighborhood. More precisely, we have found that the cosmic velocity field within the LS, modeled by FAM using the mass distribution traced by PSC$`z`$ galaxies out to $`100h^1\mathrm{Mpc}`$ is unbiased. The differences between true and FAM velocity field are random and mainly occur in high density environments which are dominated by virial motions that are not modeled correctly by FAM.
The gravity and velocity fields generated by the mass distribution beyond the scale of the Local Supercluster are well characterized by their bulk flow and shear components. Therefore, one can obtain an unbiased model velocity field by superimposing a local model velocity field within the Local Supercluster to the bulk flow and shear components of the velocity field generated by the mass distribution between 30 and $`100h^1\mathrm{Mpc}`$.
These considerations suggest that velocity models which only consider the dynamics within the Local Superclusters might have been affected by systematic errors. In particular, our work suggests that, when compared with observed velocities, they might have underestimated the value of the density parameter, $`\mathrm{\Omega }_m`$, by 15-25 %. However, the analysis of Shaya, Tully and Peebles (shaya (1995)), based on the galaxy distribution in the Local Supercluster, shows that a lower, not a larger, value of $`\mathrm{\Omega }_m`$ is found when complementing the local mass distribution with the large scale one traced by rich Abell clusters. A more precise evaluation of this bias will be performed in a future work in which we will perform the same analysis presented here using a new set of mock catalogs that are constrain to reproduce the distribution of the mass in our local Universe (see e.g. van de Weygaert & Hoffman vdwhof (1999), Mathis et al. mathis2002 (2002), Klypin et al. klypin (2003)).
Furthermore, our analysis shows that all model velocity fields of the Local Supercluster which are based on the PSC$`z`$ catalog (e.g. Branchini et al. branchini99 (1999), Schmoldt et al. schmoldt99 (1999), Valentine et al. valentine (2000), Sharpe et al. sharpe (2001)) are free from systematic biases arising from having neglected the large scale contribution from scales beyond its realm. Moreover, since the IRAS PSC$`z`$ survey is considerably deeper than $`100h^1\mathrm{Mpc}`$, it is reasonable to assume that the PSC$`z`$ catalog can be used to predict unbiased velocities well beyond our Local Supercluster that, if compared with observed galaxy peculiar velocities, can discriminate among different cosmologies characterized by different values of $`\mathrm{\Omega }_m`$, like the $`\mathrm{\Lambda }`$CDM and $`\tau `$CDM models.
The plausibility of this hypothesis has been recently confirmed by the analysis of Hoffman et al. (hoffman01 (2001)) that shows that the bulk and shear components of the external velocity field in the local universe inferred from the peculiar velocities in the Mark III catalog (Willick et al. 1997a , 1997b ) are qualitatively consistent with those expected from the mass distribution traced by IRAS PSC$`z`$ galaxies. On the other hand, the claim on the basis of the SMAC cluster peculiar velocity sample (Hudson et al. hudsv (2004)) of an extra $`225\mathrm{km}\mathrm{s}^1`$ bulk flow component generated by matter concentrations on a scale exceeding $`100h^1\mathrm{Mpc}`$ should issue some caution with respect to claims of having accounted for all external influences on the local cosmic flow.
Coupling the local velocity model provided by FAM to the large scale contribution provided by linear theory allows to obtain a model velocity field which is unbiased, nonlinear and fast to compute. This means that, for the first time, we are in the position of performing a large number of experiments aimed at studying the nonlinear evolution of cosmic structures, such as filaments and clusters, and explore the role of tidal fields during their gravitational collapse. This relates to the observation that filaments are forming as a consequence of anisotropic collapse, induced a compressional tidal force acting perpendicular to the “axis” of the filament. By tracing out the coherent paths of the compressional modes of the primordial tidal field one can identify the sites of the later nonlinear filaments (Bond et al. cosmweb (1996), see Van de Weygaert weyfoam (2002)). In turn this is directly related to cluster locations: the strong primordial tidal shear is the result of a local quadrupolar mass distribution. The corresponding overdensities tend to evolve into rich clusters, explaining the intimate link of clusters and filaments in the cosmic web.
Finally, it is worth stressing that in this work we have neglected the fact that we measure galaxy redshifts rather then positions. By means of an elegant formalism, Phelps phelps2000 (2000) demonstrated the feasibility of working out the action principle in redshift space. With respect to FAM, Nusser & Branchini (nusbran (2000)) have shown that it can be easily implemented in redshift space and Branchini, Eldar & Nusser (branchini02 (2002)) demonstrated that it performs equally well in real and redshift space. Therefore, our unbiased, nonlinear model velocity field also allows to perform an accurate correction for redshift space distortions and thus lead to a precise reconstruction of the mass distribution in real space. Mapping the mass in the local universe down to nonlinear scales and comparing it with the distribution of baryonic mass (in form of stars or diffuse, ionized gas) is of considerable astrophysical interests as it will constrain and help understanding the process of galaxy formation and evolution within the Universe
###### Acknowledgements.
The authors thank Shaun Cole for allowing the use of his $`N`$-body simulations. E.R.D. thanks W.E. Schaap for stimulating discussions. E.R.D. thanks the Universita’ di Roma tre for its hospitality while part of this work was done. E.R.D. has been supported by The National Council for Research and Technology (CONACyT, México) through a scholarship. San Crispino provided unique and inspiring guidance.
|
warning/0506/nlin0506001.html
|
ar5iv
|
text
|
# Discrete Klein-Gordon models with static kinks free of the Peierls-Nabarro potential
## I Introduction
Discrete solitons and more specifically kink-like topological excitations are ubiquitous structures that arise in numerous physical applications ranging from dislocations or ferroelectric domain walls in solids, to bubbles in DNA, or magnetic chains and Josephson junctions, among others (see, e.g., Kivshar for a recent exposition of relevant applications). The mobility of such lattice kinks is one of the key issues in many of these applications, especially since the pioneering works of peyrard ; yip which illustrated that the kinematics on the lattice is dramatically different from the continuum analog of such equations where constant speed propagation is typical. Instead, on the discrete substrate, kinks need to overcome the, so-called, Peierls-Nabarro potential (PNp), constantly radiating their energy and being eventually trapped by the lattice. The static PNp refers to the energy difference between a stable inter-site centered discrete kink and an unstable, onsite centered discrete kink. Clearly, as a kink is travelling from one site to the next, it “wobbles” over this potential energy landscape boesch . However, even though clearly travelling is intimately connected with overcoming the static PNp without “radiating” energy KW , this connection is relatively subtle and the inter-dependence of these two features (static PNp and travelling) still remains elusive flach . Typically, discrete kinks traveling with finite velocity have only been obtained for a discrete set of velocities yzolo which makes the motion unstable with respect to perturbations. There exists a class of more exotic exact solutions (the so-called “nanoptera”) where the kink propagates together with a plane wave having the same velocity yzolo .
While the travelling problem is extremely interesting in its own right, in the present work, we will start by examining the construction of discrete models with PNp-free kinks, using a simplified (quasi)static approach. Two classes of discrete models where static kink can be placed anywhere with respect to the lattice have been previously derived: one conserving energy SpeightKleinGordon and another one conserving momentum PhysicaD . In both cases the static kink solution can be obtained from a two-body nonlinear map. In the present paper we demonstrate that, in general, a discrete version of the first integral of the static continuum Klein-Gordon field plays the role of this nonlinear map. Thus we derive a wide class of such models including the two above-mentioned classes as special cases. The advantage of this approach is that the kinks are no longer (typically) trapped by the lattice. Instead they can be accelerated by even weak external fields. However, a note of caution should be added here. While one might naively expect that such solutions would be intimately connected with slow travelling, it has been demonstrated numerically that travelling solutions (when they can be found as e.g. in yzolo ; karpan for Klein-Gordon lattices, using the methods of flesh ) have a sharp lower bound in their wave speed aigner . The existence of such a threshold illustrates the fact that one should be particularly careful in trying to infer features of the travelling problem from such “static” considerations. On the other hand, as the recent work of Barashenkov, Oxtoby and Pelinovsky demonstrates dima , discretizations without PNp are much more natural candidates for possessing travelling solutions for a isolated wave speeds (not close to zero).
The presentation of our results will be structured as follows. Section II will contain the setup and notations used for the Klein-Gordon models. Section III will present the general methodology for obtaining static PNp-free discretizations. Section IV will illustrate the connection to previously reported models. Section V will focus on the special case example of the $`\varphi ^4`$ model, for which our numerical observations will be presented in section VI. Finally in section VII, we will summarize our findings and present our conclusions.
## II Setup
We consider the Lagrangian of the Klein-Gordon field,
$$L=_{\mathrm{}}^{\mathrm{}}\left[\frac{1}{2}\varphi _t^2\frac{1}{2}\varphi _x^2V(\varphi )\right]𝑑x,$$
(1)
and the corresponding equation of motion,
$$\varphi _{tt}=\varphi _{xx}V^{}(\varphi )D(x).$$
(2)
Topological solitons (kinks) are possible only if $`V(\varphi )`$ has at least two minima $`\varphi _{01}`$ and $`\varphi _{02}`$, where $`V^{}(\varphi _{0i})=0`$ and $`V^{\prime \prime }(\varphi _{0i})>0`$. Obviously, $`\varphi =\varphi _{01}`$ and $`\varphi =\varphi _{02}`$ are the stationary solutions to Eq. (2). We will study the properties of kinks that interpolate between these two stationary solutions. Our considerations allow one to treat the cases when other minima appear in between the two minima, $`\varphi _{01}`$ and $`\varphi _{02}`$, connected by the kink.
Equation (2) will be discretized on the lattice $`x=nh`$, where $`n=0,\pm 1,\pm 2\mathrm{}`$, and $`h`$ is the lattice spacing.
For brevity, when possible, we will use the notations
$$\varphi _{n1}l,\varphi _nm,\varphi _{n+1}r.$$
(3)
We would like to construct a nearest-neighbor discrete analog to Eq. (2) of the form
$$\ddot{m}=D(C,l,m,r),$$
(4)
where $`C>0`$ is a parameter related to the lattice spacing $`h`$ as $`C=1/h^2`$, such that in the continuum limit $`(C\mathrm{})`$, $`D(C,l,m,r)D(x)=\varphi _{xx}V^{}(\varphi )`$.
Note that in this context, the “standard” discretization emerges in the form: $`D(C,l,m,r)=C(l2m+r)V^{}(m)`$. Generalizations of this model will be discussed in the form
$$\ddot{m}=C(l2m+r)B(l,m,r),$$
(5)
where $`B(l,m,r)`$ has $`V^{}(\varphi )`$ as the continuum limit.
We will characterize a model as PNp-free if a static kink can be placed anywhere with respect to the lattice (continuum, rather than discrete, set of equilibrium solutions). This is equivalent to demanding that the kink must have an neutral direction, or (from Noether’s theory arnold ) a Goldstone translational mode. It is natural to categorize this definition of PNp-free model as “static” or “quasi-static”, in the sense that it does not involve the kinematic or dynamical properties of the model.
On the other hand, one can demand the absence of PNp at finite kink velocities. This can be transformed to the demand that the discrete model supports the exact traveling wave solutions and this demand can be called “dynamic” definition; see e.g. flach for such travelling wave examples, where the “static” definition of the PNp clearly fails.
In this paper we aim to construct the models PNp-free in the static sense as a first (yet nontrivial) step towards understanding the nature of the discrete travelling problem (see also the comments above).
We will also focus on the existence of physically motivated conserved quantities for the derived models. Hamiltonian models are energy-conserving models and the models with $`dM/dt=0`$, where
$`M={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\dot{\varphi }_n\left(\varphi _{n+1}\varphi _{n1}\right),`$ (6)
will be called momentum-conserving models. As was shown in PhysicaD , the discrete model of Eq. (4) conserves the momentum of Eq. (6), if it can be presented in the form
$`\ddot{m}={\displaystyle \frac{(m,r)(l,m)}{rl}}.`$ (7)
This can be verified by calculating
$`{\displaystyle \frac{dM}{dt}}={\displaystyle \underset{n}{}}\ddot{\varphi }_n(\varphi _{n+1}\varphi _{n1})`$
$`={\displaystyle \underset{n}{}}[(\varphi _n,\varphi _{n+1})(\varphi _{n1},\varphi _n)]=0,`$ (8)
where we have used the fact that the terms $`\dot{\varphi }_n(\dot{\varphi }_{n+1}\dot{\varphi }_{n1})`$ cancel out due to telescopic summation.
## III Static PNp-free discretization
Our aim here will be to discretize Eq. (2) in a symmetric way, so that the static kink solution can be found from a reduced first-order difference equation. According to SpeightKleinGordon , if we achieve that, then we are going to have a one-parameter family of solutions with the possibility to place equilibrium kinks anywhere with respect to the lattice (and hence, PNp-free in the static sense).
The first integral of the steady state problem in Eq. (2), $`\varphi _x\sqrt{2V(\varphi )}=0`$ (with zero integration constant), can be written in the form
$$w(x)g(\varphi _x)g\left(\sqrt{2V(\varphi )}\right)=0,$$
(9)
where $`g`$ is a continuous function.
Our plan will then be the following:
* discretize the first-order differential equation of Eq. (9) using a first order difference scheme $`w(l,m)=0`$.
* Then express the right-hand side of Eq. (2) as a sum of terms containing derivatives, e.g., $`dw/dx`$, $`dw/d\varphi `$, etc.
* As a result, discretizations of such terms, e.g., $`dw/dx\sqrt{C}[w(m,r)w(l,m)]`$, vanish for $`w(l,m)=0`$ (or otherwise stated: the construction of the equilibrium solution is converted to a first order difference problem). Then, the static kink (PNp-free, by construction) solutions for the obtained discrete model can be found from this two-site problem.
In the following, we will consider a particular case of Eq. (9) with $`g(\xi )=\xi ^2`$, for which we introduce the notation
$`u(x)\varphi _x^22V(\varphi )=0,`$ (10)
and the following two-site discrete analog
$`u(l,m)C(ml)^22V(l,m)=0.`$ (11)
We will also use the shorthand notations,
$$u_l=u(l,m),u_m=u(m,r).$$
(12)
We have assumed that the Klein-Gordon field supports kink solutions. Then, at least for the case of weak discreteness, Eq. (11) also supports static kinks because it is nothing but a discretization of the first integral of static version of Eq. (2) (see also SpeightKleinGordon ).
The next step is then to find a discretization of the right-hand side of Eq. (2), $`D(x)`$, which vanishes when Eq. (11) is fulfilled.
One simple possibility comes from the following finite difference
$`D_1(l,m,r){\displaystyle \frac{u_mu_l}{rl}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{du}{d\varphi }}=D(x).`$ (13)
One can also consider, more generally, continuous functions $`q(u_l,h)`$ such that $`q(0,h)=0`$ and, in the continuum limit, $`q(u,0)=u`$ and $`\frac{dq}{du}(u,0)=1`$. For example, one can take $`q=(e^{hu}1)/h`$ or $`q=u+_{n>1}A_nh^{n1}u^n`$ with constant $`A_n`$, etc. Then,
$$\frac{1}{2}\frac{dq}{d\varphi }\left(\frac{dq}{du}\right)^1=D(x).$$
(14)
Discretizing the left-hand side of Eq. (14) we obtain
$$D_2=\frac{1}{2}\frac{q(u_m,h)q(u_l,h)}{rl}\left[\frac{1}{q^{}(u_l)}+\frac{1}{q^{}(u_m)}\right].$$
(15)
Inspired by SpeightKleinGordon , we note that, in the continuum limit,
$`{\displaystyle \frac{v(m,r)}{rm}}{\displaystyle \frac{v(l,m)}{ml}}{\displaystyle \frac{dv}{d\varphi }}v{\displaystyle \frac{\varphi _{xx}}{\varphi _x^2}},`$ (16)
and find
$`D_3{\displaystyle \frac{u_m}{rm}}{\displaystyle \frac{u_l}{ml}}+\sqrt{2V(l,m,r)}\times `$
$`\left({\displaystyle \frac{\sqrt{C(rm)^2u_m}}{rm}}{\displaystyle \frac{\sqrt{C(ml)^2u_l}}{ml}}\right)`$
$`{\displaystyle \frac{du}{d\varphi }}u{\displaystyle \frac{\varphi _{xx}}{\varphi _x^2}}+\sqrt{2V}\left({\displaystyle \frac{d\sqrt{2V}}{d\varphi }}\sqrt{2V}{\displaystyle \frac{\varphi _{xx}}{\varphi _x^2}}\right)`$
$`=D(x).`$ (17)
Since the expressions for $`D_i(l,m,r)`$ given by Eqs. (13),(15) and (17) tend to $`D(x)`$ in the continuum limit, one can write the following discrete analog to the Klein-Gordon equation Eq. (2)
$$\ddot{m}=\underset{i}{}b_iD_i(l,m,r),\mathrm{where}\underset{i}{}b_i=1.$$
(18)
Then, by construction, any structure derived from the two-site problem of Eq. (11) is a static solution of Eq. (18) and hence, the latter is the static PNp-free discrete model.
The model of Eq. (18) can be generalized in a number of ways. For example, function $`D_3`$, Eq. (17), can be modified choosing different functions $`V(l,m,r)`$ to discretize $`V(\varphi )`$. Then, the modified $`\stackrel{~}{D}_3`$ can be added to the linear combination in the right-hand side of Eq. (18).
The model of Eq. (18) can be also generalized by appending terms which disappear in the continuum limit and ones that vanish upon substituting $`u_l=0`$ and $`u_m=0`$. For example, the derivative $`df(u)/d\varphi `$ can be discretized as $`2[f(u_m)f(u_l)]/(rl)`$ or as $`2f^{}(u_l/2+u_m/2)(u_mu_l)/(rl)`$. If then we have difference of such terms in the equation of motion, then in the continuum limit they will cancel out.
Any term in the right-hand side of Eq. (18) can be further modified by multiplying by a continuous function $`e(C,l,m,r)`$, whose continuum limit is unity (see e.g. Saxena for such an example, also discussed in more detail below).
Generally speaking, the discrete PNp-free Klein-Gordon models derived here do not conserve either an energy, or a momentum-like quantity. However, as it will be demonstrated below, they contain energy-conserving and momentum-conserving subclasses.
## IV Connection with Previously Reported Models
One energy-conserving PNp-free Klein-Gordon model has been derived by Speight with co-workers SpeightKleinGordon with the use of the Bogomol’nyi argument Bogom . Their model, can be written in the form of Eq. (5), with the Lagrangian
$`L={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\dot{\varphi }_n^2{\displaystyle \frac{C}{2}}{\displaystyle \underset{n}{}}\left(\varphi _n\varphi _{n1}\right)^2`$
$`{\displaystyle \underset{n}{}}\left({\displaystyle \frac{G(\varphi _n)G(\varphi _{n1})}{\varphi _n\varphi _{n1}}}\right)^2,`$
$`\mathrm{where}G^{}(\varphi )=\sqrt{V(\varphi )}.`$ (19)
The static kink solution can then be derived from the lattice Bogomol’nyi equation SpeightKleinGordon , which can be taken in the form
$`U(l,m)=C(ml)^22\left({\displaystyle \frac{G(m)G(l)}{ml}}\right)^2=0,`$ (20)
which is a particular case of Eq. (11). The equation of motion derived from Eq. (19), written in terms of Eq. (20), is
$`\ddot{m}={\displaystyle \frac{U_m}{rm}}{\displaystyle \frac{U_l}{ml}}+\sqrt{2V(m)}\times `$
$`\left({\displaystyle \frac{\sqrt{C(rm)^2U_m}}{rm}}{\displaystyle \frac{\sqrt{C(ml)^2U_l}}{ml}}\right).`$ (21)
The right-hand side of Eq. (21) is a particular case of $`D_3(l,m,r)`$ given by Eq. (17) with $`V(l,m,r)=V(m)`$.
Momentum-conserving PNp-free models were proposed in PhysicaD and further studied in Submitted . They are the non-Hamiltonian models of the form
$`\ddot{m}=D_1(l,m,r),`$ (22)
where $`D_1`$ is given by Eq. (13). Notice that Eq. (22) can be mapped into the formulation of Eq. (7). Static kink solutions in this model can be found from Eq. (11).
If Eq. (11) is taken in the particular form of Eq. (20), then the momentum-conserving PNp-free model Eq. (22) and the energy-conserving PNp-free model Eq. (21) have exactly the same static kink solutions. It has been proved that a standard nearest-neighbor discrete Klein-Gordon model conserving both energy and momentum does not exist Submitted .
## V Application to the $`\varphi ^4`$ model
As an example, we will discretize the well-known $`\varphi ^4`$ field theory with the potential
$$V(\varphi )=\frac{1}{4}\left(1\varphi ^2\right)^2.$$
(23)
By construction, the PNp-free models derived above are written in singular form. In this form the equations are inconvenient in practical simulations and one may wish to find such particular cases when singularities disappear.
For example, for the energy-conserving PNp-free model expressed by Eqs. (19)-(21), singularity always disappears when $`G(\varphi )`$ is polynomial SpeightKleinGordon . Particularly, for the $`\varphi ^4`$ model with the potential Eq. (23), one obtains from Eq. (21) the following energy-conserving PNp-free discretization derived in SpeightKleinGordon
$`\ddot{m}=\left(C+{\displaystyle \frac{1}{6}}\right)(l+r2m)+m`$
$`{\displaystyle \frac{1}{18}}\left[2m^3+(m+l)^3+(m+r)^3\right],`$ (24)
whose static kink solution can be found from Eq. (20), which, for the $`\varphi ^4`$ potential, obtains the form
$`3\sqrt{2C}(ml)+m^2+lm+l^23=0.`$ (25)
Now let us turn to the momentum-conserving model. Substituting Eq. (11) into Eq. (22) we obtain
$`\ddot{m}=C(r2m+l)2{\displaystyle \frac{V(m,r)V(l,m)}{rl}}.`$ (26)
To remove the singularity, $`V(l,m)`$ should be taken in the symmetric form $`V(l,m)=V(m,l)`$, e.g., as
$`V(l,m)=(1/4)(\alpha /2)(m^2+l^2)+(\alpha 1/2)ml`$
$`+(\beta /2)\left(m^3+l^3\right)(\beta /2)ml\left(m+l\right)`$
$`+(\gamma /2)\left(m^4+l^4\right)+(\delta /2)ml\left(m^2+l^2\right)`$
$`\left(\gamma +\delta 1/4\right)m^2l^2,`$ (27)
with free parameters $`\alpha `$, $`\beta `$, $`\gamma `$, and $`\delta `$. In the continuum limit, when $`lm`$ and $`rm`$, Eq. (27) reduces to $`V(\varphi )`$. Substituting Eq. (27) into Eq. (26) we obtain the following momentum-conserving PNp-free $`\varphi ^4`$ model derived in Submitted
$`\ddot{m}=\left(C+\alpha \right)(l2m+r)+m`$
$`\beta (l^2+lr+r^2)+\beta m(l+r+m)`$
$`\gamma (l^3+r^3+l^2r+lr^2)\delta m(l^2+m^2+r^2+lr)`$
$`+(2\gamma +2\delta 1/2)m^2(l+r).`$ (28)
The momentum-conserving model Eq. (28) with $`\alpha =\beta =\gamma =\delta =0`$ can be written in the form
$`\ddot{m}=\left(1{\displaystyle \frac{m^2}{2C}}\right)C(l2m+r)+mm^3.`$ (29)
The following energy-conserving model, studied in Saxena ,
$$\ddot{m}=C(l2m+r)+\frac{mm^3}{1m^2/(2C)},$$
(30)
has the same continuum limit as model Eq. (29). Furthermore, it can be derived from Eq. (29) by multiplication with a factor $`e(C,l,m,r)=1/(1m^2/(2C))`$, which possesses a unit continuum limit. The model Eq. (29) is PNp-free and thus, model Eq. (30) is also PNp-free since they have the same static solutions derivable from $`C(ml)^2(1ml)^2/2=0`$. Thus, we have another example when energy-conserving and momentum-conserving PNp-free models have exactly the same static kink solutions.
It is interesting to note that the energy-conserving model of Eq. (30) cannot be constructed by the method reported in SpeightKleinGordon where discretization of the anharmonic term always involves $`\varphi _{n1}`$ and $`\varphi _{n+1}`$. More generally than it is done in SpeightKleinGordon , the problem of finding the energy-conserving PNp-free models can be formulated as follows. We need to discretize the potential energy of the Lagrangian Eq. (1) in a way that the corresponding equation of static equilibrium is satisfied when Eq. (9) is satisfied. Both energy-conserving models discussed above are the solutions of this problem.
As an example of model conserving neither energy, nor momentum we take Eq. (15) for the case of $`q(u,h)=u+Ahu^2`$ with constant $`A`$ and obtain
$$\ddot{m}=\frac{u_mu_l}{rl}\frac{(1+Ahu_l+Ahu_m)^2}{(1+2Ahu_l)(1+2Ahu_m)}.$$
(31)
This model can be obtained from the momentum-conserving model defined by Eq. (13) by multiplying by another function that reduces to unity in the continuum limit ($`h0`$). Obviously, the original momentum-conserving model and model Eq. (31) have the same static kink solutions. It can be demonstrated that these two models also have the same spectra of small amplitude vibrations and the same frequencies of kink internal modes.
## VI Numerics
In our recent work Submitted , some properties of kinks were compared for the “standard” energy-conserving $`\varphi ^4`$ discretization having PNp,
$`\ddot{m}=C(l+r2m)+mm^3,`$ (32)
with the PNp-free models conserving energy Eq. (24) and momentum Eq. (29).
It was found that the mobility of kinks in the PNp-free models is higher and also that in the momentum-conserving, PNp-free models a kink self-acceleration effect may be observed. The origin of the effect is the non-conservative (non self-adjoint) nature of the model which, however, can be noticed only for asymmetric trajectories of particles when kink passes by Submitted . If the trajectories are symmetric, there is no energy exchange with the surroundings and kink dynamics is the same as in energy-conserving models, e.g., the kink self-acceleration effect disappears. Kinks in some of the momentum-conserving models was found to have internal modes with frequencies above the phonon spectrum. Such modes do not radiate and they can have large amplitudes storing a considerable amount of energy.
Here we present/compare results for the energy-conserving PNp-free model Eq. (30) and the PNp-free model of Eq. (31), generally speaking, conserving neither energy nor momentum. For the latter model we take $`u_l`$ in the form of Eq. (11) where the $`\varphi ^4`$ potential is discretized according to Eq. (27) and, for the sake of simplicity, we set $`\alpha =\beta =\gamma =\delta =0`$. We obtain
$`\ddot{m}=[(1{\displaystyle \frac{m^2}{2C}})C(l2m+r)+mm^3]\times `$
$`{\displaystyle \frac{(1+Ahu_l+Ahu_m)^2}{(1+2Ahu_l)(1+2Ahu_m)}},`$
$`\mathrm{where}u_l=C(ml)^2(1ml)^2/2.`$ (33)
For $`A=0`$, Eq. (33) coincides with the momentum-conserving model Eq. (29).
In the model Eq. (33), the static kink solutions, phonon spectra, and frequencies of kink internal modes are $`A`$-independent. The energy-conserving model Eq. (30) has the same static kink solutions as model Eq. (33) but their spectra are different. The linear vibration spectrum of the vacuum for Eq. (33) is $`\omega ^2=2+(4C2)\mathrm{sin}^2(\kappa /2)`$ and that for Eq. (30) is $`\omega ^2=4C/(2C1)+4C\mathrm{sin}^2(\kappa /2)`$, while the one for the classical model Eq. (32) is $`\omega ^2=2+4C\mathrm{sin}^2(\kappa /2)`$.
The top panels of Fig. 1 present the boundaries of the linear vibration spectrum of the vacuum (solid lines) and the kink internal modes (dots) as the functions of lattice spacing $`h`$ for (a) the classical $`\varphi ^4`$ model of Eq. (32), (b) the PNp-free model of Eq. (30) conserving energy, and (c) the PNp-free model of Eq. (33) at $`A=0`$ conserving momentum. In PNp-free models kinks possess a zero frequency, Goldstone translational mode. Since all models presented in Fig. 1 share the same continuum $`\varphi ^4`$ limit, their spectra are very close for small $`h(<0.5)`$.
The bottom panels of Fig. 1 show the time evolution of kink velocity for the corresponding models at $`h=0.7`$ for kinks launched with different initial velocities. To boost the kink we used the semi-analytical solution for the normalized Goldstone mode, whose amplitude serves as a measure of the initial kink velocity. One can see that the mobility of kinks in the PNp-free models shown in (b) and (c) is higher than in the classical model having PNp and shown in (a). In the energy-conserving models shown in (a) and (b), the kink velocity decreases monotonically due to the energy radiation. Non-Hamiltonian momentum-conserving model in (c) shows the effect of kink self-acceleration discussed in Submitted .
It is interesting to study what happens when the parameter $`A`$ in Eq. (33) deviates from zero and the conservation law of the model (momentum conservation) disappears. We found that the effect of kink self-acceleration, which can be seen in the bottom panel of Fig. 1 (c) for $`A=0`$, remains for $`|A|<0.2`$ but the value of the kink velocity in the steady motion regime decreases with increase in $`|A|`$ as it is presented in Fig. 2. For $`|A|>0.2`$ kink self-acceleration effect disappears and kink velocity gradually decreases with time. From the above, we infer that properties such as the self-acceleration (for momentum-conserving models) or the Bogomol’nyi bounds (for energy-conserving discretizations) render such models rather special within the broader class of PNp free models. However, the critical ingredient for the more general feature of (static) PN absence exists in the form of a reduction of the second order problem into a first order.
## VII Conclusions
A general procedure for deriving discrete Klein-Gordon models whose static kinks can be placed anywhere with respect to the underlying lattice was described. Such models are called static PNp-free models. It was demonstrated that the models of this kind derived earlier SpeightKleinGordon ; PhysicaD ; Submitted ; Saxena are special cases of the wider family of models derived here.
Static kink solutions for the PNp-free models can be found from the nonlinear algebraic equation of the form $`u(l,m)=0`$, which is a discrete analog of the first integral of the static continuum Klein-Gordon equation of motion. This ensures the existence of static kink solutions at least for the regime of sufficiently weak discreteness and smooth background potential. The range of the discreteness parameter supporting stable static kinks varies according to the specific properties of the model.
In this paper we have discussed only nearest-neighbor discretizations. However, one can easily write down a PNp-free model involving more distant neighbors by replacing Eq. (18) with higher-order finite difference operators approximating Eq. (2), keeping the two-point approximation, Eq. (11), for the first integral of Eq. (10).
Discrete kinks in the static PNp-free models possess the zero-frequency translational Goldstone mode and they can (almost) freely move with at least infinitesimally small velocity. Such kinks are not trapped by the lattice and they can be accelerated by even weak external fields.
As a topic for future studies, it would be interesting to find any possible relation between models constructed here and models that support traveling kink solutions for finite kink velocity. Such connections are apparently under intense investigation dima and should provide a framework for understanding travelling in dispersive lattice systems.
|
warning/0506/math0506035.html
|
ar5iv
|
text
|
# Integrability of planar polynomial differential systems through linear differential equations.The second and third authors are partially supported by a MCYT grant number BFM 2002-04236-C02-01. The second author is partially supported by DURSI of Government of Catalonia’s Acció Integrada ACI2002-24.
## 1 Introduction
This paper deals with rational ordinary differential equations such as
$$\frac{dy}{dx}=\frac{Q(x,y)}{P(x,y)},$$
(1)
where $`Q(x,y)`$ and $`P(x,y)`$ are coprime polynomials with real coefficients. We associate to this rational equation a planar polynomial differential system by introducing an independent variable $`t`$ usually called time. Denoting by $`\dot{}=d/dt`$, we have
$$\dot{x}=P(x,y),\dot{y}=Q(x,y),$$
(2)
where $`(x,y)^2`$. This system defines the vector field $`𝒳=P(x,y)\frac{}{x}+Q(x,y)\frac{}{y}`$ over $`^2`$ and, equivalently, the $`1`$-form $`\omega =Q(x,y)dxP(x,y)dy`$. We indistinctively talk about equation (1) and system (2). Let d be the maximum degree of $`P`$ and $`Q`$. We say that system (2) is of degree d. When $`\mathrm{d}=2`$, we say that (2) is a quadratic system.
In order to simplify notation, we define $`[x,y]`$ as the ring of polynomials in two variables with real coefficients and $`(x,y)`$ as the field of rational functions in two variables with real coefficients, that is, the quotient field of the previous ring. Analogous definitions stand for $`[x]`$ and $`(x)`$.
We have an equation (1) defined in a certain class of functions, in this case, the rational functions with real coefficients $`(x,y)`$ and we consider the problem whether there is a first integral in another, possibly larger, class. For instance, as we will discuss later on, H. Poincaré stated the problem of determining when a system (2) has a rational first integral. A $`𝒞^k`$ function $`H:𝒰`$ such that it is constant on each trajectory of (2) and it is not locally constant is called a first integral of system (2) of class $`k`$ and the equation $`H(x,y)=c`$ for a fixed $`c`$ gives a set of trajectories of the system, but in an implicit way. When $`k1`$, these conditions are equivalent to $`\omega dH=0`$ and $`H`$ not locally constant. The problem of finding such a first integral and the functional class it must belong to is what we call the integrability problem.
To find an integrating factor or an inverse integrating factor for system (2) is closely related to finding a first integral for it. When considering the integrability problem we are also addressed to study whether an (inverse) integrating factor belongs to a certain given class of functions.
###### Definition 1
Let $`𝒲`$ be an open set of $`^2`$. A function $`\mu :𝒲`$ of class $`𝒞^k(𝒲)`$, $`k>1`$, that satisfies the linear partial differential equation
$$\omega d\mu =\mu d\omega ,$$
(3)
is called an integrating factor of system (2) on $`𝒲`$.
It has been shown than an easier function to find which also gives additional properties for a differential system (2) is the inverse of an integrating factor, that is, $`V=1/\mu `$, which is called inverse integrating factor.
We note that $`\{V=0\}`$ is formed by orbits of system (2). The function $`\mu =1/V`$ defines on $`𝒲\{V=0\}`$ an integrating factor of system (2), which allows the computation of a first integral of the system on $`𝒲\{V=0\}`$. The first integral $`H`$ associated to the inverse integrating factor $`V`$ can be computed through the integral $`H(x,y)=\omega /V,`$ and the condition (3) for $`\mu =1/V`$ ensures that this line integral is well defined.
The inverse integrating factors play an important role in two of the most difficult open problems of qualitative theory of planar polynomial vector fields, which are the center problem and $`16^{th}`$ Hilbert problem. In , it has been noticed that for many polynomial differential systems with a center at the origin there is always an inverse integrating factor $`V`$ globally defined in all $`^2`$, which is usually a polynomial. However, the first integral for a polynomial differential system with a center at the origin can be very complicated.
We say that a function $`f(x,y)`$ is an invariant for a system (2) if $`\omega df=kf`$ with $`k(x,y)`$ a polynomial of degree lower or equal than $`\mathrm{d}1`$, where d is the degree of the system. This polynomial $`k(x,y)`$ is called the cofactor of $`f(x,y)`$. In the previous equality and all along this paper we use the convention of identifying the space of functions over $`^2`$ and the space of $`2`$–forms over $`^2`$. In case $`f(x,y)=0`$ defines a curve in the real plane, this definition implies that $`\omega df`$ equals zero on the points such that $`f(x,y)=0`$. In case $`f(x,y)`$ is a polynomial we say that $`f(x,y)=0`$ is an invariant algebraic curve for system (2).
Let us consider $`f(x,y)=0`$ an invariant algebraic curve for system (2), we will always assume that $`f(x,y)`$ is an irreducible polynomial in $`[x,y]`$. Otherwise, it can be shown that each of its factors is an invariant algebraic curve for system (2). We will denote by $`n`$ the degree of the polynomial $`f(x,y)`$.
In , G. Darboux gives a method for finding an explicit first integral for a system (2) in case that $`\mathrm{d}(\mathrm{d}+1)/2+1`$ different irreducible invariant algebraic curves are known, where d is the degree of the system. In this case, a first integral of the form $`H=f_1^{\lambda _1}f_2^{\lambda _2}\mathrm{}f_r^{\lambda _r},`$ where $`f_i(x,y)=0`$ is an invariant algebraic curve for system (2) and $`\lambda _i`$ not all of them null, for $`i=1,2,\mathrm{},r`$, $`r`$, can be defined in the open set $`^2\mathrm{\Sigma }`$, where $`\mathrm{\Sigma }=\{(x,y)^2(f_1f_2\mathrm{}f_r)(x,y)=0\}.`$ The functions of this type are called Darboux functions. We remark that, particularly, if $`\lambda _i`$ , $`i=1,2,\mathrm{},r`$, $`H`$ is a rational first integral for system (2). In this sense J. P. Jouanoulou , showed that if at least $`\mathrm{d}(\mathrm{d}+1)+2`$ different irreducible invariant algebraic curves are known, then there exists a rational first integral.
The main fact used to prove Darboux’s theorem (and Jouanoulou’s improvement) is that the cofactor corresponding to each invariant algebraic curve is a polynomial of degree $`\mathrm{d}1`$. Invariant functions can also be used in order to find a first integral for the system. This observation permits a generalization of Darboux’s theory which is given in , where, for instance, non-algebraic invariant curves with an algebraic cofactor for a polynomial system of degree $`4`$ are presented. In our work, we give other families of systems with such invariant curves.
C. Christopher, in , studies the multiplicity of an invariant algebraic curve and gives the definition for exponential factor, which is a particular case of invariant for system (2).
###### Definition 2
Let $`h,g`$ be two coprime polynomials. The function $`e^{h/g}`$ is called an exponential factor for system (2) if for some polynomial $`k`$ of degree at most $`\mathrm{d}1`$, where d is the degree of the system, the following relation is fulfilled: $`\omega d\left(e^{h/g}\right)=ke^{h/g}.`$ As before, we say that $`k(x,y)`$ is the cofactor of the exponential factor $`e^{h/g}`$.
We note that an exponential factor $`e^{h/g}`$ does not define an invariant curve, but the next proposition, proved in , gives the relationship between both notions.
###### Proposition 3
If $`F=e^{h/g}`$ is an exponential factor and $`g`$ is not a constant, then $`g=0`$ is an invariant algebraic curve, and $`h`$ satisfies the equation $`\omega dh=(hk_g+gk_F)`$ where $`k_g`$ and $`k_F`$ are the cofactors of $`g`$ and $`F`$, respectively.
All these previous results are closely related to a result of M. F. Singer which represents an important progress in the resolution of the integrability problem when considering first integrals for a system (2) in the class of Liouville functions. Roughly speaking, we can define a Liouville function or a function which can be expressed by means of quadratures, as a function constructed from rational functions using composition, exponentiation, integration, and algebraic functions. A precise definition of this class of functions is given in .
###### Theorem 4
Let us consider the polynomial $`1`$–form $`\omega =QdxPdy`$ related to system (2). System (2) has a Liouville first integral if, and only if, $`\omega `$ has an inverse integrating factor of the form $`V=\mathrm{exp}\xi `$, where $`\xi `$ is a closed rational $`1`$–form.
We notice that the conditions on the function $`V`$ given in this Theorem can be restated as $`d\omega =\xi \omega `$ and $`d\xi =0`$.
Taking into account Theorem 4, C. Christopher gives the following result, which precises the form of the inverse integrating factor.
###### Theorem 5
If system (2) has an inverse integrating factor of the form $`\mathrm{exp}\xi `$ with $`d\xi =0`$ and $`\xi =\xi _1dx+\xi _2dy`$ where $`\xi _i`$, $`i=1,2`$, are rational functions in $`x`$ and $`y`$, then there exists an inverse integrating factor of system (2) of the form
$$V=\mathrm{exp}\{D/E\}C_i^{l_i},$$
where $`D`$, $`E`$ and the $`C_i`$ are polynomials in $`x`$ and $`y`$, and $`l_i`$.
Theorem 5 states that the search for Liouvillian first integrals can be reduced to the search of invariant algebraic curves and exponential factors.
In , H. Poincaré stated the following problem concerning the integration of an equation (1): Give conditions on the polynomials $`P`$ and $`Q`$ to recognize when there exists a rational first integral. As the same H. Poincaré noticed, a sufficient condition to solve this problem consists on finding an upper bound for the degree of the invariant algebraic curves for a given system (2). From Darboux’s result, it is known that for every polynomial vector field, there exists an upper bound for the possible degrees of irreducible invariant algebraic curves. However, it is a hard problem to explicitly determine such an upper bound. Some bounds have been given under certain conditions on the invariant curves, see D. Cerveau and A. Lins Neto work , or on the local behavior of critical points, see M. Carnicer’s work .
In this sense, A. Lins Neto conjectured that a polynomial system (2) of degree d with an invariant algebraic curve of degree high enough (where this bound only depends on d) would have a rational first integral. This conjecture has been shown to be false by several counterexamples. J. Moulin-Ollagnier gives a family of quadratic Lotka-Volterra systems, each with an invariant algebraic curve of degree $`2\mathrm{}`$, where $`\mathrm{}`$ is the parameter of the family, without rational first integral. A simpler example is given by C. Christopher and J. Llibre in . In a family of quadratic systems with an invariant algebraic curve of arbitrarily high degree without a Darboux first integral nor a Darboux inverse integrating factor is given. All these counterexamples exhibit a Liouvillian first integral.
The natural conjecture at this step, also given by A. Lins Neto, see , after the counterexample of J. Moulin-Ollagnier appeared, is that a polynomial system (2) of degree d with an invariant algebraic curve of degree high enough (where this bound only depends on d) has a Liouvillian first integral.
In this work we show a relationship between solutions of a class of systems (2) and linear homogeneous ordinary differential equations of order $`2`$ of the form
$$A_2(x)w^{\prime \prime }(x)+A_1(x)w^{}(x)+A_0(x)w(x)=0,$$
(4)
where $`x`$, $`w^{}(x)=dw(x)/dx`$ and $`w^{\prime \prime }(x)=dw^{}(x)/dx`$. We only consider equations (4) where $`A_i(x)[x]`$ for $`0i2`$ and $`A_2(x)0`$.
By means of a change of variable we rewrite an equation (4) as a polynomial differential system such that it has an invariant related to $`w(x)`$. In case $`w(x)`$ is a polynomial we get an invariant algebraic curve.
Moreover, we give an explicit first integral for all the systems built up by this method by means of two independent solutions of equation (4).
We give analogous results for a linear homogeneous ordinary differential equation of order $`1`$ such as
$$w^{}(x)+A(x)w(x)=0,$$
(5)
where $`x`$, $`w^{}(x)=dw(x)/dx`$ and $`A(x)(x)`$. All these results are given in Section 2.
In Section 3 we consider all the families of quadratic systems with an algebraic curve of arbitrarily high degree known until the moment of composition of this paper and we show that they all belong to the construction explained in Section 2. The families of quadratic systems with an algebraic curve of arbitrarily high degree studied in this paper are the ones appearing in and one example more first appearing in this work. This new example consists on a biparametrical family of quadratic systems, which we give an explicit expression of a first integral for, such that when one of the parameters is a natural number, say $`n`$, the system exhibits an irreducible invariant algebraic curve of degree $`n`$.
We give the explicit expression for the first integral of a certain system (2) by means of invariant functions for it, and applying the Generalized Darboux’s Theory as explained in where a new kind of first integrals, not only the Liouvillian ones as in classical theories, is described. We exemplify this result with the families of systems depending on parameters described in Section 3. We remark that the first integrals that we give in Section 2 are not, in general, of Liouvillian type. However, these first integrals are Liouvillian at the values of parameters which correspond to the systems with algebraic solutions. In the Subsection 3.4, we give an example of a $`3`$-parameter family of quadratic systems with a center at the origin which can be constructed by the method appearing in Section 2 from an equation (5).
A question suggested by these examples is whether there are polynomial systems which are not reversible nor Liouvillian integrable which have a center and can be integrated by means of Theorem 6, see Section 2. The work is related to this question as it gives an example of an analytic system, not polynomial, with a center which is not reversible nor Liouvillian integrable. All the known families of polynomial vector fields with a center at the origin are either Liouvillian integrable or reversible, see for the definition of reversibility. In , Żoła̧dek classifies all the reversible cubic systems with a center. The reversible systems may have a first integral not given by Liouville functions or no explicit form of a first integral may be known. For instance the reversible system $`\dot{x}=y+x^4`$, $`\dot{y}=x`$ has a first integral composed by Airy functions, see , and no Liouvillian first integral exists. The system $`\dot{x}=y^3+x^2y^2/2`$, $`\dot{y}=x^3`$ is an example given by Moussu, see , which has a center at the origin since it is a monodromic and reversible singular point and no explicit first integral is known for this system.
Since some examples of polynomial systems, which can be integrated by the method described in Section 2, appear after a birrational transformation, another suggested open question is if all the polynomial systems with a center are birrationally equivalent to one derived from Theorem 6 or from Theorem 11.
## 2 Homogeneous linear differential equations of <br>order $`2`$ and planar polynomial systems.
Let us consider a homogeneous linear differential equation of order $`2`$:
$$A_2(x)w^{\prime \prime }(x)+A_1(x)w^{}(x)+A_0(x)w(x)=0,$$
(6)
where $`w^{}(x)=dw(x)/dx`$, $`w^{\prime \prime }(x)=dw^{}(x)/dx`$, $`A_i(x)[x]`$, $`i=0,1,2`$, and $`A_2(x)0`$.
###### Theorem 6
Given $`g(x,y)=g_0(x,y)/g_1(x,y)`$ with $`g_i(x,y)[x,y]`$, $`g_1(x,y)0`$ and $`g/y0`$, each nonzero solution $`w(x)`$ of equation (6) is related to a finite number of solutions $`y=y(x)`$ of the rational equation
$$\frac{dy}{dx}=\frac{A_0(x)g_1^2+A_1(x)g_1g_0+A_2(x)g_0^2+A_2(x)\left(g_1\frac{g_0}{x}g_0\frac{g_1}{x}\right)}{A_2(x)\left(g_0\frac{g_1}{y}g_1\frac{g_0}{y}\right)},$$
(7)
by the functional change $`dw/dx=g(x,y)w(x)`$, which implicitly defines $`y`$ as a function of $`x`$.
Proof. Let us consider equation (6) and the functional change $`dw/dx=g(x,y)w(x)`$ where $`y=y(x)`$, that is, $`y`$ is implicitly defined as a function of $`x`$. This change may also be written as $`w(x)=\mathrm{exp}(_{x_0}^xg(s,y(s))𝑑s),`$ where $`x_0`$ is any constant, and it is injective. We see that it is not necessarily bijective unless the maximum degree of $`g_1(x,y)`$ and $`g_0(x,y)`$ in the variable $`y`$ equals $`1`$. But it defines a finite number of functions $`y(x)`$.
By this functional change, equation (6) becomes
$$w(x)\left(A_0(x)+gA_1(x)+g^2A_2(x)+A_2(x)\frac{dy}{dx}\frac{g}{y}+A_2(x)\frac{g}{x}\right)=0.$$
We have that $`w(x)`$ is a nonzero solution of (6) so this equation is equivalent to the ordinary differential equation of first order (7). Therefore, each non-zero solution $`w(x)`$ of equation (6) corresponds to a finite number of solutions $`y=y(x)`$ of the planar polynomial system (7).
###### Theorem 7
We consider the $`1`$-form related to equation (7)
$`\omega `$ $`=`$ $`\left(A_0(x)g_1^2+A_1(x)g_1g_0+A_2(x)g_0^2+A_2(x)\left(g_1{\displaystyle \frac{g_0}{x}}g_0{\displaystyle \frac{g_1}{x}}\right)\right)dx`$
$`A_2(x)\left(g_0{\displaystyle \frac{g_1}{y}}g_1{\displaystyle \frac{g_0}{y}}\right)dy.`$
Let $`w(x)`$ be any nonzero solution of equation (6). Then the curve defined by $`f(x,y)=0`$, with $`f(x,y):=g_1(x,y)w^{}(x)g_0(x,y)w(x)`$ is invariant for system (7) and has the polynomial cofactor
$$\begin{array}{cc}k(x,y)=\hfill & \left(A_0(x)\frac{g_1}{y}+A_1(x)\frac{g_0}{y}\right)g_1+A_2(x)g_0\frac{g_0}{y}\hfill \\ & +A_2(x)\left(\frac{g_1}{y}\frac{g_0}{x}\frac{g_0}{y}\frac{g_1}{x}\right).\hfill \end{array}$$
Proof. Let us consider $`f(x,y)`$ as defined above and let us compute $`\omega df`$:
$`\begin{array}{ccc}\omega df\hfill & =\hfill & ({\displaystyle \frac{g_1}{x}}w^{}(x)+g_1w^{\prime \prime }(x){\displaystyle \frac{g_0}{x}}w(x)g_0w^{}(x))A_2(x)(g_0{\displaystyle \frac{g_1}{y}}\hfill \\ & & g_1{\displaystyle \frac{g_0}{y}})+({\displaystyle \frac{g_1}{y}}w^{}(x){\displaystyle \frac{g_0}{y}}w(x)\left)\right[A_0(x)g_1^2+A_1(x)g_1g_0\hfill \\ & & +A_2(x)g_0^2+A_2(x)(g_1{\displaystyle \frac{g_0}{x}}g_0{\displaystyle \frac{g_1}{x}})]\hfill \end{array}`$
$`\begin{array}{ccc}& =\hfill & g_1\left(g_0{\displaystyle \frac{g_1}{y}}g_1{\displaystyle \frac{g_0}{y}}\right)A_2(x)w^{\prime \prime }(x)+\hfill \\ & & g_1\left(A_1(x)g_0+A_0(x)g_1\right)\left({\displaystyle \frac{g_1}{y}}w^{}(x){\displaystyle \frac{g_0}{y}}w(x)\right)+\hfill \\ & & \left(A_2(x)\left({\displaystyle \frac{g_0}{x}}{\displaystyle \frac{g_1}{y}}{\displaystyle \frac{g_1}{x}}{\displaystyle \frac{g_0}{y}}\right)+A_2(x)g_0{\displaystyle \frac{g_0}{y}}\right)\left(g_1w^{}(x)g_0w(x)\right)\hfill \end{array}`$
Since $`w(x)`$ is a solution of (6), we can substitute $`A_2(x)w^{\prime \prime }(x)`$ by $`A_1(x)w^{}(x)A_0(x)w(x)`$. Therefore,
$`\begin{array}{ccc}\omega df\hfill & =\hfill & [(A_0(x){\displaystyle \frac{g_1}{y}}+A_1(x){\displaystyle \frac{g_0}{y}})g_1+\hfill \\ & & A_2(x)g_0{\displaystyle \frac{g_0}{y}}+A_2(x)({\displaystyle \frac{g_1}{y}}{\displaystyle \frac{g_0}{x}}{\displaystyle \frac{g_0}{y}}{\displaystyle \frac{g_1}{x}})]f(x,y).\hfill \end{array}`$
Then, we have that the function $`f(x,y)`$ is an invariant for system (7) and has the written polynomial cofactor.
###### Theorem 8
Let $`\{w_1(x),w_2(x)\}`$ be a set of fundamental solutions of equation (6). We define $`f_i(x,y):=g_1(x,y)w_i^{}(x)g_0(x,y)w_i(x)`$, $`i=1,2`$. Then, system (7) has a first integral $`H(x,y)`$ defined by
$$H(x,y):=\frac{f_1(x,y)}{f_2(x,y)}=\frac{g_1(x,y)w_1^{}(x)g_0(x,y)w_1(x)}{g_1(x,y)w_2^{}(x)g_0(x,y)w_2(x)}.$$
Proof. By Theorem 7, we have that $`f_i(x,y):=g_1(x,y)w_i^{}(x)g_0(x,y)w_i(x)`$, $`i=1,2`$, are invariants for system (7), both with the polynomial cofactor
$$k(x,y)=\left(A_0(x)\frac{g_1}{y}+A_1(x)\frac{g_0}{y}\right)g_1+A_2(x)g_0\frac{g_0}{y}+A_2(x)\left(\frac{g_1}{y}\frac{g_0}{x}\frac{g_0}{y}\frac{g_1}{x}\right).$$
We remark that $`f_1/f_2`$ cannot be constant since the two solutions $`w_i(x)`$, $`i=1,2`$, are independent. Therefore,
$$\omega dH=\frac{f_2(\omega df_1)f_1(\omega df_2)}{f_2^2}=\frac{f_2kf_1f_1kf_2}{f_2^2}0.$$
So, $`H(x,y)`$ is a first integral of system (7).
###### Lemma 9
The function defined by
$$q(x):=A_2(x)\mathrm{exp}\left(_{x_0}^x\frac{A_1(s)}{A_2(s)}𝑑s\right)$$
is an invariant for system (7), with cofactor $`(A_1(x)+A_2^{}(x))\left(g_0\frac{g_1}{y}g_1\frac{g_0}{y}\right)`$.
We notice that $`q(x)`$ is a product of invariant algebraic curves and exponential factors for system (6), with complex exponents.
Proof. We compute $`\omega dq`$ and we have
$$\omega dq=\omega \frac{A_1(x)+A_2^{}(x)}{A_2(x)}qdx=(A_1(x)+A_2^{}(x))\left(g_0\frac{g_1}{y}g_1\frac{g_0}{y}\right)q.$$
We notice that this algebraic cofactor has degree $`\mathrm{d}1`$ provided that system (7) has degree d.
###### Proposition 10
We use the same notation as in Theorem 6. Let $`w(x)`$ be a nonzero solution of (6) and we define $`f(x,y):=w^{}(x)g(x,y)w(x)`$ and $`q(x)`$ as in Lemma 9. The function $`V(x,y)=q(x)f(x,y)^2`$ is an inverse integrating factor of system (7).
Proof. We only need to verify that $`\omega dV+Vd\omega =0`$. We have that
$$\begin{array}{cc}d\omega =\hfill & [2A_0(x)g_1\frac{g_1}{y}+A_1(x)(g_0\frac{g_1}{y}+g_1\frac{g_0}{y})+2A_2(x)g_0\frac{g_0}{y}\hfill \\ & +2A_2(x)(\frac{g_1}{y}\frac{g_0}{x}\frac{g_0}{y}\frac{g_1}{x})+A_2^{}(x)(g_0\frac{g_1}{y}g_1\frac{g_0}{y})].\hfill \end{array}$$
Then,
$$\begin{array}{cc}\hfill \omega dV=& \omega (2qfdf+f^2dq)=2qf(\omega df)+f^2(\omega dq)\hfill \\ \hfill =& [2A_0(x)g_1\frac{g_1}{y}+2A_1(x)g_1\frac{g_0}{y}+2A_2(x)g_0\frac{g_0}{y}\hfill \\ & +2A_2(x)\left(\frac{g_1}{y}\frac{g_0}{x}\frac{g_0}{y}\frac{g_1}{x}\right)+A_1(x)g_0\frac{g_1}{y}A_1(x)g_1\frac{g_0}{y}\hfill \\ & +A_2^{}(x)(g_0\frac{g_1}{y}g_1\frac{g_0}{y})]V\hfill \\ \hfill =& [2A_0(x)g_1\frac{g_1}{y}+A_1(x)(g_0\frac{g_1}{y}+g_1\frac{g_0}{y})+2A_2(x)g_0\frac{g_0}{y}\hfill \\ & +2A_2(x)(\frac{g_1}{y}\frac{g_0}{x}\frac{g_0}{y}\frac{g_1}{x})+A_2^{}(x)(g_0\frac{g_1}{y}g_1\frac{g_0}{y})]V\hfill \\ \hfill =& Vd\omega .\hfill \end{array}$$
We remark that Theorem 8 gives, in general, a non Liouvillian first integral for the planar polynomial systems (7). In Section 3 we analyze some polynomial systems constructed from Theorem 8 that have no Liouvillian first integral.
We consider now a linear homogeneous ordinary differential equation of order $`1`$ such as
$$w^{}(x)+A(x)w(x)=0,$$
(11)
where $`x`$, $`w^{}(x)=dw(x)/dx`$ and $`A(x)=A_0(x)/A_1(x)`$ with $`A_i(x)[x]`$ and $`A_1(x)0`$. We give analogous results for this case whose proofs are not given to avoid non useful repetitions.
###### Theorem 11
Given $`g(x,y)=g_0(x,y)/g_1(x,y)`$ with $`g_i(x,y)[x,y]`$, $`g_1(x,y)0`$ and $`g/y0`$ and $`h(x)=h_0(x)/h_1(x)`$ with $`h_i(x)[x]`$ and $`h_1(x)0`$, each nonzero solution $`w(x)`$ of equation (11) is related to a finite number of solutions $`y=y(x)`$ of the rational equation
$$\frac{dy}{dx}=\frac{A_1(x)h_0(x)g_1^2A_0(x)h_1(x)g_0g_1A_1(x)h_1(x)\left(g_1\frac{g_0}{x}g_0\frac{g_1}{x}\right)}{A_1(x)h_1(x)\left(g_1\frac{g_0}{y}g_0\frac{g_1}{y}\right)},$$
(12)
by the functional change
$$w(x)=g(x,y)\mathrm{exp}\left(_0^xA(s)𝑑s\right)\left[_0^x\mathrm{exp}\left(_0^sA(r)𝑑r\right)h(s)𝑑s\right].$$
###### Theorem 12
We consider the $`1`$-form related to equation (12)
$$\begin{array}{cc}\hfill \omega =& \left[A_1(x)h_0(x)g_1^2A_0(x)h_1(x)g_0g_1A_1(x)h_1(x)\left(g_1\frac{g_0}{x}g_0\frac{g_1}{x}\right)\right]dx\hfill \\ & A_1(x)h_1(x)\left(g_1\frac{g_0}{y}g_0\frac{g_1}{y}\right)dy.\hfill \end{array}$$
Let $`w(x)`$ be any nonzero solution of equation (11), that is, for $`C\{0\}`$ we have $`w(x)=C\mathrm{exp}\left(_0^xA(s)𝑑s\right)`$.
Then, the function
$$f(x,y):=g_1w(x)g_0+g_1\mathrm{exp}\left(_0^xA(s)𝑑s\right)\left[_0^x\mathrm{exp}\left(_0^sA(r)𝑑r\right)h(s)𝑑s\right]$$
is invariant for the polynomial system (12), with the polynomial cofactor
$$\begin{array}{cc}k(x,y)=\hfill & A_0(x)h_1(x)g_1\frac{g_0}{y}+A_1(x)h_0(x)g_1\frac{g_1}{y}\hfill \\ & +A_1(x)h_1(x)\left(\frac{g_0}{y}\frac{g_1}{x}\frac{g_1}{y}\frac{g_0}{x}\right).\hfill \end{array}$$
###### Lemma 13
The function $`q(x,y)=g_1(x,y)\mathrm{exp}\left(_0^xA(s)ds\right)`$ is an invariant for system (12) with the same polynomial cofactor as $`f(x,y)`$.
###### Theorem 14
We use the same notation as in Theorem 12 and Lemma 13. The function $`H(x,y)`$ defined by $`H(x,y):=f(x,y)/q(x,y)`$ is a first integral for system (12) and the function $`V(x,y):=A_1(x)h_1(x)g_1(x,y)q(x,y)`$ is an inverse integrating factor.
We remark that $`H(x,y)`$ is a Liouvillian function and, therefore, a system (12) has always a Liouville first integral.
In Section 3 we give an example of a 3-parameter family of quadratic systems with a center at the origin which can be constructed following Theorem 11.
## 3 Examples of families of quadratic systems.
### 3.1 Quadratic systems with invariant algebraic curves of arbitrarily high degree linear in one variable.
We first consider the examples of families of quadratic systems with algebraic solutions of arbitrarily high degree appearing in . In that work all the invariant algebraic curves linear in the variable $`y`$, that is, defined by $`f(x,y)=p_1(x)y+p_2(x),`$ where $`p_1`$ and $`p_2`$ are polynomials, are determined.
The example appearing in is a further study of an example appearing in and the example given in is also described in . We show that all these quadratic systems, with an invariant algebraic curve of arbitrary degree can be constructed by the method explained in the previous section. Moreover, we give the explicit expression of a first integral for any value of the parameter $`n`$, even in the case when $`n`$ is not a natural number. If $`n`$ is a natural number, we obtain the invariant algebraic curves of arbitrary degree and a Liouvillian first integral. However, when $`n`$ we obtain polynomial systems with a non Liouvillian first integral.
As it is shown in , all these families of systems can be written, after an affine change of variables if necessary, in the form
$$\dot{x}=\mathrm{\Omega }_1(x),\dot{y}=(2n+1)L^{}(x)\mathrm{\Omega }_1(x)\frac{n(n+1)}{2}\mathrm{\Omega }_1(x)\mathrm{\Omega }_1^{\prime \prime }(x)L(x)^2+y^2,$$
(13)
where $`\mathrm{\Omega }_1(x)`$ is any quadratic polynomial, $`L(x)`$ is any linear polynomial and $`{}_{}{}^{}=d/dx`$. We have that system (13) has an invariant curve $`f(x,y)=0`$, where $`f(x,y):=p_1(x)y+\mathrm{\Omega }_1(x)p_1^{}(x)L(x)p_1(x)`$, with a cofactor $`y+L(x)`$, where $`p_1(x)`$ is a solution of the second order linear differential equation
$$\mathrm{\Omega }_1(x)w^{\prime \prime }(x)+(\mathrm{\Omega }_1^{}(x)2L(x))w^{}(x)+\frac{n}{2}(4L^{}(x)(n+1)\mathrm{\Omega }_1^{\prime \prime }(x))w(x)=0.$$
(14)
In it is shown that, in case $`n`$, an irreducible polynomial of degree $`n`$ belonging to a family of orthogonal polynomials is a solution of equation (14). For instance, when $`\mathrm{\Omega }_1(x)=1`$, we get the Hermite polynomials, when $`\mathrm{\Omega }_1(x)=x`$, we get the Generalized Laguerre polynomials and when $`\mathrm{\Omega }_1(x)=1x^2`$, we get the Jacobi polynomials.
We consider again the general case in which $`n`$ and we define $`A_2(x):=\mathrm{\Omega }_1(x)`$, $`A_1(x):=\mathrm{\Omega }_1^{}(x)2L(x)`$ and $`A_0(x):=\frac{n}{2}(4L^{}(x)(n+1)\mathrm{\Omega }_1^{\prime \prime }(x))`$. We have the linear differential equation (14) in the same notation as in Theorem 6 and we consider $`g(x,y):={\displaystyle \frac{L(x)y}{A_2(x)}}`$.
The system obtained by the method explained in Section 2 exactly coincides with system (13). We consider a set of fundamental solutions of equation (14) $`\{w_1(x),w_2(x)\}`$ and applying Theorem 8, we have a first integral for system (13) for any value of the parameter $`n`$.
In case $`n`$ we have that $`w_1(x)`$ degenerates to a polynomial and $`w_1^{}(x)g(x,y)w_1(x)=0`$ coincides with the algebraic curve given in the work .
We explicitly give the first integral for each of the families described in and for $`n`$. We have that $`A_2(x)`$ is a non-null quadratic polynomial in this case, and depending on its number of roots, we can transform it by a real affine change of variable to one of the following forms: $`A_2(x)=1,x,x^2,1x^2,1+x^2`$.
If $`A_2(x)=1`$, we can choose $`L(x)=x`$ by an affine change of coordinates. A set of fundamental solutions $`\{w_1(x),w_2(x)\}`$ for (14) with $`n`$ is
$$\begin{array}{cc}w_1(x)=\hfill & 2^n\sqrt{\pi }\left(\frac{1}{\mathrm{\Gamma }\left(\frac{1n}{2}\right)}{}_{1}{}^{}F_{1}^{}(\frac{n}{2};\frac{1}{2};x^2)\frac{2x}{\mathrm{\Gamma }\left(\frac{n}{2}\right)}{}_{1}{}^{}F_{1}^{}(\frac{1n}{2};\frac{3}{2};x^2)\right),\hfill \\ w_2(x)=\hfill & 2^n\sqrt{\pi }\left(\frac{1}{\mathrm{\Gamma }\left(\frac{1n}{2}\right)}{}_{1}{}^{}F_{1}^{}(\frac{n}{2};\frac{1}{2};x^2)+\frac{2x}{\mathrm{\Gamma }\left(\frac{n}{2}\right)}{}_{1}{}^{}F_{1}^{}(\frac{1n}{2};\frac{3}{2};x^2)\right),\hfill \end{array}$$
where $`\mathrm{\Gamma }(x)`$ is the Euler’s–Gamma function defined by $`\mathrm{\Gamma }(x)=_0^{\mathrm{}}t^{x1}e^t𝑑t`$ and $`{}_{1}{}^{}F_{1}^{}(a;b;x)`$ is the confluent hypergeometric function defined by the series
$${}_{1}{}^{}F_{1}^{}(a;b;x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a)_k}{(b)_k}\frac{x^k}{k!},$$
with $`(a)_k=a(a+1)(a+2)\mathrm{}(a+k1)`$, the Pochhammer symbol. See for further information about these functions. So, a first integral for this system is the expression given in Theorem 8: $`H(x,y):=f_1(x,y)/f_2(x,y)`$, where
$$\begin{array}{cc}\hfill f_{1,2}(x,y)=& \pm \mathrm{\Gamma }\left(\frac{1n}{2}\right)[6(xyx^2+1){}_{1}{}^{}F_{1}^{}(\frac{1n}{2};\frac{3}{2};x^2)\hfill \\ & \mathrm{\hspace{0.17em}4}(n1)x^2{}_{1}{}^{}F_{1}^{}(\frac{3n}{2};\frac{5}{2};x^2)]+\hfill \\ & 3\mathrm{\Gamma }\left(\frac{n}{2}\right)\left[2nx{}_{1}{}^{}F_{1}^{}(1\frac{n}{2};\frac{3}{2};x^2)+(xy){}_{1}{}^{}F_{1}^{}(\frac{n}{2};\frac{1}{2};x^2)\right].\hfill \end{array}$$
When $`n`$, we have that (14) corresponds to the equation for Hermite polynomials and $`w_1(x)`$ coincides with the Hermite polynomial of degree $`n`$. The invariant algebraic curve given in corresponds to $`f_1(x,y)=0`$.
If $`A_2(x)=x`$, we choose $`L(x)=\frac{1}{2}(x\alpha )`$, where $`\alpha `$ is an arbitrary real constant, and a set of fundamental solutions for (14) is:
$$w_1(x)=\frac{(\alpha +1)_n}{\mathrm{\Gamma }(n+1)}{}_{1}{}^{}F_{1}^{}(n;\alpha +1;x),w_2(x)=x^\alpha {}_{1}{}^{}F_{1}^{}(\alpha n;1\alpha ;x).$$
The first integral for this system is $`H(x,y)=x^\alpha h_1(x,y)/h_2(x,y)`$ with
$$\begin{array}{cc}h_1(x,y)=\hfill & (2yx+\alpha )(\alpha +1){}_{1}{}^{}F_{1}^{}(n;\alpha +1;x)2nx{}_{1}{}^{}F_{1}^{}(1n;\alpha +2;x),\hfill \\ h_2(x,y)=\hfill & (2yx+\alpha )(\alpha 1){}_{1}{}^{}F_{1}^{}(\alpha n;1\alpha ;x)\hfill \\ & 2(\alpha +n)x{}_{1}{}^{}F_{1}^{}(1\alpha n;2\alpha ;x).\hfill \end{array}$$
The first integral as given in Theorem 8 is $`f_1(x,y)/f_2(x,y)`$ and we notice that $`H(x,y)=cf_1(x,y)/f_2(x,y)`$ where $`c\{0\}`$. We do not write $`c`$ in terms of the parameters of the system to simplify notation.
When $`n`$, we have that (14) corresponds to the equation of Generalized Laguerre polynomials and $`w_1(x)`$ coincides with the Generalized Laguerre polynomial $`L_n^{(\alpha )}`$. The invariant algebraic curve given in corresponds to $`f_1(x,y)=0`$, where $`f_1(x,y):=w_1^{}(x)g(x,y)w_1(x)`$.
If $`A_2(x)=x^2`$, the birrational transformation yet described in , $`x=1/X`$ and $`y=(1/X)(1/2Y)`$, makes this case equivalent to the previous one.
If $`A_2(x)=1x^2`$, we choose $`L(x)=\frac{1}{2}((\alpha +\beta )x+(\alpha \beta ))`$, where $`\alpha ,\beta `$ are two arbitrary real constants, and a set of fundamental solutions for (14) is:
$$\begin{array}{ccc}w_1(x)& =& \frac{(\alpha +1)_n}{\mathrm{\Gamma }(n+1)}{}_{2}{}^{}F_{1}^{}(n,1+\alpha +\beta +n;\alpha +1;\frac{1x}{2}),\end{array}$$
$$\begin{array}{ccc}w_2(x)& =& (1x)^\alpha {}_{2}{}^{}F_{1}^{}(\alpha n,1+\beta +n;1\alpha ;\frac{1x}{2}),\end{array}$$
where $`{}_{2}{}^{}F_{1}^{}(a_1,a_2;b;x)`$ is the hypergeometric function defined by
$${}_{2}{}^{}F_{1}^{}(a_1,a_2;b;x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1)_k(a_2)_k}{(b)_k}\frac{x^k}{k!}.$$
The first integral given in Theorem 8 is $`H(x,y)=(1x)^\alpha h_1(x,y)/h_2(x,y)`$, where
$$\begin{array}{cc}h_1=\hfill & n(1+\alpha +\beta +n)(x^21){}_{2}{}^{}F_{1}^{}(1n,2+\alpha +\beta +n;2+\alpha ;\frac{1x}{2})+\hfill \\ & (\alpha +1)((\alpha +\beta )x+(\alpha \beta )2y){}_{2}{}^{}F_{1}^{}(n,1+\alpha +\beta +n;1+\alpha ;\frac{1x}{2}),\hfill \\ h_2=\hfill & (\alpha 1)((\alpha \beta )x+(\alpha +b)+2y){}_{2}{}^{}F_{1}^{}(\alpha n,1+\beta +n;1\alpha ;\frac{1x}{2})+\hfill \\ & (\alpha +n)(1+\beta +n)(x^21){}_{2}{}^{}F_{1}^{}(1\alpha n,2+\beta +n;2\alpha ;\frac{1x}{2}).\hfill \end{array}$$
The first integral as given in Theorem 8 is $`f_1(x,y)/f_2(x,y)`$ and we notice that $`H(x,y)=cf_1(x,y)/f_2(x,y)`$ where $`c\{0\}`$. As before, we do not write $`c`$ in terms of the parameters of the system to simplify notation.
When $`n`$, we have that (14) corresponds to the equation of Jacobi polynomials and $`w_1(x)`$ coincides with the Jacobi polynomial $`P_n^{(\alpha ,\beta )}(x)`$ and the invariant algebraic curve given in corresponds to $`f_1(x,y)=0`$, where $`f_1(x,y):=w_1^{}(x)g(x,y)w_1(x)`$.
If $`A_2(x)=1+x^2`$ the complex affine change of variable $`x=\mathrm{i}X`$ makes this case equivalent to the previous one, as it is shown in .
We have re-encountered by this method all the examples appearing in from a unified point of view. In addition, in this work we have given an explicit expression of a first integral for each case and for any value of the parameter $`n`$. To this end, we have found invariants for the system and we have applied the generalization of Darboux’s method as explained in to be able to construct a first integral which is, in general, of non Liouvillian type.
### 3.2 A Lotka-Volterra system
As it has been explained in the introduction, the first counterexample to Lins Neto conjecture was given by J. Moulin-Ollagnier in . His example is a quadratic system with two invariant straight lines and an irreducible invariant algebraic curve $`f(x,y)=0`$ of degree $`2\mathrm{}`$ when $`\mathrm{}`$. This gives a family of systems depending on the parameter $`\mathrm{}`$ which have a Darboux inverse integrating factor when $`\mathrm{}`$ but no rational first integral. The method used in only shows the existence of such invariant algebraic curve but no closed formula to compute it is given. We give an explicit expression for an invariant by means of Bessel functions for any value of $`\mathrm{}\{\frac{1}{2}\}`$ which, in the particular case $`\mathrm{}`$ degenerates to the algebraic curve encountered in .
We show that after a birrational transformation, this example coincides with a system constructed with the method explained in Section 2. A birrational transformation is a rational change of variables whose inverse is also rational. This kind of transformations bring polynomial systems such (2) to polynomial systems and do not change the Liouvillian or non Liouvillian character of the first integral.
Let us consider the system appearing in but assuming that $`\mathrm{}\{\frac{1}{2}\}`$
$$\dot{x}=x\left(1\frac{x}{2}+y\right),\dot{y}=y\left(\frac{2\mathrm{}+1}{2\mathrm{}1}+\frac{x}{2}y\right).$$
(15)
We make the birrational transformation
$$x=\frac{4uv}{12\mathrm{}},y=\frac{12\mathrm{}}{4v},$$
whose inverse is
$$u=xy,v=\frac{12\mathrm{}}{4y}.$$
By this transformation, system (15) becomes
$$\dot{u}=\frac{2u}{12\mathrm{}},\dot{v}=\frac{12\mathrm{}}{4}+\frac{2\mathrm{}+1}{2\mathrm{}1}v+\frac{2u}{2\mathrm{}1}v^2.$$
(16)
We notice that the equation for the orbits satisfied by the variable $`v`$ as a function of $`u`$ is a Ricatti equation.
Let us consider $`g(u,v):=v`$ and the linear differential equation of order $`2`$ given by
$$uw^{\prime \prime }(u)+\frac{1}{2}(1+2\mathrm{})w^{}(u)\frac{1}{8}(12\mathrm{})^2w(u)=0.$$
(17)
Applying the method given in the previous section, this linear differential equation gives system (16) modulus a change of time.
A set of two fundamental solutions for equation (17) is given by
$$w_1(u)=u^{(12\mathrm{})/4}I_{\frac{1}{2}\mathrm{}}\left((12\mathrm{})\sqrt{\frac{u}{2}}\right),w_2(u)=u^{(12\mathrm{})/4}I_{\mathrm{}\frac{1}{2}}\left((12\mathrm{})\sqrt{\frac{u}{2}}\right),$$
(18)
provided that $`\mathrm{}`$ is not of the form $`\frac{1}{2}(12r)`$, with $`r`$ an integer number, because in this case $`w_1`$ and $`w_2`$ are linearly dependent. The function $`I_\nu (u)`$ is the Modified Bessel function defined by the solution of the second order differential equation
$$u^2w^{\prime \prime }(u)+uw^{}(u)(u^2+\nu ^2)w(u)=0,$$
(19)
and being bounded when $`u0`$ in any bounded range of $`\mathrm{arg}(u)`$ with $`e(u)0`$. See for further information about this function.
Hence, by Theorem 8 we have that $`H(u,v)=f_1(u,v)/f_2(u,v)`$, where $`f_i(u,v):=w_i^{}(u)vw_i(u)`$ for $`i=1,2`$, is a first integral for system (16). For a sake of simplicity we consider the following renaming of the independent variable $`u=2z^2/(12\mathrm{})^2`$. This is not a birrational transformation and that’s why we only use it to simplify notation. The function $`H`$ writes as:
$$H=\frac{(12\mathrm{})^2I_{\left(\frac{1+2\mathrm{}}{2}\right)}(z)4vzI_{\left(\frac{2\mathrm{}1}{2}\right)}(z)}{(12\mathrm{})^2I_{\left(\frac{1+2\mathrm{}}{2}\right)}(z)4vzI_{\left(\frac{2\mathrm{}1}{2}\right)}(z)}.$$
(20)
By Theorem 7 we have that $`f_i(u,v)`$, $`i=1,2`$ are invariants with the same polynomial cofactor $`k`$ for system (16), so the curve $`f(u,v)=0`$ given by $`f(u,v)=\pi z^{2\mathrm{}+1}(f_1^2(u,v)f_2^2(u,v))`$ is also an invariant. We multiply by $`\pi `$ only for esthetic reasons.
Now we assume that $`\mathrm{}`$ and we show that $`f=0`$ defines an invariant algebraic curve. To this end we use the following formulas appearing in . When $`\nu \frac{1}{2}=n`$, we define $`c(n)=n\pi \sqrt{1}/2`$ and the following relation is satisfied:
$$\begin{array}{cc}I_\nu (z)=\hfill & \frac{1}{\sqrt{z}}\mathrm{e}^{c(n)}\sqrt{\frac{2}{\pi }}\{\mathrm{sinh}(c(n)z)\underset{k=0}{\overset{\frac{2|\nu |1}{4}}{}}\frac{(|\nu |+2k\frac{1}{2})!}{(2k)!(|\nu |2k\frac{1}{2})!(2z)^{2k}}\hfill \\ & +\mathrm{cosh}(c(n)z)\underset{k=0}{\overset{\frac{2|\nu |3}{4}}{}}\frac{(|\nu |+2k+\frac{1}{2})!}{(2k+1)!(|\nu |2k\frac{3}{2})!(2z)^{2k+1}}\},\hfill \end{array}$$
(21)
where $`x`$ stands for the greatest integer $`k`$ such that $`kx`$ and $`|\nu |`$ stands for the absolute value.
From the former equation we obtain the following two equalities, with $`\nu \frac{1}{2}=n`$ and $`\mathrm{}`$,
$$\begin{array}{ccc}I_\nu ^2(z)I_\nu ^2(z)\hfill & =\hfill & \frac{2}{\pi z}\underset{k=0}{\overset{n}{}}(1)^{k+1}\frac{(2nk)!(2n2k)!}{k!((nk)!)^2}\left(\frac{1}{2z}\right)^{2(nk)},\hfill \end{array}$$
(22)
$$\begin{array}{c}I_{\mathrm{}+\frac{1}{2}}(z)I_{\mathrm{}\frac{1}{2}}(z)I_{(\mathrm{}+\frac{1}{2})}(z)I_{(\mathrm{}\frac{1}{2})}(z)=(1)^{\mathrm{}}\frac{2}{\pi z}\hfill \\ [\left(\underset{i=0}{\overset{\frac{\mathrm{}}{2}}{}}\frac{(\mathrm{}+2i)!}{(2i)!(\mathrm{}2i)!}\left(\frac{1}{2z}\right)^{2i}\right)\left(\underset{j=0}{\overset{\frac{\mathrm{}2}{2}}{}}\frac{(\mathrm{}+2j)!}{(2j+1)!(\mathrm{}2j2)!}\left(\frac{1}{2z}\right)^{2j1}\right)\hfill \\ \left(\underset{i=0}{\overset{\frac{\mathrm{}1}{2}}{}}\frac{(\mathrm{}+2i+1)!}{(2i+1)!(\mathrm{}2i1)!}\left(\frac{1}{2z}\right)^{2i1}\right)\left(\underset{j=0}{\overset{\frac{\mathrm{}1}{2}}{}}\frac{(\mathrm{}+2j1)!}{(2j)!(\mathrm{}2j1)!}\left(\frac{1}{2z}\right)^{2j}\right)].\hfill \end{array}$$
(23)
Then, we have that $`f_1(z,v)=(12\mathrm{})^2I_{\mathrm{}+\frac{1}{2}}(z)4vzI_{\mathrm{}\frac{1}{2}}(z)`$ and $`f_2(z,v)=(12\mathrm{})^2I_{(\mathrm{}+\frac{1}{2})}(z)4vzI_{(\mathrm{}\frac{1}{2})}(z)`$, and we write $`f`$ arranged in powers of $`v`$:
$`f(z,v)`$ $`=`$ $`\pi z^{2\mathrm{}+1}((12\mathrm{})^4(I_{\mathrm{}+\frac{1}{2}}^2(z)I_{(\mathrm{}+\frac{1}{2})}^2(z))`$
$`8vz(12\mathrm{})^2(I_{\mathrm{}+\frac{1}{2}}(z)I_{\mathrm{}\frac{1}{2}}(z)I_{(\mathrm{}+\frac{1}{2})}(z)I_{(\mathrm{}\frac{1}{2})}(z))`$
$`+16v^2z^2(I_{\mathrm{}\frac{1}{2}}^2(z)I_{(\mathrm{}\frac{1}{2})}^2(z))).`$
Let us consider each coefficient of $`v`$ in $`f(z,v)`$ separately and we will show that it is an even polynomial in the variable $`z`$. The coefficient in $`f(z,v)`$ of $`v^0`$ is:
$$\pi z^{2\mathrm{}+1}(12\mathrm{})^4(I_{\mathrm{}+\frac{1}{2}}^2(z)I_{(\mathrm{}+\frac{1}{2})}^2(z)),$$
which by equation (22) is an even polynomial in the variable $`z`$ of degree $`2\mathrm{}`$. The coefficient in $`f(z,v)`$ of $`v^2`$ is:
$$16\pi z^{2\mathrm{}+3}(I_{\mathrm{}\frac{1}{2}}^2(z)I_{(\mathrm{}\frac{1}{2})}^2(z)),$$
which also by equation (22) is an even polynomial in the variable $`z`$ of degree $`2\mathrm{}+2`$. Finally, the coefficient in $`f(z,v)`$ of $`v^1`$ is:
$$8\pi (12\mathrm{})^2z^{2\mathrm{}+2}(I_{\mathrm{}+\frac{1}{2}}(z)I_{\mathrm{}\frac{1}{2}}(z)I_{(\mathrm{}+\frac{1}{2})}(z)I_{(\mathrm{}\frac{1}{2})}(z)),$$
(24)
which by equation (23) is an even polynomial in the variable $`z`$ of degree $`2\mathrm{}`$.
Hence, we have that $`f(z,v)`$ is an even polynomial in the variable $`z`$ of total degree $`2\mathrm{}+4`$. When rewriting $`z=(12\mathrm{})\sqrt{u}/\sqrt{2}`$ we have that $`f(u,v)`$ is a polynomial of total degree $`\mathrm{}+2`$ which is irreducible. The fact of being irreducible is easily seen because it is a polynomial of degree two in $`v`$ and it cannot be decomposed in linear factors (the discriminant is not a polynomial raised to an even power) and the coefficients of $`v^0`$ and $`v^2`$ do not have any root in common.
Undoing the birrational change of variables we get that $`f(x,y)`$ is an irreducible polynomial of degree $`2\mathrm{}`$ given by:
$`f(x,y)`$ $`=`$ $`x^{\mathrm{}+\frac{1}{2}}y^{\mathrm{}\frac{1}{2}}[2y(I_{\mathrm{}+\frac{1}{2}}^2(z)I_{(\mathrm{}+\frac{1}{2})}^2(z))+`$
$`2\sqrt{2}\sqrt{xy}\left(I_{\mathrm{}+\frac{1}{2}}(z)I_{\mathrm{}\frac{1}{2}}(z)I_{(\mathrm{}+\frac{1}{2})}(z)I_{(\mathrm{}\frac{1}{2})}(z)\right)+`$
$`x(I_{\mathrm{}\frac{1}{2}}^2(z)I_{(\mathrm{}\frac{1}{2})}^2(z))],`$
where $`z`$ is the same variable as before, that is, $`z=\frac{(12\mathrm{})\sqrt{xy}}{\sqrt{2}}`$.
By equation (20) we can write the first integral for system (15) for any value of $`\mathrm{}\{\frac{1}{2}(12r)r\}`$:
$$H(x,y)=\frac{\sqrt{2y}I_{\left(\frac{1+2\mathrm{}}{2}\right)}(z)\sqrt{x}I_{\left(\frac{2\mathrm{}1}{2}\right)}(z)}{\sqrt{2y}I_{\left(\frac{1+2\mathrm{}}{2}\right)}(z)\sqrt{x}I_{\left(\frac{2\mathrm{}1}{2}\right)}(z)}.$$
We have studied system (15) for any value of the parameter $`\mathrm{}\{\frac{1}{2}(12r)r\}`$ giving an explicit expression for a first integral using Theorem 6 and the Generalized Darboux’s theory as explained in . This first integral is not of Liouvillian type. Moreover, we give one of its invariants with a polynomial cofactor. In the particular case $`\mathrm{}`$, this invariant is the invariant algebraic curve whose existence was proved in .
### 3.3 A new example of a family of quadratic systems with an invariant algebraic curve of arbitrarily high degree
We give another example of a family of quadratic systems with an irreducible invariant algebraic curve of degree $`2\mathrm{}`$ when $`\mathrm{}`$, where $`\mathrm{}`$ is a parameter of the family. This family also depends on the parameter $`a`$.
Let us consider the quadratic system
$$\begin{array}{ccc}\dot{x}\hfill & =\hfill & (2a1)\mathrm{}xa(2\mathrm{}1)y+2a(a\mathrm{})(2\mathrm{}1)x^22a^2(2\mathrm{}1)^2xy,\hfill \\ \dot{y}\hfill & =\hfill & y(2(2a1)\mathrm{}+2a(2a2\mathrm{}1)(2\mathrm{}1)x4a^2(2\mathrm{}1)^2y),\hfill \end{array}$$
(25)
where $`a,\mathrm{}`$ which satisfy $`a0`$, $`\mathrm{}\frac{1}{2}`$ and $`(2\mathrm{}1)a^22\mathrm{}0`$. An straightforward computation shows that system (25) has $`y=0`$ and $`yx^2=0`$ as invariant algebraic curves.
Let us consider the following birrational transformation $`x=Y`$, $`y=XY^2`$ whose inverse is $`X=y/x^2`$ and $`Y=x`$. In these new variables system (25) becomes
$$\begin{array}{ccc}\dot{X}\hfill & =\hfill & 2a(2\mathrm{}1)(X1)XY,\hfill \\ \dot{Y}\hfill & =\hfill & ((2a1)\mathrm{}+a(2\mathrm{}1)(2a2\mathrm{}X)Y2a^2(2\mathrm{}1)^2XY^2)Y.\hfill \end{array}$$
(26)
By a change of the time variable we can divide this system by $`Y`$ and the resulting system coincides with the one described in Theorem 6 taking $`A_2(X):=2X(X1)^2`$, $`A_1(X):=(2\mathrm{}2a+3X)(X1)`$, $`A_0(X):=\mathrm{}(12a)`$ and $`g(X,Y):=a(2\mathrm{}1)Y/(X1)`$. The equation $`A_2(X)w^{\prime \prime }(X)+A_1(X)w^{}(X)+A_0(X)w(X)=0`$ has the following set of fundamental solutions in this case:
$`w_1(X)`$ $`=`$ $`(X1)^{\mathrm{}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}}\mathrm{},\mathrm{};a\mathrm{};X),`$
$`w_2(X)`$ $`=`$ $`(X1)^{\mathrm{}}X^{1a+\mathrm{}}{}_{2}{}^{}F_{1}^{}(1a,{\displaystyle \frac{3}{2}}\mathrm{};2a+\mathrm{};X).`$
By Theorem 7, $`f_i(X,Y)=w_i^{}(X)g(X,Y)w_i(X)`$, $`i=1,2`$, define invariants with a polynomial cofactor for system (26). Moreover, by Theorem 8 we have a non Liouvillian first integral given by $`H(X,Y)=f_1(X,Y)/f_2(X,Y)`$.
In the particular case $`\mathrm{}`$, we notice that $`f_1(X,Y)=0`$ is a rational function. It is an easy computation to show that this rational function is a polynomial when rewritten in coordinates $`x`$ and $`y`$. This polynomial gives place to an invariant algebraic curve of degree $`2\mathrm{}`$ for system (25). That is, by undoing the birrational transformation, we deduce that $`f_1(x,y)`$ is an irreducible invariant algebraic curve for system (25), given by:
$`f_1(x,y)`$ $`=`$ $`2(a\mathrm{})(\mathrm{}+(2\mathrm{}1)ax)x^{2\mathrm{}1}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}}\mathrm{},\mathrm{};a\mathrm{};{\displaystyle \frac{y}{x^2}})`$
$`+\mathrm{}(2\mathrm{}1)x^{2\mathrm{}3}(x^2y){}_{2}{}^{}F_{1}^{}({\displaystyle \frac{3}{2}}\mathrm{},1\mathrm{};1+a\mathrm{};{\displaystyle \frac{y}{x^2}}).`$
It is easy to see that the polynomial $`f_1(x,y)`$ has degree $`2\mathrm{}`$ and the cofactor associated to the invariant algebraic curve $`f_1(x,y)=0`$ is $`\mathrm{}(2\mathrm{}1)((2a1)+4a(a\mathrm{})x4(2\mathrm{}1)a^2y)`$.
The first integral for (25) is given by $`H(x,y)=y^a\mathrm{}f_1(x,y)/h(x,y)`$, where
$$\begin{array}{cc}\hfill h(x,y)=& 2(a\mathrm{}2)[(a\mathrm{}1)x^2+\hfill \\ & (1aa(2\mathrm{}1)x)y\left]x^{72a}_2F_1\right(1a,\frac{3}{2}a;2a\mathrm{};\frac{y}{x^2})+\hfill \\ & (a1)(2a3)x^{52a}(x^2y)y{}_{2}{}^{}F_{1}^{}(2a,\frac{5}{2}a;3a+\mathrm{};\frac{y}{x^2}).\hfill \end{array}$$
We notice that when both $`a`$ and $`\mathrm{}`$ belong to the set of natural numbers, we have that $`h(x,y)=0`$ is an invariant algebraic curve different from $`f_1(x,y)=0`$. Then we have a quadratic system with a rational first integral $`H(x,y)`$ with arbitrary degree.
### 3.4 A complete family of quadratic systems with a center at the origin
In this subsection we give an example of a $`3`$-parameter family of quadratic systems with a center at the origin which can be constructed using Theorem 11. The family encountered corresponds to the reversible case, see .
The family of quadratic systems depends on twelve parameters, but up to affine transformations and positive time rescaling, we get a family of five essential parameters. We have taken a system (12) and we have chosen $`g(x,y):=y^2`$, $`h(x):=2x(dx1)/(1+ax)`$ and $`A(x):=2b/(1+ax)`$, where $`a,b,d`$ are real parameters. Using Theorem 11, we have encountered the $`3`$-parameter family of quadratic systems next described. We remark that in spite of the simplicity of the chosen polynomials $`g(x,y)`$, $`A(x)`$ and $`h(x)`$, we amazingly obtain the complete family of quadratic systems with a reversible center at the origin. We notice that other choices of the functions $`g(x,y)`$, $`A(x)`$ and $`h(x)`$ would give place to other families of polynomial systems.
Let us recall that a center is an isolated singular point of an equation (2) with a neighborhood foliated of periodic orbits. When the linear approximation of an equation (2) near a singular point has non-null purely imaginary eigenvalues, the point can be a center or a focus. To distinguish between these two possibilities is the so-called center problem. H. Poincaré gave a method to solve it by defining a numerable set of values, called Liapunov-Poincaré constants, which are all zero when the singular point is a center and at least one of them is not null when it is a focus. When these constants are computed from a family of systems, they are polynomials on the coefficients of the family. Hilbert’s Nullstellensatz ensures that there always exists a finite number of independent polynomials which generates the whole ideal made up with all these Liapunov-Poincaré polynomials. The zero-set of these independent polynomials gives place to the center subfamilies. The reader is referred to for a survey on this subject.
The computation of these center cases for the family of quadratic systems was done by Dulac for the case of complex systems and a proof for real systems is given in . We also refer Bautin who showed the existence of only three independent constants. The computation of the zero set of these three independent values gives place to four complete families of quadratic systems with a center at the origin which are described in .
Let us now consider an equation (12) such as $`(1+ax)w^{}(x)+2bw(x)=0`$ and $`g(x,y)`$ and $`h(x)`$ as formerly defined. The rational equation as constructed in Theorem 11 is
$$\frac{dy}{dx}=\frac{x+dx^2by^2}{y+axy},$$
which gives the corresponding quadratic planar system
$$\dot{x}=y+axy,\dot{y}=x+dx^2by^2.$$
(27)
We suppose that $`ab(a+b)(a+2b)(a+b+d)0`$. In case this value is zero, the origin of system (27) is still a center but with a Darboux integrating factor instead of a Darboux first integral. This particular case can also be studied by our method, but we do not write it to avoid giving examples without essential differences.
By Theorem 12 we have that $`f(x,y)`$ is an invariant of system (27) with cofactor $`2by`$, where $`f(x,y)`$ is given by
$$\begin{array}{ccc}f(x,y)\hfill & =\hfill & b(a+b)(a+2b)w(x)(a+b+d)(1+ax)^{\frac{2b}{a}}\hfill \\ & & b(a+b)(a+2b)y^2+b(a+2b)dx^2\hfill \\ & & 2b(a+b+d)x+a+b+d,\hfill \end{array}$$
(28)
with $`w(x)`$ any non-zero solution of $`(1+ax)w^{}(x)+2bw(x)=0`$, that is, $`w(x)=C(1+ax)^{\frac{2b}{a}}`$.
Choosing $`C=(a+b+d)(b(a+b)(a+2b))^1`$ we get an invariant conic. System (27) has two invariant algebraic curves, the former conic with cofactor $`2by`$ and an invariant straight line given by $`1+ax=0`$ with cofactor $`y`$. The Darboux first integral
$$H(x,y)=(1+ax)^{\frac{2b}{a}}f(x,y)$$
coincides with the first integral described in Theorem 14.
The origin of this system is a center since it is a monodromic singular point with a continuous first integral defined in a neighborhood of it. This example addresses to the thought that other families of polynomial systems of higher degree with a center at the origin can be easily obtained by this method, avoiding the cumbersome computation of Poincaré-Liapunov constants.
|
warning/0506/hep-th0506138.html
|
ar5iv
|
text
|
# hep-th/0506138 A practical model for cosmic (p,q) superstrings.
## 1 Introduction
Cosmic strings have come in and out of fashion a number times, and whilst it is clear that they do not play a lead role in structure formation they may nevertheless have observational consequences . Indeed, it is hoped that advanced LIGO and LISA will be able to detect the characteristic radiation of a string network . The context of the current revival of cosmic strings dates back to work of Witten in perturbative string theory where it was concluded that superstrings could not act as cosmic strings. Since then it has become clear that there are more possibilities for the compact dimensions of string theory, which has led to the issue of cosmic superstrings being reconsidered ; for a nice overview see . The particular networks that we are interested in are associated to the low energy dynamics of type IIB string theory which, in ten dimensions, contains two different species of string. First of all there is the fundamental, F-string which carries charge under the Neveu Schwartz-Neveu Schwartz two-form potential, secondly there is the Dirichlet, D-string carrying charge under the Ramond-Ramond two-form potential. The types of string that motivate the current work are $`(p,q)`$ strings which are a supersymmetric bound state of F and D strings. These come about because parallel F and D strings are not supersymmetric and can reduce their energy by the F string breaking on the D string with the endpoints moving off to infinity , the resulting configuration is then supersymmetric and has the following tension,
$`\mu _{p,q}=\sqrt{(\mu _p)^2+(\mu _q)^2}=\mu _F\sqrt{p^2+q^2/g_s^2}.`$ (1.1)
where $`\mu `$ refers to the tension and $`g_s`$ is the string coupling. When we consider the physics of four dimensions we note that there are more general ways to get string solutions by utilizing the branes of string theory. If we were to wrap n-1 of the spatial dimensions of an n-brane around a compact cycle then such an object would appear as a string in four dimensions, with couplings and tensions that depended upon the details of the compactification. One model which has been worked on in detail is the KKLMMT model , using a warped compactification with matter fields living on a D3 brane in the throat. The tension of objects is then affected by the warping, with the redshift factor allowing for a large range of effective string tension. In principle there can be many such throats, each with a different warping, this would then produce strings with a variety of different tensions. Modelling such situations may well follow similar lines to those put forward here, adding an extra U(1) for each species.
Modelling the evolution of a $`(p,q)`$ network involves different approaches to that of the standard Nielsen-Olesen string due to the way $`(p,q)`$ re-connect as two strings pass through each other. Abelian strings simply swap partners as they interact but this channel is not open to $`(p,q)`$ strings as this would violate charge conservation. For such partner-swapping to take place an intermediate string is created, joining together the two initial strings in much the same way as for non-Abelian gauge vortices . This raises the possibility that these networks get more tangled as they evolve, eventually freezing out as the Universe continues to expand. This would cause the string network to eventually dominate the energy density of the Universe . An important issue therefore is to decide whether these networks reach a scaling solution , meaning that the string energy density forms a constant fraction of the background energy density, with the network length scale increasing proportional to time. Analytic and numerical approaches are now being developed to approach this problem , here we present a simple model containing gauged vortices which has many of the required properties, moreover it is well suited to numerical implementation by adapting the available codes for evolving Abelian gauge vortices. This is in counter-distinction to modelling these networks with non-Abelian gauge field theories, where imposing lattice gauge symmetry is numerically intensive .
The issue of the correct field theory to use when describing the low energy physics of D-strings in four dimensions is now converging to a supersymmetric model containing: chiral fields transforming in the usual way under U(1) gauge transformations; a chiral field, representing an axion, with a shift symmetry that gets gauged under the U(1); and a vector multiplet to carry the magnetic flux of the vortex. This builds on earlier work where the axion multiplet was not included . An obvious drawback of the model put forward here is that it is not based on a supersymmetric field theory. However, the simple form of coupling between the two U(1) sectors in our model leads us to believe that it should be possible to connect two supersymmetric U(1) gauge theories to yield a theory with vortex bound states of the type described below, whether such vortices would be BPS remains to be seen.
We organize the paper by first introducing the model in section 2 and describing how the vacuum manifold depends on the parameters. The parameter space is then reduced by considering the BPS limit of the theory in section 2.1. In section 3 we construct the ansatz for the vortices and calculate their energy, we then draw our conclusions in section 4.
## 2 the model
The model is based on the simplest way to couple two independent U(1) gauge theories and has been seen before, albeit in a different guise . There it was used to construct cosmic strings which could form a condensate in their core and become superconducting. For that to work the parameters are chosen such that one of the Higgs fields acquires a non-zero vacuum expectation value (vev) while the other sits at zero; here we require the “opposite”. What we need is for the vacuum to be such that both Higgs fields acquire a vev, allowing for there to be independent string solutions following the usual arguments . Although our action is the same as those for the superconducting string models, because we use the opposite parameter regime we write it in a way more suited to our analysis
$``$ $`=`$ $`D_\mu \overline{\varphi }D^\mu \varphi 𝒟_\mu \overline{\psi }𝒟^\mu \psi {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }{\displaystyle \frac{1}{4}}_{\mu \nu }^{\mu \nu }V(|\varphi |,|\psi |).`$ (2.2)
Where we have the following gauge-covariant derivatives, field strengths and potential,
$`D_\mu `$ $`=`$ $`_\mu ieA_\mu ,𝒟_\mu =_\mu igB_\mu ,`$ (2.3)
$`F_{\mu \nu }`$ $`=`$ $`_\mu A_\nu _\nu A_\mu ,_{\mu \nu }=_\mu B_\nu _\nu B_\mu ,`$ (2.4)
$`V(|\varphi |,|\psi |)`$ $`=`$ $`{\displaystyle \frac{\lambda _1}{4}}\left(\overline{\varphi }\varphi \eta ^2\right)^2+{\displaystyle \frac{\lambda _2}{4}}\left(\overline{\psi }\psi \nu ^2\right)^2\kappa \left(\overline{\varphi }\varphi \eta ^2\right)\left(\overline{\psi }\psi \nu ^2\right).`$ (2.5)
For generic values of the coupling constants we find that the critical “points” are given by,
$`(\overline{\varphi }\varphi ,\overline{\psi }\psi )_0`$ $`=`$ $`(0,0),(\eta ^2,\nu ^2),(\eta ^2{\displaystyle \frac{2\kappa }{\lambda _1}}\nu ^2,0),(0,\nu ^2{\displaystyle \frac{2\kappa }{\lambda _2}}\eta ^2).`$ (2.6)
When constructing the superconducting strings it was the last two critical points that were of interest, as these have one broken and one unbroken U(1) symmetry, and so the parameters were chosen to make these minima. As we are going to need two broken U(1)s then we need to make the critical point $`(\eta ^2,\nu ^2)`$ a minimum. To gain further insight on the vacuum structure we complete the square as follows
$`V(|\varphi |,|\psi |)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[\sqrt{\lambda _1}(\overline{\varphi }\varphi \eta ^2)\pm \sqrt{\lambda _2}(\overline{\psi }\psi \nu ^2)\right]^2[\kappa \pm {\displaystyle \frac{1}{2}}\sqrt{\lambda _1\lambda _2}](\overline{\varphi }\varphi \eta ^2)(\overline{\psi }\psi \nu ^2),`$ (2.7)
and note that there are special values of $`\kappa `$ where the manifold of minima increases in dimension. For $`\kappa =\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$ we get a vacuum which is an ellipsoid,
$`\sqrt{\lambda _1}\overline{\varphi }\varphi +\sqrt{\lambda _2}\overline{\psi }\psi =\sqrt{\lambda _1}\eta ^2+\sqrt{\lambda _2}\nu ^2,`$ (2.8)
while for $`\kappa =+\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$ we get a vacuum which is a hyperboloid
$`\sqrt{\lambda _1}\overline{\varphi }\varphi \sqrt{\lambda _2}\overline{\psi }\psi =\sqrt{\lambda _1}\eta ^2\sqrt{\lambda _2}\nu ^2.`$ (2.9)
It is these special values of $`\kappa `$ which separate the different classes of vacuum structure. For $`\kappa <\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$, the minima are at $`(\eta ^2\frac{2\kappa }{\lambda _1}\nu ^2,0)`$, $`(0,\nu ^2\frac{2\kappa }{\lambda _2}\eta ^2)`$, which is the parameter range for superconducting strings. For $`\frac{1}{2}\sqrt{\lambda _1\lambda _2}<\kappa <+\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$, the minima are at $`(\eta ^2,\nu ^2)`$, which is the range we are interested in. For $`\kappa >+\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$ the potential is unbounded below and so is unphysical. Within our range of interest, $`\frac{1}{2}\sqrt{\lambda _1\lambda _2}<\kappa <+\frac{1}{2}\sqrt{\lambda _1\lambda _2}`$ there is another special value of $`\kappa `$, namely zero. For $`\kappa =0`$ we see that the two U(1) are entirely decoupled and so the vortices of U(1)<sub>A</sub> do not talk to those of U(1)<sub>B</sub>, in that sense we could think of two different vortex species as being neutrally bound. It is therefore natural to expect that $`\kappa =0`$ is a boundary between the two vortex species having negative or positive binding energy, i.e. forming bound states or not. To see which sign of $`\kappa `$ gives bound states consider two separated vortices, one of each species. At the $`A`$ vortex we have $`\varphi =0`$, $`|\psi |\nu `$ while at the $`B`$ vortex we $`\psi =0`$, $`|\varphi |\eta `$. Now consider the $`A`$ vortex, we note from (2.5) that for $`\kappa >0`$ we can lower the potential by reducing $`|\psi |`$, this can be achieved by bringing the vortices closer, i.e. they attract each other. Putting together all of the above arguments leads us to the following parameter regime for bound $`AB`$ vortices
$`0<\kappa <{\displaystyle \frac{1}{2}}\sqrt{\lambda _1\lambda _2}.`$ (2.10)
### 2.1 BPS considerations
Although we understand the parameter space there are still a large number of parameters, and if we are to make any progress we need to make some choices. To make further restrictions we take our inspiration from the Bogomol’nyi limit of vortices, . Adapting those arguments to our case with $`\kappa =0`$ the Hamiltonian takes the particular form of a sum of squares and boundary terms if we choose $`\lambda _1=2e^2`$ and $`\lambda _2=2g^2`$. With that, one then sees that for a state with winding number $`n`$ in U(1)<sub>A</sub> and $`m`$ in U(1)<sub>B</sub> the minimum energy configuration has
$`E(\kappa =0,n,m)`$ $`=`$ $`2\pi \eta ^2n+2\pi \nu ^2m.`$ (2.11)
Where we have made the further parameter choice $`e=g=1`$. In order to make a comparison between the bound-state energy of our system and (1.1) we introduce the following quantity,
$`ϵ_{(n,m)}=\sqrt{(2\pi \eta ^2)^2+(2\pi \nu ^2)^2}.`$ (2.12)
What we shall now do is explore the remaining parameter space ($`\eta `$, $`\nu `$ and $`\kappa `$) to see how the mass of $`(n,m)`$ vortices compares to $`(p,q)`$ strings. The reason we keep $`\eta `$ and $`\nu `$ is that, from (2.11), they are expected to govern the tension of the basic constituents with $`\mu _F2\pi \eta ^2`$ and $`\mu _F/g_s2\pi \nu ^2`$.
## 3 Vortex solutions
What we shall now do is proceed to an evaluation of the tension of some low lying bound states. To be specific we shall look at two cases: firstly we consider $`\eta =\nu `$; secondly we look at $`\nu =2\eta `$ in order to give a mass separation between the $`A`$ and $`B`$ vortices, just as fundamental and Dirichlet strings have different tensions.
The ansatz we use for a static vortex with winding numbers $`(n,m)`$ is as follows,
$`\varphi `$ $`=`$ $`\eta f(r)\mathrm{exp}(in\theta ),`$ (3.13)
$`A_\theta `$ $`=`$ $`{\displaystyle \frac{n}{e}}\alpha (r),`$ (3.14)
$`\psi `$ $`=`$ $`\nu p(r)\mathrm{exp}(im\theta ),`$ (3.15)
$`B_\theta `$ $`=`$ $`{\displaystyle \frac{m}{e}}\beta (r),`$ (3.16)
for which the equations of motion following from (2.2) are
$`f^{\prime \prime }+{\displaystyle \frac{1}{r}}f^{}{\displaystyle \frac{n^2}{r^2}}f(\alpha 1)^2{\displaystyle \frac{1}{2}}\lambda _1\eta ^2(f^21)f+\kappa \nu ^2(p^21)f`$ $`=`$ $`0,`$ (3.17)
$`p^{\prime \prime }+{\displaystyle \frac{1}{r}}p^{}{\displaystyle \frac{m^2}{r^2}}p(\beta 1)^2{\displaystyle \frac{1}{2}}\lambda _2\nu ^2(p^21)p+\kappa \eta ^2(f^21)p`$ $`=`$ $`0,`$ (3.18)
$`\alpha ^{\prime \prime }{\displaystyle \frac{1}{r}}\alpha ^{}2e^2f^2\eta ^2(\alpha 1)`$ $`=`$ $`0,`$ (3.19)
$`\beta ^{\prime \prime }{\displaystyle \frac{1}{r}}\beta ^{}2g^2p^2\nu ^2(\beta 1)`$ $`=`$ $`0.`$ (3.20)
Note that we have kept $`e`$, $`g`$, $`\lambda _1`$, $`\lambda _2`$ in these equations even though we shall only be using the values as described above, $`e=g=\frac{1}{2}\lambda _1=\frac{1}{2}\lambda _2=1`$. The asymptotic boundary conditions are that $`f(r\mathrm{})1`$, $`p(r\mathrm{})1`$, $`\alpha (r\mathrm{})1`$, $`\beta (r\mathrm{})1`$, necessary for a finite energy configuration. The boundary conditions at the vortex core depend on the winding number; if $`\varphi `$ has a non-zero winding then $`f(0)=0`$ and $`\alpha (0)=0`$ otherwise $`f^{}(0)=0`$ and $`\alpha (0)=1`$. Similar conditions hold for $`\psi `$.
### 3.1 Case I, $`\eta =\nu `$
The first example is for $`\eta =\nu `$ and we can use our understanding of the potential gained earlier to see how the profile functions should behave. In Fig. LABEL:fig:vacuum we have a representation of the vacuum with: the points $`(\pm 1,\pm 1)`$ and $`(\pm 1,1)`$ being the vacua for $`\kappa `$ in our region of interest, $`1<\kappa <1`$; the circle is the vacuum for $`\kappa =1`$; and the 45 lines are the vacuum for $`\kappa =1`$. On the diagram we have shown how the Higgs field profile functions vary as we move from one vacuum at $`x\mathrm{}`$ to another at $`x\mathrm{}`$ for a vortex with winding in the $`\varphi `$ field but not $`\psi `$. The various solid lines correspond to the vortex solutions for different values of $`\kappa `$, the straight line at $`\psi /\nu =1`$ is the $`\kappa =0`$ case where $`\varphi `$ and $`\psi `$ are decoupled. The lines below this case are for $`\kappa `$ increasing, and they approach the 45 lines as $`\kappa `$ nears unity. If we decrease $`\kappa `$ below zero then the Higgs profile functions mark out a curve in $`\varphi \psi `$ space which tends toward the circle, which is the vacuum for $`\kappa =1`$.
In order to solve (3.17-3.20) we used a numerical approach described in called successive over-relaxation. We tested our numerics against the analytic BPS solution and we are confident that the energies quoted in the appendices are correct to four decimal places. What is clear from table 1 is that vortices with higher winding numbers do indeed form bound states, with the bound state energies being lower than the sum of their constituents. Moreover, as $`\kappa `$ is increased toward unity (the upper bound on $`\kappa `$) the binding of vortices is stronger. The results of table 1 have been plotted in Fig. LABEL:fig:tension11 with the energies scaled in order to make comparison with (2.12), which has been shown with dotted lines. The plot consists of three “rays” corresponding to different values of $`m`$, with the lowest such ray being $`m=0`$ and the upper being $`m=2`$. Within each ray we show how the tension depends on $`\kappa `$, with the upper line in each ray being $`\kappa =0.1`$ and the lowest being $`\kappa =0.9`$. We can see from this plot that the bound state spectrum does not follow the $`(p,q)`$ tension spectrum of (2.12) particularly well, but for $`\kappa `$ near its upper bound it is closest. In a real network we would expect there to be fewer high winding strings as they are heavier, as such, taking $`\kappa `$ close to unity may well give a suitable approximation.
### 3.2 Case II, $`2\eta =\nu `$
In order to represent the difference in tension between a D-string and and F-string we can choose to give different values to $`\eta `$ and $`\nu `$. For the purposes of illustration we took $`\nu ^2=4\eta ^2`$ and repeated the above process to produce table 2, which contains the tension of the vortices for the first few windings. Again we see that higher winding strings are bound states, with an energy lower than the sum of their constituent parts. Fig. LABEL:fig:tension14 illustrates the data in table 2, with the dotted lines representing (2.12). We note that the specific tension formula for $`(p,q)`$ strings seems to be modelled better in this second case, where the basic strings have different tensions, for $`\kappa `$ close to its upper bound. This means that simulations using different values for $`\eta `$ and $`\nu `$ will more closely model the dynamics of a $`(p,q)`$ network.
## 4 Conclusion
In trying to model the dynamics of a network of $`(p,q)`$ strings coming from string theory one must try both analytic and numerical approaches , each with its advantages and disadvantages. The properties of inter-commutation of $`(p,q)`$ strings, i.e. the presence of three-string vertices, makes it natural model these networks with a non-Abelian gauge theory containing a rich spectrum of vortices. Unfortunately such simulations are extremely demanding in terms of cpu hours, making large scale simulations intractable. As one needs a large dynamic range in order to observe scaling behaviour this presents us with a problem; what we have presented here is a practical resolution to this impasse. By using two Abelian gauge symmetries it is a simple matter to adapt existing cosmic string code to the study of $`(p,q)`$ networks.
An important matter that this work does not reveal is the details of the interaction between two vortices of different type. In particular we have not shown that such vortices inter-commute to create string junctions; it is possible that they simply pass through each other as they would for $`\kappa =0`$. However, as string junctions certainly exist in this model we expect inter-commutation to occur.
A natural follow-up of this work would be to try to extend the proposed model incorporating the ideas of . In it was argued that D-strings can be modelled by a U(1) supersymmetric gauge field theory containing a chiral axion field as well as the usual chiral Higgs fields. If we are to have a field theory description of $`(p,q)`$ strings along the lines presented here one would need to couple together two such field theories. It is certainly possible to couple two such field theories together, however it is unclear that one would find supersymmetric bound-states of vortices as a solution.
Acknowledgements The author would like to thank Mark Hindmarsh for various suggestions and PPARC for financial support.
Appendices
## Appendix A Data for the cases $`\nu ^2=\eta ^2`$, $`\nu ^2=4\eta ^2`$.
|
warning/0506/math0506047.html
|
ar5iv
|
text
|
# Introduction
## Introduction
Among compact Riemannian manifolds, the sphere is remarkable in several ways. For example, it admits a large group of conformal transformations, and the corresponding conformal vector fields are known to contain important information about the geometry of the sphere. What is perhaps not so well-known is the fact that the action of these vector fields also determines a large part of global analysis on the sphere, notably the exact eigenvalues of most natural differential and pseudodifferential operators. Already the example of the Laplace operator on functions is perhaps a little surprising.
In , a method was developed for finding the spectrum of intertwining operators for certain representations of semisimple groups. In this calculation scheme, one first inputs the spectral data of a differential spectrum generating operator. The idea is that this spectral data is readily accessible, obtained just from quadratic Casimir data for finite-dimensional representations of compact groups. Using this information, one may generate the (much less accessible) spectral data on the intertwining operators, which occur in series parameterized by a complex number $`r`$. Typically, a subseries of differential operators occurs at values of $`r`$ which are positive integral, in a sense appropriate to the particular series.
A special situation that may occur in such a calculation is when the spectrum generating operator is closely related to one of the differential intertwinors. The simplest interpretations of “closely related” here are (1) when the operators are identical, or differ by a constant additive shift; (2) when the spectrum generating operator is a polynomial in one of the differential intertwinors. In some of these cases, there may be some many relations among the spectral data that one can generate the spectra of the operators in question with no input at all; that is, the spectra are spontaneously generated.
In this paper, we execute this process of spontaneous generation in some examples, and derive some related results. Among the central relations are a quadratic equation (9) relating “adjacent” eigenvalues of the conformal Laplacian on scalar fields, and a cubic equation (38) which serves a similar purpose for the Dirac operator on spinor fields. Once the spectra of these fundamental operators are in place, one may go on to find parameterized families of operators of various orders which are intertwining for conformal transformations and vector fields in the same sense as are the conformal Laplacian and Dirac operator.
An interesting feature of this construction is that it proves completeness of the eigenvalue list it obtains – any “wrong” eigenvalue will generate lower and lower ones, until some basic estimate is violated. And even though no information on spherical harmonics is input, the fact that eigenfunctions are spherical harmonics is a consequence of the construction. (See, for example, Corollary 6.)
Our method, combined with Beckner’s sharp Hardy-Littlewood- Sobolev inequalities on the sphere, is perfectly suited for deriving the sharp form of Gross’ entropy inequality on the sphere, estimating the integral of $`f^2\mathrm{log}(f)`$ for a positive smooth function. This argument culminates in Theorem 8 below.
Of course, some of the spectral resolutions we get here have been known for some time; for example, the spectra of the Laplacian (on functions), and of the Dirac operator on the sphere. The spectra of the other operators is less well known, but nevertheless can be obtained by specializing formulas in the literature (for example ). However, the philosophy is avoid all heavy machinery, and derive the spectra just from a couple of elementary operator commutation relations. Among the possible virtues of this sort of “primitivist” derivation is the prospect of deriving such spectra at an early point of a course in differential geometry, quantum mechanics, or relativity.
## The Laplacian and conformal Laplacian on the sphere
Let $`n2`$, and let $`\rho _0`$ be the azimuthal angle on the sphere $`S^n`$ with homogeneous coordinate functions $`x_0,x_1,\mathrm{},x_n`$. That is,
$$x_0=\mathrm{cos}\rho _0,$$
so that $`\rho (p)`$ is the angle between an indeterminate point $`p`$ and the point $`(1,0,\mathrm{},0)`$. More invariantly, we can define an azimuthal angle $`\rho (p,q)`$ with any desired point $`q`$ in place of $`(1,0,\mathrm{},0)`$. In particular, the azimuthal angle from the point where $`x_i=1`$ will be denoted $`\rho _i`$ (for each $`i`$).
The proper conformal vector fields on $`S^n`$ are generated by the $`T_i:=(\mathrm{sin}\rho _i)(/\rho _i)`$ for $`j=0,\mathrm{},n`$. The meaning of $`/\rho _i`$ is: complete $`\rho _i`$ to a coordinate system by taking a coordinate system on the latitude $`\{\rho _i=\mathrm{const}\}`$ (a copy of the sphere $`S^{n1}`$), and compute with these other coordinates held constant. The $`T_i`$ are conformal vector fields:
$$L(T_i)g=2x_ig,$$
where $`g`$ is the round metric on $`S^n`$ and $`L`$ denotes the Lie derivative.
In the more general setting of pseudo-Riemannian manifolds $`(M,g)`$, let $`d`$ be the exterior derivative on functions, and let $`\delta `$ be its formal adjoint. The operator $`\mathrm{\Delta }=\delta d`$ is the Laplacian, and the conformal Laplacian is the operator
$$D:=\mathrm{\Delta }+\frac{n2}{4(n1)}K,$$
where $`K`$ is the scalar curvature of $`g`$. This satisfies the conformal covariance relation
$$\mathrm{\Omega }^{(n+2)/2}D_{\mathrm{\Omega }^2g}f=D_g(\mathrm{\Omega }^{(n2)/2}f)$$
(1)
for $`\mathrm{\Omega }>0`$ and $`f`$ smooth functions on $`M`$. If $`(M,g)`$ admits a conformal vector field $`T`$ with $`L_Tg=2\omega g`$, then (1) immediately leads, via the local flow, to
$$D\left(T+\frac{n2}{2}m(\omega )\right)=\left(T+\frac{n+2}{2}m(\omega )\right)D,$$
(2)
where $`m(\omega )`$ is multiplication by $`\omega `$. On the other hand,
$$[D,m(\phi )]=[\mathrm{\Delta },m(\phi )]=m(\mathrm{\Delta }f)2\iota (d\phi )d,$$
(3)
where $`\iota `$ is interior multiplication.
Specializing to the round sphere $`S^n`$, let
$$U_j:=T_j+\frac{n}{2}x_j.$$
Then (2) reads
$$D(U_im(x_i))=(U_i+m(x_i))D.$$
(4)
In particular, applying this relation to the function 1 and using the fact that the scalar curvature of round $`S^n`$ is $`n(n1)`$, so that
$$D=\mathrm{\Delta }+\frac{n(n2)}{4}\text{ on }S^n,$$
(5)
we get
$$\left(\mathrm{\Delta }+\frac{n(n2)}{4}\right)\frac{n2}{2}x_i=\frac{n+2}{2}x_i\frac{n(n2)}{4},$$
that is,
$$\mathrm{\Delta }x_i=nx_i.$$
(6)
Specializing (3), we have
$$\begin{array}{cc}\hfill [D,m(x_i)]& =m(\mathrm{\Delta }x_i)2\iota (dx_i)d\hfill \\ & =nm(x_i)+2T_i\hfill \\ & =2U_i,\hfill \end{array}$$
(7)
where we have used (6) together with the fact that $`dx_i`$ corresponds to $`T_i`$ under the metric identification.
Let $`E(\lambda )`$ denote the $`\lambda `$-eigenspace of $`D`$, and suppose for some $`\lambda `$ we have $`0\phi E(\lambda )`$. Given another number $`\mu `$ and an operator $`T`$, denote the compression of $`T`$ acting $`E(\lambda )E(\mu )`$ by $`{}_{\mu }{}^{}|T|_\lambda `$.
Given $`\mu `$, if we compress (7) and (4) to operators from $`E(\lambda )`$ to $`E(\mu )`$, we have
$$\begin{array}{cc}\hfill {}_{\mu }{}^{}|U_i|_\lambda & =\frac{\mu \lambda }{2}{}_{\mu }{}^{}|m(x_i)|_\lambda ,\hfill \\ \hfill \mu {}_{\mu }{}^{}|(U_im(x_i))|_\lambda & =\lambda {}_{\mu }{}^{}|(U_i+m(x_i))|_\lambda ,\hfill \end{array}$$
(8)
from which
$$\mu (\frac{\mu \lambda }{2}1){}_{\mu }{}^{}|m(x_i)|_\lambda =\lambda (\frac{\mu \lambda }{2}+1){}_{\mu }{}^{}|m(x_i)|_\lambda .$$
For fixed $`\lambda `$, this is a quadratic equation in $`\mu `$:
$$\mu ^22\lambda \mu 2\mu +\lambda ^22\lambda =0.$$
(9)
The solutions
$$\lambda ^\pm :=\lambda +1\pm \sqrt{4\lambda +1};$$
(10)
are candidates for new eigenvalues of $`D`$. Since $`\mathrm{\Delta }`$ is a nonnegative operator, (5) gives
$$\lambda \frac{n(n2)}{4}.$$
(11)
Thus
$$\lambda ^{}<\lambda <\lambda ^+$$
unless $`n=2`$ and $`\lambda =0`$, in which case $`\lambda ^{}=\lambda `$.
Moreover, since (9) is symmetric in $`\mu `$ and $`\lambda `$, we have that $`\mu `$ is a root of the equation obtained starting with $`\lambda `$ iff $`\lambda `$ is a root of the equation obtained starting with $`\mu `$. That is,
$$(\lambda ^+)^{}=\lambda =(\lambda ^{})^+.$$
The obvious main questions are:
* If $`E(\lambda )0`$, are the spaces $`E(\lambda ^\pm )`$ necessarily nonzero?
* Can we generate all $`\mu `$ for which $`E(\mu )0`$ by starting with a single $`E(\lambda )0`$ and iterating the process (10)?
The answer to each question is yes, with the exception that $`E(\lambda ^{})`$ vanishes for what should be the bottom eigenvalue $`\lambda `$, namely $`n(n2)/4`$. In fact, both questions are answered by the same calculation: any “wrong” eigenvalue keeps generating lower and lower “wrong” eigenvalues, until we get one that violates the estimate (11).
A weakness of the calculation is that it cannot immediately tell the dimension of the eigenspaces $`E(\lambda )`$. Of course these dimensions are easily obtainable with an injection of a small amount of Lie theory, but that would violate the spirit of the present calculation. Note however, that in Corollary 6 we prove (from our point of view) the correspondence with spherical harmonics.
Stipulating the statements just above (which will be proved below), we may start either with the fact that
$$1E(0,\mathrm{\Delta })=E(n(n2)/4,D)$$
or
$$x_iE(n,\mathrm{\Delta })=E((n+2)n/4,D)\text{(from (}\text{6}\text{))}$$
and iterate (10) to compute that the eigenvalues of $`D`$ are the
$$\lambda _j:=\left(\frac{n2}{2}+j\right)\left(\frac{n}{2}+j\right)\text{ for }j=0,1,2,\mathrm{},$$
(12)
or equivalently, that the eigenvalues of $`\mathrm{\Delta }`$ are the $`j(n1+j)`$.
To begin to answer the bullet point questions above, let us look more closely at the eigenfunctions (as opposed to just eigenvalues) that our process is generating. If $`\phi E(\lambda ,D)`$, then
$$D(x_i\phi )=\lambda x_i\phi +2U_i\phi $$
(13)
by (7). By (4) and (13),
$$D(U_i\phi )=D(x_i\phi )+\lambda (U_i+x_i)\phi =(\lambda +2)U_i\phi +2\lambda x_i\phi .$$
(14)
If we would like $`(U_i+cx_i)\phi `$ to be in $`E(\mu ,D)`$, then
$$\mu (U_i+cx_i)\phi =(\lambda +2+2c)U_i\phi +(2+c)\lambda x_i\phi .$$
If the $`U_i`$ and $`x_i`$ terms on the two sides are to agree, we need
$$c(\lambda +2+2c)=\lambda (2+c).$$
This is a quadratic equation on $`c`$, with the roots
$$c^\pm =(\lambda \lambda ^{})/2.$$
The eigenvalue $`\mu `$ corresponding to $`c^+`$ (resp. $`c^{}`$) is easily computed to be $`\lambda ^+`$ (resp. $`\lambda ^{}`$). Thus we have:
###### Proposition 1
If $`\phi E(\lambda ,D)`$, then
$$\begin{array}{cc}\hfill P_i\phi :=& \left(U_i+\frac{1}{2}(\lambda \lambda ^{})x_i\right)\phi E(\lambda ^+,D),\hfill \\ \hfill M_i\phi :=& \left(U_i+\frac{1}{2}(\lambda \lambda ^+)x_i\right)\phi E(\lambda ^{},D).\hfill \end{array}$$
Of course there is a priori the danger that $`P_i\phi `$ vanishes for $`i=0,\mathrm{},n`$, or that each $`M_i\phi `$ does, even though $`\phi 0`$. We now proceed to rule out this danger, except in the case where it is expected (applying $`M_i`$ to the eigenspace where $`1`$ lives). That is, we shall compute
$$\underset{i}{}M_iP_i\phi \text{ and }\underset{i}{}P_iM_i\phi .$$
If the first of these is nonzero, it must be that $`P_i\phi 0`$ for some $`i`$; similarly, if the second is nonzero, some $`M_i\phi `$ does not vanish identically. We begin with
$$\underset{i}{}M_iP_i\phi =\underset{i}{}\left(U_i+\frac{1\sqrt{1+4\lambda ^+}}{2}x_i\right)\left(U_i+\frac{1+\sqrt{1+4\lambda }}{2}x_i\right)\phi .$$
Since
$$\sqrt{1+4\lambda ^\pm }=\pm 2+\sqrt{1+4\lambda },$$
(15)
the above is
$$\begin{array}{c}\underset{i}{}\left(U_i+\frac{3\sqrt{1+4\lambda }}{2}x_i\right)\left(U_i+\frac{1+\sqrt{1+4\lambda }}{2}x_i\right)\phi =\hfill \\ \underset{i}{}\left((U_im(x_i))+c^{}m(x_i)\right)\left((U_im(x_i))c^{}m(x_i)\right)\phi .\hfill \end{array}$$
This simplifies to
$$\begin{array}{c}\{_i(U_im(x_i))^2(c^{})^2m\left(_ix_i^2\right)\hfill \\ c^{}\left(_i\{(U_im(x_i))m(x_i)m(x_i)(U_im(x_i))\}\right)\}\phi .\hfill \end{array}$$
(16)
Now
$$\underset{i}{}x_i^2=1,$$
(17)
and commuting this relation with $`D`$,
$$\begin{array}{cc}\hfill 0& =[D,m(1)]\hfill \\ & =_i\{m(x_i)[D,m(x_i)]+[D,m(x_i)]m(x_i)]\hfill \\ & =2_i\{m(x_i)U_i+U_im(x_i)\}\hfill \end{array}$$
(18)
by (7). This shows in turn that
$$\underset{i}{}(U_i+am(x_i))^2=a^2+\underset{i}{}U_i^2$$
(19)
for any number $`a`$. On the other hand,
$$\begin{array}{cc}\hfill _i[U_i,m(x_i)]& =_iT_i\mathrm{cos}\rho _i=_i\mathrm{sin}^2\rho _i\hfill \\ & =_i(1\mathrm{cos}^2\rho _i)=1(n+1)=n.\hfill \end{array}$$
(20)
It is clear that we shall also need a simplification of $`_iU_i^2`$. This is provided by:
###### Lemma 2
$`_iT_i^2=\mathrm{\Delta }`$.
Proof: Consider one term of the sum on the left, and suppress the subscript $`i`$ for now, so that $`T_i=T`$, $`x_i=x`$. In abstract index notation, one term from the left side of the identity is
$$x^a_a(x^b_b)=(^ax)\{(^bx)_a+_a^bx\}_b.$$
(21)
Since $`^ax=(dx)^a`$, the first term is $`(dxdx)^{ab}_a_b`$. The second term involves
$$_a_bx=(\mathrm{Hess}x)_{ba}.$$
But the conformal Killing equation $`L_Tg=2xg`$ reads, in abstract index notation,
$$_aT_b+_bT_a=2xg_{ab}.$$
Since $`T_a=(dx)_a`$, this says that $`2`$Hess$`x=2xg`$, so that the quantity (21) becomes
$$(dxdx)^{ab}_a_b(^ax)x\delta _a{}_{}{}^{b}_{b}^{}=(dxdx)^{ab}_a_b(^ax)x_a.$$
Re-inserting the subscript $`i`$ and summing over it, the above quantity becomes
$$g^{ab}_a_b\frac{1}{2}(^a1)_a=g^{ab}_a_b=\mathrm{\Delta },$$
since $`_idx_idx_i`$ is the pullback of the ambient flat $`^{n+1}`$ metric, i.e. the round metric, and
$$\underset{i}{}x_i^ax_i=\frac{1}{2}^a\underset{1}{\underset{}{\left(\underset{i}{}x_i^2\right)}}.\mathrm{}$$
By the lemma and (19),
$$\left(D\frac{n(n2)}{4}\right)=\underset{i}{}U_i^2+\frac{n^2}{4},$$
so that
$$\underset{i}{}U_i^2=Dn/2.$$
Using all the identities just derived to evaluate (16), we get
$$\underset{i}{}M_iP_i\phi =\left(\lambda \frac{n}{2}+1(c^{})^2+nc^{}\right)\phi .$$
It is convenient to write this in terms of
$$\nu :=\sqrt{1+4\lambda },$$
so that
$$\lambda =\frac{\nu ^21}{4},c^{}=\frac{1+\nu }{2},$$
and
$$\underset{i}{}M_iP_i\phi =\frac{1}{2}(\nu +n1)(\nu +2)\phi .$$
(22)
This shows that some $`P_i\phi `$ is nonzero unless $`\nu =1n`$ or $`\nu =2`$; in particular it is nonzero for all positive values of $`\nu `$. But by (11), $`\nu n1`$.
The corresponding calculation with $`M_i`$ and $`P_i`$ in the other order begins with
$$\underset{i}{}P_iM_i\phi =\underset{i}{}\left(U_i+\frac{1+\sqrt{1+4\lambda ^{}}}{2}x_i\right)\left(U_i+\frac{1\sqrt{1+4\lambda }}{2}x_i\right)\phi ,$$
and (by virtue of (15)) produces a version of (16) with $`c^+`$ in place of $`c^{}`$, yielding
$$\underset{i}{}P_iM_i\phi =\left(\lambda \frac{n}{2}+1(c^+)^2+nc^+\right)\phi .$$
Since $`c^+=(\nu 1)/2`$, we get
$$\underset{i}{}P_iM_i\phi =\frac{1}{2}(\nu n+1)(\nu 2)\phi .$$
(23)
Thus some $`M_i\phi `$ is nonzero unless $`\nu =n1`$ or $`\nu =2`$. In the first case, $`\lambda =n(n2)/4`$; this was expected, since (11) and $`1E(n(n2)/4,D)`$ show that this is the bottom eigenvalue of $`D`$. Indeed, a look back at the formula for $`M_i\phi `$ in Proposition 1 shows that $`M_i1=0`$. In the second case, $`\nu =2`$, we have $`\lambda =3/4`$. By (11), this implies that $`n=3`$ and equality holds in (11); this is a special case of the situation just discussed. We have proved:
###### Proposition 3
For $`\lambda n(n2)/4`$, if $`0\phi E(\lambda ,D)`$, then $`P_i\phi `$ is a nonzero element of $`E(\lambda ^+,D)`$ for some $`i`$. If $`\lambda >n(n2)/4`$, then $`M_k\phi `$ is a nonzero element of $`E(\lambda ^{},D)`$ for some $`k`$.
The mechanism by which we have generated the eigenvalues (12) also rules out any other numbers occurring as eigenvalues. Indeed, suppose $`\lambda `$ is an eigenvalue of $`D`$ not on the list, so that (by (11)) there is some natural number $`j`$ with $`\lambda _j<\lambda <\lambda _{j+1}`$. Since the map $`\lambda \lambda ^{}`$ is strictly monotonic, we have $`\lambda _{j1}<\lambda ^{}<\lambda _j`$, $`\lambda _{j2}<\lambda ^{}<\lambda _{j1}`$ and so on, until we reach eigenvalues $`\mu (\lambda _0,\lambda _1)`$, $`\mu ^{}<\lambda _0`$, contradicting (11). The key ingredient, of course, is the assurance from Proposition 3 that each lower eigenspace in the induction is truly nonzero. We have:
###### Proposition 4
The $`\lambda _j`$ of (12) give the complete list of eigenvalues of $`D`$ on $`S^n`$. As a result, the $`j(n1+j)`$ for $`j=0,1,2,\mathrm{}`$ give the complete list of eigenvalues of $`\mathrm{\Delta }`$ on $`S^n`$.
We can also harvest the following corollary. Let $`E_j:=E(\lambda _j,D)`$.
###### Corollary 5
The span of the $`m(x_i)E_j`$ and $`Dm(x_i)E_j`$ is $`E_{j1}E_{j+1}`$ (where $`E_1=0`$).
Proof: The inclusion $``$ is immediate from the first line of (8) together with Proposition 1. For the inclusion $``$, (22,23) show that $`E_{j1}E_{j+1}`$ is contained in the sum of the $`m(x_i)E_j`$ and the
$$U_iE_jDm(x_i)E_j+m(x_i)DE_j=Dm(x_i)E_j+m(x_i)E_j.\mathrm{}$$
Note that we have not injected any information on spherical harmonics into our procedure for generating the eigenvalues. We may, however, get the interpretation of the eigenspaces as spaces of spherical harmonics as a consequence of what we have done:
###### Corollary 6
$`E_j`$ is exactly the set of restrictions from $`^{n+1}`$ to $`S^n`$ of $`j`$-homogeneous harmonic polynomials in the $`x_i`$.
Proof: First, the elements of $`E_j`$ are restrictions of $`j`$-homogeneous polynomials, since this is true of $`E_0`$, so follows from the previous corollary by induction on $`j`$. Second, if we compute the Laplacian of $`^{n+1}`$ in spherical coordinates, we get
$$\mathrm{\Delta }_{^{n+1}}=\frac{^2}{r^2}\frac{n1}{r}\frac{}{r}+\frac{1}{r^2}\mathrm{\Delta }_{S^n},$$
so that if $`\phi E_j`$, and we extend to $`^{n+1}`$ by extending the $`j`$-homogeneous polynomial formula,
$$\mathrm{\Delta }_{^{n+1}}\phi =[j^2(n1)j+j(n1+j)]\phi =0.$$
This calculation also shows that each $`j`$-homogeneous harmonic polynomial in $`^{n+1}`$ gives rise to an element of $`E_j.\mathrm{}`$
As a bonus result, we can give operators $`A_{2r}`$ which are functions of $`D`$ (or of $`\mathrm{\Delta }`$), which satisfy generalizations of the conformal covariance relation (4), namely
$$D(U_irm(x_i))=(U_i+rm(x_i))D,$$
(24)
for each $`r`$, each of which takes an eigenvalue on $`E_j:=E(\lambda _j,D)`$. That is, we can find functions $`f_r`$ on the spectrum of $`D`$ for which
$$A_{2r}|_{E_j}=f_r(\lambda _j)\mathrm{Id}_{E_j}.$$
In the usual notation of functional calculus, we write $`A_{2r}=f_r(D)`$.
Remark: Though we have shown that we have all eigenvalues for the conformal Laplacian, our derivation does not contain a proof that the span of the corresponding eigenfunctions is dense in $`L^2(S^n)`$. Of course, we have this by general elliptic theory. This sort of completeness is implicitly used later, when we describe other covariant operators, or intertwinors, as functions of a basic one (for example $`D`$, or the operator $`A_1`$ of (27), or the Dirac operator $`P`$ as it is used in Theorem 12.
First note that $`\lambda _j^\pm =\lambda _{j\pm 1}`$. Compressing (24), and letting $`\mu _j`$ be the putative eigenvalue for $`A_{2r}`$ on $`E_j`$, we have
$$\mu _{j\pm 1}(\frac{\lambda _{j\pm 1}\lambda _j}{2}r){}_{E_{j\pm 1}}{}^{}|m(x_i)|_{E_j}=\mu _j(\frac{\lambda _{j\pm 1}\lambda _j}{2}+r){}_{E_{j\pm 1}}{}^{}|m(x_i)|_{E_j}.$$
Since the $`\lambda _j`$ are known, we may generate the various $`\mu _j`$ inductively by demanding
$$\mu _{j\pm 1}\left(\lambda _{j\pm 1}\lambda _j2r\right)=\mu _j\left(\lambda _{j\pm 1}\lambda _j+2r\right).$$
(25)
Since a priori this gives two relations between adjacent $`\mu _j`$, we must check for consistency. Since $`\lambda _{j+1}\lambda _j=n+2j`$, we have this, provided we handle occurrences of vanishing $`\lambda _{j\pm 1}\lambda _j\pm ^{}2r`$ correctly. We emerge with a choice of $`\mu _j^{(2r)}`$ that is unique up to a constant (independent of $`j`$) factor:
###### Proposition 7
For fixed $`r\{n/2,n/21,\mathrm{}\}`$,
$$Z(r,j):=\frac{\mathrm{\Gamma }(n/2+j+r)}{\mathrm{\Gamma }(n/2+jr)}$$
(26)
is, up to a constant nonzero factor, the unique function of $`j\{0,1,2,\mathrm{}\}`$ that satisfies (25) and does not vanish identically. For $`r=n/2j_0`$ with $`j_0`$ a nonnegative integer, the residue of the above expression (viewed as a meromorphic function of $`r`$) at $`n/2j`$ is, up to a constant nonzero factor, the unique function of $`j\{0,1,2,\mathrm{}\}`$ that satisfies (25) and does not vanish identically.
If $`r\frac{1}{2}^+`$, then
$$Z(r,j)=(n/2+j+r1)\mathrm{}(n/2+jr).$$
If $`r\frac{1}{2}^+`$, then
$$Z(r,j)=\left(\underset{1p2r}{\overset{}{}}(n/2+jrp)\right)^1,$$
where $`{\displaystyle \stackrel{}{}}`$ is the product over nonzero factors.
The operator $`A_1`$ has eigenvalue $`Z(\frac{1}{2},j)=(n1)/2+j`$ on $`E_j`$; thus
$$A_1:=\sqrt{\mathrm{\Delta }+\left(\frac{n1}{2}\right)^2}.$$
(27)
This implies via (26) that
$$A_{2r}=\frac{\mathrm{\Gamma }(A_1+\frac{1}{2}+r)}{\mathrm{\Gamma }(A_1+\frac{1}{2}r)},r\{n/2,n/21,\mathrm{}\}.$$
Similarly, $`A_{2r}`$ (for $`r`$ outside the exceptional set given above) may be written as a function of $`A_{2q}`$ for any nonzero $`q`$ outside the exceptional set. In particular, for $`r^+`$ we get the sequence of differential operators
$$\underset{p=1}{\overset{r}{}}\left\{\mathrm{\Delta }+\left(\frac{n}{2}+p1\right)\left(\frac{n}{2}p\right)\right\},$$
also written down in , Remark 2.23.
## The entropy inequality
Equation (25), together with the sharp Hardy-Littlewood-Sobolev inequalities of Beckner on the sphere, give an argument for an optimal form of Gross’ entropy inequality on the sphere.
Let $`r`$ be the parameter of (25), and denote the $`r`$-derivative at $`r=0`$ by a prime. Differentiating (25) and normalizing so that the intertwinor $`A_0`$ is the identity, we have
$$\mu _{j\pm 1}^{}(\lambda _{j\pm 1}\lambda _j)2=\mu _j^{}(\lambda _{j\pm 1}\lambda _j)+2.$$
Thus
$$(\mu _{j\pm 1}\mu _j)^{}=\frac{4}{\lambda _{j\pm 1}\lambda _j}.$$
(28)
The $`j+1`$ and $`j1`$ relation lists are consistent, so the information in (28) is equivalent to
$$(\mu _{j+1}\mu _j)^{}=\frac{4}{\lambda _{j+1}\lambda _j}.$$
(29)
Since
$$\lambda _{j+1}\lambda _j=n+2j,$$
this gives (with $`m:=n/2`$)
$$\mu _j^{}=\mu _0^{}+\frac{2}{m}+\frac{2}{m+1}+\mathrm{}+\frac{2}{m+j1}.$$
This is in fact a formula for the eigenvalues of $`A_{2r}^{}`$.
It is convenient to pick a normalization of the series $`A_{2r}`$ with the property that $`\mu _0^{}=0`$. This is obtained by multiplying the spectral function (26) by $`\mathrm{\Gamma }(mr)/\mathrm{\Gamma }(m+r)`$ – this factor is independent of $`j`$, so for fixed $`r`$, gives a constant multiple of the $`A_{2r}`$ described by (26). With this normalization,
$$\mu _j^{}=\frac{2}{m}+\frac{2}{m+1}+\mathrm{}+\frac{2}{m+j1}.$$
Let us denote the intertwinors normalized in this way by $`B_{2r}`$. These operators also appear in Beckner’s sharp Hardy-Littlewood-Sobolev inequalities for $`r[0,n/2)`$: if $`F`$ is a positive smooth function on the sphere,
$$F^{(n2r)/2}B_{2r}F^{(n2r)/2}\left(F^n\right)^{(n2r)/n},$$
where all integrals are with respect to normalized measure. One has equality for every $`F`$ when $`r=0`$ (each side is $`F^n`$). But for $`r>0`$, equality holds exactly when $`F`$ is a constant multiple of a conformal (diffeomorphism) factor.
Now write the inequality as
$$0\left(F^n\right)^{(n2r)/n}+F^{(n2r)/2}B_{2r}F^{(n2r)/2},$$
and take $`(d/dr)|_{r=0}`$ of each side. (Of course this differentiation of the inequality depends on the fact that equality holds for every $`F`$ at $`r=0`$.) This gives
$$0\frac{2}{n}\left(F^n\right)\mathrm{log}F^n2F^n\mathrm{log}F+F^{n/2}B_{2r}^{}F^{n/2},$$
or in a slightly better form,
$$2F^n\mathrm{log}F\frac{2}{n}\left(F^n\right)\mathrm{log}F^n+F^{n/2}B_{2r}^{}F^{n/2}.$$
What we know immediately about the case of equality is that it includes at least constant multiples of the conformal factors $`F`$. But since the quantities in play are analytic in $`r`$ near $`r=0`$, we get precisely these functions.
With $`f:=F^{n/2}`$, we can rewrite as
$$\frac{4}{n}f^2\mathrm{log}f\frac{2}{n}\left(f^2\right)\mathrm{log}f^2+fA_{2r}^{}f,$$
for
$$\mathrm{eig}(A_{2r}^{},E_j)=\frac{2}{m}+\frac{2}{m+1}+\mathrm{}+\frac{2}{m+j1}.$$
It is easily verified that these eigenvalues are $``$ (but very close to) those of $`2\mathrm{log}(2A_1/(n1))`$, where $`A_1`$ is as in (27). (The harmonic sum is a certain Riemann sum for the integral defining the log.) Thus we may write
$$\frac{2}{n}f^2\mathrm{log}f\frac{1}{n}\left(f^2\right)\mathrm{log}f^2+f\left(\mathrm{log}\frac{2A_1}{n1}\right)f,$$
(30)
giving away a little sharpness. To summarize:
###### Theorem 8
For smooth positive $`f`$ on the sphere $`S^n`$, in normalized measure,
$$\frac{4}{n}f^2\mathrm{log}f\frac{2}{n}\left(f^2\right)\mathrm{log}f^2+fA_{2r}^{}f,$$
(31)
where $`H`$ takes the eigenvalue
$$\frac{2}{m}+\frac{2}{m+1}+\mathrm{}+\frac{2}{m+j1}$$
on $`j^{\underset{¯}{\mathrm{th}}}`$ order spherical harmonics. In particular the weaker statement (30) holds. Equality holds in (31) if and only if $`f^{2/n}`$ is a positive constant multiple of a conformal (diffeomorphism) factor.
## Spinor operators
Let $`P`$ be the Dirac operator on the spinor bundle $`\mathrm{\Sigma }`$. By the Lichnerowicz formula,
$$P^2=^{}+\frac{1}{4}K,$$
so that
$$P^2=^{}+\frac{1}{4}n(n1)\text{ on }S^n.$$
As a result, if $`\lambda `$ is an eigenvalue of $`P`$ on the sphere, then $`\lambda `$ is real with
$$\lambda ^2n(n1)/4.$$
(32)
The analogue of (7) is
$$[P^2,m(\omega )]=[^{},m(\omega )]=2_T+nm(\omega ).$$
(33)
The conformal covariance relation satisfied by $`P`$ is
$$P\left(L_T+\frac{n1}{2}\omega \right)=\left(L_T+\frac{n+1}{2}\omega \right)P,$$
(34)
on general pseudo-Riemannian spin manifolds, where $`L_Tg=2\omega g`$. If $`X`$ is an arbitrary smooth vector field, shows that the Lie and covariant derivatives on spinors are related by
$$L_X_X=\frac{1}{8}(dX)_{ab}\gamma ^a\gamma ^b,$$
where $`dX`$ is the exterior derivative of the 1-form corresponding to $`X`$ under the metric, and $`\gamma `$ is the fundamental tensor-spinor (a section of $`TM\mathrm{End}(\mathrm{\Sigma })`$). This assumes that the internal conformal weight 0 has been assigned to the spinor bundle; assigning internal weight $`\pm \frac{1}{2}`$ as in results in an extra term involving div$`X`$. Our proper conformal vector fields $`T_i`$ on $`S^n`$ have the $`dx_i`$ as their metric correspondents, so we may specialize (33) to
$$[P^2,m(x_i)]=2L(T_i)+nm(x_i)=:2U_i.$$
(35)
The conformal covariance relation (34) specializes to
$$P(U_i\frac{1}{2}m(x_i))=(U_i+\frac{1}{2}m(x_i))P.$$
(36)
Compressing the operators in (35) to act between two eigenspaces for $`P`$, $`E(\lambda ,P)E(\mu ,P)`$, we have
$${}_{\mu }{}^{}|U_i|_\lambda =\frac{\mu ^2\lambda ^2}{2}{}_{\mu }{}^{}|m(x_i)|_\lambda ,$$
(37)
after which (36) implies that
$$\mu (\mu ^2\lambda ^21){}_{\mu }{}^{}|m(x_i)|_\lambda =\lambda (\mu ^2\lambda ^2+1){}_{\mu }{}^{}|m(x_i)|_\lambda .$$
This is implied by
$$\mu (\mu ^2\lambda ^21)=\lambda (\mu ^2\lambda ^2+1),$$
which may be rewritten
$$(\mu +\lambda )(\mu \lambda +1)(\mu \lambda 1)=0.$$
(38)
(Note that all these equations are symmetric in $`\mu `$ and $`\lambda `$.)
The cubic equation (38) suggests the possible nonvanishing of three “adjacent” eigenspaces $`E(\lambda ,P)`$, $`E(\lambda +1,P)`$, and $`E(\lambda 1,P)`$, given $`E(\lambda ,P)0`$. Let $`0\psi E(\lambda ,P)`$, and write the Dirac operator as $`\gamma ^a_a`$ (using abstract index notation, in which repetition of an index, once up and once down, denotes a contraction). To avoid excessive super- and subscripting, denote $`m(x_i)`$ and $`U_i`$ by $`x`$ and $`U`$ for now. In analogy with (13,14), we have:
$$\begin{array}{cc}\hfill P(x\psi )& =\gamma ^a_a(x\psi )=\lambda x\psi +\gamma ^a(_ax)\psi ,\hfill \\ \hfill P(U\psi )& =P(\frac{1}{2}x\psi )+(U+\frac{1}{2}x)\lambda \psi =\lambda x\psi +\lambda U\psi +\frac{1}{2}\gamma ^a(_ax)\psi .\hfill \end{array}$$
In contrast to (13,14), however, this system doesn’t close. We need in addition:
$$P(\gamma ^a(_ax)\psi )=\gamma ^b_b(\gamma ^a(_ax)\psi )=\gamma ^b\gamma ^a((_ax)_b\psi +(_b_ax)\psi ),$$
where we have used the relation $`\gamma =0`$ between the spin connection and the fundamental tensor-spinor. By the Clifford relation
$$\gamma ^a\gamma ^b+\gamma ^b\gamma ^a=2g^{ab},$$
we may rewrite the above as
$$P(\gamma ^a(_ax)\psi )=\gamma ^a(_ax)\lambda \psi 2(^ax)_a\psi (^a_ax)\psi =2U\psi \lambda \gamma ^a(_ax)\psi ,$$
where we have used (6) to simplify $`^a_ax=\mathrm{\Delta }x`$.
Abbreviating
$$[P,m(x)]=\gamma ^a(_ax)=:y,$$
(39)
we have
$$\begin{array}{cc}\hfill A\psi & :=(U+\lambda x+\frac{1}{2}y)\psi E(\lambda +1,P),\hfill \\ \hfill S\psi & :=(U\lambda x\frac{1}{2}y)\psi E(\lambda 1,P),\hfill \\ \hfill N\psi & :=(U\frac{1}{2}x\lambda y)\psi E(\lambda ,P).\hfill \end{array}$$
(40)
Recall that we have suppressed the subscript $`i\{0,\mathrm{},n\}`$, so that we really have operators $`A_i`$, $`S_i`$, and $`N_i`$. To find whether we really get something nonzero by the above processes, we compute
$$\begin{array}{c}_iS_iA_i\psi =_i(U_i(\lambda +1)x_i\frac{1}{2}y_i)(U_i+\lambda x_i+\frac{1}{2}y_i)\psi ,\hfill \\ _iA_iS_i\psi =_i(U_i+(\lambda 1)x_i+\frac{1}{2}y_i)(U_i\lambda x_i\frac{1}{2}y_i)\psi ,\hfill \\ _iN_iN_i\psi =_i(U_i\frac{1}{2}x_i+\lambda y_i)(U_i\frac{1}{2}x_i\lambda y_i)\psi .\hfill \end{array}$$
(41)
In simplifying this, note that
$$yy=\gamma ^a\gamma ^b(_ax)_bx=(_ax)^ax=|dx|^2.$$
Thus
$$\underset{i}{}y_i^2=\underset{i}{}|dx_i|^2=\underset{i}{}\mathrm{sin}^2\rho _i=\underset{i}{}(1\mathrm{cos}^2\rho _i)=1(n+1)=n.$$
We also have
$$yx=xy=x\gamma ^a(_ax)=\frac{1}{2}\gamma ^a_a(xx),$$
so that
$$\underset{i}{}y_ix_i=\underset{i}{}x_iy_i=\frac{1}{2}\gamma ^a_a\left(\underset{i}{}x_i^2\right)=\frac{1}{2}\gamma ^a_a1=0.$$
(18), (19), and (20) are still valid (with the new meaning of $`U_i`$, and $`P^2`$ in place of $`D`$ in the intermediate steps of (18)).
In addition,
$$\begin{array}{c}Uy=U[P,x]=\frac{1}{2}[P^2,x][P,x]=\frac{1}{2}(P[P,x]+[P,x]P)[P,x]=Py^2+yPy,\hfill \\ yU=[P,x]U=\frac{1}{2}[P,x][P^2,x]=\frac{1}{2}[P,x](P[P,x]+[P,x]P)=yPy+y^2P,\hfill \end{array}$$
so that
$$[U,y]=[P,y^2],$$
and
$$\underset{i}{}[U_i,y_i]=[P,\underset{i}{}y_i^2]=[P,n]=0.$$
(42)
The argument of Lemma 2 goes through formally as written (with $``$ now involving the spin connection), and gives
$$\underset{i}{}_{T_i}^2=^{}=P^2+\frac{n(n1)}{4}.$$
With the discussion preceding (35), this gives
$$\underset{i}{}\left(U_i\frac{n}{2}x_i\right)^2=P^2+\frac{n(n1)}{4}.$$
By the analogues of (18,19,20),
$$\begin{array}{cc}\hfill _i(U_i+ax_i)^2& =a^2+_iU_i^2\hfill \\ & =a^2P^2+\frac{(n(n1)}{4}\frac{n^2}{4}\hfill \\ & =a^2P^2\frac{n}{4}.\hfill \end{array}$$
In particular,
$$\underset{i}{}U_i^2=P^2\frac{n}{4}.$$
Using all these identities, for arbitrary $`a`$ and $`b`$,
$$\underset{i}{}(U_i(a+1)x_iby_i)(U_i+ax_i+by_i)=P^2\frac{3n}{4}a^2a(n+1)+b^2n.$$
In the first line of (41), we take $`a=\lambda `$ and $`b=\frac{1}{2}`$, and apply to a $`\lambda `$-eigenspinor $`\psi `$ to get
$$\underset{i}{}S_iA_i\psi =\left(2\lambda ^2\frac{n}{2}(n+1)\lambda \right)\psi =2\left(\lambda +\frac{n}{2}\right)\left(\lambda +\frac{1}{2}\right)\psi .$$
In the second line, take $`a=\lambda `$ and $`b=\frac{1}{2}`$ to get
$$\underset{i}{}A_iS_i\psi =\left(2\lambda ^2\frac{n}{2}+(n+1)\lambda \right)\psi =2\left(\lambda \frac{n}{2}\right)\left(\lambda \frac{1}{2}\right)\psi .$$
In the third line, take $`a=\frac{1}{2}`$ and $`b=\lambda `$ to get
$$\underset{i}{}N_iN_i\psi =(n1)\left(\lambda +\frac{1}{2}\right)\left(\lambda \frac{1}{2}\right)\psi .$$
This establishes:
###### Proposition 9
If $`P`$ has a nonzero eigenspace on $`S^n`$, then the nonzero eigenspaces of $`P`$ on round $`S^n`$ are exactly the
$$F_j:=E(n/2+j,P),G_j:=E((n/2+j),P)$$
for $`j=0,1,2,\mathrm{}`$.
Proof: The last three displayed identities and (32) show that if $`E(\lambda ,P)0`$, then
$$\begin{array}{c}E(\lambda ,P)0,\hfill \\ E(\lambda 1,P)0\mathrm{unless}\lambda =n/2,\hfill \\ E(\lambda +1,P)0\mathrm{unless}\lambda =n/2.\hfill \end{array}$$
Thus if there is an eigenvalue outside the set of $`\pm (n/2+j)`$, we may generate an eigenvalue contradicting (32). Given any nonzero eigenspace $`E(\lambda ,P)`$ with $`\lambda `$ of the form $`\pm (n/2+j)`$ however, we may generate nonzero eigenspaces corresponding to all such $`\lambda .\mathrm{}`$
The beginning of the statement of the last proposition is a bit awkward; we need to assume there is some eigenspace in order to get a foothold analogous to that provided by the function 1 in the scalar case. We can get this foothold by taking a nonzero parallel spinor in $`^n`$ and stereographically injecting it (with the proper conformal weight) to $`S^n`$. Alternatively, we could do just enough elementary elliptic theory to conclude that the minimizer for the Rayleigh quotient based on $`P^2`$ provides us with an eigensection.
The analogue of Corollary 5 is
###### Corollary 10
The span of the $`m(x_i)E(\lambda ,P)`$, $`Pm(x_i)E(\lambda ,P)`$, and $`P^2m(x_i)E(\lambda ,P)`$ is $`E(\lambda +1,P)E(\lambda 1,P)E(\lambda ,P)`$.
Proof: The inclusion $``$ follows from (40), (37), and the fact (from (39), in the notation of that display) that $`{}_{\mu }{}^{}|y|_\lambda =(\mu \lambda ){}_{\mu }{}^{}|m(x)|_\lambda `$. For the inclusion $``$, note first that (41) puts $`E(\lambda +1,P)E(\lambda 1,P)E(\lambda ,P)`$ in the span of the $`A_iE(\lambda ,P)`$, $`N_iE(\lambda ,P)`$, and $`S_iE(\lambda ,P)`$. By (41), this is in the span of the $`x_iE(\lambda ,P)`$, $`U_iE(\lambda ,P)`$, and $`y_iE(\lambda ,P)`$. By (35) and (39), this is in the span of the $`m(x_i)E(\lambda ,P)`$, $`Pm(x_i)E(\lambda ,P)`$, and $`P^2m(x_i)E(\lambda ,P).\mathrm{}`$
In analogy with Proposition 7, we may seek intertwinors $`𝒜_{2k+1}`$ for complex-valued $`k`$, satisfying the intertwining relation
$$𝒜_{2k+1}(U(k+\frac{1}{2})x)=(U+(k+\frac{1}{2})x)𝒜_{2k+1}$$
(43)
(which extends (36)). In fact, such intertwinors exist and commute with $`P`$. The subfamily of these for which $`k`$ is a nonnegative integer yields a family of odd-order (in fact, order $`2k+1`$) differential operators which are polynomial in $`P`$:
###### Theorem 11
For $`k=0,1,2,\mathrm{}`$, and
$$𝒫_{2k+1}=(Pk)(Pk+1)\mathrm{}P\mathrm{}(P+k1)(P+k)=P(P^21)(P^24)\mathrm{}(P^2k^2),$$
we have
$$𝒫_{2k+1}(U(k+\frac{1}{2})x)=(U+(k+\frac{1}{2})x)𝒫_{2k+1},$$
where $`U`$ is any $`U_i`$, and $`x`$ is any $`x_i`$.
Proof: Fix $`k`$, and consider
$$\alpha (\lambda ):=(\lambda k)(\lambda k+1)\mathrm{}\lambda \mathrm{}(\lambda +k1)(\lambda +k),$$
the eigenvalue taken by $`𝒫_{2k+1}`$ on $`E(\lambda ,P)`$. In view of (37), it is enough to show that
$$\begin{array}{c}(\mu k)(\mu k+1)\mathrm{}\lambda \mathrm{}(\mu +k1)(\mu +k)\left(\frac{\mu ^2\lambda ^2}{2}(k+\frac{1}{2})\right)\stackrel{\mathrm{?}}{=}\hfill \\ (\lambda k)(\lambda k+1)\mathrm{}\lambda \mathrm{}(\lambda +k1)(\lambda +k)\left(\frac{\mu ^2\lambda ^2}{2}+(k+\frac{1}{2})\right)\hfill \end{array}$$
(44)
for $`\mu =\lambda \pm 1`$ and for $`\mu =\lambda `$. This is easily verified in each of the three cases.$`\mathrm{}`$
For arbitrary order $`k+1/2`$, we might expect a nonlocal intertwining operator. This is provided by:
###### Theorem 12
If $`k+n/2`$, the operators
$$𝒜_{2k+1}:=\mathrm{sgn}(P)^{n+1}\frac{\mathrm{\Gamma }(P+k+1)}{\mathrm{\Gamma }(Pk)}$$
(45)
satisfy the intertwining relations (43), where for each function $`f`$, the operator $`f(P)`$ takes the eigenvalue $`f(\lambda )`$ on $`E(\lambda ,P)`$.
Proof: Fix $`k`$ and let $`\alpha [\lambda ]`$ be the eigenvalue on $`E(\lambda ,P)`$ of a putative intertwinor $`𝒜_{2k+1}`$. It is immediate from (37) that $`\alpha [\lambda ]=\alpha [\lambda ]`$, and that
$$\alpha [\lambda +1](\lambda k)=(\lambda +1+k)\alpha [\lambda ].$$
(46)
Moreover, these equations for all $`P`$-eigenvalues $`\lambda `$ are sufficient for the intertwining property (the condition relating $`\alpha [\lambda 1]`$ to $`\alpha [\lambda ]`$ is the same as (46), with $`\lambda `$ shifted to $`\lambda 1`$.) Implementing the recursion this gives, we get the intertwinor
$$\mathrm{sgn}(P)\frac{\mathrm{\Gamma }(|P|+k+1)}{\mathrm{\Gamma }(|P|k)},$$
(47)
and one may check directly that this is intertwining. Using the identity
$$\mathrm{\Gamma }(z)\mathrm{\Gamma }(z)z\mathrm{sin}(\pi z)=\pi ,$$
we get
$$\frac{\mathrm{\Gamma }(P+k+1)}{\mathrm{\Gamma }(Pk)}=\frac{\mathrm{\Gamma }(P+k+1)}{\mathrm{\Gamma }(Pk)}\frac{\mathrm{sin}\pi (P+k)}{\mathrm{sin}\pi (Pk1)}.$$
The trigonometric factor in this is
$$\frac{\mathrm{sin}(\pi P)\mathrm{cos}(\pi k)+\mathrm{cos}(\pi P)\mathrm{sin}(\pi k)}{\mathrm{sin}(\pi (P1))\mathrm{cos}(\pi k)\mathrm{cos}(\pi (P1))\mathrm{sin}(\pi k)}.$$
If $`n`$ is even, $`k`$ and $`P`$ takes integral eigenvalues, so this becomes $`1`$. If $`n`$ is odd, $`P`$ takes properly half-integral eigenvalues and $`k`$ is not a proper half-integer, so the above becomes $`1`$. This shows that
$$\frac{\mathrm{\Gamma }(|P|+k+1)}{\mathrm{\Gamma }(|P|k)}=\mathrm{sgn}(P)^n\frac{\mathrm{\Gamma }(P+k+1)}{\mathrm{\Gamma }(Pk)},$$
so that the intertwinor (47) is the $`𝒜_{2k+1}`$ of (45).$`\mathrm{}`$
Remark: When $`n`$ is even, we may set $`k=1/2`$ in the last theorem to obtain the intertwinor
$$\mathrm{sgn}(P)\frac{\mathrm{\Gamma }(P+\frac{3}{2})}{\mathrm{\Gamma }(P\frac{1}{2})}=\mathrm{sgn}(P)(P+\frac{1}{2})(P\frac{1}{2})=P|P|\frac{1}{4}P/|P|.$$
(48)
This is also a formula for an intertwinor when $`n`$ is odd, though this is not covered by Theorem 12. To see this, we just feed the formula for the operator (48) into a calculation similar to (44). Similar considerations hold for the other values of $`k`$ excluded by the hypotheses of Theorem 12. Of course it is known that even and odd dimensions act very differently in many ways with respect to the spinor bundle. For example, in even dimensions the bundle $`\mathrm{\Sigma }`$ itself admits a chirality decomposition into two subbundles. In contrast, $`\mathrm{\Sigma }`$ is irreducible in odd dimensions, but each eigenspace of $`P^2`$ admits a chiral decomposition. In particular, the relation between $`E(\lambda ,P)`$ and $`E(\lambda ,P)`$ has very different interpretations in the two dimension parities. The technique we use here seems remarkably attuned to deriving Theorem 12 using only things that are common to the two dimension parities, the power on $`\mathrm{sgn}(P)`$ being the only hint of the differences.
Remark: Our technique can also be used to obtain the eigenvalue specta of $`D`$, $`P`$ and other intertwinors on the non-Riemannian hyperboloids ; i.e. on other real forms of the sphere, again manifolds with large Lie algebras of conformal vector fields. This provides some explanation, in terms of conformal geometry, of the remarkable agreement between most of the eigenvalues of the Laplace operator on the various real forms of the sphere.
Thomas Branson: Department of Mathematics, University of Iowa, Iowa City IA 52242 USA thomas-branson@uiowa.edu
Bent Ørsted: Institut for Matematiske Fag, Aarhus Universitet, 8000 Aarhus C, Denmark orsted@imf.au.dk
|
warning/0506/math-ph0506052.html
|
ar5iv
|
text
|
# Lieb-Thirring type inequalities and Gagliardo-Nirenberg inequalities for systems
## 1 Introduction
Lieb-Thirring type inequalities are well known in the context of the stability of matter in quantum mechanics. Let $`\mathrm{}=h/2\pi >0`$ and $`m>0`$ be respectively Planck’s constant and the mass constant. Given a smooth bounded nonpositive potential $`V`$ on $`^d`$, if we denote by
$$\lambda _1(V)<\lambda _2(V)\lambda _3(V)\mathrm{}\lambda _N(V)<0$$
the finite sequence of all negative eigenvalues of the Schrödinger operator
$$H_V=\frac{\mathrm{}^2}{2m}\mathrm{\Delta }+V,$$
then it is possible to bound the sum $`_{i=1}^N|\lambda _i(V)|^\gamma `$ in terms of $`V_{L^{\gamma +d/2}(^d)}`$ whatever $`N`$ is. The inequality
$$\underset{i=1}{\overset{N}{}}|\lambda _i(V)|^\gamma C_{\mathrm{LT}}(\gamma )_^d|V|^{\gamma +\frac{d}{2}}𝑑x$$
(1)
is known as the Lieb-Thirring inequality. Here we denote by $`C_{\mathrm{LT}}(\gamma )`$ the smallest possible positive constant which is independent of $`V`$. For $`\gamma =1`$, the sum $`_{i=1}^N|\lambda _i(V)|`$ is the complete ionization energy, which is the physically relevant quantity for studying the stability of matter. Considerable efforts have been made to understand further the Lieb-Thirring inequality and in particular to find the optimal value of $`C_{\mathrm{LT}}(\gamma )`$. Up to now a few facts about the sharp constant in the Lieb-Thirring inequality are known. It was proved in for $`d=1`$ and later generalized to arbitrary $`d`$ in that for $`\gamma 3/2`$ the sharp constant is given by the semiclassical constant, i.e. the constant corresponding to the limit problem when letting $`\mathrm{}0`$, after an appropriate scaling. Among many other open problems, the Lieb-Thirring conjecture asserts that in $`d=1`$
$$C_{\mathrm{LT}}(\gamma )=C_{\mathrm{LT}}^{(1)}(\gamma ):=\underset{\begin{array}{c}V𝒟()\\ V0\end{array}}{inf}\frac{|\lambda _1(V)|^\gamma }{_{}|V|^{\gamma +\frac{1}{2}}𝑑x}.$$
This has been worked out for the case $`\gamma =1/2`$ in . Also see for further results on (1).
What we study in this note is a somewhat different problem, where $`V`$ is a nonnegative, unbounded potential on $`^d`$, such that the eigenvalues of $`H_V`$ form a positive unbounded nondecreasing sequence $`(\lambda _i(V))_i^{}`$. Our main result is the
###### Theorem 1
For any $`\gamma >d/2`$, $`d^{}`$, and for any nonnegative $`VC^{\mathrm{}}(^d)`$ such that $`V^{d/2\gamma }L^1(^d)`$,
$$\underset{i^{}}{}\left[\lambda _i(V)\right]^\gamma 𝒞(\gamma )_^dV^{\frac{d}{2}\gamma }𝑑x.$$
(2)
The value of the sharp constant $`𝒞(\gamma )`$ is given by the Weyl asymptotics, i.e. by its value in the semiclassical limit:
$$𝒞(\gamma )=\left(\frac{m}{2\pi \mathrm{}^2}\right)^{d/2}\frac{\mathrm{\Gamma }(\gamma d/2)}{\mathrm{\Gamma }(\gamma )}.$$
Although the result is quite simple and arises as an immediate consequence of the Golden and Thompson inequality, it is to our knowledge new. Some partial estimates have been obtained in in case of quadratic potentials. We refer to for earlier results in this direction. The semiclassical formula stems from prescribing $`h^d`$ phase space volume to each bound state of the Schrödinger operator. Using this heuristics, we can estimate the series by
$$\underset{i^{}}{}\left[\lambda _i(V)\right]^\gamma \frac{1}{h^d}_{^d\times ^d}\left(\frac{p^2}{2m}+V(x)\right)^\gamma 𝑑p𝑑x,$$
which is easily seen to yield the right hand side of (2). If one considers $`h`$ as a parameter, then as $`h`$ gets small the two sides in the above relation are asymptotically the same. A rigorous proof of this fact relies on the Weyl asymptotics, that we will establish later in this paper for a specific potential $`V`$.
We may notice that all physical constants can be adimensionalized. A simple scaling indeed shows that for any $`i^{}`$,
$$\lambda _i^{\mathrm{},m}(V)=\lambda _i^{1,\mathrm{\hspace{0.17em}1}/2}\left(V({\scriptscriptstyle \frac{\mathrm{}}{\sqrt{2m}}})\right),$$
so that in (2),
$$𝒞(\gamma )_{|\mathrm{},m}=\left(\frac{2m}{\mathrm{}^2}\right)^{d/2}𝒞(\gamma )_{|\mathrm{}=1,m=1/2},𝒞(\gamma )_{|\mathrm{}=1,m=1/2}=\frac{\mathrm{\Gamma }\left(\gamma \frac{d}{2}\right)}{(4\pi )^{d/2}\mathrm{\Gamma }\left(\gamma \right)}.$$
From now on, we assume for simplicity that $`\mathrm{}=1,m=1/2`$.
Theorem 1 is motivated by the study of the dynamical stability of mixed states with respect to minimizers of variational problems with temperature in quantum mechanics. Inequality (1) appears in the context of atomic and molecular physics, where it is natural to consider isolated systems for which the potential $`V`$ is asymptotically zero at infinity. Computing the full ionization energy is then a completely relevant question. Requiring that $`V`$ grows at infinity makes sense in a different context, e.g., in solid state physics, where the potential is not necessarily created by the system under consideration itself, but can be imposed by external devices (for instance a doping profile) or by a given electrostatic field (applied voltage). In that case, collective effects are fundamental and it is interesting to investigate how mixed states converge in a semi-classical limit to a classical system. At the kinetic level, the behavior of the classical system is now reasonably well understood. For instance one knows in which sense special stationary solutions are stable, see . At the quantum level, many particle systems are not so well understood. A first attempt in this direction has been made in , in a nonlinear case, but the result relies on a rather weak notion of stability and the exchange term is neglected. For a linear system, we will see in Section 3 that an appropriate functional for studying the stability of a mixed state, i.e. a sequence $`(𝝂,𝝍)=(\nu _i,\psi _i)_i^{}_+^{^{}}\times (L^2(^d))^{^{}}`$ made of nonnegative ordered occupation numbers $`\nu _i`$ and wave functions $`\psi _i`$, is the free energy functional
$$[𝝂,𝝍]:=\underset{i^{}}{}\left[\beta (\nu _i)+\nu _i_^d\left(|\psi _i|^2+V|\psi _i|^2\right)𝑑x\right],$$
where $`\beta `$ is a given convex function on $`_+`$. Under the constraints
$$(\psi _i,\psi _j)_{L^2(^d)}=\delta _{ij}i,j^{},$$
the functional $``$ has a minimizer made of the sequence $`\overline{𝝍}=(\overline{\psi }_i)_i^{}`$ of the eigenfunctions counted with multiplicity, and the sequence $`\overline{𝝂}=(\overline{\nu }_i)_i^{}`$ of occupation numbers given in terms of the eigenvalues by
$$\overline{\nu }_i=(\beta ^{})^1(\lambda _i(V)).$$
However, such considerations are purely formal as long as we dont prove that $`[𝝂,𝝍]`$ is finite at least for the formal minimizer $`(𝝂,𝝍)=(\overline{𝝂},\overline{𝝍})`$. Such a property is a condition on both $`\beta `$ and $`V`$. In case
$$\beta (\nu )=\{\begin{array}{cc}(1m)^{m1}m^m\nu ^m\hfill & \text{if}\nu 0,\hfill \\ & \\ +\mathrm{}\hfill & \text{if}\nu <0,\hfill \end{array}\text{and}m(0,1),$$
(3)
and with $`\gamma =\frac{m}{1m}`$, for any $`i^{}`$, we obtain
$`\overline{\nu }_i={\displaystyle \frac{m}{1m}}\left(\lambda _i(V)\right)^{\frac{1}{m1}},`$
$`\beta (\overline{\nu }_i)+\overline{\nu }_i\lambda _i(V)=\left(\beta (\beta ^{})^1\right)(\lambda _i(V))+(\beta ^{})^1(\lambda _i(V))\lambda _i(V)=\left[\lambda _i(V)\right]^\gamma .`$
The free energy is well defined at least for the optimal mixed state if the series $`_i^{}[\beta (\overline{\nu }_i)+\overline{\nu }_i\lambda _i(V)]`$ converges, which amounts to require that $`_i^{}\left[\lambda _i(V)\right]^\gamma `$ is finite. A sufficient condition is therefore given by Theorem 1.
Section 2 is devoted to a proof of Theorem 1 based on the inequality of Golden and Thomson (Theorem 2). We also state a more general result in Theorem 3. The notion of dynamic stability will be explained in Section 3 and illustrated by several examples. In Section 3, we will come back to the constant $`C_{\mathrm{LT}}^{(1)}(\gamma )`$ which appears in the Lieb-Thirring conjecture and prove that it is related to the best constant in some special Gagliardo-Nirenberg inequalities in the standard case corresponding to $`V0`$. Such a result is not new, but we also prove that a similar result holds in the case $`V0`$ (case of Theorem 1), which is apparently new. We also relate a limiting case to the euclidean logarithmic Sobolev inequality. In Section 5 we show in which sense Lieb-Thirring type inequalities can be seen as generalizations of the Gagliardo-Nirenberg inequalities to systems. This extends to systems of orthonormal functions what has been observed in Sections 4.1 and 4.2. We formulate optimal inequalities in an abstract framework and apply the result to the standard case (Corollary 16), to the framework of Theorem 1 (Corollary 17) and to a limiting case which provides an optimal inequality of logarithmic Sobolev type for systems. Optimal constants are expressed in terms of the optimal constants for Lieb-Thirring type inequalities.
## 2 Proof of Theorem 1
The proof of Theorem 1 is straightforward. It relies on a change of variables in the definition of the $`\mathrm{\Gamma }`$ function and the following inequality due to Golden and Thompson . See for an introduction to such methods and a proof based on the Feynman-Kac formula. Here we adopt the presentation of as stated in , Theorem 9.2, p. 94.
###### Theorem 2
Let $`V`$ be in $`L_{\mathrm{loc}}^1(^d)`$ and bounded from below. Assume moreover that $`e^{tV}`$ is in $`L^1(^d)`$ for any $`t>0`$. Then
$$\mathrm{Tr}\left(e^{t(\mathrm{\Delta }+V)}\right)(4\pi t)^{\frac{d}{2}}_^de^{tV(x)}𝑑x.$$
(4)
Proof. For completeness, we give an elementary proof of this result. We do not claim originality here and we give this result only for the convenience of the reader.
Consider the Green function $`G`$ of the heat equation:
$$G(x,t):=(4\pi t)^{\frac{d}{2}}e^{\frac{|x|^2}{4t}}.$$
We will then write
$$u(,t)=e^{t\mathrm{\Delta }}f:=G(,t)f$$
if $`u`$ is a solution of $`u_t=\mathrm{\Delta }u`$ with initial data $`u(,t=0)=f`$. By Trotter’s formula, $`e^{t(\mathrm{\Delta }+V)}`$ is obtained as the strong limit of
$$\left(e^{\frac{t}{n}\mathrm{\Delta }}e^{\frac{t}{n}V}\right)^n$$
as $`n`$ goes to infinity. Then we compute the trace of this last quantity as
$`{\displaystyle _{(^d)^n}}𝑑x𝑑x_1𝑑x_2\mathrm{}𝑑x_nG({\displaystyle \frac{t}{n}},xx_1)e^{\frac{t}{n}V(x_1)}G({\displaystyle \frac{t}{n}},x_1x_2)e^{\frac{t}{n}V(x_2)}\mathrm{}`$
$`\mathrm{}G({\displaystyle \frac{t}{n}},x_nx)e^{\frac{t}{n}V(x)}.`$
With the notation $`x=x_0=x_{n+1}`$, we rewrite this as
$$_{(^d)^n}𝑑x_0𝑑x_1𝑑x_2\mathrm{}𝑑x_n\underset{j=0}{\overset{n}{}}G(\frac{t}{n},x_jx_{j+1})e^{\frac{t}{n}_{k=0}^{n1}V(x_k)}.$$
Using the convexity of $`xe^x`$, we estimate the exponential term by:
$$e^{\frac{t}{n}_{k=0}^{n1}V(x_k)}\frac{1}{n}\underset{k=0}{\overset{n1}{}}e^{tV(x_k)}.$$
This amounts to
$`\mathrm{Tr}\left(e^{\frac{t}{n}\mathrm{\Delta }}e^{\frac{t}{n}V}\right)^n`$ $``$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle _{(^d)^n}}𝑑x_0𝑑x_1𝑑x_2\mathrm{}𝑑x_n{\displaystyle \underset{j=0}{\overset{n}{}}}G({\displaystyle \frac{t}{n}},x_jx_{j+1})e^{tV(x_k)}`$
$`=(4\pi t)^{\frac{d}{2}}{\displaystyle _{(^d)^2}}e^{tV(x)}𝑑x`$
using
$$_{(^d)^{n1}}𝑑x_0𝑑x_1𝑑x_2\mathrm{}𝑑x_{k1}𝑑x_{k+1}\mathrm{}𝑑x_n\underset{j=0}{\overset{n}{}}G(\frac{t}{n},x_jx_{j+1})=G(t,x_kx_k)=(4\pi t)^{\frac{d}{2}}.$$
$`\mathrm{}`$
Proof of Theorem 1. The definition of the $`\mathrm{\Gamma }`$ function gives, for any $`\gamma >0`$ and $`\lambda >0`$,
$$\lambda ^\gamma =\frac{1}{\mathrm{\Gamma }(\gamma )}_0^+\mathrm{}e^{t\lambda }t^{\gamma 1}𝑑t.$$
The operator $`\mathrm{\Delta }+V`$ is essentially self-adjoint on $`L^2(^d)`$, and positive, since $`V`$ is nonnegative. This implies, by the functional calculus,
$$\mathrm{Tr}\left((\mathrm{\Delta }+V)^\gamma \right)=\frac{1}{\mathrm{\Gamma }(\gamma )}_0^+\mathrm{}\mathrm{Tr}\left(e^{t(\mathrm{\Delta }+V)}\right)t^{\gamma 1}𝑑t.$$
Using (4), since $`V^{\frac{d}{2}\gamma }L^1(^d)`$, we get
$`\mathrm{Tr}\left((\mathrm{\Delta }+V)^\gamma \right)`$ $``$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(\gamma )}}{\displaystyle _0^+\mathrm{}}{\displaystyle _^d}(4\pi t)^{\frac{d}{2}}e^{tV(x)}t^{\gamma 1}𝑑x𝑑t`$
$``$ $`{\displaystyle \frac{\mathrm{\Gamma }(\gamma \frac{d}{2})}{(4\pi )^{\frac{d}{2}}\mathrm{\Gamma }(\gamma )}}{\displaystyle _^d}\left[V(x)\right]^{\frac{d}{2}\gamma }𝑑x.`$
We define $`𝒞(\gamma ):=`$ $`\frac{\mathrm{\Gamma }(\gamma \frac{d}{2})}{(4\pi )^{\frac{d}{2}}\mathrm{\Gamma }(\gamma )}`$ and obtain the announced inequality.
The optimality of the constant is established by the following example. Consider the potential $`V_\epsilon 1`$ in $`(0,\epsilon ^1\pi )^d=\mathrm{\Omega }_\epsilon ^d`$, $`V_\epsilon +\mathrm{}`$ in $`\mathrm{\Omega }_\epsilon ^c`$. Such a potential can be approximated by smooth potentials $`V_\epsilon ^n`$ such that $`V_\epsilon ^n1`$ in $`\mathrm{\Omega }_\epsilon `$ and $`lim_n\mathrm{}V_\epsilon ^n(x)=+\mathrm{}`$ for any $`x\mathrm{\Omega }_\epsilon ^c`$. The eigenvalues of $`\mathrm{\Delta }+V_\epsilon `$ on $`^d`$ are the same as the ones of $`\mathrm{\Delta }+V_\epsilon `$ on $`\mathrm{\Omega }_\epsilon `$ with zero Dirichlet boundary conditions on $`\mathrm{\Omega }_\epsilon `$:
$$1+\epsilon ^2\underset{j=1}{\overset{d}{}}n_j^2,n_1,n_2\mathrm{}n_d^{},$$
so that
$$\mathrm{Tr}\left((\mathrm{\Delta }+V_\epsilon )^\gamma \right)=\underset{n_1,n_2\mathrm{}n_d^{}}{}\left(1+\epsilon ^2\underset{j=1}{\overset{d}{}}n_j^2\right)^\gamma ,$$
which behaves asymptotically as $`\epsilon `$ tends to zero as
$`{\displaystyle \underset{n_1,n_2\mathrm{}n_d^{}}{}}\begin{array}{c}\\ {\displaystyle \mathrm{}}\\ \begin{array}{c}n_j1x_jn_j\\ j=1,2,\mathrm{}d\end{array}\end{array}{\displaystyle \frac{dx}{\left(1+\epsilon ^2|x|^2\right)^\gamma }}`$ $`=`$ $`{\displaystyle \frac{1}{(2\epsilon )^d}}{\displaystyle _^d}{\displaystyle \frac{dx}{\left(1+|x|^2\right)^\gamma }}`$
$`=`$ $`{\displaystyle \frac{|S^{d1}|}{(2\epsilon )^d}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{r^{d1}}{(1+r^2)^\gamma }}𝑑r.`$
This is precisely the right hand side of Inequality (2) as can be checked using $`(\pi /\epsilon )^d=_^dV_\epsilon ^{d/2\gamma }𝑑x`$ and
$$|S^{d1}|=\frac{2\pi ^{\frac{d}{2}}}{\mathrm{\Gamma }(\frac{d}{2})}\text{and}_0^{\mathrm{}}\frac{r^{d1}}{\left(1+r^2\right)^\gamma }𝑑r=\frac{\mathrm{\Gamma }(\gamma \frac{d}{2})\mathrm{\Gamma }(\frac{d}{2})}{2\mathrm{\Gamma }(\gamma )}.$$
$`\mathrm{}`$
From the above proof, it easy to see that the result of Theorem 1 can be generalized as follows. Let $`f`$ be a nonnegative function on $`_+`$ such that
$$_0^{\mathrm{}}f(t)\left(1+t^{d/2}\right)\frac{dt}{t}<\mathrm{}$$
(7)
and define
$$F(s):=_0^{\mathrm{}}e^{ts}f(t)\frac{dt}{t}\text{and}G(s):=_0^{\mathrm{}}e^{ts}\left(4\pi t\right)^{d/2}f(t)\frac{dt}{t}.$$
(8)
Notice that if $`d`$ is even, $`(4\pi )^{(d/2)}d^{(d/2)}G/ds^{(d/2)}(s)=F(s)`$. In the case of Theorem 1, $`F(s)=s^\gamma `$ and $`G(s)=\frac{\mathrm{\Gamma }(\gamma \frac{d}{2})}{(4\pi )^{\frac{d}{2}}\mathrm{\Gamma }(\gamma )}s^{\frac{d}{2}\gamma }`$.
###### Theorem 3
Let $`V`$ be in $`L_{\mathrm{loc}}^1(^d)`$ and bounded from below. Assume moreover that $`G(V)`$ is in $`L^1(^d)`$. With $`F`$ and $`G`$ defined by (8), if $`f`$ satisfies Asssumption (7), then
$$\underset{i^{}}{}F(\lambda _i(V))=\mathrm{Tr}\left[F\left(\mathrm{\Delta }+V\right)\right]_^dG(V(x))𝑑x.$$
Proof. The above inequality follows from the definition of $`F`$:
$$\mathrm{Tr}\left[F\left(\mathrm{\Delta }+V\right)\right]=_0^+\mathrm{}\mathrm{Tr}\left(e^{t(\mathrm{\Delta }+V)}\right)f(t)\frac{dt}{t},$$
Inequality (4) and the definition of $`G`$. $`\mathrm{}`$
As an example, if we apply Theorem 3 with $`F(s)=e^s`$, $`f(s)=\delta (s1)`$ and $`G(s)=(4\pi )^{d/2}e^s`$, we get
$$\underset{i^{}}{}e^{\lambda _i(V)}\frac{1}{(4\pi )^{d/2}}_^de^{V(x)}𝑑x.$$
(9)
In the special case $`V(x)=A^2|x|^2+B`$, eigenvalues are explicitly given as
$$B+\underset{j=1}{\overset{d}{}}\left(2n_j+1\right)A,n_1,n_2\mathrm{}n_d,$$
and we can compute
$$\mathrm{Tr}\left(e^{t(\mathrm{\Delta }+V)}\right)=\underset{i^{}}{}e^{t\lambda _i(V)}=e^{Bt}\underset{j=1}{\overset{d}{}}\left(\underset{n_j^{}}{}e^{\left(2n_j+1\right)At}\right)=\frac{e^{Bt}}{[2\mathrm{sinh}(At)]^d}.$$
(10)
On the other hand,
$$\frac{1}{(4\pi )^{d/2}}_^de^{V(x)}𝑑x=\frac{e^B}{(2A)^d}.$$
Putting these estimates together in the case $`t=1`$ shows that the upper bound in (9), namely
$$\frac{1}{(4\pi )^{d/2}}\left(\frac{A}{\mathrm{sinh}A}\right)^d=\frac{_i^{}e^{\lambda _i(V)}}{_^de^{V(x)}𝑑x}\frac{1}{(4\pi )^{d/2}},$$
is achieved in the limit $`A0_+`$.
Identity (10) is also useful in the case $`F(s)=s^\gamma `$ considered in Theorem 1. Using $`\mathrm{Tr}\left((\mathrm{\Delta }+V)^\gamma \right)=\frac{1}{\mathrm{\Gamma }(\gamma )}_0^+\mathrm{}\mathrm{Tr}\left(e^{t(\mathrm{\Delta }+V)}\right)t^{\gamma 1}𝑑t`$, we obtain in the special case $`V(x)=A^2|x|^2+B`$ the identity
$$\mathrm{Tr}\left((\mathrm{\Delta }+V)^\gamma \right)=\frac{1}{\mathrm{\Gamma }(\gamma )}_0^+\mathrm{}\frac{e^{Bt}}{[2\mathrm{sinh}(At)]^d}t^{\gamma 1}𝑑t=\frac{B^\gamma }{\mathrm{\Gamma }(\gamma )}_0^+\mathrm{}\frac{e^t}{[2\mathrm{sinh}(st)]^d}t^{\gamma 1}𝑑t$$
with $`s:=B/A`$. Under the additional restriction $`\gamma >d`$, we get
$$_^dV^{\frac{d}{2}\gamma }𝑑x=B^{d\gamma }A^d\pi ^{\frac{d}{2}}\frac{\mathrm{\Gamma }(\gamma d)}{\mathrm{\Gamma }\left(\gamma \frac{d}{2}\right)}.$$
With $`𝒞(\gamma )=\frac{\mathrm{\Gamma }(\gamma \frac{d}{2})}{(4\pi )^{d/2}\mathrm{\Gamma }(\gamma )}`$, this shows that
$$\frac{\mathrm{Tr}\left((\mathrm{\Delta }+V)^\gamma \right)}{𝒞(\gamma )_^dV^{\frac{d}{2}\gamma }𝑑x}=\frac{s^d}{\mathrm{\Gamma }(\gamma d)}_0^{\mathrm{}}\frac{t^{\gamma 1}e^t}{\left(\mathrm{sinh}(st)\right)^d}dt=:q(s).$$
It is easy to check that the function $`sq(s)`$ bounded by $`1`$ and achieves $`1`$ in the limit $`s0_+`$.
## 3 Stability for the linear Schrödinger equation
In this section we come back to the physical motivation of Theorems 1 and 3 with more details than in the introduction and state a list of examples corresponding to various functions $`F`$.
### 3.1 Notations and assumptions
Let $`E[\psi ]:=_^d(|\psi |^2+V|\psi |^2)𝑑x`$ and assume that $`V`$ is a potential such that the operator $`H_V:=\mathrm{\Delta }+V`$ has an infinite nondecreasing sequence of eigenvalues $`(\lambda _i(V))_i^{}`$:
$$\lambda _i(V):=\underset{\begin{array}{c}FL^2(^d)\\ \mathrm{dim}(F)=i\end{array}}{inf}\underset{\psi F}{sup}E[\psi ].$$
Here the eigenvalues are counted with multiplicity, and to each $`\lambda _i(V)`$, $`i^{}`$, we can associate an eigenfunction $`\overline{\psi }_i`$ such that $`\overline{𝝍}:=(\overline{\psi }_i)_i^{}`$ is an orthonormal sequence:
$$(\overline{\psi }_i,\overline{\psi }_j)_{L^2(^d)}=\delta _{ij}i,j^{}.$$
As in Section 1, we also define $`\overline{\nu }_i:=(\beta ^{})^1(\lambda _i(V))`$ for any $`i^{}`$, $`\overline{𝝂}:=(\overline{\nu }_i)_i^{}`$. The free energy of the mixed state $`(𝝂,𝝍)=((\nu _i)_i^{},(\psi _i)_i^{})_+^{^{}}\times (L^2(^d))^{^{}}`$ is
$$[𝝂,𝝍]:=\underset{i^{}}{}\beta (\nu _i)+\underset{i^{}}{}\nu _iE[\psi _i]$$
for some given function $`\beta `$. If the potential $`V`$ is such that $`\mathrm{\Delta }+V`$ has an unbounded sequence of eigenvalues, it is easy to see that $`[𝝂,𝝍]`$ is defined only if $`lim_i\mathrm{}\nu _i=0`$. This allows us to re-order the sequence $`(𝝂,𝝍)`$ in such a way that $`(\nu _i)_i^{}`$ is a non increasing sequence converging to $`0`$, and we may restrict the domain of the free energy $``$ to $`S\times (L^2(^d))^{^{}}`$, where $`S`$ denotes the set of non increasing sequences in $`_+`$ converging to $`0`$, such that $`_i^{}\beta (\nu _i)`$ is absolutely convergent. Notice that whenever it is finite, $`_i^{}\beta (\nu _i)`$ is absolutely convergent by the assumption $`lim_i\mathrm{}\nu _i=0`$.
We shall say that Assumption (H) holds if $`\beta `$ is a strictly convex function with $`\beta (0)=0`$, which is $`C^1`$ on the interior of its support and if the potential $`V`$ is such that $`\mathrm{\Delta }+V`$ has an unbounded sequence of eigenvalues $`(\lambda _i(V))_i^{}`$ for which
$$\left|\underset{i^{}}{}\beta (\overline{\nu }_i)\right|<\mathrm{}\text{and}\left|\underset{i^{}}{}\overline{\nu }_i\lambda _i(V)\right|<\mathrm{},$$
where $`\overline{\nu }_i:=(\beta ^{})^1(\lambda _i(V))`$ for any $`i^{}`$. As seen in Section 1, Assumption (H) is a consequence of the Lieb-Thirring type inequalities of Theorem 1 if $`\beta (\nu )=(1m)^{m1}m^m\nu ^m`$, $`m(0,1)`$ (see below Example 2 for more details). In the framework of Theorem 3, $`F(\lambda )=\beta (\nu )\lambda \nu `$ with $`\nu =(\beta ^{})^1(\lambda )`$.
### 3.2 Minimizers of the free energy
###### Proposition 4
Assume that $`\beta `$ and $`V`$ are such that (H) holds. Then there exists a minimizer $`(\overline{𝛎},\overline{𝛙})S\times (L^2(^d))^{^{}}`$ of $``$ under the constraint
$$(\overline{\psi }_i,\overline{\psi }_j)_{L^2(^d)}=\delta _{ij}i,j^{}.$$
Moreover,
$$\overline{\nu }_i=(\beta ^{})^1(\lambda _i(V))$$
and if $`\overline{\nu }_i`$ is positive for any $`i^{}`$, the sequence $`\overline{𝛙}=(\overline{\psi }_i)_i^{}`$ is unique up to any unitary transformation which leaves all eigenspaces of $`\mathrm{\Delta }+V`$ invariant. In particular, any $`\overline{\psi }_i`$ can be multiplied by an arbitrary constant phase factor, so that we may assume that $`\overline{\psi }_i`$ is real.
To prove this result we first prove some results about finite mixed states : given any $`n^{}`$, we can define the projection $`P_n`$ of a mixed state $`(𝝂,𝝍)S\times (L^2(^d))^{^{}}`$ onto the $`n`$-finite mixed states by $`P_n[𝝂,𝝍]:=(\stackrel{~}{𝝂},𝝍)`$ with $`\stackrel{~}{\nu }_i=\nu _i`$ for any $`i=1,\mathrm{\hspace{0.17em}2},\mathrm{}n`$ and $`\stackrel{~}{\nu }_i=0`$ for any $`in+1`$. Let $`_n:=P_n`$:
$$_n[𝝂,𝝍]:=\underset{i=1}{\overset{n}{}}\left(\beta (\nu _i)+\nu _iE[\psi _i]\right).$$
Notice indeed that $`\beta (0)=0`$, so that $`\beta (\stackrel{~}{\nu }_i)=0`$ for any $`in+1`$. We may decompose $`_n`$ into an entropy and an energy term as follows.
###### Lemma 5
Under Assumption (H), for any $`(𝛎,𝛙)S\times (L^2(^d))^{^{}}`$ such that $`𝛙=(\psi _i)_i^{}`$ is an orthonormal sequence,
$$_n[𝝂,𝝍]_n[\overline{𝝂},\overline{𝝍}]=\underset{i=1}{\overset{n}{}}\left(\beta (\nu _i)\beta (\overline{\nu }_i)\beta ^{}(\overline{\nu }_i)(\nu _i\overline{\nu }_i)\right)+\underset{i=1}{\overset{n}{}}\nu _i\left(E[\psi _i]E[\overline{\psi }_i]\right).$$
Proof. An elementary computation gives
$$\beta ^{}(\overline{\nu }_i)(\nu _i\overline{\nu }_i)+\nu _iE[\overline{\psi }_i]=\lambda _i(V)(\nu _i\overline{\nu }_i)+\nu _i\lambda _i(V)=\overline{\nu }_i\lambda _i(V)=\overline{\nu }_iE[\overline{\psi }_i].$$
$`\mathrm{}`$
We are now going to study independently the two terms of $`_n[𝝂,𝝍]_n[\overline{𝝂},\overline{𝝍}]`$ and start with the entropy term.
###### Lemma 6
Assume that $`inf_{s>0}\beta ^{\prime \prime }(s)s^{2p}=:\alpha >0`$ for some $`p[1,2]`$. Then for any sequence $`(\nu _i)_i^{}_+^{^{}}`$, if $`_i^{}\beta (\nu _i)`$ and $`_i^{}\nu _i\beta ^{}(\overline{\nu }_i)`$ are absolutely convergent, then $`(\nu _i\overline{\nu }_i)_i^{}\mathrm{}^p`$ and
$$\underset{i^{}}{}\left(\beta (\nu _i)\beta (\overline{\nu }_i)\beta ^{}(\overline{\nu }_i)(\nu _i\overline{\nu }_i)\right)2^{2/p}\alpha 𝝂\overline{𝝂}_\mathrm{}^p^2\mathrm{min}\{𝝂_\mathrm{}^p^{p2},\overline{𝝂}_\mathrm{}^p^{p2}\}.$$
See for a continuous version of this inequality. We may also refer to in the case of $`\beta (\nu )=\nu \mathrm{log}\nu \nu `$ and von Neumann algebras, and to for a review of the so-called Csiszár-Kullback inequalities. For the completeness of the paper, we give a short proof of this result.
Proof. For any $`i^{}`$, let $`\zeta _i[\mathrm{min}(\nu _i,\overline{\nu }_i),\mathrm{max}(\nu _i,\overline{\nu }_i)]`$ be an intermediate nonnegative point such that
$$\underset{i^{}}{}\left(\beta (\nu _i)\beta (\overline{\nu }_i)\beta ^{}(\overline{\nu }_i)(\nu _i\overline{\nu }_i)\right)=\frac{1}{2}\underset{i^{}}{}\beta ^{}(\zeta _i)(\nu _i\overline{\nu }_i)^2\frac{\alpha }{2}\underset{i^{}}{}\zeta _i^{p2}(\nu _i\overline{\nu }_i)^2.$$
Let $`^{}`$. Using
$$\left(\underset{i}{}|\nu _i\overline{\nu }_i|^p\zeta _i^{p(p2)/2}\zeta _i^{p(2p)/2}\right)^{2/p}\underset{i}{}\zeta _i^{p2}(\nu _i\overline{\nu }_i)^2\left(\underset{i}{}\zeta _i^p\right)^{(2p)/p},$$
we get
$$\underset{i}{}\zeta _i^{p2}(\nu _i\overline{\nu }_i)^2\left(\underset{i}{}|\nu _i\overline{\nu }_i|^p\right)^{2/p}\left(\underset{i}{}\zeta _i^p\right)^{12/p}.$$
On $`=\{i^{}:\nu _i>\overline{\nu }_i\}`$ (respectively $`=\{i^{}:\nu _i<\overline{\nu }_i\}`$), we estimate $`_i\zeta _i^p`$ from above by $`_i\nu _i^p`$ (resp. by $`_i\overline{\nu }_i^p`$). Using the inequality $`(a+b)^r2^{r1}(a^r+b^r)`$ for any $`a`$, $`b0`$ and $`2/p=r1`$, we completes the proof. $`\mathrm{}`$
Next, we turn our attention to the energy term and recall a result given, for instance, in :
###### Proposition 7
Let $`V`$ be a potential such that the sequence of eigenvalues $`(\lambda _i(V))_i^{}`$ of $`H_V`$ is unbounded, and choose any $`n`$ functions $`\psi _1,\mathrm{},\psi _n`$ that are orthonormal in $`L^2(^d)`$. Then
$$\underset{i=1}{\overset{n}{}}E[\psi _i]\underset{i=1}{\overset{n}{}}\lambda _i(V).$$
We extend this property to orthogonal families:
###### Lemma 8
Assume that $`V`$ is a potential as above. For any orthogonal family $`(\varphi _i)_{1in}`$ in $`L^2(^d)`$, with $`\varphi _i^2=\nu _i`$ and $`\nu _1\mathrm{}\nu _n`$, we get
$$\underset{i=1}{\overset{n}{}}E[\varphi _i]\underset{i=1}{\overset{n}{}}\nu _i\lambda _i(V).$$
Proof. We prove this result by induction on $`n`$. The case $`n=1`$ is trivial. Suppose that the result holds for any orthogonal system of $`n`$ functions, and take $`(\varphi _i)_{1in+1}`$ with nonincreasing squared norms (or occupation numbers) $`\nu _1\mathrm{}\nu _{n+1}0`$. If $`\nu _{n+1}=0`$, then the induction assumption directly gives the result. Assume next that $`\nu _{n+1}>0`$. By Proposition 7, we have
$$\underset{i=1}{\overset{n+1}{}}E\left[\frac{\varphi _i}{\varphi _i}\right]\underset{i=1}{\overset{n+1}{}}\lambda _i(V).$$
Multiplying by $`\nu _{n+1}`$, we obtain
$$\underset{i=1}{\overset{n+1}{}}\frac{\nu _{n+1}}{\nu _i}E[\varphi _i]\underset{i=1}{\overset{n+1}{}}\nu _{n+1}\lambda _i(V),$$
hence
$$\underset{i=1}{\overset{n+1}{}}E[\varphi _i]\underset{i=1}{\overset{n}{}}\left[\frac{\nu _i\nu _{n+1}}{\nu _i}E[\varphi _i](\nu _i\nu _{n+1})\lambda _i(V)\right]+\underset{i=1}{\overset{n+1}{}}\nu _i\lambda _i(V).$$
(11)
Since $`\nu _i\nu _{n+1}`$, we can define the family $`(\mu _i\varphi _i)_{1in}`$ with $`\mu _i:=(\frac{\nu _i\nu _{n+1}}{\nu _i})^{1/2}`$, which is orthogonal. By the induction hypothesis we get
$$\underset{i=1}{\overset{n}{}}\frac{\nu _i\nu _{n+1}}{\nu _i}E[\varphi _i]=\underset{i=1}{\overset{n}{}}E\left(\mu _i\varphi _i\right)\underset{i=1}{\overset{n}{}}\mu _i\varphi _i^2\lambda _i(V)=\underset{i=1}{\overset{n}{}}(\nu _i\nu _{n+1})\lambda _i(V).$$
In Inequality (11), the first sum of the right hand side is then nonnegative. For the system of the $`n+1`$ orthogonal functions, we obtain
$$\underset{i=1}{\overset{n+1}{}}E[\varphi _i]\underset{i=1}{\overset{n+1}{}}\nu _i\lambda _i(V),$$
which completes the proof of Lemma 8. $`\mathrm{}`$
Proof of Proposition 4. By Lemma 8 we get
$$_n[𝝂,𝝍]\underset{i=1}{\overset{n}{}}\left(\beta (\nu _i)+\nu _i\lambda _i(V)\right),$$
hence a minimization of $`_n`$ under the constraint $`(\psi _i,\psi _j)_{L^2(^d)}=\delta _{ij}`$ directly gives, for any $`[𝝂,𝝍]S\times (L^2(^d))^{^{}}`$,
$$_n[𝝂,𝝍]_n[\overline{𝝂},\overline{𝝍}].$$
Assumption (H) gives the absolute convergence of the series in the definition of $`(\overline{𝝂},\overline{𝝍})`$. Suppose now that there exists $`(\stackrel{~}{𝝂},\stackrel{~}{𝝍})S\times (L^2(^d))^{^{}}`$ such that $`(\stackrel{~}{𝝂},\stackrel{~}{𝝍})<(\overline{𝝂},\overline{𝝍})`$. This implies the existence of a $`N^{}`$ such that
$$\underset{i=1}{\overset{N}{}}\left(\beta (\stackrel{~}{\nu _i})+\stackrel{~}{\nu _i}E[\stackrel{~}{\psi _i}]\right)<\underset{i=1}{\overset{N}{}}\left(\beta (\overline{\nu }_i)+\overline{\nu }_iE[\overline{\psi }_i]\right),$$
which contradicts the result on finite mixed states. $`\mathrm{}`$
### 3.3 Stability
As a consequence of the conservation of the energy $`E[]`$ under the evolution according to the Schrödinger operator $`i_tH_V`$, we obtain the conservation of the free energy. Notice here that all above computations have been done with functions taking real values and need to be adapted to the case of complex valued functions as soon as we consider solutions to the time-dependent Schrödinger equation.
###### Proposition 9
Assume (H) and consider an initial mixed state $`(𝛎,𝛙^0)_+^{^{}}\times (L^2(^d))^{^{}}`$. If $`(𝛎,𝛙(t))`$ is the mixed state where each of the components evolves according to the linear Schrödinger equation
$$i_t\psi _j=\mathrm{\Delta }\psi _j+V\psi _j,x^d,t>0$$
with initial data $`\psi _j^0`$ for any $`j^{}`$, then
$$[𝝂,𝝍(t)]=[𝝂,𝝍^0]t>0.$$
To state a dynamical stability result, we have to impose a decay property of the sequence of occupation numbers as in Lemma 8. From Lemmas 6 and 8, and Proposition 9, we deduce the
###### Corollary 10
Consider an initial mixed state $`(𝛎,𝛙^0)S\times (L^2(^d))^{^{}}`$ with a nonincreasing sequence of occupation numbers $`𝛎`$. Under the assumption of Lemma 6, if (H) is satisfied, then for any $`t>0`$,
$$2^{2/p}\alpha 𝝂\overline{𝝂}_\mathrm{}^p^2\mathrm{min}\{𝝂_\mathrm{}^p^{p2},\overline{𝝂}_\mathrm{}^p^{p2}\}+\underset{i^{}}{}\nu _i\left(E[\psi _i(t)]\lambda _i(V)\right)[𝝂,𝝍^0],$$
where both terms of the left hand side are nonnegative.
### 3.4 Examples
We conclude these comments on stability results by a list of examples of various functions $`\beta `$ and by the corresponding Lieb-Thirring type inequalities given by Theorem 3 with $`F(s)=(\beta (\beta ^{})^1)(s)+s(\beta ^{})^1(s)`$. We refer to for a similar discussion in a non quantum mechanical context.
Example 1. Let $`m>1`$ and consider $`\beta (\nu ):=(m1)^{m1}m^m\nu ^m`$. With $`\beta ^{}(\nu )=(m1)^{m1}m^{1m}\nu ^{m1}=\lambda `$ and $`m=\frac{\gamma }{\gamma 1}`$, we get: $`(\beta (\nu )+\lambda \nu )=F(\lambda )=(\lambda )^\gamma `$, which corresponds to the setting of the standard Lieb-Thirring inequality (1). The case $`\gamma (0,1)`$ is formally covered by $`\beta (\nu ):=(1m)^{m1}|m|^m\nu ^m`$ with $`m(\mathrm{},0)`$, $`m=\frac{\gamma }{\gamma 1}`$ again and $`F(s)=(s)^\gamma `$, but in this case, $`\beta `$ is not convex and the free energy $``$ cannot be defined as above.
Example 2. As seen above, for $`m<1`$ and $`\beta (\nu ):=(1m)^{m1}m^m\nu ^m`$, with $`\beta ^{}(\nu )=(1m)^{m1}m^{1m}\nu ^{m1}=\lambda `$ and $`m=\frac{\gamma }{\gamma +1}`$, we get: $`F(\lambda )=\lambda ^\gamma `$, which corresponds to the setting of Theorem 1.
Example 3. If $`\beta (\nu ):=\nu \mathrm{log}\nu \nu `$, then $`\beta ^{}(\nu )=\mathrm{log}\nu =\lambda `$. According to Theorem 3, the corresponding inequality is
$$\underset{i^{}}{}e^{\lambda _i(V)}\frac{1}{(4\pi )^{d/2}}_^de^{V(x)}𝑑x.$$
This case can formally be seen as the limit case $`m1`$ in Examples 1 and 2. Here $`F(s)=e^s`$, $`G(s)=(4\pi )^{d/2}e^s`$.
Example 4. If $`\beta (\nu ):=\nu \mathrm{log}\nu +(1\nu )\mathrm{log}(1\nu )`$, then $`\beta ^{}(\nu )=\mathrm{log}\left(\frac{\nu }{1\nu }\right)=\lambda `$ and $`F(s)=\mathrm{log}(1+e^s)`$. According to Theorem 3, the corresponding inequality is
$$\underset{i^{}}{}\mathrm{log}\left(1+e^{\lambda _i(V)}\right)_^dG(V(x))𝑑x$$
where $`G`$ is given in terms of $`F`$ by (8).
In all the above examples we have to assume that $`lim_i\mathrm{}\lambda _i(V)=+\mathrm{}`$, except in Example 1 where $`\lambda _i(V)0`$, $`lim_i\mathrm{}\lambda _i(V)=0`$ and we adopt the convention that $`\lambda _i(V)=0`$ for any $`i>N`$ if there are only $`N`$ negative eigenvalues.
## 4 Lieb-Thirring Gagliardo-Nirenberg inequalities
In this section, we will focus on consequences of Theorems 1 and 3, when one takes only partial sums, and especially when only the first eigenvalue is considered.
### 4.1 Optimal constant in the Lieb-Thirring conjecture
We begin with a remark on the connection of the best constant in the Lieb-Thirring conjecture and its extension for $`d>1`$:
$$C_{\mathrm{LT}}^{(1)}(\gamma ):=\underset{\begin{array}{c}V𝒟(^d)\\ V0\end{array}}{inf}\frac{|\lambda _1(V)|^\gamma }{_^d|V|^{\gamma +\frac{d}{2}}𝑑x}$$
with the best constant in some special Gagliardo-Nirenberg inequalities. Such a result has already been established in (also see for earlier references). We give it here for completeness and in order to insist on some interesting scaling properties.
Define the function set for the potential $`V`$ by
$$X_\gamma :=\{VL^{\gamma +\frac{d}{2}}(^d):V0,V0a.e.\}$$
and note that by density of $`𝒟(^d)`$ in $`L^{\gamma +\frac{d}{2}}(^d)`$, it holds that
$$C_{\mathrm{LT}}^{(1)}(\gamma )=\underset{\begin{array}{c}VX_\gamma \\ V0,V0a.e.\end{array}}{sup}\frac{|\lambda _1(V)|^\gamma }{_^dV^{\gamma +\frac{d}{2}}𝑑x}.$$
By density of $`𝒟(^d)`$ in $`H^1(^d)`$, we have
$$\lambda _1(V)=\underset{\begin{array}{c}uH^1(^d)\\ u0a.e.\end{array}}{inf}\frac{_^d|u|^2𝑑x+_^dV|u|^2𝑑x}{_^d|u|^2𝑑x}.$$
Let
$$q:=\frac{2\gamma +d}{2\gamma +d2}$$
and consider the optimal constant $`C_{\mathrm{GN}}(\gamma )`$ of the Gagliardo-Nirenberg inequality corresponding to the embedding of $`H^1(^d)`$ into $`L^{2q}(^d)`$:
$$C_{\mathrm{GN}}(\gamma )=\underset{\begin{array}{c}uH^1(^d)\\ u0a.e.\end{array}}{inf}\frac{u_{L^2(^d)}^{\frac{d}{2\gamma +d}}u_{L^2(^d)}^{\frac{2\gamma }{2\gamma +d}}}{u_{L^{2q}(^d)}}.$$
(12)
Notice that for $`\gamma >\mathrm{max}(0,1d/2)`$,
$$q>1\text{and}2q<\frac{2d}{d2}.$$
###### Theorem 11
Let $`d^{}`$. For any $`\gamma >\mathrm{max}(0,1\frac{d}{2})`$,
$$C_{\mathrm{LT}}^{(1)}(\gamma )=\kappa _1(\gamma )\left[C_{\mathrm{GN}}(\gamma )\right]^{\kappa _2(\gamma )},$$
where
$$\kappa _1(\gamma )=\frac{2\gamma }{d}\left(\frac{d}{2\gamma +d}\right)^{1+\frac{d}{2\gamma }}\text{and}\kappa _2(\gamma )=2+\frac{d}{\gamma }.$$
Moreover, the constant $`C_{\mathrm{LT}}^{(1)}(\gamma )`$ is optimal and achieved by a unique pair of functions $`(u,V)`$, up to multiplications by a constant, scalings and translations.
The scaling invariance can be made clear by redefining
$$\left[C_{\mathrm{LT}}^{(1)}(\gamma )\right]^{\frac{1}{\gamma }}=\underset{\begin{array}{c}VX_\gamma \\ V0,V0a.e.\end{array}}{sup}\underset{\begin{array}{c}uH^1(^d)\\ u0a.e.\end{array}}{sup}R(u,V)$$
where
$$R(u,V)=\frac{_^dV|u|^2𝑑x+_^d|u|^2𝑑x}{_^d|u|^2𝑑xV_{L^{\gamma +\frac{d}{2}}(^d)}^{1+\frac{d}{2\gamma }}}.$$
Note indeed that $`\lambda _1(V)0`$, and $`R(u,V)`$ is invariant under the transformation
$$(u,V)(u_\lambda =u(\lambda ),V_\lambda =\lambda ^2V(\lambda )),$$
i.e., $`R(u_\lambda ,V_\lambda )=R(u,V)`$ for any $`\lambda >0`$.
Proof of Theorem 11. By Hölder’s inequality,
$$_^d|V||u|^2𝑑xAu_{L^{2q}(^d)}^2\text{with}A:=V_{L^{\gamma +\frac{d}{2}}(^d)}.$$
Let $`\tau :=u_{L^{2q}(^d)}/u_{L^2(^d)}`$. The Gagliardo-Nirenberg inequality (12), namely
$$C_{\mathrm{GN}}(\gamma )u_{L^{2q}(^d)}u_{L^2(^d)}^\theta u_{L^2(^d)}^{1\theta }$$
with $`\theta =\frac{d}{2\gamma +d}`$ can be rewritten as
$$\frac{u_{L^2(^d)}}{u_{L^2(^d)}}\left[C_{\mathrm{GN}}(\gamma )\tau \right]^{\frac{1}{\theta }}.$$
Putting these estimates together, we obtain
$$R(u,V)\frac{A\tau ^2[C_{\mathrm{GN}}(\gamma )]^{\frac{2}{\theta }}\tau ^{\frac{2}{\theta }}}{A^{1+\frac{d}{2\gamma }}}.$$
An optimization on $`\tau `$ shows the bound of Theorem 11, which is independent of $`A`$, and gives the expressions of $`\kappa _1(\gamma )`$ and $`\kappa _2(\gamma )`$.
The estimate is achieved since all above inequalities can be saturated by considering
$$|V|^{\gamma +\frac{d}{2}2}V=|u|^2V=V_u=|u|^{\frac{4}{2\gamma +d2}}=|u|^{2(q1)},$$
(13)
where $`u`$ is a solution of
$$\mathrm{\Delta }u+|u|^{2(q1)}uu=0\text{in}^d.$$
Up to a scaling, these two equations are the Euler-Lagrange equations corresponding to the maximization in $`V`$ and $`u`$ respectively. Because of the second equation, the relation with the Gagliardo-Nirenberg inequality is straightforward. In other words,
$$R(u,V)R(u,V_u)=\frac{_^d|u|^{2q}𝑑x_^d|u|^2𝑑x}{_^d|u|^2𝑑x\left(_^d|u|^{2q}𝑑x\right)^{\frac{1}{\gamma }}}=\frac{_^d|u_\lambda |^{2q}𝑑x_^d|u_\lambda |^2𝑑x}{\left(_^d|u_\lambda |^{2q}𝑑x\right)^{\frac{1}{\gamma }}}$$
where $`u_\lambda =\lambda ^{\frac{1}{q1}}u(\lambda )`$ and $`\lambda ^{\left(\frac{d}{2}\frac{1}{q1}\right)}=u_{L^2(^d)}`$, so that $`u_\lambda _{L^2(^d)}=1`$,
$$R(u,V)\frac{\tau ^{2q}[C_{\mathrm{GN}}(\gamma )]^{\frac{2}{\theta }}\tau ^{\frac{2}{\theta }}}{\tau ^{\frac{2q}{\gamma }}},$$
and the result holds by optimizing in $`\tau =u_\lambda _{L^{2q}(^d)}`$. $`\mathrm{}`$
Remark. The optimal function in the Gagliardo-Nirenberg inequality (12) is given as the nonnegative solution of (13) in $`H^1(^d)`$. This solution is radial, positive, decreasing, and unique up to translations, multiplication by constants and scalings. See for instance for uniqueness results of radial solutions, and references therein for earlier related results. Optimal function are not explicitly known in general but are easy to compute numerically as well as the optimal constants, see for instance .
### 4.2 Theorem 1 and Gagliardo-Nirenberg inequalities
In this section, we adapt the remarks of Section 4.1 to the case $`V0`$ of Theorem 1. The interpolation of $`u_{L^{2q}(^d)}`$, with $`1<q<d/(d2)`$, $`d3`$, between $`u_{L^2(^d)}`$ and $`u_{L^2(^d)}`$ of the previous section is a standard case of Gagliardo-Nirenberg inequalities, but there is also another interesting case in Gagliardo-Nirenberg inequalities, which is somewhat less standard. It corresponds to the interpolation of $`u_{L^2(^d)}`$ between $`u_{L^2(^d)}`$ and $`_^d|u|^{2q}𝑑x`$ for some $`q(0,1)`$. See for a similar setting, where both cases have been taken into account. What we establish in this section is that these less standard inequalities are related to the estimate of $`[\lambda _1(V)]^\gamma `$ in terms of $`_^dV^{d/2\gamma }𝑑x`$.
Consider now a nonnegative smooth potential $`V𝒞^{\mathrm{}}(^d)`$ such that
$$\underset{|x|+\mathrm{}}{lim}V(x)=+\mathrm{}$$
and denote by $`\lambda _1(V)`$, $`\lambda _2(V)`$, …the positive eigenvalues of $`\mathrm{\Delta }+V`$. By density we may extend this set of potentials to the set
$$Y_\gamma :=\{V^{\frac{d}{2}\gamma }L^1(^d):V0,V+\mathrm{}a.e.\}.$$
Let
$$q:=\frac{2\gamma d}{2(\gamma +1)d}(0,1)$$
and define an optimal constant of a second type Gagliardo-Nirenberg inequality by
$$C_{\mathrm{GN}}^{}(\gamma )=\underset{\begin{array}{c}uH^1(^d),u0a.e.\\ _^d|u|^{2q}𝑑x<\mathrm{}\end{array}}{inf}\frac{u_{L^2(^d)}^{\frac{d}{2\gamma }}\left(_^d|u|^{2q}𝑑x\right)^{\frac{1}{2q}(1\frac{d}{2\gamma })}}{u_{L^2(^d)}}.$$
(14)
###### Theorem 12
Let $`d^{}`$. For any $`\gamma >d/2`$, there exists a positive constant $`𝒞^{(1)}(\gamma )`$ such that, for any $`VY_\gamma `$,
$$\left[\lambda _1(V)\right]^\gamma 𝒞^{(1)}(\gamma )_^dV^{\frac{d}{2}\gamma }𝑑x.$$
As in Theorem 11, the optimal value of $`𝒞^{(1)}(\gamma )`$ is such that
$$𝒞^{(1)}(\gamma )=\kappa _1(\gamma )\left[C_{\mathrm{GN}}^{}(\gamma )\right]^{\kappa _2(\gamma )},$$
where $`\kappa _1(\gamma )=\frac{(2q)^{\gamma \frac{d}{2}}(d(1q))^{\frac{d}{2}}}{(d(1q)+2q)^\gamma }\text{and}\kappa _2(\gamma )=2\gamma `$. Moreover, the constant $`𝒞^{(1)}(\gamma )`$ is achieved by a unique pair of functions $`(u,V)`$, up to multiplications by a constant, scalings and translations.
Notice that $`q<1`$, and $`2q>1`$ if and only if $`\gamma >1+d/2`$. The best constant in the above inequality is
$$C_{\mathrm{LT}}^{(1)}(\gamma ):=\underset{\begin{array}{c}VY_\gamma \\ V0,V0a.e.\end{array}}{sup}\frac{[\lambda _1(V)]^\gamma }{_^dV^{\frac{d}{2}\gamma }𝑑x}.$$
The scaling invariance can be made clear by writing
$$\left[𝒞^{(1)}(\gamma )\right]^{\frac{1}{\gamma }}=\underset{\begin{array}{c}VX_\gamma \\ V0,V0a.e.\end{array}}{sup}\underset{\begin{array}{c}uH^1(^d)\\ u0a.e.\end{array}}{sup}R(u,V)$$
where
$$R(u,V):=\frac{_^d|u|^2𝑑x\left(_^dV^{\frac{d}{2}\gamma }𝑑x\right)^{\frac{1}{\gamma }}}{_^d|u|^2𝑑x+_^dV|u|^2𝑑x}$$
is invariant under the transformation
$$(u,V)(u_\lambda =u(\lambda ),V_\lambda =\lambda ^2V(\lambda )),$$
i.e., $`R(u_\lambda ,V_\lambda )=R(u,V)`$ for any $`\lambda >0`$.
Proof of Theorem 12. By Hölder’s inequality,
$$_^du^{2q}𝑑x=_^du^{2q}V^qV^q𝑑x\left(_^dV|u|^2𝑑x\right)^q\left(_^dV^{\frac{q}{1q}}𝑑x\right)^{1q}.$$
With $`A:=\left(_^dV^{q/(1q)}𝑑x\right)^{(1q)/q}=\left(_^dV^{d/2\gamma }𝑑x\right)^{2/(2\gamma d)}`$, this means that
$$_^dV|u|^2𝑑xA\left(_^d|u|^{2q}𝑑x\right)^{\frac{1}{q}}.$$
We may therefore estimate $`R(u,V)`$ as follows:
$$R(u,V)\frac{u_{L^2(^d)}^2A^{1\frac{d}{2\gamma }}}{u_{L^2(^d)}^2+A\left(_^d|u|^{2q}𝑑x\right)^{\frac{1}{q}}}.$$
An optimization under the scaling $`\lambda \lambda ^{d/2}u(/\lambda )`$, which leaves the $`L^2(^d)`$-norm invariant, shows that
$$u_{L^2(^d)}^2+Au_{L^2(^d)}^2u_{L^{2q}(^d)}^{\frac{2d(1q)}{d(1q)+2q}}\left(_^d|u|^{2q}𝑑x\right)^{\frac{2}{d(1q)+2q}}A^{\frac{2q}{d(1q)+2q}}(\kappa _1(\gamma ))^{\frac{1}{\gamma }}$$
using $`\frac{2q}{d(1q)+2q}=1\frac{d}{2\gamma }`$. Using the Gagliardo-Nirenberg inequality (14), we get
$$u_{L^2(^d)}^2+Au_{L^2(^d)}^2\left|C_{\mathrm{GN}}^{}(\gamma )\right|^2u_{L^2(^d)}^2A^{1\frac{d}{2\gamma }}(\kappa _1(\gamma ))^{\frac{1}{\gamma }}$$
which proves that $`𝒞^{(1)}(\gamma )\kappa _1(\gamma )[C_{\mathrm{GN}}^{}(\gamma )]^{\kappa _2(\gamma )}`$. It is moreover easy to check that the equality holds in Hölder’s inequality if $`V^{\frac{d}{2}\gamma 1}`$ is proportional to $`|u|^2`$. By taking a minimizer of (14), this completes the proof of Theorem 12. $`\mathrm{}`$
Remark. Notice that optimal functions are not explicitly known, unless $`d=1`$. Solutions to the Euler-Lagrange equations have compact support and minimal ones are radially symmetric and unique up to translations, see . Also see for more details.
### 4.3 General case
We may try to generalize the approach used for power laws to general nonlinearities like the ones of Theorem 3. However, this is not as simple as when evident scaling properties are present. We may indeed write
$$𝒞_F^{(1)}=\underset{V}{sup}\frac{F(\lambda _1(V))}{_^dG(V(x))𝑑x}1,$$
where the above supremum is taken on an appropriate space. Assuming that $`F`$ is nonincreasing, we may characterize $`𝒞_F^{(1)}`$ as
$$𝒞_F^{(1)}=\underset{\begin{array}{c}V,\varphi \\ _^d|\varphi |^2𝑑x=1\end{array}}{sup}\frac{F\left({\displaystyle _^d}\left(|\varphi |^2+V|\varphi |^2\right)𝑑x\right)}{_^dG(V(x))𝑑x},$$
so that the optimal value is at least formally given by
$$𝒞_F^{(1)}=\underset{\begin{array}{c}\varphi H^1(^d)\\ _^d|\varphi |^2𝑑x=1\end{array}}{sup}\frac{F\left({\displaystyle _^d}\left(|\varphi |^2+|\varphi |^2(G^{})^1(\kappa |\varphi |^2)\right)𝑑x\right)}{_^d\left(G(G^{})^1\right)(\kappa |\varphi |^2)𝑑x}$$
where $`\kappa `$ is given in terms of $`\varphi `$ by
$$\kappa =\left(𝒞_F^{(1)}\right)^1F^{}\left(_^d\left(|\varphi |^2+|\varphi |^2(G^{})^1(\kappa |\varphi |^2)\right)𝑑x\right).$$
This indeed results of the optimization with respect to $`V`$, which amounts to
$$\kappa |\varphi |^2G^{}(V)=0.$$
This strategy is however easy to implement in one more case: $`F(s)=e^s`$. In this case,
$$𝒞_F^{(1)}=\underset{\begin{array}{c}V,\varphi \\ _^d|\varphi |^2𝑑x=1\end{array}}{sup}\frac{e^{_^d(|\varphi |^2+V|\varphi |^2)𝑑x}}{(4\pi )^{d/2}_^de^V𝑑x}.$$
The optimization with respect to $`V`$ gives
$$V=\mathrm{log}(|\varphi |^2)$$
up to an additive constant such that $`_^de^V𝑑x=_^d|\varphi |^2𝑑x=1`$, which plays no role because its contribution to the numerator and to the denominator cancel. Summing up the inequality is therefore simply equivalent to the usual logarithmic Sobolev inequality: for any $`\varphi H^1(^d)`$ such that $`_^d|\varphi |^2𝑑x=1`$,
$$_^d|\varphi |^2\mathrm{log}(|\varphi |^2)𝑑x+\mathrm{log}\left(\frac{(4\pi )^{d/2}}{𝒞_F^{(1)}}\right)_^d|\varphi |^2𝑑x.$$
From standard results on logarithmic Sobolev inequalities, see for instance , it is known that optimal functions $`\varphi `$ are gaussian, which allows to determine the value of $`𝒞_F^{(1)}`$:
$$𝒞_F^{(1)}=\left(\frac{2}{e}\right)^d.$$
We will see later an alternative approach which allows to state the following interpolation inequality.
###### Proposition 13
Under the assumptions of Theorem 3, if $`F`$ and $`G`$ are related by (8), if $`F^{}`$ and $`G^{}`$ are invertible and if we define
$$\beta (s):=_0^s(F^{})^1(t)𝑑t\text{and}H(s):=_s^0(G^{})^1(t)𝑑t,$$
then for any $`\varphi H^1(^d)`$, the following interpolation inequality holds:
$$_^d|\varphi |^2𝑑x+\beta \left(_^d|\varphi |^2𝑑x\right)_^dH(|\varphi |^2)𝑑x.$$
This result will appear as a simple consequence of Theorem 15, where we take $`\nu _1=_^d|\varphi |^2𝑑x`$ and $`\nu _i=0`$ for any $`i2`$. We will see that the result holds not only in the framework of Theorem 3 but also in the case where $`lim_i\mathrm{}\lambda _i(V)<\mathrm{}`$ as it is the case for standard Lieb-Thirring inequalities.
### 4.4 Further results
The analogue of the Lieb-Thirring conjecture does not hold in the context of Theorem 1, i.e. for potentials such that $`lim_i\mathrm{}\lambda _i(V)=+\mathrm{}`$.
###### Proposition 14
With the notations of Sections 1 and 4.2, for any $`d^{}`$ and $`\gamma >d/2`$,
$$𝒞^{(1)}(\gamma )<𝒞(\gamma ).$$
Moreover, if $`F`$ and $`G`$ satisfy the assumptions of Theorem 3, then
$$n\underset{V}{sup}\frac{_{1in}F(\lambda _i(V))}{_^dG(V(x))𝑑x}=:𝒞^{(n)}(\gamma )$$
forms a strictly increasing sequence.
Proof. The infimum $`𝒞^{(1)}(\gamma )`$ is achieved by a function $`u_{}`$ with support in a ball $`B(0,R)`$ for some $`R>0`$, and a potential $`V_{}=cu_{}^{4/(d2(\gamma +1))}`$ in $`B(0,R)`$ for some constant $`c>0`$, and $`V_{}=+\mathrm{}`$ outside. The sequence of eigenvalues of $`\mathrm{\Delta }+V_{}`$ is therefore given by the one of $`\mathrm{\Delta }+V_{}`$ in $`B(0,R)`$ with zero Dirichlet boundary conditions on $`B(0,R)`$. It is then straightforward to realize that
$$\underset{i^{}}{}\left[\lambda _i(V_{})\right]^\gamma >\left[\lambda _1(V_{})\right]^\gamma =𝒞^{(1)}(\gamma ).$$
The general case follows from similar reasons. $`\mathrm{}`$
## 5 Interpolation inequalities
Assume that $`V`$ is a potential on $`^d`$ such that the operator $`\mathrm{\Delta }+V`$ has an infinite sequence $`(\lambda _i(V))_i^{}`$ of eigenvalues. Let $`F`$ and $`G`$ be two functions such that the inequality
$$\underset{i^{}}{}F(\lambda _i(V))=\mathrm{Tr}\left[F\left(\mathrm{\Delta }+V\right)\right]_^dG(V(x))𝑑x$$
(15)
holds (see for instance Theorem 3 for sufficient conditions). Let $`\overline{\lambda }:=lim_i\mathrm{}\lambda _i(V)`$ and assume that
$$\mathrm{Spectrum}(\mathrm{\Delta }+V)(\mathrm{},\overline{\lambda })=\{\lambda _i(V):i^{}\}.$$
Note that this includes the standard case of Lieb-Thirring inequalities, which corresponds to $`\overline{\lambda }=0`$ when $`V`$ is such that $`\mathrm{\Delta }+V`$ has infinitely many eigenvalues, and the case considered in Theorems 1 and 3: $`\overline{\lambda }=+\mathrm{}`$.
Define $`\sigma (s):=F^{}(s)`$ and $`\beta (s):=_0^s\sigma ^1(t)𝑑t`$. We may notice that
$$F(s)=_s^{\overline{\lambda }}\sigma (t)𝑑t=_s^{\overline{\lambda }}(\beta ^{})^1(t)𝑑t.$$
We assume that $`F`$ is convex on $`(\mathrm{},\overline{\lambda })`$ and $`C^1`$ on $`(\mathrm{},\overline{\lambda })`$ whenever it takes finite values. This implies that $`\beta `$ is $`C^1`$, convex and we get
$$F(s)=\underset{\nu >0}{\mathrm{min}}\left[\beta (\nu )+\nu s\right].$$
Note indeed that, at a formal level,
$$\frac{d}{ds}\left(\left[\beta (\nu )+\nu s\right]_{|\nu =(\beta ^{})^1(s)}\right)=(\beta ^{})^1(s)=\sigma (s).$$
Inequality (15) can therefore be rewritten as
$$\underset{i^{}}{}\nu _i_^d\left(|\psi _i|^2+V|\psi _i|^2\right)𝑑x+\underset{i^{}}{}\beta (\nu _i)+_^dG(V(x))𝑑x0$$
for any sequence of nonnegative occupation numbers $`(\nu _i)_i^{}`$ and any sequence $`(\psi _i)_i^{}`$ of orthonormal $`L^2(^d)`$ functions.
Let us proceed as in Section 4 and optimize on $`V`$ for fixed $`𝝂=(\nu _i)_i^{}`$, $`𝝍=(\psi _i)_i^{}`$. Assume further that $`G^{}`$ is invertible. Let
$$K[𝝂,𝝍]:=\underset{i^{}}{}\nu _i|\psi _i|^2\text{and}\rho :=\underset{i^{}}{}\nu _i|\psi _i|^2,$$
and define
$$H(s):=\left[G(G^{})^1(s)+s(G^{})^1(s)\right].$$
It is straightforward to check as above that
$$\frac{dH}{ds}(s)=(G^{})^1(s),$$
and write
$$H(s)=_s^0(G^{})^1(t)𝑑t$$
provided $`(G^{})^1`$ is integrable on a neighborhood of $`s=0_+`$.
The optimal potential $`V`$ has to satisfy
$$G^{}(V)+\rho =0,$$
so that
$$\underset{i^{}}{}\nu _i_^dV|\psi _i|^2𝑑x+_^dG(V(x))𝑑x=_^dH(\rho (x))𝑑x$$
Summarizing our computations, we have proved that (15) can be rephrased as
###### Theorem 15
Under the above notations and assumptions, the following inequality holds:
$$K[𝝂,𝝍]+\underset{i^{}}{}\beta (\nu _i)_^dH(\rho )𝑑x$$
(16)
with $`\rho =_i^{}\nu _i|\psi _i|^2`$, where $`(\nu _i)_i^{}`$ is any nonnegative sequence of occupation numbers and $`(\psi _i)_i^{}`$ is any sequence of orthonormal $`L^2(^d)`$ functions.
Written with such a generality, the result is maybe not as striking as when it applies to the various examples of Section 3, for which all the assumptions made above can be verified. To keep the generality of our result, we will not try to give sufficient conditions on $`\beta `$ and $`V`$ for which all these assumptions can be established and prefer to state three applications corresponding for the function $`\beta `$ to $`\beta (\nu )=Const\nu ^m`$ with $`m(\mathrm{},0)(1,+\mathrm{})`$, $`\beta (\nu )=Const\nu ^m`$ with $`m(0,1)`$ and $`\beta (\nu )=\nu \mathrm{log}\nu \nu `$.
Example 1. Let $`m>1`$, which corresponds to the setting of the standard Lieb-Thirring inequality (1), and consider $`\beta (\nu ):=c_m\nu ^m`$, $`c_m:=(m1)^{m1}m^m`$, $`m=\frac{\gamma }{\gamma 1}`$, $`F(s)=(s)^\gamma `$ and $`G(s)=C_{\mathrm{LT}}(\gamma )(s)^{\gamma +d/2}`$. Define
$$q:=\frac{2\gamma +d}{2\gamma +d2}\text{and}𝒦^1:=q\left[C_{\mathrm{LT}}(\gamma )\left(\gamma +\frac{d}{2}\right)\right]^{q1}.$$
###### Corollary 16
With the above notations, for any $`m(1,+\mathrm{})`$, the following optimal inequality holds:
$$K[𝝂,𝝍]+c_m\underset{i^{}}{}\nu _i^m𝒦_^d\rho ^q𝑑x.$$
Using the scaling invariance, we can reformulate this result as follows. If we replace $`\psi _i(x)`$ by $`\lambda ^{d/2}\psi _i(x/\lambda )`$ and $`\nu _i`$ by $`\lambda ^{d(11/q)}\nu _i`$, the right hand side of the above inequality is invariant. An optimization of the left hand side shows that
$$\left(K[𝝂,𝝍]\right)^\theta \left(\underset{i^{}}{}\nu _i^m\right)^{(1\theta )}_^d\rho ^q𝑑x,$$
where $`\theta =\frac{d}{2(\gamma 1)+d}`$ and $``$ can be explicitly computed in terms of $`𝒦`$, $`d`$ and $`\gamma `$.
The case $`m=\frac{\gamma }{\gamma 1}(\mathrm{},0)`$, which corresponds to $`\gamma (0,1)`$ and $`\beta (\nu ):=c_m\nu ^m`$, $`c_m:=(1m)^{m1}|m|^m`$ is formally covered with the same constants, but does not enter in our framework for infinite systems (see Example 1, Section 3.4). Notice that $`q`$ varies in the range $`(1,1+\frac{d}{2})`$ for $`m>1`$ and $`(1+\frac{d}{2},\frac{d}{d2})`$ if $`m<0`$. The case $`\gamma =1`$, $`q=1+\frac{d}{2}`$ is not covered.
Example 2. If $`m(0,1)`$, which corresponds to the setting of Theorem 1, and $`\beta (\nu ):=c_m\nu ^m`$, $`c_m:=(1m)^{m1}m^m`$, $`m=\frac{\gamma }{\gamma +1}`$, $`F(\lambda )=\lambda ^\gamma `$ and $`G(s)=𝒞(\gamma )s^{d/2\gamma }`$. Define
$$q:=\frac{2\gamma d}{2(\gamma +1)d}(0,1)\text{and}𝒦^1:=q\left[𝒞(\gamma )\left(\gamma \frac{d}{2}\right)\right]^{q1}.$$
Notice that due to the restriction $`\gamma >d/2`$, the range of $`m`$ is reduced to the interval $`(\frac{d}{d+2},1)`$.
###### Corollary 17
With the above notations, for any $`m(\frac{d}{d+2},1)`$, the following optimal inequality holds:
$$K[𝝂,𝝍]+𝒦_^d\rho ^q𝑑xc_m\underset{i^{}}{}\nu _i^m.$$
Using the scaling invariance, we can also reformulate this result as follows. If we replace $`\psi _i(x)`$ by $`\lambda ^{d/2}\psi _i(x/\lambda )`$ but dont change $`\nu _i`$, the right hand side of the above inequality is of course invariant. An optimization of the left hand side shows that
$$\left(K[𝝂,𝝍]\right)^\theta \left(_^d\rho ^q𝑑x\right)^{(1\theta )}\underset{i^{}}{}\nu _i^m,$$
where $`\theta =\frac{d}{2(\gamma +1)}`$ and $``$ can be explicitly computed in terms of $`𝒦`$, $`d`$ and $`\gamma `$.
Example 3. If $`\beta (\nu ):=\nu \mathrm{log}\nu \nu `$, then $`\beta ^{}(\nu )=\mathrm{log}\nu =\lambda `$, $`F(s)=e^s`$ and $`G(s)=(4\pi )^{d/2}e^s`$. Inequality (16) is a logarithmic Sobolev inequality for systems:
###### Corollary 18
With the above notations, the following optimal inequality holds:
$$K[𝝂,𝝍]+\underset{i^{}}{}\nu _i\mathrm{log}\nu _i_^d\rho \mathrm{log}\rho dx+\frac{d}{2}\mathrm{log}(4\pi )_^d\rho 𝑑x.$$
As above, an optimization under a scaling preserving the $`L^2`$ norm of $`\psi `$ and leaving $`\nu _i`$ invariant allows to write
$$_^d\rho \mathrm{log}\rho dx\underset{i^{}}{}\nu _i\mathrm{log}\nu _i+\frac{d}{2}\mathrm{log}\left(\frac{e}{2\pi d}\frac{K[𝝂,𝝍]}{_^d\rho 𝑑x}\right)_^d\rho 𝑑x.$$
Note that we immediately recover the Gagliardo-Nirenberg inequalities of Section 3 by taking $`\nu _1=1`$, $`\nu _i=0`$ for any $`i2`$, in case of Examples 1 and 2, but with a priori non optimal constants, at least in the case of Example 2. The proof of Proposition 13 follows for the same reason.
Remark. We notice that the limit case $`\gamma =0`$ for $`d3`$ is not covered, even as a limit case. For $`\nu _i=1`$, for $`i=1`$, $`2`$,…$`N`$, and $`\nu _i=0`$ otherwise, a Sobolev type inequality for orthonormal functions has been given in in the case which corresponds to the critical Sobolev embedding $`H^1(^d)L^{2d/(d2)}(^d)`$. By taking the occupation numbers $`\nu _i`$ into account, we always achieve optimal inequalities which are related in a natural way to some corresponding optimal Lieb-Thirring inequalities as long as $`\gamma `$ is positive. To a large extend, this improves the known results for orthonormal and sub-orthonormal systems .
Acknowledgments. J.D. and E.P. are partially supported by ECOS-Conicyt under grants no. C02E06 and no. C02E08 and by European Programs HPRN-CT # 2002-00277 & 00282. M.L. is partially supported by U.S. National Science Foundation grant DMS 03-00349. P.F. is partially supported by ECOS-Conicyt under grant no. C02E08, Fondecyt Grant 1030929 and FONDAP de Matemáticas Aplicadas.
© 2005 by the authors. This paper may be reproduced, in its entirety, for non-commercial purposes.
|
warning/0506/hep-th0506150.html
|
ar5iv
|
text
|
# Acknowledgements
## Acknowledgements
We wish to thank Dan Freedman for discussions and helpful suggestions and Juan Maldacena for enlightening comments.
This work has been supported in part by INFN, PRIN prot. 2003023852\_008 and the European Commission RTN program MRTN–CT–2004–005104.
|
warning/0506/cond-mat0506362.html
|
ar5iv
|
text
|
# Critical temperature and giant isotope effect in presence of paramagnons.
## Abstract
We reconsider the long-standing problem of the effect of spin fluctuations on the critical temperature and isotope effect in a phonon-mediated superconductor. Although the general physics of the interplay between phonons and paramagnons had been rather well understood, the existing approximate formulas fail to describe the correct behavior of $`T_c`$ for general phonon and paramagnon spectra. Using a controllable approximation, we derive an analytical formula for $`T_c`$ which agrees well with exact numerical solutions of the Eliashberg equations for a broad range of parameters. Based on both numerical and analytical results, we predict a strong enhancement of the isotope effect when the frequencies of spin fluctuation and phonons are of the same order. This effect may have important consequences for near-magnetic superconductors such as MgCNi<sub>3</sub>.
In the last decade a large number of superconductors were discovered in which enhanced spin fluctuations (SF) play a role in the superconductivity, $`e.g.,`$ Sr<sub>2</sub>RuO$`_4,`$ MgCNi$`_3,`$ $`\epsilon `$-Fe, ZrZn$`_2,`$ and others, bringing about new and interesting physics. However, understanding such materials, even at an intuitive level, has been hindered by the lack of a simple formula that would approximate the full Eliashberg theory in a compact analytical form, as the conventional McMillan formula (MMF) does. As a result, uncritical generalizations of the latter have been used as a substitute, despite the fact that, as we will show below, some of them are too approximate or outright incorrect. In this Letter we present an analogue of the MMF, derived in a controllable way and tested against numerical solutions of full Eliashberg equations, including interaction with SF (paramagnons). We point out the possibility of a giant phonon isotope effect induced by SF. We will also apply this theory, as an example, to a nearly-ferromagnetic superconductor, MgCNi$`_3.`$
The understanding that SF are pair breakers in conventional superconductors is nearly as old as the BCS theory itself BSGJS . Moreover, it was soon realized that strong coupling manifests itself in a nontrivial way in the presence of SF Rietschel ; VIS . In a number of papers numerical solutions of the Eliashberg equations were presented, incorporating phonon $`\alpha ^2F_p(\omega )`$ as well as SF $`\alpha ^2F_s(\omega )`$ spectral functions (see, e.g., Refs. WCNC ). However, solving the full Eliashberg equation is not always an option, and does not provide as much physical insight as analytical treatment. An analytical tool comparable to the famed MMF is needed.
Retrospectively, one can realize that the overwhelming success of the MMF is due to three facts: (a) it can be derived analytically using simple approximations, (b) it includes Coulomb repulsion effects, (c) it has three universal adjustable parameters, which, after little tuning, produce an expression which is surprisingly accurate for a large range of phonon frequencies and coupling strengths. Compared to the BCS equation, the MMF includes three essential pieces of additional physics: effective mass renormalization, logarithmic reduction of the Coulomb repulsion, and proper (logarithmic) averaging of the phonon frequency. All three effects can be derived analytically in some approximations. In fact, it is known that the functional form of the MMF can be derived in two different ways. One, known as the square-well model AllenMitr , uses the Matsubara representation, where the coupling with the phonons is parametrized in terms of the matrix $`\lambda (n,n^{}).`$ The model assumes two different approximations for the same function $`\lambda (n,n^{}),`$ depending on whether it is used in the equation for the mass renormalization $`Z`$ or in the one for the gap function $`\varphi `$:
$`\lambda _Z(n,n^{})`$ $`=`$ $`\lambda _p\mathrm{\Theta }(\omega _p\left|\omega _{nn^{}}\right|)`$ (1)
$`\lambda _\varphi (n,n^{})`$ $`=`$ $`\lambda _p\mathrm{\Theta }(\omega _p\left|\omega _n\right|)\mathrm{\Theta }(\omega _p\left|\omega _n^{}\right|).`$
This models leads to an equation for the critical temperature, $`T_c,`$
$$T_c=a\omega _{\mathrm{log}}\mathrm{exp}\{b(1+\lambda _Z)/[\lambda _\varphi \mu ^{}(1+c\lambda _Z)],$$
(2)
where the theoretical parameters are $`a=1.14,`$ $`b=c=1`$, $`\lambda _Z=\lambda _\varphi =\lambda _p=2_0^{\mathrm{}}\omega ^1\alpha ^2F_p(\omega )𝑑\omega `$ and $`\lambda _p\mathrm{ln}\omega _{\mathrm{log}}=2_0^{\mathrm{}}\omega ^1\mathrm{ln}\omega \alpha ^2F_p(\omega )𝑑\omega `$. The renormalized Coulomb potential is reduced from its bare value $`\mu `$ as $`\mu ^{}=\mu /(1+\mu \mathrm{ln}\frac{\omega _C}{\omega _{\mathrm{log}}}),`$ where $`\omega _C`$ characterizes the frequency cutoff of the Coulomb interaction. The MMF formula is given by Eq.2 with optimized parameters $`a=1/1.2,`$ $`b=1.04,`$ and $`c=0.62.`$
SF, as opposed to phonons, induce repulsion for singlet pairs. However, they contribute to the mass renormalization just the same. Therefore the first instinct is to let $`\lambda _Z=\lambda _p+\lambda _s,`$ where $`\lambda _s`$ describes the SF, and $`\lambda _\varphi =\lambda _p\lambda _s.`$ Eq.2 with this modification and standard $`a,b`$ and $`c`$ is the one routinely used in the literature for materials with SF (e.g., Refs.MazinFe ; helge ; someJarlborg ).
Obviously, using two different approximations for the same physical function $`\lambda (n,n^{})`$ depending on whether it appears in the first or second Eliashberg equation cannot be justified by any logic. It appears that the MMF formula can be fortuitously derived in this way, but, as we will see below, this approach fails when SF are included. An alternative derivation of the MMF utilizes the real frequency axis formalism kmmProblema . The one-mode approximation is used, which assumes an Einstein phonon at a frequency $`\omega _p,`$ i.e., $`\alpha ^2F(\omega )=\lambda _p\omega _p\delta (\omega \omega _p)/2.`$ The Eliashberg equations are then solved iteratively. After the first iteration one obtains kmmProblema
$$T_c=1.14\omega _p\mathrm{exp}\left\{\frac{1}{2}\frac{1+\lambda _p}{\lambda _p\mu ^{}[1+0.5\frac{\lambda _p}{1+\lambda _p}]}\right\},$$
(3)
which is similar to the square well formula Eq.2 with $`a=1.14/\sqrt{e}=1/1.44`$ (note that this value of $`a`$ is much closer to the optimized one), $`b=1,`$ and $`c=0.5/(1+\lambda _p)`$. This approach is a controllable approximation with a concrete physical meaning. However, it has never been applied to superconductors with SF.
On the contrary, several attempts to apply the square well model to SF have been reported. In Refs. VIS ; WCNC the following expression was derived (for $`\mu ^{}=0`$):
$`T_c`$ $`=`$ $`1.14\omega _p^\nu \omega _s^{1\nu }\mathrm{exp}\{(1+\lambda _p+\lambda _s)/(\lambda _p\lambda _s)],`$ (4)
$`\mathrm{with}`$ $`\nu `$ $`=\lambda _p/(\lambda _p\lambda _s)(a)`$
$`\mathrm{or}`$ $`\nu `$ $`={\displaystyle \frac{\lambda _p^2}{\lambda _p\lambda _s}}\left[\lambda _p\lambda _s+{\displaystyle \frac{\lambda _p\lambda _s}{1+\lambda _p+\lambda _s}}\mathrm{ln}{\displaystyle \frac{\omega _p}{\omega _s}}\right]^1(b)`$
where the choice (a) is due to Carbotte et al. WCNC , and (b) to Vonsovskii et al. VIS . Unfortunately, neither authors give details of their derivations, so we do not know what was different in their models. We were not able to reproduce either result. The latest paper utililizing the square well model (in the weak coupling limit) is that by Shimahara Shim . Our own result for the square well model reduces to that of Ref.Shim in the weak limit, and reads
$`T_c`$ $`=`$ $`1.14\omega _p\mathrm{exp}[{\displaystyle \frac{1+\lambda _s+\lambda _p}{\lambda _p\frac{\lambda _s(1+\lambda _s)}{1+\lambda _s+\lambda _s\mathrm{ln}\frac{\omega _s}{\omega _p}}}}],\omega _s\omega _p`$ (5)
$`T_c`$ $`=`$ $`1.14\omega _s\mathrm{exp}[{\displaystyle \frac{1+\lambda _s+\lambda _p}{\frac{\lambda _p(1+\lambda _p)}{1+\lambda _p\lambda _p\mathrm{ln}\frac{\omega _p}{\omega _s}}\lambda _s}}],\omega _s\omega _p`$ (6)
Unlike Eq.4, Eqs.5, 6 reduce to the McMillan form upon substitution $`\omega _s\omega _C\omega _p,`$ $`\lambda _s\mu ,`$ as it should.
Given the controversy about the square-well model, it is desirable to have a derivation in a controllable approximation, such as the real frequency axis formalism of Refs.kmmProblema . Assuming an Einstein phonon at a frequency $`\omega _p`$ and an “Einstein” paramagnon at $`\omega _s`$, $`2\alpha ^2F(\omega )=\lambda _p\omega _p\delta (\omega \omega _p)\lambda _s\omega _s\delta (\omega \omega _s),`$ we obtain the following iterative solution of the Eliashberg equations.
$`T_c`$ $`=`$ $`1.14\omega _p^{\frac{\lambda _p}{\lambda _p\lambda _s}}\omega _s^{\frac{\lambda _s}{\lambda _p\lambda _s}}\mathrm{exp}(K)`$ (7)
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{1+\lambda _p+\lambda _s}{\lambda _p\lambda _s\mu ^{}(1K\frac{\lambda _p\lambda _s}{1+\lambda _p+\lambda _s})}}\right\}`$
$`K`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _p\lambda _s}{\left(\lambda _p\lambda _s\right)^2}}\left[1+{\displaystyle \frac{\omega _p^2+\omega _s^2}{\omega _p^2\omega _s^2}}\mathrm{ln}{\displaystyle \frac{\omega _s}{\omega _p}}\right].`$
For $`\omega _p\omega _s`$, $`K=1/2,`$ and at $`\mu ^{}=0`$, Eq.7 reduces to Eq.4 with $`\nu =\lambda _p/(\lambda _p\lambda _s)`$.
As usual, the ultimate test for any approximation is numerical calculations. We solved the Eliashberg equations for a variety of model $`\alpha ^2F(\omega )`$ including SF and compare them with the proposed analytical formulas. In Fig.1 we show this comparison for the simplest ”one-mode” approximation, one phonon and one paramagnon (we have verified that other model spectra lead to similar results). As we can see, while the Eq.7, as well as its simplified version Eq.4(a), describe the numerical results rather well when $`\omega _s`$ and $`\omega _p`$ are comparable, the latter fails at $`\omega _s\omega _p,`$ and both fail at $`\omega _s\omega _p.`$ Both effects can be easily understood: Eq.4 includes $`\omega _s`$ in a negative power in all regimes, thus leading to a total suppression of superconductivity at $`\omega _s\mathrm{}.`$ In reality, in this limit the negative effect of the SF is renormalized down logarithmically in the same spirit as the Coulomb repulsion.
Eqs. 7, 4 diverge at $`\omega _s0`$. This is due to the fact that the derivations above assume that $`\omega _s,\omega _p\pi T_c`$. It is possible to treat this regime separately. If $`\omega _sT_c,`$ the SF act as static magnetic defects, and the standard theory of the magnetic pair-breakingAG can be applied. In the Matsubara representation, at $`\omega _s=0`$ one needs only to keep the term with $`n=m`$ in $`\lambda _s(\omega _n\omega _m)`$. Then the equations reduce to the standard form AG ; AllenMitr with the pair-breaking parameter $`\gamma (1/2\tau _P)/\pi T_c=\lambda _s`$. In the weak coupling limit, $`T_c`$ is
$$T_c=T_{c0}\mathrm{exp}[\psi (1/2)\psi (1/2+\gamma )],$$
(8)
where $`T_{c0}=T_c(\lambda _s=0).`$ One important difference exists between pair-breaking by SF with $`\omega _s=0`$ and by magnetic impurities: in the former case the pair-breaking parameter $`\gamma `$ now does not depend on $`T_c`$. This has consequences for the isotope effect, as we will see below.
For small, but finite $`\omega _s\pi T_c`$ summation of $`\lambda _s(nm)`$ over $`nm`$ provides the expression for the pair breaking rate in Eq.8: $`\gamma =\lambda _s\frac{T_c}{2\omega _s}\mathrm{coth}\frac{T_c}{2\omega _s}`$. This result coincides with Eq. 5.8 of Ref.Millis for dynamical pair breaking in anisotropic superconductors if the anisotropy parameter $`g`$ (as defined in Ref. Millis ) is set to -1. When $`\omega _s`$ increases, $`T_c`$ drops sharply with a complete loss of superconductivity at $`\omega _s=\omega _s^{}=e^CT_{c0}/2\gamma `$ (where $`e^C1.78`$). However, the condition $`\omega _s\pi T_c`$ used in the derivation of Eq.8 is lost well before $`\omega _s^{}`$ (in fact, at $`\omega _s\omega _s^{}/2)`$.
One can take into account the strong coupling effects in the square-well model, resulting in a renormalization $`\gamma \gamma /(1+\lambda _p)=\frac{\lambda _s}{1+\lambda _p}\frac{T_c}{2\omega _s}\mathrm{coth}\frac{T_c}{2\omega _s}.`$ As the comparison with numerical calculations shows (Fig. 1), this approximation underestimates $`T_c`$. However, it illustrates why $`T_c`$ flattens out at a finite value smaller than $`T_{c0}`$, when $`\omega _s0`$, instead of raising as Eq.7 suggests.
We also show in Fig.1 that both Eq.4(b) and the square well model, Eqs.5 and 6, disagree qualitatively with the numerical results in the whole range of $`\omega _s`$.
We will now turn to the isotope effect. Looking at Eq.4, one observes that the isotope coefficient, $`\beta =\nu /2=\lambda _p/2(\lambda _p\lambda _s)>0.5,`$ is always enhanced compared to its BCS value and is independent of the SF frequencies. Clearly, this should hold approximately in the range of the applicability of this formula, $`\omega _s\omega _p\pi T_c.`$ Indeed, the more accurate Eq.7 yields for $`\beta `$
$`\beta `$ $`=`$ $`0.5{\displaystyle \frac{\lambda _p}{\lambda _p\lambda _s}}\left[1{\displaystyle \frac{\lambda _s}{\lambda _p\lambda _s}}F\left({\displaystyle \frac{\omega _s^2}{\omega _p^2}}\right)\right]`$ (9)
$`F(r)`$ $`=`$ $`(r^22r\mathrm{ln}r1)/(r1)^2.`$
The second term here is the correction to Eq.4. It can be of either sign, since with growing $`r`$ the $`F(r)`$ monotonically grows from -1 to 1, and $`F(1)=0.`$ As discussed, Eq.4 itself becomes invalid at $`\omega _s<\pi T_c.`$ As $`\omega _s0`$, according to Eq.8, $`\beta =0.5`$ (note that in the case of magnetic impurities $`\beta >0.5`$ due to the dependence of $`\gamma `$ on $`T_c`$ AllenMitr ). Therefore, the isotope effect has to have a maximum at some $`0<\omega _s<\omega _p,`$ and $`\beta _{\mathrm{max}}>\lambda _p/2(\lambda _p\lambda _s).`$
This is confirmed by numerical calculations, which do show that the maximum isotope effect for given $`\lambda _s,\lambda _p`$ is achieved close to $`\omega _s\omega _p`$ and is not far from $`\lambda _p/2(\lambda _p\lambda _s).`$ This is a very import result, and we emphasize it again: if superconductivity is depressed by spin-fluctuations, the total isotope effect increases compared to its BCS value.
We shall now apply this formalism to a superconductor where $`T_c`$ is believed to be substantially suppressed by SF, MgCNi<sub>3</sub> we ; kitaicy ; helge , which has attracted substantial interest not because of its relatively modest critical temperature, $`T_c8`$ K, but because of its unusual antiperovskite crystal structure and proximity to ferromagnetic instability. The latter was first pointed out by Rosner et al. helge , who believed in such strong coupling with SF that they proposed a $`p`$ wave superconductivity. Singh and Mazin we also came to the conclusion that SF should play a role in superconductivity of MgCNi$`_3,`$ but, based on their frozen phonon calculation, they deduced a large electron-phonon coupling constant ($`\lambda _p1)`$ due to the bond-bending Ni phonons. They reconciled this relatively large $`\lambda _p`$ with a modest $`T_c`$ within a scenario of s-wave phonon-induced superconductivity depressed by SF. Later this scenario was re-invented by Shan et al. kitaicy , who proved the s-symmetry of the order parameter by tunneling experiments. This point has been since confirmed by several groups and seems to be well established.
Singh and Mazin’s we prediction of the Ni phonon playing the major role in the electron-phonon coupling in MgNiC<sub>3</sub> was based on a limited number of calculations at a high-symmetry point in the Brillouin zone, and therefore was more an educated guess than a quantitative argument. A quantitative analysis was provided by Ignatov et al. serega , who performed linear-response calculations of the phonon frequencies and their coupling with electrons for the whole Brillouin zone. They found a gigantic coupling for the Ni bond-bending modes, and the most strongly coupled modes (the mode considered by Singh and Mazin was not among them) actually unstable. In other words, they found a set of double-well type instabilities involving mostly Ni atoms. This was verified by EXAFS measurements serega . Ignatov et al. serega estimated the total electron-phonon coupling constant as 1.5 and the logarithmically averaged phonon frequency as 131 K.
Thus, the scenario of Ref. we was modified in Ref. serega in the sense that electron-phonon coupling and superconductivity were coming from highly anharmonic predominantly Ni modes, but not exactly the simple rotations of the Ni<sub>6</sub> octahedra considered in Ref. we . Unfortunately, strong anharmonicity of these modes makes it impossible to evaluate their coupling with electrons in the linear response calculations, but it is obviously strong. However, one can estimate the electron-phonon and electron-paramagnon coupling indirectly from experimental data. Indeed, specific heat renormalization, from different reports, ranges from 2.6 to 3.1 (see Ref. Dresden and refs. therein), implying that the sum $`\lambda _p+\lambda _s`$ varies between 1.6 and 2.1. Wälte et al. Dresden estimated $`\omega _p143`$ K, smaller than, but comparable to the calculation in Ref.serega , $`\omega _s25`$ K and the mass renormalization due to paramagnons as $`1+\lambda _s1.43.`$ Then, using MMF, $`\mu ^{}=0.13,`$ and $`T_c=6.8`$ K, as measured for their samples, they deduced $`\lambda _p`$=1.91.
However, there are several problems with this derivation. First of all, as shown above, the proper formula is Eq.4. Using this formula instead of Eq.2, and keeping all their other parameters, we get a much more reasonable number, $`\lambda _p`$=1.61, not far from the value of 1.51 obtained in Ref. serega . However, the SF model adopted in Ref. Dresden cannot be considered as proven. It is based on the disputable assumption that the upturn of the specific heat quotient at low temperature and high magnetic field is due to the paramagnon contribution to specific heat, but there many other explanations of this effect. 25 K seems to be unrealistically soft. Also, low $`T_c`$ and high residual resistance cast doubt on the sample quality in this study.
Here we adopt a different approach: we adopt the calculated values $`\lambda _p`$=1.5 and $`\omega _p=131`$ K, in the harmonic approximation, and total mass renormalization $`1+\lambda _p+\lambda _s=2.85,`$ so that $`\lambda _s=0.35.`$ The results of the numerical solution of the Eliashberg equations with the $`\alpha ^2F(\omega )`$ function calculated by Ignatov et al.serega and $`\mu ^{}=0.12`$ are shown in Fig .2, together with the curve calculated from Eq.7. This way, we find $`\omega _s50`$ K, which, we believe, is a more realistic number than 25 K. The corresponding total isotope effect coefficient is 0.75.
This may sound in agreement with the recent experiment by Klimczuk and Cava Cava , who have measured the isotope effect to be 0.54 on carbon only. If the total isotope effect is 0.75, this suggests a seemingly reasonable Ni isotope effect of 0.21, suggesting that Ni phonons couple with the electrons twice weaker than C ones. Unfortunately, the first-principles calculations suggest that the Ni modes couple with electrons at least an order of magnitude stronger than the C modes (there is hardly any C character present at the Fermi level). In the moment, the only way to reconcile this with the measurements of Ref. Cava is to assume that the observed isotope effect is not a result of the frequency shift of the C modes, but of some subtle changes in the crystal structure induced by the isotope substitution. Such a possibility is suggested by an earlier study Cavaold , where it was found that (i) $`T_c`$ depends on the lattice parameter at a rate of $`310`$ K/Å, which translates an error of $`\pm 0.0015`$ Å in the lattice parameterCava into an error of $`\pm 0.46`$ K in $`T_c,`$ larger than the isotope shift of 0.3 K, and (ii) that two samples with the same lattice parameter and the same neutron-measured C content have $`T_c`$ differing by 0.71 K. A possible explanation is that, given the proximity of MgCNi<sub>3</sub> to a ferromagnetic instability, crystallographic defects may induce local magnetic moments which, in turn, work as pair-breakers. The concentration of such defects, even for the same net C content, may depend on the sample preparation and, possibly, on isotope substitution.
Therefore further studies of the isotope effect both on C and on Ni are necessary, in particular combined with accurate measurements of the isotope shift of the phonon modes.
We acknowledge support from the NWO-RFBR grant 047.016.005 and from the NSF DMR grants 0342290 and 023188.
|
warning/0506/math0506569.html
|
ar5iv
|
text
|
# Caculus of Variation and the 𝐿²-Bergman Metric on Teichmüller Space
## 1 Introduction
We present a geometric analytic approach to the $`L^2`$-Bergman metric on Teichmüller space of Riemann surfaces in this paper.
Teichmüller space $`𝒯_g`$ is the space of conformal structures on a compact, smooth, oriented, closed Riemann surface $`\mathrm{\Sigma }`$ of genus $`g1`$, where two conformal structures $`\sigma `$ and $`\rho `$ are equivalent if there is a biholomorphic map between $`(\mathrm{\Sigma },\sigma )`$ and $`(\mathrm{\Sigma },\rho )`$ in the homotopy class of the identity map. When $`g2`$, Teichmüller space $`𝒯_g`$ is naturally a complex manifold of complex dimension $`3g3>1`$, and the cotangent space at $`\mathrm{\Sigma }`$ is identified with $`QD(\mathrm{\Sigma })`$, the space of holomorphic quadratic differentials. ().
Since $`𝒯_g`$ is a complex manifold, it is natural to study its metric geometry. There are several interesting metrics defined on Teichmüller space, all have advantages and disadvantages. These metrics reflect different perspectives of Teichmüller space. Among those metrics, two are named after S. Bergman. One of which, we still call it the Bergman metric, comes from the Bergman kernel function of a complex manifold: as a bounded complex domain $`𝒯_g`$ carries an invariant Kählerian Bergman metric, defined by the line element
$`ds^2=\frac{^2logK(z,z)}{z\overline{z}}dz_id\overline{z}_j`$,
where $`K(z,\zeta )`$ is the Bergman kernel. This Bergman metric is complete ().
The main object of this paper is the other metric which sometimes bears the name of Bergman. It is a Weil-Petersson type metric on Teichmüller space, i.e., it is obtained from duality by a $`L^2`$ inner product. In order to distinguish from the first Bergman metric, it may be appropriate to call this metric the $`L^2`$-Bergman metric.
From the classical Riemann surface theory, the period map $`p:\mathrm{\Sigma }J_\mathrm{\Sigma }`$ embeds the surface $`\mathrm{\Sigma }`$ to its Jacobian $`J_\mathrm{\Sigma }`$. The pullback metric of the flat metric on $`J_\mathrm{\Sigma }`$ via this period map thus defines the so-called canonical metric or Bergman metric on $`\mathrm{\Sigma }`$, denoted by $`\rho _B`$. This metric $`\rho _B`$ is of nonpositive Gaussian curvature, and when $`g2`$, the curvature vanishes if and only if the surface is hyperelliptic and only at $`2g+2`$ Weierstrass points ( ), in other words, the Gaussian curvatures characterize hyperelliptic surfaces. There is a unique canonical metric in every conformal structure.
The induced $`L^2`$-Bergman cometric is defined on $`QD(\mathrm{\Sigma })`$ by $`L^2`$-norm
$`\varphi _B^2=_\mathrm{\Sigma }\frac{|\varphi |^2}{\rho _B}`$,
thus we obtain a metric on Teichmüller space by duality. This is a Riemannian, Hermitian metric, invariant under the mapping class group.
This metric has been studied by Haberman and Jost who showed that it is incomplete (). Roughly speaking, with respect to the $`L^2`$-Bergman metric, boundary points of the moduli space $`_g`$ corresponding to pinching a nonseparating curve on the surface are at infinite distance from the interior, while boundary points of $`_g`$ corresponding to pinching a separating curve on the surface are at finite distance from the interior. In a sense, the $`L^2`$-Bergman metric detects topology of the surface.
One of the motivations of this study is to compare the $`L^2`$-Bergman metric with more intensively studied Weil-Petersson metric. These two metrics are both defined from duality from $`L^2`$ inner products, and they are both incomplete. However, the $`L^2`$-Bergman metric does not depend on the uniformization theorem. The difference between hyperbolic metric (constant curvature $`1`$) and the canonical metric on the surface results in different behavior of the induced $`L^2`$ metrics on Teichmüller space. The Weil-Petersson metric is of negative curvature (, ), we are yet to understand the curvature properties of the $`L^2`$-Bergman metric.
In this paper, we take a variational approach to the study of the $`L^2`$-Bergman metric. To do so, we fix a conformal structure $`(\sigma ,z)`$ with conformal coordinates $`z`$. For each canonical metric $`\rho `$ on the surface $`\mathrm{\Sigma }`$, one obtains a quadratic differential $`\varphi (z)dz^2`$ which is the Hopf differential of the unique harmonic map from $`\sigma `$ to $`\rho `$. This quadratic differential is holomorphic with respect to the conformal structure $`(\sigma ,z)`$, therefore an element of the space $`QD(\mathrm{\Sigma })`$. We thus obtain a map $`\varphi `$ between Teichmüller space $`𝒯_g`$ and $`QD(\mathrm{\Sigma })`$, sending $`\rho `$ to $`\varphi (z)dz^2`$. We show that this map is a global homeomorphism, hence it provides global coordinates to $`𝒯_g`$. The following theorem is an analog to Wolf’s theorem in the case of hyperbolic metrics ().
###### Theorem 1.1.
The map $`\varphi :𝒯_gQD(\mathrm{\Sigma })`$ is a homemorphism.
We note that, in the case of hyperbolic metrics, the injectivity of the map $`\varphi `$ is a direct application of Bochner’s identities and maximum principle, as seen in , relying on the fact that hyperbolic metric is of constant curvature $`1`$. In the case of canonical metrics, this is rather difficult since canonical metric has varied curvatures.
With the homeomorphism theorem in mind, we then consider a family of harmonic maps between canonical metrics on a surface and show that the second variation of an energy functional is the $`L^2`$-Bergman metric of two infinitesimal cotangent vectors on Teichmüller space. In the case of varying target metrics, we find:
###### Theorem 1.2.
Let $`w(t):(\mathrm{\Sigma },\sigma (z)|dz|^2)(\mathrm{\Sigma },\rho (t)|dw|^2)`$ be a family of harmonic maps between canonical metrics on surface $`\mathrm{\Sigma }`$, where $`\rho (0)=\sigma `$, for $`|t|<ϵ`$ small. Then the second variation of the energy functional of $`w(t)`$, at $`t=0`$, is given by the $`L^2`$-Bergman metric of infinitesimal holomorphic quadratic differentials (up to a constant).
Similar result holds in the case of varying domain metrics:
###### Theorem 1.3.
Let $`w(s):(\mathrm{\Sigma },\sigma (s))(\mathrm{\Sigma },\rho )`$ be a family of harmonic maps between canonical metrics on surface $`\mathrm{\Sigma }`$, where $`\sigma (0)=\rho `$, for $`|s|<ϵ`$ small. Then the second variation of the energy functional of $`w(t)`$, at $`t=0`$, is given by the $`L^2`$-Bergman metric of infinitesimal holomorphic quadratic differentials (up to a constant).
This paper is organized as follows. We introduce the preliminaries in section 2, then prove the homeomorphism theorem 1.1 in section 3. Section Four is devoted to a variational approach to the study of the $`L^2`$-Bergman metric, where we prove theorem 1.2 (varying the target metric) in $`\mathrm{\S }4.1`$ and theorem 1.3 (varying the domain metric) in $`\mathrm{\S }4.2`$.
The author owes a great debt to, and wishes to thank, Xiaodong Wang and Mike Wolf for helpful discussions over the topic.
## 2 Preliminaries
On a compact Riemann surface $`\mathrm{\Sigma }`$ of genus $`g>1`$, the dimension of the space of Abelian differentials of the first kind, or holomorphic one forms, is $`g`$. There is a natural pairing of Abelian differentials defined on this space:
$`<\mu ,\nu >=\frac{\sqrt{1}}{2}_\mathrm{\Sigma }\mu \overline{\nu }`$
Let $`\{\omega _1,\omega _2,\mathrm{},\omega _g\}`$ be a basis of Abelian differentials, normalized with respect to the $`A`$-cycles of some symplectic homology basis $`\{A_i,B_i\}_{1ig}`$, i.e., $`_{A_i}\omega _j=\delta _{ij}`$. Thus the period matrix $`\mathrm{\Omega }_{ij}=_{B_i}\omega _j`$. One finds that, since not all Abelian differentials vanish at the same point according to Riemann-Roch, the period matrix is then symmetric with positive definite imaginary part: $`Im\mathrm{\Omega }_{ij}=<\omega _i,\omega _j>`$ ().
The canonical metric $`\rho _B`$ on surface $`\mathrm{\Sigma }`$ is the metric associated to the $`(1,1)`$ form given by
$`\frac{\sqrt{1}}{2}_{i,j=1}^g(Im\mathrm{\Omega })_{ij}^1\omega _i(z)\overline{\omega }_j(\overline{z})`$.
It is not hard to see that this metric is the pull-back of the Euclidean metric from the Jacobian variety $`J(\mathrm{\Sigma })`$ via the period map ().
###### Remark 2.1.
It is easy to see that the area of the surface $`\mathrm{\Sigma }`$ with respect to the canonical metric is a constant, i.e., $`_\mathrm{\Sigma }\rho _B=g`$. Sometimes the canonical metric is also refered to $`\frac{\rho _B}{g}`$ to unify the surface area.
It is known that, when $`g2`$, the Gaussian curvature $`K_c`$ satisfies $`K_c0`$ (, ), and $`K_c(p)=0`$ for some $`p\mathrm{\Sigma }`$ if and only if $`\mathrm{\Sigma }`$ is hyperelliptic and $`p`$ is one of the $`2g+2`$ classical Weierstrass points of $`\mathrm{\Sigma }`$ ().
The Weil-Petersson cometric on Teichmüller space is defined on the space of holomorphic quadratic differentials $`QD(\mathrm{\Sigma })`$ by the $`L^2`$-norm:
$`\varphi _{WP}^2={\displaystyle _\mathrm{\Sigma }}{\displaystyle \frac{|\varphi |^2}{\sigma }}𝑑z𝑑\overline{z}`$ (1)
where $`\sigma |dz|^2`$ is the hyperbolic metric on $`\mathrm{\Sigma }`$. By duality, we obtain a Riemannian metric on the tangent space of $`𝒯_g`$.
The $`L^2`$-Bergman metric on $`𝒯_g`$ is similarly defined by duality from the $`L^2`$-norm
$`\varphi _B^2=_\mathrm{\Sigma }\frac{|\varphi |^2}{\rho _B}`$.
We now introduce harmonic maps between canonical metrics on a surface as much of our analysis will focus on the techniques of harmonic maps.
For a Lipschitz map $`w:(\mathrm{\Sigma },\sigma |dz|^2)(\mathrm{\Sigma },\rho |dw|^2)`$, where $`\sigma |dz|^2`$ and $`\rho |dw|^2`$ are metrics on $`\mathrm{\Sigma }`$, and $`z`$ and $`w`$ are conformal coordinates on $`\mathrm{\Sigma }`$, one follows some notations of Sampson () to define
$`(z)=\frac{\rho (w(z))}{\sigma (z)}|w_z|^2,(z)=\frac{\rho (w(z))}{\sigma (z)}|w_{\overline{z}}|^2`$.
Then the energy density of $`w`$ is simply $`e(w)=+`$, and the total energy is then given by
$`E(w,\sigma ,\rho )=_\mathrm{\Sigma }e\sigma |dz|^2`$,
which depends on the target metric and conformal structure of the domain. The map $`w`$ is called harmonic if it is a critical point of this energy functional, i.e., it satisfies Euler-Lagrange equation:
$`w_{z\overline{z}}+\frac{\rho _w}{\rho }w_zw_{z\overline{z}}=0`$.
The $`(2,0)`$ part of the pullback $`w^{}\rho `$ is the so-called Hopf differential:
$`\varphi (z)dz^2=(w^{}\rho )^{(2,0)}=\rho w_z\overline{w}_zdz^2`$.
It is routine to check that $`w`$ is harmonic if and only if $`\varphi dz^2QD(\mathrm{\Sigma })`$, and $`w`$ is conformal if and only if $`\varphi =0`$.
One also finds that
$`(z)(z)={\displaystyle \frac{\varphi \overline{\varphi }}{\sigma ^2}}={\displaystyle \frac{|\varphi |^2}{\sigma ^2}}`$ (2)
and the Jacobian functional is $`J(z)=(z)(z)`$.
Now assume both $`\sigma `$ and $`\rho `$ are canonical metrics on surface $`\mathrm{\Sigma }`$ (then they represent two different conformal structures unless they are biholomorphic). Since the target surface $`(\mathrm{\Sigma },\rho )`$ has negative Gaussian curvatures almost everywhere, with possibly finitely many flat points, the classical theory of harmonic maps guarantees that there is a unique harmonic map $`w:(\mathrm{\Sigma },\sigma )(\mathrm{\Sigma },\rho )`$ in the homotopy class of the identity, moreover, this map $`w`$ is a diffeoemorphism with $`J>0`$ and $`>0`$ (, , , , , ).
## 3 A Homeomorphism
The method of harmonic maps has been a great computational tool in Teichmüller theory (see ) . In the case of hyperbolic metrics on a compact Riemann surface, the second variation of the energy of the harmonic map $`w=w(\sigma ,\rho )`$, with respect to the domain metric $`\sigma `$ (or target metric $`\rho `$) at $`\sigma =\rho `$, yields the Weil-Petersson metric on $`𝒯_g`$ (, ). In our case of canonical metrics on a Riemann surface, we prove a homeomorphism theorem, the theorem 1.1, to link Teichmüller space of canonical metrics to the space $`QD(\mathrm{\Sigma })`$.
To define this map, we fix a point $`\sigma `$ in Teichmüller space, with conformal coordinates $`z`$. Thus $`\sigma `$ is a conformal structure on surface $`\mathrm{\Sigma }`$. For each canonical metric $`\rho `$ on $`\mathrm{\Sigma }`$, we obtain the unique harmonic map $`w(\sigma ,\rho )`$ in the homotopy class of the identity map, since $`\rho `$ is of nonpositive curvature, with only possibly finitely many flat points on $`\mathrm{\Sigma }`$. The associated Hopf differential of the harmonic map $`w(\sigma ,\rho )`$ is then given by $`\varphi (z)dz^2=\rho (w(z))w_z\overline{w}_zdz^2`$. Therefore the map $`\varphi :𝒯_gQD(\mathrm{\Sigma })`$ which sends $`\rho `$ to $`\varphi (z)dz^2`$ is well defined.
###### Remark 3.1.
Sampson considered this map in the case that of hyperbolic metrics (), and showed that it is continuous and one-to-one. Later Wolf showed the map is actually a homeomorphism (). The condition of constant Gaussian curvature of hyperbolic metric is essential in the argument of proving this map is injective.
We start with a technical lemma, which is only slightly different than the case of hyperbolic metrics, as shown in . For Hopf differential $`\varphi `$ corresponding to metric $`\rho `$, we define $`\varphi =_\mathrm{\Sigma }|\varphi |𝑑z𝑑\overline{z}`$. We need to show $`\varphi `$ is approximately the total energy of the harmonic map $`w`$ in a large scale, i.e.,
###### Lemma 3.2.
$`\varphi \mathrm{}`$ if and only if $`E(\rho )\mathrm{}`$
###### Proof.
This harmonic map $`w(z)`$ is naturally quasiconformal, and we write its Beltrami differential as $`\nu =\frac{w_{\overline{z}}}{w_z}`$, and $`|\nu |<1`$.
We abuse our notation to write $`\sigma `$ as the domain canonical metric, and $`dA=\sigma (z)dzd\overline{z}`$ is the area element of the domain surface. Recall from section two, we have density functions $`(z)=\frac{\rho (w(z))}{\sigma (z)}|w_z|^2`$ and $`(z)=\frac{\rho (w(z))}{\sigma (z)}|w_{\overline{z}}|^2`$. The total energy is $`E(\rho )=_\mathrm{\Sigma }((z)+(z))𝑑A`$, and the Jacobian determinant of the map is $`J(z)=(z)(z)`$. Note that $`(z)>(z)0`$.
It is not hard to see that
$`|\varphi |^2=\sigma ^2`$, and $`|\nu |^2=\frac{}{}<1`$.
Therefore, we now have
$`\varphi `$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}|\varphi |𝑑z𝑑\overline{z}={\displaystyle _\mathrm{\Sigma }}|\nu |\sigma 𝑑z𝑑\overline{z}`$
$`=`$ $`{\displaystyle _\mathrm{\Sigma }}|\nu |𝑑A<{\displaystyle _\mathrm{\Sigma }}𝑑A`$
$``$ $`{\displaystyle _\mathrm{\Sigma }}(+)𝑑A=E(\rho ).`$
For the opposite direction, we find $`\sqrt{}`$ since $`<`$, and therefore,
$`E(\rho )`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}(+)𝑑A={\displaystyle _\mathrm{\Sigma }}J(z)𝑑A+2{\displaystyle _\mathrm{\Sigma }}𝑑A`$
$``$ $`Area(\mathrm{\Sigma },\sigma )+2{\displaystyle _\mathrm{\Sigma }}\sqrt{}𝑑A`$
$`=`$ $`g+2{\displaystyle _\mathrm{\Sigma }}{\displaystyle \frac{|\varphi |}{\sigma }}𝑑A=g+2\varphi .`$
Here we used the fact that $`Area(\mathrm{\Sigma },\sigma )=g`$, as pointed out in remark 2.1. This completes the proof of the lemma. ∎
We now start to prove the homeomorphism theorem.
###### Proof.
(Proof of theorem 1.1): It is clear that this map is continuous because of the uniqueness of harmonic map $`w`$ in the homotopy class of the identity map. We want to show this map is a local diffeoemorphism and proper.
We firstly notice that the map $`\varphi `$ is a local diffeomorphism. To see this, we consider a sufficiently small neighborhood of $`\sigma `$ and a family of harmonic maps $`w(t):(\mathrm{\Sigma },\sigma )(\mathrm{\Sigma },\rho (t))`$ between canonical metrics near $`t=0`$, where $`\rho (t)`$ is a family of canonical metrics with $`\rho (0)=\sigma `$. It is easy to see that $`w(0)=z`$, the identity map. Associated Hopf differentials of this family are given by $`\varphi (t)dz^2=\rho (t)w_z(t)\overline{w}_z(t)dz^2`$ with $`\varphi (0)=0`$. we take $`t`$-derivative on $`\varphi (t)`$ at $`t=0`$ to find that
$`\frac{d\varphi (t)}{dt}|_{t=0}=\rho (0)w_z(0)\frac{d\overline{w}_z(t)}{dt}|_{t=0}=\sigma \frac{d\overline{w}_z(t)}{dt}|_{t=0}`$.
This shows that $`\frac{dw(t)}{dt}|_{t=0}`$ is conformal, provided that $`\frac{d\varphi (t)}{dt}|_{t=0}=0`$. So the map $`d\varphi `$ is nonsingular, and $`\varphi `$ is a local diffeomorphism by applying inverse function theorem.
We then apply a slightly rearranged argument of Wolf () (on hyperbolic metrics) to show map $`\varphi `$ is proper.
Given that $`𝒯_g`$ and $`QD(\mathrm{\Sigma })`$ are finite dimensional spaces, and from lemma 3.2, it suffices to show the energy function $`E(\rho )`$ is a proper map from Teichmüller space to $`\mathrm{}`$ (see theorem 2.7.1, ). In other words, we need to show the set $`B=\{\rho 𝒯_g:E(\rho )K\}`$ is a compact subset of $`𝒯_g`$. Without loss of generality, we assume $`id:(\mathrm{\Sigma },\sigma )(\mathrm{\Sigma },\rho )`$ is harmonic, or we can choose $`w^{}\rho `$ to represent the equivalency class $`[\rho ]`$.
Consider a geodesic ball $`B(x_0,\delta )`$ for some $`x_0`$ in domain surface $`(\mathrm{\Sigma },\sigma )`$, where positive constant $`\delta <min\{1,inj_\sigma (\mathrm{\Sigma })^2\}`$, where $`inj_\sigma (\mathrm{\Sigma })`$ is the injectivity radius of $`\mathrm{\Sigma }`$ with respect to the metric $`\sigma `$. Notice that a harmonic map between surfaces does not depend on the choice of metrics on the domain, but on the choice of conformal structures of the domain surface. Therefore, we can choose $`\sigma `$ to be hyperbolic in this argument, and then introduce polar coordinates $`(r,\theta )`$ in the hyperbolic disk $`B(x_0,\delta )`$ so that $`\sigma =dr^2+sinh^2(r)d\theta ^2`$.
For $`r<\sqrt{\delta }<1`$, we have $`sinh(r)<2r`$ and then
$`_\delta ^\sqrt{\delta }\frac{dr}{sinh(r)}>\frac{1}{2}_\delta ^\sqrt{\delta }\frac{dr}{r}=\frac{1}{4}|log\delta |`$.
Now considering the annulus $`A(x_0)=A(x_0,\delta ,\sqrt{\delta })`$ centered at $`x_0`$ of inner and outer radii $`\delta `$ and $`\sqrt{\delta }`$, respectively, in domain metric $`\sigma `$, we apply the upper bound of the energy to find
$`{\displaystyle _{A(x_0)}\frac{}{\theta }_\rho ^2\frac{dr}{sinh(r)}𝑑\theta }`$ $``$ $`{\displaystyle _{A(x_0)}\frac{}{r}_\rho ^2}+{\displaystyle \frac{1}{sinh^2(r)}}{\displaystyle \frac{}{\theta }}_\rho ^2sinh(r)drd\theta `$
$``$ $`2{\displaystyle _\mathrm{\Sigma }}e(z)𝑑A=2E(\rho )2K.`$
Thus there exists $`\delta <r<\sqrt{\delta }`$ such that
$`_0^{2\pi }\frac{}{\theta }_\rho ^2𝑑\theta \frac{8K}{|log\delta |}`$.
For this $`r`$ and two points $`x_3`$ and $`x_4`$ on the boundary of the disk $`B_\sigma (x_0,r)`$, and two points $`x_1`$ and $`x_2`$ in $`B_\sigma (x_0,\delta )`$, we now have,
$`d_\rho (w(x_1),w(x_2))`$ $`=`$ $`d_\rho (x_1,x_2)d_\rho (x_3,x_4)`$
$``$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{}{\theta }}_\rho 𝑑\theta 2\pi \sqrt{{\displaystyle _0^{2\pi }}{\displaystyle \frac{}{\theta }}_\rho ^2𝑑\theta }`$
$``$ $`{\displaystyle \frac{4\sqrt{2K}\pi }{|log\delta |}}.`$
Applying the Courant-Lebesgue lemma, we conclude that the energy $`E(\rho )`$ is proper, hence so is map $`\varphi `$.
We have showed that the map $`\varphi `$ is a proper local diffeomorphism between Teichmüller space and $`QD(\mathrm{\Sigma })`$. It is clear that $`𝒯_g`$ is path-connected and $`QD(\mathrm{\Sigma })`$ is a simply connected Hausdorff space. Standard theory of covering map between manifolds implies such a local homeomorphism is actually a global homeomorphism. This completes the proof of theorem 1.1. ∎
## 4 Variations of the Energy
In the previous section, we showed that the map $`\varphi :𝒯_gQD(\mathrm{\Sigma })`$ is a homeomorphism, thus its inverse map $`\varphi ^1`$ provides coordinates for any canonical metric $`\rho 𝒯_g`$. We will study these coordinates in this section, i.e., we apply variational approach to derive the infinitesimal $`L^2`$-Bergman norm on Teichmüller space. We will separate the cases where either target canonical metrics are varying or domain canonical metrics are varying, in subsections 4.1 and 4.2, respectively.
### 4.1 Varying the Target
To prove theorem 1.2, we need to develop some infinitesimal calculations for the variation of a harmonic map between canonical metrics. This technique is an analog to that of Wolf’s on the case of hyperbolic metrics, which plays an important role in studying the Weil-Petersson geometry of Teichmüller space.
Now we consider a family of harmonic maps $`w(t):(\mathrm{\Sigma },\sigma )(\mathrm{\Sigma },\rho (t))`$ between canonical metrics, where $`w(t)`$ varies real analyticaly in $`t`$ for $`|t|<ϵ`$, and $`\rho (t)`$ is a family of canonical metrics with $`\rho (0)=\sigma `$, therefore $`w(0)=z`$. Associated Hopf differentials are given by $`\varphi (t)dz^2=\rho (t)w_z(t)\overline{w}_z(t)dz^2`$ with $`\varphi (0)=0`$.
For $`t=(t^\alpha ,t^\beta )`$, denote $`\varphi _\alpha =\frac{d\varphi (t)}{dt^\alpha }|_{t=0}`$ and $`\varphi _\beta =\frac{d\varphi (t)}{dt^\beta }|_{t=0}`$ as infinitesimal holomorphic quadratic differentials.
Recall that the holomorphic and antiholomorphic functions of this family of harmonic maps are
$`(t)=\frac{\rho (w(t))}{\sigma (z)}|w_z(t)|^2,(t)=\frac{\rho (w(t))}{\sigma (z)}|w_{\overline{z}}(t)|^2`$.
We denote $`_\alpha =\frac{d(t)}{dt^\alpha }|_{t=0}`$ and $`_\alpha =\frac{d(t)}{dt^\alpha }|_{t=0}`$, also $`_{\alpha \overline{\beta }}=\frac{d^2}{dt^\alpha \overline{dt^\beta }}|_{t=0}(t)`$, and we assign similiar meaning for $`_{\alpha \overline{\beta }}`$ and $`E_{\alpha \overline{\beta }}`$.
We also write $`K(t)=K(\rho (t))=\frac{1}{2}\mathrm{\Delta }_\rho log\rho `$ as the Gaussian curvature of the metric $`\rho (t)`$, and assign obvious meaning to $`K_\alpha `$ and $`K_{\alpha \overline{\beta }}`$. Since $`K(\sigma )0`$ and is negative everywhere except possibly finitely many points, it is not hard to see that the operator $`\mathrm{\Delta }_\sigma +2K(\sigma )`$ is invertible on $`(\mathrm{\Sigma },\sigma )`$, and we denote $`D_B=2(\mathrm{\Delta }_\sigma +2K(\sigma ))^1`$.
###### Lemma 4.1.
For this family of harmonic maps $`w(t)`$, the following holds:
* $`(0)=1`$ and $`(0)=0`$;
* $`_\alpha 0`$, $`_\alpha =D_B(K_\alpha )`$, and $`_\mathrm{\Sigma }_\alpha \sigma =0`$;
* $`_{\alpha \overline{\beta }}=\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2}`$;
* $`_{\alpha \overline{\beta }}=D_B(K_{\alpha \overline{\beta }})+D_B(K_{\overline{\beta }}D_B(K_\alpha ))+D_B(K_\alpha D_B(K_{\overline{\beta }}))`$
$`D_B(K(\sigma )\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2})\frac{1}{2}D_B(\mathrm{\Delta }_\sigma (D_B(K_\alpha )D_B(K_{\overline{\beta }})))`$.
###### Proof.
* This is true since the map $`w(t)`$ is the identity map at time $`t=0`$.
* Recalling formula (2):
$`(t)(t)=\frac{\varphi (t)\overline{\varphi }(t)}{\sigma ^2}`$,
we take $`t`$-derivative at $`t=0`$, to find that
$`_\alpha (0)+(0)_\alpha =\frac{\varphi _\alpha \overline{\varphi }(0)+\varphi (0)\overline{\varphi }_\alpha }{\sigma ^2}`$.
The righthand side is zero as $`\varphi (0)=0`$. Therefore $`_\alpha 0`$.
We notice that $`_\mathrm{\Sigma }((t)(t))\sigma =_\mathrm{\Sigma }J(t)\sigma =g`$ is independent of the parameter $`t`$. Therefore
$`_\mathrm{\Sigma }_\alpha \sigma =_\mathrm{\Sigma }_\alpha \sigma =0`$.
From standard Bochner identities, we have
$`\mathrm{\Delta }_\sigma log=2K(\sigma )2K(\rho )().`$ (3)
Therefore $`\mathrm{\Delta }_\sigma log(t)=2K(\sigma )2K(\rho (t))((t)(t))`$ and
$`\mathrm{\Delta }_\sigma _\alpha `$ $`=`$ $`\mathrm{\Delta }_\sigma {\displaystyle \frac{_\alpha }{(0)}}`$
$`=`$ $`2K_\alpha ((0)(0))2K(\rho (0))(_\alpha _\alpha )`$
$`=`$ $`2K_\alpha 2K(\sigma )_\alpha .`$
We now obtain $`(\mathrm{\Delta }_\sigma +2K(\sigma ))_\alpha =2K_\alpha `$ and then
$`_\alpha =2(\mathrm{\Delta }_\sigma +2K(\sigma ))^1(K_\alpha )=D_B(K_\alpha )`$.
* To calculate next variation, we consider formula $`(t)(t)=\frac{\varphi (t)\overline{\varphi }(t)}{\sigma ^2}`$ again. We find
$`_{\alpha \overline{\beta }}(0)+_{\overline{\beta }}_\alpha +_\alpha _{\overline{\beta }}+(0)_{\alpha \overline{\beta }}=\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2}`$.
Therefore $`_{\alpha \overline{\beta }}=\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2}`$.
* We take second $`t`$-derivative from (3) to find
$`\mathrm{\Delta }_\sigma (_{\alpha \overline{\beta }}_\alpha _{\overline{\beta }})`$ $`=`$ $`2K_{\alpha \overline{\beta }}2K_\alpha _{\overline{\beta }}2K_{\overline{\beta }}_\alpha `$
$``$ $`2K(\sigma )(_{\alpha \overline{\beta }}_{\alpha \overline{\beta }}),`$
then we obtain that
$`(\mathrm{\Delta }_\sigma +2K(\sigma ))(_{\alpha \overline{\beta }})`$ $`=`$ $`\mathrm{\Delta }(_\alpha _{\overline{\beta }})2K_{\alpha \overline{\beta }}2K_\alpha _{\overline{\beta }}2K_{\overline{\beta }}_\alpha `$
$``$ $`2K(\sigma )(_{\alpha \overline{\beta }}_{\alpha \overline{\beta }}).`$
Now we apply formulas $`_\alpha =D_B(K_\alpha )`$, and $`_{\overline{\beta }}=D_B(K_{\overline{\beta }})`$, and $`_{\alpha \overline{\beta }}=\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2}`$ to above equation to complete the proof of this lemma.
###### Remark 4.2.
It is very interesting to compare our situation with the case of varations of a harmonic map between hyms on the surface. If we assume all metrics on the surface are hyperbolic with constant Gaussian curvature, under the same notations, then we have the following comparison: (i) and (iii) in lemma 4.1 hold; (ii) also holds except furthermore, $`_\alpha 0`$, i.e., the holomorphic energy reaches its minimum at time zero; (iv) of the lemma takes the form of $`_{\alpha \overline{\beta }}=D(\frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2})`$, where $`D=2(\mathrm{\Delta }_\sigma 2)^1`$ is a compact, self-adjoint operator. Operator $`D_B`$ in (iv) of lemma 4.1 is not self-adjoint for $`L^2`$ functions, while $`K(\sigma )D_B`$ is, and coincides with operator $`D`$ when $`K1`$.
We now consider the variations of the corresponding total energy $`E(t)`$ of the family $`w(t)`$ near $`t=0`$, i.e., we show theorem 1.2 in following equivalent form:
###### Theorem 4.3.
$`\frac{d^2}{dt^\alpha \overline{dt^\beta }}|_{t=0}E(t)=2<\varphi _\alpha ,\varphi _\beta >_B`$.
###### Proof.
The total energy is
$`E(t)`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}((t)+(t))\sigma `$
$`=`$ $`{\displaystyle _\mathrm{\Sigma }}((t)(t))\sigma +2{\displaystyle _\mathrm{\Sigma }}(t)\sigma `$
$`=`$ $`{\displaystyle _\mathrm{\Sigma }}J(t)\sigma +2{\displaystyle _\mathrm{\Sigma }}(t)\sigma `$
$`=`$ $`g+2{\displaystyle _\mathrm{\Sigma }}(t)\sigma g,`$
since $`E(0)=g`$ is equal to the area of the surface. Thus $`E(t)`$ reaches its global minimum $`g`$ at $`t=0`$ from (i) of lemma 4.1.
From (ii) of lemma 4.1, it is easy to see that $`t=0`$ is a critical point of $`E(t)`$ as
$`E_\alpha =_\mathrm{\Sigma }(_\alpha +_\alpha )\sigma =0`$.
and
$`E_{\alpha \overline{\beta }}`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}(_{\alpha \overline{\beta }}+_{\alpha \overline{\beta }})\sigma `$
$`=`$ $`2{\displaystyle _\mathrm{\Sigma }}_{\alpha \overline{\beta }}\sigma =2{\displaystyle _\mathrm{\Sigma }}{\displaystyle \frac{\varphi _\alpha \overline{\varphi }_\beta }{\sigma ^2}}\sigma `$
$`=`$ $`2<\varphi _\alpha ,\varphi _\beta >_B.`$
### 4.2 Varying the Domain
In this subsection, we consider a family of harmonic maps between fixed target metric and varying domain metrics.
Again, since the target metric is negatively curved (except possibly finitely many flat points), we have the existence and uniqueness of a harmonic map in the homotopy class of the identity, and this map is a diffeomorphism. In other words, let $`w(s):(\mathrm{\Sigma },\sigma (s)|dz|^2)(\mathrm{\Sigma },\rho |dw|^2)`$ be this family of harmonic maps near the identity map, where $`w(s)`$ varies real analytically in $`s`$ for $`|s|<ϵ`$, and $`\sigma (s)`$ is a family of canonical metrics with $`\sigma (0)=\rho `$, therefore $`w(0)=z`$. Associated Hopf differentials are given by $`\varphi (s)dz(s)^2=\rho w_z(s)\overline{w}_z(s)dz(s)^2`$ with $`\varphi (0)=0`$.
For $`s=(s^a,s^b)`$, similar to last subsection, we denote $`\varphi _a=\frac{\varphi (s)}{s^a}|_{s=0}`$ and $`\varphi _b=\frac{\varphi (s)}{s^b}|_{s=0}`$, and assign similar meanings to $`_a`$, and $`_{a\overline{b}}`$, etc.
Let $`K(s)=\frac{1}{2}\mathrm{\Delta }_{\sigma (s)}log\sigma (s)`$ be the Gaussian curvature of the surface $`(\mathrm{\Sigma },\sigma (s))`$ and denote $`K_a=\frac{K(s)}{s^a}|_{s=0}`$. Again, since $`K(\rho )0`$ and is negative everywhere except possibly finitely many points, the operator $`\mathrm{\Delta }_\rho +2K(\rho )`$ is invertible on $`(\mathrm{\Sigma },\rho )`$, and we denote $`D_B^{}=2(\mathrm{\Delta }_\rho +2K(\rho ))^1`$. This operator $`D_B^{}`$ is not self-adjoint for $`L^2`$ functions.
We firstly calculate the variations of these two density functions $`(s)`$ and $`(s)`$. It is interesting to notice the difference with the case of varying the target showed in lemma 4.1.
###### Lemma 4.4.
For this family of harmonic maps $`w(s)`$, the following holds:
* $`(0)=1`$ and $`(0)=0`$;
* $`_a0`$, $`_a=D_B^{}(K_a)`$, and $`_\mathrm{\Sigma }_a\sigma =0`$;
* $`_{a\overline{b}}=\frac{\varphi _a\overline{\varphi }_b}{\sigma ^2}`$.
###### Proof.
* It is true since the map $`w(0)`$ is the identity map.
* We take $`s^a`$-derivative of $`(s)(s)=\frac{\varphi (s)\overline{\varphi }(s)}{\sigma ^2(s)}`$ to find
$`_a(0)+(0)_a=\frac{\varphi _a\overline{\varphi }(0)+\varphi (0)\overline{\varphi }_a}{\sigma ^2(0)}+[\frac{(\frac{1}{\sigma ^2(s)})}{s^a}]|_{s=0}|\varphi (0)|^2`$,
and this implies $`_a=0`$, for $`\varphi (0)=0`$.
Therefore $`_\mathrm{\Sigma }_a\sigma =_\mathrm{\Sigma }_a\sigma =0`$.
To calculate $`_a`$, recalling formula (3):
$`\mathrm{\Delta }_{\sigma (s)}log(s)=2K(\sigma (s))2K(\rho )((s)(s))`$,
we find that
$`{\displaystyle \frac{}{s^a}}|_{s=0}(\mathrm{\Delta }_{\sigma (s)})log(0)+\mathrm{\Delta }_\rho (_a)`$ $`=`$ $`2K_a2K(\rho )(_a_a)`$
$`=`$ $`2K_a2K(\rho )(_a).`$
Therefore $`(\mathrm{\Delta }_\rho +2K(\rho ))(_a)=2K_a`$, and so $`_a=D_B^{}(K_a)`$.
* For the second variation of $`(s)`$, we have
$`_{a\overline{b}}`$ $`=`$ $`_{a\overline{b}}(0)+_{\overline{b}}_a+_a_{\overline{b}}+(0)_{a\overline{b}}`$
$`=`$ $`{\displaystyle \frac{\varphi _a\overline{\varphi }_b}{\rho ^2}}+[{\displaystyle \frac{(\frac{1}{\sigma ^2(s)})}{\overline{s}^b}}]|_{s=0}[\varphi _a\overline{\varphi }(0)+\varphi (0)\overline{\varphi }_a]`$
$`+`$ $`[{\displaystyle \frac{(\frac{1}{\sigma ^2(s)})}{s^a}}]|_{s=0}[\varphi _b\overline{\varphi }(0)+\varphi (0)\overline{\varphi }_b]+[{\displaystyle \frac{^2(\frac{1}{\sigma ^2(s)})}{s^a\overline{s}^b}}]|_{s=0}[\varphi (0)\overline{\varphi }(0)]`$
$`=`$ $`{\displaystyle \frac{\varphi _a\overline{\varphi }_b}{\rho ^2}}.`$
Now we show the equivalent form of theorem 1.3:
###### Theorem 4.5.
$`\frac{^2}{s^a\overline{s^b}}|_{s=0}E(s)=2<\varphi _a,\varphi _b>_B`$.
###### Proof.
The total energy is now
$`E(s)`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}((s)+(s))\sigma (s)`$
$`=`$ $`{\displaystyle _\mathrm{\Sigma }}J(s)\sigma (s)+2{\displaystyle _\mathrm{\Sigma }}(s)\sigma (s)`$
$`=`$ $`g+2{\displaystyle _\mathrm{\Sigma }}(s)\sigma (s)g,`$
where $`E(0)=g`$ reaches the global minimum.
Together with $`(0)=0`$, we then find
$`E_a=2\frac{}{s^a}|_{s=0}[_\mathrm{\Sigma }(s)\sigma (s)]=0`$,
and then $`s=0`$ is also a critical point of $`E(s)`$.
Now we consider $`\frac{^2}{s^a\overline{s^b}}|_{s=0}E(s)`$ from lemma 4.4. We apply $`(0)=_a=_{\overline{b}}=0`$ to find
$`{\displaystyle \frac{^2}{s^a\overline{s^b}}}|_{s=0}E(s)`$ $`=`$ $`2{\displaystyle \frac{^2}{s^a\overline{s}^b}}|_{s=0}\{{\displaystyle _\mathrm{\Sigma }}(s)\sigma (s)\}`$
$`=`$ $`2{\displaystyle _\mathrm{\Sigma }}{\displaystyle \frac{\varphi _a\overline{\varphi }_b}{\rho ^2}}\rho `$
$`=`$ $`2<\varphi _a,\varphi _b>_B.`$
###### Remark 4.6.
Teichmüller space is a complex manifold (when $`g2`$), so it has its own complex structure. For Riemannian metrics on this complex manifold, it is ideal that metrics are compatible with the complex structure. Ahlfors () showed that the Weil-Petersson metric is Kählerian. From the definition, we know that the $`L^2`$-Bergman metric is an Hermitian metric, yet it is unknown if the $`L^2`$-Bergman metric is actually Kählerian.
|
warning/0506/hep-ph0506265.html
|
ar5iv
|
text
|
# Low-energy photon-photon collisions to two loops revisited
## 1 Introduction
We consider the process $`\gamma \gamma \pi ^0\pi ^0`$ in the framework of chiral perturbation theory (ChPT ) . The one-loop calculation of the scattering amplitude was performed in Refs. , and the two-loop amplitude was worked out in . Because the effective Lagrangian at order $`p^6`$ was not available at that time, the ultraviolet divergences were evaluated in the $`\overline{\mathrm{MS}}`$scheme, then dropped and replaced with a corresponding polynomial in the external momenta. The three new counterterms which enter at this order in the low-energy expansion were estimated with resonance saturation. Whereas such a procedure is legitimate from a technical point of view, it does not make use of the full information provided by chiral symmetry. The evaluation of the two-loop amplitude involving charged pions was performed later by Burgi .
Over the last ten years, considerable progress has been made in this field, both in theory and experiment. As for theory, the Lagrangian at order $`p^6`$ has been constructed , and its divergence structure has been determined . This provides an important check on the above calculations: adding the counterterm contributions from the $`p^6`$ Lagrangian to the $`\overline{\mathrm{MS}}`$amplitude evaluated in and in must provide a scale independent result. Also in the theory, improved techniques to evaluate the two-loop diagrams that occur in these amplitudes have been developed . The improvement arises mainly in the evaluation of diagrams with four external legs, where the techniques of Ref. allow one to extract the ultraviolet divergences by use of simple recursion relations. We are now able to present the final result for the two-loop amplitudes in a rather compact form (in Refs. , the result was presented partly in numerical form only, because the algebraic expressions were too long to be published).
Concerning experiment, quadrupole polarizabilities for the neutral pions have recently been determined from data on $`\gamma \gamma \pi ^0\pi ^0`$ . Further, the charged pion polarizabilities $`(\alpha \beta )_{\pi ^+}`$ have been determined at the Mainz Microtron MAMI , with a result that is at variance with the two-loop calculation presented in . Last but not least, there is an ongoing experiment by the COMPASS collaboration at CERN to measure the charged pion and kaon polarizabilities .
In view of these developments, and because the two-loop expressions for the polarizabilities had never been checked, we decided to recalculate these amplitudes, using the improved techniques of Ref. to evaluate the integrals, and invoking the chiral Lagrangian at order $`p^6`$ . As the calculation in the case of neutral pions involves considerably less diagrams, and because the Fortran code for these amplitudes is still available to us for checks, we have decided to start the program with a re-evaluation of these amplitudes. This is the main purpose of the present work. The evaluation of the corresponding expressions for the charged pions and for the kaons is underway and will be presented elsewhere .
The article is organized as follows. Section 2 contains the necessary kinematics of the process $`\gamma \gamma \pi ^0\pi ^0`$. To make the article self contained, we summarize in Section 3 the necessary ingredients from the effective Lagrangian framework. In Section 4, we display the Feynman diagrams and discuss their evaluation. Section 5 contains a concise representation of the two Lorentz invariant amplitudes that describe the scattering matrix element. In Section 6, we compare the present work with the previous calculation , while Section 7 contains explicit expressions for the dipole and quadrupole polarizabilities valid at next-to-next-to-leading order in the chiral expansion, together with a numerical analysis and a comparison with an evaluation from data on $`\gamma \gamma \pi ^0\pi ^0`$ . The summary and an outlook are given in Section 8. Finally, several technical aspects of the calculation are relegated to the Appendices.
## 2 Kinematics
The matrix element for the reaction
$$\gamma (q_1)\gamma (q_2)\pi ^0(p_1)\pi ^0(p_2)$$
(2.1)
is given by
$$\pi ^0(p_1)\pi ^0(p_2)\mathrm{out}|\gamma (q_1)\gamma (q_2)\mathrm{in}=i(2\pi )^4\delta ^{(4)}\left(P_fP_i\right)T^N,$$
(2.2)
with
$`T^N`$ $`=`$ $`e^2ϵ_1^\mu ϵ_2^\nu V_{\mu \nu },`$
$`V_{\mu \nu }`$ $`=`$ $`i{\displaystyle 𝑑xe^{i(q_1x+q_2y)}\pi ^0(p_1)\pi ^0(p_2)\mathrm{out}|Tj_\mu (x)j_\nu (y)|\mathrm{\hspace{0.17em}0}}.`$ (2.3)
Here $`j_\mu `$ is the electromagnetic current, and $`\alpha =e^2/4\pi 1/137`$. We consider real photons, $`q_i^2=0`$, with $`ϵ_iq_i=0`$. The decomposition of the correlator $`V_{\mu \nu }`$ into Lorentz invariant amplitudes reads
$`V_{\mu \nu }`$ $`=`$ $`A(s,t,u)T_{1\mu \nu }+B(s,t,u)T_{2\mu \nu }+C(s,t,u)T_{3\mu \nu }+D(s,t,u)T_{4\mu \nu },`$
$`T_{1\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}sg_{\mu \nu }q_{1\nu }q_{2\mu },`$
$`T_{2\mu \nu }`$ $`=`$ $`2s\mathrm{\Delta }_\mu \mathrm{\Delta }_\nu \nu ^2g_{\mu \nu }2\nu (q_{1\nu }\mathrm{\Delta }_\mu q_{2\mu }\mathrm{\Delta }_\nu ),`$
$`T_{3\mu \nu }`$ $`=`$ $`q_{1\mu }q_{2\nu },`$
$`T_{4\mu \nu }`$ $`=`$ $`s(q_{1\mu }\mathrm{\Delta }_\nu q_{2\nu }\mathrm{\Delta }_\mu )\nu (q_{1\mu }q_{1\nu }+q_{2\mu }q_{2\nu }),`$
$`\mathrm{\Delta }_\mu `$ $`=`$ $`(p_1p_2)_\mu ,`$ (2.4)
where
$`s`$ $`=`$ $`(q_1+q_2)^2,t=(p_1q_1)^2,u=(p_2q_1)^2,\nu =tu`$ (2.5)
are the standard Mandelstam variables. The tensor $`V_{\mu \nu }`$ satisfies the Ward identities
$$q_1^\mu V_{\mu \nu }=q_2^\nu V_{\mu \nu }=0.$$
(2.6)
The amplitudes $`A`$ and $`B`$ are analytic functions of the variables $`s,t`$ and $`u`$, symmetric under crossing $`(t,u)(u,t)`$. The amplitudes $`C`$ and $`D`$ do not contribute to the process considered here, because $`ϵ_iq_i=0`$.
It is useful to introduce in addition the helicity amplitudes
$`H_{++}`$ $`=`$ $`A+2(4M_\pi ^2s)B,H_+={\displaystyle \frac{8(M_\pi ^4tu)}{s}}B.`$ (2.7)
The helicity components $`H_{++}`$ and $`H_+`$ correspond to photon helicity differences $`\lambda =0,2`$, respectively. With our normalization of states $`𝐩_\mathrm{𝟏}|𝐩_\mathrm{𝟐}=2(2\pi )^3p_1^0\delta ^{(3)}(𝐩_\mathrm{𝟏}𝐩_\mathrm{𝟐})`$, the differential cross section for unpolarized photons in the centre-of-mass system is
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}^{\gamma \gamma \pi ^0\pi ^0}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2s}{64}}\beta (s)H(s,t),H(s,t)=|H_{++}|^2+|H_+|^2,`$ (2.8)
with $`\beta (s)=\sqrt{14M_\pi ^2/s}`$. The relation between the helicity amplitudes $`M_{+\pm }`$ in Ref. and the amplitudes used here is
$`M_{++}(s,t)=2\pi \alpha H_{++}(s,t),M_+(s,t)=16\pi \alpha B(s,t).`$ (2.9)
The physical regions for the reactions $`\gamma \gamma \pi ^0\pi ^0`$ and $`\gamma \pi ^0\gamma \pi ^0`$ are displayed in Fig. 1, where we also indicate with dashed lines the nearest singularities in the amplitudes $`A`$ and $`B`$. These singularities are generated by two-pion intermediate states in the $`s,t`$ and $`u`$ channel.
## 3 The effective Lagrangian and its low energy constants
The effective Lagrangian consists of a string of terms. Here, we consider QCD with two flavours, in the isospin symmetry limit $`m_u=m_d=\widehat{m}`$. At next-to-next-to-leading order (NNLO), one has
$`_{\mathrm{eff}}=_2+_4+_6.`$ (3.1)
The subscripts refer to the chiral order. The expression for $`_2`$ is
$`_{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`{\displaystyle \frac{F^2}{4}}D_\mu UD^\mu U^{}+M^2(U+U^{}),`$
$`D_\mu U`$ $`=`$ $`_\mu Ui(QUUQ)A_\mu ,Q={\displaystyle \frac{e}{2}}\mathrm{diag}(1,1),`$ (3.2)
where $`e`$ is the electric charge, and $`A_\mu `$ denotes the electromagnetic field. The quantity $`F`$ denotes the pion decay constant in the chiral limit, and $`M^2`$ is the leading term in the quark mass expansion of the pion (mass)<sup>2</sup>, $`M_\pi ^2=M^2(1+O(\widehat{m}))`$. Further, the brackets $`\mathrm{}`$ denote a trace in flavour space. In Eq. (3), we have retained only the terms relevant for the present application, i.e., we have dropped additional external fields. We choose the unitary $`2\times 2`$ matrix $`U`$ in the form
$`U`$ $`=`$ $`\sigma +i\pi /F,\sigma ^2+{\displaystyle \frac{\pi ^2}{F^2}}=\mathrm{𝟏}_{2\times 2},\pi =\left(\begin{array}{cc}\pi ^0& \sqrt{2}\pi ^+\\ \sqrt{2}\pi ^{}& \pi ^0\end{array}\right).`$ (3.5)
The Lagrangian at NLO has the structure
$`_4={\displaystyle \underset{i=1}{\overset{10}{}}}l_iK_i={\displaystyle \frac{l_1}{4}}D_\mu UD^\mu U^{}^2+\mathrm{},`$ (3.6)
where $`l_i`$ denote low energy couplings (LECs), not fixed by chiral symmetry. At NNLO, one has
$`_6={\displaystyle \underset{i=1}{\overset{57}{}}}c_iP_i.`$ (3.7)
For the explicit expressions of the polynomials $`P_i`$, we refer the reader to Refs. . The vertices relevant for $`\gamma \gamma \pi ^0\pi ^0`$ involve $`l_1,\mathrm{},l_6`$ from $`_4`$ and $`c_{29},\mathrm{},c_{34}`$ from $`_6`$.
The couplings $`l_i`$ and $`c_i`$ absorb the divergences at order $`p^4`$ and $`p^6`$, respectively,
$`l_i`$ $`=`$ $`(\mu c)^{d4}\left\{l_i^r(\mu ,d)+\gamma _i\mathrm{\Lambda }\right\},`$
$`c_i`$ $`=`$ $`{\displaystyle \frac{(\mu c)^{2(d4)}}{F^2}}\left\{c_i^r(\mu ,d)\gamma _i^{(2)}\mathrm{\Lambda }^2(\gamma _i^{(1)}+\gamma _i^{(L)}(\mu ,d))\mathrm{\Lambda }\right\},`$
$`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2(d4)}},\mathrm{ln}c={\displaystyle \frac{1}{2}}\left\{\mathrm{ln}4\pi +\mathrm{\Gamma }^{}(1)+1\right\}.`$ (3.8)
The physical couplings are $`l_i^r(\mu ,4)`$ and $`c_i^r(\mu ,4)`$, denoted by $`l_i^r,c_i^r`$ in the following. The coefficients $`\gamma _i`$ are given in , and $`\gamma _i^{(1,2,L)}`$ are tabulated in . In order to compare the present calculation with the result of , we will use the scale independent quantities $`\overline{l}_i`$ introduced in ,
$`l_i^r`$ $`=`$ $`{\displaystyle \frac{\gamma _i}{32\pi ^2}}(\overline{l}_i+l),`$ (3.9)
where the chiral logarithm is $`l=\mathrm{ln}(M_\pi ^2/\mu ^2)`$. We will use
$`\overline{l}_1`$ $`=`$ $`0.4\pm 0.6,\overline{l}_2=4.3\pm 0.1,\overline{l}_3=2.9\pm 2.4,\overline{l}_4=4.4\pm 0.2,`$ (3.10)
and
$`\overline{l}_\mathrm{\Delta }\overline{l}_6\overline{l}_5=3.0\pm 0.3.`$ (3.11)
The constants $`c_i^r`$ occur in the combinations
$`a_1^r`$ $`=`$ $`4096\pi ^4\left(c_{29}^rc_{30}^r+c_{34}^r\right),`$
$`a_2^r`$ $`=`$ $`256\pi ^4\left(8c_{29}^r+8c_{30}^r+c_{31}^r+c_{32}^r+2c_{33}^r\right),`$
$`b^r`$ $`=`$ $`128\pi ^4\left(c_{31}^r+c_{32}^r+2c_{33}^r\right).`$ (3.12)
Their values have been estimated by resonance exchange e.g. in Ref. , see also , where $`c_{34}^r`$ has been determined from a chiral sum rule. For the present application, we simply take the values obtained in ,
$`a_1^r(M_\rho )+8b^r(M_\rho )`$ $`=`$ $`14\pm 5,`$
$`a_2^r(M_\rho )2b^r(M_\rho )`$ $`=`$ $`7\pm 3,`$
$`b^r(M_\rho )`$ $`=`$ $`3\pm 1;M_\rho =770\text{MeV}.`$ (3.13)
In the numerical evaluations discussed later on, we use $`\mu =M_\rho `$. As mentioned already in Ref. , varying this scale between 500 MeV and 1 GeV leads to a negligible change of e.g. the cross section for the reaction $`\gamma \gamma \pi ^0\pi ^0`$ below 400 MeV. Finally, we will use $`F_\pi =92.4`$ MeV (see for a recent update of this value), and $`M_\pi =135`$ MeV.
## 4 Evaluation of the diagrams
The lowest-order contributions are the one-loop diagrams displayed in Fig. 2. They have been evaluated for the first time in Refs. , where it was noticed that the sum of these two amplitudes is ultraviolet finite, because there are no contributions from the effective Lagrangian at order $`p^4`$ at this order. The two-loop diagrams are displayed in the Figs. 3,5 and 6. The two-loop diagrams in Fig. 3 may be generated according to the scheme indicated in Fig. 4, where the shadowed blob denotes the $`d`$-dimensional elastic $`\pi \pi `$-scattering amplitude at one-loop accuracy, with two pions off-shell. As is discussed in Appendix B, the one-loop integrals in the $`\pi \pi `$ amplitude may be represented in a dispersive manner. This allows one to reduce the two-loop integrals in Fig. 3 to the one-loop ones, where one has to perform at the end an integration over a dispersive parameter.
Two further diagrams are displayed in Fig. 5. The first one - called “acnode” in the literature - may again be evaluated by use of a dispersion relation, see Appendix B. The second one is trivial to evaluate, because it is a product of one-loop diagrams. The remaining diagrams at order $`p^6`$ are shown in Fig. 6.
The evaluation of the diagrams was done in the following manner.
1. We have performed the integration over the loop momenta in the $`d`$-dimensional regularization scheme, in particular using the procedure described in Ref. and in Appendix B, and invoking FORM .
2. We have then checked numerically that the amplitude satisfies the Ward identities (2.6) in $`d`$ dimensions, in the unphysical region, where the amplitudes are real.
3. We have verified that the counterterms from the Lagrangian $`_6`$ remove all ultraviolet divergences, which is a very non-trivial check on our calculation.
4. We have checked that the (ultra-violet finite) amplitude so obtained is scale independent.
5. Finally, we have numerically verified that the three lowest partial waves of the helicity non-flip amplitude $`H_{++}`$ carry the proper one-loop $`\pi \pi `$ phase, in conformity with unitarity.
We note that the steps (3) and (4) could not be performed in Refs. , because the counterterms at order $`p^6`$ were not yet available .
## 5 The two-loop amplitudes
We give the expression for the amplitudes $`A`$ and $`B`$ by using the same notation as in , and refer the reader to Appendix C of this reference for the one-loop integrals $`\overline{J}(s),\stackrel{=}{J}(s),\overline{G}(s),\stackrel{=}{G}(s)`$ and $`\overline{H}(s)`$ that occur. We have
$`A`$ $`=`$ $`{\displaystyle \frac{4\overline{G}_\pi (s)}{sF_\pi ^2}}(sM_\pi ^2)+U_A+P_A+O(E^4),`$
$`B`$ $`=`$ $`U_B+P_B+O(E^2).`$ (5.1)
The quantity $`\overline{G}_\pi (s)`$ stands for $`\overline{G}(s)`$, evaluated with the physical pion mass. The unitary parts $`U_{A(B)}`$ contain $`s,t`$ and $`u`$-channel cuts, and $`P_{A(B)}`$ are linear polynomials in $`s`$. We find
$`U_A={\displaystyle \frac{2}{sF_\pi ^4}}\overline{G}(s)\left[(s^2M_\pi ^4)\overline{J}(s)+C(s,\overline{l}_i)\right]+{\displaystyle \frac{\overline{l}_\mathrm{\Delta }}{24\pi ^2F_\pi ^4}}(sM_\pi ^2)\overline{J}(s)`$
$`+{\displaystyle \frac{(\overline{l}_25/6)}{144\pi ^2sF_\pi ^4}}(s4M_\pi ^2)\{\overline{H}(s)+4[s\overline{G}(s)+2M_\pi ^2(\stackrel{=}{G}(s)3\stackrel{=}{J}(s))]d_{00}^2\}`$
$`+\mathrm{\Delta }_A(s,t,u),`$
$`C(s,\overline{l}_i)={\displaystyle \frac{1}{48\pi ^2}}\{2(\overline{l}_1{\displaystyle \frac{4}{3}})(s2M_\pi ^2)^2+{\displaystyle \frac{1}{3}}(\overline{l}_2{\displaystyle \frac{5}{6}})(4s^28sM_\pi ^2+16M_\pi ^4)`$
$`3M_\pi ^4\overline{l}_3+12M_\pi ^2(sM_\pi ^2)\overline{l}_412sM_\pi ^2+15M_\pi ^4\},`$
$`d_{00}^2={\displaystyle \frac{1}{2}}(3\mathrm{cos}^2\theta 1),`$ (5.2)
$`U_B={\displaystyle \frac{(\overline{l}_25/6)\overline{H}(s)}{288\pi ^2F_\pi ^4s}}+\mathrm{\Delta }_B(s,t,u).`$ (5.3)
The expressions for $`\mathrm{\Delta }_{A(B)}`$ are displayed in the Appendices C and D.
The polynomial parts are
$`P_A`$ $`=`$ $`{\displaystyle \frac{1}{(16\pi ^2F_\pi ^2)^2}}\left[a_1M_\pi ^2+a_2s\right],`$
$`a_1`$ $`=`$ $`a_1^r+{\displaystyle \frac{1}{18}}\left\{4l^2+l\left(8\overline{l}_2+12\overline{l}_\mathrm{\Delta }{\displaystyle \frac{4}{3}}\right){\displaystyle \frac{20}{3}}\overline{l}_2+12\overline{l}_\mathrm{\Delta }+{\displaystyle \frac{110}{9}}\right\},`$
$`a_2`$ $`=`$ $`a_2^r{\displaystyle \frac{1}{18}}\left\{l^2+l\left(2\overline{l}_2+12\overline{l}_\mathrm{\Delta }{\displaystyle \frac{4}{3}}\right){\displaystyle \frac{5}{3}}\overline{l}_2+12\overline{l}_\mathrm{\Delta }+{\displaystyle \frac{697}{144}}\right\},`$ (5.4)
$`P_B`$ $`=`$ $`{\displaystyle \frac{b}{(16\pi ^2F_\pi ^2)^2}},`$
$`b`$ $`=`$ $`b^r{\displaystyle \frac{1}{36}}\left[l^2+l\left(2\overline{l}_2+{\displaystyle \frac{2}{3}}\right){\displaystyle \frac{1}{3}}\overline{l}_2+{\displaystyle \frac{393}{144}}\right],`$
$`l`$ $`=`$ $`\mathrm{log}{\displaystyle \frac{M_\pi ^2}{\mu ^2}}.`$ (5.5)
The constants $`a_1^r,a_2^r`$ and $`b^r`$ are displayed in terms of the LECs at order $`p^6`$ in Eq. (3.12). Using the fact that the bare couplings $`c_i`$ displayed in Eq. (3) are scale independent, one indeed finds that the above expressions for the amplitudes $`A,B`$ are scale independent as well.
## 6 Comparison with the previous calculation
We can now compare the amplitudes $`A,B`$ with the earlier calculation, presented in Section 7 of Ref. . In that reference, the amplitudes were evaluated with a different techniques. Furthermore, the Lagrangian $`_6`$ was not available in those days, and an important ingredient to check the final result was, therefore, missing. We can make the following observations.
1. The amplitudes $`A`$ and $`B`$ consist of a part with explicit analytic expressions, and additional terms $`\mathrm{\Delta }_{A,B}`$, that are given in the Appendices C and D of the present work in the form of integrals over Feynman parameters. These latter terms were given only in numerical form in .
2. The explicit analytic expressions agree with the previous calculation, except for the coefficient of the single logarithm in $`a_2`$ in Eq. (5). The factor 2/3 in is replaced by - 4/3 here. As the present amplitude is scale independent, we conclude that it is the result Eq. (5) which is correct<sup>1</sup><sup>1</sup>1Burgi provides in his thesis work the isospin $`I=0,2`$ amplitudes, and the one for the charged pions. Subtracting the latter from the former reveals that his calculation agrees with the statement just made.. This mistake does not affect the algebraic expressions for the polarizabilities discussed below, for which we fully agree with Ref. .
3. We can compare the quantities $`\mathrm{\Delta }_{A,B}`$ in numerical form only. For this purpose, we have made two checks. First, we have evaluated the cross section for the reaction $`\gamma \gamma \pi ^0\pi ^0`$ below a centre-of-mass energy of 400 MeV, using the same values for the LECs as in . It agrees with the previous one within a fraction of a percent - the difference would not be visible in Fig. 5 of Ref. , and we do not, therefore, reproduce that plot here. Second, we have re-evaluated the two-loop contributions to the polarizabilities presented in column 4 of Table 3 in . The numbers (0.17, - 0.31) in the old calculation become (0.17, - 0.30) here.
4. To summarize, we confirm the previous result up to the coefficient in one of the chiral logarithms, and up to minute changes in the numerical values of $`\mathrm{\Delta }_{A,B}`$. Numerically, the results in are not affected in any significant manner by these modifications, whose effect is by far smaller than the uncertainties generated by the (not precisely known) values of the low energy constants.
## 7 Pion polarizabilities: dipole and quadrupole
The dipole and quadrupole polarizabilities are defined through the expansion of the helicity amplitudes at fixed $`t=M_\pi ^2`$,
$`{\displaystyle \frac{\alpha }{M_\pi }}H_+(s,t=M_\pi ^2)=(\alpha _1\pm \beta _1)_{\pi ^0}+{\displaystyle \frac{s}{12}}(\alpha _2\pm \beta _2)_{\pi ^0}+𝒪(s^2).`$ (7.1)
Because we have at our disposal the helicity amplitudes at two-loop order, we can work out the polarizabilities to the same accuracy. It turns out that all relevant integrals can be performed in closed form. We discuss the results in the remaining part of this Section.
### 7.1 Chiral expansion
Using the same notation as in , we find for the dipole polarizabilities
$$(\alpha _1\pm \beta _1)_{\pi ^0}=\frac{\alpha }{16\pi ^2F_\pi ^2M_\pi }\left\{c_{1\pm }+\frac{M_\pi ^2d_{1\pm }}{16\pi ^2F_\pi ^2}+O(M_\pi ^4)\right\},$$
(7.2)
with
$`c_{1+}`$ $`=`$ $`0,c_1=1/3,`$
$`d_{1+}`$ $`=`$ $`8b^r{\displaystyle \frac{1}{648}}(144l(l+2\overline{l}_2)+96l+288\overline{l}_2+113+\mathrm{\Delta }_+),`$
$`d_1`$ $`=`$ $`a_1^r+8b^r+{\displaystyle \frac{1}{648}}(144l(3\overline{l}_\mathrm{\Delta }1)+36(8\overline{l}_13\overline{l}_312\overline{l}_4+12\overline{l}_\mathrm{\Delta })`$
$`+43+\mathrm{\Delta }_{}),`$
$`\mathrm{\Delta }_+`$ $`=`$ $`136431395\pi ^2,\mathrm{\Delta }_{}=3559+351\pi ^2.`$ (7.3)
We have split off the numbers 113 and 43, respectively, to illustrate that these expressions completely agree with the ones displayed in Eq. (8.14) of Ref. , where, as already mentioned, no explicit expressions for the remainders $`\mathrm{\Delta }_{+,}`$ were worked out. For the quadrupole polarizabilities, we obtain
$$(\alpha _2\pm \beta _2)_{\pi ^0}=\frac{\alpha }{16\pi ^2F_\pi ^2M_\pi ^3}\left\{c_{2\pm }+\frac{M_\pi ^2d_{2\pm }}{16\pi ^2F_\pi ^2}+O(M_\pi ^4)\right\},$$
(7.4)
with
$`c_{2+}`$ $`=`$ $`0,c_2=156/45,`$
$`d_{2+}`$ $`=`$ $`{\displaystyle \frac{5009}{27}}+{\displaystyle \frac{13453\pi ^2}{720}}+{\displaystyle \frac{16\overline{l_2}}{45}},`$ (7.5)
$`d_2`$ $`=`$ $`12a_2^r24b^r+{\displaystyle \frac{1}{960}}\left(1280l(16\overline{l}_\mathrm{\Delta })+192161811\pi ^2\right)`$ (7.6)
$`{\displaystyle \frac{4(52\overline{l}_1+5\overline{l}_2+3\overline{l}_378\overline{l}_4+105\overline{l}_\mathrm{\Delta })}{45}}.`$
### 7.2 Numerical results
For numerical evaluations of the polarizabilities we use the values of the LECs given in Section 3.
The results are displayed in Table 1, where we also quote the results from dispersive calculations. The following comments are in order.
1. The slight difference in the value of $`(\alpha _1+\beta _1)_{\pi ^0}`$ with the one reported in is due to the updated values of $`\overline{l}_2`$ and of the pion decay constant $`F_\pi `$ used here.
2. Our results for the dipole polarizabilities as well as for the quadrupole polarizabilities $`(\alpha _2\beta _2)_{\pi ^0}`$ agree with the results of the recent investigations performed in and within the uncertainties quoted.
3. The prediction for the quadrupole polarizability $`(\alpha _2+\beta _2)_{\pi ^0}`$ is positive, in contrast to the result reported in Ref. . The ChPT expression contains as the only LEC $`\overline{l}_2`$, known rather accurately from $`\pi \pi `$ scattering . We come back to this point in the following subsection, where we also discuss the uncertainties quoted in the Table for the ChPT calculation.
4. We plot the helicity amplitudes in Fig. 7. It illustrates the fact that the helicity flip amplitude $`H_+`$ is quite flat at this order, in contrast to the non-flip amplitude $`H_{++}`$, see the values of the quadrupole polarizabilities in Table 1.
5. For a comparison of the ChPT - predictions of the dipole polarizabilities with calculations performed before 1994, and for additional information on these quantities, we refer the interested reader to Refs. and .
### 7.3 Estimating the uncertainties
To estimate the uncertainties in the prediction of the polarizabilities, we first note that the helicity non-flip amplitude $`H_{++}`$ starts out at order $`p^4`$. We have therefore, for this quantity, a leading and next-to-leading order calculation at our disposal. For the corresponding polarizabilities $`(\alpha _1\beta _1)_{\pi ^0}`$ and $`(\alpha _2\beta _2)_{\pi ^0}`$, we thus simply add in quadrature the uncertainties generated by the order $`p^4`$ and $`p^6`$ LECs (see Section 3). The resulting numbers are given in column 2 of Table 1. They do not incorporate an estimate of the higher order contributions.
On the other hand, the helicity flip amplitude $`H_+`$ starts out at order $`p^6`$, and we have determined here only its leading order term. According to Eq. (2.7), this amplitude is proportional to $`B(s,t,u)`$, which is an analytic function of the variables $`s,\nu `$ at the Compton threshold and can be, therefore, expanded in a Taylor series,
$`B(s,t,u)=U+Vs+W\nu ^2+𝒪(s^2,\nu ^4,s\nu ^2).`$ (7.7)
The relation to the polarizabilities is
$`(\alpha _1+\beta _1)_{\pi ^0}=8\alpha M_\pi U,(\alpha _2+\beta _2)_{\pi ^0}=96\alpha M_\pi V.`$ (7.8)
The Taylor coefficients themselves have a chiral expansion of the form
$`U`$ $`=`$ $`{\displaystyle \frac{1}{(16\pi ^2F_\pi ^2)^2}}\left[U_0+{\displaystyle \frac{M_\pi ^2U_1}{16\pi ^2F_\pi ^2}}+𝒪(M_\pi ^4)\right],`$
$`V`$ $`=`$ $`{\displaystyle \frac{1}{(16\pi ^2F_\pi ^2)^2M_\pi ^2}}\left[V_0+{\displaystyle \frac{M_\pi ^2V_1}{16\pi ^2F_\pi ^2}}+𝒪(M_\pi ^4)\right].`$ (7.9)
Whereas LECs from order $`p^6`$ do contribute to $`U_0`$, the leading term $`V_0`$ is a pure loop effect, because $`V_0/M_\pi ^2`$ is not analytic in the pion mass and thus cannot receive contributions from polynomial counterterms. To illustrate this point, and to estimate the size of $`V_1`$, we consider the vector meson exchange amplitudes worked out in . The contribution from $`\omega `$ exchange is dominant and given by
$`B_\omega (s,t,u)={\displaystyle \frac{C_\omega }{2}}\left[{\displaystyle \frac{1}{M_\omega ^2t}}+{\displaystyle \frac{1}{M_\omega ^2u}}\right];C_\omega =0.67\mathrm{GeV}^2.`$ (7.10)
In the language of ChPT , this amplitude starts out at order $`p^6`$. It does contribute to $`U_0`$ \- this term is included in the resonance exchange estimates for the $`𝒪(p^6)`$ LECs in (3). For this reason, we calculate the uncertainties in $`(\alpha _1+\beta _1)_{\pi ^0}`$ as before, with a result that is given in the second row of Table 1. Again, it doe not incorporate an estimate of higher order contributions.
Finally, we come to the estimate of the uncertainty in $`(\alpha _2+\beta _2)_{\pi ^0}`$. In agreement with what is said above, resonance exchange does not contribute to $`V_0`$. On the other hand, there is no reason why this term should dominate the contribution from $`V_1`$. Indeed, if we use (7.10) to also estimate effects from order $`p^8`$, we find with $`M_\omega =782`$ MeV the value $`V_1=2.2`$. The corresponding contribution to the quadrupole polarizability $`(\alpha _2+\beta _2)_{\pi ^0}`$ is $`0.2510^4\mathrm{fm}^5`$ \- of the order needed to bring the ChPT calculation into agreement with the analysis of Ref. .
This result illustrates that the discrepancy between the chiral prediction for the quadrupole polarizability $`(\alpha _2+\beta _2)_{\pi ^0}`$ at order $`p^6`$ and the dispersion analysis in Ref. is of no significance, because the terms neglected may well be much larger than the leading order term, which would only dominate for very small values of the pion mass. On the other hand, to obtain a reliable estimate of $`(\alpha _2+\beta _2)_{\pi ^0}`$ in the framework of ChPT , one needs to perform a reliable calculation of the relevant couplings at order $`p^8`$. This is outside the scope of the present work. For this reason, we do not quote an uncertainty for $`(\alpha _2+\beta _2)_{\pi ^0}`$ in Table 1.
## 8 Summary and outlook
1. We have recalculated the two-loop expression for the amplitude $`\gamma \gamma \pi ^0\pi ^0`$ in the framework of chiral perturbation theory. We have made use of the techniques developed in Ref. , and of the effective Lagrangian $`_6`$ available now .
2. The method has allowed us to evaluate the dipole and quadrupole polarizabilities in closed form. \[As far as we are aware, the quadrupole polarizabilities have never been calculated in ChPT before.\] The two Lorentz invariant amplitudes $`A`$ and $`B`$ are presented as a sum over multiple integrals over Feynman parameters whose numerical evaluation poses no difficulty. This is in contrast to Ref. , where part of the amplitudes, denoted by $`\mathrm{\Delta }_{A,B}`$, were published in numerical form only.
3. Our result agrees with the earlier calculation up the the coefficient in one of the chiral logarithms in the amplitude $`A`$, and up to minute differences in the numerical values of the remainder $`\mathrm{\Delta }_{A,B}`$. The induced changes in the numerics of the cross section and of the dipole polarizabilities are far below the uncertainties generated by the (not precisely known) values of the low energy constants.
4. The values for the dipole and quadrupole polarizabilities are presented in Table 1 and confronted with recent evaluations from data on $`\gamma \gamma \pi ^0\pi ^0`$. There is reasonable agreement for the dipole polarizabilities. As for the quadrupole ones, the combination $`(\alpha _2\beta _2)_{\pi ^0}`$ related to the helicity non-flip amplitude agrees with within the uncertainties quoted. On the other hand, the sum $`(\alpha _2+\beta _2)_{\pi ^0}`$ \- related to the helicity flip amplitude - differs in sign from the one in Ref. . We have shown why this does not contradict the predictions of ChPT : this quantity is a two-loop effect, and one expects from order $`p^8`$ (three loops) substantial corrections to the leading order result. We have indeed identified $`\omega `$-exchange as an important contribution at this order.
5. It would be instructive to improve the estimates for the LECs $`c_{29},\mathrm{},c_{34}`$ in the sense that in these estimates, the constraints from the asymptotics of QCD should be respected.
The corresponding calculation of the charged pion polarizabilities is in progress .
This work was completed while M.A.I. visited the University of Bern. It is a pleasure to thank G. Colangelo, J. Bijnens and B. Moussallam for useful discussions, and S. Bellucci for useful remarks concerning the manuscript. This work was supported by the Swiss National Science Foundation, by RTN, BBW-Contract No. 01.0357, and EC-Contract HPRN–CT2002–00311 (EURIDICE). M.A.I. also appreciates the partial support by the Russian Fund of Basic Research under Grant No. 04-02-17370. Also, partial support by the Academy of Finland, grant 54038, is acknowledged.
## Appendix A Notation
In order to simplify the expressions, we set the pion mass equal to one in all Appendices,
$`M_\pi =1.`$ (1.11)
We use the following notation for $`d`$-dimensional one-loop and two-loop integrals,
$`\mathrm{}`$ $`=`$ $`{\displaystyle \frac{d^dl}{(2\pi )^di}(\mathrm{})},\mathrm{}={\displaystyle \frac{d^dl_1}{(2\pi )^di}\frac{d^dl_2}{(2\pi )^di}(\mathrm{})}.`$ (1.12)
In particular,
$`{\displaystyle \frac{1}{[zl^2]^n}}`$ $`=`$ $`F_n[z],n1,`$
$`F_n[z]`$ $`=`$ $`z^{w+2n}C(w){\displaystyle \frac{\mathrm{\Gamma }(n2w)}{\mathrm{\Gamma }(n)}},w={\displaystyle \frac{d}{2}}2.`$ (1.13)
The measures in the integration over Feynman parameters are defined by
$`d^2x=dx_2dx_3,d^3x=dx_1dx_2dx_3.`$ (1.14)
In dispersion relations, we use the $`d`$-dimensional measure
$`[d\sigma ]`$ $`=`$ $`{\displaystyle \frac{C(w)\mathrm{\Gamma }(3/2)}{\mathrm{\Gamma }(3/2+w)}}\left({\displaystyle \frac{\sigma }{4}}1\right)^w\beta d\sigma ,`$
$`C(w)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^{2+w}}},\beta =\sqrt{14/\sigma }.`$ (1.15)
## Appendix B The acnode
The evaluation of the two-loop vertex and box diagrams (5) and (6) in Fig. 3 is described in Ref. . It is based on a $`d`$-dimensional dispersive representation of the fish-type diagram,
$`J(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{\left[1l^2\right]\left[1(l+p)^2\right]}}=C(w)\mathrm{\Gamma }(w){\displaystyle \underset{0}{\overset{1}{}}}𝑑x[1p^2x(1x)]^w`$ (2.16)
$`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}{\displaystyle \frac{[d\sigma ]}{\sigma p^2}};1.5<\omega <0,`$
where the measure $`[d\sigma ]`$ is given in Appendix A. This representation allows one to reduce the two-loop vertex and box integrals to the one-loop case, with a final integration over the dispersion parameter $`\sigma `$. The ultraviolet divergences can be extracted by invoking recursion relations . While the acnode was treated in a different manner in , we evaluate it here analogously to the vertex and box diagrams just mentioned. This results in considerable simplifications in the numerical programs. For this purpose, we invoke a dispersion relation for the function $`I(m,n;s)`$ defined by
$`I(m,n;s)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{1}{}}}𝑑x\left[1sx(1x)\right]^m\left[x(1x)\right]^n,1<m<0.`$
$`I(m,n;s)`$ is analytic in the complex $`s`$-plane, cut along the real axis for Re $`s4`$. To evaluate its absorptive part, we observe that the imaginary part of the first factor of the integrand in Eq. (LABEL:Imn), evaluated at the upper rim of the cut, is
$`\mathrm{Im}[1sx(1x)]^m=\mathrm{sin}(\pi m)\left[sx(1x)1\right]^m,`$
$`s>4x_{}<x<x_+,x_\pm =(1/2)\left(1\pm \sqrt{14/s}\right).`$ (2.18)
The integrand is symmetric around $`x=1/2`$, so we may restrict the integration from $`x_{}`$ to $`1/2`$ in the evaluation of the absorptive part of $`I(m,n;s)`$. The substitution $`x=(1/2)\left(1\sqrt{(14/s)(1u)}\right)`$ generates a hypergeometric function, and we arrive at the dispersion relation
$`I(m,n;s)`$ $`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\sigma \rho (m,n;\sigma )}{\sigma s}},`$
$`\rho (m,n;\sigma )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(3/2)\delta ^{1/2+m}}{4^n\mathrm{\Gamma }(3/2+m)\mathrm{\Gamma }(m)}}_2F_1({\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{2}}+m+n;{\displaystyle \frac{3}{2}}+m;\delta ),`$
$`\delta `$ $`=`$ $`{\displaystyle \frac{\sigma }{4}}1.`$ (2.19)
In particular, we find
$`\rho (m,0;\sigma )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(3/2)\beta \delta ^m}{\mathrm{\Gamma }(3/2+m)\mathrm{\Gamma }(m)}},`$
$`\rho (m,1;\sigma )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(3/2)\beta \delta ^m}{4\mathrm{\Gamma }(5/2+m)\mathrm{\Gamma }(m)}}\left(1+m+{\displaystyle \frac{2}{\sigma }}\right),`$
$`\rho (m,2;\sigma )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(3/2)\beta \delta ^m}{16\mathrm{\Gamma }(7/2+m)\mathrm{\Gamma }(m)}}\left(2+3m+m^2+{\displaystyle \frac{4(1+m)}{\sigma }}+{\displaystyle \frac{12}{\sigma ^2}}\right),`$
with $`\beta `$ given by Eq. (1.15).
The dispersion relation (2.19) allows us to evaluate the integral that occurs in the evaluation of the acnode diagram (1) in Fig. 5,
$`A_N^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{4l_1^\mu l_2^\nu V_LV_R}{D_1D_2D_3D_4D_5}},`$
$`V_L`$ $`=`$ $`(l_2+q_2l_1)^21,`$
$`V_R`$ $`=`$ $`(l_1+q_1l_2)^21,`$
$`D_1`$ $`=`$ $`1l_1^2,D_2=1(l_1+q_1)^2,`$
$`D_3`$ $`=`$ $`1l_2^2,D_4=1(l_2+q_2)^2,`$
$`D_5`$ $`=`$ $`1(l_1l_2p_1+q_1)^2.`$ (2.21)
We combine $`1/(D_1D_2)`$ and $`1/(D_3D_4)`$ by using $`x_1`$ and $`x_2`$ as Feynman parameters, respectively. By shifting $`l_1`$ and $`l_2`$ and again dropping terms that vanish upon contraction with the polarization vectors, one obtains
$`A_N^{\mu \nu }`$ $`=`$ $`{\displaystyle \underset{0}{\overset{1}{}}}d^2x{\displaystyle \frac{l_2^\nu }{[1l_2^2]^2}}{\displaystyle \frac{l_1^\mu P_N(l_i,p_i,q_i)}{[1l_1^2]^2[1(l_1r)^2]}},`$
$`r`$ $`=`$ $`l_2q,q=x_1q_1+x_2q_2p_1.`$ (2.22)
Here, $`P_N(l_i,p_i,q_i)`$ is a polynomial in the momenta indicated. The integration over $`l_1`$ is performed by using the dispersion relation (2.19) with $`s=r^2`$. Then we proceed in a manner which is similar to the case of the box diagram described in Ref. . The final expression can be written as a combination of the integrals
$`A(i,k,m,n)`$ $`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}[d\sigma ]{\displaystyle \underset{0}{\overset{1}{}}}d^3x\left({\displaystyle \frac{\sigma }{4}}1\right)^{1i}\sigma ^k(1x_3)^mF_n[z_{\mathrm{acn}}],`$
$`z_{\mathrm{acn}}`$ $`=`$ $`x_3^2+(1x_3)\sigma x_3(1x_3)a,`$
$`a`$ $`=`$ $`x_1x_2s+x_1(t1)+x_2(u1),`$ (2.23)
where $`F_n[z]`$ is defined in (A). The integrals $`A(i,k,m,n)`$ are convergent at $`w=0`$ in the case
$`i=1,k=0,m1,n4,`$
$`i=1,k1,m1,n3,`$
$`i=2,k0,m0,n3.`$
To single out the divergent part in the remaining integrals, we invoke recursion relations in the following manner. We perform a partial integration in $`x_3`$, and use
$$𝑑x_3(1x_3)^m=\frac{(1x_3)^{m+1}}{m+1}.$$
Then we express $`(1x_3)\sigma `$ through $`z_{\mathrm{acn}}`$ and obtain
$`(m+3n+\omega )A(i,k,m,n)=`$
$`\text{Div}(i,k,n)n\left[A(i,k,m,n+1)(1+a)A(i,k,m+2,n+1)\right],`$ (2.24)
where
$`\text{Div}(i,k,n)`$ $`=`$ $`4^{3nk+\omega }{\displaystyle \frac{C^2(w)\mathrm{\Gamma }(3/2)}{\mathrm{\Gamma }(3/2+\omega )}}{\displaystyle \frac{\mathrm{\Gamma }(n2\omega )}{\mathrm{\Gamma }(n)}}\times `$
$`B(5/2i+\omega ,n+k4+i2\omega ),`$
$`B(x,y)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(x)\mathrm{\Gamma }(y)}{\mathrm{\Gamma }(x+y)}}.`$ (2.25)
The divergences in $`Div(i,k,n)`$ can be worked out straightforwardly.
The integral $`A(1,0,0,3)`$ must be considered separately. We write
$`A(1,0,0,3)`$ $`=`$ $`D(0,3)+{\displaystyle _4^{\mathrm{}}}[d\sigma ]{\displaystyle \underset{0}{\overset{1}{}}}d^3x\{F_3[z_{\mathrm{acn}}]F_3[y]\},`$
$`y`$ $`=`$ $`x_3^2+(1x_3)\sigma .`$ (2.26)
The divergent quantity $`D(0,3)`$ is worked out in Appendix C.1 of Ref. , whereas the integral on the right-hand side is convergent at $`d=4`$.
This concludes our discussion of the acnode integral (2.22).
## Appendix C The quantities $`\mathrm{\Delta }_A`$ and $`\mathrm{\Delta }_B`$
Here we display the expressions for the quantities $`\mathrm{\Delta }_{A(B)}`$ in Eqs. (5.2) and (5.3).
$`\mathrm{\Delta }_A(s,t,u)={\displaystyle \frac{1}{(4\pi F_\pi )^4}}\left\{\left({\displaystyle \frac{689}{162}}{\displaystyle \frac{4\pi ^2}{9}}\right)+{\displaystyle \frac{15043}{64800}}s\right\}`$
$`+{\displaystyle \frac{1}{(4\pi F_\pi )^4}}{\displaystyle \frac{1}{288}}\left\{F_A^{\mathrm{acn}}(s,t,u)+F_A^{\mathrm{ver}}(s)+F_A^{\mathrm{box}}(s,t,u)\right\},`$ (3.27)
$`\mathrm{\Delta }_B(s,t,u)={\displaystyle \frac{1}{(4\pi F_\pi )^4}}\left\{{\displaystyle \frac{8329}{43200}}+\left({\displaystyle \frac{2987}{1350}}{\displaystyle \frac{2\pi ^2}{9}}\right){\displaystyle \frac{1}{s}}\right\}`$
$`+{\displaystyle \frac{1}{(4\pi F_\pi )^4}}{\displaystyle \frac{1}{288}}\left\{F_B^{\mathrm{acn}}(s,t,u)+F_B^{\mathrm{ver}}(s)+F_B^{\mathrm{box}}(s,t,u)\right\},`$ (3.28)
where
$`F_I^{\mathrm{acn}}`$ $`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}𝑑\sigma \beta {\displaystyle \underset{0}{\overset{1}{}}}d^3x\left\{\left[{\displaystyle \frac{P_{I;\mathrm{acn}}^{(0)}}{y}}+{\displaystyle \frac{P_{I;\mathrm{acn}}^{(1)}}{\sigma }}\right]{\displaystyle \frac{1}{z_{\mathrm{acn}}}}+{\displaystyle \frac{P_{I;\mathrm{acn}}^{(2)}}{z_{\mathrm{acn}}^2}}\right\},`$
$`F_I^{\mathrm{ver}}`$ $`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\sigma \beta }{\sigma }}{\displaystyle \underset{0}{\overset{1}{}}}d^2x{\displaystyle \frac{P_{I;\mathrm{ver}}}{z_{\mathrm{ver}}}},`$
$`F_I^{\mathrm{box}}`$ $`=`$ $`{\displaystyle \underset{4}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\sigma \beta }{\sigma }}{\displaystyle \underset{0}{\overset{1}{}}}d^3x{\displaystyle \underset{n=1}{\overset{2}{}}}\left\{P_{I;\mathrm{box}_+}^{(n)}D_{\mathrm{box}_+}^{(n)}+P_{I;\mathrm{box}_{}}^{(n)}D_{\mathrm{box}_{}}^{(n)}\right\};I=A,B,`$
and
$`D_{\mathrm{box}_\pm }^{(n)}`$ $`=`$ $`{\displaystyle \frac{1}{zn_{\mathrm{box};\mathrm{t}}}}\pm {\displaystyle \frac{1}{z_{\mathrm{box};\mathrm{u}}^n}}.`$
Here $`P_I`$ are polynomials in $`s,\nu =tu`$ and in $`x_i`$. Their explicit expressions are given in Appendix D. The quantity $`z_{\mathrm{acn}}`$ is displayed in Eq. (2.23), $`y`$ is given in (B), and
$`z_{\mathrm{ver}}`$ $`=`$ $`\sigma (1x_3)+x_3^2y_2,y_2=1sx_2(1x_2),`$
$`z_{\mathrm{box};\mathrm{t}}`$ $`=`$ $`B_\mathrm{t}A_\mathrm{t}x_1,`$
$`A_\mathrm{t}`$ $`=`$ $`x_2x_3\left[s(1x_2)x_3+(1t)(1x_3)\right],`$
$`B_\mathrm{t}`$ $`=`$ $`A_\mathrm{t}+z_{\mathrm{ver}},`$
$`z_{\mathrm{box};\mathrm{u}}`$ $`=`$ $`z_{\mathrm{box};\mathrm{t}}|_{tu}.`$ (3.30)
The acnode integrals are easy to evaluate numerically in the physical region for the reaction $`\gamma \gamma \pi \pi `$, because branch points occur at $`t=4,u=4`$ only.
On the other hand, the vertex and box integrals contain branch points at $`s=4`$. In order to evaluate these integrals at $`s4`$, we invoke dispersion relations in the manner described in .
## Appendix D The polynomials $`P_A`$ and $`P_B`$
Here, we display the polynomials $`P_{A(B)}`$ that occur in the expressions $`\mathrm{\Delta }_{A(B)}`$ in Appendix C. We use the abbreviations
$`x_+`$ $`=`$ $`x_1+x_22x_1x_2,x_{}=x_1x_2,`$
$`x_{123}`$ $`=`$ $`(1+x_32x_2x_3)(1x_3+2x_1x_2x_3).`$ (4.31)
### D.1 The polynomials $`P_A`$
$`P_{\mathrm{A};\mathrm{acn}}^{(0)}=\mathrm{\hspace{0.17em}192}x_3(1x_3)(sx_+\nu x_{}),`$
$`P_{\mathrm{A};\mathrm{acn}}^{(1)}=\mathrm{\hspace{0.17em}6}s\nu x_{}(1x_3)`$
$`\times \left[1+8x_3^3+3x_3^44x_+(x_3+3x_3^3+2x_3^4)+2x_{}^2(1+2x_3^33x_3^4)\right]`$
$`6s^2(1x_3)[x_+^2(1+2x_3+14x_3^3+7x_3^4)+6x_{}^2x_3^4`$
$`x_+(1+2x_{}^2+4x_3^3(2+x_{}^2)+3x_3^4(1+2x_{}^2))]`$
$`+6\nu ^2x_{}^2(1x_3)\left[12x_3+2x_3^3+5x_3^412x_+x_3^4\right]`$
$`+12s[2x_+^2(1x_3)^3(1+2x_3+3x_3^2)`$
$`+x_3(2+33x_319x_3^219x_3^3+15x_3^4)`$
$`+x_+(4+2x_363x_3^2+25x_3^3+51x_3^435x_3^5)`$
$`+2x_{}^2x_3(2+9x_37x_3^2+5x_3^33x_3^4)]`$
$`+12\nu x_{}(1x_3)[42x_35x_3^24x_3^3+19x_3^4`$
$`+2x_+(1+8x_3^321x_3^4)]`$
$`48x_3[4+4x_3+6x_3^219x_3^3+10x_3^4`$
$`+x_+(23x_319x_3^2+41x_3^321x_3^4)],`$
$`P_{\mathrm{A};\mathrm{acn}}^{(2)}=2s\nu x_{}(1x_3)^2[118x_38x_3^224x_3^39x_3^4`$
$`+12x_+(4+x_3+3x_3^3+2x_3^4)2x_{}^2(11+8x_3+8x_3^2+6x_3^39x_3^4)]`$
$`2s^2(1x_3)^2[x_+^2(3518x_3+34x_3^3+21x_3^4)+18x_{}^2x_3^4`$
$`x_+(11+16x_3^3+9x_3^4+2x_{}^2(112x_3^3+9x_3^4))]`$
$`+2\nu ^2x_{}^2(1x_3)^2\left[13+18x_32x_3^3+(1536x_+)x_3^4\right]`$
$`+4s(1x_3)[x_3(22+11x_3+83x_3^2101x_3^3+45x_3^4)`$
$`2x_+^2(1x_3)^2\left(11(1+x_3+x_3^2)9x_3^3\right)`$
$`+x_+(60178x_3+43x_3^2173x_3^3+233x_3^4105x_3^5)`$
$`+2x_{}^2(2x_3)x_3(11(1+x_3+x_3^2)+9x_3^3)]`$
$`+4\nu x_{}(1x_3)^2[60+74x_3+9x_3^256x_3^3+57x_3^4`$
$`+x_+(22+104x_3^3126x_3^4)]`$
$`+16x_3(1x_3)[60+18x_3108x_3^2+93x_3^330x_3^4`$
$`+x_+(2283x_3+205x_3^2187x_3^3+63x_3^4)],`$
$`P_{\mathrm{A};\mathrm{ver}}=128sx_2^2(12x_2)x_3^4\left[32x_2(215x_3)x_318x_3^215x_2^2x_3^2\right]`$
$`32s^2x_2^2(12x_2)(3042x_2+7x_2^2)x_3^6,`$
$`P_{\mathrm{A};\mathrm{box}_+}^{(1)}=96sx_2x_3^2[6+9x_3+20x_2x_322x_1x_2x_3`$
$`2x_3^240x_2x_3^2+6x_1x_2x_3^24x_2^2x_3^2+52x_1x_2^2x_3^28x_1^2x_2^2x_3^2`$
$`15x_3^3+58x_2x_3^3+14x_1x_2x_3^312x_2^2x_3^372x_1x_2^2x_3^32x_1^2x_2^2x_3^3`$
$`+8x_1x_2^3x_3^3+4x_1^2x_2^3x_3^3+20x_3^448x_2x_3^436x_1x_2x_3^4+20x_2^2x_3^4`$
$`+76x_1x_2^2x_3^4+10x_1^2x_2^2x_3^424x_1x_2^3x_3^424x_1^2x_2^3x_3^4+24x_1^2x_2^4x_3^4]`$
$`16s^2x_2^2x_3^4[12+12x_112x_224x_1x_2+12x_1^2x_2+30x_1x_3`$
$`24x_2x_390x_1x_2x_3+27x_1^2x_2x_3+30x_2^2x_3+66x_1x_2^2x_3`$
$`24x_1^2x_2^2x_310x_1^3x_2^2x_324x_3^2+6x_1x_3^2+66x_2x_3^218x_1x_2x_3^2`$
$`+45x_1^2x_2x_3^248x_2^2x_3^2+48x_1x_2^2x_3^2120x_1^2x_2^2x_3^228x_1^3x_2^2x_3^2`$
$`24x_1x_2^3x_3^2+36x_1^2x_2^3x_3^2+56x_1^3x_2^3x_3^2]`$
$`+48\nu ^2x_2^3x_3^4[4+8x_18x_1^2+2x_1x_3+7x_1^2x_3+10x_2x_3`$
$`30x_1x_2x_3+30x_1^2x_2x_310x_1^3x_2x_3+6x_3^214x_1x_3^2+3x_1^2x_3^2`$
$`12x_2x_3^2+12x_1x_2x_3^2+12x_1^2x_2x_3^212x_1^3x_2x_3^2+24x_1x_2^2x_3^2`$
$`48x_1^2x_2^2x_3^2+24x_1^3x_2^2x_3^2]`$
$`+384x_2x_3^2[6+9x_3+19x_2x_319x_1x_2x_3+7x_3^250x_2x_3^2`$
$`2x_1x_2x_3^2+52x_1x_2^2x_3^215x_3^3+44x_2x_3^3+16x_1x_2x_3^3`$
$`60x_1x_2^2x_3^3+8x_3^416x_2x_3^416x_1x_2x_3^4+32x_1x_2^2x_3^4],`$
$`P_{\mathrm{A};\mathrm{box}_+}^{(2)}=24s\nu ^2(1x_1)^2x_2^3(1+x_2+x_1x_22x_1x_2^2)x_3^6x_{123}`$
$`+96sx_2x_3^4x_{123}[7x_27x_1x_2+14x_1x_2^2+3x_3+6x_2x_3`$
$`+6x_1x_2x_312x_1x_2^2x_33x_3^22x_2x_3^22x_1x_2x_3^2+4x_1x_2^2x_3^2]`$
$`48s^2x_2^2x_3^4x_{123}[22x_1x_2+6x_1x_2x_1^2x_2+3x_3+3x_1x_3`$
$`6x_1x_2x_34x_3^24x_1x_3^2+2x_2x_3^2+8x_1x_2x_3^2`$
$`+2x_1^2x_2x_3^24x_1x_2^2x_3^24x_1^2x_2^2x_3^2+4x_1^2x_2^3x_3^2]`$
$`24s^3x_2^3(1+x_12x_1x_2)(1+x_1x_2x_1^2x_2)x_3^6x_{123}`$
$`48\nu ^2(1x_1)^2x_2^3(1x_3)x_3^4(1+4x_3)x_{123}+2x_2(1x_3)^2x_3^4x_{123},`$
$`P_{\mathrm{A};\mathrm{box}_{}}^{(1)}=\mathrm{\hspace{0.17em}32}s\nu x_2^2x_3^4[6+6x_1+12x_1x_218x_1^2x_2`$
$`15x_1x_3+12x_2x_3+33x_1^2x_2x_330x_1x_2^2x_3+21x_1^2x_2^2x_3`$
$`10x_1^3x_2^2x_3+12x_3^2+3x_1x_3^224x_2x_3^230x_1x_2x_3^2`$
$`+18x_1^2x_2x_3^2+6x_2^2x_3^2+54x_1x_2^2x_3^248x_1^2x_2^2x_3^2+4x_1^3x_2^2x_3^2]`$
$`96\nu x_2^2x_3^3[9+11x_1+12x_314x_1x_3+24x_2x_3`$
$`56x_1x_2x_3+28x_1^2x_2x_3+15x_3^227x_1x_3^258x_2x_3^2`$
$`+84x_1x_2x_3^2+8x_1^2x_2x_3^2+60x_1x_2^2x_3^280x_1^2x_2^2x_3^2`$
$`16x_3^3+28x_1x_3^3+32x_2x_3^320x_1x_2x_3^342x_1^2x_2x_3^3`$
$`64x_1x_2^2x_3^3+84x_1^2x_2^2x_3^3],`$
$`P_{\mathrm{A};\mathrm{box}_{}}^{(2)}=48s\nu (1x_1)x_2^2x_3^4x_{123}`$
$`\times [2+3x_3+3x_2x_3+3x_1x_2x_36x_1x_2^2x_34x_3^2`$
$`2x_2x_3^22x_1x_2x_3^2+4x_1x_2^2x_3^2]`$
$`+24s^2\nu (1x_1^2)x_2^3(2x_2x_1x_2)x_3^6x_{123}`$
$`+\nu (1x_1)x_2^2(1x_3)x_3^4(12x_3)x_{123}`$
$`+24\nu ^3(1x_1)^3x_2^4x_3^6x_{123}.`$
### D.2 The polynomials $`P_B`$
$`P_{\mathrm{B};\mathrm{acn}}^{(0)}=96x_3(1x_3)\left[{\displaystyle \frac{\nu }{s}}x_{}x_+\right],`$
$`P_{\mathrm{B};\mathrm{acn}}^{(1)}={\displaystyle \frac{6\nu }{s}}x_{}\left[42x_3+21x_3^219x_3^3+3x_3^4+5x_3^5\right]`$
$`{\displaystyle \frac{3\nu ^2}{s}}x_{}^2(1x_3)[1+2x_3+2x_3^3+x_3^4]`$
$`{\displaystyle \frac{24}{s}}x_3(2x_3^2)[28x_3+5x_3^2]+3s(1x_3)`$
$`\times \left[x_+^2(1x_3)^3(1+x_3)x_+(1+8x_3^3+3x_3^4)+6x_{}^2x_3^4\right]`$
$`+6[x_+(1x_3)^2(42x_33x_3^2+13x_3^3)`$
$`+x_3(233x_3+19x_3^2+19x_3^315x_3^4)],`$
$`P_{\mathrm{B};\mathrm{acn}}^{(2)}={\displaystyle \frac{2\nu }{s}}x_{}(1x_3)`$
$`\times \left[34+152x_343x_3^2+21x_3^327x_3^4+15x_3^5\right]`$
$`{\displaystyle \frac{\nu ^2}{s}}x_{}^2(1x_3)^2\left[3518x_32x_3^3+3x_3^4\right]`$
$`+{\displaystyle \frac{8}{s}}x_3(1x_3)\left[3467x_3+95x_3^256x_3^3+15x_3^4\right]`$
$`s(1x_3)^2[x_+^2(1x_3)^2(13+8x_3+3x_3^2)`$
$`+x_+(11+16x_3^3+9x_3^4)18x_{}^2x_3^4]`$
$`+\nu x_{}(1x_3)^2\left[11+16x_3^3+9x_3^4+12x_+(43x_3x_3^4)\right]`$
$`+2(1x_3)[x_3(22+11x_3+83x_3^2101x_3^3+45x_3^4)`$
$`+x_+(34108x_3+65x_3^2+73x_3^3103x_3^4+39x_3^5)],`$
$`P_{\mathrm{B};\mathrm{ver}}=16sx_2^2(12x_2)x_3^5\left[33x_2^2x_36(45x_3)2x_2(4+9x_3)\right]`$
$`1536x_2^2(12x_2)(1x_3)(12x_3)x_3^4,`$
$`P_{\mathrm{B};\mathrm{box}_+}^{(1)}={\displaystyle \frac{24\nu ^2}{s}}x_2^3(1x_3)x_3^4\left[4+8x_110x_32x_1x_33x_1^2x_3\right]`$
$`{\displaystyle \frac{1152}{s}}x_2(1x_3)^2x_3^2(12x_3+2x_3^2)`$
$`24sx_2^2x_3^4[4+12x_14x_212x_314x_1x_3+10x_2x_3`$
$`14x_1x_2x_3+17x_1^2x_2x_314x_1^2x_2^2x_3+8x_3^22x_1x_3^26x_2x_3^2`$
$`+14x_1x_2x_3^23x_1^2x_2x_3^2+4x_1x_2^2x_3^212x_1^2x_2^2x_3^2+12x_1^2x_2^3x_3^2]`$
$`48x_2(1x_3)x_3^2[6+19x_3+16x_1x_2x_327x_3^22x_2x_3^2`$
$`18x_1x_2x_3^228x_1x_2^2x_3^2+20x_3^332x_1x_2x_3^3+64x_1x_2^2x_3^3],`$
$`P_{\mathrm{B};\mathrm{box}_+}^{(2)}={\displaystyle \frac{24\nu ^2}{s}}(1x_1)^2x_2^3(1x_3)^3x_3^4(1+4x_3)`$
$`+{\displaystyle \frac{1}{s}}x_2(1x_3)^4x_3^4`$
$`24sx_2^2(1x_3)^2x_3^4[22x_1x_2+6x_1x_2x_1^2x_2`$
$`+3x_3+3x_1x_36x_1x_2x_34x_3^24x_1x_3^2+2x_2x_3^2`$
$`+8x_1x_2x_3^2+2x_1^2x_2x_3^24x_1x_2^2x_3^24x_1^2x_2^2x_3^2+4x_1^2x_2^3x_3^2]`$
$`12s^2x_2^3(1+x_12x_1x_2)(1+x_1x_2x_1^2x_2)(1x_3)^2x_3^6`$
$`+12\nu ^2(1x_1)^2x_2^3(1x_2x_1x_2+2x_1x_2^2)(1x_3)^2x_3^6`$
$`+48x_2(1x_3)^2x_3^4[7x_27x_1x_2+14x_1x_2^2+3x_3`$
$`+6x_2x_3+6x_1x_2x_312x_1x_2^2x_33x_3^22x_2x_3^2`$
$`2x_1x_2x_3^2+4x_1x_2^2x_3^2],`$
$`P_{\mathrm{B};\mathrm{box}_{}}^{(1)}={\displaystyle \frac{48\nu }{s}}x_2^2(1x_3)^2x_3^3(5+11x_128x_34x_1x_3)`$
$`+48\nu x_2^2x_3^4[2+2x_1+8x_1x_26x_33x_1x_32x_2x_3`$
$`4x_1x_2x_3+3x_1^2x_2x_312x_1x_2^2x_33x_1^2x_2^2x_3+4x_3^2`$
$`+x_1x_3^2+2x_2x_3^210x_1x_2x_3^2+18x_1x_2^2x_3^2],`$
$`P_{\mathrm{B};\mathrm{box}_{}}^{(2)}={\displaystyle \frac{144\nu }{s}}(1x_1)x_2^2(1x_3)^3x_3^4(12x_3)`$
$`+{\displaystyle \frac{12\nu ^3}{s}}(1x_1)^3x_2^4(1x_3)^2x_3^6`$
$`12s\nu (1x_1^2)x_2^3(2+x_2+x_1x_2)(1x_3)^2x_3^6`$
$`+24\nu (1x_1)x_2^2(1x_3)^2x_3^4[2+3x_3+3x_2x_3+3x_1x_2x_3`$
$`6x_1x_2^2x_34x_3^22x_2x_3^22x_1x_2x_3^2+4x_1x_2^2x_3^2].`$
|
warning/0506/math0506050.html
|
ar5iv
|
text
|
# Simple subalgebras of simple special Jordan algebras
## 1 Introduction
This paper provides a classification of simple subalgebras in finite-dimensional special simple Jordan algebras over algebraically closed field $`F`$ with characteristic unequal to 2. More precisely, we will determine the canonical forms for any simple subalgebras of special simple Jordan algebras, and the number of conjugate classes corresponding to the given simple Jordan subalgebra. In particular, a Jordan algebra of any type can be realized as a Jordan subalgebra of symmetric or symplectic matrices of an appropriate order.
In 1987 N. Jacobson determined the orbits under the orthogonal group $`O(n)`$ of the subalgebras of the Jordan algebra of $`n\times n`$ real symmetric matrices ().
The paper significantly relies on the description of maximal subalgebras of finite-dimensional special simple Jordan algebras obtained by M. Racine in 1974 ().
In this paper we consider simple Jordan algebras presented in canonical matrix realizations, that is, $`H(F_n)`$ can be viewed as the algebra of all symmetric matrices in $`F_n`$ with respect to the ordinary transpose; $`F_n^{(+)}`$ is the set of all matrices of order $`n`$ under the circle product $`AB=\frac{AB+BA}{2}`$; $`H(F_{2n},j)`$ where $`j`$ denotes a symplectic involution consists of all matrices of order $`2n`$ of the form
$$\left(\begin{array}{cc}A& B\\ C& A^t\end{array}\right),$$
where $`B`$, $`C`$ are any skew-symmetric matrices of order $`n`$, and $`A`$ is any matrix of order $`n`$. If $`f`$ is a non-singular symmetric bilinear form on a vector space $`V`$, then $`𝒥=FV`$ is a Jordan algebra of the type $`J(f,1)`$.
Throughout the paper we assume that the base field $`F`$ is algebraically closed with characteristic not two.
## 2 Subalgebras
### 2.1 Matrix subalgebras
Let $`𝒥`$ be a simple Jordan algebra of the type $`F_{\frac{n}{2}}^{(+)}`$ where $`n`$ is even. Then it can always be presented as a subalgebra of $`H(F_n)`$ as follows
$$\left\{\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)\right\},$$
$`(1)`$
where $`A`$ is any symmetric matrix of order $`\frac{n}{2}`$ and $`B`$ is any skewsymmetric matrix of order $`\frac{n}{2}`$.
###### Lemma 2.1.
Any automorphism of a Jordan algebra of the form (1) is induced by an automorphism of $`H(F_n)`$.
###### Proof.
Any automorphism of $`𝒥`$ can be extended to an automorphism or antiautomorphism of a special universal enveloping algebra $`U(𝒥)`$ which is isomorphic to $`F_{\frac{n}{2}}F_{\frac{n}{2}}`$ (see ). Notice that in this particular case the associative enveloping algebra $`S(𝒥)`$ is isomorphic to $`U(𝒥)`$ because from the explicit form (1) $`S(𝒥)`$ consists of all matrices of the form:
$$\left\{\left(\begin{array}{cc}X& Y\\ Y& X\end{array}\right)\right\},$$
where $`X`$ and $`Y`$ are any matrices of order $`\frac{n}{2}`$. Since any automorphism of $`F_{\frac{n}{2}}F_{\frac{n}{2}}`$ either induces non-trivial automorphisms of these ideals or sends one ideal onto another, it can be lifted up to an inner automorphism of the entire matrix algebra $`F_n`$. Consequently, for any antiautomorphism of $`F_{\frac{n}{2}}F_{\frac{n}{2}}`$ we can choose an automorphism (not necessarily non-trivial) of $`F_{\frac{n}{2}}F_{\frac{n}{2}}`$ such that their composition induces non-trivial antiautomorphisms of simple ideals. Therefore, any (Jordan) automorphism of $`𝒥`$ can be written as follows:
$$\phi (X)=Q^1XQ$$
or
$$\phi (X)=Q^1X^tQ,$$
$`(2)`$
for some non-singular matrix $`Q`$.
The next step is to prove that $`\phi `$ is orthogonal. In other words, all we have to show is that for any automorphism $`\phi `$ of $`𝒥`$, we can choose $`Q`$ such that (2) holds and $`Q^tQ=I`$ where $`I`$ is the identity matrix. Since $`𝒥`$ is a subalgebra of $`H(F_n)`$, for each $`X`$ in $`𝒥`$, $`(Q^1XQ)^t=Q^1XQ`$, $`Q^tX(Q^1)^t=Q^1XQ`$, $`QQ^tX=XQQ^t`$. Denote $`B=QQ^t`$. Next we are going to show that $`B`$ is actually a scalar multiple of the identity matrix. We are given that $`BX=XB`$ where $`X`$ is any matrix of the form (1). Let us write $`B`$ as follows:
$$B=\left(\begin{array}{cc}B_1& B_2\\ B_3& B_4\end{array}\right)$$
where $`B_i`$ are matrices of order $`\frac{n}{2}`$. By performing the matrix multiplication, we obtain the $`B_2=B_3=0`$, and $`B_1=B_4=\alpha I`$, for some non-zero $`\alpha `$. Since the ground field $`F`$ is algebraically closed, we can choose $`\beta F`$ such that $`\alpha =\beta ^2`$. Set $`Q^{}=\beta ^1Q`$. Obviously, $`Q^{}`$ determines the same automorphism as $`Q`$ does, and $`Q^tQ^{}=I`$.
The Lemma is proved. ∎
###### Lemma 2.2.
Let $`𝒥F_{\frac{n}{2}}^{(+)}`$ be a subalgebra of $`H(F_n)`$. Then, by an appropriate automorphism of $`H(F_n)`$, $`𝒥`$ can always be reduced to the form (1).
###### Proof.
Since $`𝒥`$ has the type $`F_{\frac{n}{2}}^{(+)}`$, by some (not necessarily orthogonal) automorphism $`\phi `$ of $`F_n^{(+)}`$ we can always bring $`𝒥`$ to the following form (see and )
$$\left\{\left(\begin{array}{cc}X& 0\\ 0& X^t\end{array}\right)\right\}$$
$`(3)`$
where $`X`$ is any matrix of order $`\frac{n}{2}`$. Then, it is easily seen that $`\theta (Y)=S^1YS`$, where $`S=\left(\begin{array}{cc}I& iI\\ \frac{1}{2}I& \frac{i}{2}I\end{array}\right)`$, $`I`$ is the identity matrix, $`i^2=1`$, sends each element of the form (3) into the algebra of the form (1). Therefore, by $`\chi =\theta \phi `$ we can bring $`𝒥`$ to the form (1).
Next we will show that $`\chi `$ is actually an orthogonal automorphism. Notice that $`\chi `$ sends $`H(F_n)`$ onto a Jordan subalgebra of $`F_n^{(+)}`$ which consists of all matrices symmetric with respect to the following involution: $`j^{}=\chi t\chi ^1`$ where $`t`$ is the standard transpose involution. This involution can be rewritten as follows $`j^{}(X)=C^1X^tC`$ for some non-singular symmetric matrix $`C`$ of order $`n`$. It follows from the above considerations that any matrix of the form (1) is symmetric with respect to $`j^{}`$. Equivalently, for any $`Y`$ of the form (1), $`C^1Y^tC=Y`$, $`Y^tC=CY`$, $`YC=CY`$ because $`Y`$ is symmetric. As proved in the previous Lemma, $`C=\alpha I`$ for some non-zero $`\alpha `$. Therefore, $`j^{}=t`$, and $`\chi (H(F_n))=H(F_n)`$, and $`\chi `$ is actually an automorphism of $`H(F_n)`$. Hence, the Lemma is proved. ∎
###### Lemma 2.3.
Let $`𝒜`$ be a special simple matrix Jordan algebra, and $`𝒥`$ be a proper simple subalgebra of $`𝒜`$. Denote a maximal subalgebra which contains $`𝒥`$ as $`M`$. Next, consider a Wedderburn splitting $`M=SR`$ where $`S`$ is a semisimple algebra, $`R`$ is the radical. Then, there exists an automorphism $`\phi `$ of $`𝒜`$ such that $`\phi (𝒥)S`$.
###### Proof.
Let 1 be the identity element of $`𝒜`$, and $`1𝒥`$. According to , if $`𝒥`$ is special and the degree of $`𝒥`$ is not divisible by the characteristic, then $`𝒥`$ is conjugate under an inner automorphism $`T`$ of $`M`$ to some subalgebra of $`S`$, and $`T`$ is a composition of the standard automorphisms $`T_{x,y}`$ that can be represented in associative terms as follows
$$T_{x,y}(a)=tat^1,$$
$`(4)`$
where $`t=u^{\frac{1}{2}}(1xy)(1+yx)`$, $`u=(1xy)(1+yx)(1+xy)(1yx),`$ $`x,yM`$.
Let $`x`$, $`y`$ be symmetric with respect to an involution $`j`$ of $`𝒜`$: $`j(x)=x`$, $`j(y)=y`$. Then, it is obvious that
$$j(u)=u,\text{and}j(t)=t^1.$$
$`(5)`$
If $`𝒜=F_n^{(+)}`$, then from the explicit form (4) $`T_{x,y}`$ is easily extendable to $`𝒜`$. If $`𝒜=H(F_n)`$, then, because of (5), $`T_{x,y}`$ is orthogonal, therefore, extendable to $`𝒜`$. If $`𝒜=H(F_{2n},j)`$, then, because of (5), $`T_{x,y}`$ is symplectic, therefore, extendable to $`𝒜`$.
If $`𝒥`$ is special and the degree of $`𝒥`$ is divisible by characteristic, then $`T`$ is a generalized inner automorphism (see ), that is, $`T`$ is a composition of an automorphisms $`T_{x_1,\mathrm{},x_n,m}`$ of the form
$$T_{x_1,\mathrm{},x_n,m}=U_v^1(I+V_{x_1,\mathrm{},x_n,m}+U_{x_1}\mathrm{}U_{x_n}U_m)(I+V_{m,x_n\mathrm{},x_1}+U_mU_{x_n}\mathrm{}U_{x_1})$$
where $`v,x_iM`$, $`mR`$. In associative terms operators take the form:
$$U_v(a)=vav,$$
$$U_{x_i}(a)=x_iax_i,$$
$$V_{x_1,\mathrm{},x_n,m}(a)=x_1\mathrm{}x_nma+amx_n\mathrm{}x_1.$$
Hence, if all $`x_i`$, $`m`$ and $`a`$ are symmetric with respect to an involution of $`𝒜`$, then $`j(T_{x_1,\mathrm{},x_n,m}(a))=T_{x_1,\mathrm{},x_n,m}(a)`$. Therefore, $`T_{x_1,\mathrm{},x_n,m}`$ as well as $`T`$ is extendable to $`𝒜`$. The Lemma is proved. ∎
###### Lemma 2.4.
Let $`𝒥F_n^{(+)}`$ be a subalgebra of $`H(F_{2n},j)`$. Then, by an appropriate automorphism of $`H(F_{2n},j)`$, $`𝒥`$ can always be reduced to the following form
$$\left\{\left(\begin{array}{cc}X& 0\\ 0& X^t\end{array}\right)\right\}$$
$`(6)`$
###### Proof.
Since $`𝒥`$ has the type $`F_n^{(+)}`$, by some automorphism $`\phi `$ of $`F_{2n}^{(+)}`$, $`𝒥`$ can be brought to the form (6) (see and ). Notice that $`\phi `$ sends $`H(F_{2n},j)`$ onto a Jordan subalgebra of $`F_{2n}^{(+)}`$ which consists of all matrices symmetric with respect to the following involution: $`j^{}=\phi j\phi ^1`$. This involution can be rewritten as follows $`j^{}(Y)=C^1Y^tC`$ for some non-singular skew-symmetric matrix $`C`$ of order $`2n`$. It follows from the above considerations that any matrix of the form (6) is symmetric with respect to $`j^{}`$. Equivalently, for any $`Y`$ of the form (6), $`C^1Y^tC=Y`$, $`Y^tC=CY`$. Acting in the same manner as above, we can show that $`C=\alpha \left(\begin{array}{cc}0& I_n\\ I_n& 0\end{array}\right)`$ for some non-zero $`\alpha `$, where $`I_n`$ denotes the identity matrix of order $`n`$. Therefore, $`\phi (H(F_n,j))=H(F_n,j)`$, and $`\phi `$ is an automorphism of $`H(F_n,j)`$. Hence, the Lemma is proved. ∎
###### Definition 2.5.
Subalgebras $`𝒥_1`$ and $`𝒥_2`$ of a Jordan algebra $`𝒜`$ are said to be equivalent if there exists an automorphism $`\phi `$ of $`𝒜`$ such that either $`𝒥_1=\phi (𝒥_2)`$ or $`𝒥_2=\phi (𝒥_1)`$.
###### Definition 2.6.
Let $`𝒥`$ be a subalgebra of $`𝒜`$. Then the set $`C(𝒥)`$ of all subalgebras equivalent to $`𝒥`$ in $`𝒜`$ is said to be a conjugate class of $`𝒥`$.
Canonical realizations of simple subalgebras
Let $`𝒜`$ be a simple Jordan algebra, and $`𝒥`$ be a simple subalgebra of $`𝒜`$. All realizations listed below we will call canonical.
1. Let $`𝒜=F_n^{(+)}`$
(1.1) $`𝒥F_m^{(+)}`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,X^t,\mathrm{},X^t,0,\mathrm{},0)\}`$ where $`X`$ is any matrix of order $`m`$.
(1.2) $`𝒥H(F_m)`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,0,\mathrm{},0)\}`$ where $`X`$ is any symmetric matrix of order $`m`$.
(1.3) $`𝒥H(F_{2m},j)`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,0,\mathrm{},0)\}`$ where $`X`$ is any symplectic matrix of order $`2m`$.
2. Let $`𝒜=H(F_n)`$
(2.1) $`𝒥F_m^{(+)}`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,0\mathrm{},0)\}`$ where $`X`$ is of the form (1) in which $`A`$ and $`B`$ are of order $`m`$.
(2.2) $`𝒥H(F_m)`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,0,\mathrm{},0)\}`$ where $`X`$ is any symmetric matrix of order $`m`$.
(2.3) $`𝒥H(F_{2m},j)`$, $`𝒥=\{\text{diag}(X,\mathrm{},X,0\mathrm{},0)\}`$
$$X=\left(\begin{array}{cccc}A& B& C& D\\ B& A& D& C\\ C& D& A& B\\ D& C& B& A\end{array}\right)$$
where $`A`$ is a symmetric matrix of order $`m`$, $`B`$, $`C`$, $`D`$ are skew-symmetric matrices of order $`m`$.
3. Let $`𝒜=H(F_{2n},j)`$
(3.1) $`𝒥F_m^{(+)}`$,
$$𝒥=\{\text{diag}(\underset{k}{\underset{}{X,\mathrm{},X,}}\underset{l}{\underset{}{X^t,\mathrm{},X^t,}}\underset{s}{\underset{}{0,\mathrm{},0}},X^t,\mathrm{},X^t,X,\mathrm{},X,0,\mathrm{},0)\}$$
where $`k+l+s=n`$, $`X`$ is any matrix of order $`m`$.
(3.2) $`𝒥H(F_m)`$, $`𝒥=\{\text{diag}(\underset{k}{\underset{}{X,\mathrm{},X,}}\underset{l}{\underset{}{0,\mathrm{},0}},X,\mathrm{},X,0,\mathrm{},0)\}`$ where $`k+l=n`$, $`X`$ is any symmetric matrix of order $`m`$.
(3.3) $`𝒥H(F_{2m},j)`$, $`𝒥=\left\{\left(\begin{array}{cc}A& B\\ C& A^t\end{array}\right)\right\}`$
$$A=\left(\begin{array}{cccccccc}X& & & & & & & \\ & \mathrm{}& & & & 0& & \\ & & X& & & & & \\ & & & X& Y& & & \\ & & & Z& X^t& & & \\ & & 0& & & \mathrm{}& & \\ & & & & & & X& Y\\ & & & & & & Z& X^t\end{array}\right),$$
$$B=\left(\begin{array}{cccccccc}Y& & & & & & & \\ & \mathrm{}& & & & 0& & \\ & & Y& & & & & \\ & & & & & & & \\ & & & & & & & \\ & & 0& & & 0& & \\ & & & & & & & \end{array}\right),C=\left(\begin{array}{cccccccc}Z& & & & & & & \\ & \mathrm{}& & & & 0& & \\ & & Z& & & & & \\ & & & & & & & \\ & & & & & & & \\ & & 0& & & 0& & \\ & & & & & & & \end{array}\right)$$
where $`X`$ is any matrix of order $`m`$, $`Y`$,$`Z`$ are skew-symmetric matrices of order $`m`$.
###### Definition 2.7.
Let $`𝒥`$ and $`𝒥^{}`$ be two proper subalgebras of $`𝒜`$, and $`𝒥^{}`$ be given in the canonical realization. If $`𝒥`$ is equivalent to $`𝒥^{}`$, then $`𝒥^{}`$ is said to be the canonical form of $`𝒥`$.
###### Theorem 2.8.
Let $`𝒜`$ be a simple matrix Jordan algebra. Then, any simple matrix subalgebra of $`𝒜`$ has a unique canonical form as above.
###### Proof.
Let $`𝒥`$ be any proper simple matrix subalgebra of $`𝒜`$. In particular, the degree of $`𝒥3`$. Denote the identity of $`𝒜`$ as 1.
The proof of the Theorem consists of three cases.
Case 1 $`𝒜=F_n^{(+)}`$
1.1 Let $`𝒥`$ be of the type $`F_m^{(+)}`$ for some $`m<n`$. Since any Jordan algebra of this type has precisely two non-equivalent irreducible representations both of which have degree $`m`$ (see ), $`𝒥`$ is equivalent to the subalgebra in the canonical realization (1.1). If $`1𝒥`$, then the last zero in (1.1) is omitted. Next we are going to show that $`𝒥`$ has a unique canonical form. Equivalently, if $`l`$ and $`k`$ are the number of $`X`$-blocks and $`X^t`$-blocks in (1.1), then $`|lk|`$ is an invariant for $`𝒥`$. Indeed, let $`S(𝒥)`$ be a simple algebra. Then, either $`l`$ or $`k`$ is zero, and $`|lk|=\text{rk}e`$ where $`e`$ is the identity of $`𝒥`$. If $`S(𝒥)`$ is a non-simple semisimple algebra, then $`S(𝒥)=_1_2`$ where $`_i`$ are simple ideals with the identity elements $`e_i`$. Hence, $`|lk|=|\text{rk}e_1\text{rk}e_2|`$ which is invariant for $`𝒥`$. Denote $`k_𝒜(𝒥)=|lk|`$.
1.2 Let $`𝒥`$ be of the type $`H(F_m)`$ for some $`mn`$. Then, it follows from the uniqueness of the irreducible representation of $`H(F_m)`$ (see ) that $`𝒥`$ is equivalent to the subalgebra in the canonical realization (1.2). If $`1𝒥`$, then the last zero in (1.2) is omitted.
1.3 The proof of the case when $`𝒥H(F_{2m},j)`$, $`2m<n`$, is exactly the same as the previous proof. In particular, $`𝒥`$ of the type $`H(F_{2m},j)`$ is equivalent to the subalgebra in the canonical realization (1.3). Obviously, the canonical form is unique.
Case 2 $`𝒜=H(F_n)`$
Here, our main goal is to determine the canonical form of any simple matrix Jordan subalgebra of $`H(F_n)`$. Let $`M`$ be a maximal subalgebra of $`H(F_n)`$. According to , $`M`$ is isomorphic to one of the following:
1. $`H(F_k)H(F_l)`$, $`k+l=n`$,
2. $`F_k^{(+)}H(F_l)R`$, $`2k+l=n`$, $`R`$ is the radical (if $`l=0`$, then $`MF_{\frac{n}{2}}^{(+)}R`$)
3. $`J(f,1)`$ only if $`n=2^m`$ and either $`dimJ(f,1)=2(m+1)`$, $`m`$ is even, or $`dimJ(f,1)=2m+1`$, $`m`$ is odd.
First, assume that $`𝒥`$ is a simple matrix subalgebra of $`H(F_n)`$ such that $`1𝒥`$. Then there exists a maximal subalgebra $`M`$ such that $`𝒥M`$. Since $`\text{deg}𝒥3`$, $`M`$ cannot be of the type 3. If $`M`$ contains a non-zero radical, that is, $`M=SR`$, where $`S`$ a semisimple algebra, $`R`$ the radical, then by Lemma 2.3 we can assume that $`𝒥S`$. If $`S=S_1S_2`$ where $`S_i`$ non-trivial simple ideals, we can choose three orthogonal idempotents (see ): $`e`$, $`e^t`$, $`ff^t`$, $`1=e+e^t+ff^t`$ such that
$$S_1=ff^tH(F_n)ff^t,S_2=eF_ne+e^tF_n^te^t$$
$`(7)`$
Since $`ff^t`$ is an element of $`H(F_n)`$, by an automorphism $`\phi `$ of $`H(F_n)`$, it can be reduced to the following form:
$$\phi (ff^t)=\left(\begin{array}{cc}I_l& 0\\ 0& 0\end{array}\right)$$
where $`I_l`$ is the identity matrix of order $`l`$. Since $`e`$ and $`e^t`$ are orthogonal to $`ff^t`$, they take the forms:
$$\phi (e)=\left(\begin{array}{cc}0& 0\\ 0& K\end{array}\right),\phi (e^t)=\left(\begin{array}{cc}0& 0\\ 0& K^t\end{array}\right),$$
where $`K`$ is a matrix of order $`nl`$. Therefore, according to (7),
$$\phi (S)=\left\{\left(\begin{array}{cc}X& 0\\ 0& Y\end{array}\right)\right\},\phi (S_1)=\left\{\left(\begin{array}{cc}X& 0\\ 0& 0\end{array}\right)\right\},$$
$$\phi (S_2)=\left\{\left(\begin{array}{cc}0& 0\\ 0& Y\end{array}\right)\right\},$$
$`(8)`$
where $`X`$ is any symmetric matrix of order $`l`$, $`Y`$ is a symmetric matrix of order $`2k=nl`$ which is also an element of a subalgebra of the type $`F_k^{(+)}`$.
In the case when $`M`$ is semisimple, that is, $`M=S=S_1S_2`$, there exist two orthogonal idempotents such that
$$S=eH(F_n)e+fH(F_n)f,e+f=1.$$
Acting in the same manner as above we can reduce $`S`$ to (8).
Therefore, we can define two homomorphisms $`\pi _1`$, $`\pi _2`$ as projections on $`S_1`$ and $`S_2`$, respectively. Since $`1𝒥`$, $`\pi _1(𝒥)\{0\}`$, $`\pi _2(𝒥)\{0\}`$. This implies that $`𝒥\pi _1(𝒥)H(F_l)`$, $`l<n`$, and $`𝒥\pi _2(𝒥)H(F_{2k})`$, $`2k<n`$. Therefore, we can reduce the problem of finding the canonical form of $`𝒥`$ to the case of all symmetric matrices of order less than $`n`$. However, the above reduction does not work in the case when $`S`$ is simple, that is, $`M=F_{\frac{r}{2}}^{(+)}R`$, $`rn`$. Hence we can conclude that as soon as the given simple subalgebra $`𝒥`$ is in the maximal subalgebra $`M`$ which has a non-simple semisimple factor $`S`$, the problem can be reduced to the case of symmetric matrices of a lower order. This process stops only if at some step either $`\pi _i(𝒥)MF_{\frac{r}{2}}^{(+)}R`$, or $`\pi _i(𝒥)`$ coincides with $`S_i`$. Without any loss of generality, we can assume that $`r=n`$, that is, $`𝒥MF_{\frac{n}{2}}^{(+)}R`$.
All we need to reach our goal is to determine the canonical form of $`𝒥`$ which is covered by a maximal subalgebra of the type $`F_{\frac{n}{2}}^{(+)}R`$. Notice that there is an isomorphic imbedding $`\theta `$ of $`F_{\frac{n}{2}}^{(+)}`$ into $`H(F_n)`$ such that $`\theta (A+iB)=\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)`$, where $`A`$ is a symmetric matrix of order $`\frac{n}{2}`$, $`B`$ is a skew-symmetric matrix of order $`\frac{n}{2}`$, $`i^2=1`$.
2.1. Let us assume that $`𝒥`$ has the type $`F_m^{(+)}`$ where $`n=2ml`$. We know that by an appropriate automorphism $`\psi `$ of $`F_{\frac{n}{2}}^{(+)}`$, we can reduce $`\theta ^1(𝒥)`$ to the following canonical form:
$$\psi (\theta ^1(𝒥))=\left\{\left(\begin{array}{cccccc}X& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& X& 0& \mathrm{}& 0\\ 0& \mathrm{}& 0& X^t& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& 0& \mathrm{}& X^t\end{array}\right)\right\},$$
where $`X`$ is any matrix of order $`m`$. Then, $`X`$ can be written as $`A+iB`$ for an appropriate symmetric $`A`$ and skew-symmetric $`B`$. Therefore, $`\theta (\psi (\theta ^1(𝒥)))`$ has the following representation in $`H(F_n)`$:
$$\theta (\psi (\theta ^1(𝒥)))=\left\{\left(\begin{array}{cccccccc}A& & & 0& B& & & 0\\ & A& & & & B& & \\ & & \mathrm{}& & & & \mathrm{}& \\ 0& & & A& 0& & & B\\ & & & & & & & \\ B& & & 0& A& & & 0\\ & B& & & & A& & \\ & & \mathrm{}& & & & \mathrm{}& \\ 0& & & B& 0& & & A\end{array}\right)\right\}$$
By Lemma 2.1, $`\theta \psi \theta ^1`$ (an automorphism of the algebra of the form (1)) can be extended to an automorphism of $`H(F_n)`$. Finally, by interchanging the $`k`$-th and $`(\frac{n}{2}+k)`$-th columns, and $`k`$-th and $`(\frac{n}{2}+k)`$-th rows, $`1k\frac{n}{2}`$, and the columns and rows inside the block (if necessary), we can achieve the following block-diagonal form:
$$\left\{\left(\begin{array}{ccccc}A& B& & & \\ B& A& & 0& \\ & & \mathrm{}& & \\ & 0& & A& B\\ & & & B& A\end{array}\right)\right\}$$
$`(9)`$
As a result any subalgebra of $`H(F_n)`$ of the type $`F_m^{(+)}`$ can be brought to the canonical form (2.1). This canonical form is obviously unique.
2.2. Let $`𝒥`$ be of the type $`H(F_m)`$. Acting in the same manner as before, $`𝒥`$ can be brought to the unique canonical form as follows
$$\theta (\psi (\theta ^1(𝒥)))=\left\{\left(\begin{array}{cccccc}X& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& X& 0& \mathrm{}& 0\\ 0& \mathrm{}& 0& X& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& 0& \mathrm{}& X\end{array}\right)\right\},$$
$`(10)`$
where $`X`$ is a symmetric matrix of order $`m`$.
2.3. Let $`𝒥`$ be of the type $`H(F_{2m},j)`$, $`n=4kl`$. Like in the previous cases, by an appropriate automorphism $`\psi `$ of $`F_{\frac{n}{2}}^{(+)}`$, $`\theta ^1(𝒥)`$ can be brought to the following block-diagonal form:
$$\psi (\theta ^1(𝒥))=\left\{\left(\begin{array}{cccccc}X& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& X& 0& \mathrm{}& 0\\ 0& \mathrm{}& 0& X& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& 0& \mathrm{}& X\end{array}\right)\right\},$$
where $`X`$ is a symplectic matrix of order $`2m`$. If we represent $`X`$ as the sum of symmetric and skew-symmetric matrices as follows:
$$X=\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)+\left(\begin{array}{cc}C& D\\ D& C\end{array}\right)$$
where all matrices have order $`m`$; $`A`$ is symmetric, $`B`$, $`C`$,$`D`$ are skew-symmetric, then $`\theta `$ induces the following representation of $`𝒥`$ in $`H(F_n)`$
$$\theta (\psi (\theta ^1(𝒥)))=\left\{\left(\begin{array}{cccccc}A& B& & C& D& \\ B& A& & D& C& \\ & & \mathrm{}& & & \mathrm{}\\ & & & & & \\ C& D& & A& B& \\ D& C& & B& A& \\ & & \mathrm{}& & & \mathrm{}\end{array}\right)\right\}$$
Similarly, by Lemma 2.1, $`\theta \psi \theta ^1`$ (an automorphism of the algebra of the form (1)) can be extended to an automorphism of $`H(F_n)`$.
By interchanging appropriate blocks, we can reduce it to the canonical form:
$$\left\{\left(\begin{array}{ccccccccc}A& B& C& D& & & & & \\ B& A& D& C& & & & & \\ C& D& A& B& & & 0& & \\ D& C& B& A& & & & & \\ & & & & \mathrm{}& & & & \\ & & & & A& B& C& D& \\ & & 0& & B& A& D& C& \\ & & & & C& D& A& B& \\ & & & & D& C& B& A& \end{array}\right)\right\}$$
$`(11)`$
From the explicit form (11), the canonical form of $`𝒥`$ of the type $`H(F_{2m},j)`$ is uniquely determined.
If $`1𝒥`$, then $`\text{rk}(e)=k<n`$ where $`e`$ is the identity element of $`𝒥`$, and by an appropriate automorphism of $`H(F_n)`$ $`𝒥`$ can be brought to the form:
$$\left\{\left(\begin{array}{cc}X& 0\\ & \\ 0& 0\end{array}\right)\right\}$$
$`(12)`$
where $`X`$ is a symmetric matrix of order $`k`$. Let $`\pi `$ denote a mapping that sends each matrix of the form (12) to the first block of order $`k`$. Clearly, $`𝒥\pi (𝒥)`$. Acting in the same manner as before we can bring $`\pi (𝒥)`$ to the canonical form in $`H(F_k)`$. As a result, the original subalgebra $`𝒥`$ also takes the unique canonical form.
Case 3 $`𝒜=H(F_{2n},j)`$
Since the proof of this case is not much different from the proof of the case of $`H(F_n)`$, we will omit most details. According to classification (see ), any maximal subalgebra $`M`$ in $`H(F_{2n},j)`$ is isomorphic to one of the following:
1. $`H(F_{2k},j)H(F_{2l},j)`$, $`k+l=n`$,
2. $`H(F_{2k},j)F_l^{(+)}R`$, $`k+l=n`$. If $`k=0`$, then $`M=F_n^{(+)}R`$.
3. $`J(f,1)`$ only if $`n=2^m`$ and either $`dimJ(f,1)=2(m+1)`$, $`m`$ is even, or $`dimJ(f,1)=2m+1`$, $`m`$ is odd.
First we assume that $`𝒥`$ is a simple matrix subalgebra of $`H(F_{2n},j)`$ such that $`1𝒥`$. Let $`M`$ be a maximal subalgebra which contains $`𝒥`$, $`𝒥M`$. By Lemma 2.3, $`𝒥S`$. If $`S`$ is a non-simple semisimple algebra, then $`𝒥`$ can be projected into the simple components of $`S`$. Hence, the problem will be reduced to the case of symplectic matrices of order less than $`2n`$. This reduction stops only when either the image of $`𝒥`$ can be covered by the maximal subalgebra with a simple Wedderburn factor $`S`$ or the image of $`𝒥`$ coincides with one of the simple components of $`S`$.
Next we look into the case when $`𝒥M`$, where $`M`$ has a simple Wedderburn factor $`S`$. There is no loss in generality if we assume that $`M=SR`$, $`SF_n^{(+)}`$. By Lemma 2.4, $`S`$ can be brought to the form (5). Notice that any automorphism of $`F_n^{(+)}`$ of the form $`\phi (X)=C^1XC`$ can be extended to an automorphism of $`H(F_{2n},j)`$ as follows:
$$\overline{\phi }(X)=\overline{C}^1X\overline{C},\overline{C}=\left(\begin{array}{cc}C& 0\\ 0& (C^1)^t\end{array}\right)$$
$`(13)`$
3.1 If $`𝒥H(F_m)`$, $`mn`$, then acting by some automorphism of the form (13), it can be reduced to (3.2). This canonical form is obviously uniquely determined.
3.2 If $`𝒥F_m^{(+)}`$, $`mn`$, then by some automorphism of the form (13) it can be brought to
$$\left\{\left(\begin{array}{cccccc}X& & & & & \\ & \mathrm{}& & & 0& \\ & & X^t& & & \\ & & & & & \\ & & & X^t& & \\ & 0& & & \mathrm{}& \\ & & & & & X\end{array}\right)\right\}$$
$`(14)`$
where $`X`$ is an arbitrary matrix of order $`m`$. This is the canonical form (3.1). For further purposes, we introduce the following characteristic of the canonical form. Let $`l`$ and $`k`$ denote the number of $`X`$-blocks and $`X^t`$-blocks in the top left block of order $`n`$ in (14). Set $`k_𝒜(𝒥)=|lk|`$.With some effort it can be shown that in this case the canonical form is also unique. All we have to show is that any two canonical forms $`𝒥_1`$ and $`𝒥_2`$ of the same type with $`k_𝒜(𝒥_1)k_𝒜(𝒥_2)`$ are not conjugate under symplectic automorphism, or, equivalently, automorphism of $`H(F_{2n},j)`$. For clarity, let $`𝒥_1=\text{diag}\{\underset{n}{\underset{}{X,X^t,\mathrm{},X}},X^t,X,\mathrm{}X^t\}`$ and $`𝒥_2=\text{diag}\{\underset{n}{\underset{}{Y,Y,\mathrm{}Y^t}},Y^t,Y^t,\mathrm{},Y\}`$ where $`X`$ and $`Y`$ are any matrices of order $`m`$. Let $`𝒮`$ stand for the subalgebra of $`H(F_{2n},j)`$ of the form (3). Obviously, $`𝒥_1S`$, $`𝒥_2S`$. Next we are going to show that for any automorphism $`\phi `$ of $`H(F_{2n},j)`$ such that $`\phi (𝒥_1)=𝒥_2`$ we can always find a symplectic automorphism $`\psi `$ that can be restricted to $`𝒮`$ and $`\psi (𝒥_1)=𝒥_2`$. Let $`C`$ be a non-singular matrix that determines $`\phi `$. Then, for any $`A𝒥_1`$ there exists $`B𝒥_2`$ such that
$$C^1AC=B,AC=CB.$$
$`(15)`$
Set $`C=(C_{ij})_{i,j=1,s}`$ where $`C_{ij}`$ is a square matrix of order $`m`$. By performing a matrix multiplication in (15) we obtain a series of equations:
$$XC_{ij}=C_{ij}Y,X^tC_{kl}=C_{kl}Y$$
where $`(i,j),(k,l)I\times I`$, $`I=\{1,\mathrm{},s\}`$. Since $`X`$ and $`Y`$ can be any matrices of order $`m`$, $`C_{ij}`$ can not be degenerate. Therefore, $`Y=C_{ij}^1XC_{ij}`$, $`Y=C_{kl}^1X^tC_{kl}`$. Hence, the matrix $`\overline{C}=\text{diag}\{\underset{n}{\underset{}{C_{ij},C_{kl},\mathrm{},C_{ij}}},(C_{ij}^t)^1,(C_{kl}^t)^1,\mathrm{},(C_{ij}^t)^1\}`$ determines an automorphism $`\psi `$ of $`H(F_{2n},j)`$ such that $`\psi (𝒥_1)=𝒥_2`$. Besides $`\psi `$ can be restricted to $`𝒮`$, therefore, induces an automorphism of a subalgebra of the type $`F_n^{(+)}`$. However we have already showed (case 1.1) that the two canonical forms in $`F_n^{(+)}`$ with $`k_𝒜(𝒥_1)k_𝒜(𝒥_2)`$ are not conjugate.
3.3 If $`𝒥H(F_{2m},j)`$, $`mn`$, then , it can be reduced to (3.3). Let $`l`$ and $`k`$ denote the number of $`X`$-blocks and $`X^t`$-blocks, respectively, in $`A`$. Set $`k_𝒜(𝒥)=|lk|`$.Then, we are going to show that any two canonical forms $`𝒥_1`$ and $`𝒥_2`$ of the same type with $`k_𝒜(𝒥_1)k_𝒜(𝒥_2)`$ are not conjugate under the automorphism of $`H(F_{2m},j)`$. Assume the contrary, that is, there exists an automorphism of $`H(F_{2m},j)`$ such that $`\phi (𝒥_1)=𝒥_2`$. Next we can choose $`𝒮_1𝒥_1`$, $`𝒮_1F_m^{(+)}`$ such that $`k_𝒜(𝒮_1)=k_𝒜(𝒥_1)`$. Similarly, we can select $`𝒮_2𝒥_2`$, $`𝒮_2F_m^{(+)}`$ such that $`k_𝒜(𝒮_2)=k_𝒜(𝒥_2)`$. By Lemma 2.2 there exists $`\psi `$, $`\psi :𝒥_2𝒥_2`$, $`\psi (\phi (𝒮_1))=𝒮_2`$. From the explicit from of $`𝒥_2`$, $`\psi `$ can be extended to an automorphism of $`H(F_{2m},j)`$, $`\psi \phi :H(F_{2m},j)H(F_{2m},j)`$. It follows that $`𝒮_1`$ and $`𝒮_2`$ have the same canonical forms, in particular, $`k_𝒜(𝒮_1)=k_𝒜(𝒮_2)`$, a contradiction.
If $`1𝒥`$, then in order to find the canonical form of $`𝒥`$ we use the same approach as in the case of $`H(F_n)`$.
The Theorem is proved.
Let $`e`$ denote the identity element of $`𝒥F_m^{(+)}`$, and $`\rho `$ stand for the natural representation of $`𝒥`$ in $`F^n`$, $`m<n`$. Obviously, $`\rho `$ induces the representation of $`S(𝒥)`$ in $`F^n`$. If $`S(𝒥)`$ is a non-simple semisimple associative algebra, that is, $`S(𝒥)=_1_2`$ where $`_1`$, $`_2`$ are isomorphic simple ideals, then $`\rho =\rho _1\rho _2`$ where $`\rho _i`$ is a representation of $`_i`$ in the corresponding invariant subspace of $`F^n`$. Then $`k_𝒜(𝒥)=|\text{deg}\rho _1(I_1)\text{deg}\rho _2(I_2)|`$.
###### Theorem 2.9.
Let $`𝒜`$ be a Jordan algebra of any of the following types: $`F_n^{(+)}`$, $`H(F_n)`$ or $`H(F_{2n},j)`$, $`n3`$, and $`𝒥`$, $`𝒥^{}`$ be proper simple matrix subalgebras of $`𝒜`$. If $`𝒥^{}`$ has the same type as $`𝒥`$ does, then $`𝒥^{}C(𝒥)`$ if and only if
(1) $`\text{rk}(e)=\text{rk}(e^{})`$;
(2) $`k_𝒜(𝒥)=k_𝒜(𝒥^{}),`$ in the case when $`𝒥F_m^{(+)}`$ or $`H(F_{2m},j)`$, for some $`m<n`$, and $`𝒜F_n^{(+)}`$ or $`H(F_{2n},j)`$.
###### Proof.
First it should be noted that the degree of $`𝒥3`$. The case of $`𝒥`$ of the degree 2 will be considered later in the text.
The case of $`F_n^{(+)}`$
In this case we assume that $`𝒥`$ and $`𝒥^{}`$ are subalgebras of $`F_n^{(+)}`$ which is as usual the set of all matrices of order $`n`$ closed under the Jordan multiplication. This case breaks into the following subcases.
(1) Let $`𝒥`$ be of the type $`F_m^{(+)}`$ for some $`m<n`$. First we assume that $`S(𝒥)`$ is a simple algebra. Equivalently, $`k_𝒜(𝒥)=\text{rk}(e)`$. Let $`𝒥^{}`$ be as given in the conditions of the Theorem. If $`𝒥^{}𝒞(𝒥)`$, then there exists an automorphism $`\phi `$ of $`F_n^{(+)}`$ which maps $`𝒥^{}`$ onto $`𝒥`$. It follows that $`\phi (e^{})=e,`$ therefore, $`\text{rk}(e^{})=\text{rk}(e)`$. Besides, $`\phi (S(𝒥^{}))=S(𝒥)`$. Hence, $`S(𝒥^{})`$ is also simple, $`k(𝒥^{})=\text{rk}(e^{})`$. It follows that $`k(𝒥^{})=\text{rk}(e^{})=\text{rk}(e)=k(𝒥)`$.
Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$ and $`k(𝒥)=k(𝒥^{})`$, then $`k(𝒥^{})=k(𝒥)=\text{rk}(e)=\text{rk}(e^{})`$, because $`k(𝒥)=\text{rk}(e)`$. Therefore, $`k(𝒥^{})=\text{rk}(e^{})`$, that is, $`S(𝒥^{})`$ is also simple, and $`𝒥`$, $`𝒥^{}`$ have the same canonical forms. This implies that $`𝒥^{}𝒞(𝒥)`$.
Now we assume that $`S(𝒥)`$ is a non-simple semisimple subalgebra. Let $`𝒥^{}`$ be another subalgebra which satisfies the conditions of the Theorem. If $`𝒥^{}𝒞(𝒥)`$, then there exists an automorphism $`\phi `$ of $`F_n^{(+)}`$ which maps $`𝒥^{}`$ onto $`𝒥`$. Therefore, $`𝒥^{}`$ and $`𝒥`$ have equivalent representations in $`F^n`$, and so do $`S(𝒥^{})`$ and $`S(𝒥)`$. Consequently, either $`\text{deg}\rho _1(_1)=\text{deg}\rho _1(_1^{})`$ and $`\text{deg}\rho _2(_2)=\text{deg}\rho _2(_2^{})`$ or $`\text{deg}\rho _1(_1)=\text{deg}\rho _2(_2^{})`$ and $`\text{deg}\rho _2(_2)=\text{deg}\rho _1(_1^{})`$. Equivalently, $`|\text{deg}\rho _1(_1)\text{deg}\rho _2(_2)|=|\text{deg}\rho _1(_1^{})\text{deg}\rho _2(_2^{})|`$, that is, $`k_𝒜(𝒥)=k_𝒜(𝒥^{})`$.
Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$ and $`k_𝒜(𝒥)=k_𝒜(𝒥^{})`$, then $`𝒥`$ and $`𝒥^{}`$ have the same canonical forms. Therefore, these subalgebras are conjugate by some automorphism of $`F_n^{(+)}`$, and $`𝒥^{}𝒞(𝒥)`$.
(2) Let $`𝒥`$ be of the type $`H(F_m)`$ for some $`mn`$.
Suppose that $`𝒥^{}`$ is another subalgebra of $`F_n^{(+)}`$ which has the type $`H(F_m)`$. If $`𝒥^{}`$ is conjugate to $`𝒥`$ under some automorphism $`\phi `$ of $`F_n^{(+)}`$ then $`\phi (e^{})=e`$ and $`\text{rk}(e^{})=\text{rk}(e)`$. In other words, the canonical form of $`𝒥^{}`$ is exactly the same as that of $`𝒥`$. Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$, then $`𝒥`$ and $`𝒥^{}`$ have the same canonical forms. Therefore, $`𝒥^{}𝒞(𝒥)`$.
(3) Let $`𝒥`$ be of the type $`H(F_{2m},j)`$ for some $`mn`$. The proof of this case is exactly the same as the previous proof.
The case of $`H(F_n)`$
Suppose that $`𝒥`$ and $`𝒥^{}`$ are two subalgebras of $`H(F_n)`$ that satisfy the conditions of the Theorem.
(1) Let $`𝒥`$ as well as $`𝒥^{}`$ be of the type $`F_m^{(+)}`$. Assume that $`𝒥^{}𝒞(𝒥)`$. It follows that there exists an automorphism of $`H(F_n)`$ such that $`\phi (𝒥^{})=𝒥`$. Hence, $`\text{rk}(e^{})=\text{rk}(e)`$. Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$, then $`𝒥`$ and $`𝒥^{}`$ have the same canonical form. Therefore, $`𝒥^{}𝒞(𝒥)`$.
(2) Now let both $`𝒥`$ and $`𝒥^{}`$ have the type $`H(F_k)`$ (or $`H(F_{2k},j)`$). If $`𝒥^{}𝒞(𝒥)`$, then there exists an automorphism $`\phi `$ of $`H(F_n)`$ that sends $`𝒥^{}`$ onto $`𝒥`$, $`\phi (𝒥^{})=𝒥`$. Consequently, $`\text{rk}(e^{})=\text{rk}(e)`$.
Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$, then they have the same canonical form. Therefore, $`𝒥^{}𝒞(𝒥)`$.
The case of $`H(F_{2n},j)`$
Suppose that $`𝒥`$ and $`𝒥^{}`$ are two subalgebras of $`H(F_{2n},j)`$ that satisfy the conditions of the Theorem.
(1) Let $`𝒥`$ as well as $`𝒥^{}`$ be of the type $`F_m^{(+)}`$, $`m<n`$. Assume that $`𝒥^{}𝒞(𝒥)`$. It follows that there exists an automorphism of $`H(F_{2n},j)`$ such that $`\phi (𝒥^{})=𝒥`$. Hence, $`\text{rk}(e^{})=\text{rk}(e)`$. Since $`𝒥^{}`$ and $`𝒥`$ are conjugate in $`H(F_{2n},j)`$, they have the same canonical forms in $`H(F_{2n},j)`$. Therefore, $`k_𝒜(𝒥)=k_𝒜(𝒥^{})`$.
Conversely, if all conditions hold true, then $`𝒥`$ and $`𝒥^{}`$ have the same canonical forms. Therefore, $`𝒥^{}𝒞(𝒥)`$.
(2) Now let both $`𝒥`$ and $`𝒥^{}`$ have the type $`H(F_m)`$, $`m<n`$. If $`𝒥^{}𝒞(𝒥)`$, then there exists an automorphism of $`H(F_{2n},j)`$ that sends $`𝒥^{}`$ onto $`𝒥`$, $`\phi (𝒥^{})=𝒥`$. Consequently, $`\text{rk}(e^{})=\text{rk}(e)`$.
Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$, then they have the same canonical forms. Therefore, $`𝒥^{}𝒞(𝒥)`$.
(3) Now let both $`𝒥`$ and $`𝒥^{}`$ have the type $`H(F_{2m},j)`$, $`m<n`$. If $`𝒥^{}𝒞(𝒥)`$, then there exists an automorphism of $`H(F_{2n},j)`$ that sends $`𝒥^{}`$ onto $`𝒥`$, $`\phi (𝒥^{})=𝒥`$. Consequently, $`\text{rk}(e^{})=\text{rk}(e)`$, $`k_𝒜(𝒥)=k_𝒜(𝒥^{})`$.
Conversely, if $`\text{rk}(e^{})=\text{rk}(e)`$ and $`k_𝒜(𝒥)=k_𝒜(𝒥^{})`$, then they have the same canonical forms. Therefore, $`𝒥^{}𝒞(𝒥)`$
The Theorem is proved.
###### Corollary 2.10.
If $`m`$ is any number such that $`mn`$, and $`n=mk+r`$, $`0r<m`$, then there exist subalgebras of $`F_n^{(+)}`$ of the type $`H(F_m)`$. Moreover, there are precisely $`k`$ conjugate classes corresponding to $`H(F_m)`$. If $`2mn`$, and $`n=2mk+r`$, $`0r<m`$ then $`F_n^{(+)}`$ has subalgebras of the type $`H(F_{2m},j)`$, and the number of conjugate classes corresponding to $`H(F_{2m},j)`$ is equal to $`k`$. Finally, if $`m<n`$, and $`n=mk+r`$, $`0r<m`$ then there exist subalgebras of $`F_n^{(+)}`$ of the type $`F_m^{(+)}`$, and, moreover, the number of conjugate classes is given by $`_{j=1}^k[\frac{j}{2}].`$
###### Corollary 2.11.
If $`m`$ is any number such that $`m<n`$, and $`n=mk+r`$, $`0r<m`$, then there exist subalgebras of $`H(F_n)`$ of the type $`H(F_m)`$. Moreover, there are precisely $`k`$ conjugate classes corresponding to $`H(F_m)`$. If $`2mn`$, and $`n=2mk+r`$, $`0r<m`$ then $`H(F_n)`$ has subalgebras of the type $`F_m^{(+)}`$, and the number of conjugate classes corresponding to $`F_m^{(+)}`$ is equal to $`k`$. Finally, if $`4mn`$, and $`n=4mk+r`$, $`0r<m`$ then there exist subalgebras of $`H(F_n)`$ of the type $`H(F_{2m},j)`$, and, moreover, the number of conjugate classes is $`k`$.
###### Corollary 2.12.
If $`m`$ is any number such that $`mn`$, and $`n=mk+r`$, $`0r<m`$, then there exist subalgebras of $`H(F_{2n},j)`$ of the type $`H(F_m)`$. Moreover, there are precisely $`k`$ conjugate classes corresponding to $`H(F_m)`$. If $`mn`$, and $`n=mk+r`$, $`0r<m`$ then $`H(F_{2n},j)`$ has subalgebras of the type $`F_m^{(+)}`$, and the number of conjugate classes corresponding to $`F_m^{(+)}`$ is equal to $`_{j=1}^k[\frac{j}{2}].`$ Finally, if $`m<n`$, and $`n=mk+r`$, $`0r<m`$ then there exist subalgebras of $`H(F_n,j)`$ of the type $`H(F_m,j)`$, and, moreover, the number of conjugate classes is $`_{j=1}^k[\frac{j}{2}].`$
### 2.2 Subalgebras of the type $`J(f,1)`$
First we recall a few facts concerning Clifford algebras over a field of characteristic not 2(see ). Let $`𝒥=FV`$ where $`V=\text{span}x_1,\mathrm{},x_{2m}`$, and $`f`$ a non-degenerate symmetric bilinear form on $`V`$. Then, $`C(V,f)`$ is a central simple associative algebra with a unique canonical involution ”—” such that it fixes elements from $`V`$. In this case the imbedding of $`𝒥`$ into $`C(V,f)^{(+)}`$ we will call canonical of the first type. Next, let $`𝒥=FV`$ where $`V=\text{span}x_1,\mathrm{},x_{2m+1}`$, and $`V_0=\text{span}x_1,\mathrm{},x_{2m}`$. Then, $`C(V,f)`$ is isomorphic to a tensor product of $`C(V_0,f)`$ and the two-dimensional center $`E`$ of $`C(V,f)`$. Moreover, $`E=F[z]`$ where $`z=x_1x_2\mathrm{}x_{2m+1}`$. In other words, $`C(V,f)=_1_2`$, $`_iC(V_0,f)`$. Note that $`FV𝒥/_iC(V,f)/_iC(V_0,f)`$. This imbedding of $`𝒥=FV`$ into $`C(V_0,f)^{(+)}`$ we will call canonical of the second type.
Let $`𝒜`$ be a simple matrix Jordan algebra, and $`𝒥`$ be a subalgebra of the type $`J(f,1)`$. According to , $`𝒥`$ of the type $`J(f,1)`$ is maximal in $`𝒜`$ if and only if one of the following cases hold
1. $`𝒜=(C(V_0,f),)`$, $`𝒥=FV`$ where $`dimV=2m+1`$ and $`m`$ is odd.
2. $`𝒜=H(C(V_0,f),)`$, $`𝒥=FV`$ where $`dimV=2m+1`$, $`m`$ is even.
3. $`𝒜=H(C(V,f),)`$, $`𝒥=FV`$ where $`dimV=2m`$
Next we recall that if $`dimV=2m`$, and $`m0,1(\text{mod}\mathrm{\hspace{0.17em}4})`$ then $`dimH(C(V,f),)=2^{m1}(2^m+1)`$. If $`dimV=2m`$ and $`m2,3(\text{mod}\mathrm{\hspace{0.17em}4})`$ then $`dimH(C(V,f),)=2^{m1}(2^m1)`$. If $`dimV=2m+1`$ and $`m0(\text{mod}\mathrm{\hspace{0.17em}4})`$ then $`dimH(C(V_0,f),)=2^{m1}(2^m+1)`$. If $`dimV=2m+1`$ and $`m2(\text{mod}\mathrm{\hspace{0.17em}4})`$ then $`dimH(C(V_0,f),)=2^{m1}(2^m1)`$.
Canonical realizations of $`J(f,1)`$
Let $`𝒜`$ be a simple matrix Jordan algebra, and $`𝒥=FV`$ is a subalgebra of $`𝒜`$.
1.1 $`𝒜=F_n^{(+)}`$, $`n=2^ml+r`$, $`dimV=2m`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is a matrix of order $`2^m`$, and if $`\pi _i`$ denotes the projection on the $`i`$th non-zero block, then $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the first type.
1.2 $`𝒜=F_n^{(+)}`$, $`n=2^ml+r`$, $`dimV=2m+1`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is a matrix of order $`2^m`$, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the second type.
1.3 $`𝒜=F_n^{(+)}`$, $`n=2^ml+r`$, $`dimV=2m+1`$,
$$𝒥=\{\text{diag}(\underset{s}{\underset{}{X,\mathrm{},X,}}\underset{k}{\underset{}{X^t,\mathrm{},X^t}}0,\mathrm{},0)\}$$
where $`s+k=l`$, $`X`$ is a matrix of order $`2^m`$, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the second type.
2.1 $`𝒜=H(F_n)`$, $`n=2^ml+r`$, $`dimV=2m`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is a symmetric matrix of order $`2^m`$, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the first type.
2.2 $`𝒜=H(F_n)`$, $`n=2^{m+1}l+r`$, $`dimV=2m`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is of the form (1) in which $`A`$ and $`B`$ are of order $`2^m`$. If $`𝒮`$ denotes the algebra of the form (1), then $`\pi _i(𝒥)𝒮`$ is a canonical imbedding of the first type.
2.3 $`𝒜=H(F_n)`$, $`n=2^{m+1}l+r`$, $`dimV=2m+1`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is of the form (1) in which $`A`$ and $`B`$ are of order $`2^m`$. If $`𝒮`$ denotes the entire algebra of the form (1), then $`\pi _i(𝒥)𝒮`$ is a canonical imbedding of the second type.
2.4 $`𝒜=H(F_n)`$, $`n=2^ml+r`$, $`dimV=2m+1`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},0,\mathrm{},0)\}$$
where $`X`$ is a symmetric matrix, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the second type.
3.1 $`𝒜=H(F_{2n},j)`$, $`n=2^ml+r`$, $`dimV=2m`$,
$$𝒥=\{\text{diag}(\underset{l}{\underset{}{X,\mathrm{},X}},\underset{k}{\underset{}{0,\mathrm{},0}},\underset{l}{\underset{}{X,\mathrm{},X}},\underset{k}{\underset{}{0,\mathrm{},0)}},\}$$
where $`k+l=n`$, $`X`$ is a symmetric matrix of order $`2^m`$, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the first type.
3.2 $`𝒜=H(F_{2n},j)`$, $`n=2^ml+r`$, $`dimV=2m`$, $`𝒥`$ has a canonical form (3.3), and if $`\pi _i`$ denotes the projection of $`𝒥`$ into $`i`$th simple component (of the type $`H(F_{2^m},j)`$) of (3.3), then $`\pi _i(𝒥)H(F_{2^m},j)`$ is a canonical imbedding of the first type.
3.3 $`𝒜=H(F_{2n},j)`$, $`n=2^ml+r`$, $`dimV=2m+1`$, $`𝒥`$ has a canonical form (3.3) where $`\pi _i(𝒥)H(F_{2^m},j)`$ is a canonical imbedding of the second type.
3.4 $`𝒜=H(F_{2n},j)`$, $`n=2^ml+r`$, $`dimV=2m+1`$,
$$𝒥=\{\text{diag}(\underset{s}{\underset{}{X,\mathrm{},X,}}\underset{k}{\underset{}{X^t,\mathrm{},X^t}}0,\mathrm{},0,\underset{s}{\underset{}{X^t,\mathrm{},X^t,}}\underset{k}{\underset{}{X,\mathrm{},X}}0,\mathrm{},0)\}$$
where $`s+k=l`$, $`X`$ is a matrix of order $`2^m`$, and $`\pi _i(𝒥)F_{2^m}^{(+)}`$ is a canonical imbedding of the second type.
###### Theorem 2.13.
Let $`𝒜`$ be a simple matrix Jordan algebra, and $`𝒥`$ be a subalgebra of $`𝒜`$ of the type $`J(f,1)`$. Then, $`𝒥`$ has a unique canonical form as above.
###### Proof.
Let $`𝒥=FV`$. Then the following cases occur.
Case $`𝒜=F_n^{(+)}`$
1.1 Let $`dimV=2m`$. Then $`U(𝒥)C(V,f)`$ is a simple algebra. In particular, $`S(𝒥)U(𝒥)`$.
If $`S(𝒥)=𝒜`$, then $`n=2^m`$, $`𝒜U(𝒥)`$. Therefore, the imbedding of $`𝒥`$ into $`𝒜`$ is equivalent to the imbedding of $`FV`$ into $`C(V,f)^{(+)}`$. Therefore, this is a canonical imbedding of the first type.
If $`S(𝒥)𝒜`$, then $`S(𝒥)`$ is a proper simple associative subalgebra of $`F_n`$. Therefore, $`S(𝒥)`$ can be reduced to
$$\{\text{diag}(\underset{l}{\underset{}{Y,\mathrm{},Y}},0,\mathrm{},0)\}$$
$`(16)`$
where the order of $`Y`$ is $`2^m`$, and $`n=2^ml+r`$. As a result, $`𝒥`$ also takes the canonical form 1.1.
1.2 Let $`dimV=2m+1`$. Then $`U(𝒥)C(V,f)`$, and $`U(𝒥)=_1_2`$, $`_iC(V_0,f)`$. Hence $`S(𝒥)`$ is isomorphic to either $`C(V,f)`$ or $`C(V_0,f)`$.
If $`S(𝒥)=𝒜`$, then the imbedding of $`𝒥`$ into $`𝒜`$ is the canonical of the second type.
If $`S(𝒥)_i`$, and $`S(𝒥)𝒜`$, then $`S(𝒥)`$ is a proper simple associative subalgebra of $`F_n`$. Therefore, $`S(𝒥)`$ can be reduced to (16). As a result, $`𝒥`$ takes the canonical form 1.2.
Finally, if $`S(𝒥)=_1_2`$, then $`𝒥`$ takes the canonical form 1.3
Case $`𝒜=H(F_n)`$
Let $`M`$ be the maximal subalgebra of $`H(F_n)`$ such that $`𝒥MH(F_n)`$. Then, the following cases occur.
1. $`M=SR`$ where $`S=S_1S_2`$ a semisimple factor, $`R`$ the radical. Then, we reduce the problem to the case of symmetric matrices of a lower dimension (see section 2.1).
2. $`M=S`$ where $`S=S_1S_2`$. Like in the previous case we can reduce the problem to the case of symmetric matrices of a lower dimension.
3. $`M=SR`$ where $`SF_{\frac{n}{2}}^{(+)}`$, $`R`$ the radical.
4. $`M=FW`$ where $`W`$ is a finite-dimensional vector space.
After a series of reductions of the form 1 and 2, the image of $`𝒥`$ becomes a subalgebra of
$$\left(\begin{array}{ccccc}𝒜_1& & & 0& \\ & \mathrm{}& & & \\ & & 𝒜_i& & \\ & & & \mathrm{}& \\ & 0& & & 𝒜_k\end{array}\right)$$
where $`𝒜_iH(F_{n_i})`$. Let $`\pi _i`$ be the projection of $`𝒥`$ into $`𝒜_i`$. To simplify our notations we denote $`\pi _i(𝒥)`$ as $`𝒥^{}`$, and the maximal subalgebra of $`𝒜_i`$ which covers $`𝒥^{}`$ as $`M_i`$, $`𝒥^{}M_i𝒜_i=H(F_{n_i})`$.
Case 1. Let $`dimV=2m`$ and $`m0,1(\text{mod}4)`$. Then we have the following cases:
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{n_i}`$, then $`n_i=2^m`$, $`F_{n_i}C(V,f)`$, and the imbedding of $`𝒥^{}`$ into $`F_{n_i}^{(+)}`$ is equivalent to the imbedding of $`FV`$ into $`C(V,f)^{(+)}`$, that is, canonical of the first type. If $`S(𝒥^{})F_{n_i}`$, then $`H(S(𝒥^{}))H(F_{n_i})`$ is a proper subalgebra of $`H(F_{n_i})`$. Hence, $`n_i=2^ml+r`$, and $`H(S(𝒥^{}))`$ can be reduced to (16) in which $`X`$ denotes a symmetric matrix of order $`2^m`$. Then, $`𝒥^{}`$ takes the canonical form 2.1.
(b) Let $`M_i=SR`$ where $`SF_{\frac{n_i}{2}}^{(+)}`$. By using $`\theta `$-isomorphism (see section 2.1) we obtain that $`\theta ^1(𝒥^{})F_{\frac{n_i}{2}}^{(+)}`$. If $`S(\theta ^1(𝒥^{}))=F_{\frac{n_i}{2}}^{(+)}`$, then $`n_i=2^{m+1}`$ and the imbedding of $`\theta ^1(𝒥^{})`$ into $`F_{\frac{n_i}{2}}^{(+)}`$ is the canonical imbedding of the first type. In particular, $`\theta ^1(𝒥^{})H(F_{\frac{n}{2}})`$. As a result $`𝒥^{}`$ takes the canonical form 2.1 in which $`l=2`$ and no zeros. If $`S(𝒥^{})F_{\frac{n_i}{2}}^{(+)}`$, then $`S(\theta ^1(𝒥^{}))`$ is a proper simple subalgebra of $`F_{\frac{n_i}{2}}`$, therefore, takes the form (16) and $`n_i=2^{m+1}l+r`$. Hence $`𝒥`$ takes the canonical form 2.1.
Case 2 Let $`dimV=2m`$, $`m2,3(\text{mod}4)`$.
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{n_i}`$, then $`n_i=2^m`$, $`F_{n_i}C(V,f)`$, $`𝒥^{}H(F_{n_i})`$. Hence we have the following commutative diagram:
$$\begin{array}{ccc}𝒥^{}=FV& \stackrel{id}{}& 𝒥^{}=FV\\ \sigma & & \eta \\ U(𝒥^{})& \stackrel{\phi }{}& F_{n_i}^{(+)}\end{array}$$
where $`\sigma =\phi \eta `$. Therefore, $`\sigma (𝒥^{})=\phi (\eta (𝒥^{}))`$ is symmetric with respect to the canonical involution ”—” which is symplectic in this particular case. On the other hand, $`\sigma (𝒥^{})`$ is also symmetric with respect to $`j^{}=\phi t\phi ^1`$. By the uniqueness of ”—”, $`j^{}`$ equals to ”—”. However it is not possible because $`dimH(C(V,f),)=\frac{2^m(2^m1)}{2}\frac{2^m(2^m+1)}{2}=dimH(C(V,f),j^{})`$. If $`S(𝒥^{})F_{n_i}`$, then $`H(S(𝒥^{}))H(F_{n_i})`$ is a proper subalgebra of $`H(F_{n_i})`$. Hence $`n_i=2^ml+r`$, and $`H(S(𝒥^{}))`$ can be reduced to (16) where $`X`$ denotes a symmetric matrix of order $`2^m`$. Let $`\pi _{ij}`$ denote the projection on $`j`$th non-zero block of (16). Then the imbedding $`\pi _{ij}(𝒥^{})\pi _{ij}(H(S(𝒥^{})))`$ is similar to the above imbedding, which is not possible.
(b) Let $`M_i=SR`$ where $`SF_{\frac{n_i}{2}}^{(+)}`$. Then $`\theta ^1(𝒥^{})F_{\frac{n_i}{2}}^{(+)}`$. Since $`S(\theta ^1(𝒥^{}))U(𝒥^{})`$, then $`n_i=2^{m+1}l+r`$, and $`S(\theta ^1(𝒥^{}))`$ can be reduced to (16) in which $`X`$ is any matrix of order $`2^m`$. Hence $`\pi _{ij}(\theta ^1(𝒥^{}))F_{2^m}^{(+)}`$ is a canonical imbedding of the first type, and $`𝒥`$ has the canonical form 2.2.
Case 3 Let $`dimV=2m+1`$ where $`m`$ is odd.
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{n_i}`$, then $`n_i=2^m`$, $`F_{n_i}C(V_0,f)`$. Therefore, the imbedding of $`𝒥^{}`$ into $`F_{n_i}^{(+)}`$ is equivalent to the imbedding of $`FV`$ into $`C(V_0,f)^{(+)}`$ which is canonical imbedding of the second type. Since $`m`$ is odd, $`𝒥^{}`$ is a maximal subalgebra in $`F_{n_i}^{(+)}`$. However, $`𝒥^{}H(F_{n_i})`$, hence, $`𝒥^{}`$ cannot be maximal. This case is not possible. If $`S(𝒥^{})F_{n_i}`$, then $`H(S(𝒥^{}))H(F_{n_i})`$ is a proper subalgebra of $`H(F_{n_i})`$, therefore, can be reduced to (16). However, the imbedding of $`\pi _{ij}(𝒥^{})`$ into $`F_{2^m}^{(+)}`$ is as shown above. Hence this case is also not possible.
(b)Let $`M_i=SR`$ where $`SF_{\frac{n_i}{2}}^{(+)}`$. Acting in the same manner as in case 2(b) we will come to the canonical form 2.3.
Case 4. Let $`dimV=2m+1`$ and $`m0(\text{mod}4)`$. Acting in the same manner as in previous cases we will reduce $`𝒥^{}`$ to the canonical form 2.4.
Case 5. Let $`dimV=2m+1`$ and $`m2(\text{mod}4)`$. Acting in the same manner as in previous cases we will reduce $`𝒥^{}`$ to the canonical form 2.3.
Case $`𝒜=H(F_{2n},j)`$
Let $`M`$ be the maximal subalgebra of $`H(F_{2n},j)`$ such that $`𝒥MH(F_{2n},j)`$. Then, the following cases occur.
1. $`M=SR`$ where $`S=S_1S_2`$ a semisimple factor, $`R`$ the radical. Then, we reduce the problem to the case of symplectic matrices of a lower dimension (see section 2.1).
2. $`M=S`$ where $`S=S_1S_2`$. Like in the previous case we can reduce the problem to the case of symplectic matrices of a lower dimension.
3. $`M=SR`$ where $`SF_n^{(+)}`$, $`R`$ the radical.
4. $`M=FW`$ where $`W`$ is a finite-dimensional vector space.
After a series of reductions of the form 1 and 2, the image of $`𝒥`$ becomes a subalgebra of the algebra in the canonical form (3.3) in which the $`i`$th component has order $`2n_i`$. Let $`\pi _i`$ denote the projection of $`𝒥`$ into the $`i`$th simple component of (3.3).
Case 1. Let $`dimV=2m`$ and $`m0,1(\text{mod}4)`$.
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{2n_i}`$, $`F_{2n_i}C(V,f)`$, $`2n_i=2^m`$. Acting in the same manner as in case 2(a), we can show that this situation is not possible. Likewise if $`S(𝒥^{})F_{2n_i}`$ then we can reduce this case to the case just considered. Therefore, it also never occurs.
(b) Let $`M_i=SR`$ where $`SF_{n_i}^{(+)}`$. Then $`𝒥^{}F_{n_i}^{(+)}`$, therefore, $`S(𝒥^{})`$ can be brought to (16), and $`\pi _{ij}(𝒥^{})F_{2^m}^{(+)}`$ is the canonical imbedding of the first type. Finally the original subalgebra takes the form 3.1
Case 2 Let $`dimV=2m`$, $`m2,3(\text{mod}4)`$.
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{2n_i}`$, then $`2n_i=2^m`$, $`F_{2n_i}C(V,f)`$, $`𝒥^{}F_{2n_i}^{(+)}`$ is the canonical imbedding of the first type. If $`S(𝒥^{})F_{2n_i}`$, then $`H(S(𝒥^{}),j)H(F_{2n_i},j)`$ is a proper subalgebra of $`H(F_{2n_i},j)`$, that is, $`n_i=2^ml+r`$, and $`H(S(𝒥^{}),j)`$ can be reduced to (3.3) in which each component has order $`2^m`$. Then, $`𝒥`$ takes the canonical form 3.2.
(b) Let $`M_i=SR`$ where $`SF_{n_i}^{(+)}`$. This case also lead us to the canonical form 3.2.
Case 3 Let $`dimV=2m+1`$ where $`m`$ is odd.
(a) Let $`M_i=FW`$. If $`S(𝒥^{})=F_{2n_i}`$, then $`2n_i=2^m`$, $`F_{2n_i}C(V_0,f)`$. Therefore, the imbedding of $`𝒥^{}`$ into $`F_{2n_i}^{(+)}`$ is equivalent to the imbedding of $`FV`$ into $`C(V_0,f)^{(+)}`$ which is canonical imbedding of the second type. Since $`m`$ is odd, $`𝒥^{}`$ is a maximal subalgebra in $`F_{2n_i}^{(+)}`$. However, $`𝒥^{}H(F_{2n_i},j)`$, hence, $`𝒥^{}`$ cannot be maximal. This case is not possible. If $`S(𝒥^{})F_{2n_i}`$, then $`H(S(𝒥^{}),j)H(F_{2n_i},j)`$ is a proper subalgebra of $`H(F_{2n_i},j)`$, therefore, can be reduced to (3.3). Let $`\pi _{ij}`$ denote the projection of $`𝒥^{}`$ into the $`j`$th simple component of (3.3). However, the imbedding of $`\pi _{ij}(𝒥^{})`$ into $`F_{2^m}^{(+)}`$ is as shown above. Hence this case is also not possible.
(b) Let $`M_i=SR`$ where $`SF_{n_i}^{(+)}`$. Then $`𝒥^{}F_{n_i}^{(+)}`$, therefore, $`S(𝒥^{})`$ can be brought to (16), and $`\pi _{ij}(𝒥^{})F_{2^m}^{(+)}`$ is the canonical imbedding of the second type. Finally the original subalgebra takes the form 3.4
Case 4. Let $`dimV=2m+1`$ and $`m0(\text{mod}4)`$. Acting in the same manner as in previous cases we will reduce $`𝒥^{}`$ to the canonical form 3.1.
Case 5. Let $`dimV=2m+1`$ and $`m2(\text{mod}4)`$. Acting in the same manner as in previous cases we will reduce $`𝒥^{}`$ to the canonical form 3.2.
###### Corollary 2.14.
Let $`𝒜`$ be a simple matrix Jordan algebra of degree $`3`$, and $`𝒥=FV`$ where either $`dimV=2m`$ or $`dimV=2m+1`$. Then, $`𝒥`$ is a subalgebra of $`𝒜`$ if and only if
1. $`2^mn`$,
2. $`2^{m+1}n`$, in the case when $`𝒜=H(F_n)`$ and $`m2,3(\text{mod}\mathrm{\hspace{0.17em}4})`$.
The author uses this opportunity to thank her supervisor Prof. Bahturin for his helpful cooperation, many useful ideas and suggestions.
|
warning/0506/hep-ph0506320.html
|
ar5iv
|
text
|
# ULB-TH/05-14. Leptogenesis and Dark Matter related ?
## 1 Introduction
According to the Concordance Model, ordinary matter in the form of baryons represents only $`\mathrm{\Omega }_B5\%`$ of the energy density of the Universe. The rest is apparently shared between Dark Matter and Dark Energy, with $`\mathrm{\Omega }_{DM}25\%`$ and $`\mathrm{\Omega }_{DE}70\%`$ respectively . Dark energy is supposed to be responsible for the accelerated expansion of the Universe but otherwise its true nature eludes us. The dark matter problem is almost mundane in comparison. We have a plethora of well-motivated and well-understood particle physics candidates, the most acclaimed currently being a neutralino, and we know of the existence of at least one component of dark matter, in the form of light neutrinos.
In the present paper we would like to address a nagging puzzle related to dark matter. This is the apparently coincidental fact that the energy density in baryons and that of dark matter are nearly the same
$$\mathrm{\Omega }_b/\mathrm{\Omega }_{dm}1/5.$$
(1)
This similitude is generally not addressed by scenarios predicting the existence of dark matter, nor a fortiori by those concerned with baryogenesis. Yet, although the ratio (1) is constant today this was not the case for all the history of the universe and, at least for conventional dark matter and baryon matter generation mechanisms, (1) is a puzzle.
By way of introduction, it is instructive to have a look at leptogenesis, the simplest mechanism which establishes a relation between dark matter (in the form of neutrinos) and the abundance of baryons. Leptogenesis fixes the ratio of baryon to cosmic background neutrino number densities (assuming the neutrino asymmetry itself is negligible) and requires the neutrinos to be light Majorana particles. It is well appreciated that neutrinos are too light to be the dominant form of dark matter but this is not our main concern here. More to the point is the fact that the constraints from leptogenesis on neutrinos masses are rather loose (the range $`0.001<m_\nu <\mathrm{\hspace{0.33em}0.1}`$ eV is claimed in but the range could be much broader, see ). Yet $`3<\mathrm{\Omega }_b/\mathrm{\Omega }_\nu <\mathrm{\hspace{0.33em}70}`$, where the lower bound comes from large scales structure formation ($`m_\nu 0.7eV`$) while the upper bounds comes from neutrino oscillations ($`m_\nu 0.03eV`$) . This is surprising, since leptogenesis has nothing to say about the baryon to neutrino mass ratio. Yet the ratio of baryon to neutrino energy densities are almost similar.
The above discussion illustrate a shortcoming of most attempts (including ours) to explain (1) i.e. that one has to understand both the particle number density ratio and the particle mass ratio. A most straightforward explanation could be that dark matter is made of antibaryons, albeit of course of an exotic, neutral and stable form, that could compensate the baryon number of ordinary matter. This is not in contradiction with nucleosynthesis or CMB fluctuations, since these observations constrain only the number of protons and neutrons (and their bound states). In such a scheme one would automatically get (1) of $`𝒪(1)`$ with the mass of dark and visible matter related to the scale of QCD. Of course we know too much about strong interactions and it seems difficult to make this idea consistent with observations. (There has been however and interesting recent attempt in this direction .)<sup>1</sup><sup>1</sup>1Yet another possibility would be to hide ordinary antibaryons into primordial black holes but this idea raises further issues, not the least being to find a mechanism responsible for the separation of matter and anti-matter. Also, primordial black hole have problems of their own (see for a recent discussion).
This lengthly introduction brings us to the much less ambitious path that will be ours. The main idea goes back to old works of Barr et al and Kaplan and more recent inputs of Kuzmin and Kitano and Low . This approach allows to fix the ratio of particle densities. The mass of dark matter particles then comes as a prediction to be tested.
## 2 Matter Genesis
The basic setup assumes that there is an asymmetry in the dark sector related to the baryon asymmetry of the universe. Both baryon matter and dark matter then owe their existence to a single mechanism, a sort of matter genesis.
The different existing scenarios (see ) differ in the implementation of this very idea, however there are some similarities in the conditions to be satisfied. Here we outline the version of that inspired us. By necessity, there is a dark sector, composed of a set of new particles. The visible sector, which consists of, among other things, baryons, and the dark sector communicate with each other but the interactions are suppressed at low energies. The lightest of these particles is protected from decay by some discrete symmetry, analogous to R-parity. This lightest particle cannot be produced thermally in the Universe. If it were, the tiny asymmetry in the dark sector would be drowned by numbers. This last condition motivates the introduction of a particle in the dark sector that we call the messenger particle. This particle is strongly interacting and in thermal equilibrium in the early universe. Because it is strongly interacting, it stays in thermal equilibrium even when it becomes non-relativistic and that messengers and their antiparticles begin to annihilate. The situation in the dark sector at this point is like that for ordinary baryons in the visible sector. Baryons and messengers both survive to annihilation thanks to a tiny asymmetry in their respective sector. In the visible sector, neutrons decay into protons and the chain ends. In the dark sector, the messengers decay into the lightest stable particle, that should better be electrically neutral.
There are presumably many possible concrete realization of this scenario. Ours differs from those pre-existing in the literature on the following points. First our prejudice will be that the mechanism responsible for matter genesis is leptogenesis. Then dark matter will then be made of light, $`m`$ few $`GeV`$, right-handed Majorana neutrinos. Last our model is based on an extension of the Standard Model (SM) which has been proposed for other purposes. The model is very constrained and, we agree, not the nicest model one would dream of. However we believe that there are some lessons to be drawn from it.
As we shall discuss, the main drawback of this model and its siblings, will be that, at the end of the day, it does not look very natural. Then, the mass of dark matter particles will come in as a constraint, not a prediction, but this was to be anticipated from the discussion in the introduction. Finally, the kind of dark matter of the type we consider would escape all attempts of detection. The messenger particle could be observed in high energy colliders, since it is a strongly interacting particle, similar to a (very very) heavy quark.
## 3 The Model
We have chosen to concentrate on a specific extension of the Standard Model that was proposed many years ago in as an alternative to the SM way of giving mass to the quarks and leptons and is known in the literature as the ”universal see-saw model”. The gauge group is $`SU(2)_L\times SU(2)_R\times U(1)_{BL}`$. The left and right-handed quarks $`Q_{R,L}`$ and leptons $`L_{R,L}`$ are respectively $`SU(2)_L`$ and $`SU(2)_R`$ doublets and, in the simplest framework, there are two Brout-Englert-Higgs (BEH) doublets,
$$\varphi _L(2,1,1)$$
and
$$\varphi _R(1,2,1).$$
To give mass to the quarks and leptons, one introduces a set of $`SU(2)`$ singlet Weyl fermions and a Majorana fermion $`N`$:
$$U(1,1,4/3)D(1,1,2/3)E(1,1,2)N(1,1,0).$$
Note the unusual $`BL`$ charge assignment of these fields. The BEH bosons, for instance, have a non-zero $`BL`$ charge, and there is a completely neutral field $`N`$. The latter will play the role of the heavy Majorana particle, analogous to the heavy right-handed Majorana neutrinos in standard leptogenesis scenarios.
This model looks nice but, unfortunately, we will need to complicate it a bit further. In particular we need to implement a discrete symmetry to protect the dark sector. We follow in that an old proposal of Babu et al . First we add two BEH scalars in the adjoint, whose purpose will become clear later on:
$$\mathrm{\Delta }_L(3,1,2)\mathrm{\Delta }_R(1,3,2).$$
Then we impose the following $`Z_4`$ symmetry
$`D_LD_LQ_RiQ_RL_RiL_R`$
$`\varphi _Ri\varphi _R\mathrm{\Delta }_R\mathrm{\Delta }_RN_RN_R,`$
all other fields transforming trivially under $`Z_4`$. The first effect of this symmetry is to forbid a Dirac mass term for the $`D`$ field and Yukawa couplings to the $`N`$ (would be neutrino Dirac mass terms). The allowed Yukawa couplings and mass terms then take the form
$`_y`$ $`=`$ $`h_d\overline{Q}_L\varphi _LD_R+h_u\overline{Q}_L\stackrel{~}{\varphi }_LU_R+h_e\overline{L}_L\varphi _LE_R`$ (2)
$`+`$ $`\lambda L_L^TC^1\tau _2\stackrel{}{\tau }\stackrel{}{\mathrm{\Delta }}_LL_L`$ (3)
$`+`$ $`M_U\overline{U}_LU_R+M_E\overline{E}_LE_R+M_N\overline{N^c}N`$ (4)
$`+`$ $`(LR)+h.c..`$ (5)
This seems utterly complicated but the interesting things come with symmetry breaking. Let us write $`v_{L,R}`$ the vev of $`\varphi _{L,R}`$ and $`\kappa _{L,R}`$ the vev of the triplets. Then the neutrino fields are all pure Majorana
$$\lambda \kappa _L\overline{\nu _L^c}\nu _L+\lambda ^{}\kappa _R\overline{\nu _R^c}\nu _R+M_N\overline{N^c}N.$$
The up-like quarks and charged leptons get their mass from mixing with the heavy Dirac singlets
$$(\overline{f}\overline{F})\left(\begin{array}{cc}0& hv_L\\ hv_R& M\end{array}\right)\left(\begin{array}{c}f\\ F\end{array}\right)\frac{h^2v_lv_R}{M}\overline{f}f+M\overline{F}F,$$
where $`f=e,u`$ and $`F=U,E`$ thus following the usual ”universal see-saw” pattern.
The twist is in the down-like quark sector. Because there is no Dirac mass term for the $`D`$ field, mixing is maximal
$$h_dv_L\overline{d_L}D_R+h_dv_R\overline{D}_Ld_R+h.c.=h_dv_L\overline{d}^{}d^{}+h_dv_R\overline{D}^{}D^{},$$
and the role of the ”light” and ”heavy” right-handed down-like fields are so to speak exchanged. The $`D^{}`$ particle, which couples to $`SU(2)_R`$ gauge bosons, will be our strongly interacting messenger particle. It is supposed to be lighter than the singlet fermions.
The $`\nu _R`$ will get their mass from the vev of the $`SU(2)_R`$ adjoint scalar field. In the sequel, we assume that $`m_{\nu _R}m_D^{}M_N`$. (The mass of the $`U`$ and $`E`$ are not very much constrained. We will only request that the $`E,U`$ disappear before the electroweak phase transition.)
Finally, after left-right symmetry breaking, there is a residual $`Z_2`$ symmetry. The heavy Majorana field $`N`$, the heavy down-like quark $`D^{}`$, the Majorana neutrino $`\nu _R`$ as well as the charged boson fields $`W_R^\pm `$, $`\varphi _R^\pm `$ and $`\mathrm{\Delta }_R^\pm `$ are all odd under $`Z_2`$. All together, they constitute the dark sector of our model.
### 3.1 Initial B-L asymmetry
We will assume that the initial $`BL`$ asymmetry is provided by the out-of-equilibrium, CP violating decay of the heavy singlet Majorana fields $`N`$. For definiteness, we assume that decay takes place after left-right symmetry breaking. The abundance of $`N`$’s could be thermal or they could be created during reheating after inflation. Note that these fields are odd under the $`Z_2`$ symmetry and are thus the grandfather of our dark matter particles. The decay process is supposed to be dictated by higher scale interactions but we can parameterize it by dimension six effective operators like
$$\frac{1}{\mathrm{\Lambda }^2}\overline{N}E\overline{D}U+h.c.,$$
where the $`D`$ particle is the mass eigenstate, odd under the $`Z_2`$ symmetry (since there should be no confusion at this point, we drop the prime on the $`D`$). Assuming CP violation, these decay processes may sequestrate a $`BL`$ asymmetry between the dark and visible sectors
$$n_{BL}^{vis}=n_{BL}^{dark}=q_{BL}^D(n_Dn_{\overline{D}}),$$
where
$$n_Dn_{\overline{D}}=n_{\overline{U}}n_U=n_{\overline{E}}n_E=ϵn_N,$$
with
$$ϵ=(\mathrm{\Gamma }_{N\overline{E}\overline{U}D}\mathrm{\Gamma }_{NE\overline{D}U})/\mathrm{\Gamma }_N.$$
### 3.2 Annihilation of messenger particles
After sequestration of a $`BL`$ asymmetry in the dark sector, the Universe contains $`U`$, $`E`$ and $`D`$ particles on top of the usual Standard Model fermions.
In the visible sector, the $`E`$ and $`U`$ are in thermal and chemical equilibrium with the Standard Model fermions, and all together they carry a $`Z_2`$-even $`BL`$ asymmetry. Eventually, we will require the $`E`$ and $`U`$ disappears through annihilation and decay before the electroweak phase transition, leaving only SM degrees of freedom behind. As in standard leptogenesis scenarios, baryon number violating processes that are in equilibrium give birth to a non-zero baryon asymmetry
$$n_B=Cn_{BL}^{vis}=Cq_{BL}^D(n_Dn_{\overline{D}}).$$
(6)
The constant of proportionality $`C=25/79`$ is calculated in the standard way , taking into account that the $`BL`$ charge is shared between the visible and the dark sector.
In the dark sector, the messenger particles $`D`$ carry a $`Z_2`$-odd $`BL`$ asymmetry. They are heavy, $`M_Dv_R`$, strongly interacting particles and when the temperature of the universe drops below their mass, they annihilate into light quarks but a small asymmetry survives
$$n_Dn_{\overline{D}}n_Dϵn_N.$$
It is crucial that we require that there are essentially no $`\nu _R`$ in the universe at this level since we want to obtain a relation between the baryon asymmetry and the density of dark matter. As we will see in section 3.4, this condition constrains the scale of left-right symmetry breaking.
It is also crucial that the messenger particles are strongly interacting so as to leave only the asymmetry as a remnant.
### 3.3 Decay of messengers into $`\nu _R`$
The dominant $`D`$ decay channel is
$$Du+e+\nu _R^c,$$
through the exchange of a $`W_R`$. If the messenger particles were to decay before the electroweak phase transition, baryon number violating processes in equilibrium would completely erase the asymmetry (6). Indeed the $`\nu _R`$ carry no $`BL`$ charge in our framework and all the $`BL`$ that was sequestrated in the dark sector is released in the $`u`$ and $`e`$ degrees of freedom.
If $`D`$ decay takes place after electroweak symmetry breaking, the final $`B`$ asymmetry is given by (6) plus the contribution from the $`D`$ decay into baryons
$$n_B^{fin}=\left(q_B^u\frac{25}{79}q_{BL}^D\right)n_D.$$
(7)
The density of dark matter is simply equal to
$$n_{dm}=n_{\nu _R}=n_D.$$
Taking the ratio we obtain
$$\frac{\mathrm{\Omega }_B}{\mathrm{\Omega }_{DM}}=\left(q_B^u\frac{25}{79}q_{BL}\right)\frac{m_b}{m_{\nu _R}}0.5\frac{m_b}{m_{\nu _R}},$$
which implies that $`m_{\nu _R}3GeV`$. As expected, the mass of the dark matter particle is of order of the proton mass.
This scenario, the main features of which are summarized in Figure 1, is quite involved. The main element is that a $`BL`$ asymmetry is sequestrated in a sector insensitive to $`B+L`$ violating processes, at least as long as they are active, and is eventually released. In the present model, this is possible thanks to an exact discrete symmetry which differentiate the dark and the visible sector.
### 3.4 Summary of constraints
There are several constraints to put on scales and couplings for the above scenario to work. They are summarized in the present section.
First, the messenger particles $`D`$ have to decay after EW symmetry breaking to protect the baryon number from erasure. Moreover, since the $`D`$ decay products contributes to the baryon number, $`D`$ decay should take place before nucleosynthesis. From this we get
$$h_d^5v_R>\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}10}^{21}TeV,$$
(8)
where $`h_d`$ is the $`D`$ Yukawa coupling. This is quite a nasty constraint, since it require a rather small Yukawa coupling to be satisfied.
Second, in order to get a ratio of baryon and dark matter number density of $`𝒪(1)`$, we require $`D`$ to decay after the completion of $`D\overline{D}`$ annihilation. This implies, using for the temperature of annihilation interactions freeze-out $`T_A^{fo}=M_D/x_f`$ ($`x_f=𝒪(20)`$ see ) :
$$h_d^3v_R>g_{}^{1/2}\times 10^{15}TeV.$$
(9)
Third, the $`D`$ asymmetry produced in $`N`$ decay must be larger than the $`D\overline{D}`$ relic from freeze-out. Since $`D`$ annihilates through strong interactions, using the same arguments than , we obtain :
$$h_dv_R<\left(\frac{3GeV}{m_{\nu _R}}\right)\times 10^4TeV.$$
(10)
Finally, the abundance $`\nu _R`$ produced after reheating at $`T_{RH}`$ must be negligible compared to the abundance from $`D`$ decay. Assuming the $`\nu _R`$ are produced essentially through $`SU(2)_R`$ gauge bosons and taking $`T_{RH}>M_D`$, we get
$$h_d^3v_R>\mathrm{\hspace{0.33em}10}^{23}TeV.$$
(11)
All together, these constraints yield a parameters space reduced to
$$10^7TeV<v_R<\mathrm{\hspace{0.33em}\hspace{0.33em}10}^{11}TeV\text{and}10^7<h_d<\mathrm{\hspace{0.33em}\hspace{0.33em}10}^5.$$
(12)
This region is showed in Figure 2. There is a small but non-vanishing region where all the constraints can be met. In particular, the messenger $`D`$ particles are rather light, with a mass $`𝒪(TeV)`$, compared to the scale of left-right symmetry breaking. This result is consistent with the results of Kitano and Low .
## 4 Observational implications ?
Our dark matter candidate is, by construction, rather light $`m_{\nu _R}GeV`$ and abundant. Its cross-section is, by necessity, very small. This is essentially because our right-handed neutrinos must be non-thermal relics, with nearly the same number density as baryons. We had to pay a heavy price to achieve this result. First, the discrete symmetry of our model is not particularly natural. Second, the Yukawa coupling of the messenger particle is quite small. Last, the mass of the dark matter candidate is fixed by hand.
On the observational side, we expect our right-handed neutrinos to be present in the core of the Galaxy where they could annihilate with each other producing a heavy $`Z_R`$ boson, or be co-annihilated with right-handed quarks or leptons. Unfortunately the cross-section is way too small, $`\sigma v<\mathrm{\hspace{0.33em}10}^{32}pb`$, to give any observable signal.<sup>2</sup><sup>2</sup>2By way of comparison, the cross-section needed to reach the sensitivity of INTEGRAL signals would be $`𝒪(10100pb)`$ for a dark matter candidate with mass of $`𝒪(GeV)`$ (see for more details about the INTEGRAL signal and it’s correlation with light dark matter annihilation). We expect this conclusion to be generic for dark matter candidates related to the baryon asymmetry of the Universe, although we have no general proof.
Our dark matter candidate and the messenger have otherwise similar characteristics as in the model discussed in . In particular, baring other explanations, light right-handed neutrinos might be of interest to explain the apparent suppression in the power spectrum on small scales, having a free-steaming length $`0.1`$ Mpc.
The only hope to detect something in our model is by the production at a collider of the strongly interacting messenger particle, analog to a very heavy quark. Our messenger has a mass range between $`1TeV`$ and $`10^6TeV`$, corresponding to a life time between $`10^2s`$ and $`10^{10}s`$. As already underlined in , at least at the very lower part of this mass range, such a particle could be produced at the LHC.
## 5 Conclusion
We have discussed a mechanism of matter genesis, based on a left-right symmetric extension of the Standard Model, the basic idea being that both a baryonic and dark matter asymmetry have to be generated at some stage in the history of the Universe. Our dark matter candidate is a stable right-handed neutrino with mass $`\mathrm{\hspace{0.17em}3}GeV`$. The idea, which has been proposed by several authors, is quite attractive. However we found it quite difficult to realize. Although one should perhaps not try to draw a general conclusion from our model, the introduction of realistic gauge and Yukawa couplings shows that such a scenario is doomed to be very constrained.
This being said, the main drawback of the whole approach is still that such a candidate dark matter is essentially undetectable. On the theoretical side, we should also pause and ask what has been gained. We have a very contrived model, with a discrete symmetry, many new degrees of freedom and new interactions and yet all we can do is to relate the baryon and dark matter particle densities. The mass of the dark matter particle has still be fixed by hand.
By way of conclusion we would like to mention a recent attempt which could confront this difficulty. This mechanism could arise in the context of scalar-tensor theories of gravity coupled to matter. Since the mass of matter fields depends generically on the vev of a scalar field, the presence of matter induces an effective potential. For concreteness, suppose that the coupling of $`\phi `$ to matter is such that
$$V(\phi )=m_be^{\alpha \phi }n_b+m_{dm}e^{\beta \phi }n_{dm},$$
with $`\alpha ,\beta >0`$. Then
$$\mathrm{\Omega }_b/\mathrm{\Omega }_{dm}=\beta /\alpha $$
(13)
at the minimum of the potential (which depends on the density of ordinary and dark matter). If the couplings are of the same order, one gets a dynamical relaxation of the ratio (1). This idea is all nice and well, but again poses problems of its own. Baryons masses are varying, there is an extremely light scalar field with gravitational coupling, etc. The authors in have proposed to add an extra potential term to cure these issues ($`\phi `$ then behaves as a chameleon, changing mass in function of its environment) but the potential needs some fine tuning so as not to ruin (13). This model is thus not very satisfying but the idea is seductive. At any rate, explaining the apparent coincidence of (1) is a challenge worth pursuing.
## Acknowledgments
We thank Fu-Sin Ling, Emmanuel Nezri and Jean-Marie Frère for helpful discussions. This work is supported in part by IISN, la Communauté Française de Belgique (ARC), and the belgian federal government (IUAP-V/27).
|
warning/0506/cond-mat0506082.html
|
ar5iv
|
text
|
# Superconductivity in C6Ca explained.
## Abstract
Using density functional theory we demonstrate that superconductivity in C<sub>6</sub>Ca is due to a phonon-mediated mechanism with electron-phonon coupling $`\lambda =0.83`$ and phonon-frequency logarithmic-average $`\omega =24.7`$ meV. The calculated isotope exponents are $`\alpha (\mathrm{Ca})=0.24`$ and $`\alpha (\mathrm{C})=0.26`$. Superconductivity is mostly due C vibrations perpendicular and Ca vibrations parallel to the graphite layers. Since the electron-phonon couplings of these modes are activated by the presence of an intercalant Fermi surface, the occurrence of superconductivity in graphite intercalated compounds requires a non complete ionization of the intercalant.
Graphite intercalated compounds (GICs) were first synthesized in 1861 Schaffautl but only from the 30s a systematic study of these systems began. Nowadays a large number of reagents can be intercalated in graphite ($`100`$)DresselhausRev . Intercalation allows to change continuously the properties of the pristine graphite system, as it is the case for the electrical conductivity. The low conductivity of graphite can be enhanced to obtain even larger conductivities than Copper Foley . Moreover at low temperatures, intercalation can stabilize a superconducting stateDresselhausRev . The discovery of superconductivity in other intercalated structures like MgB<sub>2</sub>Nagamatsu and in other forms of doped Carbon (diamond) Ekimov has renewed interest in the field.
The first discovered GIC superconductors were alkali-intercalated compoundsHannay (C<sub>8</sub>A with A= K, Rb, Cs with T$`{}_{c}{}^{}<`$ 1 K). Synthesis under pressure has been used to obtain metastable GICs with larger concentration of alkali metals (C<sub>6</sub>K, C<sub>3</sub>K, C<sub>4</sub>Na, C<sub>2</sub>Na) where the highest T<sub>c</sub> corresponds to the largest metal concentration, T<sub>c</sub>(C<sub>2</sub>Na)=5 K Belash . Intercalation involving several stages have also been shown to be superconductingAlexander ; Outti (the highest T<sub>c</sub> = 2.7 K in this class belongs to KTl<sub>1.5</sub>C<sub>4</sub>). Intercalation with rare-earths has been tried, C<sub>6</sub>Eu, C<sub>6</sub>Cm and C<sub>6</sub>Tm are not superconductors, while recently it has been shown that C<sub>6</sub>Yb has a T<sub>c</sub> = 6.5 K Weller . Most surprising superconductivity on a non-bulk sample of C<sub>6</sub>Ca was also discoveredWeller . The report was confirmed by measurements on bulk C<sub>6</sub>Ca poly-crystalsGenevieve and a $`T_c=11.5`$ K was clearly identified. At the moment C<sub>6</sub>Yb and C<sub>6</sub>Ca are the GICs with the highest T<sub>c</sub>. It is worthwhile to remember that elemental Yb and Ca are not superconductors.
Many open questions remain concerning the origin of superconductivity in GICs. (i) All the aforementioned intercalants act as donors respect to graphite but there is no clear trend between the number of carriers transferred to the Graphene layers and T<sub>c</sub>DresselhausRev . What determines T<sub>c</sub>? (ii) Is superconductivity due to the electron-phonon interaction Mazin or to electron correlation Csanyi ? (iii) In the case of a phonon mediated pairing which are the relevant phonon modes Mazin ? (iv) How does the presence of electronic donor states (or interlayer states) affect superconductivity DresselhausRev ; Csanyi ; Mazin ?
Two different theoretical explanations has been proposed for superconductivity in C<sub>6</sub>Ca. In Csanyi it was noted that in most superconducting GICs an interlayer state is present at E<sub>f</sub> and a non-conventional excitonic pairing mechanismAllender has been proposed. On the contrary Mazin Mazin suggested an ordinary electron-phonon pairing mechanism involving mainly the Ca modes with a 0.4 isotope exponent for Ca and 0.1 or less for C. However this conclusion is not based on calculations of the phonon dispersion and of the electron-phonon coupling in C<sub>6</sub>Ca. Unfortunately isotope measurements supporting or discarding these two thesis are not yet available.
In this work we identify unambiguously the mechanism responsible for superconductivity in C<sub>6</sub>Ca. Moreover we calculate the phonon dispersion and the electron-phonon coupling. We predict the values of the isotope effect exponent $`\alpha `$ for both species.
We first show that the doping of a graphene layer and an electron-phonon mechanism cannot explain the observed T<sub>c</sub> in superconducting GICs. We assume that doping acts as a rigid shift of the graphene Fermi level. Since the Fermi surface is composed by $`\pi `$ electrons, which are antisymmetric respect to the graphene layer, the out-of-plane phonons do not contribute to the electron-phonon coupling $`\lambda `$. At weak doping $`\lambda `$ due to in-plane phonons can be computed using the results of ref. Piscanec . The band dispersion can be linearized close to the K point of the hexagonal structure, and the density of state per two-atom graphene unit-cell is $`N(0)=\beta ^1\sqrt{8\pi \sqrt{3}}\sqrt{\mathrm{\Delta }}`$ with $`\beta =14.1`$ eV and $`\mathrm{\Delta }`$ is the number of electron donated per unit cell (doping). Only the E<sub>2g</sub> modes near $`\mathrm{\Gamma }`$ and the A$`{}_{1}{}^{}{}_{}{}^{}`$ mode near K contribute:
$$\lambda =N(0)\left[\frac{2g_𝚪^2_F}{\mathrm{}\omega _𝚪}+\frac{1}{4}\frac{2g_𝐊^2_F}{\mathrm{}\omega _𝐊}\right]=0.34\sqrt{\mathrm{\Delta }}$$
(1)
where the notation is that of ref. Piscanec . Using this equation and typical values of $`\mathrm{\Delta }`$ Pietronero the predicted T<sub>c</sub> are order of magnitudes smaller than those observed. As a consequence superconductivity in C<sub>6</sub>Ca and in GICs cannot be simply interpreted as doping of a graphene layer, but it is necessary to consider the GIC’s full structure.
The atomic structureGenevieve of CaC<sub>6</sub> involves a stacked arrangement of graphene sheets (stacking AAA) with Ca atoms occupying interlayer sites above the centers of the hexagons (stacking $`\alpha \beta \gamma `$). The crystallographic structure is R3̄m Genevieve where the Ca atoms occupy the 1a Wyckoff position (0,0,0) and the C atoms the 6g positions (x,-x,1/2) with x$`=1/6`$. The rombohedral elementary unit cell has 7 atoms, lattice parameter 5.17 $`\mathrm{\AA }`$ and rombohedral angle $`49.55^o`$. The lattice formed by Ca atoms in C<sub>6</sub>Ca can be seen as a deformation of that of bulk Ca metal. Indeed the fcc lattice of the pure Ca can be described as a rombohedral lattice with lattice parameter 3.95 $`\mathrm{\AA }`$ and angle $`60^o`$. Note that the C<sub>6</sub>Ca crystal structure is not equivalent to that reported in Weller which has a stacking $`\alpha \beta `$. In Weller the structure determination was probably affected by the non-bulk character of the samples.
Density Functional Theory (DFT) calculations are performed using the PWSCF/espresso codePWSCF and the generalized gradient approximation (GGA) PBE . We use ultrasoft pseudopotentialsVanderbilt with valence configurations 3s<sup>2</sup>3p<sup>6</sup>4s<sup>2</sup> for Ca and 2s<sup>2</sup>2p<sup>2</sup> for C. The electronic wavefunctions and the charge density are expanded using a 30 and a 300 Ryd cutoff. The dynamical matrices and the electron-phonon coupling are calculated using Density Functional Perturbation Theory in the linear responsePWSCF . For the electronic integration in the phonon calculation we use a $`N_k=6\times 6\times 6`$ uniform k-point meshfootnotemesh and and Hermite-Gaussian smearing of 0.1 Ryd. For the calculation of the electron-phonon coupling and of the electronic density of states (DOS) we use a finer $`N_k=20\times 20\times 20`$ mesh. For the $`\lambda `$ average over the phonon momentum q we use a $`N_q=4^3`$ $`𝐪`$points mesh. The phonon dispersion is obtained by Fourier interpolation of the dynamical matrices computed on the $`N_q`$ points mesh.
The DFT band structure is shown in figure 1(b). Note that the $`\mathrm{\Gamma }\chi `$X direction and the L$`\mathrm{\Gamma }`$ direction are parallel and perpendicular to the graphene layers. The K special point of the graphite lattice is refolded at $`\mathrm{\Gamma }`$ in this structure. For comparison we plot in 1(c) the band structure of C<sub>6</sub>Ca and with Ca atoms removed (C<sub>6</sub>) and the structure C<sub>6</sub>Ca with C<sub>6</sub> atoms removed (Ca). The size of the red dots in fig. 1(b) represents the percentage of Ca component in a given band (Löwdin population). The Ca band has a free electron like dispersion as in fcc Ca. From the magnitude of the Ca component and from the comparison between fig. 1(b) and (c) we conclude that the C<sub>6</sub>Ca bands can be interpreted as a superposition of the Ca and of the C<sub>6</sub> bands. At the Fermi level, one band originates from the free electron like Ca band and disperses in all the directions. The other bands correspond to the $`\pi `$ bands in C<sub>6</sub> and are weakly dispersive in the direction perpendicular to the graphene layers. The Ca band has been incorrectly interpreted as an interlayer-band Csanyi not associated to metal orbitals.
More insight on the electronic states at E<sub>f</sub> can be obtained calculating the electronic DOS. The total DOS, fig. 1(a), is in agreement with the one of ref. Mazin and at E<sub>f</sub> it is $`N(0)=1.50`$ states/(eV unit cell). We also report in fig. 1(a) the atomic-projected density of state using the Löwdin populations, $`\rho _\eta (ϵ)=\frac{1}{N_k}_{𝐤n}|\varphi _\eta ^L|\psi _{𝐤n}|^2\delta (ϵ_{𝐤n}ϵ)`$. In this expression $`|\varphi _\eta ^L=_\eta [𝐒^{1/2}]_{\eta ,\eta ^{}}|\varphi _\eta ^{}^a`$ are the orthonormalized Löwdin orbitals, $`|\varphi _\eta ^{}^a`$ are the atomic wavefunctions and $`S_{\eta ,\eta ^{}}=\varphi _\eta ^a|\varphi _\eta ^{}^a`$. The Kohn and Sham energy bands and wavefunctions are $`ϵ_{𝐤n}`$ and $`|\psi _{𝐤n}`$. This definition leads to projected DOS which are unambiguously determined and are independent of the method used for the electronic structure calculation. At E<sub>f</sub> the Ca 4s, Ca 3d, Ca 4p, C 2s, C 2p<sub>σ</sub> and C 2p<sub>π</sub> are 0.124, 0.368, 0.086, 0.019, 0.003, 0.860 states/(cell eV), respectively. Most of C DOS at E<sub>f</sub> comes from C 2p<sub>π</sub> orbitals. Since the sum of all the projected DOSs is almost identical to the total DOS, the electronic states at E<sub>f</sub> are very well described by a superposition of atomic orbitals. Thus the occurrence of a non-atomic interlayer-state, proposed in ref. Csanyi , is further excluded. From the integral of the projected DOSs we obtain a charge transfer of 0.32 electrons (per unit cell) to the Graphite layers ($`\mathrm{\Delta }=0.11`$).
The phonon dispersion ($`\omega _{𝐪\nu }`$) is shown in fig. 2. For a given mode $`\nu `$ and at a given momentum $`𝐪`$, the radii of the symbols in fig.2 indicate the square modulus of the displacement decomposed in Ca and C in-plane ($`xy`$, parallel to the graphene layer) and out-of-plane ($`z`$, perpendicular to the graphene layer) contributions. The corresponding phonon density of states (PHDOS) are shown in fig. 3 (b) and (c). The decomposed PHDOS are well separated in energy. The graphite modes are weakly dispersing in the out-of-plane direction while the Ca modes are three dimensional. However the Ca<sub>xy</sub> and the Ca<sub>z</sub> vibration are well separated contrary to what expected for a perfect fcc-lattice. One Ca<sub>xy</sub> vibration is an Einstein mode being weakly dispersive in all directions.
The superconducting properties of C<sub>6</sub>Ca can be understood calculating the electron-phonon interaction for a phonon mode $`\nu `$ with momentum $`𝐪`$:
$$\lambda _{𝐪\nu }=\frac{4}{\omega _{𝐪\nu }N(0)N_k}\underset{𝐤,n,m}{}|g_{𝐤n,𝐤+𝐪m}^\nu |^2\delta (ϵ_{𝐤n})\delta (ϵ_{𝐤+𝐪m})$$
(2)
where the sum is over the Brillouin Zone. The matrix element is $`g_{𝐤n,𝐤+𝐪m}^\nu =𝐤n|\delta V/\delta u_{𝐪\nu }|𝐤+𝐪m/\sqrt{2\omega _{𝐪\nu }}`$, where $`u_{𝐪\nu }`$ is the amplitude of the displacement of the phonon and $`V`$ is the Kohn-Sham potential. The electron-phonon coupling is $`\lambda =_{𝐪\nu }\lambda _{𝐪\nu }/N_q=0.83`$. We show in fig.3 (a) the Eliashberg function
$$\alpha ^2F(\omega )=\frac{1}{2N_q}\underset{𝐪\nu }{}\lambda _{𝐪\nu }\omega _{𝐪\nu }\delta (\omega \omega _{𝐪\nu })$$
(3)
and the integral $`\lambda (\omega )=2_{\mathrm{}}^\omega 𝑑\omega ^{}\alpha ^2F(\omega ^{})/\omega ^{}`$. Three main contributions to $`\lambda `$ can be identified associated to Ca<sub>xy</sub>, C<sub>z</sub> and C<sub>xy</sub> vibrations.
A more precise estimate of the different contributions can be obtained noting that
$$\lambda =\frac{1}{N_q}\underset{𝐪}{}\underset{i\alpha j\beta }{}[𝐆_𝐪]_{i\alpha ,j\beta }[𝐂_{𝐪}^{}{}_{}{}^{1}]_{j\beta ,i\alpha }$$
(4)
where $`i,\alpha `$ indexes indicate the displacement in the Cartesian direction $`\alpha `$ of the $`i^{\mathrm{th}}`$ atom, $`[𝐆_𝐪]_{i\alpha ,j\beta }=_{𝐤,n,m}4\stackrel{~}{g}_{i\alpha }^{}\stackrel{~}{g}_{j\beta }\delta (ϵ_{𝐤n})\delta (ϵ_{𝐤+𝐪m})/[N(0)N_k]`$, and $`\stackrel{~}{g}_{i\alpha }=𝐤n|\delta V/\delta x_{𝐪i\alpha }|𝐤+𝐪m/\sqrt{2}`$. The $`𝐂_𝐪`$ matrix is the Fourier transform of the force constant matrix (the derivative of the forces respect to the atomic displacements). We decompose $`\lambda `$ restricting the summation over $`i,\alpha `$ and that over $`i,\beta `$ on two sets of atoms and Cartesian directions. The sets are C<sub>xy</sub>, C<sub>z</sub>, Ca<sub>xy</sub>, and Ca<sub>z</sub>. The resulting $`𝝀`$ matrix is:
$$𝝀=\begin{array}{cc}& \begin{array}{cccc}\mathrm{C}_{xy}& \mathrm{C}_z& \mathrm{Ca}_{xy}& \mathrm{Ca}_z\end{array}\\ \begin{array}{c}\mathrm{C}_{xy}\\ \mathrm{C}_z\\ \mathrm{Ca}_{xy}\\ \mathrm{Ca}_z\end{array}& \left(\begin{array}{cccc}0.12& 0.00& 0.00& 0.00\\ 0.00& 0.33& 0.04& 0.01\\ 0.00& 0.04& 0.27& 0.00\\ 0.00& 0.01& 0.00& 0.06\end{array}\right)\end{array}$$
(5)
The off-diagonal elements are negligible. The Ca out-of-plane and C in-plane contributions are small. For the in-plane C displacements, eq. 1 with $`\mathrm{\Delta }=0.11`$ gives $`\lambda _{\mathrm{C}_{xy},\mathrm{C}_{xy}}=0.11`$. Such a good agreement is probably fortuitous given the oversimplified assumptions of the model. The main contributions to $`\lambda `$ come from Ca in-plane and C out-of-plane displacements. As we noted previously the C out-of-plane vibration do not couple with the C $`\pi `$ Fermi surfaces. Thus the coupling to the C out-of-plane displacements comes from electrons belonging to the Ca Fermi surface. Contrary to what expected in an fcc lattice, the Ca<sub>xy</sub> phonon frequencies are smaller than the Ca<sub>z</sub> ones. This can be explained from the much larger $`\lambda `$ of the Ca in-plane modes.
The critical superconducting temperature is estimated using the McMillan formulamcmillan :
$$T_c=\frac{\omega }{1.2}\mathrm{exp}\left(\frac{1.04(1+\lambda )}{\lambda \mu ^{}(1+0.62\lambda )}\right)$$
(6)
where $`\mu ^{}`$ is the screened Coulomb pseudopotential and $`\omega =24.7`$ meV is the phonon frequencies logarithmic average. We obtain T$`{}_{c}{}^{}=11`$K, with $`\mu ^{}=0.14`$. We calculate the isotope effect by neglecting the dependence of $`\mu ^{}`$ on $`\omega `$. We calculate the parameter $`\alpha (\mathrm{X})=\frac{d\mathrm{log}T_c}{dM_\mathrm{X}}`$ where X is C or Ca. We get $`\alpha (\mathrm{Ca})=0.24`$ and $`\alpha (\mathrm{C})=0.26`$. Our computed $`\alpha (\mathrm{Ca})`$ is substantially smaller than the estimate given in ref. Mazin . This is due to the fact that only $`40\%`$ of $`\lambda `$ comes from the coupling to Ca phonon modes and not $`85\%`$ as stated in ref.Mazin .
In this work we have shown that superconductivity in C<sub>6</sub>Ca is due to an electron-phonon mechanism. The carriers are mostly electrons in the Ca Fermi surface coupled with Ca in-plane and C out-of-plane phonons. Coupling to both modes is important, as can be easily inferred from the calculated isotope exponents $`\alpha (\mathrm{Ca})=0.24`$ and $`\alpha (\mathrm{C})=0.26`$. Our results suggest a general mechanism for the occurrence of superconductivity in GICs. In order to stabilize a superconducting state it is necessary to have an intercalant Fermi surface since the simple doping of the $`\pi `$ bands in graphite does not lead to a sizeable electron-phonon coupling. This condition occurs if the intercalant band is partially occupied, i. e. when the intercalant is not fully ionized. The role played in superconducting GICs by the intercalant Fermi surface has been previously suggested by Jishi . More recently a correlation between the presence of a band, not belonging to graphite, and superconductivity has been observed in Csanyi . However the attribution of this band to an interlayer state not derived from intercalant atomic orbitals is incorrect.
We acknowledge illuminating discussions with M. Lazzeri,G. Loupias, M. d’Astuto, C. Herold and A. Gauzzi. Calculations were performed at the IDRIS supercomputing center (project 051202).
|
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.